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# Virtually Haken surgeries on once-punctured torus bundles
## 1. Introduction
A compact 3-manifold $`M`$ is Haken if it is irreducible, and contains an orientable, essential surface. $`M`$ is virtually Haken if it is finitely covered by a Haken manifold. One of the central problems in 3-dimensional topology is Waldhausenโs Conjecture, which states that every closed, irreducible 3-manifold with infinite fundamental group is virtually Haken.
A knot manifold is an irreducible, orientable, compact 3-manifold whose boundary is a single torus. A knot manifold $`M`$ is small if every closed incompressible surface in $`M`$ is boundary parallel. A knot manifold is hyperbolic if its interior admits a complete hyperbolic metric of finite volume. A slope on a torus $`T`$ is a non-trivial isotopy class of simple closed curves on $`T`$. If $`M`$ is a 3-manifold, $`T`$ is a torus component of $`M`$, and $`\alpha `$ is a slope on $`T`$, then $`M(\alpha )`$ denotes the manifold obtained from $`M`$ by Dehn filling along a simple closed curve representing the slope $`\alpha `$. We say that a knot manifold has Property VH if $`M(\alpha )`$ is virtually Haken for all but finitely many slopes $`\alpha `$ on $`M`$. An important special case of Waldhausenโs conjecture is the following:
###### Conjecture 1.1.
Let $`M`$ be a hyperbolic knot manifold. Then $`M`$ has Property VH.
Examples of small knot manifolds satisfying Property VH have been given in , and .
Let $`F`$ be homeomorphic to a torus with an open disk removed. We choose a basepoint, $`pF`$, for $`\pi _1F`$. Let $`x`$ and $`y`$ be generators of $`\pi _1F`$, as pictured in Figure 1. Let $`D_x`$ and $`D_y`$ represent Dehn twists along simple closed curves in $`intF`$ which are isotopic to loops representing $`x`$ and $`y`$ respectively. Let $`_1^1`$ denote the mapping class group of $`F`$; that is the group of orientation-preserving automorphisms of $`F`$ which restrict to the identity on $`F`$, modulo isotopies which fix every point of $`F`$. It is well-known that $`_1^1`$ is generated by $`D_x`$ and $`D_y`$. There is a natural map $`\varphi :_1^1SL_2()`$; we sometimes use the notation $`\varphi (f)=f_{}`$.
Let $`H_3=\varphi ^1<D_x,D_y^3>`$, and let $`H_4=\varphi ^1<D_x,D_y^4>`$ It is a fact (see Section 7) that $`H_3`$ and $`H_4`$ are both finite-index subgroups of $`_1^1`$.
###### Theorem 1.2.
Let $`M`$ be an orientable, atoroidal 3-manifold which fibers over $`S^1`$, whose fiber, $`F`$, is a compact, orientable surface of genus 1, with a single boundary component. Let $`f:FF`$ be the monodromy, and suppose that the mapping class of $`f`$ lies in the subgroup $`H_i`$, where $`i=3`$ or $`4`$. Then there are slopes $`\beta _i^1,\beta _i^2`$ and an integer $`N`$ such that $`M(\alpha )`$ is virtually Haken whenever $`I(\alpha ,\beta _i^1)>N`$ and $`I(\alpha ,\beta _i^2)>1`$.
Since $`H_3`$ and $`H_4`$ have finite index, it follows that every hyperbolic punctured torus bundle has infinitely many virtually Haken surgeries, a result which was first proved by Baker (see also ). The slopes $`\beta _i^j`$ are computable (see Sections 6-8).
For a given monodromy $`f`$, let $`\beta _i^j=\beta _i^j(f)`$ be slopes as in Theorem 1.2. Let $`๐`$ be the set of mapping classes $`f`$ of $`F`$ such that $`fH_3H_4`$ and $`\{\beta _3^1,\beta _3^2\}\{\beta _4^1,\beta _4^2\}=\mathrm{}`$.
###### Corollary 1.3.
If $`f๐`$, then $`M`$ has Property VH.
We shall show that there are infinitely many commensurability classes of once-punctured torus bundles whose monodromies lie in the class $`๐`$. Thus we obtain:
###### Theorem 1.4.
There are infinitely many pair-wise non-commensurable once-punctured torus bundles which have Property VH.
Remarks:
1. In the case where $`fH_4`$, Theorem 1.2 is a corollary of Theorem 1.3 in . The proof we give here is completely different.
2. Many of the bundles in $`๐`$ have no exceptional surgeries, and so the techniques of and cannot be applied.
It appears that (in some sense) โmostโ monodromies $`f`$ have a power which lies in $`๐`$, and thus most punctured torus bundles are commensurable to ones with Property VH. We have verified this on a computer, for monodromies of low complexity.
To make these statements precise, we introduce a complexity function on monodromies. Recall that every element $`gSL_2()`$ can be written uniquely as a positive word in $`D_x^1`$ and $`D_y`$, times $`\pm Id`$. We define the โcomplexityโ of $`g`$ to be the length of this word. Similarly, if $`f_1^1`$, then we define the complexity of $`f`$ to be the complexity of $`f_{}`$ in $`SL_2()`$, and if $`M`$ is a punctured-torus bundle with monodromy $`f`$, the complexity of $`M`$ is defined to be the minimum of the complexity of $`f`$ among all monodromies of bundles $`N`$ which are bundle equivalent to $`M`$. Using a computer, we can show:
###### Theorem 1.5.
a. Every mapping class of complexity at most 5 has a power which lies in $`๐`$. Thus every once-punctured torus bundle of complexity at most 5 is commensurable to one with Property VH.
b.There are 745 once-punctured torus bundles of complexity at most 12; with at most 36 exceptions, these are all commensurable to bundles with Property VH.
IDEA OF PROOF
Our proof of Theorem 1.2 is inspired by the arguments of Cooper and Walsh (), who show that every every fibered knot in a $`/2`$ homology sphere admits infinitely many virtually Haken surgeries. The idea is to replace $`M`$ with a finite cover $`\stackrel{~}{M}`$ in which the fiber has multiple boundary components. Then one hopes to find a non-separating surface in $`\stackrel{~}{M}`$ which is not a fiber, so that the techniques of may be applied to find an essential surface in a cyclic cover of $`\stackrel{~}{M}`$. Finally one must show that certain slopes on $`M`$ lift to the ultimate cover.
Our argument diverges from that of in that we choose $`\stackrel{~}{M}`$ to have three or four boundary components instead of two. Since we have two covers to work with, we obtain two surfaces, which is the key to proving Property VH. Furthermore, the surfaces which we construct are disjoint from one of the boundary components of $`\stackrel{~}{M}`$, and thus cannot be fibers. Thus we avoid a number of issues in involving semi-bundle structures.
However, we encounter several new issues. First, it is possible that our surfaces may become fibers after Dehn filling, and to rule this out requires the computation of an Alexander polynomial. Secondly, and more importantly, we require that the boundary components of the non-separating surface in $`\stackrel{~}{M}`$ must all project to the same slope in $`M`$. To arrange this, we must develop techniques for constructing surfaces, and computing slopes, explicitly. Some of these techniques (those in Section 3) may be applied to any bundle, but some (those in Section 4) exploit special features of the genus 1 mapping class group.
A QUESTION
Although some of the methods in this paper apply only to punctured torus bundles, it is conceivable Conjecture 1.1 can be attacked along broadly similar lines. A key step would be to answer the following (presumably difficult) question:
Question: Let $`M`$ be a knot manifold. Is there a finite cover $`\stackrel{~}{M}`$ of $`M`$ which contains a non-separating surface, which is disjoint from some component of $`\stackrel{~}{M}`$, and whose boundary curves all project to the same (embedded) slope on $`M`$?
By the results of this paper, the answer to the question is yes if $`M`$ is a punctured torus bundle.
PLAN OF PAPER
Section 2 fixes a choice of basis for $`H_1(M)`$, when $`M`$ is any 3-manifold fibering over $`S^1`$. In Section 3, we show how to compute โalgebraicโ boundary slopes of non-separating surfaces in bundles. In Section 4, we prove the main theorem. In Sections 5-7, we give methods for explicit computations of slopes. Section 8 is devoted to the example of the figure eight knot exterior. In Section 9, we prove Theorem 1.4. Finally, in Section 10, we discuss our computer-generated data.
ACKNOWLEDGMENT
Thanks are due to Genevieve Walsh for pointing out an error in a previous version.
## 2. Framing convention
Let $`f:FF`$ be an automorphism of a compact, orientable surface, and let $`M=F\times [0,1]/(x,0)=(fx,1)`$ be the 3-manifold fibering over $`S^1`$ with monodromy $`f`$. Let $`\lambda _1,\mathrm{},\lambda _k`$ be the boundary components of $`F`$, and suppose that $`f`$ acts trivially on $`F`$, so $`M`$ has torus boundary components $`T_1,\mathrm{},T_k`$.
We wish to fix a framing for $`M`$ (see Figure 1). For $`i=1,\mathrm{},k`$, we fix a point $`p_i\lambda _i`$, and let the meridian, $`\mu _iT_i`$, be the suspension of the point $`p_i`$. The orientation of $`\mu _i`$ is chosen so that the map from $`[0,1]`$ (with standard orientation) to $`\mu _i`$ given by $`t(p_i,1t)`$ is orientation-preserving. We let the longitude, $`\lambda _i`$, of $`T_i`$ be given by $`\lambda _i\times \{1\}`$. We orient $`\lambda _i`$ so that $`I(\mu _i,\lambda _i)=1`$, where $`I(.,.)`$ is the standard intersection pairing on $`H_1(M)`$. (see Figure 1).
Given a surface $`S`$ properly embedded in $`M`$, we may specify the homology classes of the boundary curves of $`S`$ by a vector $`(\alpha _1,\mathrm{},\alpha _k)`$, where each $`\alpha _i`$ is an ordered pair of integers. We shall refer to this as the vector of โalgebraic boundary slopesโ of $`S`$. For example if we say that $`S`$ has algebraic boundary slopes $`((1,2),(0,0),(0,3))`$, we mean that $`[ST_1]=[\mu _1+2\lambda _1]`$ in $`H_1(T_1)`$, that $`[ST_2]=[0]H_1(T_2)`$, and that $`[ST_3]=3[\lambda _3]`$ in $`H_1(T_3)`$.
## 3. Homology and boundary slopes of bundles
Unless otherwise specified, all homology groups in this paper will be taken with $``$ coefficients.
We begin by recalling some well-known facts. If $`M`$ is a manifold which fibers over $`S^1`$, with monodromy $`f:FF`$, then there is a corresponding Wang exact sequence
$$\begin{array}{ccccccccccc}\mathrm{}& & H_j(M)& \stackrel{\theta }{}& H_{j1}(F)& \stackrel{f_{}Id}{}& H_{j1}(F)& \stackrel{i_{}}{}& H_{j1}(M)& & \mathrm{}\end{array}$$
where $`\theta `$ is the map induced by intersection with $`F\times \{0\}`$. There is also a relative version of this sequence, which fits with the exact sequence of a pair into the following commutative diagram:
$`\begin{array}{ccccccccc}\mathrm{}H_j(M,M)& & H_{j1}(F,F)& & H_{j1}(F,F)& & H_{j1}(M,M)& & \mathrm{}\\ & & & & & & & & \\ \mathrm{}H_{j1}(M)& & H_{j2}(F)& & H_{j2}(F)& & H_{j2}(M)& & \mathrm{}\\ & & & & & & & & \\ \mathrm{}H_{j1}(M)& & H_{j2}(F)& & H_{j2}(F)& & H_{j2}(M)& & \mathrm{}\end{array}`$
In the case where $`M`$ is a 3-manifold, one may use the Wang sequence to compute that $`RankH_2(M,M,)=1+Rank(fix(f_{}),)`$, where $`fix(f_{})`$ denotes the set of vectors in $`H_1(F,F)`$ which are fixed by $`f_{}`$.
Suppose now that $`f`$ acts trivially on $`F`$. As described in the previous section, for each component of $`M`$, there is a canonical meridian, given by the suspension of a point on $`F`$, and longitude, given by intersection with $`F`$. This gives a canonical basis for $`H_1(M)`$, and we have corresponding projections $`\mu :H_1(M)H_0(F)`$, which sends each longitude to 0, and $`\lambda :H_1(M)H_1(F)`$, which sends each meridian to 0. We will fix a preferred component $`\mathrm{}`$ of $`F`$, and let $`\pi _{\mathrm{}}`$ be the projection map from $`H_1(F)`$ to $`H_1(\mathrm{})`$.
We may define an injective map $`\eta :H_2(M,M)H_1(F,F)\times `$ by
$$\eta [R]=([RF\times \{0\}],\pi _{\mathrm{}}\lambda ([R])).$$
If $`[\delta ]fix(f_{})H_1(F,F)`$, then $`[f\delta \delta ]=0H_1(F,F)`$, and so the closed loop $`f\delta \delta `$ is homologous to some class $`[x]`$ in $`Image(i_{}:H_1(F)H_1(F)))`$. Note that $`[x]`$ is unique, up to adding copies of $`[F]`$. Thus we may define $`\psi :fix(f_{})\times H_1(F)`$ uniquely, by requiring
$`i_{}\psi ([\delta ],m)`$ $`=`$ $`[f\delta \delta ],\text{ and}`$
$`\pi _{\mathrm{}}\psi ([\delta ],m)`$ $`=`$ $`m[\mathrm{}].`$
Finally, let $`\varphi :H_2(M,M)H_1(M)`$ come from the sequence of the pair.
We have the following diagrams:
$$\begin{array}{ccc}H_2(M,M)& \stackrel{\eta }{}& fix(f_{})\times \\ \varphi & & \psi & & \\ H_1(M)& \stackrel{\lambda }{}& H_1(F)\end{array}$$
and
$$\begin{array}{ccc}H_2(M,M)& & H_1(F,F)\\ & & & & \\ H_1(M)& \stackrel{\mu }{}& H_0(F),\end{array}$$
where in the second diagram, all unlabeled maps come from the exact sequence of the pair, and the Wang sequence.
###### Lemma 3.1.
The above diagrams commute.
###### Proof.
For the second diagram, note that the map $`\mu `$ agrees with the map $`\theta `$ from the Wang sequence fpr $`M`$. Thus the second diagram fits into the larger diagram given at the beginning of the section, which commutes.
For the first diagram, suppose we are given a class $`[R]H_2(M,M)`$. A 2-chain homologous to $`R`$ may be constructed as follows. We let $`\delta =R(F\times \{0\})`$, and let $`\eta [R]=([\delta ],m)`$. By definition of the map $`\psi `$, we have $`[\delta f\delta ]=i_{}\psi ([\delta ],m)`$. Thus there is a map $`g:XF`$, where $`X`$ is an orientable surface with $`X=_0X_1X`$, such that $`g|_{_0X}`$ is the immersed curve $`\delta f\delta `$ and such that $`g_1XF`$, with $`i_{}[g_1X]=i_{}\psi ([\delta ],m)`$; since $`i_{}[F]=0`$, we may also assume that $`g_1X\mathrm{}=\mathrm{}`$. Let $`_{00}X_0X`$ be the union of arcs which map to $`\delta `$, and let $`_{01}X_0X`$ be the union of arcs which map to $`f\delta `$. Let $`\sigma `$ be a properly embedded collection of separating arcs in $`X`$, with $`\sigma =_{00}X_{01}X`$. Let $`h:X\sigma \{0,1\}`$ be a continuous map, such that $`h(_{00}X)=0`$ and $`h(_{01}X)=1`$.
Let $`Y`$ be an orientable surface obtained from $`\overline{X\sigma }(\sigma \times [0,1])`$, by identifying $`_{00}X`$ and $`_{01}X`$ according to the map $`f`$, and identifying $`\sigma \times [0,1]`$ with $`\overline{X\sigma }(X\sigma )`$ in the obvious way.
Then we may construct a map $`j:YM`$, by the rule $`j(x)=(gx,hx)`$, if $`xX\sigma `$, and $`j(x,t)=(gx,t)`$ if $`(x,t)\sigma \times [0,1]`$. By construction, $`\eta ([jY])=([\delta ],0)`$, and so $`\eta ([jY]+m[F])=([\delta ],m)=\eta [R]`$. Thus $`[jY]+m[F]=[R]`$. Also, by inspection, $`\lambda \varphi [jY]=[g_1X]=\psi ([\delta ],0)`$, and so
$`\lambda \varphi [R]`$ $`=`$ $`\lambda \varphi ([jY]+m[F])`$
$`=`$ $`\psi ([\delta ],0)+m[F]`$
$`=`$ $`\psi ([\delta ],m)`$
$`=`$ $`\psi \eta ([jY]+m[F])`$
$`=`$ $`\psi \eta [R].`$
###### Corollary 3.2.
Suppose $`R`$ is a properly embedded, orientable, non-separating surface in $`M`$, and let $`\delta =R(F\times \{0\})`$. Then the algebraic boundary slopes of $`R`$ satisfy:
$`\mu ([R])`$ $`=`$ $`[\delta F]H_0(F),`$
$`\lambda ([R])`$ $`=`$ $`[f\delta \delta ]+k[F]H_1(F)\text{ for some integer }k.`$
The following corollary employs notation introduced in Section 2.
###### Corollary 3.3.
Suppose there is an arc $`\delta `$ properly embedded in $`F`$, with $`[\delta F]=a_i[p_i]H_0(F)`$, and $`[f\delta \delta ]=b_i[\lambda _i]H_1(F)`$. Then there is a non-separating, orientable surface $`R`$ properly embedded in $`M`$, with algebraic boundary slopes $`((a_1,b_1),\mathrm{},(a_k,b_k))`$.
###### Proof.
By the Wang exact sequence, there is a class $`[R]H_2(M,M)`$, such that $`[RF\times \{0\}]=\theta [R]=[\delta ]H_1(F,F)`$. Then, by the previous corollary, we have $`\mu ([R])=[\delta F]H_0(F)`$, and $`\lambda ([R])=[f\delta \delta ]+k[F]H_1(F)`$. Adding a multiple of $`[F]`$ to $`[R]`$, we get $`\lambda [R]=[f\delta \delta ]`$, and the corollary follows. โ
## 4. Proof of Theorem 1.2
We shall make use of the following, which can be proved by straightforward applications of the methods of .
###### Theorem 4.1.
For any compact, orientable surface, $`S`$, there is a positive integer $`n=n(S)`$, depending only on the topological type of $`S`$, such that the following is true. Let $`M`$ be any compact, orientable, irreducible, atoroidal 3-manifold, with two torus boundary components, containing a properly embedded, orientable, incompressible, non-separating surface homeomorphic to $`S`$, which is not a fiber in a fibration of $`M`$, with algebraic boundary slopes $`(\delta _1,\delta _2)`$, where $`\delta _i0`$. Then, if $`|I(\alpha _1,\delta _1)|=|I(\alpha _2,\delta _2)|>n(S)`$, the manifold obtained by Dehn filling $`M`$ along $`\alpha _1`$ and $`\alpha _2`$ is virtually Haken.
Remark: By choosing the function $`n(S)`$ appropriately, we may assume that, whenever $`S^{}`$ is obtained by compressing $`S`$, we have $`n(S^{})<n(S)`$ .
We are now ready to prove the main theorem.
###### Proof.
(Of Theorem 1.2)
It is well-known that, if $`f,g_1^1`$, and $`f_{}=g_{}`$, then $`fg^1`$ is isotopic to a power of a Dehn twist along a peripheral curve in $`F`$. Therefore, if $`f_{}=g_{}`$, then $`M_f`$ is bundle equivalent to $`M_g`$, and so we may reduce to the case where $`fJ_i=<D_y^i,D_x>`$.
Case a: $`fJ_3`$.
In this case, $`f`$ can be written as a word $`W`$ in $`D_x`$ and $`D_y^3`$. Let $`m=m_3`$ denote the exponent sum of $`D_x`$ in $`W`$. Let $`\beta =\beta _3^1`$ be the slope $`(3,m)/gcd(3,m)`$ on $`M`$.
Let $`\theta :\widehat{F}F`$ be the 3-fold cyclic cover of $`F`$ dual to the curve $`x`$ (see Fig. 1). Since $`fH_3`$, then $`f`$ lifts to an automorphism $`\widehat{f}:\widehat{F}\widehat{F}`$ which acts trivially on $`\widehat{F}`$. Let $`\pi :\widehat{M}M`$ be the 3-fold cover of $`M`$ induced by the lift $`\widehat{f}:\widehat{F}\widehat{F}`$.
If $`g`$ is a map between two surfaces with boundary, we let $`g_{\mathrm{}}`$ be the induced map on $`H_1`$ rel. boundary. Let $`\lambda _1,\lambda _2,\lambda _3`$ be the components of $`\widehat{F}`$ (with pre-image orientations induced from $`\lambda `$), and let $`p_i\lambda _i`$ be the pre-images of $`pF`$. Let $`\delta _i`$ be an arc connecting $`\lambda _i`$ and $`\lambda _{i+1}`$, as pictured in Figure 2a. Let $`\delta =\delta _1\delta _2`$. Then it is easy to see that $`[\delta ]`$ is a non-zero class in $`H_1(\widehat{F},\widehat{F})`$ which is fixed by the lifts of $`D_x`$ and $`(D_y)^3`$, and therefore by $`\widehat{f}_{\mathrm{}}`$. Moreover, $`[f\delta \delta ]=m[\lambda _2]H_1(F)`$.
Let $`\tau `$ be the covering transformation of $`\widehat{F}`$ such that $`\tau (\delta _1)=\delta _2`$.
###### Lemma 4.2.
There is a non-separating surface $`RM`$ whose boundary slopes are given by $`((0,0),(3,m),(3,m))`$.
###### Proof.
We have $`[f(\delta \tau \delta )(\delta \tau \delta )]=m[\lambda _3]m[\lambda _2]H_1(F)`$, and $`[(\delta \tau \delta )F]=0[p_1]3[p_2]+3[p_3]H_0(F)`$. The lemma now follows from Corollary 3.3. โ
Let $`T_i`$ be the component of $`\widehat{M}`$ containing $`p_i`$. By attaching annuli to $`R`$, we may assume that $`R`$ is disjoint from $`T_1`$. Let $`\alpha `$ be a slope on $`M`$, let $`\widehat{\alpha }_i`$ be the lift of $`\alpha `$ to $`T_i`$, and let $`S`$ be an incompressible surface in $`\widehat{M}(\widehat{\alpha }_1)`$ obtained by compressing $`R`$. Then $`S`$ has boundary slopes which project to curves of slope $`\beta `$ in $`M`$.
Let $`N`$ be the infinite cyclic cover of $`\widehat{M}(\widehat{\alpha }_1)`$ which is dual to $`S`$, and let $`\mathrm{\Delta }(s)`$ be the Alexander polynomial associated to this infinite cyclic cover.
Let $`\theta :\stackrel{~}{F}\widehat{F}`$ be the infinite cyclic cover dual to $`\delta \tau \delta `$. We let $`\widehat{x}_i=y^{i1}xy^{(i1)}\widehat{F}`$, and let $`\widehat{y}=y^3\widehat{F}`$. There is a $`[s^{\pm 1}]`$-module decomposition $`H_1(\stackrel{~}{F},\theta ^1p_1)H_1(\stackrel{~}{F})+<[\widehat{x}_1]>`$.
The monodromy $`\widehat{f}`$ acts on the module $`H_1(\stackrel{~}{F},\theta ^1p_1)`$ by a matrix $`f_{}`$, and since $`\widehat{f}`$ fixes $`p_1`$, then $`Im(Idf_{})H_1(\stackrel{~}{F})`$. The map $`Idf_{}`$ has a kernel containing $`<[\widehat{x}_1\widehat{x}_2^1]>`$, and so there is an induced map $`\overline{Idf_{}}:H_1(\stackrel{~}{F},\theta ^1p_1)/<[\widehat{x}_1\widehat{x}_2^1]>H_1(\stackrel{~}{F})H_1(\stackrel{~}{F},\theta ^1p_1)`$.
The group $`H_1(\stackrel{~}{F},\theta ^1p_1)`$ is a free $`[s^{\pm 1}]`$-module, with basis $`=([\widehat{y}],[\widehat{x}_1\widehat{x}_2^1],[\widehat{x}_1^2\widehat{x}_3],[\widehat{x}_1])`$. We let $`[f_{}]`$ be the matrix representative for $`f_{}`$ with respect to $``$. Then $`\overline{Idf_{}}`$ can be represented by a matrix $`[\overline{Idf_{}}]`$, which is obtained by deleting the second column and fourth row of $`Id[f_{}]`$.
###### Lemma 4.3.
a. The matrix $`[f_{}]`$ is given by matrix $`W([D_x],[D_y^3])`$, where
$`[D_x]=\left(\begin{array}{cccc}s& 0& 0& 0\\ s^1& 1& ss^1& 0\\ s^1& 0& s^1& 0\\ 0& 0& 0& 1\end{array}\right),`$
and
$`[D_y^3]=\left(\begin{array}{cccc}1& 0& 1+s+s^2& 1\\ 0& 1& 0& 0\\ 0& 0& 1& 0\\ 0& 0& 0& 1\end{array}\right),`$
b. The Alexander polynomial $`\mathrm{\Delta }(s)`$ is the determinant of the following matrix:
$`B=\left(\begin{array}{cccc}0& & & \\ 0& & [\overline{Idf_{}}]^T& \\ 1s& & & \\ p& 0& qs^1& 0\end{array}\right).`$
Using Lemma 4.3 and a computer, the Alexander polynomial $`\mathrm{\Delta }(s)`$ can be easily computed.
###### Proof.
Part a is a computation, which can be done with Fox derivatives. We leave this to the reader.
For part b, let $`\stackrel{~}{M}`$ be the $``$-cover of $`\widehat{M}`$ which is dual to $`R`$. We begin by computing a presentation for $`H_1(\stackrel{~}{M})`$ as a $`[s^{\pm 1}]`$\- module.
We choose the point $`p_1`$ as a basepoint for $`\pi _1\widehat{M}`$. We have $`\pi _1\widehat{F}\pi _1\widehat{M}`$, and we let $`\widehat{t}=t\pi _1\widehat{M}`$. We have the following presentation for $`\pi _1(\widehat{M})`$:
$`\pi _1\widehat{M}=<\widehat{y},\widehat{x}_1,\widehat{x}_1\widehat{x}_2^1,\widehat{x}_1^2\widehat{x}_3,\widehat{t}|R_1,R_2,R_3,R_4>,`$
$`R_1=\widehat{t}\widehat{y}\widehat{t}^1\widehat{y}^1(\widehat{y}\widehat{f}_{}\widehat{y}^1),`$
$`R_2=\widehat{t}\widehat{x}_1\widehat{t}^1\widehat{x}_1^1(\widehat{x}_1\widehat{f}_{}\widehat{x}_1^1)`$
$`R_3=\widehat{t}\widehat{x}_1\widehat{x}_2^1\widehat{t}^1(\widehat{x}_1\widehat{x}_2^1)^1(\widehat{x}_1\widehat{x}_2^1\widehat{f}_{}(\widehat{x}_1\widehat{x}_2^1)^1),`$
$`R_4=\widehat{t}\widehat{x}_1^2\widehat{x}_3\widehat{t}^1(\widehat{x}_1^2\widehat{x}_3)^1(\widehat{x}_1^2\widehat{x}_3\widehat{f}_{}(\widehat{x}_1^2\widehat{x}_3)^1).`$
For any element $`w\pi _1\stackrel{~}{M}\pi _1M`$, it will be convenient to let $`[w]`$ denote the the image of $`w`$ in $`H_1(\stackrel{~}{M})`$, and to let $`s=\widehat{x}_1`$. Then, as a $`[s^{\pm 1}]`$-module, $`H_1(\stackrel{~}{M})`$ has an ordered generating set $`^{}=([\widehat{t}],[\widehat{y}],[\widehat{x}_1\widehat{x}_2^1],[\widehat{x}_1^2\widehat{x}_3])`$.
The relators of $`H_1(\stackrel{~}{M})`$ (as a $`[s^{\pm 1}]`$-module) may be obtained from the relators $`R_1,\mathrm{},R_4`$ of $`\pi _1\stackrel{~}{M}`$. We have:
$`R_1`$ $`=`$ $`\widehat{t}\widehat{y}\widehat{t}^1\widehat{y}^1(\widehat{y}\widehat{f}_{}\widehat{y}^1),`$
$``$ $`[\widehat{y}\widehat{f}_{}\widehat{y}^1]=0`$
$`R_2`$ $`=`$ $`\widehat{t}\widehat{x}_1\widehat{t}^1\widehat{x}_1^1(\widehat{x}_1\widehat{f}_{}\widehat{x}_1^1)`$
$``$ $`(1s)[\widehat{t}]+[\widehat{x}_1\widehat{f}_{}\widehat{x}_1^1]=0`$
$`R_3`$ $`=`$ $`\widehat{t}(\widehat{x}_1\widehat{x}_2^1)\widehat{t}^1(\widehat{x}_1\widehat{x}_2^1)^1`$
$``$ $`0=0`$
$`R_4`$ $`=`$ $`\widehat{t}(\widehat{x}_1^2\widehat{x}_3)\widehat{t}^1(\widehat{x}_1^2\widehat{x}_3)^1((\widehat{x}_1^2\widehat{x}_3)\widehat{f}_{}(\widehat{x}_1^2\widehat{x}_3)^1)`$
$``$ $`[(\widehat{x}_1^2\widehat{x}_3)\widehat{f}_{}(\widehat{x}_1^2\widehat{x}_3)^1]=0`$
Thus a presentation matrix for $`H_1(\stackrel{~}{M})`$, in terms of $`^{}`$, is given by:
$`A=\left(\begin{array}{cccc}0& & & \\ 0& & \overline{Idf_{}}^T& \\ 1s& & & \end{array}\right).`$
The presentation matrix for $`H_1(N)`$ is obtained from $`A`$ by adding a single relator of the form $`[t^p(x^1yxy^1)^q]=0`$, where $`p,q`$. This yields the relator $`p[\widehat{t}]qs^1[x_1x_2^1]=0`$. Thus the presentation matrix for $`H_1(N)`$ is:
$`\left(\begin{array}{cccc}0& & & \\ 0& & \overline{Idf_{}}^T& \\ 1s& & & \\ p& 0& qs^1& 0\end{array}\right),`$
and so $`\mathrm{\Delta }(s)`$ is the determinant of this matrix. โ
###### Corollary 4.4.
There is a slope $`\beta _3^2`$ such that, if $`I(\alpha ,\beta _3^2)>1`$, then $`S`$ is not a fiber in a fibration of $`\widehat{M}(\widehat{\alpha }_1)`$.
###### Proof.
It is well-known that, if an infinite cyclic cover of a compact 3-manifold is dual to a fiber in a fibration, then the corresponding Alexander polynomial is monic. By Lemma 4.3, $`\mathrm{\Delta }(s)=pDetM_1+qDetM_2`$ for some matrices $`M_1`$ and $`M_2`$ with entries in $`[s^{\pm 1}]`$. Thus the leading term of $`\mathrm{\Delta }(s)`$ is $`pn_1+qn_2`$, for some integers $`n_1,n_2`$. We let $`\beta _3^2=(n_2,n_1)/gcd(n_2,n_1)`$. Then if $`|I(\alpha ,\beta _3^2)|1`$, $`\mathrm{\Delta }(s)`$ is non-monic. โ
Let $`n`$ be a positive integer. By Thurstonโs hyperbolic Dehn surgery theorem, we may assume that $`n`$ is chosen large enough so that $`\widehat{M}(\widehat{\alpha }_1)`$ is hyperbolic whenever $`I(\alpha ,\beta _3^1)>n`$. We also assume that $`n`$ is larger than the integer $`n(R)`$ given by Theorem 4.1, and hence also bigger than $`n(S)`$.
Suppose $`|I(\alpha ,\beta _3^1)|=k>n`$, and that $`|I(\alpha ,\beta _3^2)|>1`$. Since $`k>n`$, then $`\widehat{M}(\widehat{\alpha }_1)`$ is hyperbolic, and since $`|I(\alpha ,\beta _3^2)|>1`$ then, by Corollary 4.4, $`S`$ is not a fiber in a fibration of $`\widehat{M}(\widehat{\alpha }_1)`$. Let $`\widehat{\alpha }_2,\widehat{\alpha }_3`$ be the lifts of $`\alpha `$ to the components of $`\widehat{M}(\widehat{\alpha }_1)`$. Then $`|I(\widehat{\alpha }_2,S)|=|I(\widehat{\alpha }_3,S)|=k>nn(S)`$, so by Theorem 4.1, the Dehn filling of $`\widehat{M}(\widehat{\alpha }_1)`$ along $`\widehat{\alpha }_2`$ and $`\widehat{\alpha }_3`$ is virtually Haken. Since this manifold covers $`M(\alpha )`$, then $`M(\alpha )`$ is virtually Haken.
Case b: $`fJ_4`$.
The argument is similar to the argument for Case a. Let $`\beta =\beta _4=(2,m)/gcd(2,m)`$, where $`m=m_4`$ is the exponent sum of $`D_x`$ in the word $`W(D_x,D_y^4)`$ which represents $`f`$. We let $`\widehat{F}`$ be the 4-fold cyclic cover dual to $`x`$, and let $`\widehat{M}`$ be the corresponding cover of $`M`$. We fix a basepoint $`p_1\widehat{M}`$, and let $`p_i`$ be the translate of $`p_1`$ by the covering translation corresponding to $`y^{i1}`$. We let $`T_i`$ be the component of $`\widehat{M}`$ containing $`p_i`$.
We define arcs $`\delta _1`$ and $`\delta _2`$ in $`\widehat{F}`$ as pictured in Figure 3, and let $`\delta =\delta _1\delta _2`$. In this case, the non-separating surface $`R`$ corresponding to $`[\delta ][\tau ^2\delta ]`$ has boundary slopes $`((0,0),(2,m),(0,0),(2,m))((0,0),(0,0),(0,0),(0,0))`$. By compressing $`R`$, we obtain a properly embedded, orientable, non-separating, incompressible surface $`S`$ in $`\widehat{M}(\widehat{\alpha }_1,\widehat{\alpha }_3)`$, whose boundary curves all project to curves of slope $`\beta `$ on $`M`$.
Let $`N`$ be the infinite cyclic cover of $`\widehat{M}(\widehat{\alpha }_1,\widehat{\alpha }_3)`$ dual to $`S`$, and let $`\mathrm{\Delta }(s)`$ be the corresponding Alexander polynomial. Let $`\theta :\stackrel{~}{F}\widehat{F}`$ be the infinite cyclic cover of $`\widehat{F}`$ dual to $`\delta \tau ^2\delta `$. We let $`\widehat{x}_i=y^{i1}xy^{(i1)}\widehat{F}`$, and let $`\widehat{y}=y^4\widehat{F}`$.
The automorphism $`\widehat{f}`$ induces an automorphism $`f_{}`$ of $`H_1(\stackrel{~}{F},\theta ^1p_1p_3)`$. There is an induced map
$$\overline{Idf_{}}:H_1(\stackrel{~}{F},\theta ^1(p_1p_3))/<[\widehat{x}_1\widehat{x}_2^1],[\widehat{x}_3\widehat{x}_4^1]>H_1(\stackrel{~}{F})H_1(\stackrel{~}{F},\theta ^1(p_1p_3)).$$
The $`[s^{\pm 1}]`$-module $`H_1(\stackrel{~}{F},\theta ^1(p_1p_3))`$ is free, with basis
$$=([\widehat{y}^4],[\widehat{x}_1\widehat{x}_3],[\widehat{x}_1],[\widehat{y}^2],[\widehat{x}_1\widehat{x}_2^1],[\widehat{x}_3\widehat{x}_4^1]).$$
We let $`[f_{}]`$ be the matrix representing $`f_{}`$, in terms of the basis $``$. Then $`\overline{Idf_{}}`$ is represented by a matrix $`[\overline{Idf_{}}]`$ obtained from $`Id[f_{}]`$ by deleting the 5th and 6th columns, and the 3rd and 4th rows.
###### Lemma 4.5.
a. The matrix $`[f_{}]`$ is given by $`W([D_x],[D_y^4]`$, where
$`[D_x]=\left(\begin{array}{cccccc}s& 0& 0& 0& 0& 0\\ (1+s^1)& s^1& 0& s^1& 0& 0\\ 0& 0& 1& 0& 0& 0\\ 0& 0& 0& 1& 0& 0\\ s^1& 1s^1& 0& s^1& 1& 0\\ 0& 0& 0& 0& 0& 1\end{array}\right),`$
and
$`[D_y^4]=\left(\begin{array}{cccccc}1& 1+s& 1& 0& 0& 0\\ 0& 1& 0& 0& 0& 0\\ 0& 0& 1& 0& 0& 0\\ 0& 0& 0& 1& 0& 0\\ 0& 0& 0& 0& 1& 0\\ 0& 0& 0& 0& 0& 1\end{array}\right),`$
b. The Alexander polynomial $`\mathrm{\Delta }(s)`$ is the determinant of the following matrix:
$`B=\left(\begin{array}{ccccc}0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0\\ 1s& 0& 0& 0& 0\\ p& 0& 0& 0& qs\\ p& 0& 0& qs^1& 0\end{array}\right)+\left(\begin{array}{ccccc}0& & & & \\ 0& & \overline{Idf_{}}^T& & \\ 0& & & & \\ 0& & & & \\ 0& 0& 0& 0& 0\end{array}\right)\left(\begin{array}{ccccc}1& 0& 0& 0& 0\\ 0& 1& 0& 0& 0\\ 0& 0& 1& 0& 0\\ 0& 0& 0& p& 0\\ 0& 0& 0& 0& 1\end{array}\right),`$
and $`\mathrm{\Delta }(s)`$ is divisible by $`q`$.
###### Proof.
The proof of part a is an elementary application of Fox calculus.
For part b, let $`\stackrel{~}{M}`$ be the infinite cyclic cover of $`\stackrel{~}{M}`$ dual to $`S`$. We choose an ordered generating set $`^{}=(\widehat{t},\widehat{y}^4,\widehat{x}_1\widehat{x}_3,\widehat{x}_1\widehat{x}_2^1,\widehat{x}_3\widehat{x}_4^1)`$ for the $`[s^{\pm 1}]`$-module $`H_1(\stackrel{~}{M})`$.
The relators for $`\pi _1\widehat{M}`$ give the following relations for $`H_1(\stackrel{~}{M})`$:
$`R_1`$ $`=`$ $`\widehat{t}\widehat{y}\widehat{t}^1\widehat{y}^1(\widehat{y}\widehat{f}_{}\widehat{y}^1),`$
$``$ $`[\widehat{y}\widehat{f}_{}\widehat{y}^1]=0`$
$`R_2`$ $`=`$ $`\widehat{t}\widehat{x}_1\widehat{x}_2^1\widehat{t}^1(\widehat{x}_1\widehat{x}_2^1)^1(\widehat{x}_1\widehat{x}_2^1\widehat{f}_{}(\widehat{x}_1\widehat{x}_2^1)^1)`$
$``$ $`0=0`$
$`R_3`$ $`=`$ $`\widehat{t}\widehat{x}_1\widehat{x}_3\widehat{t}^1(\widehat{x}_1\widehat{x}_3)(\widehat{x}_1\widehat{x}_3\widehat{f}_{}(\widehat{x}_1\widehat{x}_3)^1)`$
$``$ $`[\widehat{x}_1\widehat{x}_3\widehat{f}_{}(\widehat{x}_1\widehat{x}_3)^1]=0`$
$`R_4`$ $`=`$ $`\widehat{t}\widehat{x}_1\widehat{t}^1\widehat{x}_1^1(\widehat{x}_1\widehat{f}_{}\widehat{x}_1^1)`$
$``$ $`(1s)[\widehat{t}]+[\widehat{x}_1\widehat{f}_{}\widehat{x}_1^1]=0`$
$`R_5`$ $`=`$ $`\widehat{t}\widehat{x}_3\widehat{x}_4^1\widehat{t}^1(\widehat{x}_3\widehat{x}_4^1)(\widehat{x}_3\widehat{x}_4^1\widehat{f}_{}(\widehat{x}_3\widehat{x}_4^1)^1)`$
$``$ $`0=0`$
The module $`H_1(N)`$ has two additional relators, $`R_6:[t^p(x^1yxy^1)^q]=0`$, and $`R_7:[y^2t^py^{}2y^2(x^1yxy^1)^qy^2]=0`$. The relator $`R_6`$ can be written as: $`p[\widehat{t}]qs^1[\widehat{x}_1\widehat{x}_2^1]=0`$. For $`R_7`$, we have:
$`0`$ $`=`$ $`p[y^2ty^2]+q[y^2(x^1yxy^1)y^2]`$
$`=`$ $`p[(y^2ty^2t^1)t]qs[x_1x_2^1]`$
$`=`$ $`p[\widehat{t}]+p[y^2f_{}y^2]qs[x_1x_2^1]`$
$`=`$ $`p[\widehat{t}]+p[\overline{Idf_{}}y^2]qs[x_1x_2^1].`$
Thus the presentation matrix for $`H_1(N)`$, with respect to $``$, is:
$`\left(\begin{array}{ccccc}0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0\\ 1s& 0& 0& 0& 0\\ p& 0& 0& 0& qs\\ p& 0& 0& qs^1& 0\end{array}\right)+\left(\begin{array}{ccccc}0& & & & \\ 0& & \overline{Idf_{}}& & \\ 0& & & & \\ 0& & & & \\ 0& 0& 0& 0& 0\end{array}\right)\left(\begin{array}{ccccc}1& 0& 0& 0& 0\\ 0& 1& 0& 0& 0\\ 0& 0& 1& 0& 0\\ 0& 0& 0& p& 0\\ 0& 0& 0& 0& 1\end{array}\right)`$
We then observe that the last column of $`\overline{Idf_{}}`$ is all 0โs. Therfore the last column of the presentation matrix for $`H_1(N)`$ has a single non-zero entry $`qs`$. Thus $`\mathrm{\Delta }(s)`$ is divisible by $`q`$. โ
###### Corollary 4.6.
Let $`\beta _4^2=(1,0)`$. If $`I(\alpha ,\beta _4^2)>1`$, then $`S`$ is not a fiber in a fibration of $`\widehat{M}(\widehat{\alpha }_1,\widehat{\alpha }_3)`$.
###### Proof.
By Lemma 4.5, $`\mathrm{\Delta }(s)`$ is divisible by $`q`$. If $`|q|=|I(\alpha ,(1,0))|>1`$, then $`\mathrm{\Delta }(s)`$ is non-monic. โ
Now an application of Theorem 4.1 shows that $`\widehat{M}(\widehat{\alpha }_1,\widehat{\alpha }_2,\widehat{\alpha }_3,\widehat{\alpha }_4)`$ has a Haken cyclic cover, provided that $`|I(\alpha ,\beta _4^2)|>1`$ and $`|I(\alpha ,\beta _4^1)|`$ is large enough. โ
## 5. Computations: framings
In the proof of Theorem 1.2, we used the fact that, if $`f,g_1^1`$, and $`f_{}=g_{}SL_2()`$, then $`M_f`$ is bundle-equivalent to $`M_g`$. In this section, we shall show how to compute the effect of this equivalence on the framings.
Let $`D_\lambda `$ be a Dehn twist about a peripheral curve in $`F`$. Suppose $`f_1^1`$, and that $`f_{}=Id`$. Then $`f`$ is equivalent in $`_1^1`$ to $`D_\lambda ^n`$ for some integer $`n`$. We define the twist of $`f`$ by the formula $`t(f)=n`$.
###### Lemma 5.1.
Suppose that, $`f,g_1^1`$, with $`f_{}=g_{}`$. Then there is a homeomorphism $`h:M_fM_g`$, such that, with respect to the standard framings on $`M_f`$ and $`M_g`$, $`h(1,0)=(1,t(fg^1))`$ and $`h(0,1)=(0,1)`$.
###### Proof.
Since $`fg^1=D_\lambda ^{t(fg^1)}`$, then $`fg^1`$ is isotopic to the identity, by an isotopy which twists $`t(fg^1)`$ times around the boundary of $`F`$. Using this isotopy, one may construct a bundle-equivalence between $`M_f`$ and $`M_g`$, and verify that the effect on the framings is as claimed. โ
Thus, given elements $`f,g_1^1`$, in terms of $`D_x`$ and $`D_y`$, with $`f_{}=g_{}`$, we require a method for computing the twist $`t(fg^1)`$.
###### Lemma 5.2.
Let $`f,g_1^1`$ be given as words, $`W_f,W_g`$ in $`D_x`$ and $`D_y`$, and suppose that $`f_{}=g_{}`$. Then $`t(fg^1)`$ is equal to 1/12 of the total sum of the exponents of $`D_x`$ and $`D_y`$ in the word $`fg^1=W_f(D_x,D_y)W_g^1(D_x,D_y)`$.
###### Proof.
We use the following well-known presentation for $`SL_2()`$:
$$SL_2()<a,b,\tau |\tau =(ab)^3=(aba)^2,\tau ^2=[\tau ,a]=[\tau ,b]=id>,$$
where the map sending $`D_x`$ to $`a`$ and $`D_y`$ to $`b`$ is an isomorphism. Thus if $`f_{}=g_{}`$, then $`fg^1=W(D_x,D_y)`$ is a product of conjugates of the elements $`R_1=(D_xD_y)^3(D_xD_yD_x)^2`$, $`R_2=(D_xD_y)^3D_x(D_xD_y)^3D_x^1`$, $`R_3=(D_xD_y)^3D_y(D_xD_y)^3D_y^1`$, and $`R_4=(D_xD_y)^6`$. By computing the effect of these automorphisms on $`\pi _1F`$, one may check directly that each one is trivial in $`_1^1`$ except for $`R_4`$, and that $`t(R_4)=1`$. Therefore $`t(fg^1)`$ is equal to the (signed) number of conjugates of $`R_4`$ in $`W_f(D_x,D_y)W_g^1(D_x,D_y)`$. Since the sum of the exponents is 12 on $`R_4`$, and zero on each of $`R_1,R_2`$ and $`R_3`$, then we see that $`t(fg^1)`$ is simply 1/12 of the the total sum of the exponents of $`W_fW_g^1`$. โ
## 6. Computations: subgroups of $`SL_2()`$
The following lemma is known. It can be proved directly, using standard combinatorial group theory algorithms (as implemented for example on the program GAP). We shall give a proof based on the Euclidean algorithm.
###### Lemma 6.1.
The subgroups $`H_3`$ and $`H_4`$ have finite index in $`_1^1`$.
###### Proof.
We shall prove the statement for $`H_4`$, the proof for $`H_3`$ being entirely analogous.
Given an ordered pair of relatively prime, non-zero integers $`(m,n)`$, we may generate a sequence $`(m_0,n_0),\mathrm{},(m_k,n_k)`$ recursively, as follows:
Let $`(m_0,n_0)=(m,n)`$. Suppose $`(m_i,n_i)`$ has been defined. If $`|n_i\pm m_i|<|n_i|`$, then let $`(m_{i+1},n_{i+1})=(m_i,n_i\pm m_i)`$; if $`|n_i+m_i||n_i|`$ and $`|n_im_i||n_i|`$ and $`|m_i\pm 4n_i|<|m_i|`$, then let $`(m_{i+1},n_{i+1})=(m_i\pm 4n_i,n_i)`$; if neither of these conditions holds, then terminate the sequence at $`(m_i,n_i)`$.
Claim: For any pair of relatively prime, non-zero integers $`(m,n)`$, the above rule defines a finite sequence terminating in $`(\pm 1,0)`$ or $`(0,\pm 1)`$.
Proof of claim: Suppose that $`(m_i,n_i)`$ has been defined, that neither $`m_i`$ nor $`n_i`$ is zero, and that $`|n_i+m_i|`$ and $`|n_im_i|`$ are both as big as $`|n_i|`$. Then $`|m_i|2|n_i|`$. Since $`m`$ and $`n`$ are relatively prime, it follows that $`m_i`$ and $`n_i`$ are relatively prime, and thus we have strict inequality $`|m_i|>2|n_i|`$, and so $`|m_i\pm 4n_i|<|m_i|`$. So the sequence continues until we reach $`(j,0)`$ or $`(0,j)`$. In this case, both $`m`$ and $`n`$ are divisible by $`j`$, so we have $`j=\pm 1`$. This proves the claim.
Now, $`SL_2()`$ acts on the hyperbolic plane, and there is an induced action on the circle at infinity, which is identified with $`\{\mathrm{}\}`$.
Identifying the rational number $`m/n`$ with the vector $`\left(\begin{array}{c}m\\ n\end{array}\right)`$, and $`\mathrm{}`$ with $`\left(\begin{array}{c}1\\ 0\end{array}\right)`$, the action of $`<D_x,D_y^4>`$ on $`\mathrm{}`$ is given by:
$`\left(\begin{array}{cc}1& \pm 1\\ 0& 1\end{array}\right)\left(\begin{array}{c}m\\ n\end{array}\right)=\left(\begin{array}{c}m\pm n\\ n\end{array}\right)`$
and $`\left(\begin{array}{cc}1& 0\\ \pm 4& 1\end{array}\right)\left(\begin{array}{c}m\\ n\end{array}\right)=\left(\begin{array}{c}m\\ \pm 4m+n\end{array}\right).`$ Therefore, by the claim, the orbit $`<D_x,D_y^4>\{0,\mathrm{}\}`$ is dense in $`\{\mathrm{}\}`$. Thus the domain of discontinuity of $`<D_x,D_y^4>`$ is empty, and so $`<D_x,D_y^4>`$ is a finite-index subgroup of $`SL_2()`$. Thus $`H_4`$ has finite index in $`_1^1`$. โ
In fact, it can be checked (for example on GAP), that $`H_3`$ has index 8, and $`H_4`$ has index 12.
Computation of exponent sums:
Given an element $`g_1^1`$, let $`m_i`$ be the smallest positive integer such that $`g^{m_i}H_i`$. Thus there is word $`W`$ (not necessarily unique) such that $`g_{}^{m_i}=W(D_x,D_y^i)`$. Let $`n_i`$ be the exponent sum of $`D_x`$ in $`W`$. The numbers $`m_i`$ and $`n_i`$ can be computed on GAP, by using Reidemeister-Schreier style algorithms for subgroup presentations.
The computation of $`m_i`$ is quite straightforward. The computation of $`W_i`$ (and hence $`n_i`$) requires a slightly more complicated, but standard, procedure. The idea is to have GAP compute a presentation for the subgroup of $`SL_2()`$ generated by $`D_x`$, $`D_y^i`$ and $`g_{}^{n_i}`$. After simplifying, GAP finds that the generator $`g_{}^{n_i}`$ is redundant, and returns the word $`W`$. Details can be found in the source code at www.math.buffalo.edu/ jdmaster.
## 7. Computations: slopes
Given an arbitrary $`f_1^1`$, we may associate slopes $`\beta _3^1,\beta _3^2`$ and $`\beta _4^1,\beta _4^2`$ for the boundary of a cyclic cover $`M_{f^m}`$, as follows.
To compute the slopes $`\beta _i^1`$, We first compute an integer $`m`$ such that $`f^mH_3H_4`$. We then compute words $`W_3`$ and $`W_4`$ such that $`f_{}^m=W_i(D_x,D_y^i)`$, as described in the previous section, and let $`n_i`$ be the exponent sum of $`D_x`$ in $`W_i`$. We then let $`g_i=W_i(D_x,D_y^i)_1^1`$. Since $`g_iJ_i`$, then by the proof of Theorem 1.2, we see that associated to the bundle $`M_{g_i}`$ are slopes $`\beta _3^1=(3,n_3)/gcd(3,n_3)`$, $`\beta _4^1=(2,n_4)/gcd(2,n_4)`$, and $`\beta _4^2=(1,0)`$.
We use Lemma 5.2 to compute $`t(f^mg_i^1)`$, and then use Lemma 5.1 to compute that
$`\beta _3^1`$ $`=`$ $`(3,n_33(t(f^mg_3^1)))/gcd(3,n_3),`$
$`\beta _4^1`$ $`=`$ $`(2,n_42(t(f^mg_4^1)))/gcd(2,n_4),`$
$`\beta _4^2`$ $`=`$ $`(1,t(f^mg_4^1)).`$
To compute $`\beta _3^2`$, we first use Lemmas 4.3 and 4.5 to compute the Alexander polynomial for the relevant cover of the manifold $`M_{g_3}`$. Then, as in the proofs of Corollaries 4.4 and 4.6, we obtain the slope $`\beta _3^2`$. Finally, using Lemma 5.1, we compute the slope $`\beta _3^2`$ for the manifold $`M_{f^m}`$.
If $`\{\beta _3^1,\beta _3^2\}\{\beta _4^1,\beta _4^2\}=\mathrm{}`$, then we report โsuccessโ, meaning that $`M_f`$ has a finite cover satisfying property VH.
## 8. Example
Let $`f=D_x^1D_y`$, so $`M_f`$ is the figure-eight knot exterior. We shall show that $`f^{12}`$ has property VH. Using GAP, we compute: $`m=12`$, so $`f^{12}H_3H_4`$; also $`W_3(a,b)=(a^1baba^1bab)^3`$ and $`W_4(a,b)=(a^2b^1a^1b^1a^1)^4`$, and so $`n_3=0`$ and $`n_4=16`$. We have $`f_{}^{12}=W_3(D_x,D_y^3)=W_4(D_x,D_y^4)`$. Letting $`g_i=W_i(D_x,D_y^i)`$, then the slopes associated to $`M_{g_i}`$ are $`\beta _3^1=(1,0)`$, $`\beta _4^1=(1,8)`$, and $`\beta _4^2=(1,0)`$.
The Alexander polynomial, $`\mathrm{\Delta }_3(s)`$, for the relevant cover of $`M_{g_3}`$ is given by:
$`\mathrm{\Delta }_3(s)`$ $`=`$ $`qs^4+3qs^3+2qs^2+2qsq+qs^12qs^22qs^33qs^4qs^5.`$
This polynomial is non-monic whenever $`|q|>1`$. Thus, associated to $`M_{g_3}`$ is the slope $`\beta _3^2=(1,0)`$. We use Lemma 5.2 to compute that $`t(f^{12}g_3^1)=3`$ and $`t(f^{12}g_4^1)=4`$.
Then
$`\beta _3^1`$ $`=`$ $`(3,n_33(t(f^mg_3^1)))/gcd(3,n_3),`$
$`=`$ $`(1,3)`$
$`\beta _3^2`$ $`=`$ $`(1,t(f^mg_3^1))`$
$`=`$ $`(1,3)`$
$`\beta _4^1`$ $`=`$ $`(2,n_42(t(f^mg_4^1)))/gcd(2,n_4),`$
$`=`$ $`(1,12)`$
$`\beta _4^2`$ $`=`$ $`(1,t(f^mg_4^1))`$
$`=`$ $`(1,4)`$
Since $`\{\beta _3^1,\beta _3^2\}\{\beta _4^1,\beta _4^2\}=\mathrm{}`$, then $`M_{f^{12}}`$ has property VH.
## 9. An infinite family with Property VH
###### Proof.
(of Theorem 1.4)
Let $`f_n=(D_x^1D_y)^{12}D_y^{12n}`$. It may easily be checked that the induced map on $`H_1(F)`$ has trace bigger than 3, and so the corresponding bundle $`M_{f_n}`$ is atoroidal. For $`f_0`$, a computation (see previous section) gives $`\beta _3^1=(1,3),\beta _3^2=(1,3),\beta _4^1=(1,12)`$ and $`\beta _4^2=(1,4)`$.
The element $`f_n`$ is equivalent in $`SL_2()`$ to the element
$$g_n=(D_x^1D_y^3D_xD_y^3D_x^1D_y^3D_xD_y^3)^3D_y^{12n}.$$
One checks that the Alexander polynomial for the relevant cover of $`M_{g_n}`$ is equivalent mod $`n`$ to the Alexander polynomial for the relevant cover of the bundle with monodromy $`D_x^1D_y^3D_xD_y^3D_x^1D_y^3D_xD_y^3`$. One computes that that the leading coefficient of the latter polynomial is $`q`$, and thus the leading coefficient of the former polynomial is divisible by $`q`$. Thus for the manifold $`M_{g_n}`$ we have slopes $`\beta _3^1=\beta _3^2=(1,0)`$. One computes $`t(g_nf_n^1)=3`$, and so for $`f_n`$, we have $`\beta _3^1=\beta _3^2=(1,3)`$.
Also, $`f_n`$ is equivalent in $`SL_2()`$ to the element $`h_n=(D_x^2D_y^4D_x^1D_y^4D_x^1)^4D_y^{12n}`$. For the manifold $`M_{h_n}`$ we have $`\beta _4^2=(1,0)`$, and we compute $`\beta _4^1=(1,8)`$. We compute $`t(h_nf_n^1)=4`$, and so for $`f_n`$, we have $`\beta _4^1=(1,12)`$, and $`\beta _4^2=(1,4)`$.
Since $`\{\beta _3^1,\beta _3^2\}\{\beta _4^1,\beta _4^2\}=\mathrm{}`$, then for all $`n>0`$, the manifold $`M_{f_n}`$ is finitely covered by a bundle with property VH, by Theorem 1.2. Furthermore, the manifolds $`M_{f_n}`$ are all obtained by doing surgery on the same hyperbolic knot $`KM_{f_0}`$. Therefore, by Section 3, there are infinitely many non-commensurable manifolds in the family $`\{M_{f_n}\}`$. โ
## 10. Computer results
For every monodromy of complexity at most 5, the computer verified that a โsuccessโ criterion was met, and so the associated bundle is commensurable to one with property VH. The data for monodromies of complexity at most 5 is given below. Since we are only interested in bundles up to commensurability, we have left out monodromies which are proper powers and monodromies with negative trace. Also, we have only included words up to cyclic permutations. The โnโ in column two is the smallest positive integer such that $`f^nH_3H_4`$.
monodromy f n $`t(f^ng_3^1)`$ $`[\beta _3^1,\beta _3^2]`$ $`t(f^ng_4^1)`$ $`[\beta _4^1,\beta _4^2]`$ $`M_{f^n}`$ has VH? $`D_x^1D_y`$ 12 -3 \[(1,3),(1,3)\] 4 \[(1, -12),(1,-4)\] Yes $`D_x^2D_y`$ 12 2 \[(1,-6), (1,-2)\] 3 \[(1, -15),(1,-3)\] Yes $`D_x^1D_y^2`$ 12 -2 \[(1,6),(1,2)\] -3 \[(1,3),(1,3)\] Yes $`D_x^3D_y`$ 6 2 \[(1,-6), (1,-2)\] 1 \[(1, -5),(1,-1)\] Yes $`D_x^2D_y^2`$ 4 -1 \[(1, 1),(1,1)\] -2 \[(1, 2),(1,2)\] Yes $`D_x^1D_y^3`$ 6 0 \[(1, -2),(1,0)\] -1 \[(1, 5),(1,1)\] Yes $`D_x^4D_y`$ 4 1 \[(1, -5 ), (1,-1)\] -2 \[(1, 2),(1,2)\] Yes $`D_x^3D_y^2`$ 12 0 \[(1, -4),(1,0)\] -9 \[(1, 9),(1,9)\] Yes $`D_x^2D_y^3`$ 4 0 \[( 3, -8),(1,0)\] -1 \[(1, 1),(1,1)\] Yes $`D_x^1D_y^4`$ 4 1 \[(1, -3),(1,-3)\] 0 \[(1, -2),(1,0)\] Yes $`D_x^2D_yD_x^1D_y`$ 4 0 \[(3, -10),(1,-1)\] 0 \[(1, 0),(1,0)\] Yes $`D_x^1D_yD_x^1D_y^2`$ 12 -6 \[(1, 10),(1,6)\] 0 \[(1, 0),(1,0)\] Yes
From the table we see that every bundle of complexity at most five is commensurable to a bundle with Property VH.
We considered all cyclically reduced, primitive, positive words on $`D_x^1`$ and $`D_y`$ of length at most 12, representing hyperbolic monodromies. There are 745 of these. We verified that all but 36 of the associated bundles are finitely covered by a bundle with Property VH. Here we have not distinguished conjugate classes in $`SL_2()`$, so these words do not all correspond to distinct bundles. The monodromy of smallest complexity which we cannot handle is $`D_x^3D_y^3`$, for which $`\beta _3^1=\beta _4^1=(1,3)`$.
The routine runs quickly on words of rather large size. For example, for the monodromy $`f=D_x^{11}D_y^3D_xD_y^6D_x^4(D_yD_x^1)^4D_y`$, a few secondsโ computation gives $`f^4H_3H_4`$, with associated slopes $`\beta _3^1=(1,17),\beta _3^2=(1,15)`$, and $`\beta _4^1=(1,18),\beta _4^2=(1,18)`$. Thus $`M_{f^4}`$ has Property VH.
Mathematics Department
SUNY at Buffalo
Buffalo, NY 14260-2900
jdmaster@buffalo.edu
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# Thomas rotation and Thomas precession
## 1 Introduction
The โparadoxicโ phenomenon of Thomas precession has given rise to much discussion ever since the publication of Thomasโ seminal paper (Thomas 1927) in which he made a correction by a factor 1/2 to the angular velocity of the spin of an electron moving in a magnetic field. Let us mention here that in the literature there seems to be no standard agreement as to the usage of the terms โThomas precessionโ and โThomas rotationโ. As explained in more detail in Section 10 below, we prefer to use the term Thomas precession to refer to the continuous change of direction, with respect to an inertial frame, of a gyroscopic vector moving along a world line. Thomas rotation, on the other hand, will refer to the spatial rotation experienced by a gyroscopic vector having moved along a โclosedโ world line, and having returned to its initial frame of reference (see Section 9).
One of the most studied cases (see e.g. Costella et al. 2001, Kennedy 2002) is the fact that the application of three successive Lorentz boosts (with the relative velocities adding up to zero) results, in general, in a spatial rotation: the discrete Thomas rotation (see Section 4 for details). The same fact is often described as โthe composition of two Lorentz boosts is equivalent to a boost and a spatial rotationโ. We prefer to use three Lorentz boosts instead (with the relative velocities adding up to zero), in order to return to the initial frame of reference, in accordance with our terminology of Thomas rotation. Describing the mathematical structure of discrete Thomas rotations has motivated A.A. Ungar to build the comprehensive theory of gyrogroups and gyrovector spaces (Ungar 2001).
The other case typically under consideration comes from the original observation of Thomas: the continuous change of direction, with respect to an inertial frame, of a gyroscopic vector moving along a circular orbit. This phenomenon has been subject to considerations from various points of view (Muller 1992 (Appendix), Philpott 1996, Rebilas 2002 (Appendix), Herrera & Di Prisco 2002, Rhodes & Shemon 2003). The considerations usually involve, either explicitly or implicitly, the viewpoint of the orbiting โairplaneโ, i.e. a rotating observer. This might lead us to believe (see Herrera & Di Prisco 2002) that the calculated angle of rotation depends on the definition of the rotating observer (and this could lead to an experimental checking of what the โrightโ definition of a rotating observer is). From our treatment below, however, it will be clear the Thomas rotation is an absolute fact, independent of the rotating (or, any other) observer.
It is also interesting to note that new connections between quantum mechanical phenomena and Thomas rotation have recently been pointed out (Lรฉvay 2004).
As it is well known, the theory of special relativity contradicts our common sense notions about space and time in many respects. Early day โparadoxesโ were usually based on our intuitive assumption of absolute simultaneity. With the resolution of paradoxes such as the โtwin paradoxโ or the โtunnel paradoxโ it has become common knowledge that the concept of time must be handled very carefully. As it is also well known, the theory of special relativity implies, besides the non-existence of absolute time, the non-existence of absolute space. An expression such as โa point in spaceโ simply does not have an absolute meaning, just as the expression โan instant in timeโ. However, this fact seems to be given less attention to and even overlooked sometimes. The fact that the space vectors of any observer are usually represented as vectors in $`^3`$ leads one to forget that these spaces really are different. This conceptual error lead e.g. to the โvelocity addition paradoxโ (Mocanu 1992). The spaces of two different inertial observers are, of course, connected via the corresponding Lorentz boost, and the non-transitivity of Lorentz boosts (which, in fact, gives rise to the notion of Thomas rotation) gave the correct explanation of this โparadoxโ (Ungar 1989, Matolcsi & Goher 2001).
To grasp the essence of the concepts related to Thomas rotation, let us mention that in some sense this intriguing phenomenon is analogous to the well known twin paradox. Consider two twins in an inertial frame. One of them remains in that frame for all times, while the other goes for a trip in spacetime, and later returns to his brother. It is well-known that different times have passed for the two twins: the traveller is younger than his brother. What may be surprising is that the space of the traveller when he arrives, although he experienced no torque during his journey, will be rotated compared to the space of his brother; this is, in fact, the Thomas rotation. This analogy is illuminating in one more respect: until the traveller returns to the original frame of reference it makes no sense to ask โhow much younger is the traveller compared to his brother?โ and โby what angle is the travellerโs gyroscope rotated compared to that of his brother?โ Different observers may give different answers. When the traveller returns to his brother, these questions suddenly make perfect sense, and there is an absolute answer (independent of who the observer is) as to how much younger and how much rotated the traveller is.
Of course, an arbitrtary inertial frame can observe the brothers continuously, and can tell, at each of the frameโs instants, what difference he sees between the ages of the brothers. More explicitly, as it is well known, given a world line, an arbitrary inertial frame can tell the relation between the frameโs time and the proper time of the world line. This relation depends on the inertial frame: different inertial frames establish different relations.
Similarly, an arbitrary inertial frame, observing the two brothers, can tell at each frame-instant what difference he sees between the directions of the gyroscopes of the brothers. Different inertial frames establish different relations.
This philsophy makes a clear distinction between Thomas rotation and Thomas precession connected to a world line:
โ Thomas rotation refers to an absolute fact (independent of who observes it), which makes sense only for two equal local rest frames (if such exist) of the world line,
โ Thomas precession refers to a relative fact (i.e. depending on who observes the motion), which makes sense with respect to an arbitrary inertial frame.
In this paper we use the formalism of (Matolcsi 1993) to give a concise and rigorous treatment of the discrete and circular-path Thomas rotations. The Thomas rotation as well as the Thomas precession (with respect to certain inertial observers) along a circular world line are calculated. Our basic concept here is that special relativistic spacetime has a four-dimensional affine structure, and coordinatization (relative to some observer) is, in many cases, unnecesary in the description of physical phenomena. In fact, coordinates can sometimes lead to ambiguities in concepts and definitions, and bear the danger of leading us to overlook the fact that absolute space does not exist.
As well as providing a clear overview of the appearing concepts, the coordinate-free formulation of special relativity enables us to give simple calculations. The indispensable Fermi-Walker equation is also straightforward to derive in our formalism.
## 2 Fundamental notions
In this section some notions and results of the special relativistic spacetime model as a mathematical structure (Matolcsi 1993, 1998, 2001) will be recapitulated. As the formalism slightly differs from the usual textbook treatments of special relativity (but only the formalism: our treatment is mathematically equivalent to the usual treatments), we will point out several relations between textbook formulae and those of our formalism.
Special relativistic spacetime is an oriented four dimensional affine space $`M`$ over the vector space $`๐`$; the spacetime distances form an oriented one dimensional vector space $`๐`$, and an arrow oriented Lorentz form $`๐\times ๐๐๐`$, $`(๐ฑ,๐ฒ)๐ฑ๐ฒ`$ is given.
An absolute velocity $`๐ฎ`$ is a future directed element of $`\frac{๐}{๐}`$ for which $`๐ฎ๐ฎ=1`$ holds (absolute velocity corresponds to four-velocity in usual terminology).
For an absolute velocity $`๐ฎ`$, we define the three dimensional spacelike linear subspace
$$๐_๐ฎ:=\{๐ฑ๐๐ฎ๐ฑ=0\};$$
(1)
then
$$๐
_๐ฎ:=1+๐ฎ๐ฎ:๐๐_๐ฎ,๐ฑ๐ฑ+๐ฎ(๐ฎ๐ฑ)$$
(2)
is the projection onto $`๐_๐ฎ`$ along $`๐ฎ`$. The restriction of the Lorentz form onto $`๐_๐ฎ`$ is positive definite, so $`๐_๐ฎ`$ is a Euclidean vector space (this will correspond to the space vectors of an inertial observer with velocity $`๐ฎ`$).
The history of a classical material point is described by a differentiable world line function $`r:๐M`$ such that $`\dot{r}(๐ฌ)`$ is an absolute velocity for all proper time values $`๐ฌ`$. The range of a world line function โ a one dimensional submanifold โ is called a world line.
An observer $`๐`$ is an absolute velocity valued smooth map defined in a connected open subset of $`M`$. (This is just a mathematical definition; it may sound unfamiliar at first, but considering that something that an observer calls a โfixed space-pointโ is, in fact, a world line in spacetime, this definition will make perfect โphysicalโ sense). A maximal integral curve of $`๐`$ โ a world line โ is a space point of the observer, briefly a $`๐`$-space point; the set of the maximal integral curves of $`๐`$ is the space of the observer, briefly the $`๐`$-space.
An observer having constant value is called inertial. An inertial observer will be referred to by its constant velocity. The space points โ the integral curves โ of an inertial observer with absolute velocity $`๐ฎ`$ are straight lines parallel to $`๐ฎ`$. The $`๐ฎ`$-space point containing the world point $`x`$ is the straight line $`x+\mathrm{๐ฎ๐}`$, where $`\mathrm{๐ฎ๐}:=\{\mathrm{๐ฎ๐ญ}๐ญ๐\}`$.
In order to arrive at the analogue of the coordinate system corresponding to an inertial observer we need to specify the time-syncronization of the observer. Of course, the standard syncroniztion is used: according to the standard synchronization of $`๐ฎ`$, two world points $`x`$ and $`y`$ are simultaneous if and only if $`๐ฎ(yx)=0`$. Thus, simultaneous world points form a hyperplane parallel to $`๐_๐ฎ`$; such a hyperplane is an $`๐ฎ`$-instant, their set is $`๐ฎ`$-time. The $`๐ฎ`$-instant containing the world point $`x`$ is the hyperplane $`x+๐_๐ฎ`$.
An inertial observer together with its standard synchronization is called a standard inertial frame. Note that a standard inertial frame is an exactly defined object in our framework, it does not refer to any coordinates, coordinate axes, it contains an inertial observer and its standard synchronization only.
The space vector between two $`๐ฎ`$-space points (straight lines in spacetime) is the world vector between $`๐ฎ`$-simultaneous world points of the straight lines in question; in formula, $`๐ฎ`$-space, endowed with the subtraction
$$(x+\mathrm{๐ฎ๐})(y+\mathrm{๐ฎ๐}):=๐
_๐ฎ(xy)$$
(3)
becomes a three dimensional affine space over $`๐_๐ฎ`$ (this fact shows that $`๐_๐ฎ`$ does indeed correspond to the space vectors of the observer $`๐ฎ`$).
This is a crucial point: the space vectors of the standard inertial frame $`๐ฎ`$ are elements of $`๐_๐ฎ`$, so the space vectors of different inertial frames form different three dimensional vector spaces.
The time passed between to $`๐ฎ`$-instants (hyperplanes in spacetime) is the time passed between them in an arbitrary $`๐ฎ`$-space point. In formula, $`๐ฎ`$-time, endowed with the subtraction
$$(x+๐_๐ฎ)(y+๐_๐ฎ):=๐ฎ(xy)$$
(4)
becomes a one dimensional affine space over $`๐`$.
If $`r`$ is a world line function, then the standard inertial frame with velocity value $`\dot{r}(๐ฌ)`$ is called the local rest frame corresponding to $`r`$ at $`๐ฌ`$.
In usual treatments the coordinates distinguish a certain inertial frame (the โrestโ frame) and any other inertial frame is considered through its relative velocity with respect to the rest frame (and the coordinates with respect to the new frame are given via the corresponding Lorentz transformation). The main feature of our approach is the systematic use of absolute velocities for characterizing standard inertial frames (this perfectly reflects the principle of relativity: no inertial frame can be distinguished compared to other inertial frames). Among several advantages, such as clarity of many concepts appearing in the theory of relativity, it often results in highly simplified and clear formulae.
## 3 Relative velocity and relative acceleration
Let $`r`$ be a world line function $`r`$ (describing the history of a classical material point). A standard inertial frame with absolute velocity $`๐ฎ`$ gives a correspondence between $`๐ฎ`$-time $`t`$ and the proper time $`๐ฌ`$ of the world line function $`r`$: if $`t_0`$ is the $`๐ฎ`$-instant of the world point $`r(0)`$, then, according to (4), $`๐ญ:=(r(๐ฌ)+๐_๐ฎ)(r(0)+๐_๐ฎ)=๐ฎ(r(๐ฌ)r(0))`$; therefore
$$\frac{d๐ญ}{d๐ฌ}=๐ฎ\dot{r}(๐ฌ).$$
(5)
As a consequence, the proper time, too, can be given as a function of $`๐ฎ`$-time, and
$$\frac{d๐ฌ}{d๐ญ}=\frac{1}{๐ฎ\dot{r}(๐ฌ(๐ญ))}.$$
(6)
The inertial frame observes the history of the material point as a motion, assigning $`๐ฎ`$-space points to $`๐ฎ`$-instants : $`r_๐ฎ(๐ญ):=r(๐ฌ(๐ญ))+\mathrm{๐ฎ๐}`$. Then, according to (3) and the previous equality, the relative velocity is (for the sake of brevity we omit the variable $`๐ญ`$ from the expressions)
$$๐ฏ_๐ฎ:=r_๐ฎ^{}=\underset{๐ก0}{lim}\frac{r_๐ฎ(๐ญ+๐ก)r_๐ฎ(๐ญ)}{๐ก}=\frac{\dot{r}(๐ฌ)}{๐ฎ\dot{r}(๐ฌ)}๐ฎ$$
(7)
and the relative acceleration is
$$๐_๐ฎ:=r_๐ฎ^{\prime \prime }=\frac{1}{(๐ฎ\dot{r}(๐ฌ))^2}\left(\ddot{r}(๐ฌ)+\frac{\dot{r}(๐ฌ)(๐ฎ\ddot{r}(๐ฌ)}{๐ฎ\dot{r}(๐ฌ)}\right)$$
(8)
where the derivative according to $`๐ฎ`$-time is denoted by a prime.
It is worth mentioning that
$$๐ฎ\dot{r}(๐ฌ)=\frac{1}{\sqrt{1|๐ฏ_๐ฎ|^2}}=:\gamma _๐ฎ,$$
(9)
the well-known relativistic factor.
## 4 Lorentz boosts and discrete Thomas rotations
As we emphasized, the space vectors of different standard inertial frames form different three dimensional vector spaces; for the absolute velocities $`๐ฎ`$ and $`๐ฎ^{}`$, $`๐_๐ฎ`$ and $`๐_๐ฎ^{}`$ are different vector spaces. A natural correspondence can be given between them, the Lorentz boost from $`๐ฎ`$ to $`๐ฎ^{}`$ (Matolcsi 1993, 2001),
$$๐(๐ฎ^{},๐ฎ):=1+\frac{(๐ฎ^{}+๐ฎ)(๐ฎ^{}+๐ฎ)}{1๐ฎ^{}๐ฎ}2๐ฎ^{}๐ฎ$$
(10)
which is a Lorentz form preserving linear map on $`๐`$, such that $`B(๐ฎ^{},๐ฎ)๐ฎ=๐ฎ^{}`$. This is the absolute form (which appears implicitly in Rowe 1984, too) of the usual Lorentz boost. It is clear from the given formula that this absolute form depends on two absolute velocities. The explicit matrix form of a textbook Lorentz boost depends on a single relative velocity but, in fact, it also refers to two inertial observers (one of which is the โrest frameโ, not appearing explicitly in the formulae).
The vector $`๐ช^{}`$ in the space of the inertial frame $`๐ฎ^{}`$ is called physically equal to the vector $`๐ช`$ in the space of the inertial frame $`๐ฎ`$ if $`๐ช^{}=๐(๐ฎ^{},๐ฎ)๐ช`$; we say also that $`๐ช`$ boosted from $`๐_๐ฎ`$ to $`๐_๐ฎ^{}`$ equals $`๐ช^{}`$. This Lorentz boost gives sense to the usual tacit assumption that the corresponding coordinate axes of different inertial frames are parallel. The coordinate axes defined by the vectors $`๐_i`$ in $`๐_๐ฎ`$ are parallel to the axes defined by the vectors $`๐_i^{}`$ in $`๐_๐ฎ^{}`$ if $`๐_i^{}=๐(๐ฎ^{},๐ฎ)๐_i`$ $`(i=1,2,3)`$. (The parallelism of frame axes is usually a nagging problem in standard treatments; see the discussion in the Introduction of Kennedy 2002.)
To be physically equal is a symmetric relation: $`๐(๐ฎ^{},๐ฎ)^1=๐(๐ฎ,๐ฎ^{})`$, so if $`๐ช^{}`$ is physically equal to $`๐ช`$, then $`๐ช`$ is physically equal to $`๐ช^{}`$.
On the other hand, to be physically equal is not transitive: the product of two Lorentz boosts, in general, is not a Lorentz boost (as it is well known): we have
$$๐(๐ฎ^{\prime \prime },๐ฎ^{})๐(๐ฎ^{},๐ฎ)=๐(๐ฎ^{\prime \prime },๐ฎ)\text{iff}๐ฎ,๐ฎ^{},๐ฎ^{\prime \prime }\text{are coplanar},$$
(11)
(which is equivalent to the standard formalism: the relative velocity of $`๐ฎ^{\prime \prime }`$ with respect to $`๐ฎ`$ and the relative velocity of $`๐ฎ^{}`$ with respect to $`๐ฎ`$ are collinear.)
In an equivalent formulation,
$$๐_๐ฎ(๐ฎ^{},๐ฎ^{\prime \prime }):=๐(๐ฎ,๐ฎ^{\prime \prime })๐(๐ฎ^{\prime \prime },๐ฎ^{})๐(๐ฎ^{},๐ฎ)$$
(12)
is the identity transformation if and only if $`๐ฎ,๐ฎ^{},๐ฎ^{\prime \prime }`$ are coplanar. Note that $`๐_๐ฎ(๐ฎ^{},๐ฎ^{\prime \prime })๐ฎ=๐ฎ`$ and the restriction of $`๐_๐ฎ(๐ฎ^{},๐ฎ^{\prime \prime })`$ onto $`๐_๐ฎ`$ is a rotation, called the discrete Thomas rotation corresponding to $`๐ฎ`$, $`๐ฎ^{}`$ and $`๐ฎ^{\prime \prime }`$.
Thus if $`๐ช^{}`$ is physically equal to $`๐ช`$ and $`๐ช^{\prime \prime }`$ is physically equal to $`๐ช^{}`$, then $`๐ช`$ need not be physicllay equal to $`๐ช^{\prime \prime }`$. This is why the Thomas rotation appears to be โparadoxicโ.
In other words, a vector $`๐ช`$ boosted from $`๐_๐ฎ`$ to $`๐_๐ฎ^{}`$ yields $`๐ช^{}`$ and then $`๐ช^{}`$ boosted from $`๐_๐ฎ^{}`$ to $`๐_{๐ฎ^{\prime \prime }}`$ yields $`๐ช^{\prime \prime }`$, and lastly $`๐ช^{\prime \prime }`$ boosted from $`๐_{๐ฎ^{\prime \prime }}`$ back to $`๐_๐ฎ`$, results in a vector rotated from the original $`๐ช`$.
## 5 Compasses
A boost, as defined above, does not mean a real transport of vectors from an observer space into another one. Nevertheless, it can be related to such a transport in the following situation.
A compass (a needle fixed to a central point) can be described in spacetime as a vector attached to a material point; more precisely, as a pair of functions $`(r,๐ณ)`$ where $`r`$ is a world line function (the history of the material point) and $`๐ณ`$ is a vector valued function (describing the direction of the needle) defined on the proper time of $`r`$, $`๐ณ:๐๐`$, such that
โ it is always spacelike according to the corresponding local rest frame of the world line, i.e. $`\dot{r}๐ณ=0`$,
โ the magnitude of $`๐ณ`$, $`|๐ณ|`$ is constant.
Thus the needle of the compass passes continuously from the space of one local rest frame to that of another one. The compass is conceived to be locally inertial if $`๐ณ`$ is physically constant along $`r`$ (keeps direction in itself) i.e. the values of $`๐ณ`$ are boosted continuously corresponding to the absolute velocities of the world line. This means that if $`๐ก`$ is a โsmallโ time period, then $`๐ณ(๐ฌ+๐ก)`$ in $`๐_{\dot{r}(๐ฌ+๐ก)}`$ is โnearlyโ physically equal to $`๐ณ(๐ฌ)`$ in $`๐_{\dot{r}(๐ฌ)}`$, more precisely
$$\underset{๐ก0}{lim}\frac{๐ณ(๐ฌ+๐ก)๐(\dot{r}(๐ฌ+๐ก),\dot{r}(๐ฌ))๐ณ(๐ฌ)}{๐ก}=0.$$
(13)
Because $`\dot{r}๐ณ=0`$, we can replace $`\left(\dot{r}(๐ฌ+๐ก)+\dot{r}(๐ฌ)\right)๐ณ(๐ฌ)`$ with $`\left(\dot{r}(๐ฌ+๐ก)\dot{r}(๐ฌ)\right)๐ณ(๐ฌ)`$, so
$$๐(\dot{r}(๐ฌ+๐ก),\dot{r}(๐ฌ))๐ณ(๐ฌ)=๐ณ(๐ฌ)+\frac{\left(\dot{r}(๐ฌ+๐ก)+\dot{r}(๐ฌ)\right)\left(\dot{r}(๐ฌ+๐ก)\dot{r}(๐ฌ)\right)๐ณ(๐ฌ)}{1\dot{r}(๐ฌ+๐ก)\dot{r}(๐ฌ)}$$
(14)
and the above limit becomes $`\dot{๐ณ}\dot{r}(\ddot{r}๐ณ)=0`$, from which, taking into account again $`\dot{r}๐ณ=0`$, we get the well known Fermi-Walker equation along $`r`$
$$\dot{๐ณ}=\dot{r}(\ddot{r}๐ณ)\ddot{r}(\dot{r}๐ณ)=(\dot{r}\ddot{r})๐ณ.$$
(15)
Note that the Lorentz boosts in terms of absolute velocities yielded this equation in an extremely brief and simple way (in contrast to the usual deductions, see e.g. Mรธller 1972).
If $`๐ณ`$ is any vector satisfying the Fermi-Walker equation along $`r`$, then $`(\dot{r}๐ณ)\dot{}=0`$, so $`\dot{r}๐ณ`$ is constant; if $`๐ณ(๐ฌ_0)`$ is spacelike according to $`\dot{r}(๐ฌ_0)`$ for one proper time value $`๐ฌ_0`$, then $`๐ณ(๐ฌ)`$ is spacelike according to $`\dot{r}(๐ฌ)`$ for all $`๐ฌ`$ ($`๐ณ`$ is always spacelike according to the corresponding local rest frame of $`r`$). Moreover, then $`\dot{๐ณ}๐ณ=0`$, so the magnitude of $`๐ณ`$ is constant.
Let us introduce another term. Let $`r`$ be world line function. We call a function $`๐ณ:๐๐`$ a gyroscopic vector on $`r`$ if $`๐ณ`$ satisfies the Fermi-Walker equation along $`r`$ and a value of $`๐ณ`$ is spacelike according to the corresponding local rest frame of $`r`$. Obviously, if $`๐ณ`$ is a gyroscopic vector along $`r`$, then $`(r,๐ณ)`$ is a locally inertial compass. It is well known and easily verifiable that if $`๐ณ_1`$ and $`๐ณ_2`$ are gyroscopic vectors on the same world line, then $`๐ณ_1๐ณ_2`$ is constant (which corresponds to the fact that โnon-rotatingโ vectors retain their relative angle).
## 6 Circular world line
Take a standard inertial frame with velocity value $`๐ฎ_c`$. A circular motion with respect to this frame can be given by
โ its centre $`q_c`$ in $`๐ฎ_c`$-space,
โ its angular velocity, an antisymmetric linear map $`0๐:๐_{๐ฎ_c}\frac{๐_{๐ฎ_c}}{๐}`$ (usually one considers angular velocity as a spatial axial vector which, in fact, corresponds to an antisymmetric tensor),
โ its initial position with respect to the centre, a vector $`0๐ช`$ in $`๐_{๐ฎ_c}`$, orthogonal to the kernel of $`๐`$ such that $`|๐๐ช|<1.`$
This motion has the form
$$๐ญq_c+e^{๐ญ๐}๐ช=q_c+๐ช\mathrm{cos}\omega ๐ญ+\frac{๐๐ช}{\omega }\mathrm{sin}\omega ๐ญ$$
(16)
where $`\omega :=|๐|=\sqrt{\frac{1}{2}\mathrm{Tr}๐^{}๐}`$. Note that we have
$$๐^2๐ช=\omega ^2๐ช,|๐๐ช|=\omega \rho $$
(17)
where $`\rho :=|๐ช|`$.
The relative velocity of this motion equals $`e^{๐ญ๐}๐๐ช`$ which has the magnitude $`\omega \rho `$. Thus, we infer from (5) and (9), that the relation between the proper time $`๐ฌ`$ of the world line and the $`๐ฎ_c`$-time $`๐ญ`$ is $`๐ญ=๐ฌ\lambda `$, where
$$\lambda :=\frac{1}{\sqrt{1\omega ^2\rho ^2}}.$$
(18)
Then we easily derive that this motion comes from the world line function
$$๐ฌr(๐ฌ)=o+๐ฌ\lambda ๐ฎ_c+e^{๐ฌ\lambda ๐}๐ช$$
(19)
where $`o`$ is a world point of the centre $`q_c`$ (which is a straight line in spacetime). Then
$$\dot{r}(๐ฌ)=\lambda (๐ฎ_c+e^{๐ฌ\lambda ๐}๐๐ช),\ddot{r}(๐ฌ)=\lambda ^2\omega ^2e^{๐ฌ\lambda ๐}๐ช.$$
(20)
Note that $`๐ฎ_c`$ is the absolute velocity of the centre and $`๐ฎ_0:=\lambda (๐ฎ_c+๐๐ช)`$ is the โinitialโ absolute velocity of the world line.
## 7 Gyroscopic vectors on a circular world line
Introducing the variable $`๐ญ:=\lambda ๐ฌ`$ ($`๐ฎ_c`$-time) and the function $`\widehat{๐ณ}(๐ญ):=๐ณ(๐ญ/\lambda )`$, then omitting the โhatโ for brevity, we get the Fermi-Walker diferential equation (15) along the above circular world line in the form
$$๐ณ^{}(๐ญ)=\lambda ^2\omega ^2\left((๐ฎ_c+e^{๐ญ๐}๐๐ช)(e^{๐ญ๐}๐ช)\right)๐ณ(๐ญ).$$
(21)
In the sequel we find it convenient to consider $`๐`$ as defined on the whole of $`๐`$ in such a way that $`๐๐ฎ_c=0`$. Then $`๐`$ will be a Lorentz antisymmetric linear map on the whole of $`๐`$, thus $`e^{๐ญ๐}`$ will preserve the Lorentz form (it will be a Lorentz transformation) for which $`e^{๐ญ๐}๐ฎ_c=๐ฎ_c`$ holds.
Then we infer that $`๐(๐ญ):=e^{๐ญ๐}๐ณ(๐ญ)`$ satisfies the autonomouos linear differential equation
$$๐^{}(๐ญ)=\mathrm{๐๐}(๐ญ)$$
(22)
where
$$๐:=๐+\lambda ^2\omega ^2(๐ฎ_c+๐๐ช)๐ช=\lambda ^2๐+\lambda ^2\omega ^2๐ฎ_c๐ช$$
(23)
where the latter equality relies on the simple fact that
$$(๐๐ช)๐ช=\rho ^2๐.$$
(24)
As a consequence โ since $`๐(0)=๐ณ(0)`$ โ, we get the solution of the Fermi-Walker differential equation in the form
$$๐ณ(๐ญ)=e^{๐ญ๐}e^{\mathrm{๐ญ๐}}๐ณ(0).$$
(25)
Let us investigate the properties of
$$๐
(๐ญ):=e^{๐ญ๐}e^{\mathrm{๐ญ๐}}$$
(26)
which we call the Fermi-Walker operator at $`๐ญ=\lambda ๐ฌ`$, $`๐ฌ`$ being a proper time point of the circular world line function.
Since $`๐`$ is an antisymmetric linear map, $`e^{\mathrm{๐ญ๐}}`$ is a Lorentz transformation. It is trivial that $`\mathrm{๐๐ฎ}_0=0`$, thus the restriction of $`e^{\mathrm{๐ญ๐}}`$ onto the three dimensional Euclidean space $`๐_{๐ฎ_0}`$ is a rotation.
We know that the restriction of the Lorentz transformation $`e^{๐ญ๐}`$ onto the Euclidean vector space $`๐_{๐ฎ_c}`$ is a rotation.
Thus $`e^{๐ญ๐}e^{\mathrm{๐ญ๐}}`$, as a product of two Lorentz transformations, is a Lorentz transformation, too. Its restriction onto $`๐_{๐ฎ_0}`$ is a Euclidean structure preserving linear bijection from $`๐_{๐ฎ_0}`$ onto $`๐_{\dot{r}(๐ญ)}`$. This can be conceived as a spatial rotation only if $`\dot{r}(๐ญ)=๐ฎ_0`$ (otherwise it acts between different Euclidean spaces).
## 8 Thomas rotation on the circular world line
The absolute velocity of the circular world line is periodic, $`\dot{r}(\frac{2\pi }{\omega })=\dot{r}(0)=๐ฎ_0`$. Since $`e^{\frac{2\pi }{\omega }๐}`$ is the identity map, we have for the corrresponding Fermi-Walker operator
$$๐
\left(\frac{2\pi }{\omega }\right)=e^{\frac{2\pi }{\omega }๐}$$
(27)
whose restriction onto the Euclidean vector space $`๐_{๐ฎ_0}`$ is a rotation, called the Thomas rotation on the circular world line (19).
The angle of the Thomas rotation is $`2\pi \frac{2\pi }{\omega }|๐|`$ where $`|๐|`$ is the magnitude of $`๐`$; $`|๐|:=\sqrt{|\mathrm{๐๐}_1|^2+|\mathrm{๐๐}_2|^2}`$ where $`๐_1`$ and $`๐_2`$ are arbitrary $`๐ฎ_0`$-spacelike unit vectors orthogonal to the kernel of $`๐`$ such that $`๐_1๐_2=0`$.
It is trivial from (23) that if $`๐๐_{๐ฎ_c}`$ is in the kernel of $`๐`$ โ i.e. $`๐๐=0`$ and $`๐ช๐=0`$ โ, then $`๐`$ is in the kernel of $`๐`$, too. Therefore, the intersection $`E_{๐ฎ_0}E_{๐ฎ_c}\mathrm{Ker}๐\mathrm{Ker}๐`$ is 1-dimensional.
This means that we can choose $`๐_1:=\frac{๐ช}{|๐ช|}`$ and $`๐_2:=\lambda \left(\omega \rho ๐ฎ_c+\frac{๐๐ช}{\omega \rho }\right)`$ (it is easy to verify that all conditions imposed on $`๐_1`$ and $`๐_2`$ are satisfied). Thus, $`\mathrm{๐๐}_1=\lambda \omega ๐_2`$, $`\mathrm{๐๐}_2=\lambda \omega ๐_1`$, which implies that $`|๐|=\lambda \omega `$.
As a consequence, the Thomas angle on the circular world line equals
$$2\pi \left(1\frac{1}{\sqrt{1\omega ^2\rho ^2}}\right)$$
(28)
which is the well known result (Thomas 1927).
It is worth noting that the value of a gyroscopic vector after a whole revolution equals the original one if and only if the gyroscopic vector is parallel to the kernel of $`๐`$ i.e. is orthogonal to the plane of rotation in the space of the centre.
## 9 Generalizations
Besides deriving the Thomas angle on the circular world line in a short and transparent way, our method gives the Thomas rotation itself and allows us a deeper insight into the nature of gyroscopic vectors in general.
Let $`r`$ be an arbitrary world line function. The solutions of the corresponding Fermi-Walker equation with various initial values give us a Fermi-Walker operator $`๐
(๐ฌ_2,๐ฌ_1)`$, a Lorentz transformation for all proper time points $`๐ฌ_1`$ and $`๐ฌ_2`$ such that
$$\dot{r}(๐ฌ_2)=๐
(๐ฌ_2,๐ฌ_1)\dot{r}(๐ฌ_1)$$
(29)
and
$$๐ณ(๐ฌ_2)=๐
(๐ฌ_2,๐ฌ_1)๐ณ(๐ฌ_1)$$
(30)
for an arbitrary gyroscopic vector $`๐ณ`$ on $`r`$.
Thus the restriction of $`๐
(๐ฌ_2,๐ฌ_1)`$ onto $`๐_{\dot{r}(๐ฌ_1)}`$ โ the space vectors of the local rest frame at $`๐ฌ_1`$ โ is a Euclidean structure preserving linear bijection onto $`๐_{\dot{r}(๐ฌ_2)}`$ โ the space vectors of the local rest frame at $`๐ฌ_2`$.
In particular, if $`\dot{r}(๐ฌ_2)=\dot{r}(๐ฌ_1)`$, the restriction of $`๐
(๐ฌ_2,๐ฌ_1)`$ onto $`๐_{\dot{r}(๐ฌ_1)}`$ is a rotation, which we call the Thomas rotation on the world line $`r`$, corresponding to the proper time points $`๐ฌ_1`$ and $`๐ฌ_2`$.
It is worth noting: a Thomas rotation on a world line for two proper time values has a meaning only if the corresponding absolute velocites are equal. Thus no Thomas rotation can be defined on a world line if all its absolute velocity values are different.
## 10 Thomas precession with respect to an inertial frame
Now, let $`๐ณ`$ be a gyroscopic vector on the world line function $`r`$. An inertial frame $`๐ฎ`$ observes $`๐ณ`$ by boosting it continuously to its own space, i.e. giving the function $`๐ณ_๐ฎ:๐๐_๐ฎ`$ such that
$$๐ณ_๐ฎ(๐ญ):=๐(๐ฎ,\dot{r}(๐ฌ(๐ญ)))๐ณ(๐ฌ(๐ญ)).$$
(31)
Then, omitting $`๐ญ`$ as previously, we infer that
$$\begin{array}{cc}\hfill ๐ณ_๐ฎ^{}=\frac{1}{๐ฎ\dot{r}(๐ฌ)}\left(\right(& \frac{d}{d๐ฌ}๐(๐ฎ,\dot{r}(๐ฌ)))๐ณ(๐ฌ)+๐(๐ฎ,\dot{r}(๐ฌ))\dot{๐ณ}(๐ฌ))\hfill \\ \hfill =\frac{1}{๐ฎ\dot{r}(๐ฌ)}\left(\right(& \frac{d}{d๐ฌ}๐(๐ฎ,\dot{r}(๐ฌ)))๐(\dot{r}(๐ฌ),๐ฎ)๐ณ_๐ฎ+\hfill \\ & ๐(๐ฎ,\dot{r}(๐ฌ))(\dot{r}(๐ฌ)\ddot{r}(๐ฌ))๐(\dot{r}(๐ฌ),๐ฎ)๐ณ_๐ฎ))\hfill \end{array}$$
(32)
Omitting $`๐ฌ`$ for the sake of brevity, we get immediately that the second term above equals
$$๐ฎ\left(\ddot{r}+\frac{\dot{r}(๐ฎ\ddot{r})}{1๐ฎ\dot{r}}\right).$$
(33)
As concerns the first term, a straightforward calculation yields that it equals
$$\frac{\dot{r}\ddot{r}}{1๐ฎ\dot{r}}๐ฎ\ddot{r}2๐ฎ\left(\frac{\dot{r}(๐ฎ\ddot{r})}{1๐ฎ\dot{r}}\right)$$
(34)
Taking into account (7) and (8), finally we obtain the known result
$$๐ณ_๐ฎ^{}=\frac{\gamma _๐ฎ^2}{1+\gamma _๐ฎ}(๐ฏ_๐ฎ๐_๐ฎ)๐ณ_๐ฎ.$$
(35)
Thus the inertial frame $`๐ฎ`$ sees the gyroscopic vector $`๐ณ`$ โ which keeps direction in itself โ precessing, the angular velocity of precession is the antisymmetric linear map (depending on $`๐ฎ`$-time)
$$\mathrm{\Omega }_๐ฎ:=\frac{\gamma _๐ฎ^2}{1+\gamma _๐ฎ}๐ฏ_๐ฎ๐_๐ฎ=\frac{\gamma _๐ฎ1}{|๐ฏ_๐ฎ|^2}๐ฏ_๐ฎ๐_๐ฎ:๐_๐ฎ\frac{๐_๐ฎ}{๐}.$$
(36)
Call attention to the fact: the same gyroscopic vector precesses differently to different inertial frames.
## 11 Thomas precessions corresponding to a circular world line
Let us consider the circular world line described in Section 6).
Let us take the standard inertial frame of the centre i.e. the one with absolute velocity $`๐ฎ_c`$. Then equalities in (20), (7) and (8) yield
$$๐ฏ_{๐ฎ_c}(๐ญ)=e^{๐ญ๐}๐๐ช,๐_{๐ฎ_c}(๐ญ)=\omega ^2e^{๐ญ๐}๐ช.$$
(37)
Then $`๐ฏ_{๐ฎ_c}(๐ญ)๐_{๐ฎ_c}(๐ญ)=\omega ^2e^{๐ญ๐}((๐๐ช)๐ช)e^{๐ญ๐}=\omega ^2\rho ^2๐`$ because of (24). Since $`\omega ^2\rho ^2=|๐ฏ_{๐ฎ_c}|^2`$, the angular velocity of the Thomas precession with respect to the โcentral frameโ $`๐ฎ_c`$ is constant in $`๐ฎ_c`$-time, equalling
$$\left(1\frac{1}{\sqrt{1\omega ^2\rho ^2}}\right)๐.$$
(38)
Usual treatments consider exclusively this precession (Mรธller, โฆโฆ) in connection with the circular world line i.e. the Thomas precession with respect to the central frame. Of course, there are other possibilities, too.
For instance, let us take the standard inertial frame in which the gyroscopic vector is at rest initially i.e. the one with absolute velocity $`๐ฎ_0=\lambda (๐ฎ_c+๐๐ช)`$.
Then
$$๐ฎ_0\dot{r}(๐ฌ)=\lambda ^2(1\omega ^2\rho ^2\mathrm{cos}\omega \lambda ๐ฌ).$$
(39)
Consequently, now the $`๐ฎ_0`$-time $`๐ญ`$ and the proper time $`๐ฌ`$ have the relation $`๐ญ=\lambda ^2๐ฌ\lambda \omega ^2\rho ^2\mathrm{sin}\omega \lambda ๐ฌ`$. Then in view of (7), we find
$$๐ฏ_{๐ฎ_0}=\lambda \frac{(๐ฎ_c+e^{\lambda ๐ฌ๐}๐๐ช)(1\omega ^2\rho ^2)}{1\omega ^2\rho ^2\mathrm{cos}\omega \lambda ๐ฌ}\lambda (๐ฎ_c+๐๐ช)$$
(40)
and a similar, more complicated formula gives $`๐_{๐ฎ_0}`$, too; as a consequence, the angular velocity of the Thomas precession with respect to the inertial frame $`๐ฎ_0`$ depends rather intricately on $`๐ฎ_0`$-time. For instance, if $`n`$ is an arbitrary natural number, then
โ for $`๐ฎ_0`$-instants given by $`\lambda ๐ฌ=\frac{2n\pi }{\omega }`$, the value of the relative velocity is zero, so the angular velocity of Thomas precession has zero value, too;
โ for $`๐ฎ_0`$-instants given by $`\lambda ๐ฌ=\frac{(2n1)\pi }{\omega }`$, the relative velocity equals
$`\frac{2\lambda }{1+\omega ^2\rho ^2}(\omega ^2\rho ^2๐ฎ_c+๐๐ช)`$ and the relative acceleration is $`\frac{(1\omega ^2\rho ^2)\omega ^2}{(1+\omega ^2\rho ^2)^2}๐ช`$, so the angular velocity of Thomas precession has value
$$\frac{\lambda }{(1+\omega ^2\rho ^2)\rho ^2}(\omega ^2\rho ^2๐ฎ_c๐๐ช)๐ช=\frac{\lambda }{(1+\omega ^2\rho ^2)}(๐\omega ^2๐ฎ_c)๐ช.$$
(41)
## 12 Discussion
The systematic use of absolute velocities instead of relative ones gives us a nice form of the Lorentz boosts which results in extremely brief and simple derivation of
โ the discrete Thomas rotation due to successive Lorentz boosts,
โ the Fermi-Walker equation,
โ the Thomas rotation on a circular world line,
โ Thomas rotations in general,
โ the Thomas precession with respect to an inertial frame,
and it allows us a deeper insight into the nature of Thomas rotations and Thomas precessions. It is an important fact that the Thomas rotation is absolute i.e. independent of reference frames while the Thomas precession is relative i.e. refers to inertial frames. It is emphasized again that the same gyroscope shows different precessions to different inertial frames.
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# Quantum phases of a Feshbach-resonant atomic Bose gas in one dimension
## Introduction
Trapped dilute cold atomic gases are one of the most exciting fields in condensed matter physics.A An important recent development in this area is the application of Feshbach resonances. The energy difference between the molecular state and the two-atom continuum, known as the detuning $`\nu `$, can be experimentally tuned by means of a magnetic field. Therefore, by sweeping the magnetic field from positive to negative detuning through the Feshbach resonance, it is actually possible to form molecules in the atomic gas.JBA
Depending on the quantum statistics of atoms, the low temperature properties of a dilute atomic gas with an $`s`$-wave Feshbach resonance will be much different. In the case of fermonic atoms, a crossover from a superfluid (SF) of the Bardeen-Cooper-Schrieffer type to the Bose-Einstein condensation (BEC) is predicted.OG On the other hand, a quantum phase transition can occur for bosonic atoms by changing the value of magnetic field detuning.RPW There, two thermodynamically distinct phases exist at zero temperature: the โatomic superfluidโ (ASF) phase with both atomic BEC and molecular BEC and the โmolecular superfluidโ (MSF) with molecular BEC only.
In the present letter, we study the possible quantum phases for an atomic Bose gas with an $`s`$-wave Feshbach resonance in a one-dimensional optical lattice. It can be described by the Hamiltonian
$`H`$ $`=`$ $`{\displaystyle \underset{i}{}}(t_mb_{i+1}^{}b_i+t_aa_{i+1}^{}a_i+.c.)`$ (1)
$`{\displaystyle \underset{i}{}}[\mu n_{a,i}+(2\mu \nu )n_{b,i}]`$
$`+{\displaystyle \underset{i}{}}{\displaystyle \frac{U_a}{2}}n_{a,i}(n_{a,i}1)+{\displaystyle \frac{U_m}{2}}n_{b,i}(n_{b,i}1)`$
$`+{\displaystyle \underset{i}{}}U(b_i^{}a_ia_i+.c.)+U_{ab}{\displaystyle \underset{i}{}}n_{a,i}n_{b,i}.`$
Here $`n_{a(b),i}`$ is the atom (molecule) density operator and $`\mu `$ is the chemical potential. This system can also be considered a kind of binary mixtures with two types of bosons. However, a main distinction between the present case and the binary mixture studied previouslyD ; CH is that a Josephson coupling between atoms and molecules, the $`U`$ term in Eq. (1), is absent in the latter. As we shall see, this will change the global phase diagram dramatically.
Our main results can be summarized in Fig. 1. For on-site interactions only, i.e. the Hamiltonian given by Eq. (1), the Josephson coupling dominates the low-energy physics in most of the parameter region, and the resulting phase exhibits strong superfluid fluctuations in both atomic and molecular sectors. This is because the phase fluctuations of both sectors are locked by the Josephson coupling. This phase is, in fact, a one-dimensional ($`1d`$) analog of the ASF phase in three dimensions ($`3d`$). (But we do not have real condensates here.) Therefore, it will be referred to as $`1d`$ ASF. The $`1d`$ ASF will collapse for sufficiently strong attractions between atoms and molecules ($`U_{ab}<0`$). When the density difference between atoms and molecules is close to some rational number, depending on the values of parameters as shown in Fig. 1 (b), two additional phases may emerge by including nearest-neighbor repulsions between atoms (molecules): the two component Luttinger liquid ($`2`$-LL) where both the atomic and molecular sectors are gapless, and the inter-channel charge density wave (I-CDW) where the relative density fluctuations between atoms and molecules are frozen at low energy.foot1
## The continuum theory
We now outline below the derivation of our results. For simplicity, we shall consider the case where both the densities of atoms and molecules, $`\rho _a=n_{a,i}/a_0`$ and $`\rho _b=n_{b,i}/a_0`$, respectively, are incommensurate with the lattice. (Here $`a_0`$ is the lattice spacing. $`\rho _a`$ and $`\rho _b`$ satisfy the constraint: $`\rho _a+2\rho _b=\rho _0`$ where $`\rho _0`$ is the density of bare atoms.) That is, we do not consider the possibility of the SF-Mott insulator transition. Then, in terms of the โbosonizationโ formulaG : $`a_ie^{i\sqrt{\pi }\theta _a(x)}_{n=\mathrm{}}^{\mathrm{}}e^{i2\pi n\rho _ax}e^{i\sqrt{4\pi }n\varphi _a(x)}`$ and $`b_ie^{i\sqrt{\pi }\theta _m(x)}_{n=\mathrm{}}^{\mathrm{}}e^{i2\pi n\rho _bx}e^{i\sqrt{4\pi }n\varphi _m(x)}`$, the low-energy physics of $`H`$ \[Eq. (1)\] can be described by the following effective Hamiltonian:
$`H_{eff}`$ $`=`$ $`{\displaystyle \frac{v_a}{2}}{\displaystyle ๐x\left[K_a\left(_x\theta _a\right)^2+\frac{1}{K_a}\left(_x\varphi _a\right)^2\right]}`$ (2)
$`+{\displaystyle \frac{v_m}{2}}{\displaystyle ๐x\left[K_m\left(_x\theta _m\right)^2+\frac{1}{K_m}\left(_x\varphi _m\right)^2\right]}`$
$`+g_1{\displaystyle ๐x\mathrm{cos}[\sqrt{\pi }(\theta _m2\theta _a)]}`$
$`+g_2{\displaystyle ๐x_x\varphi _a_x\varphi _m}`$
$`+g_3{\displaystyle ๐x\mathrm{cos}[\sqrt{4\pi }(\varphi _a\varphi _m)+2\pi \delta x]},`$
where $`\delta =\rho _a\rho _b`$. $`v_{a/m}`$ and $`K_{a/m}`$ are sound velocities and Luttinger liquid (LL) parameters, respectively. In Eq (2), only those terms which may become the most relevant in the renormalization group (RG) sense are retained. The values of $`v_{a/m}`$, $`K_{a/m}`$, $`g_1`$, $`g_2`$, and $`g_3`$ depend on the short-distance physics. In general, they must be extracted from numerics or experiments. $`K_{a/m}1`$ in the weak coupling regime. On the other hand, $`K_{a/m}=1`$ in the Tonks limit, i.e. $`U_{a(m)}/t_{a(m)}+\mathrm{}`$. Therefore, for on-site interactions only, $`1K_{a/m}<+\mathrm{}`$.G The value of $`K_a`$ ($`K_m`$) can be further decreased by including nearest-neighbor repulsions between atoms (molecules).
The $`g_3`$ term is a Umklapp process. It can be neglected when $`\delta `$ is not close to zero (incommensurate filling). On the other hand, when $`\delta `$ is close to zero (commensurate filling), the $`g_3`$ term can affect low-energy physics and one may no longer neglect it in Eq. (2). When $`a_0|\delta |`$ is close to some rational number $`k/l`$ where $`k=1,2,\mathrm{}`$, $`l=2,3,\mathrm{}`$, and $`k`$ and $`l`$ are co-prime with one another, other Umklapp processes must be taken into account, and one must include the following term in Eq. (2):
$$\stackrel{~}{g}_3๐x\mathrm{cos}[\sqrt{4\pi }l(\varphi _a\varphi _m)+2\pi l\delta x].$$
(3)
We shall focus on Eq. (2) in the following and discuss the effects of the $`\stackrel{~}{g}_3`$ term later.
## Incommensurate filling
When $`\delta `$ is incommensurate with the lattice, one may neglect the $`g_3`$ term in Eq. (2). To analyze the effects of the $`g_1`$ and $`g_2`$ terms on low-energy physics, we resort to the RG method. A perturbative calculation up to the one-loop order already shows that the $`g_1`$ and $`g_2`$ terms alone do not form a closed operator algebra in the sense of operator product expansion (OPE). One must include the term $`_x\theta _a_x\theta _m`$ in $`H_{eff}`$ \[Eq. (2)\]. To simplify the analysis, we consider the case $`v_a=v_m=v_0`$. (The effects of velocity anisotropy will be discussed later.) Further, we rescale $`g_1`$ by $`a_0^2g_1g_1`$. The one-loop RG equations can be obtained by setting $`\lambda _3(l)=0`$ in Eqs. (10) โ (15). By solving these scaling equations, one may obtain the following results: Within the weak-coupling region, i.e. $`\pi |g_1|/(2v_0),|g_2|/(2v_0)1`$, the $`g_1`$ term is relevant in the regime $`D_1<2+\frac{\sqrt{2}\pi |g_1|}{v_0}`$, while it becomes irrelevant for $`D_1>2+\frac{\sqrt{2}\pi |g_1|}{v_0}`$, where $`D_1`$ is the scaling dimension of the $`g_1`$ term, defined by
$$D_1=\frac{1}{K_a}+\frac{1}{4K_m}.$$
(4)
Therefore, there are two zero-temperature phases. It turns out that the phase transition between these two phases is of the KT type. We note that for on-site interactions only, i.e. $`1K_{a/m}<+\mathrm{}`$, the $`g_1`$ term is always relevant from Eq. (4).
In the regime where the $`g_1`$ term is relevant, it is convenient to define new bosonic fields: $`\theta _\pm \theta _a\pm \frac{1}{2}\theta _m`$ and $`\varphi _\pm \frac{1}{2}\varphi _a\pm \varphi _m`$. Then, the value of $`\theta _{}`$ is pinned and a gap $`\mathrm{\Delta }_0[\pi |g_1|/(2v_0)]^{2D_1}`$ is opened for the $`\theta _{}`$ sector. That is, the phase fluctuations of atoms and molecules are locked by the Josephson coupling. By integrating out the gapped sector, the low-energy effective Hamiltonian describing the gapless ($`\theta _+`$) sector takes the form of LLs:
$$H_+=\frac{v}{2}๐x\left[K\left(_x\theta _+\right)^2+\frac{1}{K}\left(_x\varphi _+\right)^2\right].$$
(5)
For $`\pi |g_1|/(2v_0),|g_2|/(2v_0)1`$, $`v`$ and $`K`$ can be related to the short distance variables with the help of the one-loop RG equations, yielding
$`{\displaystyle \frac{v}{v_0}}`$ $`=`$ $`\sqrt{({\displaystyle \frac{K_a^{}}{4}}+K_m^{}+2\lambda _4^{})\left({\displaystyle \frac{1}{K_a^{}}}+{\displaystyle \frac{1}{4K_m^{}}}+2\lambda _2^{}\right)},`$
$`K`$ $`=`$ $`\sqrt{{\displaystyle \frac{K_a^{}/4+K_m^{}+2\lambda _4^{}}{1/K_a^{}+1/(4K_m^{})+2\lambda _2^{}}}},`$ (6)
where $`K_a^{}=K_a+\frac{1}{2D_1}`$, $`K_m^{}=K_m+\frac{1}{4(2D_1)}`$, $`\lambda _2^{}=\frac{g_2}{2v_0}+\frac{1}{2K_aK_m(2D_1)}`$, and $`\lambda _4^{}=\frac{1}{2(2D_1)}`$. Equation (6) indicates that the $`1d`$ ASF will become unstable provided that the inequality is satisfied: $`\frac{1}{K_a^{}}+\frac{1}{4K_m^{}}<2\lambda _2^{}`$. This is because $`v0`$, indicating an instability of this system. In other words, the system will collapse for sufficiently strong attraction between atoms and molecules.
The $`1d`$ ASF can be characterized by the single-particle Green functions of atoms and molecules:
$`\mathrm{\Psi }_a(๐)\mathrm{\Psi }_a^{}(0)=A_1\left({\displaystyle \frac{a_0}{r}}\right)^{\alpha _1}+\mathrm{},`$
$`\mathrm{\Psi }_m(๐)\mathrm{\Psi }_m^{}(0)=A_2\left({\displaystyle \frac{a_0}{r}}\right)^{\alpha _2}+\mathrm{},`$ (7)
where $`๐=(\tau ,x)`$, $`r=\sqrt{(v_+\tau )^2+x^2}`$, $`\mathrm{\Psi }_a(x)=a_i/\sqrt{a_0}`$, $`\mathrm{\Psi }_m(x)=b_i/\sqrt{a_0}`$, $`A_{1,2}`$ are nonuniversal constants, and
$$\alpha _2=4\alpha _1=\frac{1}{2K}.$$
(8)
We would like to stress that this exact relation between $`\alpha _1`$ and $`\alpha _2`$ \[Eq. (8)\] results from the fact that the phase fluctuations of atoms and molecules are locked by the Josephson coupling, and is a characteristic of the $`1d`$ ASF. On the other hand, the density correlation functions of atoms and molecules at nonzero momenta decay exponentially due to the presence of the gap $`\mathrm{\Delta }_0`$. Equation (8) implies that both the atomic and molecular sectors exhibit the behavior of a $`1d`$ SF as long as $`K>1/4`$. For $`1K_{a/m}<+\mathrm{}`$, this condition is satisfied, which can be verified from Eq. (6). It gives further support to our claim that this phase is a $`1d`$ analog of the ASF in $`3d`$.
In the regime where the $`g_1`$ term is irrelevant, the excitation spectrum consists of two branches of gapless excitations with linear dispersion relations. This phase, which has been thoroughly discussed in Ref. CH, , will be referred to as the $`2`$-LLs. We just mention that in the present situation this phase can exist only by including sufficiently strong repulsions between atoms or molecules. This is because $`D_1<2`$ for on-site interactions only.
## Commensurate filling
When $`\delta `$ is close to zero, the $`g_3`$ term must be retained. For simplicity, we shall consider $`\delta =0`$. Let us first define the scaling dimension of the $`g_3`$ term:
$$D_2=K_a+K_m.$$
(9)
The $`g_1`$ ($`g_3`$) term will be relevant if $`D_1<2`$ ($`D_2<2`$). Accordingly, a tree-level scaling analysis suggests the phase diagram in the $`K_a`$-$`K_m`$ space as shown in Fig. 2 (a). $`D_1<2`$ and $`D_2>2`$ corresponds to the $`1d`$ ASF, $`D_1>2`$ and $`D_2<2`$ will be referred to as the I-CDW, and $`D_1,D_2>2`$ corresponds to the $`2`$-LL. However, a competition between the two relevant operators, the $`g_1`$ and $`g_3`$ terms, occur when $`D_1,D_2<2`$. To determine whether Region IV in Fig. 2 (a) corresponds to a new phase or not, we employ the one-loop RG analysis.
By integrating out the fast modes, the one-loop RG equations are given by
$`{\displaystyle \frac{dK_a(l)}{dl}}`$ $`=`$ $`2\left[\lambda _1^2(l)K_a^2(l)\lambda _3^2(l)\right],`$ (10)
$`{\displaystyle \frac{dK_m(l)}{dl}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[\lambda _1^2(l)4K_m^2(l)\lambda _3^2(l)\right],`$ (11)
$`{\displaystyle \frac{d\lambda _1(l)}{dl}}`$ $`=`$ $`\left[2{\displaystyle \frac{1}{K_a(l)}}{\displaystyle \frac{1}{4K_m(l)}}\right]\lambda _1(l)`$ (12)
$`{\displaystyle \frac{\lambda _1(l)\lambda _4(l)}{K_a(l)K_m(l)}},`$
$`{\displaystyle \frac{d\lambda _2(l)}{dl}}`$ $`=`$ $`{\displaystyle \frac{\lambda _1^2(l)}{K_a(l)K_m(l)}}2\lambda _3^2(l),`$ (13)
$`{\displaystyle \frac{d\lambda _3(l)}{dl}}`$ $`=`$ $`[2K_a(l)K_m(l)]\lambda _3(l)`$ (14)
$`2K_a(l)K_m(l)\lambda _2(l)\lambda _3(l),`$
$`{\displaystyle \frac{d\lambda _4(l)}{dl}}`$ $`=`$ $`2K_a(l)K_m(l)\lambda _3^2(l)\lambda _1^2(l),`$ (15)
with the initial values: $`K_a(0)=K_a`$, $`K_m(0)=K_m`$, $`\lambda _{1(3)}(0)=\pi g_{1(3)}/(2v_0)`$, $`\lambda _2(0)=g_2/(2v_0)`$, and $`\lambda _4(0)=0`$. By solving Eqs. (10) โ (15), one may find three kinds of behaviors of the RG flow of $`\lambda _1(l)`$ and $`\lambda _3(l)`$: (i) $`\lambda _1(l)`$ flows to strong coupling while $`\lambda _3(l)`$ flows to zero. This is the $`1d`$ ASF. (ii) $`\lambda _3(l)`$ flows to strong coupling while $`\lambda _1(l)`$ flows to zero. This is the I-CDW. (iii) Both $`\lambda _1(l)`$ and $`\lambda _3(l)`$ flow to zero. This is the $`2`$-LL. Thus, Region IV in Fig. 2 (a) shrinks to a transition line between the $`1d`$ ASF and I-CDW in the $`K_a`$-$`K_m`$ space as shown in Fig. 2 (b). Both the phase transition between the $`1d`$ ASF and $`2`$ LL and that between the I-CDW and $`2`$-LL belong to the KT type. The transition between the $`1d`$ ASF and I-CDW is of second order. Further, all these transition lines coincide at two tricritical points, the point $`A`$ and $`B`$ in Fig. 2 (b). The very reason why there is no way for both $`\lambda _1(l)`$ and $`\lambda _3(l)`$ flowing to strong coupling simultaneously is that the operators $`\mathrm{cos}\sqrt{\pi }(\theta _m2\theta _a)`$ and $`\mathrm{cos}\sqrt{4\pi }(\varphi _a\varphi _m)`$ are exclusive to one another, that is, the field configurations which minimize one perturbation term do not minimize the other. The interplay between these two competing relevant operators then produces a novel quantum phase transition.
In the I-CDW, it is convenient to define new bosonic fields: $`\stackrel{~}{\varphi }_+=\frac{1}{2}(\varphi _a+\varphi _m)`$, $`\stackrel{~}{\varphi }_{}=\varphi _a\varphi _m`$, $`\stackrel{~}{\theta }_+=\theta _a+\theta _m`$, and $`\stackrel{~}{\theta }_{}=\frac{1}{2}(\theta _a\theta _m)`$. The $`\stackrel{~}{\varphi }_+`$ and $`\stackrel{~}{\varphi }_{}`$ fields describe the in-phase and out-of-phase density fluctuations, respectively. Due to the relevant perturbation $`\mathrm{cos}\sqrt{4\pi }\stackrel{~}{\varphi }_{}`$, the value of $`\stackrel{~}{\varphi }_{}`$ is pinned and a gap is opened for the $`\stackrel{~}{\varphi }_{}`$ sector, while the $`\stackrel{~}{\varphi }_+`$ sector is still gapless. On account of this, both the single-particle Green functions of atoms and molecules decay exponentially. On the other hand, the $`2\pi \rho _{a/b}`$ parts of the density fluctuations for atoms and molecules are enhanced:
$$\rho _{a(b)}(x)\rho _{a(b)}(0)|_{2\pi \rho _{a(b)}}\left(\frac{a_0}{|x|}\right)^\gamma ,$$
(16)
with $`\gamma <2`$ for $`\pi |g_1|/(2v_0),|g_2|/(2v_0)1`$. Here $`\rho _{a/b}(x)=n_{a/b,i}/a_0`$.
To understand the nature of the ground state of the I-CDW, a simple picture can be obtained in the limit of strong atom-molecule interactions ($`g_3`$), where the potential energy (the $`g_3`$ term) dominates over quantum fluctuations. In this case, atoms and molecules form a regular lattice (Wigner crystal of hard-core bosons). For $`g_3>0`$, the energy of the repulsion between atoms and molecules
$`g_3\mathrm{cos}(\sqrt{4\pi }\varphi _{})=g_3\mathrm{cos}[\sqrt{4\pi }(\varphi _a\varphi _m+\sqrt{\pi }/2)],`$
is minimized by a relative phase shift of $`\sqrt{\pi }/2`$ between atoms and molecules, which corresponds to a shift of the atom (or molecule) lattice by half-a-period.SMH Thus, the I-CDW respects the symmetry of translation by one site, $`a_ia_{i+1}`$ and $`b_ib_{i+1}`$, but spontaneously breaks the reflection symmetry about the origin, $`a_ia_i`$ and $`b_ib_i`$ (or $`\varphi _{a/m}\varphi _{a/m}`$ and $`\theta _{a/m}\theta _{a/m}`$).
## Experimental signatures
We suggest a few methods to detect the above results experimentally. The most important distinction between the 2-LL and the $`1d`$ ASF is the behaviors of the density correlators. For the 2-LLs, the $`2\pi \rho _{a(b)}`$ part of the density correlators of atoms (molecules) exhibits power-law decay. However, the corresponding sector decays exponentially for the $`1d`$ ASF. Therefore, a threshold behavior will be observed at the momentum $`k=2\pi \rho _{a(b)}`$ in a time-of-flight measurement for atoms (molecules) in the $`1d`$ ASF, while such a behavior does not exist in the 2-LL. The other unique feature of the ASF is that the ground state is a coherent state formed by hybridizing the atoms and molecules. Therefore, by suddenly changing the detuning, some kind of Rabi oscillation will be observed between the atomic and molecular condensates. Such a change of detuning can be achieved by applying magnetic field pulses to the ASF.KH Although there is no true condensate in $`1d`$ due to strong phase fluctuations, a similar oscillation can also be observed between the densities of atoms and molecules, which is strongly damped by the phase fluctuations. Further support to the $`1d`$ ASF is the examination of Eq. (8), which can be achieved by a measurement of single-particle Green functions of atoms and molecules through the absorption line shape.
## Discussions
Finally, two points should be addressed here. First of all, when $`a_0|\delta |`$ is close to some rational number, one must replace the $`g_3`$ term in Eq. (2) by the $`\stackrel{~}{g}_3`$ term \[Eq. (3)\], with the scaling dimension $`D_2l`$. A tree-level scaling analysis suggests the phase diagram in the $`K_a`$-$`K_m`$ space as shown in Fig. 3. We note that there is no direct transition between the $`1d`$ ASF and I-CDW. Moreover, the I-CDW can be reached only with very small values of $`K_{a/m}`$. Next, a RG study shows that small velocity anisotropy, i.e. $`v_av_m`$, just shifts the phase boundary, and does not affect our conclusions. This is because the OPEโs between the additional operators arising from velocity anisotropy and the already existing operators do not generate themselves.
When this work was finished, we were aware of two recent papers dealing with a related system โ a Feshbach-resonant atomic Fermi gas in $`1d`$.RO Their results about the charge sector are similar to ours for the $`1d`$ ASF.
###### Acknowledgements.
The work of Y.-W. Lee is supported by the National Science Council of Taiwan under grant NSC 93-2112-M-029-007. The work of Y.L. Lee is supported by the National Science Council of Taiwan under grant NSC 93-2112-M-018-009.
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# Symplectic reflection algebras and non-homogeneous ๐-Koszul property
## 1 Introduction
Let $`V`$ be a finite dimensional complex vector space which is endowed with a symplectic 2-form $`\omega `$. Let $`\mathrm{\Gamma }`$ be a finite subgroup of Sp$`(V)`$ and $`(TV)\mathrm{\#}\mathrm{\Gamma }`$ the smash product of the tensor algebra $`TV`$ of $`V`$ with the group algebra $`\mathrm{\Gamma }`$ of $`\mathrm{\Gamma }`$. For any $`g\mathrm{\Gamma }`$, introduce the subspaces $`M_g:=\mathrm{Im}(\mathrm{Id}g)`$ and $`L_g:=\mathrm{Ker}(\mathrm{Id}g)`$, so that one has $`V=M_gL_g`$ and
$$\mathrm{\Lambda }^2(V)=(\mathrm{\Lambda }^2(M_g))(M_gL_g)(\mathrm{\Lambda }^2(L_g)).$$
(1.1)
Recall that the integer $`a(g):=dimM_g`$ is even, and that $`g`$ is called a *symplectic reflection* if this dimension is 2. Define the $``$-linear map $`\psi _g:\mathrm{\Lambda }^2(V)`$ in order to coincide with $`\omega `$ on $`\mathrm{\Lambda }^{a(g)}(M_g)\mathrm{\Lambda }^{2a(g)}(L_g)`$ and to vanish on the other components of (1.1). Clearly, $`\psi _g=0`$ if $`g`$ is neither $`Id`$ nor a symplectic reflection.
Writing $`\psi =_{g\mathrm{\Gamma }}\psi _gg`$, we define a skew-symmetric $``$-bilinear pairing $`\psi :V\times V\mathrm{\Gamma }`$ which is $`\mathrm{\Gamma }`$-equivariant. Then the *symplectic reflection algebra* is the $`\mathrm{\Gamma }`$-algebra H<sub>ฯ</sub> defined by
$$H_\psi =(TV\mathrm{\#}\mathrm{\Gamma })/I(xyyx\psi (x,y);x,yV),$$
Actually, for any map $`m:\mathrm{\Gamma }`$ which is constant on any conjugation class, $`m\psi `$ is $`\mathrm{\Gamma }`$-equivariant, and H<sub>mโ
ฯ</sub> is also called a symplectic reflection algebra.
The algebra H<sub>mโ
ฯ</sub> is naturally filtered, and there is a natural graded algebra morphism $`H_{0\psi }=(SV)\mathrm{\#}\mathrm{\Gamma }gr(H_{m\psi })`$. The first fundamental feature of H<sub>mโ
ฯ</sub> is that this morphism is an *isomorphism* : it is the so-called PBW property for H<sub>mโ
ฯ</sub>. The aim of this paper is to show that such a property makes sense in a more general context than the quadratic one, and that the generalized PBW property holds for a new class of H<sub>mโ
ฯ</sub>. Let us describe briefly this new class.
Now $`p`$ is an integer with $`2pdimV`$, $`\mathrm{\Gamma }`$ is a finite subgroup of GL$`(V)`$, and $`\varphi :\mathrm{\Lambda }^pVk`$ is a $`\mathrm{\Gamma }`$-invariant linear map (playing the role of $`\omega `$). We have an analogous decomposition (1.1) for $`\mathrm{\Lambda }^p(V)`$, and analogous definitions for the $`\psi _g`$โs, $`\psi `$, $`m\psi `$, and H<sub>mโ
ฯ</sub>. For the relations of H<sub>mโ
ฯ</sub>, we replace the 2-tensor $`xyyx`$ by any totally skew-symmetric tensor of $`p`$ variables, while $`m\psi `$ is applying on these variables. We shall prove the following result (Corollary 4.5 below) which states a PBW property for the new class.
###### Theorem 1.1
Under the previous notations and assumptions, the natural graded algebra morphism $`H_{0\psi }gr(H_{m\psi })`$ is an isomorphism.
The PBW property in the case of any filtered algebra with (inhomogeneous) quadratic relations is well understood in the setting due to Braverman and Gaitsgory (see also Polishchuk and Positselski ). This setting brings out the fundamental role of the Koszul property (in Priddyโs sense) of the homogeneous quadratic algebra which is obtained by forgetting the non-quadratic part.
The first author has extended the Koszul property to any algebra with $`N`$-homogeneous relations ($`N`$ is fixed $`2`$. So an analog of the PBW theorem for any filtered algebra with $`N`$-inhomogeneous relations is naturally the first step for our proof of Theorem 1.1. We shall prove such a PBW theorem under very large assumptions, including the fact that the ground field $``$ has to be replaced by the group algebra $`\mathrm{\Gamma }`$ which is a non-commutative ring! Fortunately, this ring is semi-simple, and our formalism for Koszul and PBW properties works out even for more general rings, the von Neumann regular rings. Notice that the extension of the Koszul property to any semi-simple ground ring was already performed in in the quadratic case. Our PBW theorem is the following (Theorem 3.4 below).
###### Theorem 1.2
*(PBW theorem in the $`N`$-case)* Assume that $`k`$ is a von Neumann regular ring, $`V`$ is a $`k`$-$`k`$-bimodule, $`N`$ is an integer $`2`$, and $`P`$ is a sub-$`k`$-$`k`$-bimodule of $`F^N`$, where $`F^n=_{0in}V^i`$ for any $`n0`$. Set $`U=T(V)/I(P)`$ and $`A=T(V)/I(R)`$, where $`R=\pi (P)`$ and $`\pi `$ is the projection of $`F^N`$ onto $`V^N`$ modulo $`F^{N1}`$.
Assume that the graded left $`k`$-module Tor$`{}_{3}{}^{A}(k_A,_Ak)`$ is concentrated in degree $`N+1`$ (this property is a part of the Koszul property of $`A`$). Then the conditions
$$PF^{N1}=0,$$
$$(PV+VP)F^NP,$$
imply that the PBW property holds, i.e., the natural algebra morphism $`Agr(U)`$ is an isomorphism.
This theorem suggests to extend the terminology โ$`A`$ is Koszulโ to โ$`U`$ is Koszulโ (see Definition 3.9 below). It is natural since J. L. Koszul introduced his resolution for the polynomial algebra in the *filtered* context of the enveloping algebra of a Lie algebra, and used the classical PBW property as a trick to carry over the exactness of his resolution to the standard complex .
## 2 $`N`$-Koszul algebras over von Neumann regular rings
A ring $`k`$ is said to be *von Neumann regular* if for every $`xk`$, there exists $`yk`$ such that $`xyx=x`$. In this text, von Neumann regular rings will be used only through the following characterization (in which left can be replaced by right) .
###### Proposition 2.1
A ring $`k`$ is von Neumann regular if and only if all left $`k`$-modules are flat.
Von Neumann regular rings are exactly the rings having weak dimension 0 . So a semi-simple ring (i.e., a ring having left or right global dimension 0) is von Neumann regular. An infinite product $`k`$ of fields is von Neumann regular, but is not semi-simple. The same holds for $`k=End(V)`$, where $`V`$ is an infinite-dimensional vector space.
Throughout this section, $`k`$ is a von Neumann regular ring, and $`V`$ denotes a graded $`k`$-$`k`$-bimodule which is *concentrated in degree 1*. The tensor power (over $`k`$) $`V^n`$ for $`n=0,1,\mathrm{}`$ is a graded $`k`$-$`k`$-bimodule which is concentrated in degree $`n`$. The direct sum $`T(V)=_{n0}V^n`$ is naturally a $``$-graded ring, and a $`k`$-$`k`$-bimodule whose left or right actions coincide with the products in the ring by elements of $`V^0=k`$. We sum up these properties by saying that $`T(V)`$ is a connected $``$-graded $`k`$-$`k`$-algebra.
If $`a`$ and $`b`$ are in $`T(V)`$, their product in $`T(V)`$ is denoted by $`ab`$. For any sub-$`k`$-$`k`$-bimodules $`E`$ and $`F`$ of $`T(V)`$, $`EF`$ denotes the sub-$`k`$-$`k`$-bimodule formed by finite sums of products $`ab`$, $`aE`$, $`bF`$. The two-sided ideal $`I(E)`$ of $`T(V)`$ generated by $`E`$ is such that
$$I(E)=T(V)ET(V)=\underset{i,j0}{}V^iEV^j.$$
On the other hand, if $`EV^i`$ and $`FV^j`$, the canonical homomorphism $`E_kFV^{(i+j)}`$ is injective since the right $`k`$-module $`E`$ and the left $`k`$-module $`F`$ are *flat*. So the $`k`$-$`k`$-bimodules $`E_kF`$ and $`EF`$ will be identified. The following result is known (, chap.I, 2, n 6) and is again a consequence of flatness.
###### Lemma 2.2
If $`E`$ and $`E^{}`$ are sub-$`k`$-$`k`$-bimodules of $`V^i`$, and if $`F`$ and $`F^{}`$ are sub-$`k`$-$`k`$-bimodules of $`V^j`$, the following formulas hold :
(i) $`(EF)(EF^{})=E(FF^{})`$,
(ii) $`(EF)(E^{}F)=(EE^{})F`$,
(iii) $`(E^{}F)(EF^{})=E^{}F^{}`$ if moreover $`E^{}E`$ and $`F^{}F`$.
Now fix an integer $`N2`$ and a sub-$`k`$-$`k`$-bimodule $`R`$ of $`V^N`$. Then the two-sided ideal $`I(R)`$ is graded by the $`k`$-$`k`$-bimodules
$$I(R)_n=\underset{i+N+j=n}{}V^iRV^j,n0.$$
One has $`I(R)_n=0`$ if $`0nN1`$. In this section, we are interested in the connected $``$-graded $`k`$-$`k`$-algebra $`A=T(V)/I(R)`$. The gradation of $`A`$ is formed by the $`k`$-$`k`$-bimodules $`A_n=V^n/I(R)_n`$. One has $`A_n=V^n`$ if $`0nN1`$. The natural projection $`ฯต:AA_0=k`$ makes $`k`$ as being an $`A`$-$`A`$-bimodule, denoted by $`{}_{A}{}^{}k_{A}^{}`$. We also use $`{}_{A}{}^{}k`$ and $`k_A`$ for the left and right associated $`A`$-modules.
The categories of left $`A`$-modules, right $`A`$-modules, $`A`$-$`A`$-bimodules are respectively denoted by $`A`$-Mod, Mod-$`A`$, $`A`$-Mod-$`A`$. When the objects are graded and the arrows are homogeneous of degree 0, the categories are respectively denoted by $`A`$-grMod, grMod-$`A`$, $`A`$-grMod-$`A`$. For example, the object $`{}_{A}{}^{}k`$ is graded (as concentrated in degree 0), so that $`ฯต`$ is an arrow of $`A`$-grMod. Any object $`M`$ of $`A`$-grMod has a free resolution in $`A`$-grMod, i.e., a resolution by graded-free left $`A`$-modules. Accordingly, for any natural number $`n`$, Tor$`{}_{n}{}^{A}(k_A,M)`$ is a graded left $`k`$-module. Koszul property of $`A`$ will be defined from the objects Tor$`{}_{n}{}^{A}(k_A,_Ak)`$, $`n`$. Let us begin to compute these objects for $`n=0,1,2`$.
Actually, for computing Tor$`{}_{n}{}^{A}(k_A,_Ak)`$, it is enough to have a *flat* resolution of $`{}_{A}{}^{}k`$ in $`A`$-grMod, i.e., a resolution in $`A`$-grMod formed by flat left $`A`$-modules (an object of $`A`$-grMod which is flat in $`A`$-Mod is flat in $`A`$-grMod). On the other hand, it is well-known (and easy to prove) that, for any homomorphism from a ring $`k`$ to a ring $`A`$ and any flat left $`k`$-module $`E`$, the left $`A`$-module $`A_kE`$ is flat (, chap.I, 2, n 7). Thus, since $`k`$ is von Neumann regular, $`A_kE`$ is a *flat* left $`A`$-module (not free in general!) for any left $`k`$-module $`E`$. So we search resolutions in $`A`$-grMod whose objects are of the type $`A_kE`$, $`E`$ object of $`k`$-grMod.
###### Lemma 2.3
Fix the natural numbers $`n`$ and $`m`$, and let $`E`$ be a sub-$`k`$-$`k`$-bimodule of $`V^m`$. The $`k`$-$`k`$-bimodules $`A_n_kE`$ and $`V^nE/I(R)_nE`$ are naturally isomorphic.
Proof. As $`E`$ is flat in $`k`$-Mod, one has the natural exact sequence of $`k`$-$`k`$-grMod
$$0I(R)_n_kEV^n_kEA_n_kE0.$$
But $`V^n_kE`$ is identified to the sub-$`k`$-$`k`$-bimodule $`V^nE`$ of $`T(V)`$, and the image of the injective map is identified to the sub-$`k`$-$`k`$-bimodule $`I(R)_nE`$ of $`T(V)`$.
The arrow $`ฯต:A_Ak`$ of $`A`$-grMod has $`A_1=_{n1}A_n`$ as kernel. The inclusion $`VA`$ defines an injective (by flatness!) natural arrow $`A_kVA_kA`$, which is composed with the multiplication $`\mu :A_kAA`$ to define an arrow $`\delta _1:A_kVA`$ of $`A`$-grMod. In degree 0, $`\delta _1`$ vanishes, and in degree $`n1`$, it is identified by Lemma 2.3 to the canonical map
$$\frac{V^n}{I(R)_{n1}V}\frac{V^n}{I(R)_n},$$
which is surjective. So im$`(\delta _1)=\mathrm{ker}(ฯต)`$.
The natural injections $`RV^N=V^{(N1)}_kVA_kV`$ define a natural injection $`A_kRA_k(A_kV)(A_kA)_kV`$, which is composed with $`\mu _k1_V`$ to define an arrow $`\delta _2:A_kRA_kV`$ of $`A`$-grMod. In degrees $`<N`$, $`\delta _2`$ vanishes, and in degree $`nN`$, it is identified to the canonical map
$$\frac{V^{(nN)}R}{I(R)_{nN}R}\frac{V^n}{I(R)_{n1}V},$$
(2.1)
so that one has
$$(\mathrm{im}(\delta _2))_n=\frac{V^{(nN)}R+I(R)_{n1}V}{I(R)_{n1}V}.$$
On the other hand, for $`n1`$, one has
$$(\mathrm{ker}(\delta _1))_n=\frac{I(R)_n}{I(R)_{n1}V}.$$
Comparing the two equalities, we get im$`(\delta _2)=\mathrm{ker}(\delta _1)`$. Clearly, $`\mathrm{ker}(\delta _2)`$ vanishes in degrees $`N`$. Choose an object $`E_3`$ of $`k`$-grMod living in degrees $`N+1`$, and a surjective arrow $`E_3\mathrm{ker}(\delta _2)`$ in $`k`$-grMod. Extending the latter by $`A`$-linearity, one gets a surjective arrow $`A_kE_3\mathrm{ker}(\delta _2)`$, which is composed with the inclusion $`\mathrm{ker}(\delta _2)A_kR`$ to define an arrow $`\delta _3:A_kE_3A_kR`$ of $`A`$-grMod. Finally
$$A_kE_3\stackrel{\delta _3}{}A_kR\stackrel{\delta _2}{}A_kV\stackrel{\delta _1}{}A0$$
(2.2)
is the beginning of a flat resolution of $`{}_{A}{}^{}k`$ (via $`ฯต:Ak`$) in $`A`$-grMod.
Applying the functor $`k_A_A`$ to (2.2), the sequence
$$E_3RVk0$$
of $`k`$-grMod is obtained, in which all arrows are vanishing (remember that the arrows in $`k`$-grMod are of degree 0). Therefore, one has the isomorphisms in $`k`$-grMod : Tor$`{}_{0}{}^{A}(k_A,_Ak)k`$, Tor$`{}_{1}{}^{A}(k_A,_Ak)V`$, Tor$`{}_{2}{}^{A}(k_A,_Ak)R`$. Moreover Tor$`{}_{3}{}^{A}(k_A,_Ak)`$ lives in degrees $`N+1`$.
The three first Torโs are concentrated in only one degree (0, 1, $`N`$ respectively). Roughly speaking, the Koszul property means that each Tor$`{}_{n}{}^{A}(k_A,_Ak)`$ is concentrated in the lowest possible degree when $`n=0,1,2,3,\mathrm{}`$ We have to examine now when Tor$`{}_{3}{}^{A}(k_A,_Ak)`$ is concentrated in degree $`N+1`$.
###### Proposition 2.4
The graded left $`k`$-module Tor$`{}_{3}{}^{A}(k_A,_Ak)`$ is concentrated in degree $`N+1`$ if and only if the graded left $`A`$-module $`\mathrm{ker}(\delta _2)`$ is generated in degree $`N+1`$.
Proof. By dimension shifting applied to flat resolutions (, p.47), one has
$$\mathrm{Tor}_3^A(k_A,_Ak)=\mathrm{ker}(k_A_A\mathrm{ker}(\delta _2)\stackrel{k_A_Ai}{}k_A_A(A_kR)R),$$
where $`i:\mathrm{ker}(\delta _2)A_kR`$ is the inclusion. Since $`\mathrm{ker}(\delta _2)`$ lives in degrees $`N+1`$, $`k_A_A\mathrm{ker}(\delta _2)`$ lives in degrees $`N+1`$. But $`R`$ is concentrated in degree $`N`$ and $`k_A_Ai`$ preserves the degrees. Thus Tor$`{}_{3}{}^{A}(k_A,_Ak)=k_A_A\mathrm{ker}(\delta _2)`$. It suffices to prove the following.
###### Lemma 2.5
Fix $`n`$. Let $`M=_{in}M_i`$ be an object of $`A`$-grMod, living in degrees $`n`$. Then $`M`$ is generated in degree $`n`$ if and only if the graded left $`k`$-module $`k_A_AM`$ is concentrated in degree $`n`$.
Proof. The graded left $`k`$-module $`k_A_AM`$ is canonically isomorphic to $`M/A_1.M`$. Assuming $`M=A.M_n`$, one sees that $`M_i=0`$ if $`i<n`$, and for $`i>n`$, $`M_i=A_{in}.M_n`$ is contained in $`(A_1.M)_i`$, hence $`M/A_1.M`$ vanishes in degrees $`n`$. Conversely, assuming that $`M_i=(A_1.M)_i`$ for $`i>n`$, one has $`M_{n+1}=A_1.M_n`$, and inductively $`M_i=A_{in}.M_n`$ for any $`i>n`$, so $`M=A.M_n`$.
Identifying $`\delta _2`$ in degree $`n`$ with (2.1) provides
$$(\mathrm{ker}(\delta _2))_n=\frac{(V^{(nN)}R)(I(R)_{n1}V)}{I(R)_{nN}R},nN+1.$$
(2.3)
It follows
$$(\mathrm{ker}(\delta _2))_n=(V^{(nN)}R)(RV^{(nN)}+\mathrm{}+V^{(nN1)}RV),N+1n2N1,$$
and in particular $`(\mathrm{ker}(\delta _2))_{N+1}=W_{N+1}`$, where the following notation is used
$$W_n=\underset{i+N+j=n}{}V^iRV^j,nN.$$
In our choice of the graded left $`k`$-module $`E_3`$, we can assume that $`(E_3)_{N+1}=W_{N+1}`$, so that $`\delta _3`$ is injective in degree $`N+1`$. Thus $`\mathrm{ker}(\delta _3)`$ lives in degrees $`N+2`$, and
$$(\mathrm{Tor}_3^A(k_A,_Ak))_{N+1}=W_{N+1}.$$
Coming back to our initial question, we have to know when $`\mathrm{ker}(\delta _2)`$ is generated in degree $`N+1`$. It is equivalent to saying that for every $`nN+2`$, one has
$$(\mathrm{ker}(\delta _2))_n=A_{nN1}.W_{N+1}.$$
(2.4)
So a partial condition is obtained for $`N+2n2N1`$ :
$$(V^{(nN)}R)(RV^{(nN)}+\mathrm{}+V^{(nN1)}RV)=V^{(nN1)}W_{N+1}.$$
(2.5)
The set of equalities (2.5) when $`N+2n2N1`$ is called the *extra condition*, and is denoted by (ec). That condition does not occur if $`N=2`$ (hence its name), and it reveals the jump of degrees between generators and relations. It is remarkable that (ec) will be sufficient to take into account the other jumps of degrees which will appear in the inductive process.
Examine now (2.4) for $`n2N`$. Since $`W_{N+1}`$ is a left sub-$`k`$-module of $`V_kR`$, $`A.W_{N+1}`$ is the left sub-$`A`$-module of $`A.(V_kR)A_kR`$ which is the image of the composite
$$A_kW_{N+1}\stackrel{1_Ai}{}A_kV_kR\stackrel{\delta _11_R}{}A_kR,$$
where $`i:W_{N+1}V_kR`$ is the inclusion. Using Lemma 2.3, the maps of this composite become in each degree $`nN+1`$, the canonical following ones
$$\frac{V^{(nN1)}W_{N+1}}{I(R)_{nN1}W_{N+1}}\frac{V^{(nN)}R}{I(R)_{nN1}VR}\frac{V^{(nN)}R}{I(R)_{nN}R}.$$
Therefore
$$A_{nN1}.W_{N+1}=\frac{V^{(nN1)}W_{N+1}+I(R)_{nN}R}{I(R)_{nN}R},$$
and joining that with (2.3), we get the answer to our initial question.
###### Theorem 2.6
The left $`k`$-module Tor$`{}_{3}{}^{A}(k_A,_Ak)`$ is concentrated in degree $`N+1`$ if and only if (ec) and the following relations
$$(V^{(nN)}R)(I(R)_{n1}V)=V^{(nN1)}W_{N+1}+I(R)_{nN}R,n2N,$$
(2.6)
are satisfied.
A triple $`(E,F,G)`$ of sub-$`k`$-$`k`$-bimodules of $`T(V)`$ is said to be *distributive* if $`E(F+G)=EF+EG`$, and the latter equality is called a *distributivity relation*. Show that, if (ec) holds, (2.6) is a distributivity relation for each $`n`$. In fact, setting
$$E=V^{(nN)}R,F=I(R)_{nN}V^N,G=V^{(n2N+1)}I(R)_{2N2}V,$$
the left-hand side of (2.6) is $`E(F+G)`$. Using Lemma 2.2, (iii), one has
$$EF=I(R)_{nN}R.$$
The part (i) of this lemma provides
$$EG=V^{(n2N+1)}[(V^{(N1)}R)(RV^{(N1)}+\mathrm{}+V^{(N2)}RV)],$$
and (ec) for $`n=2N1`$ implies
$$EG=V^{(n2N+1)}[V^{(N2)}W_{N+1}]=V^{(nN1)}W_{N+1},$$
so the right-hand side of (2.6) is $`EF+EG`$.
Theorem 2.6 will be enough to state the PBW theorem, but we want to give further indications about the definition of the Koszul property. The details are left to the reader. Assume now that Tor$`{}_{3}{}^{A}(k_A,_Ak)`$ is concentrated in degree $`N+1`$. In particular, (ec) holds. The inclusion $`W_{N+1}A_kR`$ defines a natural injection $`A_kW_{N+1}A_k(A_kR)(A_kA)_kR`$, which is composed with $`\mu _k1_R`$ to define a more suitable arrow $`\delta _3:A_kW_{N+1}A_kR`$ of $`A`$-grMod. The beginning of the flat resolution is now
$$A_kW_{N+1}\stackrel{\delta _3}{}A_kR\stackrel{\delta _2}{}A_kV\stackrel{\delta _1}{}A0.$$
(2.7)
In degrees $`<N+1`$, $`\delta _3`$ vanishes, and for $`nN+1`$, one has
$$(\mathrm{ker}(\delta _3))_n=\frac{(V^{(nN1)}W_{N+1})(I(R)_{nN}R)}{I(R)_{nN1}W_{N+1}}.$$
Clearly, $`(\mathrm{ker}(\delta _3))_n=0`$ if $`N+1n2N1`$, and
$$(\mathrm{ker}(\delta _3))_{2N}=(V^{(N1)}W_{N+1})(RR).$$
Lemma 2.2 (iii) implies that $`RR=(RV^N)(V^NR)`$. On the other hand, it is easy to check that (ec) implies that
$$(V^{(N1)}R)(RV^{(N1)})=W_{2N1}.$$
So, using Lemma 2.2 (ii), we get $`(\mathrm{ker}(\delta _3))_{2N}=W_{2N}`$. In other words, Tor$`{}_{4}{}^{A}(k_A,_Ak)`$ lives in degrees $`2N`$ and
$$(\mathrm{Tor}_4^A(k_A,_Ak))_{2N}=W_{2N}.$$
The next question would be to know when Tor$`{}_{4}{}^{A}(k_A,_Ak)`$ is concentrated in degree $`2N`$.
More generally, in our inductive process to define the Koszul property, the successive degrees of the Torโs are given by the values of the so-called *jump map* $`\zeta :`$, where
$$\zeta (2q)=qN,\zeta (2q+1)=qN+1,q\mathrm{integer}0.$$
Let $`n`$ be $`3`$. Assume that for any $`3in`$, Tor$`{}_{i}{}^{A}(k_A,_Ak)`$ is concentrated in degree $`\zeta (i)`$. In particular, (ec) holds. Then Tor$`{}_{n+1}{}^{A}(k_A,_Ak)`$ lives in degree $`\zeta (n+1)`$, and Tor$`{}_{n+1}{}^{A}(k_A,_Ak)`$ is concentrated in degree $`\zeta (n+1)`$ if and only if a certain sequence of distributivity relations hold. This sequence of distributivity relations is written down in Theorem 2.11 (iii) of , where the case of a field $`k`$ was treated (but the statement holds for any von Neumann regular ring thanks to Lemma 2.2). Anyway, the following definition becomes natural.
###### Definition 2.7
The connected $``$-graded $`k`$-$`k`$-algebra $`A`$ is said to be Koszul if for any $`n3`$, the graded left $`k`$-module Tor$`{}_{n}{}^{A}(k_A,_Ak)`$ is concentrated in degree $`\zeta (n)`$.
When $`A`$ is Koszul, the inductive process constructs a flat resolution in $`A`$-grMod of $`{}_{A}{}^{}k`$, which is called the *Koszul resolution*. Actually, independently of the Koszul property of $`A`$, a *Koszul complex* $`๐`$ in $`A`$-grMod is easily constructed : the objects are the $`A_kW_{\zeta (n)}`$, $`n0`$, and the arrow $`A_kW_{\zeta (n+1)}A_kW_{\zeta (n)}`$ is the extension by left $`A`$-linearity of the inclusion $`W_{\zeta (n+1)}A_kW_{\zeta (n)}`$. If $`A`$ is Koszul, $`๐`$ is the Koszul resolution. Conversely, if the complex $`๐`$ is exact in any (homological) degree $`>0`$, then $`๐`$ is a resolution of $`{}_{A}{}^{}k`$ via $`ฯต:Ak`$, and $`A`$ is Koszul since all arrows of the complex $`k_A_A๐`$ are vanishing.
We shall need a change of rings result. Let $`k^{}`$ be a ring and $`kk^{}`$ a ring morphism. Assume that the $`k`$-$`k`$-bimodule $`V`$ is a $`k^{}`$-$`k`$-bimodule whose left action extends left action by $`k`$ and right action is the same. Assume also that $`R`$ is a left sub-$`k^{}`$-module of $`V^N`$, so that $`A`$ is a $`k^{}`$-$`k`$-bimodule. Set $`V^{}=V_kk^{}`$. One has the sequence of $`k^{}`$-$`k^{}`$-bimodule morphisms
$$V^{}_k^{}V^{}(V_kk^{})_k^{}V^{}V_k(k^{}_k^{}V^{})V_k(V_kk^{})(V_kV)_kk^{}.$$
More generally, $`V^{_kn}_kk^{}V^{{}_{}{}^{}_{k^{}}^{}n}`$ for any $`n`$. As $`k`$ is von Neumann regular, $`V^{_kn}_kk^{}`$ is considered as included in $`T_k(V)_kk^{}`$, and $`T_k(V)_kk^{}`$ is the direct sum of the $`V^{_kn}_kk^{}`$, $`n0`$. So we have a natural isomorphism of $``$-graded $`k^{}`$-$`k^{}`$-algebras
$$T_k(V)_kk^{}T_k^{}(V^{})$$
which sends $`I(R)_kk^{}`$ (considered as included in $`T_k(V)_kk^{}`$) onto $`I(R^{})`$, where $`R^{}`$ is the image of $`R_kk^{}`$. The connected $``$-graded $`k^{}`$-$`k^{}`$-algebra
$$A^{}=\frac{T_k^{}(V^{})}{I(R^{})}$$
is then canonically isomorphic to $`A_kk^{}`$ since
$$\frac{T_k(V)_kk^{}}{I(R)_kk^{}}\frac{T_k(V)}{I(R)}_kk^{}.$$
###### Proposition 2.8
With above notations and assumptions, assume that $`k^{}`$ is von Neumann regular and that the connected $``$-graded $`k`$-$`k`$-algebra $`A`$ with $`N`$-homogeneous relations is Koszul. Then the connected $``$-graded $`k^{}`$-$`k^{}`$-algebra $`A^{}`$ has $`N`$-homogeneous relations and is Koszul.
Proof. The Koszul complex $`๐`$ of $`A`$ is exact in degree $`>0`$. Then the complex $`๐_kk^{}`$ is exact in degree $`>0`$, and its objects are naturally the $`A_kW_{\zeta (n)}^{}`$, $`n0`$, where $`W_{\zeta (n)}^{}W_{\zeta (n)}_kk^{}`$. But
$$A_kW_{\zeta (n)}^{}A_k(k^{}_k^{}W_{\zeta (n)}^{})(A_kk^{})_k^{}W_{\zeta (n)}^{}A^{}_k^{}W_{\zeta (n)}^{},$$
so that $`๐_kk^{}`$ is isomorphic to the Koszul complex of $`A^{}`$ (which makes sense since $`k^{}`$ is von Neumann regular). Thus $`A^{}`$ is Koszul.
## 3 PBW theorem
Our PBW theorem is the generalization to the $`N`$-case of a result (due to Braverman and Gaitsgory , see also Polishchuk and Positselski ) concerning non-homogeneous quadratic algebras over a field. It will be more convenient for us to follow along the lines of .
Throughout this section, $`k`$ is a von Neumann regular ring, $`V`$ is a $`k`$-$`k`$-bimodule, and $`T(V)`$ is the connected $``$-graded $`k`$-$`k`$-algebra introduced in the previous section. Actually, in this section, $`T(V)`$ will be seen as a connected $``$-*filtered* $`k`$-$`k`$-algebra: the ring $`T(V)`$ is filtered by the sub-$`k`$-$`k`$-bimodules
$$F^n=\underset{0in}{}V^i,n0,$$
and the left or right $`k`$-actions coincide with the products in the ring by elements of $`F^0=k`$.
Our data are now an integer $`N2`$ and a sub-$`k`$-$`k`$-bimodule $`P`$ of $`F^N`$. Then the two-sided ideal $`I(P)`$ is filtered by the $`k`$-$`k`$-bimodules
$$I(P)^n=I(P)F^n,n0.$$
We are interested in the $``$-filtered $`k`$-$`k`$-algebra $`U=T(V)/I(P)`$. The filtration of $`U`$ is formed by the $`k`$-$`k`$-bimodules $`U^n=F^n/I(P)^n`$. Warning: $`U^0`$ can vanish (if $`P=k`$ for example), but the PBW property will avoid this case.
Although $`I(P)=_{i,j0}V^iPV^j`$, $`I(P)^n`$ may contain *strictly* the sum $`_{i+N+jn}V^iPV^j`$ in a non-trivial way. The following example is well-known.
Example 3.1 Take $`n=N=2`$, $`k`$ a field of characteristic 0, and $`V`$ finite-dimensional. Let $`f:V\times VV`$ be an alternate bilinear map, and let $`P`$ be the subspace of $`F^2`$ spanned by all elements
$$r_{xy}=xyyxf(x,y),x,yV.$$
Set $`c_{xyz}=[r_{xy},z]+[r_{yz},x]+[r_{zx},y]`$, where $`[,]`$ is the usual commutator in $`T(V)`$. Clearly, $`c_{xyz}PV+VP`$. Jacobi identity for $`[,]`$ shows that $`c_{xyz}`$ belongs to $`V^2`$, hence to $`I(P)^2`$. On the other hand,
$$c_{xyz}=r_{f(x,y)z}r_{f(y,z)x}r_{f(z,x)y}J(x,y,z),$$
where $`J(x,y,z)=f(f(x,y),z)+f(f(y,z),x)+f(f(z,x),y)V`$. Since $`PV=0`$, $`c_{xyz}`$ belongs to $`P`$ if and only if $`J(x,y,z)=0`$. Thus, if $`(V,f)`$ is not a Lie algebra (i.e., if $`f`$ does not satisfy the Jacobi identity), $`I(P)^2`$ contains strictly $`P`$.
The PBW theorem gives a โsimpleโ (as the Jacobi identity) sufficient condition in order to have
$$I(P)^n=\underset{i+N+jn}{}V^iPV^j,\mathrm{for}\mathrm{any}n0.$$
(3.1)
By convention, the right-hand side of (3.1) vanishes when $`n<N`$. So if (3.1) holds, $`I(P)^n=0`$ (and $`U^n=F^n`$) when $`n<N`$. Our aim is now to express more conceptually (3.1).
Introduce the $``$-graded $`k`$-$`k`$-algebra $`gr(U)`$ associated to $`U`$, and denote by $`(gr(U)_n)`$ its gradation. One has $`gr(U)_n=U^n/U^{n1}`$. Using $`I(P)^{n1}=I(P)^nF^{n1}`$ and canonical isomorphisms, we identify $`gr(U)_n`$ to the $`k`$-$`k`$-bimodule $`F^n/(I(P)^n+F^{n1})`$. The product of the algebra $`gr(U)`$ is natural in these identifications. On the other hand, let $`\pi :F^NV^N`$ be the projection associated to $`F^N=V^NF^{N1}`$. Then $`R=\pi (P)`$ is a sub-$`k`$-$`k`$-bimodule of $`V^N`$, and we consider, as in the previous section, the connected $``$-graded $`k`$-$`k`$-algebra $`A=T(V)/I(R)`$.
For each $`n0`$, let $`\varphi _n:V^ngr(U)_n`$ be the composite of canonical maps
$$V^nF^nF^n/(I(P)^n+F^{n1}).$$
Since $`P+F^{N1}=RF^{N1}`$, one has $`\varphi _N(R)=0`$. Thus the surjective algebra morphism $`\varphi =_{n0}\varphi _n:T(V)gr(U)`$ defines a surjective morphism of $``$-graded $`k`$-$`k`$-algebras
$$p:Agr(U).$$
Using the canonical isomorphism
$$A_n=\frac{V^n}{I(R)_n}\frac{V^nF^{n1}}{I(R)_nF^{n1}},$$
one identifies $`p_n:A_ngr(U)_n`$ to the canonical map
$$\frac{F^n}{I(R)_n+F^{n1}}\frac{F^n}{I(P)^n+F^{n1}}.$$
Thus $`p`$ is an isomorphism if and only if $`I(P)^nI(R)_n+F^{n1}`$ for any $`n0`$. The condition โ$`p`$ is an isomorphismโ is called the *PBW property* (for $`U`$).
###### Proposition 3.2
The PBW property is equivalent to (3.1) (actually, (3.1) for $`nN1`$ suffices since $`I(P)^nI(P)^{N1}`$ when $`0nN2`$).
Proof. Introduce the notations
$$J^n=\underset{i+N+jn}{}V^iPV^j,n0,$$
(3.2)
with convention $`J^n=0`$ if $`n<N`$. One has $`J^nI(P)^n`$. Assume firstly that $`I(P)^n=J^n`$ for any $`n`$. From that and $`PR+F^{N1}`$, one draws
$$I(P)^n\underset{i+N+j=n}{}V^iRV^j+F^{n1},$$
thus $`p`$ is an isomorphism.
Conversely, assume that $`p`$ is an isomorphism. Prove equalities (3.1) by induction on $`n0`$. Since $`I(P)^0I(R)_0+F^1`$, one has $`I(P)^0=0`$, hence (3.1) for $`n=0`$. Let $`n`$ be $`1`$, and assume (3.1) for $`n1`$. Since $`I(P)^nI(R)_n+F^{n1}`$ and $`RP+F^{N1}`$, one has
$$I(P)^n\underset{i+N+j=n}{}V^iPV^j+F^{n1},$$
where the sum is supposed to vanish if $`n<N`$ (because $`I(R)_n=0`$ in this case). Let $`a`$ be in $`I(P)^n`$, and write down $`a=b+c`$, $`b_{i+N+j=n}V^iPV^j`$, $`cF^{n1}`$. Clearly, $`c=ab`$ belongs to $`I(P)F^{n1}=I(P)^{n1}`$, hence to $`J^{n1}J^n`$ by induction hypothesis. But $`bJ^n`$, thus $`aJ^n`$.
###### Proposition 3.3
Keeping notations (3.2), the PBW property is equivalent to
$$J^nF^{n1}=J^{n1},\mathrm{for}\mathrm{any}nN.$$
(3.3)
Proof. If $`p`$ is an isomorphism, (3.1) shows that $`J^nF^{n1}=I(P)^nF^{n1}=I(P)^{n1}=J^{n1}`$. Conversely, assume (3.3). Fix $`nN1`$. Since $`I(P)=_{i,j0}V^iPV^j`$, one has
$$I(P)^n=(\underset{iN}{}J^i)F^n=\underset{iN}{}(J^iF^n).$$
For $`Nin`$, $`J^iJ^n`$ hence $`J^iF^nJ^n`$. Next $`J^{n+1}F^n=J^n`$ because of (3.3) and $`n+1N`$. Moreover $`J^{n+2}F^n=J^{n+2}F^{n+1}F^n=J^{n+1}F^n=J^n`$. A straightforward induction shows that $`J^iF^n=J^n`$ for any $`in+1`$. Thus $`I(P)^nJ^n`$, and we conclude by the previous proposition.
For $`n=N`$ and $`n=N+1`$, (3.3) is respectively equivalent to
$$PF^{N1}=0,$$
(3.4)
$$(PV+VP)F^NP.$$
(3.5)
(For $`n=N+1`$, one can replace $`P+PV+VP`$ by $`PV+VP`$ in the non-trivial inclusion of (3.3).)
When $`N=2`$ and $`k`$ is a field, (3.4) and (3.5) are the conditions (I) and (J) of Braverman and Gaitsgory. Their PBW theorem asserts that if $`A`$ is Koszul, (I) and (J) are sufficient in order to have the PBW property. This theorem extends to any $`N2`$ and any von Neumann regular ring $`k`$, as follows.
###### Theorem 3.4
*(PBW theorem in the $`N`$-case)* Assume that $`k`$ is a von Neumann regular ring, $`V`$ is a $`k`$-$`k`$-bimodule, $`N`$ is an integer $`2`$, and $`P`$ is a sub-$`k`$-$`k`$-bimodule of $`F^N`$, where $`F^n=_{0in}V^i`$ for any $`n0`$. Set $`U=T(V)/I(P)`$ and $`A=T(V)/I(R)`$, where $`R=\pi (P)`$ and $`\pi `$ is the projection of $`F^N`$ onto $`V^N`$ modulo $`F^{N1}`$.
Assume that the graded left $`k`$-module Tor$`{}_{3}{}^{A}(k_A,_Ak)`$ is concentrated in degree $`N+1`$. Then (3.4) and (3.5) imply that the PBW property holds, i.e., the natural algebra morphism $`p:Agr(U)`$ is an isomorphism.
Before proving the theorem, let us introduce the arrows $`\phi ^{i,i+N1}`$ which are a generalization of the arrows $`\psi ^{i,i+1}`$ of . From now on, we assume (3.4). As $`\mathrm{ker}(\pi )=F^{N1}`$, $`\pi `$ realizes an isomorphism from $`P`$ onto $`R`$. For any $`xR`$, let $`y`$ be the element of $`P`$ such that $`\pi (y)=x`$, and set $`y=x\phi (x)`$. That defines a $`k`$-$`k`$-linear map $`\phi :RF^{N1}`$. Then $`P`$ can be described only by $`R`$ and $`\phi `$:
$$P=\{x\phi (x);xR\}.$$
(3.6)
Fix $`i1`$, $`j0`$. Using the identifications introduced in Section 2 just before Lemma 2.2, the $`k`$-$`k`$-linear map
$$1_{V^{(i1)}}_k\phi _k1_{V^j}:V^{(i1)}_kR_kV^jV^{(i1)}_kF^{N1}_kV^j$$
is identified to a $`k`$-$`k`$-linear map $`V^{(i1)}RV^jT(V)`$ which is denoted by $`\phi ^{i,i+N1}`$ (actually it arrives in $`V^{(i1)}F^{N1}V^j`$ which is naturally embedded in $`T(V)`$).
###### Proposition 3.5
Assume that (3.4) holds. Then (3.5) is equivalent to
$$(\phi ^{1,N}\phi ^{2,N+1})(W_{N+1})P.$$
(3.7)
Proof. Let $`a`$ in $`W_{N+1}`$. Then $`a\phi ^{1,N}(a)PV`$ and $`a\phi ^{2,N+1}(a)VP`$, so $`(\phi ^{1,N}\phi ^{2,N+1})(a)`$ belongs to $`(PV+VP)F^N`$. Conversely, if $`x`$ belongs to $`(PV+VP)F^N`$, decompose
$$x=\underset{i}{}(x_i\phi (x_i))v_i+\underset{j}{}v_j^{}(x_j^{}\phi (x_j^{})),$$
with $`x_i`$, $`x_j^{}`$ in $`R`$, $`v_i`$, $`v_j^{}`$ in $`V`$. Since $`xF^N`$, $`_ix_iv_i+_jv_j^{}x_j^{}=0`$, hence
$$x=(\phi ^{1,N}\phi ^{2,N+1})(\underset{i}{}x_iv_i),$$
where $`_ix_iv_iW_{N+1}`$. Thus $`(\phi ^{1,N}\phi ^{2,N+1})(W_{N+1})=(PV+VP)F^N`$.
*Proof of Theorem 3.4.* According to Proposition 3.3, it suffices to prove that $`J^nF^{n1}J^{n1}`$ by induction on $`nN`$. Assume $`nN+2`$ and the inclusion true for $`n1`$. Let $`x`$ be in $`J^nF^{n1}`$. In the sequel, the symbol $``$ means equality modulo $`J^{n1}`$. One can find $`x_i`$ in $`V^{(i1)}RV^{(nNi+1)}`$ for $`1inN+1`$ such that
$$x\underset{i=1}{\overset{nN+1}{}}(x_i\phi ^{i,i+N1}(x_i)).$$
Since $`xF^{n1}`$, the sum of the $`x_i`$โs vanishes, so
$$x_{nN+1}(V^{(nN)}R)(I(R)_{n1}V).$$
But Tor$`{}_{3}{}^{A}(k_A,_Ak)`$ is concentrated in degree $`N+1`$. Theorem 2.6 shows that there exist $`y_i`$ in $`V^{(i1)}RV^{(n2Ni+1)}R`$ for $`1in2N+1`$ and $`y_{nN}`$ in $`V^{(nN1)}W_{N+1}`$ such that
$$x_{nN+1}=\underset{i=1}{\overset{n2N+1}{}}y_i+y_{nN}.$$
Note that if $`n<2N`$, the latter equality reduces to $`x_{nN+1}=y_{nN}`$, so that the other $`y_i`$โs are considered as vanishing in this case (actually, (ec) and (2.6) can be put together in the statement of Theorem 2.6). Using (3.7), we get
$$\phi ^{nN+1,n}(x_{nN+1})\underset{i=1}{\overset{n2N+1}{}}\phi ^{nN+1,n}(y_i)+\phi ^{nN,n1}(y_{nN}).$$
Fix the index $`i`$, $`1in2N+1`$. Decompose
$$\phi ^{i,i+N1}:V^{(i1)}RV^{(n2Ni+1)}RV^{(i1)}F^{N1}V^{(n2Ni+1)}R$$
as $`\phi ^{i,i+N1}=_{j=0}^{N1}\phi _j^{i,i+N1}`$, where
$$\phi _j^{i,i+N1}:V^{(i1)}RV^{(n2Ni+1)}RV^{(n2N+j)}R.$$
But the two $`R`$โs in the tensor product $`E=V^{(i1)}RV^{(n2Ni+1)}R`$ do not overlap. Thus the equality
$$\phi ^{i,i+N1}\phi ^{nN+1,n}=\underset{j=0}{\overset{N1}{}}\phi ^{n2N+j+1,nN+j}\phi _j^{i,i+N1}$$
holds *on E*, so that *on E*, one has
$$\phi ^{i,i+N1}\phi ^{nN+1,n}=\underset{j=0}{\overset{N1}{}}(\mathrm{id}\phi ^{n2N+j+1,nN+j})\phi _j^{i,i+N1}(\mathrm{id}\phi ^{i,i+N1})\phi ^{nN+1,n}$$
and we get $`\phi ^{nN+1,n}(y_i)\phi ^{i,i+N1}(y_i)`$. Finally
$$x\underset{i=1}{\overset{n2N+1}{}}\phi ^{i,i+N1}(x_i+y_i)\underset{i=n2N+2}{\overset{nN1}{}}\phi ^{i,i+N1}(x_i)\phi ^{nN,n1}(x_{nN}+y_{nN})$$
in which the second sum does not occur if $`N=2`$.
Set $`x_i^{}=x_i+y_i`$ for $`1in2N+1`$ and $`i=nN`$. Set $`x_i^{}=x_i`$ for $`n2N+2inN1`$. Then for $`1inN`$, $`x_i^{}`$ belongs to $`V^{(i1)}RV^{(nNi)}V`$, and $`_{i=1}^{nN}x_i^{}=0`$. Introducing
$$x^{}=\underset{i=1}{\overset{nN}{}}(x_i^{}\phi ^{i,i+N1}(x_i^{})),$$
we get $`xx^{}`$ and $`x^{}(J^{n1}V)(F^{n2}V)`$. Lemma 2.2 gives $`(J^{n1}V)(F^{n2}V)=(J^{n1}F^{n2})V`$, therefore $`x^{}`$ belongs to $`J^{n2}V`$ by induction hypothesis. Thus $`x^{}`$ and $`x`$ are in $`J^{n1}`$.
###### Proposition 3.6
Assume that (3.4) holds. Let $`\phi :RF^{N1}`$ be the $`k`$-$`k`$-linear map such that $`P=\{x\phi (x);xR\}`$. Decompose $`\phi =_{j=0}^{N1}\phi _j`$, $`\phi _j:RV^j`$. Then (3.5) is equivalent to all following relations
$$(\phi _{N1}^{1,N}\phi _{N1}^{2,N+1})(W_{N+1})R,$$
(3.8)
$$\left(\phi _j(\phi _{N1}^{1,N}\phi _{N1}^{2,N+1})+\phi _{j1}^{1,N}\phi _{j1}^{2,N+1}\right)(W_{N+1})=0,1jN1,$$
(3.9)
$$\phi _0(\phi _{N1}^{1,N}\phi _{N1}^{2,N+1})(W_{N+1})=0.$$
(3.10)
Proof. Apply Proposition 3.5. Let $`x`$ be in $`W_{N+1}`$ and $`X=(\phi ^{1,N}\phi ^{2,N+1})(x)`$. Since the two projections in $`PRF^{N1}`$ are respectively $`\pi `$ and $`\phi \pi `$, then $`X`$ belongs to $`P`$ if and only if $`\pi (X)R`$ and $`X=\pi (X)\phi \pi (X)`$. The component of degree $`j`$, $`0jN`$, of $`X`$ is $`(\phi _{j1}^{1,N}\phi _{j1}^{2,N+1})(x)`$ (vanishing if $`j=0`$). In particular
$$\pi (X)=(\phi _{N1}^{1,N}\phi _{N1}^{2,N+1})(x).$$
And the component of degree $`jN1`$ of $`\phi \pi (X)`$ is $`\phi _j(\phi _{N1}^{1,N}\phi _{N1}^{2,N+1})(x)`$.
Example 3.7 In Example 3.1, (3.4) holds and $`\phi _0=0`$, so (3.10) holds. As elements of $`W_3`$ are totally skew-symmetric, (3.8) is easily checked, and (3.9) is equivalent to the Jacobi identity for $`f`$.
Example 3.8 (Down-up algebras ) Here $`k`$ is a field, $`\alpha `$, $`\beta `$, $`\gamma `$ are in $`k`$ with $`\beta 0`$, $`U=U(\alpha ,\beta ,\gamma )`$ is the associative $`k`$-algebra with two generators $`d`$, $`u`$ and following relations
$$d^2u=\alpha dud+\beta ud^2+\gamma d,$$
$$du^2=\alpha udu+\beta u^2d+\gamma u.$$
Taking $`d`$ and $`u`$ of degree 1, $`U`$ is a non-homogeneous cubic algebra, and the homogeneous cubic algebra $`A`$ is defined by relations $`r_1=r_2=0`$, where $`r_1=d^2u\alpha dud\beta ud^2`$ and $`r_2=du^2\alpha udu\beta u^2d`$. Then $`A`$ is the AS-regular algebra of global dimension 3 which is cubic of type S<sub>1</sub> . In particular, $`A`$ is Koszul . Moreover , $`W_4=RVVR`$ is one-dimensional and a generator is
$$w=r_1u\beta r_2d=\beta ur_1+dr_2.$$
Let us show the PBW property. Firstly, $`PF^2=0`$ is clear. Use Proposition 3.6. One has $`\phi _0=\phi _2=0`$ and $`\phi _1`$ is defined by $`\phi _1(r_1)=\gamma d`$, $`\phi _1(r_2)=\gamma u`$. Thus (3.8) and (3.10) hold, whereas (3.9) comes from the calculation
$$(\phi _1^{1,3}\phi _1^{2,4})(w)=\gamma du\beta \gamma ud(\beta u\gamma d+d\gamma u)=0.$$
The fact that $`p:Agr(U)`$ is an isomorphism allows us to deduce properties of $`U`$ from those of $`A`$. For example, using , $`U(\alpha ,\beta ,\gamma )`$ is a noetherian domain of Gelfand-Kirillov dimension 3 (see also ).
###### Definition 3.9
Assume that $`k`$ is a von Neumann regular ring, $`V`$ is a $`k`$-$`k`$-bimodule, $`N`$ is an integer $`2`$, and $`P`$ is a sub-$`k`$-$`k`$-bimodule of $`F^N`$, where $`F^n=_{0in}V^i`$ for any $`n0`$. Set $`U=T(V)/I(P)`$ and $`A=T(V)/I(R)`$, where $`R=\pi (P)`$ and $`\pi `$ is the projection of $`F^N`$ onto $`V^N`$ modulo $`F^{N1}`$. Then $`U`$ is said to be Koszul if the graded algebra (with $`N`$-homogeneous relations) $`A`$ is Koszul and if the PBW property holds.
Note that if $`U`$ is Koszul, the integer $`N`$ is uniquely determined, and we shall say that $`U`$ is $`N`$-Koszul. Definition 3.9 generalizes the definition of the graded situation: if the relations of $`U`$ are all $`N`$-homogeneous, then $`U=A`$ is graded and the PBW property holds trivially.
Remark 3.10 Let us show how $`\phi _0`$ measures the obstruction for $`k`$ to be an $`U`$-module. Assume firstly that $`\phi _0=0`$. It is easy to check that the natural projection $`q:T(V)k`$ vanishes on the ideal $`I(P)`$, so that we get a natural $`k`$-$`k`$-linear map $`ฯต_U:Uk`$. Then the left $`U`$-module $`k`$ deduced from $`ฯต_U`$ satisfies the following properties :
(1) $`u.1=0`$ for any $`uV^n`$ and $`1nN1`$,
(2) $`\overline{x}.1=0`$ for any $`xR`$, where $`\overline{x}`$ denotes the class of $`x`$ in $`U^N`$.
Conversely, assume that $`k`$ is a $`U`$-module satisfying (1) and (2). Then it is immediate that $`\phi _0=0`$.
## 4 $`N`$-inhomogeneous algebras associated to finite groups
Throughout this section, $`k`$ is an algebraically closed field of characteristic 0, $`V`$ a $`k`$-vector space of dimension $`n`$, and $`\mathrm{\Gamma }`$ is a finite subgroup of $`GL(V)`$. The group algebra $`K=k[\mathrm{\Gamma }]`$ is a semi-simple ring, hence is von Neumann regular. Fix $`p`$, $`1<pn`$.
Given a $`k`$-linear map $`\psi :\mathrm{\Lambda }_k^p(V)K`$, put
$$H_\psi =(T_k(V)\mathrm{\#}\mathrm{\Gamma })/I(\mathrm{Alt}(v_1,\mathrm{},v_p)\psi (v_1,\mathrm{},v_p);v_1,\mathrm{},v_pV),$$
where Alt stands for the anti-symmetrization in $`T_k^p(V)`$. Consider the $`K`$-$`K`$-bimodule $`E=V_kK`$, with left $`\mathrm{\Gamma }`$-action given by $`g:vag(v)(ga)`$, and right $`\mathrm{\Gamma }`$-action given by $`vav(ag)`$, where $`ga`$ and $`ag`$ stand for the product in the group algebra. Then $`T_K(E)T_k(V)\mathrm{\#}\mathrm{\Gamma }`$, so that $`H_\psi `$ is identified to the $``$-filtered $`K`$-$`K`$-algebra $`T_K(E)/I(P)`$ where $`P`$ is the sub-$`K`$-$`K`$-bimodule of $`T_K(E)`$ generated by the following elements of $`T_K^p(E)K`$:
$$\mathrm{Alt}(v_1,\mathrm{},v_p)\psi (v_1,\mathrm{},v_p)\mathrm{with}v_1,\mathrm{},v_pV.$$
So $`H_\psi `$ has relations of degree $`N=p`$. Using notations of the previous section, $`U=H_\psi `$ and $`A=H_{\psi =0}`$. Clearly, $`A`$ is obtained from the $`N`$-homogeneous $`k`$-$`k`$-algebra
$$๐=T_k(V)/I(\mathrm{Alt}(v_1,\mathrm{},v_p);v_1,\mathrm{},v_pV)$$
by the natural change of rings $`kK`$. It is known that $`๐`$ is Koszul . Thus Proposition 2.8 shows that $`A`$ is Koszul.
###### Lemma 4.1
In the notation of (3.8)-(3.10), in $`T_K^{p+1}E`$, we have $`W_{p+1}=(_k^{p+1}V)_kK`$ as a $`K`$-subbimodule in $`T_k^{p+1}V_kK`$.
Proof. According to , the algebra $`๐`$ is Koszul. Furthermore, if $``$ denotes the space of relations of $`๐`$, then the space of relations of $`A`$ equals $`R=_kK`$. Hence we get $`W_q=_{i+p+j=q}E^iRE^j=(_k^qV)_kK,`$ as a $`K`$-subbimodule in $`T_k^qV_kK.`$ $`\mathrm{}`$
Now, we may write the map $`\psi :\mathrm{\Lambda }_k^p(V)K`$ in the form $`\psi =_{g\mathrm{\Gamma }}\psi _gg,`$ where $`\psi _g:\mathrm{\Lambda }_k^p(V)k`$ are certain linear maps. Further, let $`\mathrm{\Gamma }`$ act on $`K`$ by conjugation.
###### Lemma 4.2
The algebra $`H_\psi `$ is Koszul if and only if $`\psi :\mathrm{\Lambda }_k^p(V)K`$ is $`\mathrm{\Gamma }`$-equivariant and, for any $`g\mathrm{\Gamma }`$ and $`v_1,\mathrm{},v_{p+1}V,`$ in $`V`$ one has
$$\underset{i=1}{\overset{p+1}{}}(1)^i\psi _g(v_1,\mathrm{},v_{i1},v_{i+1},\mathrm{},v_{p+1})(\mathrm{Id}(1)^pg)(v_i)=0.$$
(4.1)
Proof. First of all, we claim that if the PBW-property holds then the map $`\psi `$ must be $`\mathrm{\Gamma }`$-equivariant. To see this, let $`e_\rho k[\mathrm{\Gamma }]`$ denote the central idempotent corresponding to an irreducible representation $`\rho `$ of $`\mathrm{\Gamma }`$. Let $`\mathrm{\Gamma }`$ act on $`K`$ via the adjoint action and on the vector space $`\mathrm{Hom}_k(T_k^pV,T_k^pV_kK)`$ via the corresponding induced action. The map $`\mathrm{Alt}:v_1,\mathrm{},v_p\mathrm{Alt}(v_1,\mathrm{},v_p)1`$ is clearly $`\mathrm{\Gamma }`$-equivariant, hence, in $`\mathrm{Hom}_k(T_k^pV,T_k^pV_kK)`$, we have $`e_\rho (\mathrm{Alt})=0`$, for every nontrivial irreducible representation $`\rho `$ of $`\mathrm{\Gamma }`$. Now, if the map $`\psi `$ is not $`\mathrm{\Gamma }`$-equivariant, then there exists a nontrivial irreducible representation $`\rho `$ such that $`e_\rho (\psi )0`$. We conclude that there exist $`v_1,\mathrm{},v_pV`$ such that $`\left(e_\rho (\mathrm{Alt}\psi )\right)(v_1,\mathrm{},v_p)=\left(e_\rho (\psi )\right)(v_1,\mathrm{},v_p)0`$. This means that in the algebra $`H_\psi `$ we have a relation $`a=0`$, where $`a:=\left(e_\rho (\psi )\right)(v_1,\mathrm{},v_p)K`$ is a nonzero element. Thus, the canonical map $`KH_\psi `$ is not injective and PBW-property fails. Actually, denoting by $`(F^n)_{n0}`$ the natural filtration of $`T_K(E)`$, it is elementary that $`PF^{p1}=0`$ (i.e., the condition (3.4)) holds if and only if $`\psi :\mathrm{\Lambda }_k^p(V)K`$ is $`\mathrm{\Gamma }`$-equivariant.
The rest of the argument is very similar to the proof of \[12, formula (2.3)\]. Assume that $`\psi `$ is $`\mathrm{\Gamma }`$-equivariant. Since $`\phi =\phi _0=\psi `$, the criterion of Proposition 3.6 (i.e., the condition (3.5)) reduces to the following equation in $`V`$:
$$(\psi ^{1,p}\psi ^{2,p+1})(\mathrm{Alt}(v_1,\mathrm{},v_p))=0,v_1,\mathrm{},v_{p+1}V.$$
Explicitly, writing $`ฯต(\sigma )`$ for the sign of permutation $`\sigma `$, the equation reads
$$\underset{\sigma ๐_{p+1}}{}ฯต(\sigma )\left[\psi (v_{\sigma (1)},\mathrm{},v_{\sigma (p)})v_{\sigma (p+1)}v_{\sigma (1)}\psi (v_{\sigma (2)},\mathrm{},v_{\sigma (p+1)})\right]=0.$$
Rewriting this expression one obtains the following condition:
$`0`$ $`={\displaystyle \underset{i=1}{\overset{p+1}{}}}(1)^{p+1i}{\displaystyle \underset{\tau ๐_p}{}}ฯต(\tau )\psi (v_{\tau (1)},\mathrm{},v_{\tau (i1)},v_{\tau (i+1)},\mathrm{},v_{\tau (p+1)})v_i`$
$`{\displaystyle \underset{i=1}{\overset{p+1}{}}}(1)^{i1}{\displaystyle \underset{\tau ๐_p}{}}ฯต(\tau )v_i\psi (v_{\tau (1)},\mathrm{},v_{\tau (i1)},v_{\tau (i+1)},\mathrm{},v_{\tau (p+1)})`$
$`={\displaystyle \underset{i=1}{\overset{p+1}{}}}(1)^i[\psi (v_{\tau (1)},\mathrm{},v_{\tau (i1)},v_{\tau (i+1)},\mathrm{},v_{\tau (p+1)}),v_i]_\pm ,`$
where we use the notation $`[a,v]_\pm =av(1)^pva,`$ for any $`aK,vV`$.
Further, for any $`vV`$ and $`g\mathrm{\Gamma }`$, in $`T_k(V)\mathrm{\#}\mathrm{\Gamma }`$, we have
$$[v,g]_\pm =(v1)g(1)^pg(v1)=vg(1)^pg(v)g=(v(1)^pg(v))g.$$
Therefore, writing $`\psi =_{g\mathrm{\Gamma }}\psi _gg,`$ the last displayed formula reads
$`0`$ $`={\displaystyle \underset{i=1}{\overset{p+1}{}}}(1)^i\left({\displaystyle _{g\mathrm{\Gamma }}}\psi _g(v_1,\mathrm{},v_{i1},v_{i+1},\mathrm{},v_{p+1})(v_i(1)^pg(v_i))g\right)`$
$`={\displaystyle \underset{g\mathrm{\Gamma }}{}}\left({\displaystyle \underset{i=1}{\overset{p+1}{}}}(1)^i\psi _g(v_1,\mathrm{},v_{i1},v_{i+1},\mathrm{},v_{p+1})(\mathrm{Id}(1)^pg)(v_i)\right)g.`$
Thus, the coefficient in front of each element $`g`$ in the second line of the formula must vanish, and the Lemma follows. $`\mathrm{}`$
To simplify notation, we write $`=_k`$ and $`\mathrm{\Lambda }^i()=\mathrm{\Lambda }_k^i(),`$ etc.
Next, fix $`g\mathrm{\Gamma }`$. Write $`M=M_g:=\mathrm{Image}(\mathrm{Id}(1)^pg),`$ and $`L=L_g:=\mathrm{Ker}(\mathrm{Id}(1)^pg),`$ and put $`m:=dimM,l:=dimL`$. Since $`g`$ is an element of $`GL(V)`$ of finite order, hence, semisimple, we have a direct sum decomposition $`V=ML`$. Thus, we have $`m+l=n`$ and we have a natural isomorphism
$$\mathrm{\Lambda }^p(V)=\mathrm{\Lambda }^p(ML)_{i+j=p}\mathrm{\Lambda }^i(M)\mathrm{\Lambda }^j(L).$$
(4.2)
Note also that the space $`\mathrm{\Lambda }^m(M)`$ is 1-dimensional.
###### Lemma 4.3
A linear map $`\varphi :\mathrm{\Lambda }^p(V)=\mathrm{\Lambda }^p(ML)k`$ satisfies the equation
$$\underset{i=1}{\overset{p+1}{}}(1)^i\varphi (v_1,\mathrm{},v_{i1},v_{i+1},\mathrm{},v_{p+1})(\mathrm{Id}(1)^pg)(v_i)=0,v_1,\mathrm{},v_{p+1}V.$$
if and only if $`\varphi =\varphi ^{m,pm}`$, where $`\varphi ^{m,pm}`$ is the restriction of $`\varphi `$ to $`\mathrm{\Lambda }^m(M)\mathrm{\Lambda }^{pm}(L)`$. In particular, if $`\varphi `$ does not vanish, then $`pm`$.
Proof. Recall that, for any $`d`$-dimensional vector space $`E`$ and an integer $`0pd`$, one has a Koszul exact sequence (a version of Koszul complex):
$`K^p(E):\mathrm{\hspace{0.33em}0}\mathrm{\Lambda }^p(E^{})\stackrel{_E^p}{}`$ $`\mathrm{\Lambda }^{p+1}(E^{})E\mathrm{\Lambda }^{p+2}(E^{})S^2(E)`$ (4.3)
$`\mathrm{}\mathrm{\Lambda }^d(E^{})S^{dp}(E)k0.`$
The first differential $`_E^p:\mathrm{\Lambda }^p(E^{})\mathrm{\Lambda }^{p+1}(E^{})E=\mathrm{Hom}(\mathrm{\Lambda }^{p+1}E^{},E)`$ is defined by the formula
$$(_E^p\alpha )(v_1,\mathrm{},v_{p+1}):=\underset{i=1}{\overset{p+1}{}}(1)^i\alpha (v_1,\mathrm{},v_{i1},v_{i+1},\mathrm{},v_{p+1})v_i,$$
for any $`v_1,\mathrm{},v_{p+1}E.`$ The other differentials are given by a similar formula. Thus, from (4.3) we see that
The map $`_E^p`$ is injective unless $`p=dimE`$. (4.4)
The Koszul exact sequence (4.3) is compatible with tensor products. In particular, for $`V=ML`$ we have $`_V^p=_M\mathrm{Id}_L\pm \mathrm{Id}_M_L`$. In more detail, using the direct sum decomposition (4.2), the differential
$$_V^p:\mathrm{\Lambda }^p(M^{}L^{})\mathrm{\Lambda }^{p+1}(M^{}L^{})(ML)$$
may be written as a sum $`_V^p=_{r+s=p}_V^{r,s}`$, of the following components
$$_V^{r,s}:\mathrm{\Lambda }^r(M^{})\mathrm{\Lambda }^s(L^{})_{i+j=p+1}\mathrm{\Lambda }^i(M)\mathrm{\Lambda }^j(L)(ML).$$
With this notation, for any $`r,s`$, one has a Leibniz type formula:
$$_V^{r,s}=_M^r\mathrm{Id}_L^s+(1)^r\mathrm{Id}_M^r_L^s,$$
(4.5)
where the two summunds are interpreted as maps
$$_M^r\mathrm{Id}_L^s:\mathrm{\Lambda }^r(M^{})\mathrm{\Lambda }^s(L^{})\mathrm{\Lambda }^{r+1}(M^{})\mathrm{\Lambda }^s(L^{})M,$$
resp.,
$$\mathrm{Id}_M^r_L^s:\mathrm{\Lambda }^r(M^{})\mathrm{\Lambda }^s(L^{})\mathrm{\Lambda }^r(M^{})\mathrm{\Lambda }^{s+1}(L^{})L.$$
The target spaces in these formulas are viewed as being subspaces of the vector space $`_{i+j=p+1}\mathrm{\Lambda }^i(M)\mathrm{\Lambda }^j(L)(ML).`$
We can now complete the proof of the Lemma. Observe that, in terms of differential $`_V^p`$, equation in the statement of the Lemma reads $`(\mathrm{Id}(1)^pg)(_V^p(\varphi ))=0`$. Now, the map $`(\mathrm{Id}(1)^pg)`$ takes the subspace $`LLM=V`$ to zero, and restricts to an invertible operator $`MM`$. Therefore, writing
$$\pi _M:_{i+j=p+1}\mathrm{\Lambda }^i(M)\mathrm{\Lambda }^j(L)(ML)_{i+j=p+1}\mathrm{\Lambda }^i(M)\mathrm{\Lambda }^j(L)M$$
for the natural projection, we obtain $`(\mathrm{Id}(1)^pg)(_V^p(\varphi ))=0\pi _M(_V^p(\varphi ))=0`$. Further, it is clear that we have
$$\pi _M(_M^r\mathrm{Id}_L^s)=(_M^r\mathrm{Id}_L^s),\text{and}\pi _M(\mathrm{Id}_M^r_L^s)=0.$$
Thus, using formula (4.5), we conclude that
$$(\mathrm{Id}(1)^pg)(_V^p(\varphi ))=0(_M^r\mathrm{Id}_L^{pr})(\varphi ^{r,pr})=0,0rp.$$
(4.6)
Next, observe that the map $`_M^r`$, hence the map $`(_M^r\mathrm{Id}_L^{pr})`$, is injective for all $`r<dimM`$, by (4.4). Therefore, the vanishing condition on the right of (4.6) holds if and only if $`r=dimM=m`$. Furthermore, in the latter case, the space $`\mathrm{\Lambda }^rM`$ is 1-dimensional. It follows that the map $`\varphi ^{m,pm}`$ may be written as a tensor product $`\varphi _M\varphi _L`$ of linear maps $`\varphi _M:\mathrm{\Lambda }^mMk`$ and $`\varphi _L:\mathrm{\Lambda }^{pm}Lk`$. The Lemma is proved. $`\mathrm{}`$
Now, for any $`g\mathrm{\Gamma }`$ we put $`a(g):=dimM_g`$. Recall that $`M_g=\mathrm{Im}(\mathrm{Id}(1)^pg)`$ and $`L_g=\mathrm{Ker}(\mathrm{Id}(1)^pg)`$. Lemmas 4.2 and 4.3 yield the following result.
###### Theorem 4.4
The $``$-filtered $`K`$-$`K`$-algebra $`H_\psi `$ is Koszul if and only if $`\psi =_{g\mathrm{\Gamma }}\psi _gg,`$ : $`\mathrm{\Lambda }_k^p(V)K`$ is $`\mathrm{\Gamma }`$-equivariant and for any $`g\mathrm{\Gamma }`$, $`\psi _g`$ vanishes in all components $`\mathrm{\Lambda }^i(M_g)\mathrm{\Lambda }^{pi}(L_g)`$ of $`\mathrm{\Lambda }^p(V)`$ unless if $`i=a(g)`$.
In particular, $`H_\psi `$ is Koszul if $`\psi `$ has the form $`\psi =_{g\mathrm{\Gamma }}(\psi _g^{}\psi _g^{\prime \prime })g,`$ where
$`\psi ^{}:g\psi _g^{}\mathrm{\Lambda }^{a(g)}\left(M_g\right)^{},\text{resp.,}`$
$`\psi ^{\prime \prime }:g\psi _g^{\prime \prime }\mathrm{\Lambda }^{pa(g)}\left(L_g\right)^{},`$
are some $`\mathrm{\Gamma }`$-equivariant maps (where $`\mathrm{\Gamma }`$ acts on itself by conjugation). $`\mathrm{}`$
The following special case of the Theorem arises in many interesting examples. Let $`\varphi :\mathrm{\Lambda }^pVk`$ be a $`\mathrm{\Gamma }`$-invariant linear function. For each $`g\mathrm{\Gamma }`$ let
$`\psi _g=\varphi \mathrm{in}\mathrm{\Lambda }^{a(g)}(M_g)\mathrm{\Lambda }^{pa(g)}(L_g)`$
$`\psi _g=0\mathrm{in}\mathrm{any}\mathrm{other}\mathrm{component}\mathrm{\Lambda }^i(M_g)\mathrm{\Lambda }^{pi}(L_g).`$
Note that, in general, $`\psi _g\varphi `$ unless $`M_g=V`$ or 0.
Now put $`\psi =_{g\mathrm{\Gamma }}\psi _gg\mathrm{Hom}_\mathrm{\Gamma }(\mathrm{\Lambda }^pV,K)`$.
###### Corollary 4.5
With the above notations, the algebra $`H_\psi `$ is Koszul.
Symplectic reflection algebras appear as a special case of this construction where $`p=2`$ and $`\varphi `$ is a symplectic 2-form on $`V`$. Notice that the previous corollary still holds if $`\psi `$ is multiplied by any map $`m:\mathrm{\Gamma }`$ which is constant on any conjugation class.
## 5 Koszul complex
Notations and assumptions are those of Definition 3.9. So $`U`$ is a filtered $`N`$-Koszul algebra over the von Neumann regular ring $`k`$. We are interested in defining the Koszul complex of $`U`$ for bimodules. Braverman and Gaitsgory defined it in the case $`N=2`$ and $`\phi =\phi _1`$ as a subcomplex of the bar resolution (see 5.4 in ). But it seems hard to proceed in the same manner when $`N>2`$, because $`URU`$ is not naturally included in $`UU^2U`$.
Having in mind the situation of Section 4, we limit ourselves to the case $`N2`$ and $`\phi =\phi _0`$. Our method is to construct an $`N`$-differential on $`(UW_nU)_{n0}`$, and then to make an adequate contraction in order to get the Koszul complex as in . Recall the meaning of the notation $`W_n`$ for any $`n0`$:
$$W_n=\underset{i+N+j=n}{}V^iRV^j.$$
Denote by $`\mu :UUU`$ the multiplication of $`U`$, and denote by $`\mu |_{UV}`$ its restriction to $`UV`$. Define the $`U`$-$`k`$ linear map
$$d_l:UW_nUV^{(n1)}$$
as the restriction of $`\mu |_{UV}1_{V^{(n1)}}`$ to $`UW_n`$. Actually, $`d_l`$ maps into $`UW_{n1}`$. In fact, for any $`wW_n`$, one can write down $`w=_jv_jw_j`$ where $`v_jV`$ and $`w_jW_{n1}`$, so that $`d_l(aw)=_jav_jw_j`$ for any $`aU`$.
For any $`nN`$, $`W_n`$ is included in $`W_NW_{nN}`$, and one can write down
$$w=\underset{j}{}w_{j,N}w_{j,nN}^{}$$
where $`w_{j,N}W_N=R`$ and $`w_{j,nN}^{}W_{nN}`$, so
$$d_l^N(aw)=\underset{j}{}a\overline{w_{j,N}}w_{j,nN}^{}.$$
Here $`\overline{w_{j,N}}`$ is the class of $`w_{j,N}`$ in $`U^N`$. But, in $`U`$, one has $`\overline{w_{j,N}}=\phi (w_{j,N})k`$. Thus
$$d_l^N(aw)=\underset{j}{}a\phi (w_{j,N})w_{j,nN}^{}=a\underset{j}{}\phi (w_{j,N})w_{j,nN}^{},$$
so that $`d_l^N=1_U\phi ^{1,N}`$. Recall that $`\phi ^{1,N}:RV^{(nN)}V^{(nN)}`$ denotes $`\phi 1_{V^{(nN)}}`$. Notice that $`d_l`$ is an $`N`$-differential if and only if $`\phi =0`$.
Define analogously the $`k`$-$`U`$-linear map $`d_r:W_nUW_{n1}U`$. Then $`d_r^N=\phi ^{nN+1,n}1_U`$, where $`\phi ^{nN+1,n}=1_{V^{(nN)}}\phi `$.
For the sake of simplicity, the $`U`$-$`U`$-linear maps $`d_l1_U`$ and $`1_Ud_r`$ will be also denoted by $`d_l`$ and $`d_r`$ respectively. Fix a primitive $`N`$-root of unity $`q`$ (we enlarge $`k`$ if necessary). Define $`d:UW_nUUW_{n1}U`$ by $`d=d_lq^{n1}d_r`$. Explicitly we have
$$\mathrm{}\stackrel{d_ld_r}{}UW_NU\stackrel{d_lq^{N1}d_r}{}\mathrm{}\stackrel{d_lqd_r}{}UVU\stackrel{d_ld_r}{}UU.$$
(5.1)
###### Lemma 5.1
One has $`d^N=0`$, i.e., $`(UW_nU)_{n0}`$ endowed with $`d`$ is an $`N`$-complex.
Proof. Since $`d_l`$ and $`d_r`$ are commuting, we have in $`UW_nU`$
$$\underset{i=nN}{\overset{i=n1}{}}(d_lq^id_r)=d_l^Nd_r^N=1_U(\phi ^{1,N}\phi ^{nN+1,n})1_U.$$
Proposition 3.6 shows that the PBW condition (3.5) reduces in our case to the following
$$(\phi ^{1,N}\phi ^{2,N+1})(W_{N+1})=0.$$
It implies that $`(\phi ^{1,N}\phi ^{2,N+1})(W_{N+1}V)=(\phi ^{2,N+1}\phi ^{3,N+2})(VW_{N+1})=0`$, hence $`(\phi ^{1,N}\phi ^{3,N+2})(W_{N+2})=0`$. A straightforward induction provides $`(\phi ^{1,N}\phi ^{nN+1,n})(W_n)=0`$, so $`d^N:UW_nUUW_{nN}U`$ is vanishing. $`\mathrm{}`$
The associated graded of the $`N`$-complex $`(UW_nU)_{n0}`$ is the Koszul $`N`$-complex of $`A`$ in $`A`$-grMod-$`A`$ as defined in . Therefore, since $`A`$ is Koszul and the filtration is exhaustive and bounded below, the adequate contraction $`(UW_{\zeta (n)}U)_{n0}`$ is exact in degrees $`>0`$. Recall that $`\zeta :`$, where
$$\zeta (2q)=qN,\zeta (2q+1)=qN+1,q\mathrm{integer}0.$$
Explicitly, the adequate contraction is the following complex
$$\mathrm{}\stackrel{d^{N1}}{}UW_{N+1}U\stackrel{d}{}UW_NU\stackrel{d^{N1}}{}UVU\stackrel{d}{}UU.$$
(5.2)
Here $`d=d_ld_r`$ and $`d^{N1}=d_l^{N1}+d_l^{N2}d_r+\mathrm{}+d_ld_r^{N2}+d_r^{N1}`$ make sense on the genuine $`k`$. Together with $`\mu :UUU`$, $`(UW_{\zeta (n)}U)_{n0}`$ is a resolution of $`U`$ in the category of filtered bimodules $`U`$-filtMod-$`U`$, which is called the *Koszul resolution* of $`U`$.
If $`n`$ is odd, one has
$`d(a{\displaystyle \underset{i}{}}v_1^i\mathrm{}v_{\zeta (n)}^ib)`$
$`={\displaystyle \underset{i}{}}\left(av_1^iv_2^i\mathrm{}v_{\zeta (n)}^ibav_1^i\mathrm{}v_{\zeta (n1)}^iv_{\zeta (n)}^ib\right),`$
and if $`n`$ is even:
$`d^{N1}(a{\displaystyle \underset{i}{}}v_1^i\mathrm{}v_{\zeta (n)}^ib)`$
$`={\displaystyle \underset{i}{}}(av_1^i\mathrm{}v_{N1}^iv_N^i\mathrm{}v_{\zeta (n)}^ib+av_1^i\mathrm{}v_{N2}^iv_{N1}^i\mathrm{}v_{\zeta (n1)}^iv_{\zeta (n)}^ib`$
$`+\mathrm{}+av_1^i\mathrm{}v_{\zeta (nN+1)}^iv_{\zeta (nN+2)}^i\mathrm{}v_{\zeta (n)}^ib).`$
Denoting the Koszul resolution of $`U`$ by $`๐ฆ(U)`$, $`๐ฆ(U)`$ is a $`k`$-split ($`U^e`$)-projective resolution if $`k`$ is semi-simple . For any $`U`$-$`U`$-bimodule $`M`$, let us define the Hochschild homology and cohomology of the $`k`$-$`k`$-algebra $`U`$ (here $`k`$ may be non-commutative!) by
$$HH_{}(U,M)=\mathrm{Tor}_{}^{U^e/k}(M,U),HH^{}(U,M)=\mathrm{Ext}_{U^e/k}^{}(U,M).$$
So if $`k`$ is semi-simple, one has
$$HH_{}(U,M)=H_{}(M_{U^e}๐ฆ(U)),HH^{}(U,M)=H^{}(\mathrm{Hom}_{U^e}(๐ฆ(U),M)).$$
Let us specialize to the situation of Theorem 4.4, i.e. $`N=p`$ and $`U=H_\psi `$ which is supposed Koszul. Here $`W_n`$ is identified to $`_K^nE`$. So for any $`v_1,\mathrm{}v_n`$ in $`V`$, Alt$`(v_1,\mathrm{},v_n)`$ is identified to the wedge product $`v_1\mathrm{}v_n`$ *defined over* $`K=k[\mathrm{\Gamma }]`$. In the following formulas, the tensor products are also *defined over* $`K`$. So, if $`n`$ is odd:
$`d(av_1\mathrm{}v_{\zeta (n)}b)`$
$`={\displaystyle \underset{j=1}{\overset{\zeta (n)}{}}}((1)^{j1}av_jv_1\mathrm{}\widehat{v_j}\mathrm{}v_{\zeta (n)}b`$
$`(1)^{\zeta (n)j}av_1\mathrm{}\widehat{v_j}\mathrm{}v_{\zeta (n)}v_jb).`$
Note that when $`p`$ is even, $`(1)^{\zeta (n)j}=(1)^{j1}`$. Now if $`n`$ is even:
$`d^{p1}(av_1\mathrm{}v_{\zeta (n)}b)={\displaystyle \underset{1j_1<\mathrm{}<j_{p1}\zeta (n)}{}}`$
$`((1)^{j_11+\mathrm{}+j_{p1}(p1)}av_{j_1}\mathrm{}v_{j_{p1}}v_1\mathrm{}\widehat{v}_{j_1}\mathrm{}\widehat{v}_{j_{p1}}\mathrm{}v_{\zeta (n)}b`$
$`+(1)^{j_11+\mathrm{}+j_{p2}(p2)+\zeta (n)j_{p1}}av_{j_1}\mathrm{}v_{j_{p2}}v_1\mathrm{}\widehat{v}_{j_1}\mathrm{}\widehat{v}_{j_{p1}}\mathrm{}v_{\zeta (n)}v_{j_{p1}}b`$
$`+\mathrm{}+(1)^{\zeta (n)j_{p1}+\mathrm{}+\zeta (n)j_1(p2)}av_1\mathrm{}\widehat{v}_{j_1}\mathrm{}\widehat{v}_{j_{p1}}\mathrm{}v_{\zeta (n)}v_{j_1}\mathrm{}v_{j_{p1}}b),`$
which is reduced to the following when $`p`$ is even:
$`d^{p1}(av_1\mathrm{}v_{\zeta (n)}b)={\displaystyle \underset{1j_1<\mathrm{}<j_{p1}\zeta (n)}{}}(1)^{j_1+\mathrm{}+j_{p1}+\frac{p}{2}}`$
$`(av_{j_1}\mathrm{}v_{j_{p1}}v_1\mathrm{}\widehat{v}_{j_1}\mathrm{}\widehat{v}_{j_{p1}}\mathrm{}v_{\zeta (n)}b`$
$`av_{j_1}\mathrm{}v_{j_{p2}}v_1\mathrm{}\widehat{v}_{j_1}\mathrm{}\widehat{v}_{j_{p1}}\mathrm{}v_{\zeta (n)}v_{j_{p1}}b`$
$`+\mathrm{}av_1\mathrm{}\widehat{v}_{j_1}\mathrm{}\widehat{v}_{j_{p1}}\mathrm{}v_{\zeta (n)}v_{j_1}\mathrm{}v_{j_{p1}}b),`$
If $`p=2`$, the formulas for $`n`$ odd and $`n`$ even are the same and we recover in this case the differential as defined in , formula (2.6).
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# Single-copy entanglement in critical quantum spin chains
## Abstract
We introduce the single-copy entanglement as a quantity to assess quantum correlations in the ground state in quantum many-body systems. We show for a large class of models that already on the level of single specimens of spin chains, criticality is accompanied with the possibility of distilling a maximally entangled state of arbitrary dimension from a sufficiently large block deterministically, with local operations and classical communication. These analytical results โ which refine previous results on the divergence of block entropy as the rate at which EPR pairs can be distilled from many identically prepared chains, and which apply to single systems as encountered in actual experimental situations โ are made quantitative for general isotropic translationally invariant spin chains that can be mapped onto a quasi-free fermionic system, and for the anisotropic XY model. For the XX model, we provide the asymptotic scaling of $`(1/6)\mathrm{log}_2(L)`$, and contrast it with the block entropy. The role of superselection rules on single-copy entanglement in systems consisting of indistinguishable particles is emphasized.
Quantum phase transitions of second order are accompanied with a divergent length scale: this is the classical correlation length, the characteristic length associated with the two-point correlation function . Recently, it has increasingly become clear that one should expect additional insight in the scaling of quantum correlations present in the ground state of a many-body system at or close to a quantum phase transition by expressing them in terms of entanglement properties . Entanglement, after all, plays a fundamental role in quantum phase transitions at zero temperature. The theory of entanglement in turn โ developed in the quantum information context โ provides tools to characterize and quantify genuine quantum correlations in contrast to correlations that occur in states that can be prepared with mere local preparations together with classical communication (LOCC). In particular, one finds that in one-dimensional non-critical harmonic or quantum spin systems , the degree of entanglement of a block of $`L`$ systems, quantified in terms of the entropy of the reduction, typically saturates for large block size, with higher-dimensional โentropy-area lawsโ . In contrast, in critical spin systems or in fermionic systems, the entropy of a reduction has logarithmic corrections as $`L\mathrm{}`$ . These findings are consistent with expectations from conformal field theory . Such a behavior of the block entropy has also been related to the performance of DMRG simulations of ground state properties. This von-Neumann entropy of a block quantifies the rate at which one can asymptotically distill maximally entangled qubit pairs under LOCC, when one has infinitely many identically prepared many-body systems at hand .
Yet, in several contexts, in particular for condensed-matter systems, this asymptotic notion of entanglement implicitly referring to joint operations on many identical systems may not always be the most appropriate one. Instead, one may ask: does a single specimen of a critical infinite system already contain an infinite amount of entanglement? This will be the central question of this paper. We introduce the single-copy entanglement to quantify the quantum correlations in critical and non-critical many-body systems. More specifically, compared to the divergence of the block entropy, we ask the stronger question whether a single spin chain already contains an arbitrary amount of entanglement, such that from a single specimen a maximally entangled state of arbitrary dimension can be distilled.
We will make the argument quantitative by analytically considering a general framework of translationally invariant quantum spin models. As examples in which criticality is in one-to-one correspondence with a divergent single-copy entanglement, we consider isotropic spin models, as well as the XY-model. For the isotropic XY model we establish the exact asymptotic scaling behavior of $`(1/6)\mathrm{log}_2(L)`$, and relate it to the block entropy . The results can also be conceived as statements concerning the divergence of fine-grained entanglement .
Single-copy entanglement. โ Let us consider a one-dimensional quantum spin system, associated with a Hilbert space $`=(^2)^n`$, with a translationally invariant Hamiltonian. We distinguish a block of length $`L`$ of consecutive systems of the chain. So we have a bi-partioning $`n|L`$, the whole system being in a pure state $`\rho =|\psi \psi |`$.
There are several meaninful definitions of single-copy entanglement. We will primarily be concerned with the question: running a physical device once, a maximally entangled state of what dimension can be distilled from a single specimen with certainty? Hence, the state has a single-copy entanglement $`E_1(\rho )=\mathrm{log}_2(M)`$, with respect to the bi-partitioning $`n|L`$, if $`\rho `$ can be deterministically transformed under LOCC into $`|\psi _M\psi _M|`$, i.e., a maximally entangled state with state vector $`|\psi _M=(|1,1+\mathrm{}+|M,M)/\sqrt{M}`$, so if
$$\rho |\psi _M\psi _M|\text{ under LOCC}.$$
(1)
This is the non-asymptotic analogue of the entropy of entanglement of the reduction associated with a block of length $`L`$. Denote with $`\alpha _1^{},\mathrm{},\alpha _{2^L}^{}`$ the non-increasingly ordered eigenvalues of the reduced state with respect to a block of length $`L`$, then Eq. (1) holds true if and only if $`_{k=1}^K\alpha _k^{}K/M`$ for all $`1KM`$, so obviously, if and only if $`\alpha _1^{}1/M`$. In other words, the transformation is possible if the reduction is more mixed in the sense of majorization than the reduction of the maximally entangled state of dimension $`M\times M`$. Given $`\alpha _1^{}`$, the single-copy entanglement is nothing but $`E_1(\rho )=\mathrm{log}_2((\alpha _1^{})^1)`$. A variant is one allowing for probabilistic protocols. For a state $`\rho `$ we say that
$`E_\text{p}(\rho )=sup{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}p_k\mathrm{log}_2(M_k),`$
such that $`\rho `$ can be transformed under LOCC into the ensemble $`\{(p_k,|\psi _{M_k}\psi _{M_k}|):k=0,1,\mathrm{}\}`$. This is the average entanglement that can be distilled, allowing for maximally entangled states of different dimension with certain probabilities . This rate is then the solution of a linear program . By definition, we have that $`E_1(\rho )E_\text{p}(\rho )S(\text{tr}_{n\backslash L}[\rho ])`$, where the last inequality follows from the fact that the entropy of a reduction bounds the rate of any (asymptotic) distillation protocol.
Finally, note that for the single-copy entanglement, superselection rules (SSR) play a crucial role, notably a SSR with respect to particle number conservation. In the presence of SSR, $`E_1^{\text{SSR}}`$ has to be understood as referring to a probabilistic transformation distilling entanglement-EPR states with state vector $`|\psi ^{\text{SSR}}=(|0,1|1,0+|1,0|0,1)/\sqrt{2}`$ under LOCC and SSR. We will now consider the behavior of the single-copy entanglement in the limit of large $`L`$ for critical and non-critical spin chains.
Single-copy entanglement in general quantum spin chains. โ We start from the general set of translationally invariant quantum spin systems that is as in Ref. mapped onto a fermionic quadratic Hamiltonian under a Jordan-Wigner transformation. This model embodies a large class of spin models, including the anisotropic and isotropic XY-models as important special cases. Hence, the Hamiltonian is
$`H={\displaystyle \underset{q=0,1}{}}{\displaystyle \underset{k,j=1}{\overset{n}{}}}\left({\displaystyle \frac{B_{jk}}{2}}{\displaystyle \frac{A_{jk}}{4}}\right)\widehat{\sigma }_k^{q+1}\left[{\displaystyle \underset{\genfrac{}{}{0pt}{}{ik+q}{lj+1q}}{}}\widehat{\sigma }_i^3\widehat{\sigma }_l^3\right]\widehat{\sigma }_j^{q+1}`$
where $`\widehat{\sigma }_k^1,\widehat{\sigma }_k^2,\widehat{\sigma }_k^3`$ denote the Pauli operators associated with site $`k=1,\mathrm{},n`$, equivalent with the fermionic Hamiltonian
$$H=\underset{j,k=1}{\overset{n}{}}\left[\widehat{a}_j^{}A_{jk}\widehat{a}_k+\widehat{a}_j^{}B_{jk}\widehat{a}_k^{}\widehat{a}_jB_{jk}\widehat{a}_k\right]$$
The fermionic operators obey $`\{\widehat{a}_j,\widehat{a}_k\}=0`$ and $`\{\widehat{a}_j^{},\widehat{a}_k\}=\delta _{j,k}`$. The Hamiltonians are related via a Jordan-Wigner transformation leading to the Hermitian Majorana operators, $`\widehat{m}_{2i1}=(_{j<i}\widehat{\sigma }_j^3)\widehat{\sigma }_i^1`$ and $`\widehat{m}_{2i}=(_{j<i}\widehat{\sigma }_j^3)\widehat{\sigma }_i^2`$, where $`\widehat{a}_j=(\widehat{m}_{2j1}i\widehat{m}_{2j})/2`$. Translational invariance, periodic boundary conditions, and Hermicity are inherited by $`A_j,B_j`$ satisfying $`A_j=A_j`$, and $`B_j=B_j`$ for $`j=1n,\mathrm{},n1`$. For simplicity, we assume that there exists a $`w`$ such that $`A_j=B_j=0`$ for $`j>w`$. This model will be our starting point. For all isotropic instances, and also for the full XY model we will be able to identify when the single-copy entanglement is indeed logarithmically divergent. Then, one may distill a maximally entangled state of any dimension from a single specimen of the chain with certainty, containing in this sense an โinfinite single-copy entanglementโ . We will make use of the powerful methods of Toeplitz determinants . This path is yet in our instance complicated by the fact that we do not only consider isotropic models, and that in contrast to the block entropy the largest eigenvalue cannot straightforwardly be expressed as an integral of a Toeplitz determinant. The starting point, yet, is the familiar one for assessing spin systems: The ground state of this system is a fermionic Gaussian, i.e., quasi-free, state and is completely specified by the second moments of the Majorana operators. These operators satisfy $`\widehat{m}_j=\widehat{m}_j^{}`$ and $`\{\widehat{m}_j,\widehat{m}_k\}=2\delta _{j,k}`$. The second moments can be collected in a correlation matrix $`\gamma ^{2n\times 2n}`$, $`\text{tr}[\rho \widehat{m}_j\widehat{m}_k]=\delta _{j,k}+i\gamma _{j,k}`$. This matrix is skew-symmetric. The entanglement properties of the block of length $`L`$ can now be inferred from a principal submatrix $`\gamma _L^{2L\times 2L}`$ of the correlation matrix $`\gamma `$. We consider the entries of $`\gamma _L`$ in the limit of an infinite chain $`n\mathrm{}`$. Then, $`\gamma _L`$ is a block Toeplitz matrix, the $`l`$-th row, $`l=1,\mathrm{},L`$, being given by $`(M_{l1},M_{l2},\mathrm{},M_0,\mathrm{},M_{lL})`$, with $`2\times 2`$-blocks $`M_{L1},\mathrm{},M_{1L}`$ that are found to be
$$M_l=\left[\begin{array}{cc}0& t_l\\ t_l& 0\end{array}\right],t_l=\frac{1}{2\pi }_0^{2\pi }g(k)e^{ilk}๐k.$$
For no anisotropy, i.e., $`B_j=0`$ for all $`j`$, this matrix is a tensor product of a symmetric matrix and a unit skew-symmetric one. In generality, we have for this model, $`g(k):=\mathrm{\Lambda }(k)/|\mathrm{\Lambda }(k)|`$, $`\mathrm{\Lambda }(k):=A_0+2_{j=1}^wA_j\mathrm{cos}(jk)4i_{j=1}^wB_j\mathrm{sin}(jk)`$. This matrix $`\gamma _L`$ can be brought into a standard normal form $`\mathrm{\Gamma }_L`$ of a skew-symmetric matrix with an $`OO(2L)`$ preserving the anticommutation relations,
$$\mathrm{\Gamma }_L=O\gamma _LO^T,\mathrm{\Gamma }_L=\underset{l=1}{\overset{L}{}}\left[\begin{array}{cc}0& \mu _l\\ \mu _l& 0\end{array}\right].$$
This defines the quantities $`\mu _1,\mathrm{},\mu _L[0,1]`$. Such normal mode decompositions have been employed both to evaluate correlation functions and the block entropy . From now on we will be concerned with the largest eigenvalue of the reduction of a block of length $`L`$. All eigenvalues $`\alpha _1^{},\mathrm{},\alpha _{2L}^{}`$ of the reduction are given by $`\{\alpha _1^{},\mathrm{},\alpha _{2L}^{}\}=\{_{l=1}^L(1\pm \mu _l)/2\}`$. We will be looking at the behavior of the largest eigenvalue $`\alpha _1^{}`$ for large $`L`$. This largest eigenvalue is given by $`\alpha _1^{}=_{l=1}^L(1/2+\mu _l/2)`$, or
$$\alpha _1^{}=det[(\mathrm{๐}_L+|T_L|)/2],$$
$`|T_L|=(T_L^TT_L)^{1/2}`$, where $`T_L`$ is the $`L\times L`$ Toeplitz matrix, with $`l`$-th row being given by $`(t_{l+1},t_{l+2},\mathrm{},t_0,\mathrm{},t_{Ll})`$. The numbers $`\mu _1,\mathrm{},\mu _L`$ are the singular values of $`T_L`$. This matrix $`T_L`$, satisfying $`|T_L|\mathrm{๐}_L`$, is generally not symmetric, as a consequence of the anisotropy of the model. Moreover, in contrast to the matrix $`T_L`$ itself, $`\mathrm{๐}_L+|T_L|`$ is not Toeplitz. In order to show that the single-copy entanglement is logarithmically divergent, we will make use of appropriate bounds that retain this property: whenever the $`A_0,\mathrm{},A_w`$, $`B_0,\mathrm{},B_w`$ are such that one can prove that the sequence of $`L\times L`$-Toeplitz matrices $`T_L`$ satisfies
$$\mathrm{log}|det[T_L]|=\mathrm{\Omega }(\mathrm{log}(L))$$
(2)
(using Landau notation ) using a Fisher-Hartwig-statement , then one can indeed conclude that $`E_1=\mathrm{\Omega }(\mathrm{log}(L))`$, i.e, the single-copy entanglement diverges at least logarithmically with increasing block length $`L`$. This follows from the following chain,
$``$ $`\mathrm{log}det[(\mathrm{๐}_L+|T_L|)/2]{\displaystyle \frac{1}{2}}\mathrm{log}det[(\mathrm{๐}_L+T_L^TT_L)/2]`$
$``$ $`{\displaystyle \frac{1}{4}}\mathrm{log}det[T_L^TT_L]={\displaystyle \frac{1}{2}}\mathrm{log}|det[T_L]|`$
, where we also have made use of the concavity of the logarithm. So, whenever Eq. (2) holds, for appropriate length of the block $`L`$, a maximally entangled pair of any dimension can be distilled from a single specimen of the spin chain.
Isotropic models. โ This case of $`B_0,\mathrm{},B_w=0`$ is particularly transparent. Here, the asymptotics in $`L`$ of the determinants $`det[M_{x,L}]`$ of the $`L\times L`$-Toeplitz matrices $`M_{x,L}:=ix\mathrm{๐}+(1x^2)^{1/2}T_L`$ is known for all $`x(0,1)`$, using a Fisher-Hartwig statement. This small detour to infer about $`\mathrm{log}|det[T_L]|`$ โ corresponding to the case $`x=0`$ โ is needed as the Fisher-Hartwig-conjecture has not been proven yet for this case. In general, one can identify the asymptotic behavior of determinants of Toeplitz matrices by investigating the so-called symbol, see footnote . The symbol associated with the Toeplitz matrices $`M_{x,L}`$ is given by
$$G_x(k)=ix+(1x^2)^{1/2}g(k),$$
with $`g`$ as defined above. For this class of isotropic models, an explicit factorization of the symbol is known , see footnote . It follows hence from proven instances of the Fisher-Hartwig conjecture that there exists a $`c>0`$ and an $`x_0(0,1)`$ such that
$$\mathrm{log}|det[M_{x,L}]|=c_x\mathrm{log}(L)+o(\mathrm{log}(L))$$
with $`c_x>c`$ for all $`x(0,x_0)`$, whenever the function $`g`$ is discontinuous in $`[0,2\pi ]`$, where the jumps reflect the Fermi surface. From this โ and using that $`T_L`$ has real eigenvalues โ it follows that the system has a logarithmically divergent single-copy entanglement if the system is critical . For example, for the XX model this analysis immediately delivers a logarithmically divergent single-copy entanglement, whenever the system is critical.
Anisotropic XY-model. โ For the XY-model we can conclude that the single-copy entanglement is logarithmically divergent if and only if the system is critical. For this model, we have that $`A_0=1,A_1=a/2`$ and $`B_1=B_1=\gamma a/4`$, and $`0`$ elsewhere. For $`\gamma =0`$, we obtain the XX-model (the isotropic XY-model), for $`a=1`$, $`\gamma =1`$ the critical Ising model. Along the line $`\gamma [1,1]`$, $`a=1`$ the anisotropic model is critical. Then, we encounter a generally non-symmetric matrix $`T_L`$. The associated symbol is given by
$`g(k)={\displaystyle \frac{a\mathrm{cos}(k)1+ia\gamma \mathrm{sin}(k)}{((a\mathrm{cos}(k)1)^2+\gamma ^2a^2\mathrm{sin}^2(k))^{1/2}}}.`$
For $`\gamma 0`$ and $`1/a(0,1)`$, the symbol is continuous, and one finds a saturating block entropy (and hence a saturating single-copy entanglement). Along the critical line $`a=1`$, $`\gamma (1,1)`$, in turn, we can identify the explicit factorization of the discontinuous symbol. There is a single discontinuity at $`k_1=0`$ , and in the terms of footnote we find $`\beta _1=1/2`$, so that $`g(k)`$ can be decomposed as
$$g(k)=\varphi (k)t_{1/2}(k),$$
where $`\varphi `$ is a continuously differentiable function. For the case of a single discontinuity and $`\alpha _1=0`$, the Fisher-Hartwig conjecture has been proven for any $`\beta _1`$ with $`\mathrm{}(\beta _1)<5/2`$ , including our case at hand. Hence, we find $`\mathrm{log}|det[T_L]|=\mathrm{\Omega }(\mathrm{log}(L))`$, and hence $`E_1=\mathrm{\Omega }(\mathrm{log}(L))`$. Together with the result of the subsequent section this shows that the single-copy entanglement of the XY-model is logarithmically divergent exactly if the model is critical. Note that this implies also a less technical alternative proof of the logarithmic divergence of the block entropy in the critical XY model.
Scaling of single-copy entanglement in the XX-model. โ In the light of these findings, it is interesting to see how the exact asymptotic behavior is compared to that of the block entropy, including prefactors. We make this specific for the isotropic XX-model, where now $`T_L=T_L^T`$. The technicality when evaluating $`\alpha _1^{}=det[(\mathrm{๐}_L+|T_L|)/2]`$ that we encounter here is that the function $`f:`$, $`f(x):=\mathrm{log}_2((1+|x|)/2)`$, is not analytic. So before we can exploit Fisher-Hartwig-type results, we have to approximate $`\alpha _1^{}`$ with sequences based on functions with appropriate continuity properties. We can take any functions $`f_{}:\times _+`$ which are analytic on $`\{z:\mathrm{}(z)<\delta \}`$ for a $`\delta >0`$, such that on the real axis $`lim_{\delta 0}f_{}(x,\delta )=f(x,0)`$ for $`x`$. Take, e.g.,
$$f_{}(z,\delta ):=\mathrm{log}(1/2+(z^2+\delta ^2)^{1/2}/2).$$
We are then in the position to identify the asymptotic behavior of the single-copy entanglement. This can be done similarly to Ref. using the characteristic polynomial $`F:`$ of $`T_L`$ defined as $`F(\lambda ):=det[\lambda \mathrm{๐}_LT_L]`$: the function $`F`$ is meromorphic, and all zeros are in the interval $`[1,1]`$. One can hence write
$`d_{}=\underset{\delta 0}{lim}\underset{\epsilon 0}{lim}{\displaystyle \frac{1}{2\pi i}}{\displaystyle ๐zf_{}(z,\delta )\frac{F^{}(z)}{F(z)}}`$ (3)
where the integration path is chosen to enclose the interval $`[1,1]`$, with path from $`(1\delta +i\epsilon ,1+\delta +i\epsilon )`$, towards the negative real numbers along a circle segment with radius $`\delta /2`$, then $`(1+\delta i\epsilon ,1\delta i\epsilon )`$, and again along a circle segment to $`1\delta +i\epsilon `$, such that $`lim_{\delta 0}d_{}=d`$. The symbol of $`\lambda \mathrm{๐}T_L`$ with factorization as in Eq. (4) for the XX-model is known , see footnote . Using a Fisher-Hartwig statement, we find that the linear terms in $`L`$ do not contribute, using Cauchyโs theorem and using that $`lim_{\delta 0}f_{}(\pm 1,\delta )=0`$, and finally arrive at
$`d`$ $`=`$ $`\mathrm{log}(L){\displaystyle \frac{2}{\pi ^2}}{\displaystyle _1^1}{\displaystyle \frac{\mathrm{log}_2[(1+|x|)/2]}{1x^2}}๐x+o(\mathrm{log}(L)).`$
This in turn finally implies that whenever $`1/a[1,1]`$ and the XX-model is critical, we observe the scaling behavior
$`E_1={\displaystyle \frac{1}{6}}\mathrm{log}_2(L)+o(\mathrm{log}(L)),`$
independent of $`a`$; it saturates in the non-critical case. This result is astonishing: the single-copy entanglement does not only diverge, but has up to a factor of two the same asymptotic behavior as the entropy of entanglement scaling as $`S=(1/3)\mathrm{log}_2(L)+o(\mathrm{log}(L))`$. Half of the asymptotically distillable entanglement is hence already available on the single-shot level.
Outlook and summary. โ Finally, let us comment on the crucial role of SSR for the single-copy entanglement. This is relevant, e.g., when assessing the single-copy entanglement in the hard-core limit of the Bose-Hubbard model (infinite repulsion energy) . There, the Hamiltonian is isomorphic to the XX-model, via the mapping $`\widehat{\sigma }_j^1=\widehat{b}_j+\widehat{b}_j^{}`$, $`\widehat{\sigma }_j^2=i(\widehat{b}_j\widehat{b}_j^{})`$, and $`\widehat{\sigma }_j^3=12\widehat{b}_j^{}\widehat{b}_j`$ for each site $`j`$. Yet, the concept of entanglement is different due to the presence of a particle number conservation SSR in the former case: Transformations under LOCC have to be replaced by those under LOCC+SSR. The single-copy entanglement in the above sense can however still be efficiently evaluated in $`L`$; and these superselection rules must be respected when assessing single-copy entanglement.
In this paper, we have fleshed out the notion of single-copy entanglement in quantum spin chains. Such a notion is the appropriate one when one is not interested in the entanglement properties of an asymptotic supply of a identically prepared many-body systems, but of single specimens. It is the hope that these findings also serve as a guideline when assessing entanglement in actual experimental situations, let it be in condensed-matter systems or in systems of ultracold atoms in optical lattices.
Acknowledgements. โ We would like to thank J.I. Cirac, T. Cubitt, M.B. Plenio, D. Schlingemann, R.F. Werner, and M.M. Wolf for discussions. This work has been supported by the DFG (SPP 1116, SPP 1078), the EU (QUPRODIS), the EPSRC, and the European Research Councils (EURYI).
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# Blow-up of regular submanifolds in Heisenberg groups and applications
## 1 Introduction
In recent years, several efforts have been devoted to the project of developing Analysis and Geometry in stratified groups and more general Carnot-Carathรฉodory spaces with several monographs and surveys on this subject. Among them we mention , , , , , , but this list could be surely enlarged.
Our study fits into the recent project of developing Geometric Measure Theory in these spaces. Ambient of our investigations is the $`(2n`$$`+`$$`1)`$-dimensional Heisenberg group $`^n`$, which represents the simplest model of non-Abelian stratified group, , . Aim of this paper is to present an intrinsic blow-up theorem for $`C^1`$ submanifolds in the geometry of the Heisenberg group along with its applications. The main feature of this procedure is the use of natural dilations of the group, namely, a one-parameter family of group homomorphisms that are homogeneous with respect to the distance of the group. Recall that dilations in $`^n`$ are anisotropic, hence they differently act on different directions of the submanifold. The foremost directions are the so-called horizontal directions, that determine the โsub-Riemannian geometryโ of the Heisenberg group: at any point $`x^n`$ a $`2`$$`n`$-dimensional subspace $`H_x^nT_x^{2n+1}`$ is given and the family of all horizontal spaces $`H_x^n`$ forms the so-called horizontal subbundle $`H^n`$. We will defer full definitions to Section 2.
The blow-up procedure consists in enlarging the submanifold $`\mathrm{\Sigma }`$ at some point $`x\mathrm{\Sigma }`$ by intrinsic dilations and taking the intersection of the magnified submanifold with a bounded set centered at $`x`$. We are interested in studying the case when $`T_x\mathrm{\Sigma }H_x^n`$, namely, $`x`$ is a transverse point. The effect of rescaling the submanifold at a transverse point $`x`$ can be obtained by considering the behavior of $`\text{vol}_p(B_{x,r}\mathrm{\Sigma })/r^{p+1}`$ as $`r0^+`$, that heuristically is
$`{\displaystyle \frac{\text{vol}_p(B_{x,r}\mathrm{\Sigma })}{r^{p+1}}}={\displaystyle \frac{\text{vol}_p\left(l_x\delta _r(B_1\mathrm{\Sigma }_{x,r})\right)}{r^{p+1}}}={\displaystyle \frac{\text{vol}_p\left(\delta _r(B_1\mathrm{\Sigma }_{x,r})\right)}{r^{p+1}}}\alpha (x)\text{vol}_p(B_1\mathrm{\Sigma }_{x,r}).`$
Here $`\text{vol}_p`$ denotes the $`p`$-dimensional Riemannian measure restricted to $`\mathrm{\Sigma }`$, the left translation $`l_x:^n^n`$ is given by $`l_x(y)=xy`$, the dilation of factor $`r>0`$ is $`\delta _r:^n^n`$, the dilated submanifold at $`x`$ is $`\mathrm{\Sigma }_{x,r}=\delta _{1/r}\left(l_{x^1}\mathrm{\Sigma }\right)`$ and $`B_{x,r}`$ is the open ball of center $`x`$ and radius $`r`$ with respect to a fixed homogeneous distance. The meaning of $`\alpha (x)`$ will be clear in the following theorem, that makes rigorous our previous consideration and represents our first main result.
###### Theorem 1.1 (Blow-up)
Let $`\mathrm{\Sigma }`$ be a $`p`$-dimensional $`C^1`$ submanifold of $`\mathrm{\Omega }`$, where $`\mathrm{\Omega }`$ is an open subset of $`^n`$ and let $`x`$ be a transverse point. Then the following limit holds
$`\underset{r0^+}{lim}{\displaystyle \frac{\text{vol}_p(\mathrm{\Sigma }B_{x,r})}{r^{p+1}}}={\displaystyle \frac{\theta _p^\rho \left(\tau _{\mathrm{\Sigma },๐ฑ}(x)\right)}{|\tau _{\mathrm{\Sigma },๐ฑ}(x)|}}.`$ (1)
A novel object appearing in this limit is the vertical tangent $`p`$-vector $`\tau _{\mathrm{\Sigma },๐ฑ}(x)`$, introduced in Definition 2.13. Its associated $`p`$-dimensional subspace of $`๐ฅ^n`$ is a subalgebra whose image through the exponential map represents the blow-up limit of the rescaled submanifold $`\mathrm{\Sigma }_{x,r}`$ as $`r0^+`$. The $`p`$-vector $`\tau _{\mathrm{\Sigma },๐ฑ}(x)`$ in higher codimension plays the same role that the well known horizontal normal $`\nu _H`$ plays in codimension one (compare for instance with ). The metric factor $`\theta \left(\tau _{\mathrm{\Sigma },๐ฑ}(x)\right)`$, introduced in , corresponds to the measure of the intersection of $`B_1`$ with the vertical subspace associated to the vertical tangent $`p`$-vector $`\tau _{\mathrm{\Sigma },๐ฑ}(x)`$. A first consequence of Theorem 1.1 is an explicit formula to compute the $`(p`$$`+`$$`1)`$-dimensional spherical Hausdorff measure of $`p`$-dimensional $`C^1`$ submanifolds in the Heisenberg group. In fact, thanks to $`๐ฎ^{p+1}`$-negligibility of characteristic points proved in , Theorem 1.1 along with standard theorems on differentation of measures, immediately give the following result.
###### Theorem 1.2
Let $`\rho `$ be a homogeneous distance with constant metric factor $`\alpha >0`$ and let $`๐ฎ_^n^{p+1}=\alpha ๐ฎ_\rho ^{p+1}`$. Then we have
$`๐ฎ_^n^{p+1}(\mathrm{\Sigma })={\displaystyle _\mathrm{\Sigma }}|\tau _{\mathrm{\Sigma },๐ฑ}(x)|๐\text{vol}_p(x).`$ (2)
Note that in codimension one, the integral formula (2) fits into the results of in stratified groups. The connection between these results is shown in Proposition 4.18. There are several examples of homogeneous distances satisfying hypothesis of Theorem 1.2, as we show in Example 4.6. Proposition 4.5 shows a class of homogeneous distances having constant metric factor. Proposition 4.10 shows how the computation of the $`(p`$$`+1)`$-dimensional spherical Hausdorff measure of a submanifold can be easily performed in several examples, that will appear in Section 4.
Another consequence of Theorem 1.1 is the validity of an intrinsic coarea formula for vector-valued Lipschitz mappings defined on the Heisenberg group. By Sard theorem and the classical Whitney approximation theorem we can assume that a.e. level set is a submanifold of class $`C^1`$, then we apply representation formula (2). The core of the proof stands in the key relation
$`|\tau _{\mathrm{\Sigma },๐ฑ}(x)|={\displaystyle \frac{J_Hf(x)}{J_gf(x)}},`$ (3)
which surprisingly connects vertical tangent $`p`$-vector with horizontal jacobian $`J_Hf`$. The proof of (3) is given in Theorem 3.3. Thus, we can establish the following result.
###### Theorem 1.3 (Coarea formula)
Let $`f:A^k`$ be a Riemannian Lipschitz map, where $`A^n`$ is a measurable subset and $`1k<2n+1`$. Let $`\rho `$ be a homogeneous distance with constant metric factor $`\alpha >0`$. Then for every measurable function $`u:A[0,+\mathrm{}]`$ the formula
$`{\displaystyle _A}u(x)J_Hf(x)๐x={\displaystyle _^k}\left({\displaystyle _{f^1(t)A}}u(y)๐๐ฎ_^n^{p+1}(y)\right)๐t`$ (4)
holds, where $`p=2n+1k`$ and $`๐ฎ_^n^{p+1}=\alpha ๐ฎ_\rho ^{p+1}`$.
This coarea formula along with that of , which is a particular case, represent first examples of intrinsic coarea formulae for vector valued mappings defined on non-Abelian Carnot groups. It remains an interesting open question the extension of coarea formula to Lipschitz mappings with respect to a homogeneous distance. Only in the case of real-valued mappings this problem has been settled in . This question is intimately related to a blow-up theorem of โintrinsicly regularโ submanifolds. In this connection, we mention a recent work by Franchi, Serapioni and Serra Cassano , where a notion of intrinisic submanifold in $`^n`$ has been introduced in arbitrary codimension. According the their terminology, a $`k`$-codimensional $``$-regular submanifold for algebraic reasons must satisfy $`1kn`$. With this restriction it might be highly irregular, even unrectifiable in the Euclidean sense, . Nevertheless they show that an area-type formula for its $`(p`$$`+`$$`1)`$-dimensional spherical Hausdorff measure still holds. Here we wish to emphasize the difference in our approach, where we consider $`C^1`$ submanifolds, but with no restriction on their codimension.
Let us summarize the contents of the present paper. Section 2 recalls some notions. Section 3 is devoted to the proof of Theorem 1.1. In Section 4 we show the validity of Theorem 1.2, along with its applications. Precisely, in Theorem 4.9 we show how a suitable rescaling of the spherical Hausdorff measure yields an intrinsic surface measure only depending on the sub-Riemannian metric, namely, the restriction of the Riemannian metric to the horizontal subbundle. In Proposition 4.5, we single out a privileged class of homogeneous distances having constant metric factor. We present several explicit computations of $`(p`$$`+`$$`1)`$-dimensional spherical Hausdorff measure in concrete examples. As another application of Theorem 1.2, we show a lower semicontintuity result for the spherical Hausdorff measure with respect to weak convergence of regular currents. Section 5 establishes an intrinsic coarea formula for vector-valued Riemannian Lipschitz mappings on the Heisenberg group. Acknowledgments. I wish to thank Bruno Franchi, Raul Serapioni and Francesco Serra Cassano for pleasant discussions on intrinsic surface area in Heisenberg groups.
## 2 Some basic notions
The $`(2n`$$`+`$$`1)`$-dimensional Heisenberg group $`^n`$ is a simply connected Lie group whose Lie algebra $`๐ฅ^n`$ is equipped with a basis $`(X_1,\mathrm{},X_{2n},Z)`$ satisfying the bracket relations
$`[X_k,X_{k+n}]=2Z`$ (5)
for every $`k=1,\mathrm{},n`$. We will identify the Lie algebra $`๐ฅ^n`$ with the isomorphic Lie algebra of left invariant vector fields on $`^n`$, so that any $`X_j`$ also denotes a left invariant vector field of $`^n`$. In the terminology of Differential Geometry, the basis $`(X_1,\mathrm{},X_{2n},Z)`$ forms a moving frame in $`^n`$. We will say that $`(X_1,\mathrm{},X_{2n},Z)`$ is our standard frame. In particular, $`(X_1,\mathrm{},X_{2n})`$ is a horizontal frame and it spans a smooth distribution of $`2n`$-dimensional hyperplanes, called horizontal hyperplanes and denoted by $`H_x^n`$ for every $`x^n`$. The collection of all horizontal hyperplanes forms the so called horizontal subbundle, denoted by $`H^n`$. In the sequel, we will fix the unique left invariant Riemannian metric $`g`$ such that the standard frame $`(X_1,X_2,\mathrm{},X_{2n},Z)`$ forms an orthonormal basis at each point.
###### Definition 2.1
Every set of left invariant vector fields $`(Y_1,\mathrm{},Y_{2n})`$ spanning the horizontal hyperplane at the unit element of $`^n`$ will be called horizontal frame.
Recall that the exponential map $`\mathrm{exp}:๐ฅ^n^n`$ is a diffeomorphism, then it is possible to introduce a system of coordinates in all of $`^n`$.
###### Definition 2.2 (Graded coordinates)
Let $`(Y_1,\mathrm{},Y_{2n})`$ be a horizontal frame and let $`W`$ be a non horizontal left invariant vector field. The frame $`(Y_1,\mathrm{},Y_{2n},W)`$ defines a coordinate chart $`F:^{2n+1}^n`$ given by
$`F(y)=\mathrm{exp}\left(y_{2n+1}W+{\displaystyle \underset{j=1}{\overset{2n}{}}}y_jY_j\right).`$ (6)
Coordinates defined by (6) are called graded coordinates in the case $`W=Z`$ and standard coordinates in the case the standard frame $`(X_1,\mathrm{},X_{2n},Z)`$ is used. In general we will say that the coordinates are associated to the frame $`(Y_1,\mathrm{},Y_{2n},W)`$
We will assume throughout that a system of standard coordinate is fixed, if not stated otherwise.
###### Remark 2.3
Note that the horizontal frame $`(Y_1,\mathrm{},Y_{2n})`$ of Definition 2.2 may not satisfy relations (5), where $`X_i`$ are replaced by $`Y_i`$.
The standard frame with respect to standard coordinates reads as follows
$`\stackrel{~}{X}_k=_{x_k}x_{k+n}_{x_{2n+1}},\stackrel{~}{X}_{k+n}=_{x_{k+n}}+x_k_{x_{2n+1}}\text{and}\stackrel{~}{Z}=_{x_{2n+1}}`$ (7)
and the group operation is given by the following formula
$`xy=(x_1+y_1,\mathrm{},x_{2n}+y_{2n},x_{2n+1}+y_{2n+1}+{\displaystyle \underset{j=1}{\overset{n}{}}}(x_ky_{k+n}x_{k+n}y_k)).`$ (8)
A natural family of dilations which respects the group operation (8) can be defined as follows
$`\delta _r(x)=(rx_1,rx_2,\mathrm{},rx_{2n},r^2x_{2n+1})`$ (9)
for every $`r>0`$. In fact, the map $`\delta _r:^n^n`$ defined above is a group homomorphism with respect to the operation (8).
In contrast with Analysis in Euclidean spaces, where the Euclidean distance is the most natural choice, in the Heisenberg group several distances have been introduced for different purposes. However, all of them are homogeneous in the following sense. If $`\rho :^n\times ^n[0,+\mathrm{}+[`$ is a homogeneous distance, then
1. $`\rho `$ is a continuous with respect to the topology of $`^n`$,
2. $`\rho (xy,xz)=\rho (y,z)`$ for every $`x,y,z^n`$,
3. $`\rho (\delta _ry,\delta _rz)=r\rho (y,z)`$ for every $`y,z^n`$ and every $`r>0`$.
To simplify notations we write $`\rho (x,0)=\rho (x)`$, where $`0`$ denotes either the origin of $`^{2n+1}`$ or the unit element of $`^n`$. The open ball of center $`x`$ and radius $`r>0`$ with respect to a homogeneous distance is denoted by $`B_{x,r}`$. The Carnot-Carathรฉodory distance is an important example of homogeneous distance, . However, all of our computations hold for a general homogeneous distance, therefore in the sequel $`\rho `$ will denote a homogeneous distance, if not stated otherwise. Note that the Hausdorff dimension of $`^n`$ with respect to any homogeneous distance is $`2n+2`$. Next, we recall the notion of Riemannian jacobian.
###### Definition 2.4 (Riemannian jacobian)
Let $`f:MN`$ be a $`C^1`$ smooth mapping of Riemannian manifolds and let $`xM`$, where $`M`$ and $`N`$ have dimension $`d`$ and $`k`$, respectively. The Riemannian jacobian of $`f`$ at $`x`$ is given by
$`J_gf(x)=\mathrm{\Lambda }_k\left(df(x)\right),`$ (10)
where $`\mathrm{\Lambda }_k\left(df(x)\right):\mathrm{\Lambda }_d(T_xM)\mathrm{\Lambda }_k(T_{f(x)}N)`$ is the canonical linear map associated to $`df(x):T_xMT_{f(x)}N`$. The norm of $`\mathrm{\Lambda }_k\left(df(x)\right)`$ is understood with respect to the induced scalar products on $`\mathrm{\Lambda }_d(T_xM)`$ and $`\mathrm{\Lambda }_k(T_{f(x)}N)`$. We recall scalar products of $`p`$-vectors in (25).
To compute the Riemannian jacobian, we fix two orthonormal bases $`(X_1,\mathrm{},X_d)`$ and $`(E_1,\mathrm{},E_k)`$ of $`T_xM`$ and $`T_{f(x)}N`$, respectively, and we represent $`df(x)`$ with respect to these bases by the matrix
$`_{X,E}f(x)=\left[\begin{array}{cccccc}E_1,df(x)(X_1)& E_1,df(x)(X_2)& \mathrm{}& E_1,df(x)(X_d)& & \\ E_2,df(x)(X_1)& E_2,df(x)(X_2)& \mathrm{}& E_2,df(x)(X_d)& & \\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& & \\ E_k,df(x)(X_1)& E_k,df(x)(X_2)& \mathrm{}& E_k,df(x)(X_d)& & \end{array}\right].`$ (15)
Then the jacobian of the matrix $`_{X,E}f(x)`$ coincides with $`J_gf(x)`$. In the sequel, it will be useful to fix the following notation to indicate minors of a matrix.
###### Definition 2.5
Let $`G`$ be an $`m\times n`$ matrix with $`mn`$. We denote by $`G_{i_1i_2\mathrm{}i_m}`$ the $`m\times m`$ submatrix with columns $`(i_1,i_2,\mathrm{},i_m)`$. We define the minor
$`M_{i_1i_2\mathrm{}i_m}(G)=det\left(G_{i_1i_2\mathrm{}i_m}\right).`$ (16)
###### Definition 2.6 (Horizontal jacobian)
Let $`\mathrm{\Omega }`$ be an open subset of $`^n`$ and let $`x\mathrm{\Omega }`$. The horizontal jacobian of a $`C^1`$ mapping $`f:\mathrm{\Omega }^k`$ at $`x`$ is given by
$`J_Hf(x)=\mathrm{\Lambda }_k\left(df(x)_{|H_x^n}\right),`$ (17)
where $`\mathrm{\Lambda }_k\left(df(x)_{|H_x^n}\right):\mathrm{\Lambda }_k(H_x^n)\mathrm{\Lambda }_k(^k)`$.
From definition of horizontal jacobian, it follows that it only depends on the restriction of $`g`$ to the horizontal subbundle, namely, from the โsub-Riemannian metricโ. Let us consider a horizontal frame $`(Y_1,Y_2,\mathrm{},Y_{2n})`$, hence $`J_Hf(x)`$ is given by the jacobian of
$`_Yf(x)=\left[\begin{array}{ccccccc}Y_1f^1(x)& Y_2f^1(x)& \mathrm{}& Y_{2n}f^1(x)& & & \\ Y_1f^2(x)& Y_2f^2(x)& \mathrm{}& Y_{2n}f^2(x)& & & \\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& & & \\ Y_1f^k(x)& Y_2f^k(x)& \mathrm{}& Y_{2n}f^k(x)& & & \end{array}\right].`$ (22)
As a consequence, we have the formula
$`J_Hf(x)=\sqrt{{\displaystyle \underset{1i_1<i_2\mathrm{}<i_k2n}{}}\left[M_{i_1i_2\mathrm{}i_k}\left(_Yf(x)\right)\right]^2}.`$ (23)
###### Proposition 2.7
Let $`(Y_1,Y_2,\mathrm{},Y_{2n},W)`$ be an orthonormal frame with respect to a left invariant metric $`h`$ and let $`F:^{2n+1}^n`$ define coordinates with respect to this frame. Then we have $`F_{\mathrm{}}^{2n+1}=\text{vol}_{2n+1}`$, where $`\text{vol}_{2n+1}`$ denotes the Riemannian volume measure with respect to the metric $`h`$.
Proof. Let $`A`$ be a measurable set of $`^{2n+1}`$. By classical area formula and taking into account the left invariance of both $`\text{vol}_p`$ and $`F_{\mathrm{}}^{2n+1}`$ we have
$$c^{2n+1}(A)=\text{vol}_{2n+1}(F(A))=_AJ_hF(x)๐x$$
for some constant $`c>0`$. Then $`\text{ }_AJ_hF=c`$ for any measurable $`A`$. By continuity of $`xJ_hF(x)`$ we obtain that $`J_hF(x)=c`$ for any $`x^q`$. We have $`F=\mathrm{exp}L`$, with
$`L(y)=y_{2n+1}W+{\displaystyle \underset{i=1}{\overset{2n}{}}}y_jY_j`$
and $`(Y_1,Y_2,\mathrm{},Y_{2n},W)`$ is orthonormal. Since the map $`dF(0)=d\mathrm{exp}(0)L=L`$ has jacobian equal to one, then $`c=1`$ and the thesis follows. $`\mathrm{}`$
###### Remark 2.8
By previous proposition, the volume measure of a measurable subset $`A`$ of $`^n`$ corresponds to the $`(2n+1)`$-dimensional Lebesgue measure of the same subset read with respect to coordinates associated to an orthonormal frame. Here the volume measure is defined by the same left invariant metric.
Recall that $`(F_{\mathrm{}}\mu )(A)=\mu \left(F^1(A)\right)`$, where $`\mu `$ is a measure defined on the domain of $`F`$ and $`A`$ is a measurable set defined on the codomain of $`F`$. The $`d`$-dimensional spherical Hausdorff measure $`๐ฎ^d`$ is defined as
$`๐ฎ^d(A)=\underset{\epsilon >0}{sup}inf\{{\displaystyle \underset{j=1}{\overset{\mathrm{}}{}}}\text{diam}(E_j)^dA{\displaystyle E_j},\text{diam}(E_j)\epsilon \},`$ (24)
where the diameter is considered with respect to a homogeneous distance $`\rho `$ of $`^n`$ and we do not consider any dimensional factor. The $`k`$-dimensional Hausdorff measure built with respect to the Riemannian distance is denoted by $`\text{vol}_k`$ and it corresponds to the classical Riemannian volume measure with respect to the graded metric $`g`$, see for instance 3.2.46 of .
###### Definition 2.9 (Horizontal $`p`$-vectors)
For each $`x^n`$, we say that any linear combination of wedge products $`X_{j_1}(x)X_{j_2}(x)\mathrm{}X_{j_p}(x)`$, where $`1j_s2n`$ and $`j=1,\mathrm{},2n`$, is a horizontal $`p`$-vector. The space of horizontal $`p`$-vectors is denoted by $`\mathrm{\Lambda }_p(H_x^n)`$.
###### Definition 2.10 (Vertical $`p`$-vectors)
For each $`x^n`$, we say that any linear combination of wedge products $`X_{j_1}(x)X_{j_2}(x)\mathrm{}X_{j_{p1}}(x)Z(x)`$, where $`1j_s2n`$ and $`j=1,\mathrm{},2n`$, is a vertical $`p`$-vector. The space of vertical $`p`$-vectors is denoted by $`๐ฑ_p(H_x^n)`$.
For every couple of simple $`p`$-vectors $`v_1\mathrm{}v_p`$, $`w_1\mathrm{}w_p\mathrm{\Lambda }_p(T_x^n)`$, we define the scalar product induced by the left invariant Riemannian metric $`g`$ on $`T_x^n`$ as
$`v_1\mathrm{}v_p,w_1\mathrm{}w_p=det(\left(g(x)(v_i,w_j)\right),`$ (25)
see for instance 1.7.5 of for more details. This allows us to regard the space of vertical $`p`$-vectors $`๐ฑ_p(T_x^n)`$ as the orthogonal complement of the horizontal subspace $`\mathrm{\Lambda }_p(H_x^n)`$. We have the orthogonal decomposition
$`\mathrm{\Lambda }_p(T_x^n)=\mathrm{\Lambda }_p(H_x^n)๐ฑ_p(T_x^n),`$ (26)
which generalizes the case $`p=1`$, corresponding to $`T_x^n=H_x^nZ(x)`$.
###### Definition 2.11 (Vertical projection)
Let $`x^n`$ and let $`\xi \mathrm{\Lambda }_p(T_x^n)`$. The orthogonal decomposition $`\xi =\xi _H+\xi _๐ฑ`$ associated to (26) uniquely defines the vertical $`p`$-vector $`\xi _๐ฑ๐ฑ_p(T_x^n)`$. We say that $`\xi _๐ฑ`$ is the vertical projection of $`\xi `$ and that the mapping $`\pi _๐ฑ:\mathrm{\Lambda }_p(T_x^n)๐ฑ_p(T_x^n)`$, which associates $`\xi _๐ฑ`$ to $`\xi `$, is the vertical projection.
We have omitted $`x`$ in the definition of vertical projection $`\pi _๐ฑ`$.
###### Definition 2.12 (Characteristic points and transverse points)
Let $`\mathrm{\Sigma }\mathrm{\Omega }`$ be a $`C^1`$ submanifold and let $`x\mathrm{\Sigma }`$. We say that $`x\mathrm{\Sigma }`$ is a characteristic point if $`T_x\mathrm{\Sigma }H_x^n`$ and that it is a transverse point otherwise. The characteristic set of $`\mathrm{\Sigma }`$ is the subset of all characteristic points and it is denoted by $`C(\mathrm{\Sigma })`$.
Recall that a tangent $`p`$-vector to a $`p`$-dimensional submanifold $`\mathrm{\Sigma }`$ of class $`C^1`$ at $`x\mathrm{\Sigma }`$ is defined by the wedge product $`t_1t_2\mathrm{}t_p`$, where $`(t_1,\mathrm{},t_p)`$ is an orthonormal basis of $`T_x\mathrm{\Sigma }`$. We denote this simple $`p`$-vector by $`\tau _\mathrm{\Sigma }(x)`$. Notice that the tangent $`p`$-vector (which belongs to a one-dimensional space) cannot be continuously defined on all of $`\mathrm{\Sigma }`$, unless the submanifold is oriented.
###### Definition 2.13 (Vertical tangent $`p`$-vector)
Let $`\mathrm{\Sigma }\mathrm{\Omega }`$ be a $`p`$-dimensional submanifold of class $`C^1`$ and let $`x\mathrm{\Sigma }`$. A vertical tangent $`p`$-vector to $`\mathrm{\Sigma }`$ at $`x`$ is defined by $`\pi _๐ฑ(\tau _\mathrm{\Sigma })`$, where $`\tau _\mathrm{\Sigma }`$ is a tangent $`p`$-vector and $`\pi _๐ฑ`$ is the vertical projection. The vertical tangent $`p`$-vector will be denoted by $`\tau _{\mathrm{\Sigma },๐ฑ}(x)`$.
## 3 Blow-up at transverse points
This section is devoted to the proof of Theorem 1.1. In the following proposition, we give a simple characterization of characteristic points using vertical tangent $`p`$-vectors.
###### Proposition 3.1
Let $`\mathrm{\Sigma }\mathrm{\Omega }`$ be a submanifold of class $`C^1`$ and let $`x\mathrm{\Sigma }`$. Then $`xC(\mathrm{\Sigma })`$ if and only if $`\tau _{\mathrm{\Sigma },๐ฑ}(x)=0`$.
Proof. Let $`x\mathrm{\Sigma }`$ and let $`(t_1,t_2,\mathrm{},t_p)`$ be an orthonormal basis of $`T_x\mathrm{\Sigma }`$. We have the unique decomposition $`t_j=V_j+\gamma _jZ,`$ where $`V_jH_x^n`$ for every $`j=1,\mathrm{},p`$. It follows that
$`\tau =t_1t_2\mathrm{}t_p`$
$`=\left(V_1+\gamma _1Z\right)\left(V_2+\gamma _2Z\right)\mathrm{}\left(V_p+\gamma _pZ\right)`$
$`=V_1V_2\mathrm{}V_p+{\displaystyle \underset{j=1}{\overset{p}{}}}\gamma _jV_1V_2\mathrm{}V_{j1}ZV_{j+1}\mathrm{}V_p.`$
Assume that $`xC(\mathrm{\Sigma })`$. If $`V_1,V_2,\mathrm{},V_p`$ are linearly dependent, then we get
$`t_1t_2\mathrm{}t_p={\displaystyle \underset{j=1}{\overset{p}{}}}\gamma _jV_1V_2\mathrm{}V_{j1}ZV_{j+1}\mathrm{}V_p.`$
As a result, $`\pi _๐ฑ(\tau )=\tau `$ hence it is not vanishing. If $`V_1,\mathrm{},V_p`$ are linearly independent, then all wedge products of the form
$`V_1V_2\mathrm{}V_{j1}ZV_{j+1}\mathrm{}V_p`$ (27)
are non-vanishing for every $`j=1,\mathrm{},p`$. The fact that $`x`$ is transverse implies that there exists $`\gamma _{j_0}0`$, then the projection
$`\pi _๐ฑ(\tau )={\displaystyle \underset{j=1}{\overset{p}{}}}\gamma _jV_1V_2\mathrm{}V_{j1}ZV_{j+1}\mathrm{}V_p`$ (28)
is non-vanishing. Conversely, if $`\pi _๐ฑ(\tau )0`$, then (28) yields some $`\gamma _{j_1}0`$, therefore $`t_{j_1}H_x^n`$. $`\mathrm{}`$
###### Proposition 3.2
Let $`f:\mathrm{\Omega }^k`$ be of class $`C^1`$, with surjective differential at each point of $`\mathrm{\Omega }`$. Let $`\mathrm{\Sigma }`$ denote the submanifold $`f^1(0)`$ of $`\mathrm{\Omega }`$ and let $`x\mathrm{\Sigma }`$. Then $`xC(\mathrm{\Sigma })`$ if and only if $`df(x)_{|H_x^n}`$ is not surjective.
Proof. We first notice that $`\text{Ker }\left(df(x)_{|H_x^n}\right)=T_x\mathrm{\Sigma }H_x^n,`$ then we have
$`dim(H_x^nT_x\mathrm{\Sigma })=2ndim\left(\text{Im }(df(x)_{|H_x^n}\right)).`$ (29)
This last formula allows us to get our claim as follows. Assume that $`xC(\mathrm{\Sigma })`$. Then $`T_x\mathrm{\Sigma }H_x^n`$ and (29) gives
$$2n+1k=2ndim\left(\text{Im }(df(x)_{|H_x^n}\right)).$$
From this equation we conclude that $`df(x)_{|H_x^n}`$ is not surjective. Conversely, if $`df(x)_{|H_x^n}`$ is not surjective, then (29) implies
$`dim(H_x^nT_x\mathrm{\Sigma })2nk+1=dim(T_x\mathrm{\Sigma })`$
therefore $`T_x\mathrm{\Sigma }H_x^n`$, namely, $`xC(\mathrm{\Sigma })`$. $`\mathrm{}`$
###### Theorem 3.3
Let $`f:\mathrm{\Omega }^k`$ be of class $`C^1`$, with surjective differential at each point of $`\mathrm{\Omega }`$. Let $`\mathrm{\Sigma }`$ denote the submanifold $`f^1(0)`$ of $`\mathrm{\Omega }`$ and let $`x\mathrm{\Sigma }`$. Then we have
$`|\tau _{\mathrm{\Sigma },๐ฑ}(x)|={\displaystyle \frac{J_Hf(x)}{J_gf(x)}}.`$ (30)
Proof. Left invariance of Riemannian metric allows us to consider the left translated submanifold $`l_{x^1}\mathrm{\Sigma }`$. Replacing $`f`$ with $`fl_x`$ and $`\mathrm{\Omega }`$ with $`l_{x^1}\mathrm{\Omega }`$ we can assume that $`x`$ is the unit element $`0`$ of $`^n`$. Recall that $`l_x:^n^n`$ is the left translation $`l_x(y)=xy`$. If $`xC(\mathrm{\Sigma })`$, then Proposition 3.1 and Proposition 3.2 make (30) the trivial identity $`0=0`$. Assume that $`x\mathrm{\Sigma }C(\mathrm{\Sigma })`$. Then Proposition 3.2 implies that the horizotal gradients
$$_Hf^i=(X_1f^i(0),X_2f^i(0),\mathrm{},X_{2n}f^i(0))\text{for}i=1,2,\mathrm{},k$$
span a $`k`$-dimensional space of $`^{2n}`$. Let $`c_1,c_2,\mathrm{},c_k^{2n}`$ be orthogonal unit vectors generating this vector space and choose $`c_{k+1},\mathrm{},c_{2n}^{2n}`$ such that $`(c_1,c_2,\mathrm{},c_{2n})`$ is an orthonormal basis of $`^{2n}`$. These vectors allow us to define a new horizontal frame
$`Y_j={\displaystyle \underset{k=1}{\overset{2n}{}}}c_j^kX_k\text{for every}j=1,\mathrm{},2n.`$ (31)
We denote by $`C`$ the $`2n\times 2n`$ orthogonal matrix whose $`i`$-th column corresponds to the vector $`c_i`$, then by our choice of vectors $`c_j`$, we obtain $`_Yf(x)=_Xf(x)C`$ and
$`_Yf(x)C=\left[\begin{array}{ccccccc}_Hf^1,c_1& _Hf^1,c_2& \mathrm{}& _Hf^1,c_k& 0& \mathrm{}& 0\\ _Hf^2,c_1& _Hf^2,c_2& \mathrm{}& _Hf^2,c_k& 0& \mathrm{}& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& 0& \mathrm{}& 0\\ _Hf^k,c_1& _Hf^k,c_2& \mathrm{}& _Hf^k,c_k& 0& \mathrm{}& 0\end{array}\right],`$ (36)
where the symbol $`,`$ denotes the standard scalar product of $`^{2n}`$. Let us consider $`F:^{2n+1}^n`$, defining graded coordinates $`(y_1,\mathrm{},y_{2n+1})`$ associated to the frame $`(Y_1,\mathrm{},Y_{2n},Z)`$, according to Definition 2.2. Then the differential of $`f`$ at $`0`$ with respect to $`(y_1,\mathrm{},y_{2n+1})`$ can be represented by the matrix
$`_yf(0)=\left[\begin{array}{cccccccc}f_{y_1}^1(0)& f_{y_2}^1(0)& \mathrm{}& f_{y_k}^1(0)& 0& \mathrm{}& 0& f_{y_{2n+1}}^1(0)\\ f_{y_1}^2(0)& f_{y_2}^2(0)& \mathrm{}& f_{y_k}^1(0)& 0& \mathrm{}& 0& f_{y_{2n+1}}^2(0)\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& 0& \mathrm{}& 0& \mathrm{}\\ f_{y_1}^k(0)& f_{y_2}^k(0)& \mathrm{}& f_{y_k}^k(0)& 0& \mathrm{}& 0& f_{y_{2n+1}}^k(0)\end{array}\right].`$ (41)
It follows that
$$f_{y_j}^i(0)=_Hf^i,c_j$$
for every $`i,j=1,\mathrm{},2n`$. The implicit function theorem gives us a $`C^1`$ map $`\phi :A^k`$ such that $`A^p`$ is an open neighbourhood of the origin and
$`f(\phi ^1(\stackrel{~}{y}),\mathrm{},\phi ^k(\stackrel{~}{y}),y_{k+1},\mathrm{},y_{2n+1})=0`$ (42)
for every $`\stackrel{~}{y}=(y_{k+1},\mathrm{},y_{2n+1})A`$. Then we define the mapping $`\varphi :A^{2n+1}`$ as
$`\varphi (\stackrel{~}{y})=(\phi ^1(\stackrel{~}{y}),\mathrm{},\phi ^k(\stackrel{~}{y}),y_{k+1},\mathrm{},y_{2n+1}),`$ (43)
so that differentiating (42) we get
$`0=_{y_j}(f^i\varphi )={\displaystyle \underset{l=1}{\overset{k}{}}}f_{y_l}^i\phi _{y_j}^l+f_{y_j}^i`$ (44)
for every $`i=1,\mathrm{},k`$ and $`j=k+1,\mathrm{},2n+1`$. Equations (44) can be more concisely written in matrix form as follows
$`_zf\phi _{y_j}=f_{y_j},`$ (45)
where $`z=(y_1,\mathrm{},y_k)`$, the $`k\times k`$ matrix $`_zf`$ has coefficients $`f_{y_l}^i`$, where $`i,l=1,\mathrm{},k`$ and $`j=k+1,\mathrm{},2n+1`$. In order to achieve a more explicit formula for the differential of the implicit map, we explicitly write the inverse matrix of $`_zf`$ as
$`\left(_zf\right)^1={\displaystyle \frac{1}{M_{12\mathrm{}k}(_zf)}}\left[\begin{array}{cccc}C_{11}(_zf)& C_{21}(_zf)& \mathrm{}& C_{k1}(_zf)\\ C_{12}(_zf)& C_{22}(_zf)& \mathrm{}& C_{k1}(_zf)\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ C_{1k}(_zf)& C_{2k}(_zf)& \mathrm{}& C_{kk}(_zf)\end{array}\right],`$
where $`C_{ij}(_zf)`$ denotes the cofactor of $`_zf`$, which is equal to $`(1)^{i+j}det(\text{}_{ij}f)`$ and $`\text{}_{ij}f`$ is the $`(k1)\times (k1)`$ square matrix obtained by removing the $`i`$-th row and the $`j`$-th column from $`_zf`$. In view of (45) we have
$`\phi _{y_j}=\left(_zf\right)^1f_{y_j}={\displaystyle \frac{1}{M_{12\mathrm{}k}(_zf)}}\left[\begin{array}{c}_{i=1}^kC_{i1}(_zf)f_{y_j}^i\\ _{i=1}^kC_{i2}(_zf)f_{y_j}^i\\ \mathrm{}\\ _{i=1}^kC_{ik}(_zf)f_{y_j}^i\end{array}\right].`$
An elementary formula for computing the determinant of a matrix implies
$`{\displaystyle \underset{i=1}{\overset{k}{}}}C_{is}(_zf)f_{y_j}^i=M_{12\mathrm{}s1js+1\mathrm{}k}(_yf)`$
for every $`j=k+1,\mathrm{},2n+1`$. As a consequence, we get
$`\phi _{y_j}^s={\displaystyle \frac{M_{12\mathrm{}s1js+1\mathrm{}k}(_yf)}{M_{1\mathrm{}k}(_zf)}}.`$ (48)
Note that $`M_{1\mathrm{}k}(_zf)`$ corresponds to the determinant of the matrix $`_zf`$. As a consequece of (48) and of (41), we conclude that
$`\phi _{y_j}^s(0)=0`$
for every $`j=k+1,\mathrm{},2n`$. Previous considerations and expression (43) lead us to the formula
$`_{\stackrel{~}{y}}\varphi (0)=\left[\begin{array}{ccccc}0& 0& \mathrm{}& 0& \phi _{y_{2n+1}}^1(0)\\ 0& 0& \mathrm{}& 0& \phi _{y_{2n+1}}^2(0)\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& \mathrm{}& 0& \phi _{y_{2n+1}}^k(0)\\ 1& 0& \mathrm{}& 0& 0\\ 0& 1& \mathrm{}& 0& 0\\ \mathrm{}& 0& \mathrm{}& 0& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& 1& \mathrm{}\\ 0& 0& \mathrm{}& 0& 1\end{array}\right],`$ (58)
where $`_{\stackrel{~}{y}}\varphi (0)`$ is a $`(2n+1)\times p`$ matrix whose $`p\times p`$ lower block is the identity matrix. Notice that columns of (58) represent a basis of the tangent space $`T_0\mathrm{\Sigma }`$ with respect to coordinates $`(y_{k+1},\mathrm{},y_{2n+1})`$. More precisely, the set of vectors
$`(Y_{k+1}(0),Y_{k+2}(0),\mathrm{},Y_{2n}(0),{\displaystyle \frac{Z(0)+_{j=1}^kv_jY_j(0)}{\left(1+_{j=1}^kv_j^2\right)^{1/2}}})`$
form an orthonormal basis of $`T_0\mathrm{\Sigma }`$, where we have defined $`v_j=\phi _{y_{2n+1}}^j(0)`$. Then the tangent $`p`$-vector $`\tau _\mathrm{\Sigma }`$ to $`\mathrm{\Sigma }`$ at $`0`$ is given by the wedge product
$`\tau _\mathrm{\Sigma }(0)={\displaystyle \frac{Y_{k+1}(0)Y_{k+2}(0)\mathrm{}Y_{2n}(0)\left(Z(0)+_{j=1}^kv_jY_j(0)\right)}{\left(1+_{j=1}^kv_j^2\right)^{1/2}}}.`$
Obviously, $`p`$-vectors $`Y_{k+1}(0)Y_{k+2}(0)\mathrm{}Y_{2n}(0)Y_j(0)`$ are horizontal, hence they disappear in the vertical projection. It follows that
$`\tau _{\mathrm{\Sigma },๐ฑ}(0)=\pi _๐ฑ\left(\tau _{\mathrm{\Sigma },๐ฑ}(0)\right)={\displaystyle \frac{Y_{k+1}(0)Y_{k+2}(0)\mathrm{}Y_{2n}(0)Z(0)}{\left(1+_{j=1}^kv_j^2\right)^{1/2}}},`$
therefore we clearly obtain
$`|\tau _{\mathrm{\Sigma },๐ฑ}(0)|=\left(1+{\displaystyle \underset{l=1}{\overset{k}{}}}v_l^2\right)^{1/2}.`$ (59)
Due to formula (48) in the case $`j=2n+1`$ and to (41), we obtain
$`1+{\displaystyle \underset{l=1}{\overset{k}{}}}v_l^2={\displaystyle \frac{\left(M_{1\mathrm{}k}(_zf)\right)^2+_{s=1}^k\left(M_{12\mathrm{}s\mathrm{1\hspace{0.17em}2}n+1s+1\mathrm{}k}(_yf)\right)^2}{\left(M_{1\mathrm{}k}(_zf)\right)^2}}=\left({\displaystyle \frac{J_gf(0)}{J_Hf(0)}}\right)^2,`$
then (59) shows the valdity of (30) in the case $`x=0`$. Left invariance of the Riemannian metric $`g`$ leads us to the conclusion. $`\mathrm{}`$
###### Definition 3.4 (Metric factor)
Let $`\tau `$ be a vertical simple $`p`$-vector of $`\mathrm{\Lambda }_p(๐ฅ^n)`$ and let $`(\tau )`$ be the unique associated subspace, with $`L=\mathrm{exp}(\tau )`$. The metric factor of a homogeneous distance $`\rho `$ with respect to $`\tau `$ is defined by
$`\theta _p^\rho (\tau )=_{||}^p\left(F^1(LB_1)\right),`$
where $`F:^{2n+1}^n`$ defines a system of graded coordinates, $`_{||}^p`$ denotes the $`p`$-dimensional Hausdorff measure with respect to the Euclidean distance of $`^{2n+1}`$ and $`B_1`$ is the unit ball of $`^n`$ with respect to the distance $`\rho `$. Recall that the subspace associated to a simple $`p`$-vector $`\tau `$ is defined as $`\{v๐ฅ^nv\tau =0\}`$.
###### Remark 3.5
In the case of subspaces $``$ of codimension one, the notion of metric factor fits into the one introduced in . It is easy to observe that the notion of metric factor does not depend on the system of coordinates we are using. In fact, $`F_1^1F_2:^{2n+1}^{2n+1}`$ is an Euclidean isometry whenever $`F_1,F_2:^{2n+1}^n`$ represent systems of graded coordinates with respect to the same left invariant Riemannian metric.
Proof of Theorem 1.1. As in the proof of Theorem 3.3, left invariance of the Riemannian metric $`g`$ allows us to assume that $`x=0`$. For $`r_0>0`$ sufficiently small, we can suppose the existence of a function $`f:B_{r_0}^k`$ such that $`\mathrm{\Sigma }B_{r_0}=f^1(0)`$ and whose differential is surjective at every point of $`B_{r_0}`$. By Proposition 3.2, the horizontal gradients
$$_Hf^i=(X_1f^i(0),X_2f^i(0),\mathrm{},X_{2n}f^i(0))\text{for}i=1,2,\mathrm{},k$$
span a $`k`$-dimensional space of $`^{2n}`$. Now, repeating the argument in the proof of Theorem 3.3, we define the system of graded coordinates $`(y_1,\mathrm{},y_{2n+1})`$ associated to the frame $`(Y_1,\mathrm{},Y_{2n},Z)`$, where $`Y_j`$ are given by (31). The differential of $`f`$ at $`0`$ can be represented by the matrix
$`_yf(0)=\left[\begin{array}{cccccccc}f_{y_1}^1(0)& f_{y_2}^1(0)& \mathrm{}& f_{y_k}^1(0)& 0& \mathrm{}& 0& f_{y_{2n+1}}^1(0)\\ f_{y_1}^2(0)& f_{y_2}^2(0)& \mathrm{}& f_{y_k}^1(0)& 0& \mathrm{}& 0& f_{y_{2n+1}}^2(0)\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& 0& \mathrm{}& 0& \mathrm{}\\ f_{y_1}^k(0)& f_{y_2}^k(0)& \mathrm{}& f_{y_k}^k(0)& 0& \mathrm{}& 0& f_{y_{2n+1}}^k(0)\end{array}\right],`$ (64)
whose first $`k`$ columns are linearly independent. By the implicit function theorem ther eexists a $`C^1`$ mapping $`\phi :A^k`$ such that $`A^p`$ is an open neighbourhood of the origin and
$`f(\phi ^1(\stackrel{~}{y}),\mathrm{},\phi ^k(\stackrel{~}{y}),y_{k+1},\mathrm{},y_{2n+1})=0`$ (65)
for every $`\stackrel{~}{y}=(y_{k+1},\mathrm{},y_{2n+1})A`$. Proceeding as in the proof of Theorem 3.3, we define the mapping $`\varphi :A^{2n+1}`$ as
$`\varphi (\stackrel{~}{y})=(\phi ^1(\stackrel{~}{y}),\mathrm{},\phi ^k(\stackrel{~}{y}),y_{k+1},\mathrm{},y_{2n+1}),`$ (66)
and by the same computations, differentiating (65) we obtain
$`_{\stackrel{~}{y}}\varphi (0)=\left[\begin{array}{ccccc}0& 0& \mathrm{}& 0& \phi _{y_{2n+1}}^1(0)\\ 0& 0& \mathrm{}& 0& \phi _{y_{2n+1}}^2(0)\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& \mathrm{}& 0& \phi _{y_{2n+1}}^k(0)\\ 1& 0& \mathrm{}& 0& 0\\ 0& 1& \mathrm{}& 0& 0\\ \mathrm{}& 0& \mathrm{}& 0& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& 1& \mathrm{}\\ 0& 0& \mathrm{}& 0& 1\end{array}\right],`$ (76)
where $`_{\stackrel{~}{y}}\varphi (0)`$ is a $`(2n+1)\times p`$ matrix whose $`p\times p`$ lower block is the identity matrix. For each $`r<r_0`$, write the ball $`B_r`$ in terms of graded coordinates defining $`\stackrel{~}{B}_r=F^1(B_r)^{2n+1}`$. The surface $`\mathrm{\Sigma }`$ read in graded coordinates can be seen as the image of $`\varphi `$. Then we have established
$`{\displaystyle \frac{\text{vol}_p(\mathrm{\Sigma }B_r)}{r^{p+1}}}=r^{1p}{\displaystyle _{\varphi ^1(\stackrel{~}{B}_r)}}J_g\varphi (\stackrel{~}{y})๐\stackrel{~}{y}.`$ (77)
The dilation $`\delta _r`$ restricted to coordinates $`(y_1,\mathrm{},y_{2n+1})`$ gives
$`\delta _r\stackrel{~}{y}=\delta _r\left((y_{k+1},\mathrm{},y_{2n+1})\right)=(ry_{k+1},ry_{k+2},\mathrm{},ry_{2n},r^2y_{2n+1}),`$ (78)
therefore, performing a change of variable in (77) we get
$`{\displaystyle \frac{\text{vol}_p(\mathrm{\Sigma }B_{x,r})}{r^{p+1}}}={\displaystyle _{\delta _{1/r}\left(\varphi ^1(\stackrel{~}{B}_r)\right)}}J_g\varphi (\delta _r\stackrel{~}{y})๐\stackrel{~}{y}.`$ (79)
The set $`\delta _{1/r}\left(\varphi ^1(\stackrel{~}{B}_r)\right)`$ can be written as follows
$`(\delta _{1/r}\varphi \delta _r)^1(\stackrel{~}{B}_1)=\left\{\stackrel{~}{y}^p\right|({\displaystyle \frac{\phi ^1(\delta _r\stackrel{~}{y})}{r}},\mathrm{},{\displaystyle \frac{\phi ^k(\delta _r\stackrel{~}{y})}{r}},y_{k+1},\mathrm{},y_{2n+1})\stackrel{~}{B}_1\}.`$ (80)
From expressions (76) and (78) one easily gets that
$`\underset{r0^+}{lim}{\displaystyle \frac{\phi ^j(\delta _r\stackrel{~}{y})}{r}}=0`$ (81)
for every $`j=1,\mathrm{},k`$. As a result, the limit
$`\mathrm{๐}_{\delta _{1/r}\left(\varphi ^1(\stackrel{~}{B}_r)\right)}\mathrm{๐}_{\stackrel{~}{B}_1\mathrm{\Pi }}\text{as}r0^+`$ (82)
holds a.e. in $`^p`$, where we have defined
$`\mathrm{\Pi }=\{(0,\mathrm{},0,y_{k+1},\mathrm{},y_{2n+1})^{2n+1}y_j,j=k+1,\mathrm{},2n+1\}.`$
From (79), we conclude that
$`\underset{r0^+}{lim}{\displaystyle \frac{\text{vol}_p(\mathrm{\Sigma }B_{x,r})}{r^{p+1}}}=J_g\varphi (0)^p(\mathrm{\Pi }\stackrel{~}{B}_1).`$ (83)
To compute $`J_g\varphi (0)`$, we use both the canonical form of the tangent space $`T_0\mathrm{\Sigma }`$ given by (76) and the fact that our frame $`(Y_1,\mathrm{},Y_{2n},Z)`$ is orthonormal. Thus, according to Definition 2.4 the Riemannian jacobian of $`\varphi `$ at zero is given by
$`J_g\varphi (0)=\left(1+{\displaystyle \underset{j=1}{\overset{k}{}}}v_j^2\right)^{1/2},`$ (84)
where we have defined $`v_j=\phi _{y_{2n+1}}^j(0)`$ for every $`j=1,\mathrm{},2n`$. Again, following the same steps of the proof of Theorem 3.3, we get
$`|\tau _{\mathrm{\Sigma },๐ฑ}(0)|=\left(1+{\displaystyle \underset{l=1}{\overset{k}{}}}v_l^2\right)^{1/2}=\left(J_g\varphi (0)\right)^1.`$
Then (83) yields
$`\underset{r0^+}{lim}{\displaystyle \frac{\text{vol}_p(\mathrm{\Sigma }B_r)}{r^{p+1}}}={\displaystyle \frac{^p(\mathrm{\Pi }\stackrel{~}{B}_1)}{|\tau _{\mathrm{\Sigma },๐ฑ}(0)|}}.`$ (85)
The subspace $`\left(\tau _{\mathrm{\Sigma },๐ฑ}(x)\right)`$ associated to the $`p`$-vector $`\tau _{\mathrm{\Sigma },๐ฑ}(x)`$ satisfies the relation
$$\mathrm{exp}\left(\left(\tau _{\mathrm{\Sigma },๐ฑ}(x)\right)\right)=F(\mathrm{\Pi })$$
therefore the metric factor of $`\rho `$ with respect to $`\tau _{\mathrm{\Sigma },๐ฑ}(x)`$ is $`^p(\mathrm{\Pi }\stackrel{~}{B}_1)`$. This fact along with (85) implies the validity of (1) and ends the proof. $`\mathrm{}`$
## 4 Spherical Hausdorff measure of submanifolds
This section deals with various applications of Theorem 1.2. A key result to obtain this theorem is the $`๐ฎ^{Qk}`$-negligibility of characteristic points of a $`k`$-codimensional submanifold of a Carnot group of Hausdorff dimension $`Q`$, see . This result in the case of Heisenberg groups reads as follows.
###### Theorem 4.1
Let $`\mathrm{\Sigma }\mathrm{\Omega }`$ be a $`C^1`$ submanifold of dimension $`p`$. Then the set of characteristic points $`C(\mathrm{\Sigma })`$ is $`๐ฎ^{p+1}`$-negligible.
###### Remark 4.2
In order to apply the negligibility result of one has to check that the notion of characteristic point in arbitrary stratified groups coincides with our definition stated in the Heisenberg group. According to a point $`x\mathrm{\Sigma }`$ is characteristic if
$`dim\left(H_x^n\right)dim\left(T_x\mathrm{\Sigma }H_x^n\right)k1.`$ (86)
If $`x`$ is characteristic according to Definition 2.12, then
$$dim\left(T_x\mathrm{\Sigma }H_x^n\right)=p=2n+1k$$
and (86) holds. Conversely, if (86) holds, then
$$p=dim(T_x\mathrm{\Sigma })=2nk+1dim\left(T_x\mathrm{\Sigma }H_x^n\right),$$
hence $`T_x\mathrm{\Sigma }H_x^n`$.
###### Corollary 4.3
Let $`\mathrm{\Sigma }\mathrm{\Omega }`$ be a $`C^1`$ submanifold of dimension $`p`$. Then we have
$`{\displaystyle _\mathrm{\Sigma }}\theta \left(\tau _{\mathrm{\Sigma },๐ฑ}(x)\right)๐๐ฎ^{p+1}(x)={\displaystyle _\mathrm{\Sigma }}|\tau _{\mathrm{\Sigma },๐ฑ}(x)|๐\text{vol}_p(x)`$ (87)
Proof. We apply Theorem 2.10.17(2) and Theorem 2.10.18(1) of , hence from limit (1) and Theorem 4.1 the proof follows by a standard argument. $`\mathrm{}`$
###### Remark 4.4
Proof of Theorem 1.2 immediately follows from (87).
Next, we present a class of homogeneous distances in the Heisenberg group which possess constant metric factor. The standard system of graded coordinates $`F:^{2n+1}^n`$ induced by $`(X_1,\mathrm{},X_{2n},Z)`$ will be understood in the sequel. To simplify notation we will write $`x=F(\stackrel{~}{x},x_{2n+1})^n`$, with $`\stackrel{~}{x}=(x_1,\mathrm{},x_{2n})^{2n}`$.
###### Proposition 4.5
Let $`F:^{2n+1}^n`$ define standard coordintates and let $`\rho `$ be a homogeneous distance of $`^n`$ such that $`\rho (0,F()):^{2n+1}`$ only depends on $`(|\stackrel{~}{x}|,x_{2n+1})`$. Then $`\theta _p^\rho (\tau )=\theta _p^\rho (\stackrel{~}{\tau })`$ whenever $`\tau ,\stackrel{~}{\tau }`$ are vertical simple $`p`$-vectors.
Proof. Let $`\tau =U_1\mathrm{}U_{p1}Z`$ and $`\stackrel{~}{\tau }=W_1\mathrm{}W_{p1}Z`$ be vertical simple $`p`$-vectors, where it is not restrictive assuming that both $`(U_1,\mathrm{},U_{p1},Z)`$ and $`(W_1,\mathrm{},W_{p1},Z)`$ are orthonormal systems of $`๐ฅ^{2n+1}`$. Then we easily find an isometry $`J:๐ฅ^{2n+1}๐ฅ^{2n+1}`$ such that $`J\left((\tau )\right)=(\stackrel{~}{\tau })`$ and $`J(Z)=Z`$. Recall that our graded coordinates are defined by $`F=\mathrm{exp}`$, where $`:^{2n+1}๐ฅ^{2n+1}`$ is an isometry such that
$$(x_1,\mathrm{},x_{2n+1})=x_{2n+1}Z+\underset{j=1}{\overset{2n}{}}x_jX_j$$
for every $`(x_1,\mathrm{},x_{2n+1})^{2n+1}`$. Thus, defining $`\stackrel{~}{B}_1=F^1(B_1)^{2n+1}`$, we have
$`F^1(\mathrm{exp}(\stackrel{~}{\tau })B_1)=^1J\left((\tau )\right)\stackrel{~}{B}_1=\phi (^1((\tau )\left)\right)\stackrel{~}{B}_1,`$ (88)
where $`\phi =^1J:^{2n+1}^{2n+1}`$ is an Euclidean isometry such that $`\phi (e_{2n+1})=e_{2n+1}`$ and $`e_{2n+1}`$ is the $`(2n`$$`+`$$`1)`$-th vector of the canonical basis of $`^{2n+1}`$. Then $`|\stackrel{~}{x}|=|\stackrel{~}{y}|`$ whenever $`\phi (\stackrel{~}{x},t)=(\stackrel{~}{y},t)`$. As a result, the fact that $`\rho (0,F(\stackrel{~}{x},t))`$ only depends on $`(|\stackrel{~}{x}|,t)`$ easily implies that $`\phi (\stackrel{~}{B}_1)=\stackrel{~}{B}_1`$. Thus, due to (88), it follows that
$`\theta _p^\rho (\tau )=_{||}^p(^1\left((\tau )\right)\stackrel{~}{B}_1)=_{||}^p(\phi (^1((\tau )\left)\right)\stackrel{~}{B}_1)=\theta _p^\rho (\stackrel{~}{\tau }).`$ (89)
This ends the proof. $`\mathrm{}`$
###### Example 4.6
An example of homogeneous distance satisfying hypotheses of Proposition 4.5 is the gauge distance, also called Korรกnyi distance, . The gauge distance from $`x`$ to the origin is given by
$$d(x,0)=\left(|\stackrel{~}{x}|^4+16x_{2n+1}^2\right)^{1/4},$$
where $`x=(\stackrel{~}{x},x_{2n+1})`$. Then we define $`d(x,y)=d(0,x^1y)`$, for any $`x,y^n`$. Another example of homogeneous distance with this property is the โmaximum distanceโ, defined by
$$d_{\mathrm{}}(x,0)=\mathrm{max}\{|\stackrel{~}{x}|,|x_{2n+1}|^{1/2}\}.$$
Due to Proposition 4.5, both of these distances have constant metric factor.
###### Remark 4.7
The metric factor depends on the Riemannian metric $`g`$ we have fixed. Furthermore, if we divide it by the volume of the unit ball $`B_1`$ (or any other fixed subset of positive measure), then we obtain a number only depending on the restriction of the Riemannian metric to $`H^n`$.
###### Lemma 4.8
Let $`\stackrel{~}{g}`$ be a left invariant Riemannian metric such that $`\stackrel{~}{g}_{|H^n}=g_{|H^n}`$. Then $`\theta _p^\rho (\tau )/\text{vol}_p(B_1)=\stackrel{~}{\theta }_p^\rho (\tau )/\stackrel{~}{\text{vol}}_p(B_1)`$ for any simple vertical $`p`$-vector, where $`\stackrel{~}{\text{vol}}_p`$ and $`\stackrel{~}{\theta }_p^\rho (\tau )`$ are defined with respect to the metric $`\stackrel{~}{g}`$.
Proof. By hypothesis, we can choose an orthonormal frame $`(X_1,X_2,\mathrm{},X_{2n},W)`$ with respect to $`\stackrel{~}{g}`$, where $`W=\lambda Z+_{j=1}^{2n}a_jX_j`$ and $`\lambda 0`$. Let $`F,\stackrel{~}{F}:^{2n+1}^n`$ represent system of coordinates with respect to the standard basis $`(X_1,X_2,\mathrm{},X_{2n},Z)`$ and $`(X_1,X_2,\mathrm{},X_{2n},W)`$, respectively. We have $`\stackrel{~}{F}=FT`$, where $`T:^{2n+1}^{2n+1}`$ is given by the matrix
$`A=\left[\begin{array}{cccccc}1& 0& 0& \mathrm{}& 0& a_1\\ 0& 1& 0& \mathrm{}& 0& a_2\\ \mathrm{}& \mathrm{}& 1& \mathrm{}& \mathrm{}& a_3\\ \mathrm{}& \mathrm{}& 0& \mathrm{}& 0& \mathrm{}\\ 0& \mathrm{}& \mathrm{}& 0& 1& a_{2n}\\ 0& 0& \mathrm{}& 0& 0& \lambda \end{array}\right]`$ (96)
By Proposition 2.7, it follows that
$`\stackrel{~}{\text{vol}}_p=\stackrel{~}{F}_{\mathrm{}}^{2n+1}=|detA|^1F_{\mathrm{}}^{2n+1}=|\lambda |^1\text{vol}_p.`$ (97)
Let $`\tau =U_1\mathrm{}U_{p1}Z`$, where it is not restrictive to assume that the horizontal vectors $`U_1,\mathrm{},U_{p1}`$ are orthonormal with respect to both $`g`$ and $`\stackrel{~}{g}`$. We denote by $``$ the subspace $`\text{span}\{U_1,\mathrm{},U_{p1},Z\}`$ of $`๐ฅ^n`$. Recall that
$`\stackrel{~}{\theta }_p^\rho (\tau )=_{||}^p\left(\stackrel{~}{F}^1(\mathrm{exp}()B_1)\right)=_{||}^p\left(T^1\left(F^1(\mathrm{exp}()B_1)\right)\right).`$ (98)
Now we wish to determine a basis of the subspace $`F^1(\mathrm{exp}()^{2n+1}`$. The relations $`U_i=_{j=1}^{2n}c_i^jX_j`$ give rise to $`p1`$ orthonormal vectors $`c_1,\mathrm{},c_{p1}^{2n}\times \{0\}`$ with respect to the Euclidean scalar product such that
$`\text{span}\{c_1,\mathrm{},c_{p1},e_{2n+1}\}=F^1\left(\mathrm{exp}\left(\right)\right)^{2n+1}.`$ (99)
In order to compute the Euclidean jacobian of $`T^1`$ restricted to the $`p`$-dimensional subspace $`\text{span}\{c_1,\mathrm{},c_{p1},e_{2n+1}\}^{2n+1}`$ we write the matrix
$`A^1=\left[\begin{array}{cccccc}1& 0& 0& \mathrm{}& 0& \lambda ^1a_1\\ 0& 1& 0& \mathrm{}& 0& \lambda ^1a_2\\ \mathrm{}& \mathrm{}& 1& \mathrm{}& \mathrm{}& \lambda ^1a_3\\ \mathrm{}& \mathrm{}& 0& \mathrm{}& 0& \mathrm{}\\ 0& \mathrm{}& \mathrm{}& 0& 1& \lambda ^1a_{2n}\\ 0& 0& \mathrm{}& \mathrm{}& 0& \lambda ^1\end{array}\right]`$ (106)
noting that $`T^1c_j=c_j`$ for every $`j=1,\mathrm{},p1`$ and $`T^1(e_{2n+1})=\lambda ^1e_{2n+1}`$. As a result, the jacobian of $`\left(T^1\right)_|:\text{span}\{c_1,\mathrm{},c_{p1},e_{2n+1}\}^{2n+1}`$ is $`|\lambda |^1`$, hence
$`\stackrel{~}{\theta }_p^\rho (\tau )=_{||}^p\left(T^1\left(F^1(\mathrm{exp}()B_1)\right)\right)`$
$`=|\lambda |^1_{||}^p((F^1(\mathrm{exp}()B_1))=|\lambda |^1\theta _p^\rho (\tau ).`$
Joining (97) and the previous equalities, our claim follows. $`\mathrm{}`$.
###### Theorem 4.9
Let $`\mathrm{\Sigma }`$ be a $`p`$-dimensional submanifold of $`\mathrm{\Omega }`$ and let $`\stackrel{~}{g}`$ a left invariant metric such that $`\stackrel{~}{g}_{|H^n}=g_{|H^n}`$. Then
$`{\displaystyle \frac{1}{\stackrel{~}{\text{vol}}_p(B_1)}}{\displaystyle _\mathrm{\Sigma }}|\stackrel{~}{\tau }_{\mathrm{\Sigma },๐ฑ}(x)|๐\stackrel{~}{\text{vol}}_p(x)={\displaystyle \frac{1}{\text{vol}_p(B_1)}}{\displaystyle _\mathrm{\Sigma }}|\tau _{\mathrm{\Sigma },๐ฑ}(x)|๐\text{vol}_p(x)`$ (107)
The previous theorem is an immediate conseqauence of Corollary 4.3 and Lemma 4.8. Next, we apply (2) to compute the spherical Hausdorff measure of some submanifolds. We will use the following proposition.
###### Proposition 4.10
Let $`\varphi :U^{2n+1}`$ be a $`C^1`$ embedding, where $`U^p`$ is a bounded open set. Let $`F:^{2n+1}^n`$ define standard coordinates and set $`\mathrm{\Phi }=F\varphi :U^n`$, where $`\mathrm{\Sigma }=\mathrm{\Phi }(U)`$. Let $`\rho `$ be a homogeneous distance with constant metric factor $`\alpha >0`$. Then we have
$`๐ฎ_^n^{p+1}(\mathrm{\Sigma })={\displaystyle _U}\left|\pi _๐ฑ\left(\mathrm{\Phi }_{u_1}(u)\mathrm{\Phi }_{u_2}(u)\mathrm{}\mathrm{\Phi }_{u_p}(u)\right)\right|๐u,`$ (108)
for every measurable set $`A^n`$, where the norm $`||`$ is induced by the scalar product (25) on $`p`$-vectors.
Proof. By definition of Riemannian volume, formula (2) can be written with respect to $`\varphi `$ as
$`๐ฎ_^n^{p+1}(\mathrm{\Sigma })={\displaystyle _U}|\tau _{\mathrm{\Sigma },๐ฑ}(\varphi (u))|\sqrt{det\left[g(\varphi (u))(\mathrm{\Phi }_{u_i}(u),\mathrm{\Phi }_{u_j}(u))\right]}๐u,`$ (109)
where we have
$`\left|\mathrm{\Phi }_{u_1}(u)\mathrm{\Phi }_{u_2}(u)\mathrm{}\mathrm{\Phi }_{u_p}(u)\right|=\sqrt{det\left[g(\mathrm{\Phi }(u))(\mathrm{\Phi }_{u_i}(u),\mathrm{\Phi }_{u_j}(u))\right]}.`$ (110)
Therefore, taking into account the formula
$`\tau _\mathrm{\Sigma }\left(\mathrm{\Phi }(u)\right)={\displaystyle \frac{\mathrm{\Phi }_{u_1}(u)\mathrm{\Phi }_{u_2}(u)\mathrm{}\mathrm{\Phi }_{u_p}(u)}{|\mathrm{\Phi }_{u_1}(u)\mathrm{\Phi }_{u_2}(u)\mathrm{}\mathrm{\Phi }_{u_p}(u)|}},`$ (111)
the definition of vertical tangent $`p`$-vector $`\pi _๐ฑ(\tau _\mathrm{\Sigma })=\tau _{\mathrm{\Sigma },๐ฑ}`$ and joining (109), (110) and (111), formula (108) follows. $`\mathrm{}`$
###### Example 4.11
Let $`\varphi :^3^5`$, defined by $`\varphi (u)=(u_1,u_2,u_3,0,\frac{u_1^2+u_2^2+u_3^2}{2})`$. The mapping $`\varphi `$ parametrizes a 3-dimensional paraboloid $`\mathrm{\Sigma }=\mathrm{\Phi }(U)`$ of $`^5`$, where $`U`$ is an open bounded set of $`^3`$, $`\mathrm{\Phi }=F\varphi `$ and $`F:^3^2`$ represents standard coordinates, according to Definition 2.2. Using expressions (7), we have
$`\varphi _{u_1}(u)=\stackrel{~}{X}_1(\varphi (u))+\left(\varphi _3(u)+u_1\right)\stackrel{~}{T}(\varphi (u)),`$
$`\varphi _{u_2}(u)=\stackrel{~}{X}_2(\varphi (u))+\left(\varphi _4(u)+u_2\right)\stackrel{~}{T}(\varphi (u)),`$
$`\varphi _{u_3}(u)=\stackrel{~}{X}_3(\varphi (u))+\left(u_3\varphi _1(u)\right)\stackrel{~}{T}(\varphi (u)).`$
Observing that for every $`j=1,\mathrm{},2n`$, we have
$$dF(\varphi (u))\stackrel{~}{X_j}(\varphi (u))=X_j(\mathrm{\Phi }(u))H_{\mathrm{\Phi }(u)}^n\text{and}$$
$$dF(\varphi (u))\stackrel{~}{Z}(\varphi (u))=Z(\mathrm{\Phi }(u))T_{\mathrm{\Phi }(u)}^n,$$
hence we obtain
$`\mathrm{\Phi }_{u_1}(u)=X_1(\mathrm{\Phi }(u))+\left(u_3+u_1\right)T(\mathrm{\Phi }(u)),\mathrm{\Phi }_{u_2}(u)=X_2(\mathrm{\Phi }(u))+u_2T(\mathrm{\Phi }(u)),`$
$`\mathrm{\Phi }_{u_3}(u)=X_3(\mathrm{\Phi }(u))+\left(u_3u_1\right)T(\mathrm{\Phi }(u)).`$
Thus, we can compute
$`\pi _๐ฑ\left(\mathrm{\Phi }_{u_1}\mathrm{\Phi }_{u_2}\mathrm{\Phi }_{u_3}\right)`$
$`=(u_3u_1)X_1X_2Tu_2X_1X_3T+(u_3+u_1)X_2X_3T,`$
hence formula (108) yields
$`๐ฎ_^2^4(\mathrm{\Sigma })={\displaystyle _U}\sqrt{u_2^2+2(u_3^2+u_1^2)}๐u.`$
###### Example 4.12
Let $`\varphi :^2^3`$, $`\varphi (u_1,u_2)=(a_1u_1,a_2u_2,bu_1+cu_2)`$, define a hyperplane in $`^3`$, where $`a_1,a_2,b,c`$ and $`\left(J\varphi \right)^2=a_1^2a_2^2+a_1^2c^2+a_2^2b^2>0`$. Embedding the hyperplane in $`^1`$ through standard coordinates $`F:^3^1`$, we obtain
$`\mathrm{\Phi }_{u_1}(u)=a_1X_1(\mathrm{\Phi }(u))+(a_1a_2u_2+b)T(\mathrm{\Phi }(u)),`$
$`\mathrm{\Phi }_{u_2}(u)=a_2X_2(\mathrm{\Phi }(u))+(ca_1a_2u_1)T(\mathrm{\Phi }(u)),`$
where $`\mathrm{\Phi }=F\varphi `$. Then we get
$`\pi _๐ฑ\left(\mathrm{\Phi }_{u_1}(u)\mathrm{\Phi }_{u_2}(u)\right)=a_1(ca_1a_2u_1)X_1T(a_1a_2u_2+b)a_2X_2T`$
and formula (108) yields
$`๐ฎ_^1^3(\mathrm{\Pi })={\displaystyle _U}\sqrt{a_1^2(ca_1a_2u_1)^2+a_2^2(a_1a_2u_2+b)^2}๐u,`$ (112)
where $`\mathrm{\Pi }=\mathrm{\Phi }(U)`$ and $`U`$ is an open bounded set of $`^2`$.
###### Example 4.13
Let $`\varphi :^2^3`$, $`\varphi (u_1,u_2)=(u_1,u_2,\frac{u_1^2+u_2^2}{2})`$, define a paraboloid in $`^3`$. By standard coordinates $`F:^3^1`$ and arguing as in the previous examples, we have
$$\mathrm{\Phi }_{u_1}(u)=X_1(\mathrm{\Phi }(u))+(u_2+u_1)T(\mathrm{\Phi }(u))\text{and}$$
$$\mathrm{\Phi }_{u_2}(u)=X_2(\mathrm{\Phi }(u))+(u_2u_1)T(\mathrm{\Phi }(u)),$$
where $`\mathrm{\Phi }=F\varphi `$. It follows that
$`\pi _๐ฑ\left(\mathrm{\Phi }_{u_1}(u)\mathrm{\Phi }_{u_2}(u)\right)=(u_2u_1)X_1T(u_2+u_1)X_2T`$
and formula (108) yields
$`๐ฎ_^1^3(๐ซ)={\displaystyle _U}\sqrt{2u_1^2+2u_2^2}๐u`$ (113)
where $`๐ซ=\mathrm{\Phi }(U)`$ and $`U`$ is an open bounded set of $`^2`$.
###### Remark 4.14
It is curious to notice that the density of $`๐ฎ_^1^3`$ restricted to the paraboloid $`๐ซ`$, computed in (113), is proportional to the density of $`๐ฎ_^1^3`$ restricted to the horizontal projection of $`๐ซ`$ onto the plane $`F\left(\{(x_1,x_2,x_3)x_3=0\}\right)^1`$, whose density is given by (112) in the case $`a_1=a_2=1`$ and $`b=c=0`$.
###### Example 4.15
From computations of Example 4.12, one can get the 2-dimensional spherical Hausdorff measure of the line $`\mathrm{\Phi }(t)=F(at,0,bt)`$ defined on an interval $`[\alpha ,\beta ]`$. We have
$`\mathrm{\Phi }^{}(t)=aX_1(\mathrm{\Phi }(t))+bT(\mathrm{\Phi }(t))\text{and}\pi _๐ฑ(\mathrm{\Phi }^{}(t))=bT(\mathrm{\Phi }(t)),`$ (114)
then defining the submanifold $`=\mathrm{\Phi }([\alpha ,\beta ])`$, the formula $`๐ฎ_^1^2()=|b|(\beta \alpha )`$ holds.
Another consequence of (2) is the lower semicontinuity of the spherical Hausdorff measure with respect to weak convergence of regular currents. To see this, it suffices to establish the following formula
$`๐ฎ_^n^{p+1}(\mathrm{\Sigma })=\underset{\omega _c^p(\mathrm{\Omega })}{sup}{\displaystyle _\mathrm{\Sigma }}\tau _{\mathrm{\Sigma },๐ฑ},\omega ๐\text{vol}_p,`$ (115)
where $`_c^p(\mathrm{\Omega })`$ is the space of smooth $`p`$-forms with compact support in $`\mathrm{\Omega }`$ with $`|\omega |1`$. The norm of $`\omega `$ is defined making the standard frame of $`p`$-forms $`(dx_1,dx_2,\mathrm{},dx_{2n},\stackrel{~}{\theta })`$ orthonormal and extending this scalar product to $`p`$-forms exactly as we have seen in formula (25). The 1-form $`\stackrel{~}{\theta }`$ is the so called contact form
$`\stackrel{~}{\theta }=dx_{2n+1}+{\displaystyle \underset{j=1}{\overset{n}{}}}x_{j+n}dx_jx_jdx_{j+n}`$ (116)
written in standard coordinates. Note that $`(dx_1,dx_2,\mathrm{},dx_{2n},\stackrel{~}{\theta )}`$ is the dual basis of $`(\stackrel{~}{X}_1,\mathrm{},\stackrel{~}{X}_{2n},\stackrel{~}{Z})`$. Formula (115) follows from (2) observing that
$`{\displaystyle _\mathrm{\Sigma }}|\tau _{\mathrm{\Sigma },๐ฑ}|๐\text{vol}_p=\underset{\omega _c^p(\mathrm{\Omega })}{sup}{\displaystyle _\mathrm{\Sigma }}\tau _{\mathrm{\Sigma },๐ฑ},\omega ๐\text{vol}_p,`$ (117)
as one can check by standard arguments. As a consequence of these observations, we can establish the following proposition.
###### Proposition 4.16
Let $`(\mathrm{\Sigma }_m)`$ be a sequence of $`C^1`$ submanifolds of $`\mathrm{\Omega }`$ which weakly converges in the sense of currents to the $`C^1`$ submanifold $`\mathrm{\Sigma }`$. Then
$`\underset{m\mathrm{}}{lim\; inf}๐ฎ_^n^{p+1}(\mathrm{\Sigma }_m)๐ฎ_^n^{p+1}(\mathrm{\Sigma }).`$ (118)
Proof. By hypothesis
$`{\displaystyle _{\mathrm{\Sigma }_m}}\tau _{\mathrm{\Sigma }_m,๐ฑ},\omega ๐\text{vol}_p{\displaystyle _\mathrm{\Sigma }}\tau _{\mathrm{\Sigma },๐ฑ},\omega ๐\text{vol}_p,`$ (119)
for every $`\omega _c^p(\mathrm{\Omega })`$. Then (115) ends the proof. $`\mathrm{}`$
###### Remark 4.17
It is clear the importance of (118) in studying versions of the Plateau problem with respect to the geometry of Heisenberg groups.
Recall that the horizontal normal is the orthogonal projection of the normal to $`\nu (x)`$ to $`T_x\mathrm{\Sigma }`$ onto the horizontal subspace $`H_x^n`$. In the next proposition we show that in codimension one an explicit relationship can be established between vertical tangent $`2n`$-vector and horizontal normal $`\nu _H`$.
###### Proposition 4.18
Let $`\mathrm{\Sigma }`$ be a $`2n`$-dimensional submanifold of class $`C^1`$ and let $`\nu _H(x)`$ a horizontal normal at $`x\mathrm{\Sigma }`$. Then we have
$$\nu _H^j=(1)^j\tau _{\mathrm{\Sigma },๐ฑ}^j$$
where $`\nu _H=_{j=1}^{2n}\nu _H^jX_j`$ and $`\tau _{\mathrm{\Sigma },๐ฑ}=_{j=1}^{2n}\tau _{\mathrm{\Sigma },๐ฑ}^jX_1\mathrm{}X_{j1}X_{j+1}\mathrm{}X_{2n}Z.`$ In particular, the equality $`|\tau _{\mathrm{\Sigma },๐ฑ}|=|\nu _H|`$ holds.
Proof Let $`(t_1,t_2,\mathrm{},t_{2n})`$ be an orthonormal basis of $`T_x\mathrm{\Sigma }`$, where $`x`$ is a transverse point. Then
$$t_j=\underset{i=1}{\overset{2n}{}}c_j^iX_i(x)+c_j^{2n+1}Z(x)$$
where $`C=(c_j^i)`$ is a $`(2n+1)\times 2n`$ matrix, whose columns are orthonormal vectors of $`^{2n+1}`$. Then we have
$`\tau _\mathrm{\Sigma }(x)=t_1t_2\mathrm{}t_{2n}={\displaystyle \underset{j=1}{\overset{2n+1}{}}}det(\text{ฤ}^j)X_1X_2\mathrm{}X_{j1}X_{j+1}\mathrm{}Z,`$
where $`\text{ฤ}^j`$ is the $`2n\times 2n`$ matrix obtained by removing the $`j`$-th row from $`C`$. The vertical projection yields
$`\tau _{\mathrm{\Sigma },๐ฑ}(x)=\pi _๐ฑ\left(\tau _\mathrm{\Sigma }(x)\right)={\displaystyle \underset{j=1}{\overset{2n}{}}}det(\text{ฤ}^j)X_1X_2\mathrm{}X_{j1}X_{j+1}\mathrm{}Z`$ (120)
and by elementary linear algebra one can deduce that
$`{\displaystyle \underset{j=1}{\overset{2n+1}{}}}(1)^jdet(\text{ฤ}^j)c_k^j=det\left[\begin{array}{cc}C& c_k\end{array}\right]=0`$ (122)
for every $`k=1,\mathrm{},2n`$. Then the vector
$$\nu =\underset{j=1}{\overset{2n}{}}(1)^jdet(\text{ฤ}^j)X_j+(1)^{2n+1}det(\text{ฤ}^{\mathrm{\hspace{0.33em}2}n+1})Z$$
yields a unit normal to $`\mathrm{\Sigma }`$ at $`x`$. Its horizontal projection is
$`\nu _H={\displaystyle \underset{j=1}{\overset{2n}{}}}(1)^jdet(\text{ฤ}^j)X_j.`$ (123)
Formulae (120) and (123) yield the thesis. $`\mathrm{}`$
## 5 Coarea formula
This section is devoted to the proof of Theorem 1.3. Next, we recall the Riemannian coarea formula, see Section 13.4 of .
###### Theorem 5.1
Let $`f:^n^k`$ be a Riemannian Lipschitz function, with $`1k<2n+1`$. Then for any summable map $`u:^n`$, the following formula holds
$$_^nu(x)J_gf(x)๐\text{vol}_{2n+1}(x)=_^k\left(_{f^1(t)}u(y)๐\text{vol}_p(y)\right)๐t,$$
(124)
where $`p=2n+1k`$
In the previous theorem the Heisenberg group $`^n`$ is equipped with its left invariant Riemannian metric $`g`$. The terminology โRiemannian Lipschitz mapโ means that the map is Lipschitz with respect to the Riemannian distance. Proof of Theorem 1.3. We first prove (4) in the case $`f`$ is defined on all of $`^n`$ and is of class $`C^1`$. Let $`\mathrm{\Omega }`$ be an open subset of $`^n`$. In view of Riemannian coarea formula (124), we have
$`{\displaystyle _\mathrm{\Omega }}u(x)J_gf(x)๐x={\displaystyle _^k}\left({\displaystyle _{f^1(t)\mathrm{\Omega }}}u(y)๐\text{vol}_p(y)\right)๐t,`$ (125)
where $`u:\mathrm{\Omega }[0,+\mathrm{}]`$ is a measurable function. Note that in the left hand side of (125) we have used the Lebesgue measure in that, by Proposition 2.7, it coincides with the volume measure expressed in terms of standard coordinates, namely $`F_{\mathrm{}}(^{2n+1})=\text{vol}_{2n+1}`$. Now we define
$$u(x)=J_Hf(x)\mathrm{๐}_{\{Jf0\}\mathrm{\Omega }}(x)/Jf(x)$$
and use (125), obtaining
$`{\displaystyle _\mathrm{\Omega }}J_Hf(x)๐x={\displaystyle _^k}\left({\displaystyle _{f^1(t)\mathrm{\Omega }}}{\displaystyle \frac{J_Hf(x)\mathrm{๐}_{\{Jf0\}}(x)}{Jf(x)}}๐\text{vol}_p(y)\right)๐t.`$ (126)
The validity of (125) also implies that for a.e. $`t^k`$ the set of points of $`f^1(t)`$ where $`J_gf`$ vanishes is $`\text{vol}_p`$-negligible, then the previous formula becomes
$`{\displaystyle _\mathrm{\Omega }}J_Hf(x)๐x={\displaystyle _^k}\left({\displaystyle _{f^1(t)\mathrm{\Omega }}}{\displaystyle \frac{J_Hf(x)}{Jf(x)}}๐\text{vol}_p(y)\right)๐t.`$ (127)
By classical Sardโs theorem and Theorem 4.1 for a.e. $`t^k`$ the $`C^1`$ submanifold $`f^1(t)`$ has $`๐ฎ_^n^{p+1}`$-negligible characteristic points, hence Proposition 3.2 implies that
$$C_t=\{yf^1(t)\mathrm{\Omega }J_Hf(y)=0\}$$
is $`๐ฎ_^n^{p+1}`$-negligible. As a result, from formulae (30) and (2) we have proved that for a.e. $`t^k`$ the equalities
$$_{f^1(t)\mathrm{\Omega }}\frac{J_Hf(x)}{J_gf(x)}๐\text{vol}_p(y)=๐ฎ_^n^{p+1}(f^1(t)\mathrm{\Omega }C_t)=๐ฎ_^n^{p+1}(f^1(t)\mathrm{\Omega })$$
hold, therefore (127) yields
$`{\displaystyle _\mathrm{\Omega }}J_Hf(x)๐x={\displaystyle _^k}๐ฎ_^n^{p+1}(f^1(t)\mathrm{\Omega })๐t.`$ (128)
The arbitrary choice of $`\mathrm{\Omega }`$ yields the validity of (128) also for arbitrary closed sets. Then, approximation of measurable sets by closed ones, Borel regularity of $`๐ฎ_^n^{p+1}`$ and the coarea estimate 2.10.25 of extend the validity of (128) to the following one
$`{\displaystyle _A}J_Hf(x)๐x={\displaystyle _^k}๐ฎ_^n^{p+1}(f^1(t)A)๐t,`$ (129)
where $`A`$ is a measurable subset of $`^n`$. Now we consider the general case, where $`f:A^k`$ is a Lipschitz map defined on a measurable bounded subset $`A`$ of $`^3`$. Let $`f_1:^n^k`$ be a Lipschitz extension of $`f`$, namely, $`f_{1}^{}{}_{|A}{}^{}=f`$ holds. Due to the Whitney extension theorem (see for instance 3.1.15 of ) for every arbitrarily fixed $`\epsilon >0`$ there exists a $`C^1`$ function $`f_2:^n^k`$ such that the open subset $`O=\{z^nf_1(z)f_2(z)\}`$ has Lebesgue measure less than or equal to $`\epsilon `$. We wish to prove
$`\left|{\displaystyle _A}J_Hf(x)๐x{\displaystyle _^k}๐ฎ_^n^{p+1}(f^1(t)A)๐t\right|{\displaystyle _{AO}}J_Hf(x)๐x`$
$`+{\displaystyle _^k}๐ฎ_^n^{p+1}(f^1(t)AO)๐t.`$ (130)
In fact, due to the validity of (129) for $`C^1`$ mappings, we have
$$_{AO}J_Hf_2(x)๐x=_^k๐ฎ_^n^{p+1}(f_2^1(t)AO)๐t.$$
Note here that the horizontal jacobian $`J_Hf`$ is well defined on $`A`$, in that $`df`$ is well defined at density points of the domain, see for instance Definition 7 and Proposition 2.2 of . The equality $`f_{2}^{}{}_{|AO}{}^{}=f_{|AO}`$ implies that $`J_Hf_2=J_Hf`$ a.e. on $`AO`$, therefore
$$_{AO}J_Hf(x)๐x=_^k๐ฎ_^n^{p+1}(f^1(t)AO)๐t$$
holds and inequality (5) is proved. Now we observe that for a.e. $`xA`$, we have
$$J_Hf(x)\underset{i=1}{\overset{k}{}}\left(\underset{j=1}{\overset{2n}{}}\left(X_jf^i(x)\right)^2\right)^{1/2}df(x)_{|H_x^n}^k$$
therefore the estimate
$`J_Hf(x)\text{Lip}(f)^k`$ (131)
holds for a.e. $`xA`$. By virtue of the general coarea inequality 2.10.25 of there exists a dimensional constant $`c_1>0`$ such that
$`{\displaystyle _^k}๐ฎ_^n^{p+1}(f^1(t)AO)๐tc_1\text{Lip}(f)^k^{2n+2}(O).`$ (132)
The fact that the $`2n+2`$-dimensional Hausdorff measure $`^{2n+2}`$ with respect to the homogeneous distance $`\rho `$ is proportional to the Lebesgue measure, gives us a constant $`c_2>0`$ such that
$`{\displaystyle _^k}๐ฎ_^n^{p+1}(f^1(t)AO)๐tc_2\text{Lip}(f)^k^{2n+1}(O)c_2\text{Lip}(f)^k\epsilon .`$ (133)
Thus, estimates (131) and (133) joined with inequality (5) yield
$`\left|{\displaystyle _A}J_Hf(x)๐x{\displaystyle _^k}๐ฎ_^n^{p+1}(f^1(t)A)๐t\right|(1+c_2)\text{Lip}(f)^k\epsilon .`$
Letting $`\epsilon 0^+`$, we have proved that
$`{\displaystyle _A}J_Hf(x)๐x={\displaystyle _^k}๐ฎ_^n^{p+1}(f^1(t)A)๐t.`$ (134)
Finally, utilizing increasing sequences of step functions pointwise converging to $`u`$ and applying Beppo Levi convergence theorem the proof of (4) is achieved in the case $`A`$ is bounded. If $`A`$ is not bounded, then one can take the limit of (4) where $`A`$ is replaced by $`A_k`$ and $`\{A_k\}`$ is an increasing sequence of measurable bounded sets whose union yields $`A`$. Then the Beppo Levi convergence theorem concludes the proof. $`\mathrm{}`$
###### Remark 5.2
Notice that once $`f:A^k`$ in the previous theorem is considered with respect to standard coordinates it is easy to check that the locally Lipschitz property with respect to the Euclidean distance of $`^{2n+1}`$ is equivalent to the locally Lipschitz property with respect to the Riemannian distance.
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# Isomorphisms of Kac-Moody groups which preserve bounded subgroups
## 1 Introduction
Kac-Moody groups are infinite-dimensional generalizations of Chevalley groups. It is known that each automorphism of a Chevalley group (of irreducible type and over a perfect field) can be written as a product of an inner, a diagonal, a graph and a field automorphism (see Theorem 30 in ). In it was conjectured that the same statement holds for Kac-Moody groups over algebraically closed fields of characteristic 0 up to the addition of a so called sign automorphism. In this conjecture is shown to be true for Kac-Moody groups over algebraically closed fields of any characteristic. This is achieved in loc. cit. by solving the isomorphism problem for those groups. In this paper, we study the isomorphism problem for Kac-Moody groups over arbitrary fields of cardinality at least $`4`$. We restrict our attention to isomorphisms which preserve the set of bounded subgroups. In this context, a subgroup of a Kac-Moody group is called bounded if it is contained in the intersection of two finite type parabolic subgroups of opposite signs.
Throughout the paper we use Titsโ definition for Kac-Moody groups over fields . This definition does not only provide the abstract Kac-Moody group $`G`$ but also a canonical system $`(U_\alpha )_{\alpha \mathrm{\Phi }}`$ of root subgroups. The pair $`(G,(U_\alpha )_{\alpha \mathrm{\Phi }})`$ is an example of a so called twin root datum. Twin root data have been introduced by Tits in order to give suitable axioms for these pairs arising from his definition of Kac-Moody groups.
Our main result says that if two Kac-Moody groups are isomorphic via such an isomorphism, then the groups are of the same type and defined over the same ground field. Here is a precise statement.
###### Theorem.
Let $`๐=(G,(U_\alpha )_{\alpha \mathrm{\Phi }})`$ and $`๐^{}=(G^{},(U_\alpha ^{}^{})_{\alpha ^{}\mathrm{\Phi }^{}})`$ be two twin root data associated with two Kac-Moody groups of non-spherical type over fields of cardinality at least 4. Let $`\xi :GG^{}`$ be a group isomorphism which maps bounded subgroups of $`G`$ to bounded subgroups of $`G^{}`$. Then $`\xi `$ induces an isomorphism of $`๐`$ to $`๐^{}`$.
We refer to Section 2.3.2 below for the definition of an isomorphism between twin root data. Roughly speaking, it means that $`(\xi (U_\alpha ))_{\alpha \mathrm{\Phi }}`$ is โnearly $`G^{}`$-conjugateโ to $`(U_\alpha ^{}^{})_{\alpha ^{}\mathrm{\Phi }^{}}`$.
As it is the case in the paper , the present work makes crucial use of the theory of twin buildings. A group endowed with a twin root datum is indeed naturally endowed with a strongly transitive action on a twin building, and the combinatorial properties of this action turn out to be the most appropriate tool in studying Kac-Moody groups from our point of view. However, we have tried to explain each crucial building-theoretic statement in more classical terms, without making reference to the language of buildings. We hope this will help the reader who is not familiar to the theory of buildings to understand the main ideas of this paper.
As a consequence of the theorem above, we obtain the following result on automorphisms of Kac-Moody groups.
###### Corollary A.
Let $`G`$ be a Kac-Moody group over a field of cardinality at least 4. Let $`\phi `$ be an automorphism of $`G`$ which preserves the set of bounded subgroups. Then $`\phi `$ splits as a product of an inner, a diagonal, a graph, a field and a sign automorphism.
There are mainly two motivations to consider isomorphisms which preserve bounded subgroups.
The first motivation comes from the earlier work by Kac and Wang. In this paper, automorphisms of Kac-Moody groups over fields of characteristic 0 and associated with symmetrizable Cartan matrices have been studied. One of the main results of loc. cit. is that, given such a Kac-Moody group $`G`$ and its Kac-Moody algebra $`๐ค`$, then an automorphism of $`G`$ which preserves the set of $`\mathrm{Ad}_๐ค^{}`$-finite elements splits as a product as in Corollary A above, where $`๐ค^{}:=[๐ค,๐ค]`$. We recall that an element $`gG`$ is $`\mathrm{Ad}_๐ค^{}`$-finite if and only if the subgroup generated by $`g`$ is bounded (see , Theorem 2.10). Thus, Corollary A can be seen as a weaker version of Kac-Wangโs result, which remains valid for Kac-Moody groups of arbitrary type and over fields of arbitrary characteristic.
The second motivation is the fact that, in the case of a Kac-Moody group over a finite field, a subgroup is bounded if and only if it is finite (see Corollary 3.8 below). Therefore, all isomorphisms preserve bounded subgroups in this case. Consequently, we obtain the following result.
###### Corollary B.
Let $`๐=(G,(U_\alpha )_{\alpha \mathrm{\Phi }})`$ and $`๐^{}=(G^{},(U_\alpha ^{}^{})_{\alpha ^{}\mathrm{\Phi }^{}})`$ be two twin root data associated with two Kac-Moody groups over finite fields of cardinality at least 4 and let $`\xi :GG^{}`$ be an isomorphism. Then $`\xi `$ induces an isomorphism of $`๐`$ to $`๐^{}`$ unless $`G`$ and $`G^{}`$ are both finite. Moreover, any automorphism of $`G`$ splits as a product of an inner, a diagonal, a graph, a field and a sign automorphism.
Kac-Moody groups over finite fields are finitely generated and some subclasses of them are known to be finitely presented. In the recent years these groups became important in geometric group theory for several reasons (see ). In this context, B. Rรฉmy proved a factorization theorem for the automorphisms of certain Kac-Moody groups (see , Theorem 3.1). Corollary 2 covers this result as a special case.
Let us also mention the existence of exotic constructions of groups of Kac-Moody type, which are not Kac-Moody groups in the strict sense but which are also endowed with a twin root datum. For example, Rรฉmy and Ronan constructed examples of groups of Kac-Moody type defined simultaneously over different ground fields. It turns out that, provided the maximal tori are locally large enough, our methods extend also to these exotic cases, and the interested reader will have no difficulty to extend our arguments to this slightly more general situation (see also the introduction of for other remarks and results related to the isomorphism problem of exotic groups of Kac-Moody type).
The paper is organized as follows. After a preliminary section where definitions are recalled, notation is fixed and some auxiliary results are proven, we discuss in Section 3 the Levi decomposition of the intersection of two parabolic subgroups of opposite sign in a group endowed with a twin root datum (a similar but slightly less general discussion had been done in ). The key result of this paper is contained in Section 4, where we prove that the maximal bounded subgroups coincide almost always with the Levi factors of the maximal parabolic subgroups of finite type. In the next section, we use this key result to state and prove a technical version of our main result, which is valid for a larger class of groups endowed with a twin root datum. Finally, the last two sections are devoted to the proof of the main theorem above and its corollaries.
## 2 Preliminaries
The main references are , , and .
We start by fixing a general convention :
*The ordered pair $`(W,S)`$ is a Coxeter system and $`\mathrm{}`$ denotes the corresponding length function. For $`JS`$ we set $`W_J:=J`$ and we call $`J`$ spherical whenever $`W_J`$ is finite.*
### 2.1 Buildings
#### 2.1.1 Definition
A building of type $`(W,S)`$ is a set $`\mathrm{\Delta }`$, whose elements are called chambers, endowed with a map $`\delta :\mathrm{\Delta }\times \mathrm{\Delta }W`$ called the $`W`$-distance satisfying the following axioms, where $`x,y\mathrm{\Delta }`$ and $`w=\delta (x,y)`$:
$$\begin{array}{cc}\text{(Bu1)}\hfill & w=1x=y;\hfill \\ \text{(Bu2)}\hfill & \text{if }z\mathrm{\Delta }\text{ is such that }\delta (y,z)=sS,\text{ then }\delta (x,z)\{w,ws\};\hfill \\ & \text{if, furthermore, }\mathrm{}(ws)=\mathrm{}(w)+1,\text{ then }\delta (x,z)=ws;\hfill \\ \text{(Bu3)}\hfill & \text{if }sS,\text{ there exists }z\mathrm{\Delta }\text{ such that }\delta (y,z)=s\text{ and }\delta (x,z)=ws.\hfill \end{array}$$
For any two chambers $`x,y\mathrm{\Delta }`$, the natural number $`\mathrm{}(\delta (x,y))`$ is called the numerical distance between $`x`$ and $`y`$.
An isometry between subsets of buildings of type $`(W,S)`$ is a bijection preserving the $`W`$-distance.
#### 2.1.2 Apartments
Given a Coxeter system $`(W,S)`$, let $`\delta :W\times WW`$ be defined by $`\delta :(x,y)x^1y`$. In this way, we endow $`W`$ with a canonical structure of a building of type $`(W,S)`$. This building is denoted by $`๐(W,S)`$.
Any subset of a building $`(\mathrm{\Delta },\delta )`$ of type $`(W,S)`$ which is isometric to the canonical building $`๐(W,S)`$ is called an apartment. A fundamental property of buildings is that any two chambers lie in a common apartment (see , Theorem 3.7).
#### 2.1.3 Panels, residues and galleries
Given $`c\mathrm{\Delta }`$ and $`sS`$, then the set $`\left\{x\mathrm{\Delta }\right|\delta (x,c)\{1,s\}\}`$ is called an $`s`$-panel of $`\mathrm{\Delta }`$ or a panel of type $`s`$. A panel is an $`s`$-panel for some $`sS`$. More generally, for $`c\mathrm{\Delta }`$ and $`JS`$ the set
$$\mathrm{Res}_J(c):=\left\{x\mathrm{\Delta }\right|\delta (x,c)W_J\}$$
is called the $`J`$-residue of $`\mathrm{\Delta }`$ which contains $`c`$. Its rank is the cardinality of the set $`J`$; hence, residues of rank 0 are just chambers and the residues of rank 1 are panels.
It is an important fact that a $`J`$-residue is itself a building of type $`(W_J,J)`$ with the $`W_J`$-distance induced by $`\delta `$ (see , Theorem 3.5). Moreover, given a residue $`R`$ and an apartment $`\mathrm{\Sigma }`$ in a building $`\mathrm{\Delta }`$, the intersection $`R\mathrm{\Sigma }`$ is either empty or an apartment of $`R`$, and all apartments of $`R`$ arise in this way. It is common and handy to say that $`R`$ *is contained in* $`\mathrm{\Sigma }`$ whenever $`R\mathrm{\Sigma }`$ is nonempty.
A sequence of chambers such that two consecutive chambers are adjacent, namely contained in a common panel, is called a gallery. The gallery $`\gamma =(x_0,x_1,\mathrm{},x_n)`$ is called minimal if $`n=\mathrm{}(\delta (x_0,x_n))`$.
A building is called thin (resp. thick) if all of its panels have cardinality 2 (resp. at least 3). Any thin building of type $`(W,S)`$ is isomorphic to the canonical building $`(W,S)`$.
#### 2.1.4 Projections and convexity in arbitrary buildings
A fundamental property of buildings, besides the existence of apartments, is the exitence of projections onto residues. We review here the main properties of projections in arbitratry buildings. The notion of a projection can be slightly refined in the case of twin buildings; we will come back to this refinement in Section 2.2.3 below.
Let $`(W,S)`$ be a Coxeter system. We recall from that if $`J,KS`$ and $`wW`$, then there is a unique element of minimal length in the double coset $`W_JwW_K`$.
Let $`(\mathrm{\Delta },\delta )`$ be a building of type $`(W,S)`$ and let $`R_J,R_K`$ be residues of $`\mathrm{\Delta }`$ of respective type $`J`$ and $`K`$. Then the set of all $`\delta (c,d)`$ for $`cR_J`$ and $`dR_K`$ is a double coset $`W_JwW_K`$. Its minimal element is denoted by $`\delta (R_J,R_K)`$.
The set
$$\mathrm{proj}_{R_J}(R_K):=\{cR_J|dR_K\text{ such that }\delta (c,d)=\delta (R_J,R_K)\}$$
is called the projection of $`R_K`$ on $`R_J`$. It is a residue of type $`JwKw^1`$, where $`w:=\delta (R_J,R_K)`$; in particular, it is a spherical residue whenever $`J`$ or $`K`$ is spherical. Moreover, we have
$$\mathrm{proj}_{R_J}(R_K)=\{\mathrm{proj}_{R_J}(c)|cR_K\},$$
where we have written $`\mathrm{proj}_{R_J}(c)`$ for $`\mathrm{proj}_{R_J}(\{c\})`$.
If $`c`$ is a chamber and $`R`$ a residue, then $`\mathrm{proj}_R(c)`$ is a gate of $`c`$ to $`R`$. This means that for any $`xR`$ there exists a minimal gallery joining $`c`$ to $`x`$ via $`\mathrm{proj}_R(c)`$. The chamber $`\mathrm{proj}_R(c)`$ is the unique chamber of $`R`$ at minimal numerical distance from $`c`$.
A set $`๐ณ`$ of chambers in a building $`\mathrm{\Delta }`$ is called convex if the following property holds: *given chambers $`x,x^{}๐ณ`$ and a spherical residue $`R`$ containing $`x`$, then $`\mathrm{proj}_R(x^{})๐ณ`$.* For example, apartments and residues are convex sets of chambers.
#### 2.1.5 Spherical residues and spherical buildings
A building $`(\mathrm{\Delta },\delta )`$ of type $`(W,S)`$ is called spherical if $`W`$ is finite. In that case, there exists a unique element $`w_0`$ of maximal length in $`W`$. Two chambers $`x,y\mathrm{\Delta }`$ are called opposite if $`\delta (x,y)=w_0`$. Two residues $`R_J`$ and $`R_K`$ of $`\mathrm{\Delta }`$ of type $`J`$ and $`K`$ respectively are called opposite if they contain opposite chambers and if $`J=w_0Kw_0^1`$.
A residue $`R`$ of type $`J`$ in a building of arbitrary type $`(W,S)`$ is called spherical if $`J`$ is a spherical subset of $`S`$. Thus $`R`$ is a spherical building and it makes sense to talk about opposite chambers and opposite residues of $`R`$.
The following lemma is a useful criterion of sphericity in terms of projections.
###### Lemma 2.1.
Let $`(\mathrm{\Delta },\delta )`$ be a building of type $`(W,S)`$, let $`JS`$ and let $`R`$ be a $`J`$-residue. Then $`J`$ is spherical if and only if there exist $`x,yR`$ such that for every $`jJ`$, we have $`\mathrm{proj}_{\pi _j}(x)y`$ where $`\pi _j`$ denotes the $`j`$-panel containing $`y`$.
###### Proof.
This follows from , Theorem (2.16). โ
We end this subsection by recalling a celebrated fixed point theorem for finite groups acting on buildings.
###### Proposition 2.2.
Any finite group acting on a building of type $`(W,S)`$, where $`S`$ is finite, stabilizes a residue of spherical type.
###### Proof.
See , Corollary 11.9. โ
### 2.2 Twin buildings
#### 2.2.1 Definition
A twinned pair of buildings or twin building of type $`(W,S)`$ is a pair $`((\mathrm{\Delta }_+,\delta _+),(\mathrm{\Delta }_{},\delta _{}))`$ of buildings of type $`(W,S)`$, endowed with a $`W`$-codistance
$$\delta ^{}:(\mathrm{\Delta }_+\times \mathrm{\Delta }_{})(\mathrm{\Delta }_{}\times \mathrm{\Delta }_+)W$$
satisfying the following axioms, where $`ฯต\{+,\}`$, $`x\mathrm{\Delta }_ฯต`$, $`y\mathrm{\Delta }_ฯต`$ and $`w=\delta ^{}(x,y)`$:
$$\begin{array}{cc}\text{(Tw1)}\hfill & \delta ^{}(y,x)=w^1;\hfill \\ \text{(Tw2)}\hfill & \text{if }z\mathrm{\Delta }_ฯต\text{ is such that }\delta _ฯต(y,z)=sS\text{ and }\mathrm{}(ws)<\mathrm{}(w),\text{ then }\delta ^{}(x,z)=ws;\hfill \\ \text{(Tw3)}\hfill & \text{if }sS,\text{ there exists }z\mathrm{\Delta }_ฯต\text{ such that }\delta _ฯต(y,z)=s\text{ and }\delta ^{}(x,z)=ws.\hfill \end{array}$$
In the sequel, we will often use the symbol $`\mathrm{\Delta }_+`$ to denote the building $`(\mathrm{\Delta }_+,\delta _+)`$ as well as its set of chambers and similarly for $`\mathrm{\Delta }_{}`$. The meaning will be clear from the context.
A residue $`R`$ of the twin building $`\mathrm{\Delta }=(\mathrm{\Delta }_+,\mathrm{\Delta }_{},\delta ^{})`$ is a residue of $`\mathrm{\Delta }_ฯต`$, for $`ฯต=+`$ or $``$ and $`ฯต`$ is called the sign of $`R`$. Two chambers $`x`$ and $`y`$ of opposite signs are called opposite if $`\delta ^{}(x,y)=1`$. Two residues are called opposite if they are of the same type and contain opposite chambers. Given $`JS`$, then a pair of opposite residues of type $`J`$ endowed with the $`W`$-codistance induced from $`\delta ^{}`$ is itself a twin building of type $`(W_J,J)`$.
Notice that we have defined the term *opposite* at two different places, namely in Section 2.1.5 above and here in Section 2.2.1. However, this terminology is standard and coherent. Indeed, the former notion applies to chambers or residues of the same sign and lying in a common spherical residue, while the latter applies to chambers or residues of opposite signs.
An automorphism of $`\mathrm{\Delta }=(\mathrm{\Delta }_+,\mathrm{\Delta }_{},\delta ^{})`$ is by definition a pair $`\phi =(\phi _+,\phi _{})`$ of permutations of $`\mathrm{\Delta }_+`$ and $`\mathrm{\Delta }_{}`$ respectively preserving the $`W`$-distances $`\delta _+`$, $`\delta _{}`$ as well as the $`W`$-codistance $`\delta ^{}`$. Isomorphisms of twin buildings are defined similarly. We recall from , Theorem 1, that if $`\mathrm{\Delta }_+`$ and $`\mathrm{\Delta }_{}`$ are thick, then an automorphism of $`\mathrm{\Delta }`$ which fixes a pair of opposite chambers $`c,c^{}`$ and all chambers adjacent to $`c`$ is the identity.
#### 2.2.2 Reflections and twin apartments
Let $`(\mathrm{\Sigma }_+,\delta _+)`$ and $`(\mathrm{\Sigma }_{},\delta _{})`$ be two copies of the canonical building $`๐(W,S)`$ of type $`(W,S)`$ (see Section 2.1.2). Let $`\delta ^{}:\mathrm{\Sigma }_+\times \mathrm{\Sigma }_{}W:(x,y)x^1y`$; this makes sense since $`\mathrm{\Sigma }_+=\mathrm{\Sigma }_{}=W`$. Then $`\mathrm{\Sigma }(W,S):=((\mathrm{\Sigma }_+,\delta ),(\mathrm{\Sigma }_{},\delta ),\delta ^{})`$ is a thin twin building of type $`(W,S)`$, namely a twinned pair of thin buildings. It is the unique thin twin building of that type up to isomorphism.
The group $`W`$ has a faithful action on $`\mathrm{\Sigma }(W,S)`$ by automorphisms which is given by left multiplication on $`\mathrm{\Sigma }_+`$ and $`\mathrm{\Sigma }_{}`$. Every automorphism of $`\mathrm{\Sigma }`$ is of this form.
A reflection is a non-trivial element of $`W`$ which stabilizes a panel of $`\mathrm{\Sigma }(W,S)`$. Conversely, to any panel of $`\mathrm{\Sigma }(W,S)`$ corresponds a unique reflection of $`W`$ which stabilizes it. Moreover, an element of $`W`$ is a reflection if and only if it is conjugate to an element of $`S`$.
Let $`\mathrm{\Delta }=(\mathrm{\Delta }_+,\mathrm{\Delta }_{})`$ be a twin buildings of type $`(W,S)`$. A pair $`\mathrm{\Sigma }=(\mathrm{\Sigma }_+,\mathrm{\Sigma }_{})`$ of subsets of $`\mathrm{\Delta }`$ is called a twin apartment if it is isomorphic to the canonical twin building $`\mathrm{\Sigma }(W,S)`$. Given a twin apartment $`\mathrm{\Sigma }=(\mathrm{\Sigma }_+,\mathrm{\Sigma }_{})`$, the restriction of the opposition relation of $`\mathrm{\Delta }`$ to $`\mathrm{\Sigma }`$ is a one-one correspondence $`\mathrm{\Sigma }_+\mathrm{\Sigma }_{}`$. (It corresponds to the identity $`\mathrm{id}:WW`$ in the canonical twin building $`\mathrm{\Sigma }(W,S)`$.) We denote it by $`\mathrm{op}_\mathrm{\Sigma }`$.
It is a fundamental fact that, given any two chambers $`x\mathrm{\Delta }_ฯต`$ and $`y\mathrm{\Delta }_ฯต^{}`$ in a twin building $`\mathrm{\Delta }=(\mathrm{\Delta }_+,\mathrm{\Delta }_{})`$ of type $`(W,S)`$, where $`ฯต,ฯต^{}\{+,\}`$, there exists a twin apartment $`\mathrm{\Sigma }=(\mathrm{\Sigma }_+,\mathrm{\Sigma }_{})`$ such that $`x\mathrm{\Sigma }_ฯต`$ and $`y\mathrm{\Sigma }_ฯต^{}`$ (see , Lemma 2). It is common and handy to say that $`x`$ and $`y`$ are contained in $`\mathrm{\Sigma }`$, and to write $`x,y\mathrm{\Sigma }`$.
#### 2.2.3 Projections and convexity in twin buildings
Let $`(W,S)`$ be a Coxeter system and let $`J,JS`$ be spherical subsets. Then there is a unique element of maximal length in $`W_JwW_K`$ (, Lemma 9).
Let $`\mathrm{\Delta }=(\mathrm{\Delta }_+,\mathrm{\Delta }_{},\delta ^{})`$ be a twin building of type $`(W,S)`$ and let $`R_J,R_K`$ be residues of $`\mathrm{\Delta }`$ of respective type $`J`$ and $`K`$. Assume moreover that $`J`$ and $`K`$ are spherical and that $`R_J`$ and $`R_K`$ have opposite signs. Then the set of all $`\delta ^{}(c,d)`$ for $`cR_J`$ and $`dR_K`$ is a double coset $`W_JwW_K`$. Its maximal element is denoted by $`\delta ^{}(R_J,R_K)`$. The set
$$\mathrm{proj}_{R_J}(R_K):=\{cR_J|dR_K\text{ such that }\delta ^{}(c,d)=\delta ^{}(R_J,R_K)\}$$
is called the projection of $`R_K`$ on $`R_J`$. It is a residue of type $`w_J^0(JwKw^1)w_J^0`$, where $`w:=\delta ^{}(R_J,R_K)`$ and $`w_J^0`$ denotes the maximal element of $`W_K`$ (, Lemma 10). Moreover, we have
$$\mathrm{proj}_{R_J}(R_K)=\{\mathrm{proj}_{R_J}(c)|cR_K\}.$$
A set $`๐ณ`$ of chambers in a twin building is called convex if the following condition holds: *given $`x,x^{}๐ณ`$ and a spherical residue $`R`$ containing $`x^{}`$, then $`\mathrm{proj}_R(x)๐ณ`$.* For example, twin apartments are convex sets of chambers in twin buildings. Actually, the convex hull of any pair of opposite chambers is a twin apartment containing them. This implies that two opposite chambers lie in a unique common twin apartment.
Notice that we have defined the terms *projections* and *convexity* at two different places, namely in Section 2.1.4 above and here in Section 2.2.3. The point is that the notion of projections in twin buildings is a generalization of the standard notion of projections in arbitrary buildings. There will be no confusion between both. Indeed, the meaning of the symbol $`\mathrm{proj}_R(x)`$ in the context of twin buildings depends on the respective signs of the residue $`R`$ and the chamber $`x`$.
We end this subsection with a result is often useful to compute projections between residues of opposite signs using twin apartments.
###### Lemma 2.3.
Let $`\mathrm{\Delta }=(\mathrm{\Delta }_+,\mathrm{\Delta }_{},\delta ^{})`$ be a twin building and for each sign $`ฯต`$ let $`R_ฯต`$ be a spherical residue of $`\mathrm{\Delta }_ฯต`$. Let $`\mathrm{\Sigma }`$ is a twin apartment containing $`R_+`$ and $`R_{}`$ and let $`ฯต\{+,\}`$. Let $`R_ฯต^{}`$ be the residue of $`\mathrm{\Sigma }`$ opposite $`R_ฯต`$. Then the residues $`\mathrm{proj}_{R_ฯต}(R_ฯต)`$ and $`\mathrm{proj}_{R_ฯต}(R_ฯต^{})`$ are contained in $`\mathrm{\Sigma }`$ and opposite in $`R_ฯต`$ (see Section 2.1.5).
###### Proof.
This is Proposition 4 in . โ
In brief, the statement of this lemma may be written as
$$\mathrm{proj}_{R_ฯต}(R_ฯต)=\mathrm{op}_{\mathrm{\Sigma }R_ฯต}(\mathrm{proj}_{R_ฯต}(\mathrm{op}_\mathrm{\Sigma }(R_ฯต))).$$
#### 2.2.4 Parallelism
Let $`\mathrm{\Delta }=(\mathrm{\Delta }_+,\mathrm{\Delta }_{},\delta ^{})`$ be a twin building of type $`(W,S)`$. Two residues $`R_J,R_K`$ of $`\mathrm{\Delta }`$ (assumed to be spherical if they have opposite signs) are called parallel if $`\mathrm{proj}_{R_J}(R_K)=R_J`$ and $`\mathrm{proj}_{R_K}(R_J)=R_K`$.
It follows from the definitions that $`\mathrm{proj}_{R_J}(R_K)`$ and $`\mathrm{proj}_{R_K}(R_J)`$ are always parallel.
Although parallel residues need not have the same type, they are nevertheless always โalmost isometricโ in the following sense.
###### Lemma 2.4.
Let $`R_J`$ (resp. $`R_K`$) be a residue of spherical type $`J`$ (resp. $`K`$) and sign $`ฯต_J`$ (resp. $`ฯต_K`$). Assume that $`R_J`$ and $`R_K`$ are parallel. Then there exists an isomorphism $`\eta :W_JW_K`$ with $`\eta (J)=K`$ such that
$$\delta _{ฯต_J}(\mathrm{proj}_{R_J}(x),\mathrm{proj}_{R_J}(y))=\eta (\delta _{ฯต_K}(x,y))$$
for all $`x,yR_K`$. In particular, if $`x`$ and $`y`$ are opposite in $`R_K`$, then so are $`\mathrm{proj}_{R_J}(x)`$ and $`\mathrm{proj}_{R_K}(y)`$ in $`R_J`$.
The proof of Lemma 2.4 is in the same spirit as the proof of Proposition 5.15 in and is omitted here. We only mention that the isomorphism $`\eta `$ of the lemma is actually induced by the conjugation by $`\delta _{ฯต_J}(R_J,R_K)`$ if $`ฯต_J=ฯต_K`$ and by $`w_J\delta ^{}(R_J,R_K)`$ if $`ฯต_J=ฯต_K`$. However, we do not need this fact here.
###### Lemma 2.5.
The spherical residues $`R_J`$ and $`R_K`$ of $`\mathrm{\Delta }`$ are opposite if and only if $`\mathrm{proj}_{R_J}(R_K)`$ and $`\mathrm{proj}_{R_K}(R_J)`$ are opposite.
###### Proof.
Since opposite spherical residues are parallel, the implication โ$``$โ is obvious. The other implication follows from an easy computation using Lemma 9 of . โ
Next, we give a rule for the composition of projections.
###### Lemma 2.6.
As before, $`\mathrm{\Delta }`$ is a (possibly twin) building of type $`(W,S)`$. Let $`R_I,R_J,R_K`$ be residues of type $`I,J,K`$ respectively and assume that $`R_IR_J`$. Moreover, if $`\mathrm{\Delta }`$ is a twin building and if $`R_J`$ and $`R_K`$ have opposite signs, then we also assume that $`I,J,K`$ are spherical. Then we have
$$\mathrm{proj}_{R_I}(R_K)=\mathrm{proj}_{R_I}(\mathrm{proj}_{R_J}(R_K)).$$
###### Proof.
It suffices to prove the statement when the residue $`R_K`$ is reduced to a single chamber, say $`c`$ (or, in other words, when $`K=\mathrm{}`$). If $`R_J`$ and $`c`$ have the same sign, the result follows from the fact that $`\mathrm{proj}_{R_J}(c)`$ is a gate of $`c`$ to $`R_J`$. If they have opposite signs, we may reduce ourselves to the preceding case in view of Lemma 2.3. โ
The following lemma characterizes the parallelism of spherical residues in thin buildings.
###### Proposition 2.7.
Let $`(\mathrm{\Sigma },\delta )`$ be the thin building of type $`(W,S)`$. Let $`J,K`$ be spherical subsets of $`S`$ and let $`R_J,R_K`$ be residues of type $`J,K`$ respectively. Then the following statements are equivalent :
(i) $`R_J`$ and $`R_K`$ are parallel; (ii) a reflection of $`\mathrm{\Sigma }`$ stabilizes $`R_J`$ if and only if it stabilizes $`R_K`$; (iii) there exist two sequences $`R_J=R_0,R_1,\mathrm{},R_n=R_K`$ and $`T_1,\mathrm{},T_n`$ of residues of spherical type such that for each $`1in`$ the rank of $`T_i`$ is equal to $`1+\mathrm{rank}(R_J)`$, the residues $`R_{i1}`$, $`R_i`$ are contained and opposite in $`T_i`$ and moreover, we have $`\mathrm{proj}_{T_i}(R_J)=R_{i1}`$ and $`\mathrm{proj}_{T_i}(R_K)=R_i`$.
###### Proof.
The equivalence (i) $``$ (ii) is easy. The implication (iii) $``$ (i) follows from an obvious induction on $`n`$ using the fact that opposite spherical residues are parallel. It remains to prove (i) $``$ (iii). Let $`sS`$ such that $`\mathrm{}(s\delta (R_J,R_K))<\mathrm{}(\delta (R_J,R_K))`$. Clearly $`sJ`$. Let $`xR_J`$ and set $`T_1:=\mathrm{Res}_{J\{s\}}(x)`$ and $`R_1:=\mathrm{proj}_{T_1}(R_K)`$. By definition of $`T_1`$ we have $`R_JR_1=\mathrm{}`$ and so $`R_1`$ is properly contained in $`T_1`$. By Lemma 2.6 we have $`\mathrm{proj}_{R_J}(R_1)=R_J`$. Therefore $`R_1`$ and $`R_J`$ have the same rank and so they are parallel.
Let $`x^{}R_1`$ such that $`\mathrm{proj}_{R_J}(x^{})=x`$ and choose $`y`$ opposite to $`x^{}`$ in $`R_1`$ (this makes sense since $`R_1`$, being the image of $`R_K`$ under a projection, is spherical). Let now $`\pi `$ be a panel containing $`x`$ and contained in $`T_1`$. If the type of $`\pi `$ is an element of $`J`$ then $`\mathrm{proj}_\pi (y)x`$ by Lemma 2.1 and Lemma 2.4. If the type of $`\pi `$ is $`s`$ then the same inequality is still true by the definition of $`s`$ and using Lemma 2.6. Therefore, $`T_1`$ is spherical by Lemma 2.1. Since $`\delta (R_1,R_K)`$ is shorter than $`\delta (R_J,R_K)`$ by construction, the desired result follows from an immediate induction. โ
###### Corollary 2.8.
Let $`\mathrm{\Delta }`$ be a (possibly twin) building of type $`(W,S)`$ and let $`R_K`$ be a spherical residue which is maximal with respect to that property. Then, for any residue $`R_J`$ (assumed to be spherical if $`\mathrm{\Delta }`$ is a twin building and if $`R_J`$ and $`R_K`$ have opposite signs) the projection of $`R_J`$ on $`R_K`$ is properly contained in $`R_K`$ unless $`R_J`$ and $`R_K`$ are equal or opposite.
###### Proof.
Since $`\mathrm{proj}_{R_K}(R_J)`$ and $`\mathrm{proj}_{R_J}(R_K)`$ are parallel, the result clearly follows from the previous proposition, using also Lemma 2.3 if $`R_J`$ and $`R_K`$ have opposite signs. โ
###### Corollary 2.9.
Let $`\mathrm{\Delta }`$ be a twin building and let $`\mathrm{\Sigma }`$ be a twin apartment of $`\mathrm{\Delta }`$. Then the parallelism is an equivalence relation on the set of spherical residues of $`\mathrm{\Sigma }`$.
###### Proof.
This follows from Lemma 2.3 and Proposition 2.7. โ
#### 2.2.5 Twin roots
Let $`\mathrm{\Delta }=(\mathrm{\Delta }_+,\mathrm{\Delta }_{},\delta ^{})`$ be a twin building of type $`(W,S)`$.
A twin root of $`\mathrm{\Delta }`$ is the convex hull of a pair of chambers โat codistance $`1`$โ, namely a pair $`\{x,y\}`$ such that $`s:=\delta ^{}(x,y)S`$. Let $`\pi `$ be the $`s`$-panel containing $`x`$. Then any chamber $`x^{}\pi \backslash \{x\}`$ is opposite $`y`$ and determines therefore a twin apartment $`\mathrm{\Sigma }`$ which contains $`\varphi `$ and $`x^{}`$ (see Section 2.2.3). We say that $`\varphi `$ is a twin root of $`\mathrm{\Sigma }`$. The complement of $`\varphi `$ in $`\mathrm{\Sigma }`$ is also a twin root; it is actually the convex hull of $`x^{}`$ and $`\mathrm{proj}_\pi (x^{})`$. This twin root is said to be opposite to $`\varphi `$ in $`\mathrm{\Sigma }`$ and is denoted by $`\varphi `$ although its definition depends on $`\mathrm{\Sigma }`$. A residue $`R`$ of $`\mathrm{\Delta }`$ is said to be in the interior of $`\varphi `$ if it is contained in $`\mathrm{\Sigma }`$ and if $`R\mathrm{\Sigma }`$ is contained in $`\varphi `$. If $`R\varphi `$ and $`R(\varphi )`$ are both nonempty, then $`R`$ is said to be on the boundary of $`\varphi `$.
### 2.3 From groups to buildings: twin root data
#### 2.3.1 Definition
Let $`\mathrm{\Sigma }=\mathrm{\Sigma }(W,S)`$ be the canonical twin building of type $`(W,S)`$ (see Section 2.2.2) and let $`\mathrm{\Phi }(W,S)`$ be the set of all its twin roots. We have already mentioned the action of $`W`$ on $`\mathrm{\Sigma }`$ (see Section 2.2.2). Given a twin root $`\varphi \mathrm{\Phi }(W,S)`$, then all panels on the boundary of $`\varphi `$ correspond to the same reflection of $`W`$. This reflection is denoted by $`s_\varphi `$ and it permutes $`\varphi `$ and $`\varphi `$. A pair $`\{\varphi ,\psi \}`$ of twin roots of $`\mathrm{\Sigma }=(\mathrm{\Sigma }_+,\mathrm{\Sigma }_{})`$ is said to be prenilpotent if $`\varphi \psi \mathrm{\Sigma }_+`$ and $`(\varphi )(\psi )\mathrm{\Sigma }_+`$ are both nonempty; in that case, we denote by $`[\varphi ,\psi ]`$ the set of all twin roots $`\alpha `$ of $`\mathrm{\Sigma }`$ such that $`\alpha \varphi \psi `$ and $`\alpha (\varphi )(\psi )`$.
A twin root datum of type $`(W,S)`$ is a system $`๐:=(G,(U_\varphi )_{\varphi \mathrm{\Phi }(W,S)})`$ consisting of a group $`G`$ and a family of subgroups $`U_\varphi `$ which satisfies the following axioms, where $`H`$ and $`U(c)`$ denote respectively the intersection of the normalizers of all $`U_\varphi `$โs and the subgroup of $`G`$ generated by the $`U_\varphi `$โs such that $`\varphi `$ contains the chamber $`c`$ of $`\mathrm{\Sigma }`$ :
(TRD0) $`U_\varphi 1`$ for all $`\varphi \mathrm{\Phi }(W,S)`$; (TRD1) if $`\{\varphi ,\psi \}`$ is a prenilpotent pair of distinct twin roots, the commutator $`[U_\varphi ,U_\psi ]`$ is contained in the group generated by all $`U_\gamma `$โs with $`\gamma [\varphi ,\psi ]\backslash \{\varphi ,\psi \}`$; (TRD2) if $`\varphi \mathrm{\Phi }(W,S)`$ and $`uU_\varphi \backslash \{1\}`$, there exists elements $`u^{},u^{\prime \prime }`$ of $`U_\varphi `$ such that the product $`\mu (u)=u^{}uu^{\prime \prime }`$ conjugates $`U_\psi `$ onto $`U_{s_\varphi (\psi )}`$ for each $`\psi \mathrm{\Phi }(W,S)`$; (TRD3) if $`\varphi \mathrm{\Phi }(W,S)`$ and $`c`$ is a chamber of $`\mathrm{\Sigma }`$ which is not contained in $`\varphi `$, then $`U_\varphi `$ is not contained in $`U(c)`$; (TRD4) the group $`G`$ is generated by $`H`$ and the $`U_\varphi `$โs. The group $`G`$ is sometimes denoted by $`G^๐`$.
#### 2.3.2 Isomorphisms of twin root data
Let $`๐:=(G,(U_\varphi )_{\varphi \mathrm{\Phi }(W,S)})`$ and $`๐^{}:=(G^{},(U_\varphi ^{})_{\varphi \mathrm{\Phi }(W^{},S^{})})`$ be twin root data. Let $`S=S_1\mathrm{}S_n`$ be the finest partition of $`S`$ such that $`[S_i,S_j]=1`$ whenever $`1i<jn`$. Then $`๐`$ and $`๐^{}`$ are called isomorphic if there exist an isomorphism $`\phi :GG^{}`$, an isomorphism $`\pi :WW^{}`$ with $`\pi (S)=S^{}`$, an element $`xG^{}`$ and a sign $`ฯต_i`$ for each $`1in`$ such that
$$x\phi (U_\varphi )x^1=U_{ฯต_i\pi (\varphi )}^{}\text{for every twin root }\varphi \text{ with }s_\varphi W_{S_i}\text{,}$$
(1)
where we denote by $`\pi `$ the obvious bijection $`\mathrm{\Phi }(W,S)\mathrm{\Phi }(W^{},S)`$ induced by $`\pi :WW^{}`$. Thus, if $`(W,S)`$ is irreducible, then either $`x\phi (U_\varphi )x^1=U_{\pi (\varphi )}^{}`$ or $`x\phi (U_\varphi )x^1=U_{\pi (\varphi )}^{}`$ for all $`\varphi \mathrm{\Phi }(W,S)`$.
When (1) holds, we say that the isomorphism $`\phi `$ induces an isomorphism of $`๐`$ to $`๐^{}`$. In particular, this means, the isomorphism $`\phi `$ maps the union of conjugacy classes
$$\{gU_+g^1|gG\}\{gU_{}g^1|gG\}$$
to
$$\{gU_+^{}g^1|gG^{}\}\{gU_{}^{}g^1|gG\},$$
where $`U_+`$ (resp. $`U_{}`$, $`U_+^{}`$, $`U_{}^{}`$) denotes $`U(c)`$ (resp. $`U(\mathrm{op}_\mathrm{\Sigma }(c))`$, $`U(c^{})`$, $`U(\mathrm{op}_\mathrm{\Sigma }^{}(c^{}))`$) for some $`c\mathrm{\Sigma }=\mathrm{\Sigma }(W,S)`$ and some $`c^{}\mathrm{\Sigma }^{}=\mathrm{\Sigma }(W^{},S^{})`$.
A crucial fact on isomorphisms between twin root data we will need later is the following.
###### Proposition 2.10.
Let $`๐:=(G,(U_\varphi )_{\varphi \mathrm{\Phi }(W,S)})`$ and $`๐^{}:=(G^{},(U_\varphi ^{})_{\varphi \mathrm{\Phi }(W^{},S^{})})`$ be twin root data with $`S`$ and $`S^{}`$ finite and let $`\phi :GG^{}`$ be an isomorphism. Assume there exists $`gG^{}`$ such that
$$\{\phi (U_\varphi )|\varphi \mathrm{\Phi }(W,S)\}=\{gU_\varphi ^{}g^1|\varphi \mathrm{\Phi }(W^{},S^{})\}.$$
Then $`\phi `$ induces an isomorphism of $`๐`$ to $`๐^{}`$.
###### Proof.
This is Theorem (2.5) in . โ
#### 2.3.3 Twin buildings from twin root data
Let $`๐:=(G,(U_\varphi )_{\varphi \mathrm{\Phi }(W,S)})`$ be a twin root datum of type $`(W,S)`$. Let $`H`$ be the intersection of the normalizers of all $`U_\varphi `$โs and let $`N`$ be the subgroup of $`G`$ generated by $`H`$ together with all $`\mu (u)`$ such that $`uU_\varphi \backslash \{1\}`$, where $`\mu (u)`$ is as in (TRD2). Let $`c`$ be a chamber of $`\mathrm{\Sigma }=\mathrm{\Sigma }(W,S)`$ of positive sign, let $`c^{}:=\mathrm{op}_\mathrm{\Sigma }(c)`$ and let $`B_+:=H.U(c)`$ and $`B_{}:=H.U(c^{})`$.
We recall from , Proposition 4, that $`(G,B_+,N)`$ and $`(G,B_{},N)`$ are both $`BN`$-pairs of type $`(W,S)`$. Thus, we have corresponding Bruhat decompositions of $`G`$:
$$G=\underset{wW}{}B_+wB_+\text{and}G=\underset{wW}{}B_{}wB_{}.$$
For each $`ฯต\{+,\}`$, the set $`\mathrm{\Delta }_ฯต:=G/B_ฯต`$ endowed with the map $`\delta _ฯต:\mathrm{\Delta }_ฯต\times \mathrm{\Delta }_ฯตW`$ by
$$\delta _ฯต(gB_ฯต,hB_ฯต)=wB_ฯตg^1hB_ฯต=B_ฯตwB_ฯต,$$
has a canonical structure of a thick building of type $`(W,S)`$.
The twin root datum axioms imply that $`G`$ also admits Birkhoff decompositions (see Lemma 1 in ):
$$G=\underset{wW}{}B_ฯตwB_ฯต$$
for each $`ฯต\{+,\}`$. The pair $`((\mathrm{\Delta }_+,\delta _+),(\mathrm{\Delta }_{},\delta _{}))`$ of buildings admits a natural twinning by means of the $`W`$-codistance $`\delta ^{}`$ defined by
$$\delta ^{}(gB_ฯต,hB_ฯต)=wB_ฯตg^1hB_ฯต=B_ฯตwB_ฯต$$
for each $`ฯต\{+,\}`$. The triple $`\mathrm{\Delta }:=((\mathrm{\Delta }_+,\delta _+),(\mathrm{\Delta }_{},\delta _{}),\delta ^{})`$ is a twin building of type $`(W,S)`$.
We may and shall identify the chamber $`c`$ (resp. $`c^{}`$) of $`\mathrm{\Sigma }(W,S)`$ with the chamber $`B_+`$ of $`\mathrm{\Delta }_+=G/B_+`$ (resp. $`B_{}`$ of $`\mathrm{\Delta }_{}`$). We also identify $`\mathrm{\Sigma }(W,S)`$ with the unique twin apartment of $`\mathrm{\Delta }`$ containing $`c`$ and $`c^{}`$; this twin apartment is denoted by $`\mathrm{\Sigma }`$ and is called the fundamental twin apartment of $`\mathrm{\Delta }`$ (with respect to the twin root datum $`๐`$).
The diagonal action of $`G`$ on $`\mathrm{\Delta }_+\times \mathrm{\Delta }_{}`$ by left multiplication is transitive on pairs of opposite chambers and, hence, on twin apartments.
#### 2.3.4 Parabolic subgroups and root subgroups
We keep the notation of the previous subsection. We recall from the theory of $`BN`$-pairs (see , Chapter IV) that a subgroup $`P`$ of $`G`$ containing $`B_ฯต`$ or any of its conjugates is called a parabolic subgroup of sign $`ฯต`$, where $`ฯต\{+,\}`$. If $`P`$ contains $`B_ฯต`$, then there exists $`JS`$ such that $`P`$ has a Bruhat decomposition
$$P=\underset{wW_J}{}B_ฯตwB_ฯต;$$
the set $`J`$ is called the type of the parabolic subgroup $`P`$. If $`J`$ is spherical, then $`P`$ is said to be of finite type (or of spherical type). A minimal parabolic subgroup (i.e. a parabolic subgroup of type $`\mathrm{}`$) such as $`B_+`$ or $`B_{}`$ is called a Borel subgroup.
For $`ฯต\{+,\}`$, let $`P_ฯต^J`$ be the parabolic subgroup of type $`J`$ containing $`B_ฯต`$. The geometric meaning of the groups $`B_+`$, $`B_{}`$, $`P_+^J`$, $`P_{}^J`$, $`H`$ and $`N`$ is as follows:
$$B_+=\mathrm{Stab}_G(c),B_{}=\mathrm{Stab}_G(c^{}),P_+^J=\mathrm{Stab}_G(\mathrm{Res}_J(c)),P_{}^J=\mathrm{Stab}_G(\mathrm{Res}_J(c^{}))$$
and
$$N=\mathrm{Stab}_G(\mathrm{\Sigma }),H=B_+N=B_{}N=\mathrm{Fix}_G(\mathrm{\Sigma }).$$
Given a twin root $`\varphi `$ of $`\mathrm{\Sigma }`$, then $`U_\varphi `$ fixes chamberwise any panel in the interior of $`\varphi `$ and is sharply transitive on the set of twin apartments containing $`\varphi `$. Moreover, it follows then from the axioms that for each $`gG`$ and each twin root $`\varphi `$ of $`\mathrm{\Sigma }`$, the group $`U_{g(\varphi )}:=gU_\varphi g^1`$ depends only on the twin root $`g(\varphi )`$ and not on the choice of $`g`$ and $`\varphi `$. Hence, for every twin root $`\psi `$ of $`\mathrm{\Delta }`$, there is a well defined group $`U_\psi `$ which fixes chamberwise any panel in the interior of $`\psi `$. The group $`U_\psi `$ is sharply transitive on the set of twin apartments containing $`\psi `$; it is called the root subgroup associated with the twin root $`\psi `$.
#### 2.3.5 The Conditions (P1)โ(P3) and a technical lemma
Let $`๐:=(G,(U_\varphi )_{\varphi \mathrm{\Phi }(W,S)})`$ be a twin root datum. For each $`\varphi \mathrm{\Phi }(W,S)`$, we set $`L_\varphi :=U_\varphi U_\varphi `$ and $`H_\varphi :=N_{L_\varphi }(U_\varphi )N_{L_\varphi }(U_\varphi )`$. The group $`H_\varphi `$ acts on the conjugacy class $`๐_\varphi `$ of $`U_\varphi `$ in $`L_\varphi `$. We shall be interested in the following three conditions (see Theorem 5.1 below) :
(P1) for every $`\varphi \mathrm{\Phi }(W,S)`$, the group $`U_\varphi `$ is nilpoptent; (P2) for every $`\varphi \mathrm{\Phi }(W,S)`$, the group $`L_\varphi `$ is perfect; (P3) for every $`\varphi \mathrm{\Phi }(W,S)`$, the groups $`U_\varphi `$ and $`U_\varphi `$ are the only fixed points of $`H_\varphi `$ in $`๐_\varphi `$.
The following lemma gives the geometric interpretation of Condition (P3).
###### Lemma 2.11.
Let $`๐=(G,(U_\varphi )_{\varphi \mathrm{\Phi }(W,S)})`$ be a twin root datum of type $`(W,S)`$ which satisfies Condition (P3). Let $`\mathrm{\Delta }`$ be the twin building associated with $`๐`$ and let $`H`$ be as above (see Section 2.3.3). Then $`H`$ fixes no chamber outside the fundamental twin apartment of $`\mathrm{\Delta }`$.
###### Proof.
Let $`\mathrm{\Sigma }`$ be the fundamental twin apartment of $`\mathrm{\Delta }`$. Let $`\varphi \mathrm{\Phi }`$ and let $`H_\varphi `$ be the intersection of the normalizers of $`U_\varphi `$ and $`U_\varphi `$ in $`U_\varphi U_\varphi `$. The group $`H_\varphi `$ fixes $`\mathrm{\Sigma }`$ chamberwise and is therefore contained in $`H`$. Let $`\pi `$ be any panel on the boundary of $`\varphi `$. Then $`\pi `$ is stabilized by $`U_\varphi `$ and $`U_\varphi `$. The condition (P3) means precisely that the only fixed chambers of $`H_\varphi `$ in $`\pi `$ are the two elements of $`\pi \mathrm{\Sigma }`$. This implies that for any panel $`\pi `$ of $`\mathrm{\Sigma }`$, the group $`H`$ fixes no chamber in $`\pi \backslash \mathrm{\Sigma }`$. Now, an easy induction on the gallery distance from an arbitrary chamber of $`\mathrm{\Delta }`$ to $`\mathrm{\Sigma }`$ finishes the proof. โ
#### 2.3.6 Twin root data over fields
Let $`๐=(G,(U_\varphi )_{\varphi \mathrm{\Phi }(W,S)})`$ be a twin root datum of type $`(W,S)`$, let $`HG`$ be as in the previous subsection and for each $`\varphi \mathrm{\Phi }(W,S)`$, let $`๐_\varphi `$ be a field. We recall from that the twin root datum $`๐`$ is called locally split over the fields $`(๐_\varphi )_{\varphi \mathrm{\Phi }(W,S)}`$ if $`H`$ is abelian and if for each $`\varphi \mathrm{\Phi }(W,S)`$, the twin root datum $`๐_\varphi :=(HU_\varphi U_\varphi ,\{U_\varphi ,U_\varphi \})`$ is isomorphic to the natural twin root datum of $`SL_2(๐_\varphi )`$ or $`PSL_2(๐_\varphi )`$. Of course, the natural twin root datum associated to a (split) Kac-Moody group over a field $`๐`$ is locally split over $`๐`$.
Notice that if $`๐`$ is locally split over the fields $`(๐_\varphi )_{\varphi \mathrm{\Phi }(W,S)}`$, then $`๐`$ satisfies Condition (P1). Moreover, if $`๐_\varphi `$ has cardinality at least 4 for every $`\varphi \mathrm{\Phi }(W,S)`$, then Conditions (P2) and (P3) are also satisfied.
## 3 Levi decomposition in twin root data
The purpose of this section is to obtain a Levi decomposition for intersections of finite type parabolic subgroups of opposite signs in a group with twin root datum (see Proposition 3.6). In the language of buildings, this group is the stabilizer of a pair of spherical residues of opposite signs.
### 3.1 Levi decomposition of parabolic subgroups
#### 3.1.1 The setting
Let $`๐:=(G,(U_\alpha )_{\alpha \mathrm{\Phi }(W,S)})`$ be a twin root datum, let $`\mathrm{\Delta }=((\mathrm{\Delta }_+,\delta _+),(\mathrm{\Delta }_{},\delta _{}),\delta ^{})`$ be the corresponding twin building and let $`\mathrm{\Sigma }_0`$ be the fundamental twin apartment (see Section 2.3.3).
Let $`\mathrm{\Sigma }`$ be any twin apartment of $`\mathrm{\Delta }`$. Let $`c\mathrm{\Sigma }`$ be a chamber and let $`R`$ be a spherical residue of $`\mathrm{\Sigma }`$ (i.e. $`R\mathrm{\Sigma }\mathrm{}`$). Let $`\mathrm{\Phi }^\mathrm{\Sigma }`$ the set of all twin roots of $`\mathrm{\Sigma }`$ and let $`\mathrm{\Phi }^\mathrm{\Sigma }(R)`$ be the set of twin roots $`\beta `$ of $`\mathrm{\Sigma }`$ such that $`R\beta `$ and $`R(\beta )`$ are both nonempty, which means precisely that the reflection $`s_\beta `$ stabilizes $`R`$ (see Sections 2.2.2 and 2.2.5). We set
$$U^\mathrm{\Sigma }(c):=U_\varphi |\varphi \mathrm{\Phi }^\mathrm{\Sigma },U^\mathrm{\Sigma }(R):=\underset{x\mathrm{\Sigma }R}{}U^\mathrm{\Sigma }(x)$$
and
$$L^\mathrm{\Sigma }(R):=\mathrm{Fix}_G(\mathrm{\Sigma }).U_\varphi |\varphi \mathrm{\Phi }^\mathrm{\Sigma }(R).$$
(See Section 2.3.4 for the definition of $`U_\varphi `$.)
We will see in Proposition 3.1 that $`U^\mathrm{\Sigma }(c)`$ and $`U^\mathrm{\Sigma }(R)`$ are actually independent of $`\mathrm{\Sigma }`$. They will be denoted by $`U(c)`$ and $`U(R)`$ respectively.
Notice that $`L^\mathrm{\Sigma }(R)=L^\mathrm{\Sigma }(\mathrm{op}_\mathrm{\Sigma }(R))`$ since $`\mathrm{\Phi }^\mathrm{\Sigma }(R)=\mathrm{\Phi }^\mathrm{\Sigma }(\mathrm{op}_\mathrm{\Sigma }(R))`$. Moreover, if the residue $`R`$ is reduced to a chamber $`c`$ (i.e. if $`R`$ is of type $`\mathrm{}`$) then $`L^\mathrm{\Sigma }(c)=\mathrm{Fix}_G(\mathrm{\Sigma })`$, which is $`G`$-conjugate to the subgroup $`H`$ of Section 2.3.3 (see Section 2.3.4).
#### 3.1.2 Standard Levi decomposition: Levi factor and unipotent radical
The following result is the standard Levi decomposition of finite type parabolic subgroups of a group with a twin root datum. We state it in the language of buildings.
###### Proposition 3.1.
We have $`\mathrm{Stab}_G(R)=L^\mathrm{\Sigma }(R)U^\mathrm{\Sigma }(R)`$. Moreover, $`U(R)`$ is sharply transitive on the set of residues which are opposite $`R`$ in $`\mathrm{\Delta }(G)`$ and
$$L^\mathrm{\Sigma }(R)=\mathrm{Stab}_G(R)\mathrm{Stab}_G(\mathrm{op}_\mathrm{\Sigma }(R)).$$
In particular, the subgroup $`U^\mathrm{\Sigma }(R)`$ is independent of $`\mathrm{\Sigma }`$ and will be denoted by $`U(R)`$.
###### Proof.
This follows from the theorem of Section 6.2.2 in . โ
The group $`U(R)`$ is called the unipotent radical of the parabolic subgroup $`\mathrm{Stab}_G(R)`$ with respect to the twin root datum $`๐`$ and the group $`L^\mathrm{\Sigma }(R)`$ is called a Levi factor.
### 3.2 Levi decomposition of parabolic intersections
#### 3.2.1 More definitions and notation
We keep the notation of Section 3.1.1.
Let $`\mathrm{\Sigma }`$ be a twin apartment of $`\mathrm{\Delta }`$ and for each $`ฯต\{+,\}`$, let $`R_ฯต`$ be a residue of $`\mathrm{\Sigma }`$ of sign $`ฯต`$. We set
$$U^\mathrm{\Sigma }(R_+,R_{}):=U_\beta |\beta \mathrm{\Phi }^\mathrm{\Sigma }\text{ and }R_ฯต\mathrm{\Sigma }\beta ,$$
and, for $`ฯต\{+,\}`$,
$$\stackrel{~}{U}^\mathrm{\Sigma }(R_ฯต,R_ฯต):=U_\beta |\beta \mathrm{\Phi }^\mathrm{\Sigma }(R_ฯต)\text{ and }\mathrm{proj}_{R_ฯต}(R_ฯต)\beta $$
(see Sections 2.2.5 and 2.3.4). Notice that $`U^\mathrm{\Sigma }(R_+,R_{})=U^\mathrm{\Sigma }(R_{},R_+)`$ while $`\stackrel{~}{U}^\mathrm{\Sigma }(R_+,R_{})`$ need not equal $`\stackrel{~}{U}^\mathrm{\Sigma }(R_{},R_+)`$. If $`R_+`$ (resp. $`R_{}`$) is reduced to a chamber $`x_+`$ (resp. $`x_{}`$), we have $`\stackrel{~}{U}^\mathrm{\Sigma }(R_+,R_{})=\stackrel{~}{U}^\mathrm{\Sigma }(R_{},R_+)=\{1\}`$. This remains true in the case where $`R_+`$ and $`R_{}`$ are parallel (see Section 2.2.4).
See Remark 3.4 below for an interpretation of the group $`\stackrel{~}{U}^\mathrm{\Sigma }(R_ฯต,R_ฯต)`$.
#### 3.2.2 The intersection of a parabolic subgroup and unipotent radical
In order to obtain the Levi decomposition of the group $`\mathrm{Stab}_G(\{R_+,R_{}\})=\mathrm{Stab}_G(R_+)\mathrm{Stab}_G(R_{})`$ , we need a decomposition result for $`\mathrm{Stab}_G(R_+)U(R_{})`$. We start with the case where both residues are reduced to single chambers.
###### Lemma 3.2.
Let $`x_+\mathrm{\Delta }_+`$ and $`x_{}\mathrm{\Delta }_{}`$ be chambers of $`\mathrm{\Delta }`$. Let $`\mathrm{\Sigma }`$ be a twin apartment containing them both. Let $`ฯต\{+,\}`$, let $`y_ฯต:=\mathrm{op}_\mathrm{\Sigma }(x_ฯต)`$ and let $`x_ฯต=x_0,x_1,\mathrm{},x_n=y_ฯต`$ be a minimal gallery joining $`x_ฯต`$ to $`y_ฯต`$. For each $`1in`$, let $`\beta _i`$ be the twin root of $`\mathrm{\Sigma }`$ containing $`x_{i1}`$ but not $`x_i`$. Then we have
$$U(x_ฯต)\mathrm{Stab}_G(x_ฯต)=U^\mathrm{\Sigma }(x_+,x_{})=U_{\beta _1}.U_{\beta _2}\mathrm{}U_{\beta _n}.$$
In particular, the product $`U_{\beta _1}.U_{\beta _2}\mathrm{}U_{\beta _n}`$ is a group which coincides with $`U^\mathrm{\Sigma }(x_+,x_{})`$, and the latter does not depend on the twin apartment $`\mathrm{\Sigma }`$. We will denote it by $`U(x_+,x_{})`$.
###### Proof.
This follows from Lemma 1.5.2(iii) and Theorem 3.5.4 in . โ
The following lemma generalizes Lemma 3.2 to the case of spherical residues of higher rank.
###### Lemma 3.3.
Let $`R_+\mathrm{\Delta }_+`$ and $`R_{}\mathrm{\Delta }_{}`$ be spherical residues of $`\mathrm{\Delta }`$. Let $`\mathrm{\Sigma }`$ be a twin apartment intersecting them both. Then, for each $`ฯต\{+,\}`$, we have
$$U(R_ฯต)\mathrm{Stab}_G(R_ฯต)=\stackrel{~}{U}^\mathrm{\Sigma }(R_ฯต,R_ฯต).U^\mathrm{\Sigma }(R_+,R_{}).$$
In particular, if $`R_+`$ and $`R_{}`$ are parallel, then
$$U(R_+)\mathrm{Stab}_G(R_{})=\mathrm{Stab}_G(R_+)U(R_{})=U^\mathrm{\Sigma }(R_+,R_{}).$$
###### Proof.
The inclusion โ$``$โ of the first part is clear.
Let $`R_ฯต^{}:=\mathrm{op}_\mathrm{\Sigma }(R_ฯต)`$. Let us choose $`z,z^{}\mathrm{proj}_{R_ฯต^{}}(R_ฯต)\mathrm{\Sigma }`$ such that $`z`$ and $`z^{}`$ are opposite in the spherical residue $`\mathrm{proj}_{R_ฯต^{}}(R_ฯต)`$ (see Section 2.1.5). Set $`x:=\mathrm{op}_\mathrm{\Sigma }(z)`$ and $`x^{}:=\mathrm{op}_\mathrm{\Sigma }(z^{})`$. Notice that $`x`$ and $`x^{}`$ belong to $`R_ฯต\mathrm{\Sigma }`$. We also define $`y:=\mathrm{proj}_{R_ฯต}(x)`$.
We have $`U(R_ฯต)U(x^{})`$ since $`x^{}R_ฯต`$. Moreover, $`U(R_ฯต)\mathrm{Stab}_G(R_ฯต)`$ fixes $`y=\mathrm{proj}_{R_ฯต}(x)`$ because this group fixes $`x`$ and stabilizes $`R_ฯต`$; hence, $`U(R_ฯต)\mathrm{Stab}_G(R_ฯต)\mathrm{Stab}_G(y)`$. Therefore, we have $`U(R_ฯต)\mathrm{Stab}_G(R_ฯต)U(x^{})\mathrm{Stab}_G(y)=U(x^{},y)`$, where the latter equality follows Lemma 3.2.
We now choose a minimal gallery $`y=y_0,\mathrm{},y_j=\mathrm{proj}_{R_ฯต}(z^{}),\mathrm{},y_n=z^{}`$ joining $`y`$ to $`z^{}`$ (see Section 2.1.4). Hence, for all $`0in`$, we have $`y_iR_ฯต`$ if and only if $`ij`$. (Notice that $`j=0`$ or $`j=n`$ is possible.)
For each $`1in`$, let $`\beta _i`$ be the twin root of $`\mathrm{\Sigma }`$ which contains $`y_{i1}`$ but not $`y_i`$. Thus $`\{\beta _1,\mathrm{},\beta _n\}`$ is the set of twin roots containing $`x^{}`$ and $`y`$ or equivalently, $`y`$ but not $`z^{}`$ since $`z^{}=\mathrm{op}_\mathrm{\Sigma }(x^{})`$. By definition, we have $`U(x^{},y)=U_{\beta _i}|1in`$. By Lemma 3.2 this group coincides with the product $`U_{\beta _1}.U_{\beta _2}\mathrm{}U_{\beta _n}`$.
Now we observe that by the definition of $`y`$, $`y_j`$ and $`z^{}`$ and by Lemmas 2.3 and 2.4, we have
$$\mathrm{proj}_{R_ฯต^{}}(y)=\mathrm{proj}_{R_ฯต^{}}(y_j)=z^{}\text{and}\mathrm{proj}_{\mathrm{proj}_{R_ฯต}(R_ฯต)}(z^{})=\mathrm{proj}_{\mathrm{proj}_{R_ฯต}(R_ฯต)}(y_j)=y.$$
In view of the properties of projections (see Section 2.1.4), this implies
$$\begin{array}{ccc}\hfill \{\beta \mathrm{\Phi }^\mathrm{\Sigma }(R_ฯต)|\mathrm{proj}_{R_ฯต}(R_ฯต)\beta \}& =& \{\beta \mathrm{\Phi }^\mathrm{\Sigma }|y\beta \text{ and }y_j\beta \}\hfill \\ & =& \{\beta _1,\mathrm{},\beta _j\}\hfill \\ & =& \{\beta \mathrm{\Phi }^\mathrm{\Sigma }|y\beta \text{ and }\mathrm{op}_\mathrm{\Sigma }(y_j)\beta \}\hfill \end{array}$$
and
$$\begin{array}{ccc}\hfill \{\beta \mathrm{\Phi }^\mathrm{\Sigma }|R_ฯต\mathrm{\Sigma }\beta \text{ and }R_ฯต\mathrm{\Sigma }\beta \}& =& \{\beta \mathrm{\Phi }^\mathrm{\Sigma }|R_ฯต\mathrm{\Sigma }\beta \text{ and }R_ฯต^{}\mathrm{\Sigma }\beta \}\hfill \\ & =& \{\beta \mathrm{\Phi }^\mathrm{\Sigma }|y_j\beta \text{ and }z^{}\beta \}\hfill \\ & =& \{\beta _{j+1},\mathrm{},\beta _n\}\hfill \\ & =& \{\beta \mathrm{\Phi }^\mathrm{\Sigma }|y_j\beta \text{ and }x^{}\beta \}.\hfill \end{array}$$
We deduce from this, together with Lemma 3.2, that
$$\stackrel{~}{U}^\mathrm{\Sigma }(R_ฯต,R_ฯต)=U_i|\mathrm{\hspace{0.33em}1}ij=U(y,\mathrm{op}_\mathrm{\Sigma }(y_j))=U_1\mathrm{}U_j$$
and
$$U^\mathrm{\Sigma }(R_ฯต,R_ฯต)=U_i|j+1in=U(x^{},y_j)=U_{j+1}\mathrm{}U_n.$$
In summary, we have shown that
$$\begin{array}{ccc}\hfill U(R_ฯต)\mathrm{Stab}_G(R_ฯต)& & U(x^{},y)\hfill \\ & =& \left(U_{\beta _1}\mathrm{}U_{\beta _j}\right).\left(U_{\beta _{j+1}}\mathrm{}U_{\beta _n}\right)\hfill \\ & =& \stackrel{~}{U}^\mathrm{\Sigma }(R_ฯต,R_ฯต).U^\mathrm{\Sigma }(R_+,R_{}).\hfill \end{array}$$
The lemma implies that, if for $`ฯต=+`$ or $``$ we have $`\mathrm{proj}_{R_ฯต}(R_ฯต)=R_ฯต`$ (in particular, if $`R_+`$ and $`R_{}`$ are parallel) and hence $`\stackrel{~}{U}^\mathrm{\Sigma }(R_ฯต,R_ฯต)=\{1\}`$, then the group $`U^\mathrm{\Sigma }(R_+,R_{})`$ is independent of the twin apartment $`\mathrm{\Sigma }`$. In that case, we may omit the superscript $`\mathrm{\Sigma }`$ and we shall write $`U(R_+,R_{})`$ rather than $`U^\mathrm{\Sigma }(R_+,R_{})`$.
###### Remark 3.4.
The group $`\stackrel{~}{U}^\mathrm{\Sigma }(R_ฯต,R_ฯต)`$ defined in the previous lemma actually coincides with $`U(\mathrm{proj}_{R_ฯต}(R_ฯต))L^\mathrm{\Sigma }(R_ฯต)`$. This can be seen as follows. First notice that $`(L^\mathrm{\Sigma },(U_\beta )_{\beta \mathrm{\Phi }^\mathrm{\Sigma }(R_ฯต)})`$ is a twin root datum. Hence the group $`\mathrm{Stab}_{L^\mathrm{\Sigma }(R_ฯต)}(\mathrm{proj}_{R_ฯต}(R_ฯต))`$ has a Levi decomposition in $`L^\mathrm{\Sigma }(R_ฯต)`$ by Proposition 3.1. The above claim is easily deduced from that fact: actually, the group $`\stackrel{~}{U}^\mathrm{\Sigma }(R_ฯต,R_ฯต)=U(\mathrm{proj}_{R_ฯต}(R_ฯต))L^\mathrm{\Sigma }(R_ฯต)`$ is nothing but the unipotent radical of $`\mathrm{Stab}_{L^\mathrm{\Sigma }(R_ฯต)}(\mathrm{proj}_{R_ฯต}(R_ฯต))`$ with respect to the above-mentioned twin root datum. We will not need that fact here.
The following is a consequence of the proof of the previous lemma.
###### Corollary 3.5.
Let $`R_+\mathrm{\Delta }_+`$ and $`R_{}\mathrm{\Delta }_{}`$ be spherical residues. Then, for each $`ฯต\{+,\}`$ there exist chambers $`x_+R_+`$ and $`x_{}R_{}`$ such that $`U(R_ฯต)\mathrm{Stab}_G(R_ฯต)=U(x_+,x_{})`$. In particular, if all root subgroups are nilpotent (i.e. if Condition (P1) holds), then $`U^\mathrm{\Sigma }(R_+,R_{})`$ is nilpotent, where $`\mathrm{\Sigma }`$ is any twin apartment intersecting $`R_+`$ and $`R_{}`$.
###### Proof.
The first statement was proved along the way. The second statement is a consequence of the first, using also Axiom (TRD1) and Lemma 3.2. โ
#### 3.2.3 The intersection of finite type parabolics of opposite signs
We are now able to prove the Levi decomposition of intersections of finite type parabolic subgroups of opposite signs.
###### Proposition 3.6.
Let $`R_+\mathrm{\Delta }_+`$ and $`R_{}\mathrm{\Delta }_{}`$ be spherical residues of $`\mathrm{\Delta }(G)`$. Let $`\mathrm{\Sigma }`$ be a twin apartment containing $`R_+`$ and $`R_{}`$. For each sign $`ฯต`$ set $`R_ฯต^{}:=\mathrm{proj}_{R_ฯต}(R_ฯต)`$. Then, for all $`ฯต,ฯต^{}\{+,\}`$ we have
$$\begin{array}{ccc}\hfill \mathrm{Stab}_G(R_+)\mathrm{Stab}_G(R_{})& =& L^\mathrm{\Sigma }(R_ฯต^{}^{})U(R_+^{},R_{}^{})\hfill \\ & =& L^\mathrm{\Sigma }(R_ฯต^{}^{})(\stackrel{~}{U}^\mathrm{\Sigma }(R_ฯต,R_ฯต).U^\mathrm{\Sigma }(R_+,R_{}).\stackrel{~}{U}^\mathrm{\Sigma }(R_ฯต,R_ฯต)).\hfill \end{array}$$
###### Proof.
Notice first that $`\mathrm{Stab}_G(R_+)\mathrm{Stab}_G(R_{})=\mathrm{Stab}_G(R_+^{})\mathrm{Stab}_G(R_{}^{})`$. Moreover, since $`R_+^{}`$ and $`R_{}^{}`$ are parallel, we deduce from Lemma 2.3 and Proposition 2.7(ii) that $`L^\mathrm{\Sigma }(R_+^{})=L^\mathrm{\Sigma }(R_{}^{})`$.
Now the inclusion $`L^\mathrm{\Sigma }(R_ฯต^{}^{}).U(R_+^{},R_{}^{})\mathrm{Stab}_G(R_+)\mathrm{Stab}_G(R_{})`$ is clear. On the other hand, we have the following:
$$\begin{array}{ccc}\hfill \mathrm{Stab}_G(R_+^{})\mathrm{Stab}_G(R_{}^{})& =& (L^\mathrm{\Sigma }(R_ฯต^{}).U(R_ฯต^{}))\mathrm{Stab}_G(R_ฯต^{})\hfill \\ & & L^\mathrm{\Sigma }(R_ฯต^{}).\left(U(R_ฯต^{})\mathrm{Stab}_G(R_ฯต^{})\right)\hfill \\ & =& L^\mathrm{\Sigma }(R_ฯต^{}^{}).U(R_ฯต^{},R_ฯต^{}),\hfill \end{array}$$
where the last equality follows from Lemma 3.3. This proves that $`\mathrm{Stab}_G(R_+)\mathrm{Stab}_G(R_{})=L^\mathrm{\Sigma }(R_ฯต^{}^{}).U(R_ฯต^{},R_ฯต^{})`$.
Now $`L^\mathrm{\Sigma }(R_ฯต^{}^{})=L^\mathrm{\Sigma }(R_ฯต^{})`$ intersects $`U(R_ฯต^{})`$ trivially and normalizes that group by Proposition 3.1. Since $`L^\mathrm{\Sigma }(R_ฯต^{}^{})\mathrm{Stab}_G(R_ฯต^{})`$ we also see that $`L^\mathrm{\Sigma }(R_ฯต^{}^{})`$ normalizes $`\mathrm{Stab}_G(R_ฯต^{})`$. Therefore, $`L^\mathrm{\Sigma }(R_ฯต^{}^{})`$ normalizes $`U(R_ฯต^{})\mathrm{Stab}_G(R_ฯต^{})=U(R_ฯต^{},R_ฯต^{})`$ and intersects the latter group trivially. This proves the first equality of the lemma.
In order to establish the second equality, we first notice that $`\stackrel{~}{U}^\mathrm{\Sigma }(R_ฯต,R_ฯต^{})=\stackrel{~}{U}^\mathrm{\Sigma }(R_ฯต,R_ฯต)`$ by the definition of these groups. Now, (an argument as in) the proof of the previous lemma shows that
$$\begin{array}{ccc}\hfill U(R_+^{},R_{}^{})& =& \stackrel{~}{U}^\mathrm{\Sigma }(R_ฯต,R_ฯต^{}).U(R_ฯต^{},R_ฯต)\hfill \\ & =& \stackrel{~}{U}^\mathrm{\Sigma }(R_ฯต,R_ฯต).U^\mathrm{\Sigma }(R_ฯต,R_ฯต).\stackrel{~}{U}^\mathrm{\Sigma }(R_ฯต,R_ฯต),\hfill \end{array}$$
from which the conclusion follows. โ
###### Remark 3.7.
The preceding proposition is proved in , Section 6.3.4, under an additional assumption called (NILP), defined in *op. cit.*, Section 6.3. Our proof shows that this extra assumption is not necessary for the result to hold.
###### Corollary 3.8.
Let $`(W,S)`$ be a Coxeter system such that $`S`$ is finite. Let $`๐=(G,(U_\varphi )_{\varphi \mathrm{\Phi }(W,S)})`$ be a twin root datum such that each $`U_\varphi `$ is finite and such that $`H:=_{\varphi \mathrm{\Phi }(W,S)}N_G(U_\varphi )`$ is finite. Then the set of all bounded subgroups coincides with the set of all finite subgroups of $`G`$.
###### Proof.
Let $`\mathrm{\Delta }`$ be the twin building associated with $`๐`$. The fact that each finite subgroup of $`G`$ is bounded is an immediate consequence of Proposition 2.2. In order to prove that a bounded subgroup is finite, it suffices to prove that given a pair $`R_+`$, $`R_{}`$ of spherical residues of opposite signs, then the group $`\mathrm{Stab}_G(R_+)\mathrm{Stab}_G(R_{})`$ is finite. Our hypotheses imply that every Levi factor of spherical type is finite. Hence, by Proposition 3.6, it just remains to show that $`U(\mathrm{proj}_{R_+}(R_{}),\mathrm{proj}_R_{}(R_+))`$ is finite. But this follows again from our hypotheses in view of Lemma 3.2 and Corollary 3.5. โ
## 4 Maximal bounded subgroups
### 4.1 The main characterization
Let $`๐=(G,(U_\alpha )_{\alpha \mathrm{\Phi }})`$ be a twin root datum of type $`(W,S)`$. By a maximal bounded subgroup of $`G`$, we mean a bounded subgroup which is not properly contained in any other bounded subgroup of $`G`$. Let $`M`$ be such a bounded subgroup. By definition $`M`$ is the intersection of two finite type parabolic subgroups of opposite signs. The following theorem shows that there exists two canonical finite type parabolic subgroups $`P_+^M`$ and $`P_{}^M`$ such that $`M=P_+^MP_{}^M`$. Case (i) of the theorem corresponds to the case where $`P_+^M`$ and $`P_{}^M`$ are opposite; the group $`M`$ is then the common Levi factor of $`P_+^M`$ and $`P_{}^M`$. To some extent, this is the generic case (see Proposition 4.2 and Remark 4.3 below).
Before stating the theorem, we need one more notation. Let $`((\mathrm{\Delta }_+,\delta _+),(\mathrm{\Delta }_{},\delta _{}),\delta ^{})`$ be the twin building associated with $`๐`$. Given a subgroup $`MG`$, we denote by $`๐ฎ_ฯต(M)`$ the set of all spherical residues of $`\mathrm{\Delta }_ฯต`$ stabilized by $`M`$, where $`ฯต\{+,\}`$.
###### Theorem 4.1.
Let $`(G,(U_\alpha )_{\alpha \mathrm{\Phi }})`$ be a twin root datum of type (W,S) and let $`MG`$ be a maximal bounded subgroup. Then one of the following holds :
(i) for $`ฯต\{+,\}`$ the set $`๐ฎ_ฯต(M)`$ consists of a unique element $`R_ฯต`$, which is a maximal spherical residue of $`\mathrm{\Delta }_ฯต`$; moreover $`R_+`$ and $`R_{}`$ are opposite in $`\mathrm{\Delta }(G)`$; (ii) for $`ฯต\{+,\}`$ the set $`๐ฎ_ฯต(M)`$ possesses two distinguished elements $`R_ฯต`$ and $`\overline{R}_ฯต`$ such that for every $`T_ฯต๐ฎ_ฯต(M)`$ we have $`R_ฯตT_ฯต\overline{R}_ฯต`$ and $`\mathrm{proj}_{T_ฯต}(T_ฯต)=R_ฯต`$ and $`T_+`$ and $`T_{}`$ are not opposite; moreover $`\overline{R}_ฯต`$ is the only maximal spherical residue containing $`R_ฯต`$. In both cases we have $`M=\mathrm{Stab}_G(R_+)\mathrm{Stab}_G(R_{})`$.
Conversely, let $`R_+\mathrm{\Delta }_+`$ and $`R_{}\mathrm{\Delta }_{}`$ be spherical residues such that either of the following conditions holds :
(i) $`R_+`$ and $`R_{}`$ are maximal spherical and opposite; (ii) $`R_+`$ and $`R_{}`$ are parallel; moreover, for each $`ฯต\{+,\}`$ the residue $`R_ฯต`$ is properly contained in a unique maximal spherical residue $`\overline{R}_ฯต`$ and we have $`proj_{\overline{R}_ฯต}(\overline{R}_ฯต)=R_ฯต`$. Then $`M:=\mathrm{Stab}_G(R_+)\mathrm{Stab}_G(R_{})`$ is a maximal bounded subgroup.
###### Proof.
Let $`MG`$ be a maximal bounded subgroup. For $`ฯต\{+,\}`$ let $`R_ฯต๐ฎ_ฯต(M)`$.
Assume first that $`R_+`$ and $`R_{}`$ are opposite. Hence, $`M=\mathrm{Stab}_G(R_+)\mathrm{Stab}_G(R_{})`$ which implies by Proposition 3.1 that $`M`$ does not stabilize any residue properly contained in $`R_+`$ or $`R_{}`$. Moreover, the maximality of $`M`$ implies that $`R_+`$ and $`R_{}`$ are maximal spherical residues. Let now $`T_ฯต๐ฎ_ฯต(M)`$. Then $`M`$ stabilizes $`\mathrm{proj}_{R_+}(T_ฯต)`$ and $`\mathrm{proj}_R_{}(T_ฯต)`$. Since these cannot be properly contained in $`R_+`$ and $`R_{}`$ respectively, we conclude from Corollary 2.8 that $`T_ฯต=R_ฯต`$. Hence, we are in Case (i). Notice that our discussion also proves the converse statement in the case (i).
We now assume that $`R_+`$ and $`R_{}`$ are not opposite. For $`ฯต\{+,\}`$ we have $`\mathrm{proj}_{R_ฯต}(R_ฯต)๐ฎ_ฯต(M)`$. Moreover, we know by Lemma 2.5 that $`\mathrm{proj}_{R_+}(R_{})`$ and $`\mathrm{proj}_R_{}(R_+)`$ are not opposite (this follows also from the first part of the present proof). Therefore, up to replacing $`R_ฯต`$ by $`\mathrm{proj}_{R_ฯต}(R_ฯต)`$ we may assume that $`R_+`$ and $`R_{}`$ are parallel. Let $`J_+`$ and $`J_{}`$ be their respective types. Since $`R_+`$ and $`R_{}`$ are not opposite, Lemma 2.3 and Proposition 2.7 show that $`J_ฯต`$ is not a maximal spherical subset of $`S`$.
Let now $`sS\backslash J_ฯต`$ such that $`J_ฯต\{s\}`$ is spherical and let $`R_ฯต^s`$ be the residue of type $`J_ฯต\{s\}`$ containing $`R_ฯต`$. Set also $`T:=\mathrm{proj}_{R_ฯต^s}(R_ฯต)`$. We now prove that $`T=R_ฯต`$.
Assume on the contrary that $`TR_ฯต`$. Then Lemma 2.6 and Proposition 2.7 imply that $`R_ฯต`$ and $`T`$ are opposite in $`R_ฯต^s`$. Let $`\mathrm{\Sigma }`$ be a twin apartment containing $`R_+`$ and $`R_{}`$. There exists a twin root $`\alpha =(\alpha _+,\alpha _{})`$ of $`\mathrm{\Sigma }`$ such that $`T\alpha _ฯต`$, $`R_ฯต\alpha _ฯต`$ (hence the reflection $`s_\alpha `$ of $`\mathrm{\Sigma }`$ stabilizes $`R_ฯต^s`$) and $`R_ฯต\alpha _ฯต`$. Now Proposition 3.1 implies that $`U_\alpha `$ acts freely on the residues opposite $`T`$ in $`R_ฯต^s`$. In particular, we have $`U_\alpha M=\{1\}`$ since $`M`$ stabilizes $`R_ฯต`$. Therefore, the group $`U_\alpha M`$ contains $`M`$ properly, and it stabilizes $`R_ฯต^s`$ and $`R_ฯต`$ by construction. In other words, we have $`MU_\alpha M`$ is a bounded subgroup, which contradicts the maximality of $`M`$. This proves that $`T=R_ฯต`$ as claimed.
Let $`S_ฯต`$ be the set of all $`sS\backslash J_ฯต`$ such that $`J_ฯต\{s\}`$ is spherical. Let $`\overline{R}_ฯต`$ be the residue of type $`J_ฯตS_ฯต`$ containing $`R_ฯต`$. We now prove that $`\overline{R}_ฯต`$ is spherical.
To this end we consider, as before, a twin apartment $`\mathrm{\Sigma }`$ containing $`R_+`$ and $`R_{}`$. Let $`R_ฯต^{}:=\mathrm{op}_\mathrm{\Sigma }(R_ฯต)`$. Choose a chamber $`xR_ฯต`$ and a chamber $`z`$ which is opposite to $`\mathrm{proj}_{R_ฯต^{}}(x)`$ in $`R_ฯต^{}`$. Set $`y:=\mathrm{proj}_{\overline{R}_ฯต}(z)`$. Our aim is to apply the criterion of sphericity of Lemma 2.1 to $`x`$ and $`y`$. Hence, let $`jJ_ฯตS_ฯต`$ and denote by $`\pi _j`$ the $`j`$-panel containing $`x`$.
We know from Lemma 2.3 that $`R_ฯต`$ and $`R_ฯต^{}`$ are parallel. By Lemma 2.6, this implies that $`R_ฯต`$ and $`\mathrm{proj}_{\overline{R}_ฯต}(R_ฯต^{})`$ are parallel. By Lemma 2.4, we see that $`x`$ and $`\mathrm{proj}_{R_ฯต}(y)`$ are opposite in $`R_ฯต`$ and we conclude from Lemma 2.1 that if $`jJ_ฯต`$ then $`x\mathrm{proj}_{\pi _j}(y)=\mathrm{proj}_{\pi _j}(\mathrm{proj}_{R_ฯต}(y))`$.
Let us now assume that $`jS_ฯต`$. We have already proved that $`\mathrm{proj}_{R_ฯต^j}(R_ฯต)=R_ฯต`$, which implies that $`\mathrm{proj}_{R_ฯต^j}(R_ฯต^{})`$ is opposite $`R_ฯต`$ in $`R_ฯต^j`$ by Lemma 2.3. Therefore, $`x`$ and $`\mathrm{proj}_{R_ฯต^j}(y)`$ are opposite in $`R_ฯต^j`$ which implies by Lemma 2.1 that $`x\mathrm{proj}_{\pi _j}(y)=\mathrm{proj}_{\pi _j}(\mathrm{proj}_{R_ฯต^j}(y))`$.
Finally, Lemma 2.1 applied to $`x`$ and $`y`$ insures that $`\overline{R}_ฯต`$ is spherical, i.e. that $`J_ฯตS_ฯต`$ is a spherical subset of $`S`$. Moreover, it is clear by the definition of $`S_ฯต`$ that $`J_ฯตS_ฯต`$ is actually a *maximal* spherical subset of $`S`$.
By maximality of $`M`$, we have $`M=\mathrm{Stab}_G(R_+)\mathrm{Stab}_G(R_{})`$. Therefore, we see by Proposition 3.6 that $`M`$ does not stabilize any proper residue of $`R_ฯต`$ for $`ฯต\{+,\}`$. Moreover, the same result implies that if $`R`$ is a residue contained in $`\overline{R}_ฯต`$, then $`R`$ is stabilized by $`M`$ if and only if $`R`$ contains $`R_ฯต`$.
Let now $`T_ฯต๐ฎ_ฯต(M)`$. Then $`\mathrm{proj}_{R_ฯต}(T_ฯต)๐ฎ_ฯต(M)`$ and so $`\mathrm{proj}_{R_ฯต}(T_ฯต)=R_ฯต`$. Therefore, $`R_ฯต`$ and $`\mathrm{proj}_{T_ฯต}(R_ฯต)`$ are parallel. By the previous paragraph together with Proposition 2.7, this implies that $`R_ฯต=\mathrm{proj}_{T_ฯต}(R_ฯต)`$, or in other words that $`R_ฯตT\overline{R}_ฯต`$.
Going now back to the first argument in the proof of (ii) above, we conclude moreover that $`\mathrm{proj}_{T_ฯต}(T_ฯต)=R_ฯต`$. The fact that $`T_+`$ and $`T_{}`$ are not opposite is now obvious.
It remains to prove the converse statement. For (ii), let $`R_+`$, $`R_{}`$, $`\overline{R}_+`$ and $`\overline{R}_{}`$ be as in (ii) and define $`M:=\mathrm{Stab}_G(R_+)\mathrm{Stab}_G(R_{})`$. Mimicking some of the arguments above we can prove again that if $`R`$ is a residue contained in $`\overline{R}_ฯต`$, then $`R`$ is stabilized by $`M`$ if and only if $`R`$ contains $`R_ฯต`$. From this, we deduce as above that $`R_ฯตT\overline{R}_ฯต`$ whenever $`T๐ฎ_ฯต(M)`$. Therefore, if $`MM_1`$ and $`M_1`$ is bounded, then there exists $`T_ฯต๐ฎ_ฯต(M_1)`$ such that $`R_ฯตT_ฯต\overline{R}_ฯต`$. But our hypotheses then imply that $`\mathrm{proj}_{T_ฯต}(T_ฯต)=R_ฯต`$. Therefore, we have $`R_ฯต๐ฎ_ฯต(M_1)`$ namely $`M_1\mathrm{Stab}_G(R_ฯต)`$. Hence $`M=M_1`$. The proof is complete. โ
### 4.2 Two specializations
#### 4.2.1 Obstructions for Case (ii) of Theorem 4.1
In many interesting situations, only Case (i) of Theorem 4.1 occurs. The next result gives sufficient conditions on the Coxeter system $`(W,S)`$ which imply that Case (ii) never happens.
###### Proposition 4.2.
Assume that the Coxeter system $`(W,S)`$ satisfies one of the following conditions :
(R1) for all $`s,tS`$, the order of $`st`$ is not equal to 3; (R2) for any pair $`J,K`$ of spherical subsets of $`S`$ such that $`J`$ is properly contained in $`K`$ and $`K`$ is maximal spherical, there exists an $`sS\backslash K`$ such that $`J\{s\}`$ is spherical but $`K\{s\}`$ is not; (R3) for every $`jS`$, there exists a unique maximal spherical subset $`J`$ of $`S`$ such that $`jJ`$.
Then the case (ii) does not occur in the previous theorem.
###### Proof.
We keep the notation of the proof of Theorem 4.1. It is clear that if (ii) holds in that theorem, then (R2) fails by choosing $`J=J_ฯต`$ and $`K=J_ฯตS_ฯต`$ where $`ฯต=+`$ or $``$.
Now let $`\mathrm{\Sigma }`$ be a twin apartment containing $`R_+`$ and $`R_{}`$ as in the proof above, let $`R_ฯต^{}:=\mathrm{op}_\mathrm{\Sigma }(R_ฯต)`$ and let $`J_ฯต^{\prime \prime }`$ be the type of $`R_ฯต^{\prime \prime }:=proj_{\overline{R}_ฯต}(R_ฯต^{})`$. We know by Lemma 2.3 that $`R_ฯต^{\prime \prime }`$ is opposite $`R_ฯต`$ in $`\overline{R}_ฯต`$ and, moreover, that $`\overline{R}_ฯต`$ is the *unique* maximal spherical residue containing $`R_ฯต`$. On the other hand, since $`\overline{R}_+`$ and $`\overline{R}_{}`$ are not opposite, we have $`R_ฯต^{}R_ฯต^{\prime \prime }`$, and these two distinct residues are parallel. Therefore, it follows from Proposition 2.7 that there exists $`sS\backslash (J_ฯตS_ฯต)`$ such that $`J_ฯต^{\prime \prime }\{s\}`$ is spherical. In particular, $`J_ฯตS_ฯต`$ is *not* the unique maximal spherical subset of $`S`$ containing $`J_ฯต^{\prime \prime }`$. Hence, (R3) fails and furthermore, we have $`J_ฯตJ_ฯต^{\prime \prime }`$. Since $`J_ฯต`$ and $`J_ฯต^{\prime \prime }`$ are the respective types of two opposite residues of $`\overline{R}_ฯต`$, the latter inequality also implies that (R1) fails, using , Proposition 5.2.3 and the fact that there are no Moufang $`n`$-gons for odd $`n`$ greater than 3. โ
###### Remark 4.3.
1. Condition (R2) in the previous corollary is also equivalent to the following :
(R2) for every non-maximal spherical subset $`J`$ of $`S`$ there exists (at least) two distinct maximal spherical subsets $`K_1`$ and $`K_2`$ of $`S`$ containing $`J`$.
2. All affine and compact hyperbolic Coxeter diagrams satisfy Condition (R2) (notice that (R2) is empty for $`\stackrel{~}{A}_1`$ and so obviously satisfied; actually $`\stackrel{~}{A}_1`$ also satisfies (R1) and (R3)).
3. Condition (R3) in the previous corollary is also equivalent to each of the following ones :
(R3) for every maximal spherical subset $`J`$ of $`S`$ and all pairs $`j,s`$ with $`jJ`$ and $`sS\backslash J`$, the order of $`sj`$ is infinite; (R3<sup>โฒโฒ</sup>) there is a partition $`S=S_1\mathrm{}S_n`$ of $`S`$ into spherical subsets such that the order of $`st`$ is infinite whenever $`sS_i`$, $`tS_j`$ and $`ij`$.
#### 4.2.2 A group theoretic description of Case (ii)
We end this section with a lemma which is will be used in the proof of Theorem 5.1.
###### Lemma 4.4.
Let $`(G,(U_\alpha )_{\alpha \mathrm{\Phi }})`$ be a twin root datum of type (W,S) which satisfies Conditions (P1) and (P2) of Section 2.3.5. Let $`MG`$ be a maximal bounded subgroup. A necessary and sufficient condition for $`M`$ to have type (ii) in Theorem 4.1, is that
$$M=U(MM^{}),$$
where $`U`$ is a nontrivial nilpotent group and $`M^{}`$ is a maximal bounded subgroup different from $`M`$.
###### Proof.
We first prove that the condition is necessary. We keep the notation of Theorem 4.1 and assume that $`M`$ satisfies Condition (ii). Let $`\mathrm{\Sigma }`$ denote a twin apartment containing $`R_+`$ and $`R_{}`$. By Proposition 3.6, we have
$$\begin{array}{ccc}\hfill M& =& \mathrm{Stab}_G(R_+)\mathrm{Stab}_G(R_{})\hfill \\ & =& \mathrm{Stab}_G(\overline{R}_+)\mathrm{Stab}_G(R_{})\hfill \\ & =& L^\mathrm{\Sigma }(R_+).\stackrel{~}{U}^\mathrm{\Sigma }(\overline{R}_+,R_{}).U(\overline{R}_+,R_{}).\stackrel{~}{U}^\mathrm{\Sigma }(R_{},\overline{R}_+)\hfill \\ & =& L^\mathrm{\Sigma }(R_+).\stackrel{~}{U}^\mathrm{\Sigma }(\overline{R}_+,R_{}).U(\overline{R}_+,R_{}),\hfill \end{array}$$
where the last equality follows from $`\mathrm{proj}_R_{}(\overline{R}_+)=R_{}`$ which implies $`\stackrel{~}{U}^\mathrm{\Sigma }(R_{},\overline{R}_+)=\{1\}`$.
Now, set $`U:=U(\overline{R}_+,R_{})`$ and $`M^{}:=L^\mathrm{\Sigma }(\overline{R}_+)`$. By Theorem 4.1, the group $`M^{}`$ is a maximal bounded subgroup. Clearly, the group $`U`$ is nilpotent by 3.5 and nontrivial since $`\mathrm{op}_\mathrm{\Sigma }(R_{})\overline{R}_+=\mathrm{}`$ (see the proof of Proposition 4.2). Moreover, we have $`M^{}U=\{1\}`$ since $`UU(\overline{R}_+)`$. On the other hand, we also have $`L^\mathrm{\Sigma }(R_+).\stackrel{~}{U}^\mathrm{\Sigma }(\overline{R}_+,R_{})M^{}`$ by definition. Therefore, we have $`L^\mathrm{\Sigma }(R_+).\stackrel{~}{U}^\mathrm{\Sigma }(\overline{R}_+,R_{})=MM^{}`$ and hence, $`M=(MM^{}).U`$.
It remains to prove that $`MM^{}`$ normalizes $`U`$. Using again the fact that $`\stackrel{~}{U}^\mathrm{\Sigma }(R_{},\overline{R}_+)`$ is trivial, we deduce from Lemma 3.3 that $`U=U(\overline{R}_+)\mathrm{Stab}_G(R_{})`$. Now the desired conclusion follows since $`M`$ normalizes $`\mathrm{Stab}_G(R_{})`$ (because $`M\mathrm{Stab}_G(R_{})`$) and $`M^{}`$ normalizes $`U(\overline{R}_+)`$ (by Proposition 3.1).
We now show that the condition is sufficient. Assume that $`M`$ has type (i) in Theorem 4.1 and that the condition of the lemma is satisfied. Let $`R_+`$ (resp. $`R_{}`$) be the unique element of $`๐ฎ_+(M)`$ (resp. $`๐ฎ_{}(M)`$) and let $`T_+๐ฎ_+`$. Since $`M`$ and $`M^{}`$ are distinct (or since $`U`$ is nontrivial), the residues $`R_+`$ and $`T_+`$ are distinct. Now, $`MM^{}`$ stabilizes $`R:=\mathrm{proj}_{R_+}(T_+)`$, which is properly contained in $`R_+`$ by Corollary 2.8. In particular, the group $`U`$ does not act trivially on $`R_+`$ since $`M=(MM^{}).U`$ does not stabilize $`R`$.
As before, let $`\mathrm{\Sigma }`$ be a twin apartment intersecting $`R_+`$ and $`R_{}`$. Then $`M=L^\mathrm{\Sigma }(R_+)`$ and we know that $`(M,(U_\alpha )_{\alpha \mathrm{\Phi }^\mathrm{\Sigma }(R_+)})`$ is a twin root datum of spherical type $`(W_J,J)`$ for some $`JS`$. Since $`U`$ is normal in $`M`$ and since it does not act trivially on $`R_+`$, we deduce from , Theorem 5 of Chapter IV that there exists a residue $`R^{}R_+`$ such that $`U`$ contains the group $`M_1`$ generated by all $`U_\beta `$ with $`\beta \mathrm{\Phi }^\mathrm{\Sigma }(R^{})`$. But $`\beta \mathrm{\Phi }^\mathrm{\Sigma }(R^{})`$ implies $`\beta \mathrm{\Phi }^\mathrm{\Sigma }(R^{})`$ and therefore $`M_1`$ is generated by subgroups of the form $`U_\beta U_\beta `$. Since each group of the latter form is perfect by hypothesis, the group $`M_1`$ itself is perfect which contradicts the fact that it is contained in the nilpotent group $`U`$. This concludes the proof of the lemma. โ
## 5 The reduction theorem for isomorphisms which preserve bounded subgroups
In this section we state and prove a general theorem concerning isomorphisms between groups endowed with twin root data, which preserve bounded subgroups. Roughly speaking, it says that, under certain conditions, the isomorphism problem for groups with twin root data reduces to the isomorphism problem for groups of finite type. The main results of this paper will be deduced from it in the following two sections.
### 5.1 $`๐`$-rigidity of twin root data
In order to make the statement of this theorem as precise and concise as possible, we introduce some additional terminology.
Let $`๐`$ be a collection of twin root data. A twin root datum $`๐`$ is called $`๐`$-rigid if the following holds (see Section 2.3.1 for the definition of $`G^๐`$):
If $`๐^{}๐`$, then any isomorphism of $`G^๐`$ to $`G^๐^{}`$ induces an isomorphism of $`๐`$ to $`๐^{}`$.
Let $`๐=(G,(U_\alpha )_{\alpha \mathrm{\Phi }(W,S)})`$ be a twin root datum, let $`\mathrm{\Delta }`$ be the associated twin building, let $`\mathrm{\Sigma }`$ be the fundamental twin apartment and let $`c`$ be a chamber of $`\mathrm{\Sigma }`$. For each subset $`J`$ of $`S`$, we set $`L_J:=L^\mathrm{\Sigma }(\mathrm{Res}_J(c))`$ and $`๐_J:=(L_J,(U_\alpha )_{\alpha \mathrm{\Phi }^\mathrm{\Sigma }(\mathrm{Res}_J(c))})`$.
Given a collection $`๐`$ as above, then we denote by $`๐_{\mathrm{sph}}`$ the collection of all twin root data of the form $`๐_J`$ such that $`๐๐`$ is of type $`(W^๐,S^๐)`$ and $`J`$ is a maximal spherical subset of $`S^๐`$. A twin root datum $`๐๐`$ of type $`(W^๐,S^๐)`$ is called $`๐`$-locally rigid if for every every maximal spherical subset $`J`$ of $`S^๐`$, the twin root datum $`๐_J`$ is $`๐_{\mathrm{sph}}`$-rigid.
### 5.2 The result
###### Theorem 5.1.
Let $`๐`$ and $`๐^{}`$ be two twin root data which satisfy Conditions (P1), (P2) and (P3), and whose types are Coxeter systems of finite rank. Assume that $`๐`$ is $`\{๐^{}\}`$-locally rigid. If $`\xi :G^๐G^๐^{}`$ is an isomorphism which maps bounded subgroups of $`G^๐`$ to bounded subgroups of $`G^๐^{}`$, then $`\xi `$ induces an isomorphism of $`๐`$ to $`๐^{}`$.
###### Proof.
Let $`M`$ be a maximal bounded subgroup of $`G^๐`$. Then $`\xi (M)`$ is a maximal bounded subgroup of $`G^๐^{}`$ by the hypothesis on $`\xi `$. Moreover, it follows from Lemma 4.4 that $`M`$ and $`\xi (M)`$ have the same โtypeโ, where the โtypeโ of $`M`$, resp. $`\xi (M)`$ is given by either Case (i) or (ii) in Theorem 4.1.
Let now $`\mathrm{\Sigma }`$ be a twin apartment of the twin building $`\mathrm{\Delta }`$ associated with $`๐`$ and let $`\alpha `$ be a twin root of $`\mathrm{\Sigma }`$. Let also $`\mathrm{\Delta }^{}`$ be the twin building associated with $`๐^{}`$. Choose a maximal residue of spherical type $`R`$ intersecting $`\mathrm{\Sigma }`$ and such that $`\alpha \mathrm{\Phi }^\mathrm{\Sigma }(R)`$. By what we have just seen, we know that $`\xi (L^\mathrm{\Sigma }(R))`$ has the form $`L^\mathrm{\Sigma }^{}(R^{})`$ for some apartment $`\mathrm{\Sigma }^{}`$ of $`\mathrm{\Delta }^{}`$ and some maximal residue of spherical type $`R^{}`$ which intersects $`\mathrm{\Sigma }^{}`$.
Since $`๐`$ is locally rigid, the restriction of $`\xi `$ to $`M=L^\mathrm{\Sigma }(R)`$ induces an isomorphism from the twin root datum $`(L^\mathrm{\Sigma }(R),(U_\beta )_{\beta \mathrm{\Phi }^\mathrm{\Sigma }(R)})`$ to the twin root datum $`(L^\mathrm{\Sigma }^{}(R^{}),(U_\beta ^{})_{\beta \mathrm{\Phi }^\mathrm{\Sigma }^{}(R^{})})`$, where we have used superscript to denote root groups of $`๐^{}`$. We may assume without loss of generality that
$$\{\xi (U_\beta )|\beta \mathrm{\Phi }^\mathrm{\Sigma }(R)\}=\{U_\beta ^{}|\beta \mathrm{\Phi }^\mathrm{\Sigma }^{}(R^{})\}.$$
Let $`H=\mathrm{Fix}_{G^๐}(\mathrm{\Sigma })`$ and let $`H^{}:=\xi (H)`$. Since $`H=_{\beta \mathrm{\Phi }^\mathrm{\Sigma }(R)}N_{L^\mathrm{\Sigma }(R)}(U_\beta )`$ (see Section 2.3.4) we deduce $`H^{}=_{\beta \mathrm{\Phi }^\mathrm{\Sigma }^{}(R^{})}N_{L^\mathrm{\Sigma }^{}(R^{})}(U_\beta ^{})`$. This implies that $`H^{}`$ is the chamberwise stabilizer of $`\mathrm{\Sigma }^{}`$ in $`G(๐^{})`$.
Let now $`\gamma `$ be a twin root of $`\mathrm{\Sigma }`$ which does not belong to $`\mathrm{\Phi }^\mathrm{\Sigma }(R)`$. Arguing as for $`\alpha `$, we obtain that $`\xi (U_\gamma )=U_\gamma ^{}^{}`$ where $`\gamma ^{}`$ is a twin root of $`\mathrm{\Delta }^{}`$ which is contained in a twin apartment $`\mathrm{\Sigma }^{\prime \prime }`$ whose chamberwise stabilizer is $`H^{}`$. On the other hand, our hypotheses imply that $`H^{}`$ fixes a unique twin apartment chamberwise (see Lemma 2.11). In summary, we have shown that
$$\{\xi (U_\beta )|\beta \text{ is a twin root of }\mathrm{\Sigma }\}=\{U_\beta ^{}|\beta \text{ is a twin root of }\mathrm{\Sigma }^{}\}.$$
Now the conclusion follows from Proposition 2.10. โ
## 6 Kac-Moody groups over arbitrary fields
In this section and in the following one, we apply Theorem 5.1 to the case of Kac-Moody groups over fields in the strict sense.
It is known that a Kac-Moody group $`G`$ over a field $`๐`$ naturally yields a twin root datum $`๐=(G,(U_\alpha )_{\alpha \mathrm{\Phi }})`$ which is locally split over $`๐`$ (namely which is locally split over $`(๐_\alpha )_{\alpha \mathrm{\Phi }}`$, where $`๐_\alpha =๐`$ for each $`\alpha \mathrm{\Phi }`$). We have also mentioned that if $`๐`$ has cardinality at least 4, then the conditions (P1), (P2) and (P3) of Section 2.3.5 are satisfied. In order to apply Theorem 5.1 to $`G`$, it remains to discuss the local rigidity of the twin root datum $`๐`$. This is done by using the classical theorems on isomorphisms between Chevalley groups, but the arguments are slightly different according as the ground field is finite or infinite.
### 6.1 Finite fields vs. infinite fields
The following result gives a handy criterion which distinguishes between these two cases.
###### Proposition 6.1.
Let $`G`$ be a Kac-Moody group over a field $`๐`$. Then $`G`$ is finitely generated if and only if $`๐`$ is finite.
###### Proof.
Let $`๐=(G,(U_\alpha )_{\alpha \mathrm{\Phi }})`$ be the twin root datum which is naturally associated with $`G`$. Let $`H:=_{\alpha \mathrm{\Phi }}N_G(U_\alpha )`$ and let $`\mathrm{\Pi }\mathrm{\Phi }`$ be the (finite) set of simple roots in $`\mathrm{\Phi }=\mathrm{\Phi }(W,S)`$. It is known (and easy to see) that $`G`$ is generated by the set $`S_๐:=H_{\alpha \mathrm{\Pi }}U_\alpha `$. Moreover, each $`U_\alpha `$ is isomorphic to the additive group $`๐`$ and the group $`H`$ is a โsplit $`๐`$-torusโ, namely it is isomorphic to a direct product of finitely many copies of the multiplicative group $`๐^\times `$.
If $`๐`$ is finite, then $`S_๐`$ is finite, whence $`G`$ is finitely generated.
If $`G`$ is finitely generated, then $`G`$ is generated by a subset of $`S_๐`$. By , the defining relations satisfied by the elements of $`S_๐`$ in the group $`G`$ involve only the ring structure of $`๐`$. Since no infinite field is a finitely generated ring, we deduce that $`๐`$ has to be finite. โ
### 6.2 The characteristic in the case of a finite ground field
The following result will spare us to worry about the exceptional isomorphisms between finite Chevalley groups.
###### Proposition 6.2.
Let $`G`$ be an infinite Kac-Moody group over a finite field $`๐`$ of characteristic $`p`$. Let $`q`$ be a prime. Then $`p=q`$ if and only if the set of orders of finite $`q`$-subgroups of $`G`$ has no finite upper bound.
###### Proof.
Let $`๐=(G,(U_\alpha )_{\alpha \mathrm{\Phi }})`$ be the twin root datum which is naturally associated with $`G`$.
Assume first that $`p=q`$. We must show that $`G`$ possesses finite $`p`$-subgroups of arbitrary large orders. Since $`๐`$ is finite, the group $`U_\alpha `$, which is isomorphic to the additive group $`๐`$, is finite for every $`\alpha \mathrm{\Phi }`$. Since $`G`$ is infinite by assumption, we deduce that the Coxeter group $`W`$ is infinite. Let $`\mathrm{\Delta }=(\mathrm{\Delta }_+,\mathrm{\Delta }_{},\delta ^{})`$ be the twin building associated with $`๐`$ and let $`\mathrm{\Sigma }`$ be the fundamental twin apartment. Since $`\mathrm{\Delta }`$ is of non-spherical type, we can find chambers $`x_+`$ and $`x_{}`$ of $`\mathrm{\Sigma }`$ such that $`n:=\mathrm{}(\delta ^{}(x_+,x_{}))`$ is arbitrarily large. On the other hand, we know by Lemma 3.2 that the group $`U(x_+,x_{})`$ may be written as a product of the form $`U_{\beta _1}.U_{\beta _2}\mathrm{}U_{\beta _n}`$ for certain twin roots $`\beta _1,\mathrm{},\beta _n`$ of $`\mathrm{\Sigma }`$. Since $`U_{\beta _i}`$ is a finite $`p`$-group for each $`1in`$, it follows that $`U(x_+,x_{})`$ is a finite $`p`$-group of order at least $`p^n`$ which yields the desired result.
We now assume that $`pq`$. We must show that there is an upper-bound on the possible orders of finite $`q`$-subgroups of $`G`$. Let $`QG`$ be such a finite $`q`$-group. By Proposition 2.2, the group $`Q`$ is a bounded subgroup. Let $`R_+\mathrm{\Delta }_+`$ and $`R_{}\mathrm{\Delta }_{}`$ be spherical residues which are stabilized by $`Q`$. Up to replacing $`R_+`$ and $`R_{}`$ by $`\mathrm{proj}_{R_+}(R_{})`$ and $`\mathrm{proj}_R_{}(R_+)`$ respectively, we may assume that $`R_+`$ and $`R_{}`$ are parallel. Let $`\mathrm{\Sigma }_Q`$ be a twin apartment containing $`R_+`$ and $`R_{}`$. By Proposition 3.6, we have $`\mathrm{Stab}_G(R_+)\mathrm{Stab}_G(R_{})=L^{\mathrm{\Sigma }_Q}(R_+)U(R_+,R_{})`$. Hence there is a homomorphism $`f:\mathrm{Stab}_G(R_+)\mathrm{Stab}_G(R_{})L^{\mathrm{\Sigma }_Q}(R_+)`$. On the other hand, the order of every element of $`U(R_+,R_{})`$ is a power of $`p`$ by Lemma 3.2 and Corollary 3.5. Since $`pq`$ we deduce that $`Q`$ and $`f(Q)`$ are isomorphic and, hence, we may assume that $`Q`$ is contained in $`L^{\mathrm{\Sigma }_Q}(R_+)`$. Now, up to conjugation by an element of $`G`$, we may assume that $`L^{\mathrm{\Sigma }_Q}(R_+)=L_J`$ for some spherical subset $`J`$ of $`S`$. Thus $`|Q|\mathrm{max}\{|L_J||JS\text{ spherical}\}`$ (note that each $`L_J`$ is a finite Chevalley group over $`๐`$). The desired conclusion follows from the finiteness of $`S`$. โ
### 6.3 Isomorphisms of Kac-Moody groups
We are now able to apply Theorem 5.1 to Kac-Moody groups over fields. The following theorem is the main result of the introduction.
###### Theorem 6.3.
Let $`G`$ and $`G^{}`$ be infinite Kac-Moody groups over fields $`๐`$ and $`๐^{}`$ respectively, both of cardinality at least 4, and let $`๐`$ and $`๐^{}`$ be the corresponding twin root data. Let $`\xi :GG^{}`$ be a group isomorphism which maps bounded subgroups of $`G`$ to bounded subgroups of $`G^{}`$. Then $`\xi `$ induces an isomorphism of $`๐`$ to $`๐^{}`$.
###### Proof.
We have to show that $`\xi `$ induces an isomorphism of $`๐`$ to $`๐^{}`$. By Theorem 5.1, it suffices to show that $`๐`$ is $`\{๐^{}\}`$-locally rigid. Let $`(W,S)`$ and $`(W^{},S^{})`$ be the respective types of $`๐`$ and $`๐^{}`$, let $`J`$ and $`J^{}`$ be maximal spherical subsets of $`S`$ and $`S^{}`$ and let $`\phi :L_JL_J^{}`$ be a group isomorphism. We have to show that $`\phi `$ induces an isomorphism of $`๐_J`$ to $`๐_J^{}^{}`$.
By the definition of a Kac-Moody group, we know that $`L_J`$ and $`L_J^{}`$ are Chevalley groups over $`๐`$ and $`๐^{}`$ respectively. Up to replacing $`L_J`$ (resp. $`L_J^{}`$) by its derived subgroup modulo its center and $`๐_J`$ (resp. $`๐_J^{}^{}`$) by its *reduction* (see , Section 3.3 and , Section 3.13), we may assume that $`L_J`$ (resp. $`L_J^{}`$) is an adjoint Chevalley group.
Since $`G`$ and $`G^{}`$ are isomorphic, it follows from Proposition 6.1 that $`๐`$ and $`๐^{}`$ are either both finite or both infinite. Moreover, if $`๐`$ and $`๐^{}`$ are finite, then they have the same characteristic in view of Proposition 6.2. Now, the desired result is a consequence of Theorem 31 in if $`๐`$ and $`๐^{}`$ are finite and from Theorem 8.16 of otherwise. โ
The preceding theorem can be used to decompose automorphisms of a given Kac-Moody group in a product of automorphisms of five specific kinds, as mentioned in the introduction. For the precise definitions of these specific automorphisms and further comments on them, we refer the reader to Section 9 of .
###### Corollary 6.4.
Let $`G`$ be a Kac-Moody group over a field $`๐`$ of at least 4 elements and associated with a generalized Cartan matrix $`A`$ of indecomposable type. Then, any automorphism of $`G`$ which preserve bounded subgroups can be written as a product of an inner, a sign, a diagonal, a graph and a field automorphism. Furthermore, if $`G`$ is โsimply connected in the weak senseโ (see , Remark 3.7(c) p. 550) and if moreover, either $`\mathrm{char}(๐)=0`$ or every off-diagonal entry of the generalized Cartan matrix $`A`$ is prime to $`\mathrm{char}(๐)`$, then the term โgraph automorphismโ may be replaced by โdiagram automorphismโ in the preceding statement.
The proof goes along the lines of the proof of Theorem 2.7 of and is omitted here.
## 7 Kac-Moody groups over finite fields
Corollary B of the introduction is a consequence of the following two results.
###### Theorem 7.1.
Let $`๐`$ be the collection of all twin root data arising from Kac-Moody groups over finite fields of cardinality at least 4. Then any element of $`๐`$ is $`๐`$-rigid.
###### Proof.
Given any two elements $`๐`$ and $`๐^{}`$ of $`๐`$, it is clear from Corollary 3.8 that every isomorphism of $`G^๐`$ to $`G^๐^{}`$ preserves bounded subgroups. The result is thus a consequence of Theorem 6.3. โ
###### Corollary 7.2.
Let $`G`$ be a Kac-Moody group over a finite field $`๐`$ of at least 4 elements and associated with a generalized Cartan matrix $`A`$ of indecomposable type. Then, any automorphism of $`G`$ can be written as a product of an inner, a sign, a diagonal, a graph and a field automorphism. Furthermore, if $`G`$ is โsimply connected in the weak senseโ (see , Remark 3.7(c) p. 550) and if moreover, every off-diagonal entry of the generalized Cartan matrix $`A`$ is prime to $`\mathrm{char}(๐)`$, then the term โgraph automorphismโ may be replaced by โdiagram automorphismโ in the preceding statement.
As for Corollary 6.4 above, the proof is omitted.
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# 1 Introduction
## 1 Introduction
In this paper we investigate the large level asymptotics of the ReshetikhinโTuraev invariants of the $`3`$โmanifolds obtained by doing surgery on the figure $`8`$ knot with an arbitrary rational surgery coefficient. Let $`X`$ be a closed oriented $`3`$โmanifold and let $`\tau _r(X)`$ be the RTโinvariant of $`X`$ associated to $`๐ฐ๐ฉ_2()`$ at level $`r`$, some integer $`2`$. The investigations of this paper are motivated by the following conjecture.
###### Conjecture 1 (Asymptotic expansion conjecture (AEC))
There exist constants (depending on X) $`d_j`$ and $`b_j`$ for $`j=0,1,\mathrm{},n`$ and $`a_j^l`$ for $`j=0,1,\mathrm{},n`$, $`l=1,2,\mathrm{}`$ such that the asymptotic expansion of $`\tau _r(X)`$ in the limit $`r\mathrm{}`$ is given by
$$\tau _r(X)\underset{j=0}{\overset{n}{}}e^{2\pi irq_j}r^{d_j}b_j\left(1+\underset{l=1}{\overset{\mathrm{}}{}}a_j^lr^l\right),$$
where $`q_0=0,q_1,\mathrm{},q_n`$ are the finitely many different values of the ChernโSimons functional on the space of flat $`\text{SU}(2)`$โconnections on $`X`$.
Here $``$ means asymptotic expansion in the Poincarรฉ sense, which means the following: Let
$$d=\mathrm{max}\{d_0,\mathrm{},d_n\}.$$
Then, for any non-negative integer $`L`$, there is a $`c_L`$ such that
$$\left|\tau _r(X)\underset{j=0}{\overset{n}{}}e^{2\pi irq_j}r^{d_j}b_j\left(1+\underset{l=0}{\overset{L}{}}a_j^lr^l\right)\right|c_Lr^{dL1}$$
for all levels $`r`$. Of course such a condition only puts limits on the large $`r`$ behaviour of $`\tau _r(X)`$.
Given an arbitrary complex function $`f(r)`$ defined on the positive integers a little simple argument shows that there at most exists one list of numbers $`n\{0,1,\mathrm{}\}`$, $`q_0,q_1,\mathrm{},q_n[0,1[`$, $`d_j`$ and $`b_j`$ for $`j=0,1,\mathrm{},n`$ and $`a_j^l`$ for $`j=0,1,\mathrm{},n`$, $`l=1,2,\mathrm{}`$ such that $`q_0<q_1<\mathrm{}<q_n`$, $`b_j0`$ for $`j>0`$ and such that the large $`r`$ asymptotic expansion of $`f(r)`$ is given by
$$f(r)\underset{j=0}{\overset{n}{}}e^{2\pi irq_j}r^{d_j}b_j\left(1+\underset{l=1}{\overset{\mathrm{}}{}}a_j^lr^l\right).$$
This implies that if the function $`r\tau _r(X)`$ has an asymptotic expansion of the form stated in Conjecture 1 then the $`q_j`$โs, $`d_j`$โs, $`b_j`$โs and the $`a_j^l`$โs are all uniquely determined by the set of invariants $`\{\tau _r(X)\}_{r2}`$, hence they are also topological invariants of $`X`$. As stated above the AEC already includes the claim that the $`q_j`$โs are the ChernโSimons invariants. There are also conjectured topological formulae for the $`d_j`$โs, and $`b_j`$โs (see e.g. and the references given there).
In general, there should be expressions for each of the $`a_j^l`$โs in terms of sums over Feynman diagrams of certain contributions determined by the Feynman rules of the ChernโSimons theory. This has not yet been worked out in general, except in the case of an acyclic flat connection and the case of a smooth non-degenerate component of the moduli space of flat connections by Axelrod and Singer, cf. , , .
The AEC, Conjecture 1, however offers in a sense a converse point of view, where one seeks to derive the final output of perturbation theory after all cancellations have been made (i.e. collect all terms with the same ChernโSimons value). This seems actually rather reasonable in this case, since the exact invariant is known explicitly.
The AEC was proved by Andersen in in the case of mapping tori of finite order diffeomorphisms of orientable surfaces of genus at least two using the gauge-theoretic approach to the quantum invariants. Later on the AEC was proved by Hansen in for all Seifert manifolds with orientable base by supplementing the work of Rozansky with the required analytic estimates. In the AEC is futher proved for the Seifert manifolds with nonorientable base of even genus.
Using the approach of Reshetikhin and Turaev to the quantum invariants, the AEC has not yet been proved for any hyperbolic $`3`$โmanifold. It is therefore particular interesting to consider surgeries on the figure $`8`$ knot. Let $`M_{p/q}`$ be the manifold obtained by (rational) Dehn surgery on the figure $`8`$ knot with surgery coefficient $`p/q`$. Then $`M_{p/q}`$ has a hyperbolic structure if and only if $`|p|>4`$ or $`|q|>1`$, see e.g. \[25, Theorem 10.5.10\] or . (It is well-known that for $`|q|=1`$ and $`|p|\{1,2,3\}`$ $`M_{p/q}`$ is a Seifert manifold, and $`M_{\pm 4}`$ are obtained by gluing together the complement of the trefoil knot and the non-trivial $`I`$โbundle over the Klein bottle, see e.g. \[15, p. 95\]. The manifold $`M_0`$ is a mapping torus of a torus, see Appendix C.) We use here the convention of Rolfsen for surgery coefficients, cf. \[30, Chap. 9\]. In particular Dehn surgery on a knot $`K`$ in $`S^3`$ with surgery coefficient $`f`$ is equal to the boundary of the compact $`4`$โmanifold obtained by attaching a $`2`$โhandle to the $`4`$โball using the knot $`K`$ with framing $`f`$, see \[30, p. 261\]. As usual $`M_{p/q}`$ is given the orientation induced by the standard right-handed orientation of $`S^3`$.
The advantage of working with surgeries on the figure $`8`$ knot $`K`$ is that the (normalized) colored Jones polynomial $`J_K^{}(\lambda )`$ is known explicitly. In fact
$$J_K^{}(\lambda )=\underset{m=0}{\overset{\lambda 1}{}}\xi ^{m\lambda }\underset{l=1}{\overset{m}{}}(1\xi ^{\lambda l})(1\xi ^{\lambda +l}),$$
where $`\xi =\mathrm{exp}(2\pi i/r)`$ (and the product is $`1`$ for $`m=0`$). The colors $`\lambda `$ are here dimensions of irreducible representations of the quantum group associated to $`๐ฐ๐ฉ_2()`$ and the root of unity $`\xi `$, so $`\lambda =1,2,\mathrm{},r`$. By the above expression for $`J_K^{}(\lambda )`$ we have an explicit formula for the quantum invariant $`\tau _r(M_{p/q})`$ (see formula (0.17)). Although this formula is completely explicit, it is not clear from it what the leading order large $`r`$ asymptotics of $`\tau _r(M_{p/q})`$ is. In order to study this asymptotics, we observe (generalizing from Kashaevโs work) that the product in the expression for the colored Jones polynomial can be expressed in terms of a quotient of two evaluations of Faddeevโs quantum dilogarithm $`S_\gamma `$ ($`\gamma =\pi /r`$):
$$J_K^{}(\lambda )=\underset{m=0}{\overset{\lambda 1}{}}\frac{\xi ^{m\lambda }}{(1\xi ^\lambda )}\frac{S_\gamma (\pi +2\gamma (\lambda m)\gamma )}{S_\gamma (\pi +2\gamma (\lambda +m)+\gamma )}.$$
(0.1)
This follows directly from the functional equation
$$(1+e^{\sqrt{1}\zeta })S_\gamma (\zeta +\gamma )=S_\gamma (\zeta \gamma )$$
which Faddeevโs $`S_\gamma `$ satisfies. Recall that for $`\text{Re}(\zeta )<\pi +\gamma `$ we have the expression
$$S_\gamma (\zeta )=\mathrm{exp}\left(\frac{1}{4}_{C_R}\frac{e^{\zeta z}}{\mathrm{sinh}(\pi z)\mathrm{sinh}(\gamma z)z}\text{d}z\right),$$
which together with the functional equation determines $`S_\gamma `$ as a meromorphic function on $``$. For the so-called top color, i.e. $`\lambda =r`$, we obtain the sligthly simpler expression
$$J_K^{}(r)=r\underset{m=0}{\overset{r1}{}}\frac{S_\gamma (\pi (2m+1)\gamma )}{S_\gamma (\pi +(2m+1)\gamma )}.$$
Then we simply use the residue formula to convert this expression into a contour integral
$$J_K^{}(r)=_{C_r}\mathrm{tan}(\pi rx)\stackrel{~}{g}_r(x)\text{d}x,$$
where $`C_r`$ is contained in the strip $`\{z0<\text{Re}(z)<1\}`$ and encloses $`(m+\frac{1}{2})/r`$ for $`m=0,\mathrm{},r1`$, and $`\stackrel{~}{g}_r`$ is a holomorphic function in this strip expressed in terms of the above quotient of $`S_\gamma `$ functions (see formula (0.21)).
Similarly we get for the quantum invariant with the use of (0.1) and the residue theorem, now a double contour integral, since the quantum invariant also involves a sum over colors:
$$\tau _r(M_{p/q})=_{C_r\times C_r}\mathrm{cot}(\pi rx)\mathrm{tan}(\pi ry)\stackrel{~}{f}_{p,q,r}(x,y)\text{d}x\text{d}y,$$
(0.2)
where $`\stackrel{~}{f}_{p,q,r}`$ is holomorphic on the double strip $`\{z0<\text{Re}(z)<1\}^2`$ and given by some expression involving quotients of evaluations of $`S_\gamma `$ functions, and where we furthermore require of $`C_r`$ that it also encloses $`k/r`$ for $`k=1,\mathrm{},r1`$.
From this it is clear that we need to understand the small $`\gamma `$ asymptotics of $`S_\gamma `$. We have that
$$S_\gamma (\zeta )=\mathrm{exp}\left(\frac{1}{2\sqrt{1}\gamma }\text{Li}_2(e^{\sqrt{1}\zeta })+I_\gamma (\zeta )\right),$$
where $`\text{Li}_2`$ is Eulerโs dilogarithm function, and where we have certain analytic estimates on $`I_\gamma (\zeta )`$ (see Lemma 3).
Let us first explain how we use this to give a proof of the volume conjecture of Murakami and Murakami for the figure $`8`$ knot, namely that
$$\underset{r\mathrm{}}{lim}\frac{2\pi \text{Log}(J_K^{}(r))}{r}=\text{Vol}(K),$$
where the right-hand side is the hyperbolic volume of the figure $`8`$ knot, i.e. the hyperbolic volume of the complement $`S^3K`$. The basic idea in analyzing the above contour integral expression for $`J_K^{}(r)`$ is the following. In the upper half plane we approximate $`\mathrm{tan}`$ by $`\sqrt{1}`$ and by $`\sqrt{1}`$ in the lower half plane. Further we approximate $`S_\gamma `$ by the above expression involving only the dilogarithm. In Appendix B we prove the required estimates which allows us to do these approximations and we end up with the following formula for the leading order large $`r`$ asymptotics
$$J_K^{}(r)r^2_ฯต^{1ฯต}e^{r\mathrm{\Phi }(x)}\text{d}x,$$
(0.3)
where $`ฯต<1/(4r)`$ is a small positive parameter and
$$\mathrm{\Phi }(x)=\frac{1}{2\pi \sqrt{1}}\left(\text{Li}_2(e^{2\pi \sqrt{1}x})\text{Li}_2(e^{2\pi \sqrt{1}x})\right).$$
Now we simply analyze the integral on the right-hand side of (0.3) by the saddle point method. This consists of finding the stationary points of $`\mathrm{\Phi }`$ and also the directions of steepest descent, see e.g. . (In this paper we use critical point and stationary point interchangeably to mean a point in which the derivative is zero.) This analysis leads to two interesting results. Firstly, the search for stationary points leads to the hyperbolicity equation for the complement of the figure $`8`$ knot in the $`3`$โsphere. Recall that this complement can be decomposed into two ideal hyperbolic tetrahedra each parametrized by a certain complex number. This decomposition then defines a hyperbolic structure on the complement exactly when the two parameters are equal and satisfy the hyperbolicity equation.
Secondly, we find that the value of the phase function in the relevant stationary point (there is only one such point in this case) is equal to the hyperbolic volume of the knot complement (divided by $`2\pi `$), hence the leading asymptotics of $`J_K^{}(r)`$ is determined by this volume.
These phenomena were first observed by Kashaev and have been conjectured by Thurston and Yokota to be generally true for hyperbolic knots (see Remark 2).
Ultimately our asymptotic analysis leads to the following
###### Theorem 1
The leading order large $`r`$ asymptotics of the colored Jones polynomial evaluated at the top color is given by
$$J_K^{}(r)3^{1/4}r^{3/2}\mathrm{exp}\left(\frac{r}{2\pi }\text{Vol}(K)\right).$$
As a corollary we obtain the volume conjecture for the figure $`8`$ knot. We note that none of the proofs so far given in the literature for the volume conjecture for the figure $`8`$ knot have been able to see the finer details of the asymptotic behaviour, namely the polynomial part $`3^{1/4}r^{3/2}`$.
After completion of this paper we have been informed that D. Zagier has computed the full asymptotic expansion of $`J_K^{}(r)`$ by using the EulerโMaclaurin summation formula, however his techniques seem not applicable to the calculation of the large $`r`$ asymptotics of $`\tau _r(M_{p/q})`$.
Let us now return to the study of the large $`r`$ asymptotics of the quantum invariant $`\tau _r(M_{p/q})`$. We expect that an analysis of the expression (0.2) paralleling our analysis of $`J_K^{}(r)`$ should be applicable. I.e. $`\mathrm{tan}`$ and $`\mathrm{cot}`$ should be approximated by $`\pm \sqrt{1}`$ depending on the signs of $`\text{Im}(y)`$ and $`\text{Im}(x)`$ and $`\stackrel{~}{f}_{p,q,r}(x,y)`$ by an appropriate expression involving the dilogarithm for some deformation of the relevant part of $`C_r\times C_r`$. We have partial analytic results supporting this.
To be more specific we propose the following analog of (0.3) for the quantum invariant. Let $`d`$ be the inverse of $`pmodq`$ and let $`(a,b)\{0,1\}^2`$ and $`n`$. Define
$$\mathrm{\Phi }_n(x,y)=\frac{dn^2}{q}\frac{p}{4q}x^2+\frac{n}{q}xxy+\frac{1}{4\pi ^2}\left(\text{Li}_2(e^{2\pi i(x+y)})\text{Li}_2(e^{2\pi i(xy)})\right),$$
and
$$\mathrm{\Phi }_n^{a,b}(x,y)=a(x+y)+b(xy)+\mathrm{\Phi }_n(x,y).$$
(0.4)
Finally, let
$$๐ฎ=\{(x,y)\times |e^{2\pi iy}]\mathrm{},0[\}.$$
(0.5)
Thus $`๐ฎ`$ is the union of infinitely many planes, namely the ones characterized by $`x`$ and $`\text{Re}(y)\frac{1}{2}+`$.
###### Conjecture 2
There exist surfaces $`\stackrel{~}{\mathrm{\Sigma }}_{a,b}^{\mu ,\nu ,n}^2`$ for $`(a,b)\{0,1\}^2`$, $`n/|q|`$ and $`(\mu ,\nu )\{\pm 1\}^2`$ such that the leading order large $`r`$ asymptotics of the quantum invariant is given by
$`\overline{\tau }_r(M_{p/q})`$ $``$ $`Cr{\displaystyle \underset{n/|q|}{}}{\displaystyle \underset{(a,b)\{0,1\}^2}{}}{\displaystyle \underset{(\mu ,\nu )\{\pm 1\}^2}{}}\mu \nu `$ (0.7)
$`\times {\displaystyle _{\stackrel{~}{\mathrm{\Sigma }}_{a,b}^{\mu ,\nu ,n}}}\stackrel{~}{g}_n(x)e^{2\pi ir\mathrm{\Phi }_n^{a,b}(x,y)}\text{d}x\text{d}y,`$
where $`C`$ is a constant only depending on $`p`$ and $`q`$ and $`\stackrel{~}{g}_n`$ is some simple $`r`$โindependent function of $`x`$. If $`p/q0`$ the surfaces $`\stackrel{~}{\mathrm{\Sigma }}_{a,b}^{\mu ,\nu ,n}^2`$ can be chosen such that any critical point of $`\mathrm{\Phi }_n^{a,b}`$ belonging to $`\stackrel{~}{\mathrm{\Sigma }}_{a,b}^{\mu ,\nu ,n}`$ also belongs to the set $`๐ฎ`$.
Please see Conjecture 3 for the more detailed version of this conjecture, including the precise formula for $`\stackrel{~}{g}_n`$. (We have here for sign-reasons switched to the complex conjugate invariant $`\overline{\tau }_r(M_{p/q})=\overline{\tau _r(M_{p/q})}=\tau _r(M_{p/q})`$.) In case $`p/q=0`$ our results show that we have to include critical points not belonging to $`๐ฎ`$ in our surfaces $`\stackrel{~}{\mathrm{\Sigma }}_{a,b}^{\mu ,\nu ,n}`$, see Appendix C.
We now proceed by making an asymptotic analysis of the right-hand side of (0.7) using the saddle point method like in our analysis of (0.3). Thus we need to analyze integrals of the form
$`I_{a,b}^{\mu ,\nu ,n}={\displaystyle _{\stackrel{~}{\mathrm{\Sigma }}_{a,b}^{\mu ,\nu ,n}}}\stackrel{~}{g}_n(x)e^{2\pi ir\mathrm{\Phi }_n^{a,b}(x,y)}\text{d}x\text{d}y.`$ (0.8)
Again we have to determine the stationary points of $`\mathrm{\Phi }_n^{a,b}`$ and the values of $`\mathrm{\Phi }_n^{a,b}`$ in the relevant stationary points. The main idea behind the saddle point method is to deform $`\stackrel{~}{\mathrm{\Sigma }}_{a,b}^{\mu ,\nu ,n}`$ so that it contains certain stationary points of $`\mathrm{\Phi }_n^{a,b}`$ satisfying that the leading order large $`r`$ asymptotics of $`I_{a,b}^{\mu ,\nu ,n}`$ is given by integrating the integrant of the integral (0.8) over small neighborhoods of these stationary points.
If we let $`v=e^{\pi ix}`$ and $`w=e^{2\pi iy}`$, then by exponentiating the two equations for $`(x,y)`$ being a stationary point of $`\mathrm{\Phi }_n^{a,b}`$ (see Theorem 2 below) we obtain the equations
$`v^p`$ $`=`$ $`\left({\displaystyle \frac{wv^2}{1v^2w}}\right)^q,`$
$`v^2w`$ $`=`$ $`(1v^2w)(wv^2),`$ (0.9)
which are independent of the integer parameters $`a,b,n`$. To link the asymptotics to the flat connections (as proposed by the AEC) we then have to relate the relevant stationary points of the phase functions $`\mathrm{\Phi }_n^{a,b}`$ to the classical $`\text{SU}(2)`$ ChernโSimons theory on the manifolds $`M_{p/q}`$. Fortunately, this ChernโSimons theory has been given a detailed description by Kirk and Klassen using the work of Riley , on the $`\text{SL}(2,)`$ representation variety of the knot group of the figure $`8`$ knot. According to Riley the conjugacy classes of the nonabelian elements of this variety can be represented by a certain set of representations $`\rho `$ parametrized by $`\rho =\rho _{(s,u)}`$, where $`(s,u)^{}\times `$ satisfies a certain polynomial equation. (To be precise $`\rho _{(s_1,u_1)}`$ and $`\rho _{(s_2,u_2)}`$ are conjugate if and only if $`u_2=u_10`$ and $`s_2\{s_1,s_1^1\}`$ or $`u_2=u_1=0`$ and $`s_2=s_1A`$, where $`A`$ is a certain subset of $`^{}`$ consisting of $`4`$ points.)
We show that $`\rho _{(s,u)}`$ defines a $`\text{SL}(2,)`$โrepresentation of $`\pi _1(M_{p/q})`$ if and only if $`(v,w)=(s,u+1)`$ is a solution to (1) and $`v^21`$. By a result of Riley it is known that this representation is conjugate to a $`\text{SU}(2)`$โrepresentation if and only if $`(s,u)S^1\times `$. We relate the results of Kirk and Klassen on the ChernโSimons invariants of flat $`\text{SU}(2)`$โconnections on the manifolds $`M_{p/q}`$ to the asymptotic analysis of the quantum invariants $`\overline{\tau }_r(M_{p/q})`$ via a detailed analysis of the relevant critical values of the phase functions $`\mathrm{\Phi }_n^{a,b}`$. Ultimately we arrive at
###### Theorem 2
The map
$$(x,y)[\rho _{(e^{\pi ix},e^{2\pi iy}1)}]$$
gives a surjection from the set of critical points $`(x,y)`$ of the functions $`\mathrm{\Phi }_n^{a,b}`$ with $`x`$ onto the set of conjugacy classes of nonabelian $`\text{SL}(2,)`$โrepresentations of $`\pi _1(M_{p/q})`$. Moreover, $`(x,y)^2`$ is a critical point of $`\mathrm{\Phi }_n^{a,b}`$ if and only if
$`2a+{\displaystyle \frac{n}{q}}`$ $`=`$ $`y+\left({\displaystyle \frac{p}{2q}}+1\right)x+{\displaystyle \frac{i}{\pi }}\text{Log}\left(1e^{2\pi i(x+y)}\right),`$
$`2b+{\displaystyle \frac{n}{q}}`$ $`=`$ $`y+\left({\displaystyle \frac{p}{2q}}1\right)x{\displaystyle \frac{i}{\pi }}\text{Log}\left(1e^{2\pi i(xy)}\right).`$ (0.10)
Futhermore, if $`(x,y)`$ is a critical point of $`\mathrm{\Phi }_n^{a,b}`$ then $`\rho _{(e^{\pi ix},e^{2\pi iy}1)}`$ is equivalent to a $`\text{SU}(2)`$โrepresentation $`\overline{\rho }`$ of $`\pi _1(M_{p/q})`$ if and only if $`(x,y)๐ฎ`$ (see (0.5)), and in that case
$$\text{CS}(\overline{\rho })=\mathrm{\Phi }_n^{a,b}(x,y)(mod),$$
where CS is the ChernโSimons functional on the space of flat $`\text{SU}(2)`$โconnections on $`M_{p/q}`$.
Continuing our asymptotic analysis of the integral (0.8) we point out that not all critical points have to give contributions to the asymptotics. The simplest case arise when all the relevant critical points are non-degenerate. Following Conjecture 2 we claim that the only critical points giving a contribution to the leading order large $`r`$ asymptotics of $`\tau _r(M_{p/q})`$ are (some of) the ones belonging to the set $`๐ฎ`$ in (0.5). In Proposition 1 in Sec. 4 we prove that all these critical points are non-degenerate in case $`|p/q|<\sqrt{20}`$. In case $`|p/q|>\sqrt{20}`$ we have the following partial result: The surjection described in the first part of Theorem 2 restricts to a surjection $`\varphi `$ from the set of critical points belonging to $`๐ฎ`$ onto the moduli space $`_{p/q}^{}`$ of flat irreducible $`\text{SU}(2)`$โconnections on $`M_{p/q}`$. If $`p/q>\sqrt{20}`$ we prove that all the critical points belonging to $`๐ฎ๐_0`$ are non-degenerate, where $`๐_0`$ is a certain set of critical points. To be more precise the set $`๐_0`$ is either empty or it is equal to the preimage of $`\varphi `$ of one point if $`p`$ is odd and of two points if $`p`$ is even. We expect that the argument equation following from the first of the equations (1) together with the last statement in Proposition 1 should establish that $`๐_0`$ is empty.
From Conjecture 2 we see that the only relevant phase functions $`\mathrm{\Phi }_n^{a,b}`$ are the ones with $`a,b\{0,1\}`$. This, however, does not imply that only a proper subset of the conjugacy classes of $`\text{SL}(2,)`$โrepresentations of $`\pi _1(M_{p/q})`$ are in play. In fact, by using (2), we see that if $`(x,y)`$ is a critical point of $`\mathrm{\Phi }_n^{a,b}`$ then $`(x+2k,y+2l)`$ is a critical point of $`\mathrm{\Phi }_{n+pk}^{a+l+k,b+lk}`$ for any $`k,l`$. Moreover, the critical points $`(x,y)`$ and $`(x+2k,y+2l)`$ correspond to the same conjugacy class of representations by the surjection in Theorem 2. In particular all points of the moduli space of flat irreducible $`\text{SU}(2)`$โconnections on $`M_{p/q}`$ could potentially give a contribution to the asymptotics of $`\overline{\tau }_r(M_{p/q})`$.
One is only likely to succeed in using the saddle point method to calculate the large $`r`$ asymptotics of an integral of the form (0.8) in case one can deform the surface $`\stackrel{~}{\mathrm{\Sigma }}_{a,b}^{\mu ,\nu ,n}`$ so that it only contains critical point of a certain nice kind. Here we will only consider the case of non-degenerate critical points. Therefore, assume that $`(x,y)`$ is a non-degenerate critical point of $`\mathrm{\Phi }_n^{a,b}`$ belonging to $`\stackrel{~}{\mathrm{\Sigma }}_{a,b}^{\mu ,\nu ,n}`$. Then to be able to calculate that critical pointโs contribution to the large $`r`$ asymptotics of the integral in (0.8) we assume that $`(x,y)`$ is positive definite, meaning that the imaginary part of a certain โtwistedโ Hessian of $`\mathrm{\Phi }_n^{a,b}`$ in $`(x,y)`$ is positive definite, see Definition 1 for the precise definition. Ultimatively we arive at the following result.
###### Theorem 3
Assume that $`p/q0`$. IF Conjecture 2 is true and if all the critical points of the phase functions $`\mathrm{\Phi }_n^{a,b}`$ belonging to $`๐ฎ`$ are non-degenerate and positive definite, then the leading order large $`r`$ asymptotics of the quantum invariant $`\overline{\tau }_r(M_{p/q})`$ is given by
$`\overline{\tau }_r(M_{p/q})`$ $``$ $`{\displaystyle \frac{\text{sign}(q)}{4\sqrt{|q|}}}e^{\frac{3\pi i}{4}\text{sign}(pq)}{\displaystyle \underset{\overline{\rho }_{p/q}^{}}{}}e^{2\pi irCS(\overline{\rho })}b_{\overline{\rho }},`$ (0.11)
where $`_{p/q}^{}`$ is the moduli space of flat irreducible $`\text{SU}(2)`$โconnections on $`M_{p/q}`$. For each $`\overline{\rho }_{p/q}^{}`$ there are $`(a,b)\{0,1\}^2`$, $`n`$ and a stationary point $`(x,y)๐ฎ`$ for the function $`\mathrm{\Phi }_n^{a,b}`$ given by (0.4), i.e. $`(x,y)`$ solves the equations (2), such that $`\overline{\rho }`$ is equivalent to $`\rho _{(e^{\pi ix},e^{2\pi iy}1)}`$. Moreover
$$CS(\overline{\rho })=\mathrm{\Phi }_n^{a,b}(x,y)(mod),$$
and
$$b_{\overline{\rho }}=m_{\overline{\rho }}e^{\frac{\pi i\sigma _{\overline{\rho }}}{2}}\mathrm{sin}\left(\frac{\pi }{q}(x2nd)\right)|14\mathrm{cos}(2\pi x))+\frac{p}{2q}\mathrm{sinh}(2\pi \text{Im}(y))|^{\frac{1}{2}},$$
where $`m_{\overline{\rho }}`$ and $`\sigma _{\overline{\rho }}`$ are some integers.
The positive integers $`m_{\overline{\rho }}`$ arrise from the following two situations. First of all it can happen that two or more critical points of the relevant stationary phase functions correspond to the same irreducible flat $`\text{SU}(2)`$โconnection on $`M_{p/q}`$. Secondly, the same stationary point for the same phase function $`\mathrm{\Phi }_n^{a,b}`$ can belong to two or more of the surfaces $`\stackrel{~}{\mathrm{\Sigma }}_{a,b}^{\mu ,\nu ,n}`$. In fact, we expect that there is a one to one correspondence between the moduli space $`_{p/q}^{}`$ of flat irreducible $`\text{SU}(2)`$โconnections on $`M_{p/q}`$ and the set of stationary points contributing to the leading large $`r`$ asymptotics of $`\overline{\tau }_r(M_{p/q})`$. Moreover, we expect that a contributing stationary point for a relevant phase function $`\mathrm{\Phi }_n^{a,b}`$ belongs to (a deformation) of $`\mathrm{\Sigma }_{a,b}^{\mu ,\nu ,n}`$ for all $`(\mu ,\nu )\{\pm 1\}^2`$ thus causing $`m_{\overline{\rho }}=4`$ for all $`\overline{\rho }_{p/q}^{}`$, see Remark 4. The above result, Theorem 3, coincides with the one found by the second author in for the cases $`|p/q|\{1,2,3\}`$ with $`m_{\overline{\rho }}=4`$ for all $`\overline{\rho }`$.
The invariant $`\tau _r(M_0)`$ and its full asymptotic expansion have been calculated by Jeffrey . We show (see Appendix C) that Jeffreyโs result is in agreement with the AEC. In this case the reducible flat $`\text{SU}(2)`$โconnections contribute to the leading asymptotics. We find that the part of the leading order large $`r`$ asymptotics of $`\tau _r(M_0)`$ associated to the irreducible flat $`\text{SU}(2)`$โconnections on $`M_0`$ is given by the right hand side of (0.11) with all $`m_{\overline{\rho }}=4`$.
Acknowledgements
The first author thanks the Department of Mathematics and the Mathematical Sciences Research Institute, University of California, Berkeley for their hospitality during several visits, where part of this work was undertaken. The second author thanks the Universitรฉ Louis Pasteur, Strasbourg, the University of Edinburgh, and the MaxโPlanckโInstitut fรผr Mathematik, Bonn for their hospitality during this work. He was supported by the European Commission, the Danish Natural Science Research Council and the MaxโPlanckโInstitut fรผr Mathematik, Bonn.
## 2 The RT-invariant for surgeries on the figure $`8`$ knot
This section is primarily intended to introduce notation. Moreover, we present some preliminary formulas for the colored Jones polynomial of the figure $`8`$ knot and for the RTโinvariants of the $`3`$โmanifolds $`M_{p/q}`$.
Let $`t=\mathrm{exp}(2\pi \sqrt{1}/(4r))`$, $`r`$ an integer $`2`$, and let $`U_t`$ be the modular Hopf algebra considered in \[27, Sec. 8\], i.e. $`U_t`$ is a finite-dimensional factor of the quantum group $`U_\xi (๐ฐ๐ฉ_2())`$, $`\xi =t^4`$. (In , and in most literature on the subject, $`\xi `$ is denoted $`q`$, but we use in this paper $`q`$ to mean something different.) For an integer $`k`$ we let
$$[k]=\frac{t^{2k}t^{2k}}{t^2t^2}=\frac{\mathrm{sin}(\pi k/r)}{\mathrm{sin}(\pi /r)},$$
sometimes called a quantum integer. For a knot $`K`$ in $`S^3`$ we denote by $`K^0`$ the knot $`K`$ considered as a framed knot with framing zero. The colored Jones polynomial associated to $`U_t`$ of a framed oriented knot $`K`$ with color $`\lambda \{1,2,\mathrm{},r\}`$ is denoted $`J_K(\lambda )`$, and for an oriented knot $`K`$ in $`S^3`$ we let $`J_K^{}(\lambda )=J_{K^0}(\lambda )/[\lambda ]`$. Here the colors are the dimensions (as complex vector spaces) of irreducible $`U_t`$โmodules.
Let $`N_f`$ be the $`3`$โmanifold obtained by surgery on $`S^3`$ along $`K`$ with surgery coefficient $`f`$. By or the RTโinvariant (at level $`r2`$) of $`N_f`$ is
$$\tau _r(N_f)=\alpha \underset{k=1}{\overset{r1}{}}\xi ^{(k^21)f/4}[k]^2J_K^{}(k),$$
(0.12)
where $`K`$ is given an arbitrary orientation. Here $`\alpha =C^{\text{sign}(f)}๐^2`$, where
$`๐`$ $`=`$ $`\sqrt{{\displaystyle \frac{r}{2}}}{\displaystyle \frac{1}{\mathrm{sin}(\pi /r)}},`$
$`C`$ $`=`$ $`\mathrm{exp}\left({\displaystyle \frac{\sqrt{1}\pi }{4}}{\displaystyle \frac{3(2r)}{r}}\right).`$
We use here the normalization of . This is $`๐^1`$ times the normalization of and $`C^{b_1(N_f)}๐^1`$ times the normalization of , where $`b_1(N_f)`$ is the first Betti number of $`N_f`$, see \[8, Appendix A\]. (In the notation of , $`C=\mathrm{\Delta }๐^1`$.)
Let us next generalize to arbitrary rational surgery. Let $`p,q`$ be a pair of coprime integers with $`q0`$, and let $`N_{p/q}`$ be the $`3`$โmanifold obtained by surgery along $`K`$ with surgery coefficient $`p/q`$. Choose $`c,d`$ such that $`B=\left(\begin{array}{cc}p\hfill & c\hfill \\ q\hfill & d\hfill \end{array}\right)\text{SL}(2,)`$. Then (see e.g. \[8, Theorem 5.1 and the proofs of Corollary 8.3 and Theorem 8.4\]),
$$\tau _r(N_{p/q})=\left(e^{\frac{i\pi }{4}}\mathrm{exp}\left(\frac{i\pi }{2r}\right)\right)^{\mathrm{\Phi }(B)3\text{sign}(pq)}\sqrt{\frac{2}{r}}\mathrm{sin}\left(\frac{\pi }{r}\right)\underset{\lambda =1}{\overset{r1}{}}[\lambda ]J_K^{}(\lambda )\stackrel{~}{B}_{\lambda ,1},$$
where $`\mathrm{\Phi }`$ is the Rademacher Phi function, see , and $`\stackrel{~}{}`$ is the unitary representation of $`P\text{SL}(2,)`$ given by
$`\stackrel{~}{B}_{j,k}=\sqrt{1}{\displaystyle \frac{\text{sign}(q)}{\sqrt{2r|q|}}}e^{\frac{\sqrt{1}\pi }{4}\mathrm{\Phi }(B)}`$
$`\times {\displaystyle \underset{\mu =\pm 1}{}}{\displaystyle \underset{n/|q|}{}}\mu \mathrm{exp}\left({\displaystyle \frac{\sqrt{1}\pi }{2rq}}[pj^22\mu j(k+2rn\mu )+d(k+2rn\mu )^2]\right).`$
By evaluating the sum over $`\mu `$ we get
$`\tau _r(N_{p/q})=a(r){\displaystyle \underset{n/|q|}{}}\mathrm{exp}\left(2\pi ir{\displaystyle \frac{dn^2}{q}}\right)`$ (0.13)
$`\times {\displaystyle \underset{k=1}{\overset{r1}{}}}\mathrm{sin}\left({\displaystyle \frac{\pi }{q}}[2nd{\displaystyle \frac{k}{r}}]\right)\mathrm{exp}\left({\displaystyle \frac{\pi ir}{2q}}[p\left({\displaystyle \frac{k}{r}}\right)^24n{\displaystyle \frac{k}{r}}]\right)[k]J_K^{}(k),`$
where
$$a(r)=\frac{2\text{sign}(q)}{r\sqrt{|q|}}\mathrm{sin}\left(\frac{\pi }{r}\right)e^{\frac{3\pi i}{4}\text{sign}(pq)}\mathrm{exp}\left(\frac{\pi i}{2r}\left[3\text{sign}(pq)\frac{p}{q}+\text{S}\left(\frac{p}{q}\right)\right]\right).$$
Here S is the Dedekind symbol, see e.g. . We note that the quantum invariant $`\tau _r`$ is independent of the colored Jones polynomial $`J_K^{}(k)`$ for the top-color $`k=r`$.
In the remaining part of this paper $`K`$ will denote the figure $`8`$ knot unless explicitly stated otherwise. Recall that $`M_{p/q}`$ denotes the $`3`$โmanifold obtained by surgery on $`S^3`$ along $`K`$ with surgery coefficient $`p/q`$. By an $`R`$โmatrix calculation (see e.g. ) we find that
$$J_K^{}(\lambda )=\underset{m=0}{\overset{\lambda 1}{}}\xi ^{m\lambda }\underset{l=1}{\overset{m}{}}(1\xi ^{\lambda l})(1\xi ^{\lambda +l})$$
(0.14)
for $`\lambda =1,2,\mathrm{},r`$, where $`_{l=1}^m(1\xi ^{kl})(1\xi ^{k+l})=1`$ for $`m=0`$. Le and Habiro have obtained a similar formula, cf. .
###### Remark 1
Unitarity of the TQFT associated to $`U_t`$ implies that
$$\tau _r(M)=\overline{\tau _r(M)}$$
(0.15)
for any $`3`$โmanifold $`M`$, where $`\overline{}`$ means complex conjugation. This formula also follows directly from and the remarks concerning normalization following (0.12).
Since the figure $`8`$ knot is amphicheiral, $`M_{p/q}`$ and $`M_{p/q}`$ are orientation reversing homeomorphic. By (0.15) we therefore have
$$\overline{\tau _r(M_{p/q})}=\tau _r(M_{p/q}).$$
(0.16)
This formula also follows directly by (0.13) and the facts that $`J_K^{}(\lambda )`$ is real and $`S(p/q)=S(p/q)`$. That $`J_K^{}(\lambda )`$ is real follows by amphicheirality of $`K`$ but can also be seen directly from (0.14) by
$`J_K^{}(\lambda )`$ $`=`$ $`{\displaystyle \underset{m=0}{\overset{\lambda 1}{}}}{\displaystyle \underset{l=1}{\overset{m}{}}}\xi ^{\lambda /2}\xi ^{l/2}(1\xi ^{\lambda l})\xi ^{\lambda /2}\xi ^{l/2}(1\xi ^{\lambda +l})`$
$`=`$ $`{\displaystyle \underset{m=0}{\overset{\lambda 1}{}}}{\displaystyle \underset{l=1}{\overset{m}{}}}(\xi ^{(\lambda l)/2}\xi ^{(\lambda l)/2})(\xi ^{(\lambda +l)/2}\xi ^{(\lambda +l)/2})`$
$`=`$ $`{\displaystyle \underset{m=0}{\overset{\lambda 1}{}}}(4)^m{\displaystyle \underset{l=1}{\overset{m}{}}}\mathrm{sin}(\pi (\lambda l)/r)\mathrm{sin}(\pi (\lambda +l)/r).`$
By (0.13) and (0.14) we get the following preliminary formula for the RTโinvariants of the manifolds $`M_{p/q}`$:
$`\tau _r(M_{p/q})`$ $`=`$ $`{\displaystyle \frac{ia(r)}{2\mathrm{sin}(\pi /r)}}{\displaystyle \underset{n/|q|}{}}\mathrm{exp}\left(2\pi ir{\displaystyle \frac{dn^2}{q}}\right)`$ (0.17)
$`\times {\displaystyle \underset{k=1}{\overset{r1}{}}}\mathrm{exp}\left({\displaystyle \frac{\pi ir}{2q}}[p\left({\displaystyle \frac{k}{r}}\right)^24n{\displaystyle \frac{k}{r}}]\right)\mathrm{sin}\left({\displaystyle \frac{\pi }{q}}[2nd{\displaystyle \frac{k}{r}}]\right)`$
$`\times {\displaystyle \underset{m=0}{\overset{r1}{}}}{\displaystyle \frac{\xi ^{(m+1/2)k}}{(1\xi ^k)}}{\displaystyle \underset{l=0}{\overset{m}{}}}(1\xi ^{kl})(1\xi ^{k+l}).`$
## 3 A complex double contour integral formula for $`\tau _r(M_{p/q})`$
In this section we derive a complex double contour integral formula for the RTโinvariants $`\tau _r(M_{p/q})`$ by using methods similar to Kashaev . When we consider the expression for the summand in the multi sum (0.17), we see that the expression as it stands only makes sense for non-negative integers $`m`$. In order to make sense of this expression for arbitrary complex values for $`m`$, let us consider the quantum dilogarithm of Faddeev
$$S_\gamma (\zeta )=\mathrm{exp}\left(\frac{1}{4}_{C_R}\frac{e^{\zeta z}}{\mathrm{sinh}(\pi z)\mathrm{sinh}(\gamma z)z}\text{d}z\right)$$
defined on $`\mathrm{\Delta }_\gamma =\{\zeta ||\text{Re}(\zeta )|<\pi +\gamma \}`$, where $`\gamma ]0,1[`$ and $`C_R`$ is the contour $`]\mathrm{},R]+\mathrm{{\rm Y}}_R+[R,\mathrm{}[`$, where $`\mathrm{{\rm Y}}_R(t)=Re^{\sqrt{1}(\pi t)}`$, $`t[0,\pi ]`$ and $`R]0,1[`$.
The function $`S_\gamma :\mathrm{\Delta }_\gamma `$ is holomorphic and it satisfies the following well-known functional equation (see or ).
###### Lemma 1
For $`\zeta `$ with $`|\text{Re}(\zeta )|<\pi `$ we have
$$(1+e^{\sqrt{1}\zeta })S_\gamma (\zeta +\gamma )=S_\gamma (\zeta \gamma ).$$
For the sake of completeness we have given a proof in Appendix A. We use Lemma 1 to extend $`S_\gamma `$ to a meromorphic function on the complex plane $``$.
From now on we fix $`\gamma =\pi /r`$ (so $`r>3`$). By Lemma 1 we get that
$$S_\gamma (\zeta )=S_\gamma (\zeta +2\pi )\underset{j=0}{\overset{r1}{}}\left(1+e^{\sqrt{1}(\zeta +(2j+1)\pi /r)}\right).$$
If we write $`\zeta =\pi +2\pi x`$ we get that
$$\underset{j=0}{\overset{r1}{}}\left(1+e^{\sqrt{1}(\zeta +(2j+1)\pi /r)}\right)=\underset{j=0}{\overset{r1}{}}\left(1w^je^{2\pi \sqrt{1}(x+\frac{1}{2r})}\right),$$
where $`w=e^{2\pi \sqrt{1}/r}`$. Using $`1z^r=_{j=0}^{r1}(1w^jz)`$ we get that
$$S_\gamma (\pi +2\pi x)=\left(1+e^{2\pi \sqrt{1}xr}\right)S_\gamma (\pi +2\pi (x+1))$$
(0.18)
for $`x`$. Let
$$x_n=\frac{n}{r}+\frac{1}{2r},n.$$
Then $`xS_\gamma (\pi +2\pi x)`$ is analytic on $`\{x_n|n=r,r+1,\mathrm{}\}`$. If $`m`$ is a positive integer then $`\{x_n|n=mr,mr+1,\mathrm{},(m+1)r1\}`$ are poles of order $`m`$, while the points $`\{x_n|n=mr,mr+1,\mathrm{},mr+r1\}`$ are zeros of order $`m`$.
Let us use the function $`S_\gamma `$ to give another expression for $`\tau _r(M_{p/q})`$. By Lemma 1 we have that
$$\underset{l=0}{\overset{m}{}}(1\xi ^{k\pm l})=\underset{l=0}{\overset{m}{}}\frac{S_\gamma (\pi +2\gamma (k\pm l)\gamma )}{S_\gamma (\pi +2\gamma (k\pm l)+\gamma )}.$$
Therefore
$`{\displaystyle \underset{l=0}{\overset{m}{}}}(1\xi ^{kl})`$ $`=`$ $`{\displaystyle \frac{S_\gamma (\pi +2\gamma (km)\gamma )}{S_\gamma (\pi +2\gamma k+\gamma )}},`$
$`{\displaystyle \underset{l=0}{\overset{m}{}}}(1\xi ^{k+l})`$ $`=`$ $`{\displaystyle \frac{S_\gamma (\pi +2\gamma k\gamma )}{S_\gamma (\pi +2\gamma (k+m)+\gamma )}}.`$
So
$$\underset{l=0}{\overset{m}{}}(1\xi ^{kl})(1\xi ^{k+l})=(1\xi ^k)\frac{S_\gamma (\pi +2\gamma (km)\gamma )}{S_\gamma (\pi +2\gamma (k+m)+\gamma )},$$
and then by (0.17)
$$\tau _r(M_{p/q})=\beta (r)\underset{n/|q|}{}\underset{k=1}{\overset{r1}{}}\underset{m=0}{\overset{r1}{}}f_{n,r}(\frac{k}{r},\frac{m+1/2}{r}),$$
where
$$f_{n,r}(x,y)=\mathrm{sin}\left(\frac{\pi }{q}(x2nd)\right)e^{2\pi ir\left(\frac{dn^2}{q}+\frac{p}{4q}x^2\frac{n}{q}xxy\right)}\frac{S_\gamma (\pi +2\pi (xy))}{S_\gamma (\pi +2\pi (x+y))}$$
and
$`\beta (r)`$ $`=`$ $`{\displaystyle \frac{ia(r)}{2\mathrm{sin}(\pi /r)}}`$ (0.19)
$`=`$ $`{\displaystyle \frac{i\text{sign}(q)}{r\sqrt{|q|}}}e^{\frac{3\pi i}{4}\text{sign}(pq)}\mathrm{exp}\left({\displaystyle \frac{\pi i}{2r}}\left[3\text{sign}(pq){\displaystyle \frac{p}{q}}+\text{S}\left({\displaystyle \frac{p}{q}}\right)\right]\right).`$
Note that $`d`$ is equal to the inverse of $`pmodq`$ and that the functions $`f_{n,r}`$ are independent of the choice of this inverse. By the remarks following Lemma 1 the functions $`f_{n,r}`$ are holomorphic on $`\mathrm{\Omega }_r\times \mathrm{\Omega }_r`$, where
$$\mathrm{\Omega }_s=\{w\frac{1}{4s}<\text{Re}(w)<1+\frac{1}{4s}\}$$
(0.20)
for $`s]0,\mathrm{}]`$. By the residue theorem we therefore end up with
###### Lemma 2
The quantum invariants of $`M_{p/q}`$ are given by
$$\tau _r(M_{p/q})=\frac{\beta (r)r^2}{4}\underset{n/|q|}{}_{C_r^1}\mathrm{cot}(\pi rx)\left(_{C_r^2}\mathrm{tan}(\pi ry)f_{n,r}(x,y)\text{d}y\right)\text{d}x,$$
where $`\beta (r)`$ is given by (0.19) and
$$f_{n,r}(x,y)=\mathrm{sin}\left(\frac{\pi }{q}(x2nd)\right)e^{2\pi ir\left(\frac{dn^2}{q}+\frac{p}{4q}x^2\frac{n}{q}xxy\right)}\frac{S_\gamma (\pi +2\pi (xy))}{S_\gamma (\pi +2\pi (x+y))},$$
and where $`C_r^1`$ is a closed curve in $`\mathrm{\Omega }_r`$ such that the poles $`\{k/r|k=1,2,\mathrm{},r1\}`$ for $`x\mathrm{cot}(\pi rx)`$ lies inside $`C_r^1`$ and the poles $`0`$ and $`1`$ lies outside $`C_r^1`$, and $`C_r^2`$ is a closed curve in $`\mathrm{\Omega }_r`$ such that the poles $`\{(m+1/2)/r|m=0,1,\mathrm{},r1\}`$ for $`y\mathrm{tan}(\pi ry)`$ lies inside $`C_r^2`$. Both curves are oriented in the anti-clockwise direction.
Using the function $`S_\gamma `$ we can also express $`J_K^{}(r)`$ as a contour integral. By Lemma 1 we get that
$$J_K^{}(r)=r\underset{m=0}{\overset{r1}{}}\frac{S_\gamma (\pi (2m+1)\gamma )}{S_\gamma (\pi +(2m+1)\gamma )}.$$
We have here used that
$$\frac{S_\gamma (\pi +\gamma )}{S_\gamma (\pi \gamma )}=\underset{j=1}{\overset{r1}{}}\frac{S_\gamma (\pi (2j+1)\gamma )}{S_\gamma (\pi +(2j+1)\gamma )}=\underset{j=1}{\overset{r1}{}}\left(1e^{2\pi \sqrt{1}\frac{j}{r}}\right)=r.$$
If we put
$$g_r(x)=\frac{S_\gamma (\pi 2\pi x)}{S_\gamma (\pi +2\pi x)}$$
for $`x\mathrm{\Omega }_{\frac{1}{2}r}`$ (see (0.20)) we get
$$J_K^{}(r)=r\underset{m=0}{\overset{r1}{}}g_r\left(\frac{m+1/2}{r}\right),$$
and we can write this sum as the single contour integral
$$J_K^{}(r)=\frac{\sqrt{1}r^2}{2}_{C_r^2}\mathrm{tan}(\pi rx)g_r(x)\text{d}x,$$
(0.21)
where $`C_r^2`$ is given as in Lemma 2.
## 4 The large $`r`$ asymptotics of $`J_K^{}(r)`$ and $`\tau _r(M_{p/q})`$
In this section we investigate the large $`r`$ asymptotics of $`\tau _r(M_{p/q})`$ or more precisely the leading term of this asymptotics, using the saddle point method. We begin by calculating the large $`r`$ asymptotics of $`J_K^{}(r)`$ using the expression (0.21). This calculation will demonstrate the use of the saddle point method and will serve as a warm up for the more difficult considerations of the asymptotics of $`\tau _r(M_{p/q})`$ in the final part of this section.
### 4.1 Semiclassical asymptotics of the quantum dilogarithm
It is well-known that the semiclassical asymptotics, i.e. the small $`\gamma `$ asymptotics of the quantum dilogarithm $`S_\gamma `$, is given by Eulerโs dilogarithm
$$\text{Li}_2(z)=_0^z\frac{\text{Log}(1w)}{w}\text{d}w$$
(0.22)
for $`z]1,\mathrm{}[`$. Here and elsewhere Log denotes the principal logarithm. For $`\zeta `$ satisfying either $`\text{Re}(\zeta )=\pm \pi `$ and $`\text{Im}(\zeta )0`$ or $`|\text{Re}(\zeta )|<\pi `$ one can check (see Appendix A) that
$$\frac{1}{2\sqrt{1}\gamma }\text{Li}_2(e^{\sqrt{1}\zeta })=\frac{1}{4}_{C_R}\frac{e^{\zeta z}}{\mathrm{sinh}(\pi z)\gamma z^2}\text{d}z,$$
hence we have that
$$S_\gamma (\zeta )=\mathrm{exp}\left(\frac{1}{2\sqrt{1}\gamma }\text{Li}_2(e^{\sqrt{1}\zeta })+I_\gamma (\zeta )\right)$$
(0.23)
for such $`\zeta `$, where
$$I_\gamma (\zeta )=\frac{1}{4}_{C_R}\frac{e^{\zeta z}}{z\mathrm{sinh}(\pi z)}\left(\frac{1}{\mathrm{sinh}(\gamma z)}\frac{1}{\gamma z}\right)\text{d}z.$$
###### Lemma 3
If $`|\text{Re}(\zeta )|<\pi `$ then
$$|I_\gamma (\zeta )|A\left(\frac{1}{\pi \text{Re}(\zeta )}+\frac{1}{\pi +\text{Re}(\zeta )}\right)\gamma +B\left(1+e^{\text{Im}(\zeta )R}\right)\gamma ,$$
and for $`|\text{Re}(\zeta )|\pi `$ we have
$$|I_\gamma (\zeta )|2A+B\left(1+e^{\text{Im}(\zeta )R}\right)\gamma ,$$
where $`A`$ and $`B`$ are positive constants only depending on $`R`$.
A proof is given in Appendix A. On the unit circle the imaginary part of the dilogarithm is given by Clausenโs function $`\text{Cl}_2`$, i.e.
$$\text{Im}\left(\text{Li}_2(e^{i\theta })\right)=\text{Cl}_2(\theta )=\underset{n=1}{\overset{\mathrm{}}{}}\frac{\mathrm{sin}(n\theta )}{n^2}=_0^\theta \text{Log}\left|2\mathrm{sin}\left(\frac{t}{2}\right)\right|\text{d}t$$
(0.24)
for $`\theta `$. One sees that $`\text{Cl}_2`$ is increasing on $`[0,\pi /3][5\pi /3,2\pi ]`$ and decreasing on $`[\pi /3,5\pi /3]`$. In particular, $`\text{Cl}_2`$ attains its maximum value at $`\pi /3`$ and its minimum value at $`5\pi /3`$. Moreover
$$\text{Cl}_2\left(\frac{5\pi }{3}\right)=\text{Cl}_2\left(\frac{\pi }{3}\right)=2_0^{\pi /6}\text{Log}|2\mathrm{sin}(\varphi )|\text{d}\varphi =2๐\left(\frac{\pi }{6}\right)=\frac{1}{2}\text{Vol}(K),$$
(0.25)
where $`๐`$ is Lobachevskyโs function and $`\text{Vol}(K)`$ is the hyperbolic volume of the complement of the figure $`8`$ knot, see e.g. \[25, Sec. 10.4\].
### 4.2 The large $`r`$ asymptotics of $`J_K^{}(r)`$
We calculate the leading term of the large $`r`$ asymptotics of $`J_K^{}(r)`$, using the saddle point method like Kashaev . Our calculation supplements the calculation of Kashaev with the required analytic error estimates. Let
$`C_r^2=C(\epsilon )`$ $`=`$ $`[\sqrt{1}+\epsilon ,\sqrt{1}+\epsilon ]+[\sqrt{1}+\epsilon ,1\epsilon \sqrt{1}]`$
$`+[1\epsilon \sqrt{1},1\epsilon +\sqrt{1}]+[1\epsilon +\sqrt{1},\epsilon +\sqrt{1}],`$
where $`\epsilon ]0,\frac{1}{4r}[`$. We let $`C_+(\epsilon )`$ be the part of the contour $`C(\epsilon )`$ above the real axes and $`C_{}(\epsilon )`$ the part below the real axes. By (0.21) we have
$$J_K^{}(r)=\frac{\sqrt{1}r^2}{2}\left(J_+(r,\epsilon )+J_{}(r,\epsilon )\right),$$
where
$$J_\pm (r,\epsilon )=_{C_\pm (\epsilon )}\mathrm{tan}(\pi rx)g_r(x)\text{d}x.$$
Away from the real axis, the factor $`\mathrm{tan}(\pi rx)`$ can be approximated by $`\pm \sqrt{1}`$ depending on whether we are in the upper or lower half-plane. In fact we have
$$|\mathrm{tan}(\pi rx)\sqrt{1}|\{\begin{array}{cc}4e^{2\pi r\text{Im}(x)},\hfill & \text{Im}(x)\frac{1}{\pi r},\hfill \\ 2e^{2\pi r\text{Im}(x)},\hfill & r\text{Re}(x),\text{Im}(x)0\hfill \end{array}$$
(0.26)
and
$$|\mathrm{tan}(\pi rx)+\sqrt{1}|\{\begin{array}{cc}4e^{2\pi r\text{Im}(x)},\hfill & \text{Im}(x)\frac{1}{\pi r},\hfill \\ 2e^{2\pi r\text{Im}(x)},\hfill & r\text{Re}(x),\text{Im}(x)0.\hfill \end{array}$$
(0.27)
Therefore we write
$$J_\pm (r,\epsilon )=\pm \sqrt{1}_{C_\pm (\epsilon )}g_r(x)\text{d}x+_{C_\pm (\epsilon )}(\mathrm{tan}(\pi rx)\sqrt{1})g_r(x)\text{d}x.$$
The estimates on $`\mathrm{tan}(\pi rx)\pm \sqrt{1}`$ can be used (see Appendix B) to prove that
$$\left|\underset{\mu =\pm 1}{}_{C_\mu (\epsilon )}(\mathrm{tan}(\pi rx)\mu \sqrt{1})g_r(x)\text{d}x\right|K_1\frac{1}{r},$$
(0.28)
where $`K_1`$ is a constant independent of $`r`$ and $`\epsilon `$. Let now
$$\mathrm{\Phi }(x)=\frac{1}{2\pi \sqrt{1}}\left(\text{Li}_2(e^{2\pi \sqrt{1}x})\text{Li}_2(e^{2\pi \sqrt{1}x})\right).$$
(0.29)
Note that $`\mathrm{\Phi }`$ is analytic on $`D=\{x|\text{Re}(x)\}`$ but not in the points $``$, so here we see the reason for using the small deformation parameter $`\epsilon `$. We have
$`{\displaystyle _{C_\pm (\epsilon )}}g_r(x)\text{d}x`$ $`=`$ $`{\displaystyle _{C_\pm (\epsilon )}}e^{r\mathrm{\Phi }(x)}\text{d}x`$
$`+{\displaystyle _{C_\pm (\epsilon )}}\left(\mathrm{exp}\left(I_\gamma (\pi 2\pi x)I_\gamma (\pi +2\pi x)\right)1\right)e^{r\mathrm{\Phi }(x)}\text{d}x.`$
However, as we will see in Appendix B, the estimate in Lemma 3 implies that
$$\left|_{C_\mu (\epsilon )}\left(\mathrm{exp}\left(I_\gamma (\pi 2\pi x)I_\gamma (\pi +2\pi x)\right)1\right)e^{r\mathrm{\Phi }(x)}\text{d}x\right|\frac{K_2\text{Log}(r)}{r}e^{\frac{r}{2\pi }\text{Vol}(K)}$$
(0.30)
for $`\mu =\pm 1`$, where $`K_2`$ is a constant independent of $`r`$ and $`\epsilon `$. We will see below that the estimates (0.28) and (0.30) imply that the leading order large $`r`$ asymptotics of $`J_K^{}(r)`$ is given by
$$J_K^{}(r)\frac{r^2}{2}\left(_{C_{}(\epsilon )}e^{r\mathrm{\Phi }(x)}\text{d}x_{C_+(\epsilon )}e^{r\mathrm{\Phi }(x)}\text{d}x\right),$$
(0.31)
to which we can apply the saddle point method, see e.g. \[6, Chap. 5\]. First we determine the stationary points of the phase function $`\mathrm{\Phi }`$. On $`D`$ we have
$$\mathrm{\Phi }^{}(x)=\text{Log}\left(1e^{2\pi \sqrt{1}x}\right)+\text{Log}\left(1e^{2\pi \sqrt{1}x}\right).$$
If we put $`z=e^{2\pi \sqrt{1}x}`$, then $`\mathrm{\Phi }^{}(x)=0`$ implies that
$$z^2z+1=0.$$
(0.32)
The equation (0.32) has the solutions $`z_\pm =e^{\pm \sqrt{1}\pi /3}`$. We have $`1z_\pm =1/2i\sqrt{3}/2`$ which both have norm $`1`$ and are each others conjugate, so
$$\text{Log}\left(1z_+\right)+\text{Log}\left(1z_{}\right)=0.$$
We note that $`z_\pm =e^{\pm 2\pi \sqrt{1}\frac{1}{6}}`$ correspond to the $`x`$โpoints $`\pm 1/6+`$. These points are non-degenerate critical points. In fact,
$$\mathrm{\Phi }^{\prime \prime }(x)=2\pi \sqrt{1}\frac{e^{2\pi \sqrt{1}x}+1}{e^{2\pi \sqrt{1}x}1}$$
on $`D`$, so in particular $`\mathrm{\Phi }^{\prime \prime }(x_\pm )=\pm 2\pi \sqrt{3}`$ for $`x_\pm \pm 1/6+`$. The imaginary part of $`\mathrm{\Phi }(x)`$ is zero for $`x`$ and
$$\mathrm{\Phi }(x)=\frac{1}{2\pi }\text{Im}\left(\text{Li}_2(e^{2\pi \sqrt{1}x})\text{Li}_2(e^{2\pi \sqrt{1}x})\right)=\frac{1}{\pi }\text{Cl}_2(2\pi x),$$
by (0.24). Let $`x_\pm \pm 1/6+`$. By (0.25) we have $`\text{Cl}_2(2\pi x_{})=\text{Cl}_2(2\pi x_+)=\text{Cl}_2(\pi /3)=\text{Vol}(K)/2`$, i.e.
$$\mathrm{\Phi }(x_\pm )=\frac{1}{2\pi }\text{Vol}(K).$$
By Cauchyโs theorem we have
$$_{C_{}(\epsilon )}e^{r\mathrm{\Phi }(x)}\text{d}x_{C_+(\epsilon )}e^{r\mathrm{\Phi }(x)}\text{d}x=2_{C_{}(\epsilon )}e^{r\mathrm{\Phi }(x)}\text{d}x=2_{C_+(\epsilon )}e^{r\mathrm{\Phi }(x)}\text{d}x.$$
Deform $`C_{}(\epsilon )`$ to $`[\epsilon ,1\epsilon ]`$ keeping the end points fixed. This does not change the integral $`_{C_{}(\epsilon )}e^{r\mathrm{\Phi }(x)}\text{d}x`$. Let $`x_0=5/6`$. By terminology borrowed from \[6, Sec. 5.4\] the axis of the saddle point $`x_0`$ is the real axis (i.e. the directions of steepest descent are along the real axis). From the analysis of \[6, Sec. 5.7\] it follows that we can find a $`\delta >0`$ (independent of $`r`$ and $`\epsilon `$) such that $`[x_0\delta ,x_0+\delta ][1/(4r),11/(4r)]`$ and such that we have an asymptotic expansion
$$_\delta ^\delta e^{r\mathrm{\Phi }(x_0+t)}\text{d}t\frac{1}{3^{1/4}\sqrt{r}}e^{\frac{r}{2\pi }\text{Vol}(K)}\left(1+\underset{n=1}{\overset{\mathrm{}}{}}d_nr^n\right)$$
in the limit $`r\mathrm{}`$, where the $`d_n`$โs are certain complex numbers. Finally we note that
$$\left|_\epsilon ^{x_0\delta }e^{r\mathrm{\Phi }(t)}\text{d}t+_{x_0+\delta }^{1\epsilon }e^{r\mathrm{\Phi }(t)}\text{d}t\right|e^{rc},$$
where $`c=\mathrm{max}\{\text{Cl}_2(2\pi (x_0\delta ))/\pi ,Cl_2(2\pi (x_0+\delta ))/\pi \}<\text{Vol}(K)/2\pi `$, see above (0.25). We have shown
###### Lemma 4
The leading order large $`r`$ asymptotics of $`J_K^{}(r)`$ is given by
$$J_K^{}(r)3^{1/4}r^{3/2}\mathrm{exp}\left(\frac{r}{2\pi }\text{Vol}(K)\right).$$
(0.33)
In fact
$$J_K^{}(r)=3^{1/4}r^{3/2}\mathrm{exp}\left(\frac{r}{2\pi }\text{Vol}(K)\right)+\text{O}\left(r\text{Log}(r)\mathrm{exp}\left(\frac{r}{2\pi }\text{Vol}(K)\right)\right)$$
in the limit $`r\mathrm{}`$.
In particular
$$\underset{r\mathrm{}}{lim}\frac{2\pi \text{Log}(J_K^{}(r))}{r}=\text{Vol}(K)$$
as predicted by the volume conjecture of Kashaev and Murakami, Murakami and as proven by Ekholm and others, see . However, the arguments of Ekholm and others canโt see the finer details of the asymptotic behaviour (0.33), namely the polynomial part $`3^{1/4}r^{3/2}`$.
###### Remark 2
The complement $`S^3K`$ of the figure $`8`$ knot can be decomposed into two ideal hyperbolic tetrahedra each parametrized by a certain complex number. This decomposition then defines a hyperbolic structure on this complement if and only if a certain set of conditions is satisfied. In fact, if the complex parameters for the two tetrahedra are respectively $`a`$ and $`b`$, then these conditions are equivalent to $`a=b`$ and $`b^2b+1=0`$, which is equation (0.32) after substituting $`b`$ for $`z`$. We refer to \[23, Sec. 3\] for more details. This phenomenon, that one finds the hyperbolicity equation for the knot complement as the equation for the stationary points of the phase function, seems to be a general principal for hyperbolic knots as argued by Thurston and Yokota, cf. , . However, there are major unsolved analytic difficulties in their approach. Basically they conjecture that one can carry out an asymptotic analysis similar to the one we carried out above for the figure $`8`$ knot. To prove their conjecture one first has to show how to give an exact (multi-dimensional) contour integral formula for the Jones polynomial of a hyperbolic knot like our (0.21). A main part consists of choosing a correct (multi-dimensional) contour. Secondly, one has to carry out an asymptotic analysis similar to the one leading to (0.31). This analysis is relatively simple for the figure $`8`$ knot due to the fact that we have a single (one-dimensional) contour in this case. In general one gets a contour of dimension $`>1`$ and the asymptotic analysis is expected to be harder (as also illustrated by the asymptotic analysis of the double-contour integral expression for the invariant $`\overline{\tau }_r(M_{p/q})`$ in Lemma 2, see next section.)
### 4.3 The large $`r`$ asymptotics of $`\tau _r(M_{p/q})`$
In Sec. 5 we will see that the signs of the phases in the asymptotics of $`\overline{\tau }_r(M_{p/q})=\overline{\tau _r(M_{p/q})}`$ agrees with the ChernโSimons values, hence we work with this conjugate invariant. Because of (0.16) we can always obtain the asymptotic expansion of $`\tau _r(M_{p/q})`$ by complex conjugation or by replacing either $`p`$ by $`p`$ or $`q`$ by $`q`$. By Lemma 2 we have
$$\overline{\tau }_r(M_{p/q})=\beta _1(r)\underset{n/|q|}{}_{C_r^1\times C_r^2}\mathrm{cot}(\pi rx)\mathrm{tan}(\pi ry)\overline{f}_{n,r}(x,y)\text{d}x\text{d}y,$$
(0.34)
where
$$\beta _1(r)=\frac{i\text{sign}(q)r}{4\sqrt{|q|}}e^{\frac{3\pi i}{4}\text{sign}(pq)}\mathrm{exp}\left(\frac{\pi i}{2r}\left[3\text{sign}(pq)\frac{p}{q}+\text{S}\left(\frac{p}{q}\right)\right]\right)$$
(0.35)
and
$$\overline{f}_{n,r}(x,y)=\mathrm{sin}\left(\frac{\pi }{q}(x2nd)\right)e^{2\pi ir\left(\frac{dn^2}{q}\frac{p}{4q}x^2+\frac{n}{q}xxy\right)}\frac{S_\gamma (\pi +2\pi (xy))}{S_\gamma (\pi +2\pi (x+y))},$$
where $`\gamma =\pi /r`$ as usual. Let for $`k,l`$
$$\mathrm{\Omega }_{k,l}=\{(x,y)^2\text{Re}(x)+\text{Re}(y)[k,k+1],\text{Re}(x)\text{Re}(y)[l,l+1]\}.$$
For $`(x,y)\mathrm{\Omega }_{k,l}`$, we have by (0.18) and (0.23) that
$`\overline{f}_{n,r}(x,y)`$ $`=`$ $`\mathrm{sin}\left({\displaystyle \frac{\pi }{q}}(x2nd)\right)\left(1+e^{2\pi i(xy)r}\right)^l\left(1+e^{2\pi i(x+y)r}\right)^ke^{2\pi ir\mathrm{\Phi }_n(x,y)}`$
$`\times \mathrm{exp}\left(I_\gamma (\pi +2\pi (xy+l))I_\gamma (\pi +2\pi (x+yk))\right),`$
where
$$\mathrm{\Phi }_n(x,y)=\frac{dn^2}{q}\frac{p}{4q}x^2+\frac{n}{q}xxy+\frac{1}{4\pi ^2}\left(\text{Li}_2(e^{2\pi i(x+y)})\text{Li}_2(e^{2\pi i(xy)})\right).$$
We note that this expression and the above expression for $`\overline{f}_{n,r}`$ are only valid for $`(x,y)^2`$ satisfying the two conditions
(i) $`\text{Re}(x)+\text{Re}(y)\text{Im}(x)+\text{Im}(y)0,`$
(ii) $`\text{Re}(x)\text{Re}(y)\text{Im}(x)\text{Im}(y)0.`$
Observe that for $`k,l\{0,1\},`$ which corresponds to the four different $`\mathrm{\Omega }_{k,l}`$ intersecting $`C_r^1\times C_r^2`$, we have that
$$\left(1+e^{2\pi i(xy)r}\right)^l\left(1+e^{2\pi i(x+y)r}\right)^k=\underset{(a,b)F_{k,l}}{}e^{2\pi i(a(x+y)+b(xy))r},$$
where $`F_{k,l}=\{(a,b)\{0,1\}^2ak,bl\}`$. Hence, for such $`k,l`$ we have
$`\overline{f}_{n,r}(x,y)`$ $`=`$ $`{\displaystyle \underset{(a,b)F_{k,l}}{}}\mathrm{sin}\left({\displaystyle \frac{\pi }{q}}(x2nd)\right)e^{2\pi ir\mathrm{\Phi }_n^{a,b}(x,y)}`$
$`\times \mathrm{exp}\left(I_\gamma (\pi +2\pi (xy+l))I_\gamma (\pi +2\pi (x+yk))\right)`$
for $`(x,y)\mathrm{\Omega }_{k,l}`$, where
$$\mathrm{\Phi }_n^{a,b}(x,y)=a(x+y)+b(xy)+\mathrm{\Phi }_n(x,y).$$
(0.36)
Let
$$\mathrm{\Omega }_{k,l}^{\mu ,\nu }=\{(x,y)\mathrm{\Omega }_{k,l}\mu \text{Im}(x)0,\nu \text{Im}(y)0\}.$$
###### Conjecture 3
There exists surfaces $`\mathrm{\Sigma }_{k,l,a,b}^{\mu ,\nu ,n}\mathrm{\Omega }_{k,l}^{\mu ,\nu }`$ for $`(k,l)\{0,1\}^2`$, $`(a,b)F_{k,l}`$, $`(\mu ,\nu )\{\pm 1\}^2`$ and $`n/|q|`$ such that the leading order large $`r`$ asymptotics of the quantum invariant is given by
$`\overline{\tau }_r(M_{p/q})`$ $``$ $`{\displaystyle \frac{i\text{sign}(q)r}{4\sqrt{|q|}}}e^{\frac{3\pi i}{4}\text{sign}(pq)}{\displaystyle \underset{n/|q|}{}}{\displaystyle \underset{(k,l)\{0,1\}^2}{}}{\displaystyle \underset{(a,b)F_{k,l}}{}}{\displaystyle \underset{(\mu ,\nu )\{\pm 1\}^2}{}}\mu \nu `$ (0.37)
$`\times {\displaystyle _{\mathrm{\Sigma }_{k,l,a,b}^{\mu ,\nu ,n}}}\mathrm{sin}\left({\displaystyle \frac{\pi }{q}}(x2nd)\right)e^{2\pi ir\mathrm{\Phi }_n^{a,b}(x,y)}\text{d}x\text{d}y.`$
Moreover, if $`p/q0`$ the surfaces $`\mathrm{\Sigma }_{k,l,a,b}^{\mu ,\nu ,n}\mathrm{\Omega }_{k,l}^{\mu ,\nu }`$ can be chosen such that any critical point of $`\mathrm{\Phi }_n^{a,b}`$ belonging to $`\mathrm{\Sigma }_{k,l,a,b}^{\mu ,\nu ,n}`$ also belongs to the set $`๐ฎ`$ in (0.5).
The rational behind this conjecture is that we anticipate an analysis of the expression (0.34) paralleling our analysis of $`J_K^{}(r)`$ should be applicable. I.e. $`\mathrm{tan}`$ and $`\mathrm{cot}`$ should be approximated by $`\pm \sqrt{1}`$ depending on the signs of $`\text{Im}(y)`$ and $`\text{Im}(x)`$ and $`\overline{f}_{n,r}(x,y)`$ by $`_{(a,b)F_{k,l}}\mathrm{sin}\left(\frac{\pi }{q}(x2nd)\right)e^{2\pi ir\mathrm{\Phi }_n^{a,b}(x,y)}`$ for some deformation of the part of $`C_r^1\times C_r^2`$ which is contained in $`\mathrm{\Omega }_{k,l}`$. We have partial analytic results supporting this conjecture. The factor in front of the sum is simply the leading term in the large $`r`$ asymptotics of $`\beta _1(r)`$, see (0.35). In case $`p/q=0`$ our results in Appendix C show that we have to include critical points not belonging to $`๐ฎ`$ in our surfaces $`\mathrm{\Sigma }_{k,l,a,b}^{\mu ,\nu ,n}`$.
In the remaining part of this section we will compute the large $`r`$ asymptotics of the right-hand side of (0.37). The only task left is to calculate the leading term in the large $`r`$ asymptotic of the integrals
$$I=_{\mathrm{\Sigma }_{k,l,a,b}^{\mu ,\nu ,n}}\mathrm{sin}\left(\frac{\pi }{q}(x2nd)\right)e^{2\pi ir\mathrm{\Phi }_n^{a,b}(x,y)}\text{d}x\text{d}y,$$
(0.38)
where $`k,l\{0,1\}`$, $`(a,b)F_{k,l}`$, $`\mu ,\nu \{\pm 1\}`$ and $`n`$. By the properties of the surfaces $`\mathrm{\Sigma }_{k,l,a,b}^{\mu ,\nu ,n}`$ postulated in Conjecture 3 this leading asymptotics should be calculable by the saddle point method. Let us give some details. Assume that $`(x_0,y_0)`$ is a critical point of $`\mathrm{\Phi }=\mathrm{\Phi }_n^{a,b}`$ belonging to the surface $`\mathrm{\Sigma }=\mathrm{\Sigma }_{k,l,a,b}^{\mu ,\nu ,n}`$ and assume that
$$\text{Im}\left((\mathrm{\Phi }(x,y)\mathrm{\Phi }(x_0,y_0))\right)0$$
for all $`(x,y)`$ in a small neighborhood of $`(x_0,y_0)`$ in $`\mathrm{\Sigma }`$ with equality only in $`(x_0,y_0)`$. In this case we say that $`\mathrm{\Sigma }`$ pass through the critical point $`(x_0,y_0)`$ in the directions of steepest descent. Moreover, assume this is satisfied for all critical points of $`\mathrm{\Phi }`$ inside $`\mathrm{\Sigma }`$. Then the main contribution to the integral $`I`$ in the large $`r`$ limit is given by integrating over small neighborhoods of the critical points in $`\mathrm{\Sigma }`$.
Let us therefore begin by computing the stationary points of $`\mathrm{\Phi }_n^{a,b}`$. To this end, it is more convenient to work with the functions
$$\mathrm{\Psi }_n^{a,b}(x,y)=ax+by+\mathrm{\Phi }_n(x,y),$$
(0.39)
$`a,b,n`$, so $`\mathrm{\Phi }_n^{a,b}=\mathrm{\Psi }_n^{a+b,ab}`$.
Let $`a,b,n`$ and put $`\mathrm{\Psi }=\mathrm{\Psi }_n^{a,b}`$. Let $`z=e^{2\pi ix}`$ and $`w=e^{2\pi iy}`$. Then
$`2\pi i{\displaystyle \frac{\mathrm{\Psi }}{x}}(x,y)=2\pi i(ay){\displaystyle \frac{p}{2q}}2\pi ix+{\displaystyle \frac{2\pi in}{q}}+\text{Log}(1zw)\text{Log}(1zw^1),`$
$`2\pi i{\displaystyle \frac{\mathrm{\Psi }}{y}}(x,y)=2\pi i(bx)+\text{Log}(1zw)+\text{Log}(1zw^1),`$ (0.40)
where we have to assume (which we will also assume in what follows) that both $`zw`$ and $`zw^1`$ are different from $`1`$.
We will need to specify a certain square root of $`z`$, namely let $`v=e^{\pi ix}`$. Then $`\frac{\mathrm{\Psi }}{x}(x,y)=0`$ implies that
$$v^p=\left(\frac{wv^2}{1v^2w}\right)^q,$$
(0.41)
and $`\frac{\mathrm{\Psi }}{y}(x,y)=0`$ implies that
$$(1v^2w)(wv^2)=v^2w.$$
(0.42)
Both equations (0.41) and (0.42) are independent of $`a`$, $`b`$, and $`n`$. We note that $`(v,w)=(0,0)`$ is the only solution to (0.42) with $`v`$ or $`w`$ equal to zero. Note, moreover, that $`(v,w)`$ is a common solution to (0.41) and (0.42) if and only if $`(\overline{v},\overline{w})`$ is a common solution to these two equations. By writing (0.41) as
$$(1)^qv^{p+2q}=\left(\frac{1v^2w}{1v^2w}\right)^q$$
and (0.42) as
$$(1v^2w)(1v^2w)=w$$
we see that $`(v,w)`$ is a nonzero solution to (0.41) and (0.42) if and only if $`(v^1,w)`$ is such a solution. If $`p`$ is even, then $`(v,w)`$ is a solution to (0.41) and (0.42) if and only if $`(v,w)`$ is such a solution.
Let us oppositely begin with a common solution $`(v,w)^{}\times ^{}`$ to (0.41) and (0.42), where, as usual, $`^{}=\{0\}`$. Write $`v=e^{\pi ix}`$ and $`w=e^{2\pi iy}`$ with $`\text{Re}(x)]1,1]`$ and $`\text{Re}(y)]1/2,1/2]`$, i.e. $`x=\frac{1}{\pi i}\text{Log}(v)`$ and $`y=\frac{1}{2\pi i}\text{Log}(w)`$. One now easily deduce from (0.42) that there exists a unique $`b`$ such that
$$0=2\pi i(bx)+\text{Log}\left(1e^{2\pi i(x+y)}\right)+\text{Log}\left(1e^{2\pi i(xy)}\right),$$
(0.43)
and from (0.41) we deduce that there exists a unique $`n`$ such that
$$0=2\pi iy\frac{p}{q}\pi ix+\frac{2\pi in}{q}+\text{Log}\left(1e^{2\pi i(x+y)}\right)\text{Log}\left(1e^{2\pi i(xy)}\right).$$
(0.44)
That is, there exists a unique pair of integers $`b,n`$ such that $`(x,y)`$ is a stationary point of $`\mathrm{\Psi }_n^{0,b}`$.
###### Remark 3
Let us make a slight digression by giving some general remarks about the set of solutions to (0.42). Assume that $`v,w^{}`$ and let $`z=v^2`$. Then (0.42) can be written in the following two ways
$`z^2\left(w+{\displaystyle \frac{1}{w}}+1\right)z+1=0,`$
$`w^2\left(z+{\displaystyle \frac{1}{z}}1\right)w+1=0.`$ (0.45)
It is straightforward to see that if $`(z,w)`$ is a solution to (3) with $`w\{0\}`$, then $`z`$ is real and positive if $`w>0`$, $`z`$ is real and negative if $`w]\mathrm{},(3+\sqrt{5})/2][(3\sqrt{5})/2,0[`$, and $`zS^1`$ if
$$(3+\sqrt{5})/2w(3\sqrt{5})/2.$$
If $`(z,w)`$ is a solution to (3) with $`z\{0\}`$, then $`w`$ is real and negative if $`z<0`$, $`w`$ is real and positive if $`z]0,(3\sqrt{5})/2][(3+\sqrt{5})/2,\mathrm{}[`$, and $`wS^1`$ if $`z[(3\sqrt{5})/2,(3+\sqrt{5})/2]`$.
Later on we will be particularly interested in common solutions to (0.41) and (0.42) with $`vS^1`$ and $`w\{0\}`$. Assume that $`(v,w)`$ is such a solution, and write $`v=e^{\pi ix}`$, $`x]1,1]`$. In that case the 2nd of the equations in (3) is
$$w^2+(12\text{Re}(z))w+1=0.$$
(0.46)
Since $`zS^1`$ we have $`12\text{Re}(z)[1,3]`$ and since $`w`$ is real we also have $`|12\text{Re}(z)|2`$, so $`12\text{Re}(z)[2,3]`$ or $`\mathrm{cos}(2\pi x)=\text{Re}(z)[1,1/2]`$. By (0.46) we find that
$$w=w_\pm (x)=\mathrm{cos}(2\pi x)\frac{1}{2}\pm \sqrt{\mathrm{cos}^2(2\pi x)\mathrm{cos}(2\pi x)\frac{3}{4}}.$$
(0.47)
If $`w>0`$ we have already seen that $`z]0,\mathrm{}[`$. But then $`z=1`$ so $`w^2w+1=0`$ by (3) contradicting the fact that $`w`$ is real.
Taking absolute values we get from (0.41) that $`|1zw|=|wz|`$ and then from (0.42) that $`|w|=|wz|^2=(wz)(w\overline{z})=w^22\text{Re}(z)w+1`$ (still assuming that $`(v,w)S^1\times `$). Since $`w<0`$ this equation is equivalent to (0.46). In other words we have to consider the argument equation following from (0.41) to obtain information from that equation not obtainable from (0.42).
Let us now turn to the second derivative of $`\mathrm{\Psi }`$ in a critical point $`(x_0,y_0)`$, i.e. the Hessian $`H=H(x_0,y_0)`$ of $`\mathrm{\Psi }`$ in $`(x_0,y_0)`$ (which is equal to the Hessian of $`\mathrm{\Phi }_n`$ in $`(x_0,y_0)`$). Put $`(v_0,w_0)=(e^{\pi ix_0},e^{2\pi iy_0})`$ and $`z_0=v_0^2`$. By a small computation, using the fact that $`(v_0,w_0)`$ is a solution to (0.42), we find that
$`H_{11}`$ $`=`$ $`{\displaystyle \frac{^2\mathrm{\Psi }}{x^2}}(x_0,y_0)={\displaystyle \frac{1}{w_0}}w_0{\displaystyle \frac{p}{2q}},`$
$`H_{12}=H_{21}`$ $`=`$ $`{\displaystyle \frac{^2\mathrm{\Psi }}{xy}}(x_0,y_0)=z_0{\displaystyle \frac{1}{z_0}},`$
$`H_{22}`$ $`=`$ $`{\displaystyle \frac{^2\mathrm{\Psi }}{y^2}}(x_0,y_0)={\displaystyle \frac{1}{w_0}}w_0.`$ (0.48)
Let us begin by establishing under which circumstances the critical point $`(x_0,y_0)`$ is non-degenerate. By (4.3)
$$det(H)=\left(w_0+\frac{1}{w_0}\right)^2\left(z_0+\frac{1}{z_0}\right)^2+\frac{p}{2q}\left(w_0\frac{1}{w_0}\right).$$
By (3) we have $`w_0+\frac{1}{w_0}=z_0+\frac{1}{z_0}1`$ so
$$det(H)=12\left(z_0+\frac{1}{z_0}\right)+\frac{p}{2q}\left(w_0\frac{1}{w_0}\right).$$
(0.49)
If we introduce $`y_0^{}=iy_0`$ we get the formula
$$det(H)=14\mathrm{cos}(2\pi x_0)+\frac{p}{2q}\mathrm{sinh}(2\pi y_0^{}).$$
(0.50)
We are particularly interested in critical points where $`w_0<0`$ and $`z_0S^1`$ (or equivalently $`(x_0,y_0)๐ฎ`$, see (0.5)). By definition of $`z_0`$ we have $`z_01/z_0=2i\mathrm{sin}(2\pi x_0)`$ and by (4.3) we then get
$$det(H)=\left(\frac{1w_0^2}{w_0}\frac{p}{2q}\right)\left(\frac{1w_0^2}{w_0}\right)+4\mathrm{sin}^2(2\pi x_0).$$
(0.51)
Thus if $`w_0]1,0[`$ and $`p/q0`$ or if $`w_0]\mathrm{},1[`$ and $`p/q0`$ we have $`det(H)>0`$. If $`w_0=1`$ we also have $`det(H)>0`$ by (0.47). Since $`\text{Re}(y_0)\frac{1}{2}+`$ we get by (0.49) the alternative formula
$$det(H)=14\mathrm{cos}(2\pi x_0)+\frac{p}{2q}\mathrm{sinh}\left(2\pi \text{Im}(y_0)\right).$$
(0.52)
Finally we get by (0.47) and (0.49) that
$$det(H)=14\mathrm{cos}(2\pi x_0)\pm \frac{p}{q}\sqrt{\mathrm{cos}^2(2\pi x_0)\mathrm{cos}(2\pi x_0)\frac{3}{4}}.$$
(0.53)
for $`w=w_\pm (x_0)`$.
###### Proposition 1
Let the situation be as above and assume that $`(x_0,y_0)`$ is a critical point of $`\mathrm{\Psi }`$ belonging to the set $`๐ฎ`$ in (0.5). If $`|p/q|<\sqrt{20}`$ then $`(x_0,y_0)`$ is non-degenerate. If $`|p/q|>\sqrt{20}`$ then $`(x_0,y_0)`$ is non-degenerate except if
$$\text{Re}(z_0)=\mathrm{cos}(2\pi x_0)=\frac{1}{2}+\frac{4\left|\frac{p}{q}\right|\sqrt{\frac{p^2}{q^2}15}}{\frac{p^2}{q^2}16}$$
and $`w_0=w_+(x_0)`$ if $`p/q<0`$ and $`w_0=w_{}(x_0)`$ if $`p/q>0`$.
* Proof. Let $`z_0=e^{2\pi ix_0}`$ and $`w_0=e^{2\pi iy_0}`$ as above and put $`c=\mathrm{cos}(2\pi x_0)`$. Assume that $`(x_0,y_0)`$ is degenerate, i.e. that $`det(H)=0`$. Then
$$\left(16\frac{p^2}{q^2}\right)c^2+\left(\frac{p^2}{q^2}8\right)c+1+\frac{3}{4}\frac{p^2}{q^2}=0$$
by (0.53). If $`p^2/q^2<15`$ this equation does not have real solutions, so $`p^2/q^215`$. If $`p^2/q^2=16`$ we find that $`c=13/8`$ giving a contradiction. Thus $`p^2/q^2[15,16[]16,\mathrm{}[`$ and
$$c=\frac{8\frac{p^2}{q^2}\pm 2\left|\frac{p}{q}\right|\sqrt{\frac{p^2}{q^2}15}}{2\left(16\frac{p^2}{q^2}\right)}.$$
(0.54)
By this and $`c[1,1/2]`$, see below (0.46), we get that
$$\frac{3}{2}\frac{p^2}{q^2}20\pm \left|\frac{p}{q}\right|\sqrt{\frac{p^2}{q^2}15}\frac{p^2}{q^2}12$$
(0.55)
for $`p^2/q^2[15,16[`$ with the opposite inequalities for $`p^2/q^2]16,\mathrm{}[`$. For $`p^2/q^2[15,16[`$ we have $`\frac{3}{2}\frac{p^2}{q^2}20>0`$ so we have a plus in front of $`\left|\frac{p}{q}\right|\sqrt{\frac{p^2}{q^2}15}`$ in (0.55) and in that case the first inequality in (0.55) is equivalent to
$$\frac{5}{4}\left(\frac{p}{q}\right)^445\left(\frac{p}{q}\right)^2+4000,$$
(0.56)
which forces $`p^2/q^2[16,20]`$ giving a contradiction.
Therefore $`p^2/q^2>16`$ and (0.54) and $`c[1,1/2]`$ leads, as already stated, to (0.55) with the opposite inequalities. Since $`p^2/q^212>0`$ we again have a plus in front of $`\left|\frac{p}{q}\right|\sqrt{\frac{p^2}{q^2}15}`$, and the second inequality is automatically satisfied for all $`p^2/q^216`$. This time the first inequality leads to (0.56) with the opposite inequalities, so we can conclude that $`p^2/q^220`$ and in fact therefore $`|p/q|>\sqrt{20}`$. The formula for $`\mathrm{cos}(2\pi x_0)`$ is just (0.54) with the plus sign in front of $`\left|\frac{p}{q}\right|\sqrt{\frac{p^2}{q^2}15}`$.
We note that the only other solutions $`(z,w)S^1\times ]\mathrm{},0[`$ to (0.46) satisfying that $`\text{Re}(z)`$ is equal to the right-hand side of (0.54) are $`(z_0,1/w_0)`$, $`(\overline{z_0},w_0)`$ and $`(\overline{z_0},1/w_0)`$. By the remarks following (0.51) we get that among these three points only the point $`(\overline{z_0},w_0)`$ can actually satisfy, that the right-hand side of (0.51) is zero. We conjecture that also in the case $`|p/q|>\sqrt{20}`$ all critical points of $`\mathrm{\Psi }`$ which belong to the set $`๐ฎ`$ are non-degenerate. To confirm this one has probably to use the argument equation following from (0.41). By a direct check we have confirmed this conjecture in the cases $`|p/q|\{14/3,5,16/3,6,20/3,22/3,8,26/3,28/3,10,32/3\}`$.
We are now ready to calculate the contribution to the leading order large $`r`$ asymptotics of the integral (0.38) coming from a critical point $`(x,y)`$ of the phase function $`\mathrm{\Phi }_n^{a,b}`$ belonging to $`๐ฎ\mathrm{\Sigma }_{k,l,a,b}^{\mu ,\nu ,n}`$.
The integral $`I`$ in (0.38) is a double contour integral and as such is calcalculated by choosing a contour for each variable $`x`$ and $`y`$. We thus have an โinnerโ integral and an โouterโ integral and in general the contour for the inner integral can depend on where we are on the outer contour. Here we will, however, only consider the very simple case where these two contours are independent of each other.
Therefore let $`\alpha ,\beta S^1`$ and let $`\gamma _\alpha ,\gamma _\beta :I_\delta `$ be given by $`\gamma _\alpha (s)=x_0+\alpha s`$, $`\gamma _\beta (s)=y_0+\beta s`$, where $`I_\delta =[\delta ,\delta ]`$ for a sufficiently small $`\delta >0`$. According to the saddle point method the main contributions to the integral (0.38) comes from integrating over a small neighborhood of $`๐ฎ\mathrm{\Sigma }`$ in $`\mathrm{\Sigma }=\mathrm{\Sigma }_{k,l,a,b}^{\mu ,\nu ,n}`$. We are therefore lead to consider an integral of the form
$`K(x_0,y_0)`$ $`=`$ $`{\displaystyle _{\gamma _\alpha }}\left({\displaystyle _{\gamma _\beta }}\mathrm{sin}\left({\displaystyle \frac{\pi }{q}}(x2nd)\right)e^{2\pi ir\mathrm{\Psi }(x,y)}\text{d}y\right)\text{d}x`$
$`=`$ $`\alpha \beta {\displaystyle _\delta ^\delta }\left({\displaystyle _\delta ^\delta }\mathrm{sin}\left({\displaystyle \frac{\pi }{q}}(x_0+\alpha s2nd)\right)e^{2\pi ir\mathrm{\Psi }(x_0+\alpha s,y_0+\beta t)}\text{d}t\right)\text{d}s,`$
where $`\mathrm{\Psi }=\mathrm{\Psi }_n^{a,b}`$ as above. (If $`(x_0,y_0)`$ is a non-degenerate critical point on the boundary of $`\mathrm{\Sigma }`$, then one or both or the integrals $`_\delta ^\delta `$ should be replaced by $`_0^\delta `$ (or $`_\delta ^0`$) and the contribution coming from that point in Theorem 4 should be multiplied by $`\frac{1}{2}`$ or $`\frac{1}{4}`$.) By a Taylor expansion we find that
$$\mathrm{\Psi }(x_0+\alpha s,y_0+\beta t)=\mathrm{\Psi }(x_0,y_0)+\frac{1}{2}A\left(\begin{array}{c}s\\ t\end{array}\right),\left(\begin{array}{c}s\\ t\end{array}\right)+h(s,t),$$
where $`A=\text{diag}(\alpha ,\beta )H\text{diag}(\alpha ,\beta )`$ and
$$\left(\begin{array}{c}x_1\\ y_1\end{array}\right),\left(\begin{array}{c}x_2\\ y_2\end{array}\right)=x_1x_2+y_1y_2$$
for $`(x_1,y_1),(x_2,y_2)^2`$, and where $`h(s,t)`$ is a remainder term being a sum of terms of the form $`c_{m,n}s^mt^n`$, $`m,n\{0,1,2,\mathrm{},\mathrm{}\}`$, $`m+n3`$, $`c_{m,n}`$. We search for $`\alpha `$ and $`\beta `$ such that there exists a $`\delta >0`$ satisfying
$$\text{Re}\left(2\pi i(\mathrm{\Psi }(x_0+\alpha s,y_0+\beta t)\mathrm{\Psi }(x_0,y_0))\right)<0$$
for all $`(s,t)I_\delta \times I_\delta \{(0,0)\}`$. This amounts to finding $`\alpha `$ and $`\beta `$ such that
$$\text{Im}\left(A\left(\begin{array}{c}s\\ t\end{array}\right),\left(\begin{array}{c}s\\ t\end{array}\right)\right)>0$$
for all $`(s,t)^2\{(0,0)\}`$. Since
$$\text{Im}\left(A\left(\begin{array}{c}s\\ t\end{array}\right),\left(\begin{array}{c}s\\ t\end{array}\right)\right)=\text{Im}(A)\left(\begin{array}{c}s\\ t\end{array}\right),\left(\begin{array}{c}s\\ t\end{array}\right)$$
this corresponds to finding $`\alpha `$ and $`\beta `$ such that $`\text{Im}(A)`$ is positive definite. It is essential for the method that such $`\alpha `$ and $`\beta `$ exist. (If not, the analysis becomes more difficult since higher order terms in the Taylor expansion of $`\mathrm{\Psi }(x_0+\alpha s,y_0+\beta t)`$ need to be involved, and it is likely that the analysis can not be carried through in this case.) For that reason we make the following definition.
###### Definition 1
A critical point $`(x,y)`$ of the phase function $`\mathrm{\Psi }`$ is called positive definite if $`(x,y)`$ is non-degenerate and there exists $`(\alpha ,\beta )S^1\times S^1`$ such that $`\text{Im}(A)`$ is positive definite, where
$$A=\text{diag}(\alpha ,\beta )H\text{diag}(\alpha ,\beta ),$$
where $`H`$ is the Hessian of $`\mathrm{\Psi }`$ in $`(x,y)`$.
Assume that $`(x_0,y_0)`$ is positive definite and let $`(\alpha ,\beta )S^1\times S^1`$ such that $`\text{Im}(A)`$ is positive definite. Then the main contribution to the integral $`K(x_0,y_0)`$ in the limit of large $`r`$ is given by
$$K_{\text{main}}(x_0,y_0)=\alpha \beta \mathrm{sin}\left(\frac{\pi }{q}(x_02nd)\right)e^{2\pi ir\mathrm{\Psi }(x_0,y_0)}_^2e^{\pi irA\left(\begin{array}{c}s\\ t\end{array}\right),\left(\begin{array}{c}s\\ t\end{array}\right)}\text{d}s\text{d}t,$$
and this integral can be evaluated using the results of \[12, Sec. 3.4\]. In fact,
$$_^2e^{\pi irA\left(\begin{array}{c}s\\ t\end{array}\right),\left(\begin{array}{c}s\\ t\end{array}\right)}\text{d}s\text{d}t=\frac{1}{r}\left(det(iA)\right)^{1/2},$$
where $`A`$ is independent of $`r`$. Note here that the set $`S`$ of complex symmetric $`2\times 2`$โmatrices $`B`$ with $`\text{Re}(B)`$ positive definite is an open convex set in the $`3`$โdimensional complex vector space of symmetric $`2\times 2`$โmatrices. It follows that their is a unique analytic branch of $`B\left(det(B)\right)^{1/2}`$ on $`S`$ such that $`\left(det(B)\right)^{1/2}>0`$ for $`B`$ real. We have used that branch in the above result. We note that $`det(iA)=(\alpha \beta )^2det(H)`$. In conclusion we can state the following result.
###### Theorem 4
Let $`๐ฎ_{k,l,a,b}^{\mu ,\nu ,n}`$ be the set of critical points of $`\mathrm{\Phi }_n^{a,b}`$ in $`\mathrm{\Sigma }_{k,l,a,b}^{\mu ,\nu ,n}`$. IF Conjecture 3 is true and if all the critical points of $`\mathrm{\Phi }_n^{a,b}`$ in $`\mathrm{\Sigma }_{k,l,a,b}^{\mu ,\nu ,n}`$ are positive definite, $`k,l=0,1`$, $`(a,b)F_{k,l}`$, $`\mu ,\nu \{\pm 1\}`$, $`n/|q|`$, then the leading order large $`r`$ asymptotics of the quantum invariant is given by
$`\overline{\tau }_r(M_{p/q})`$ $``$ $`{\displaystyle \frac{i\text{sign}(q)}{4\sqrt{|q|}}}e^{\frac{3\pi i}{4}\text{sign}(pq)}{\displaystyle \underset{n/|q|}{}}{\displaystyle \underset{(k,l)\{0,1\}^2}{}}{\displaystyle \underset{(a,b)F_{k,l}}{}}{\displaystyle \underset{(\mu ,\nu )\{\pm 1\}^2}{}}\mu \nu `$
$`\times {\displaystyle \underset{(x,y)๐ฎ_{k,l,a,b}^{\mu ,\nu ,n}}{}}(det(H_{(x,y)}\mathrm{\Phi }_n^{a,b}))^{1/2}`$
$`\times \mathrm{sin}\left({\displaystyle \frac{\pi }{q}}(x2nd)\right)e^{2\pi ir\mathrm{\Phi }_n^{a,b}(x,y)},`$
where $`H_{(x,y)}\mathrm{\Phi }_n^{a,b}`$ is the Hessian of $`\mathrm{\Phi }_n^{a,b}`$ in $`(x,y)`$.
According to the above theorem the growth rate of the quantum invariants of $`M_{p/q}`$ is $`r^0`$ if we can prove that the values of the phase functions $`\mathrm{\Phi }_n^{a,b}`$ in the points $`S_{k,l,a,b}^{\mu ,\nu ,n}`$ are real. This we do in Sec. 5.3, see Corollary 5. This is in agreement with our computer studies of the quantum invariants of $`M_{p/q}`$. To obtain Theorem 3 from Theorem 4, hence a leading asymptotics as predicted by the AEC, we have to prove i) that the union of the sets of critical points $`๐ฎ_{k,l,a,b}^{\mu ,\nu ,n}`$ corresponds to the flat (irreducible) $`\text{SU}(2)`$โconnections on $`M_{p/q}`$ and ii) that the values of the relevant phase functions $`\mathrm{\Phi }_n^{a,b}`$ in these critical points are equal to the ChernโSimons invariants of these flat connections. This we also do in Sec. 5.3. Note that we have used the expression (0.52) for the determinant of the Hessian $`H_{(x,y)}\mathrm{\Phi }_n^{a,b}`$ in Theorem 3.
###### Remark 4
The manifolds $`M_{p/q}`$, $`|p/q|\{1,2,3\}`$, are Seifert manifolds and by (0.11) holds for these manifolds if we put $`m_{\overline{\rho }}=4`$ for all $`\overline{\rho }_{p/q}^{}`$. Moreover, the results in Appendix C show that the part of the leading order large $`r`$ asymptotics of $`\tau _r(M_0)`$ associated to the irreducible $`\text{SU}(2)`$โconnections on $`M_0`$ is equal to the right-hand side of (0.11) if we put $`m_{\overline{\rho }}=4`$ for both points $`\overline{\rho }`$ in $`_0^{}`$. We actually expect that $`m_{\overline{\rho }}=4`$ for all $`\overline{\rho }_{p/q}^{}`$ for all rational $`p/q`$. Recall here that when we calculated the leading order large $`r`$ asymptotics of $`J_K^{}(r)`$ in Sec. 4.2 we had to part the contour $`C(\epsilon )`$ into two parts $`C_\pm (\epsilon )`$ because of the tangent factor in the contour integral (0.21). The relevant stationary point was in that case on the real axes, namely it was $`x_0=5/6`$, and to obtain a contour passing through that point we had to deform $`C_{}(\epsilon )`$ to $`[\epsilon ,1\epsilon ]`$ and $`C_+(\epsilon )`$ to $`[1\epsilon ,\epsilon ]`$. In that way the stationary point $`x_0`$ gave two equal contributions. (In fact we used Cauchyโs theorem so that we only had to work with either $`C_{}(\epsilon )`$ or $`C_+(\epsilon )`$, but that of course also produced the factor $`2`$.) In the case with the quantum invariants $`\overline{\tau }_r(M_{p/q})`$ we expect a similar phenomenon with the difference that each relevant stationary point contribute with $`4`$ instead of $`2`$ equal contributions. We expect that the relevant critical points $`(x,y)`$ belong to $`[1/3,2/3]\times \{y|\text{Re}(y)=1/2\}`$. Because of the cotangent and tangent factors we had to separate our double contour integral expression for $`\overline{\tau }_r(M_{p/q})`$ in (0.34) into $`4`$ parts causing the sum $`_{(\mu ,\nu )\{\pm 1\}^2}\mu \nu `$ in the asymptotic formula for $`\overline{\tau }_r(M_{p/q})`$. The contour $`C_r^1`$ was thus separated into the part with $`\text{Im}(x)0`$ and the part with $`\text{Im}(x)0`$. The first of these parts can be deformed to $`[0,1]`$ while the second can be deformed to $`[1,0]`$. The contour $`C_r^2`$ was separated into the part with $`\text{Im}(y)0`$ and the part with $`\text{Im}(y)0`$. Both parts can be deformed to a contour containing the part of the line $`\text{Re}(y)=1/2`$ containing the $`y`$โs in the relevant stationary points. In this way we obtain as claimed that each relevant critical point contribute $`4`$ times. Like in the case of the the knot invariant $`J_K^{}(r)`$ we find that the signs $`\mu `$ and $`\nu `$ are cancelled by taking care of the orientations of the different involved contours.
Let us end this section by examine to what extend the critical points $`S_{k,l,a,b}^{\mu ,\nu ,n}`$ are positive definite, see Definition 1. Since $`A`$ is symmetric we have $`\text{Im}(A)_{ij}=\text{Im}(A_{ij})`$. Here
$`A_{11}`$ $`=`$ $`\alpha ^2H_{11},`$
$`A_{12}=A_{21}`$ $`=`$ $`\alpha \beta H_{12},`$
$`A_{22}`$ $`=`$ $`\beta ^2H_{22},`$
where $`H`$ is the Hessian in (4.3). We have $`z_0=e^{2\pi ix_0}`$ and therefore $`H_{12}=H_{21}=2i\mathrm{sin}(2\pi x_0)`$. Moreover, $`w_0\{w_\pm (x_0)\}`$ using notation from (0.47) and $`H_{11}=H_{22}\frac{p}{2q}`$, where
$$H_{22}=\frac{1}{w_\pm (x_0)}w_\pm (x_0)=2\sqrt{\mathrm{cos}^2(2\pi x_0)\mathrm{cos}(2\pi x_0)\frac{3}{4}}.$$
If $`M=\left(\begin{array}{cc}a& c\\ c& b\end{array}\right)`$ is a real symmetric matrix, then $`M`$ has real eigenvalues, and $`M`$ is positive definite if and only if both of these eigenvalues are positive, i.e. if and only if $`a`$ and $`b`$ are both positive and $`ab>c^2`$. Now assume that $`M=\text{Im}(A)`$. Since $`\text{Re}(z_0)[1,1/2]`$ we have that $`H_{12}`$ is zero if and only if $`z_0=1`$. In that case $`H_{22}=\sqrt{5}`$ which is irrational. Therefore $`H_{11}0`$ and we can always choose $`\alpha `$ and $`\beta `$ so $`a>0`$ and $`b>0`$ and hence $`ab>0=c^2`$. Therefore assume in what follows that $`H_{12}0`$. If $`H_{11}`$ and $`H_{22}`$ are both positive (meaning that $`w_0=w_{}(x_0)`$) we can let $`\alpha =\beta =e^{i\pi /4}`$ giving $`a=H_{11}`$, $`b=H_{22}`$ and $`c=0`$. If $`H_{11}`$ and $`H_{22}`$ are both negative (meaning that $`w_0=w_+(x_0)`$) then we can let $`\alpha =\beta =e^{i\pi /4}`$ giving $`a=H_{11}`$, $`b=H_{22}`$ and $`c=0`$. If $`\mu =\text{sign}(H_{11})\text{sign}(H_{22})`$ we can let $`\alpha =e^{\mu i\pi /4}`$ and $`\beta =e^{\mu i\pi /4}`$ giving $`a=|H_{11}|`$, $`b=|H_{22}|`$ and $`c=D`$, writing $`H_{12}=iD`$, $`D`$. Thus $`det(M)=det(H)`$, so if $`det(H)<0`$ this choice of $`\alpha `$ and $`\beta `$ works. The case $`\text{sign}(H_{11})\text{sign}(H_{22})`$ and $`det(H)>0`$ is more problematic, and the above analysis is actually too simplified as the examples $`|p/q|\{4/3,10/3,16/3,22/3\}`$ show.
One can in an even simpler way see that the above approach is too simplified, namely by considering the cases where $`H_{11}=0`$ or $`H_{22}=0`$. In that case $`a=0`$ or $`b=0`$ so $`\text{Im}(A)`$ is not positive definite. We have $`H_{22}=0`$ if and only if $`\mathrm{cos}(2\pi x_0)=1/2`$ and that case occurs if and only if $`p=6m+3`$ for some integer $`m`$.
Since $`\mathrm{cos}(2\pi x_0)[1,1/2]`$ the entry $`H_{11}`$ can not be zero if $`|p/q|>\sqrt{20}`$. For $`|p/q|<\sqrt{20}`$, $`H_{11}=0`$ implies that $`\mathrm{cos}(2\pi x_0)=\frac{1}{2}\sqrt{1+\frac{p^2}{16q^2}}`$. Finally note that if $`H_{11}=H_{22}=0`$ then $`p/q=0`$, but in that case $`H_{22}0`$ for critical points $`(x_0,y_0)๐ฎ`$, see Appendix C.
The solution to these problems is to let the contour for the โinnerโ integral depend on the contour for the โouterโ integral. We will not give more details here, but refer to a later paper. Note that the expression for the matrix $`A`$ in Definition 1 in this more general approach is different from the one stated in that definition.
## 5 Classical ChernโSimons theory on $`M_{p/q}`$
In this section we will describe the classical theory, that is the classical ChernโSimons theory on the manifolds $`M_{p/q}`$. The $`\text{SU}(2)`$ ChernโSimons functional is a map with values in $`/`$ defined on the set $`๐`$ of gauge equivalence classes of connections in a principal $`\text{SU}(2)`$ bundle on $`M_{p/q}`$ (all such bundles being trivializable). Inside $`๐`$ sits the moduli space $`_{p/q}`$ of flat $`\text{SU}(2)`$โconnections on $`M_{p/q}`$, that is the set of classical solutions to the $`\text{SU}(2)`$ ChernโSimons field theory. Recall that
$$_{p/q}=\text{Hom}(\pi _1(M_{p/q}),\text{SU}(2))/\text{SU}(2)$$
from which it is clear that $`_{p/q}`$ is a compact space. The ChernโSimons functional is constant on the connected components of $`_{p/q}`$, thus there are only finitely many different values on flat connections. Let $`_{p/q}^{}`$ be the subset of $`_{p/q}`$ consisting of nonabelian representations. Recall that these representations correspond to the irreducible connections, while the abelian representations correspond to the reducible connections.
The main results in this section are Theorem 7 and Theorem 8 which tie up the ChernโSimons theory to the large $`r`$ asymptotics of $`\overline{\tau }_r(M_{p/q})`$ by showing that a certain subset of the critical points of the phase functions $`\mathrm{\Phi }_n^{a,b}`$ corresponds to $`_{p/q}^{}`$. Under this correspondence, the ChernโSimons functional is taken to the phase functions $`\mathrm{\Phi }_n^{a,b}`$.
We begin by giving a description of $`_{p/q}`$ following Riley , and Kirk & Klassen .
### 5.1 The moduli space of flat $`\text{SU}(2)`$โconnections on $`M_{p/q}`$
In the following $`\pi =\pi _1\left(S^3\text{nbd}(K)\right)`$ denotes the knot group of the figure $`8`$ knot. We have a presentation
$$\pi =x,y|wx=yw,$$
(0.57)
where $`w=[x^1,y]`$, and where $`\mu =x`$ and $`\lambda =yx^1y^1x^2y^1x^1y`$ are the elements of $`\pi `$ corresponding to the meridian and the preferred longitude of $`K`$. The $`\text{SL}(2,)`$ representation variety of $`\pi `$ was analyzed by Riley , relevant to our work. Consider a group $`G`$ given by a presentation
$$G=x,y|wx=yw,$$
where $`w=x^{\epsilon _1}y^{\epsilon _2}x^{\epsilon _3}\mathrm{}y^{\epsilon _{\alpha 1}}`$, where $`\alpha `$ is odd and $`\epsilon _j=\epsilon _{\alpha j}=\pm 1`$, $`j=1,2,\mathrm{},\alpha 1`$. Such groups are denoted $`2`$โbridge kmot groups by Riley since they generalize the $`2`$โbridge knot groups. Following Riley we say that a representation $`\psi :G\text{SL}(2,)`$ is affine when the image of $`\psi `$ fixes exactly one point in $`^1`$ and not affine when this image has no fixed points. We note that if $`\psi `$ is nonabelian then $`\psi `$ is affine if and only if $`\psi (x)`$ and $`\psi (y)`$ have a common eigenvector, and $`\psi `$ is not affine if these two matrices have no common eigenvector. Let $`H`$ be some subgroup of $`\text{SL}(2,)`$. Then we will say that two representations $`\psi _1,\psi _2:G\text{SL}(2,)`$ are $`H`$โequivalent if they are conjugate to each other by a matrix in $`H`$, i.e. if there exists a matrix $`UH`$ such that $`\psi _2(\gamma )=U\psi _1(\gamma )U^1`$ for all $`\gamma G`$. In particular, we will say that $`\psi _1`$ and $`\psi _2`$ are equivalent if they are $`\text{SL}(2,)`$โequivalent. For $`(t,u)^{}\times `$ we put
$`C_0(t)`$ $`=`$ $`\left(\begin{array}{cc}t& 1\\ 0& 1\end{array}\right),D_0(t,u)=\left(\begin{array}{cc}t& 0\\ tu& 1\end{array}\right),`$
$`C_1(t)`$ $`=`$ $`\left(\begin{array}{cc}t& 1\\ 0& t^1\end{array}\right),D_1(t,u)=\left(\begin{array}{cc}t& 0\\ u& t^1\end{array}\right),`$
$`C_2(t)`$ $`=`$ $`\left(\begin{array}{cc}t& t^1\\ 0& t^1\end{array}\right),D_2(t,u)=\left(\begin{array}{cc}t& 0\\ tu& t^1\end{array}\right).`$
We note that $`C_\nu (t)`$ and $`D_\nu (t,u)`$ are elements of $`\text{SL}(2,)`$, $`\nu =1,2`$. If $`s`$ is a square root of $`t`$ and $`V(s)=\text{diag}(s,s^1)`$ then
$$V(s)C_2(t)V(s)^1=C_1(t),V(s)D_2(t,u)V(s)^1=D_1(t,u),$$
(0.61)
and
$$C_0(t)=sC_2(s),D_0(t,u)=sD_2(s,u).$$
(0.62)
Let $`W_\nu (t,u)`$ denote the matrix obtained by replacing $`x`$ and $`y`$ by respectively $`C_\nu (t)`$ and $`D_\nu (t,u)`$ in the expression for $`w`$, and let
$$\varphi (t,u)=W_{11}+(1t)W_{12},$$
where $`W=W(t,u)=W_0(t,u)`$. We let $`\rho _{(t,u)}`$ be the assignment $`xC_2(t)`$, $`yD_2(t,u)`$. We have the following $`\text{SL}(2,)`$ version of \[28, Theorem 1\].
###### Theorem 5
Let $`(s,u)^{}\times `$. None of the assignments $`\rho _{(s,u)}`$ extend to an abelian $`\text{SL}(2,)`$โrepresentation of $`G`$. The assignment $`\rho _{(s,u)}`$ extends to a nonabelian representation $`\rho _{(s,u)}:G\text{SL}(2,)`$ if and only if
$$\varphi (s^2,u)=0.$$
(0.63)
Conversely, if $`\psi :G\text{SL}(2,)`$ is a nonabelian representation, then there exists a pair $`(s,u)^{}\times `$ satisfying (0.63) such that $`\psi `$ and $`\rho _{(s,u)}`$ are equivalent. When $`\psi `$ is affine this pair is unique, and when $`\psi `$ is not affine the pair $`(s,u)`$ can only be replaced by $`(s^1,u)`$.
* Proof. The theorem follows by results of , . The assignment $`\rho _{(s,u)}`$ extends to a $`\text{SL}(2,)`$โrepresentation of $`G`$ if and only if $`W_2(s,u)C_2(s,u)=D_2(s,u)W_2(s,u)`$. The matrices $`C_2(s)`$ and $`D_2(s,u)`$ commute if and only if $`s=\pm 1`$ and $`u=0`$, and since $`C_2(\pm 1)=\pm C_2(1)`$ is different from $`D_2(\pm 1,0)=\pm D_2(1,0)`$ we have that the assignment $`\rho _{(\pm 1,0)}`$ does not extend to a $`\text{SL}(2,)`$โrepresentation of $`G`$. (We note that if $`W=W(1,0)`$ then $`W_{21}=0`$ since the matrices $`C_0(1),D_0(1,0)=I`$ and their inverses are upper triangular. But we also have that $`\varphi (1,0)=0`$ would imply that $`W_{11}=0`$ contradicting the fact that $`W`$ is invertible. Therefore $`\varphi (1,0)0`$.)
Let $`\sigma =_{j=1}^{\alpha 1}\epsilon _j`$. Then $`s^\sigma W_2(s,u)=W(s^2,u)`$. By (the proof of) \[28, Theorem 1\] we have $`W(s^2,u)C_0(s^2)=D_0(s^2,u)W(s^2,u)`$ if and only if $`\varphi (s^2,u)=0`$, so by (0.62) we find that $`\rho _{(s,u)}`$ extends to a (necessarily nonabelian) $`\text{SL}(2,)`$โrepresentation of $`G`$ if and only if $`\varphi (s^2,u)=0`$.
If $`\psi :G\text{SL}(2,)`$ is an arbitrary nonabelian representation it follows by \[29, Lemma 7\] and (0.61) that there exists a pair $`(s,u)^{}\times `$ (necessarily satisfying (0.63)) such that $`\psi `$ and $`\rho _{(s,u)}`$ are equivalent. By the above any such pair is different from $`(\pm 1,0)`$ and by \[29, Lemma 8\] and (0.61) we then get the final statement of the theorem.
If $`(s,u)^{}\times `$ with $`\varphi (s^2,u)=0`$ then $`\rho _{(s,u)}`$ is affine if and only if $`C_2(s)`$ and $`D_2(s,u)`$ have a common eigenvector. But this happens exactly when $`u=0`$ or $`u=(ss^1)^2`$.
Let us now restrict to the case where $`G`$ is the figure $`8`$ knot group $`\pi `$. Then
$$\varphi (t,u)=u^2+\left(3(t+t^1)\right)(u+1)$$
(0.64)
so in particular $`\varphi (s^2,0)=3s^2s^2=0`$ if and only if $`s^43s^2+1=0`$ i.e. if and only if $`s=\mu _1\sqrt{(3+\mu _2\sqrt{5})/2}`$ for some $`\mu _1,\mu _2\{\pm 1\}`$. If $`u=(ss^1)^2=s^2+s^22`$ then $`\varphi (s^2,u)=u^2+(3u2)(u+1)=u^2+1u^2=1`$, so we conclude that $`\rho _{(s,u)}`$ is affine if and only if $`u=0`$ and $`s^43s^2+1=0`$. Let
$$๐ฉ=\text{Hom}(\pi ,\text{SL}(2,))/\text{SL}(2,)$$
be the space of conjugacy classes of $`\text{SL}(2,)`$โrepresentations of $`\pi `$ and let $`๐ฉ_{\text{nab}}`$ be the subset consisting of classes represented by nonabelian $`\text{SL}(2,)`$โrepresentations. Moreover, let
$$\stackrel{~}{๐ฉ}=\{(s,u)^{}\times |\varphi (s^2,u)=0\},$$
(0.65)
and let $`\mathrm{\Phi }:\stackrel{~}{๐ฉ}๐ฉ`$ be the map which maps $`(s,u)`$ to the class represented by $`\rho _{(s,u)}`$. We have shown
###### Corollary 1
The image of $`\mathrm{\Phi }`$ is $`๐ฉ_{\text{nab}}`$. If we let
$$\stackrel{~}{๐ฉ}_0=\{(\mu _1\sqrt{(3+\mu _2\sqrt{5})/2},0)|\mu _1,\mu _2\{\pm 1\}\}$$
then $`\mathrm{\Phi }|_{\stackrel{~}{๐ฉ}_0}:\stackrel{~}{๐ฉ}_0๐ฉ`$ is injective and $`\mathrm{\Phi }^1(\mathrm{\Phi }(s,u))=\{(s,u),(s^1,u)\}`$ for any $`(s,u)\stackrel{~}{๐ฉ}\stackrel{~}{๐ฉ}_0`$.
Recall that $`M_{p/q}`$ denotes the closed oriented $`3`$โmanifold obtained by surgery on $`S^3`$ along the figure $`8`$ knot with rational surgery coefficient $`p/q`$. The representation $`\rho _{(s,u)}`$, $`(s,u)\stackrel{~}{๐ฉ}`$, extends to a $`\text{SL}(2,)`$โrepresentation of $`\pi _1(M_{p/q})`$ if and only if
$$\rho _{(s,u)}(\mu )^p\rho _{(s,u)}(\lambda )^q=1.$$
To analyze this criterion it is an advantage to diagonalize $`\rho _{(s,u)}(\lambda )`$. Assume in the following that $`(s,u)\stackrel{~}{๐ฉ}`$. By a rather long but completely elementary and straightforward calculation we find that
$$\rho _{(s,u)}(\lambda )=\left(\begin{array}{cc}\lambda _{11}(s,u)& \lambda _{12}(s,u)\\ 0& \lambda _{11}(s^1,u)\end{array}\right),$$
where
$$\lambda _{11}(s,u)=1+s^22s^2+s^4+u(s^2s^2)$$
and $`\lambda _{12}(s,u)=2u(1u)`$ if $`s^2=1`$ and
$$\lambda _{12}(s,u)=\frac{\lambda _{11}(s,u)\lambda _{11}(s^1,u)}{s^21}$$
for $`s^21`$. We note that $`\lambda _{11}(s^1,u)=\lambda _{11}(s,u)^1`$.
In general, if $`A=\left(\begin{array}{cc}\alpha & \beta \\ 0& \alpha ^1\end{array}\right)\text{SL}(2,)`$, then $`A`$ can be diagonalized if and only if $`A`$ is not parabolic, i.e. if and only if $`\text{tr}(A)\pm 2`$ or equivalently if and only if $`\alpha \pm 1`$ (except, of course, if $`\beta =0`$). If $`\alpha \pm 1`$ then $`(1,0)`$ is an eigenvector with eigenvalue $`\alpha `$ and $`(\beta /(\alpha \alpha ^1),1)`$ is an eigenvector with eigenvalue $`\alpha ^1`$. If $`s^2=1`$ then $`\lambda _{11}(s,u)=1`$ and $`\lambda _{12}(s,u)=\pm i2\sqrt{3}`$. If $`s^21`$ then $`\lambda _{11}(s,u)=\pm 1`$ if and only if $`\lambda _{12}(s,u)=0`$. Therefore $`\rho _{(s,u)}(\lambda )`$ is diagonalizable (or diagonal) if and only if $`s^21`$. In case $`s^21`$ and $`\lambda _{12}(s,u)0`$ we have
$$\frac{\lambda _{12}(s,u)}{\lambda _{11}(s,u)\lambda _{11}(s,u)^1}=\frac{1}{s^21}=\frac{s^1}{ss^1}.$$
We conclude that if $`s^21`$, then $`^2`$ has a basis consisting of a set of common eigenvectors for the matrices $`\rho _{(s,u)}(\mu )`$ and $`\rho _{(s,u)}(\lambda )`$, namely $`u_1=(1,0)`$ and $`u_2=(1/(s^21),1)`$. If we let $`\stackrel{~}{\rho }_{(s,u)}:\pi \text{SL}(2,)`$ be the representation $`\stackrel{~}{\rho }_{(s,u)}(\gamma )=U^1\rho _{(s,u)}(\gamma )U`$, where $`U\text{SL}(2,)`$ with $`j`$th column $`u_j`$, we therefore have
$$\stackrel{~}{\rho }_{(s,u)}(x)=\text{diag}(s,s^1),\stackrel{~}{\rho }_{(s,u)}(\lambda )=\text{diag}(\lambda _{11}(s,u),\lambda _{11}(s,u)^1).$$
(0.66)
In particular, $`\rho _{(s,u)}:\pi \text{SL}(2,)`$, $`s^21`$, extends to a representation of $`\pi _1(M_{p/q})`$ if and only if
$$s^p=\lambda _{11}(s,u)^q.$$
(0.67)
Recall here that $`\rho _{(s,u)}`$ and $`\rho _{(s^1,u)}`$ are equivalent for $`(s,u)\stackrel{~}{๐ฉ}\stackrel{~}{๐ฉ}_0`$, cf. Corollary 1. But as noted above $`\lambda _{11}(s^1,u)=\lambda _{11}(s,u)^1`$ in accordance with (0.67). For $`(s,u)\stackrel{~}{๐ฉ}_0`$ we have $`\lambda _{11}(s,u)=1`$ and $`|s|1`$, so $`\rho _{(s,u)}`$ extends to a representation of $`\pi _1(M_{p/q})`$ if and only if $`p=0`$.
A direct check shows that if $`s^2=1`$ then $`\rho _{(s,u)}`$ does not extend to a representation of $`\pi _1(M_{p/q})`$ for any rational number $`p/q`$. In fact, if $`s=\pm 1`$, then
$$\rho _{(s,u)}(\mu )^p=(\pm 1)^{|p|}\left(\begin{array}{cc}1& p\\ 0& 1\end{array}\right)$$
and
$$\rho _{(s,u)}(\lambda )^q=\left(\left(\begin{array}{cc}1& \epsilon i2\sqrt{3}\\ 0& 1\end{array}\right)\right)^q=(1)^{|q|}\left(\begin{array}{cc}1& \epsilon i2q\sqrt{3}\\ 0& 1\end{array}\right)$$
for a $`\epsilon \{1,1\}`$. On the other hand, since $`\lambda _{11}(s,u)=1`$, we have that (0.67) is satisfied if $`s=1`$ and $`q`$ is even or $`s=1`$ and both $`p`$ and $`q`$ are odd. For the following we note that if $`s^2=1`$ then $`u^2+u+1=\varphi (1,u)=0`$ so $`u`$ is not real.
We are mostly interested in the $`\text{SU}(2)`$โrepresentations of $`\pi `$. Let in the following $`_{\text{nab}}`$ be the set of conjugacy classes of nonabelian $`\text{SU}(2)`$โrepresentations of $`\pi `$.
###### Proposition 2
Let $`(s,u)\stackrel{~}{๐ฉ}`$. The representation $`\rho _{(s,u)}:\pi \text{SL}(2,)`$ is $`\text{SL}(2,)`$โequivalent to a representation $`\pi \text{SU}(2)`$ if and only if $`|s|=1`$ and $`u`$ is real. If we write $`s=e^{2\pi i\theta }`$, $`\theta ]1/2,1/2]`$, then $`u\{u_\pm \}`$, where $`u_\pm =u_\pm (\theta )`$ are the two solutions to $`\varphi (e^{4\pi i\theta },u)=0`$, i.e.
$$u_\pm (\theta )=\mathrm{cos}(4\pi \theta )\frac{3}{2}\pm \sqrt{\mathrm{cos}^2(4\pi \theta )\mathrm{cos}(4\pi \theta )\frac{3}{4}}.$$
Since $`u`$ is real we have $`\theta [1/3,1/6][1/6,1/3]`$. The representation $`\rho _{(e^{2\pi i\theta },u_\epsilon )}`$, $`\theta [1/3,1/6][1/6,1/3]`$ and $`\epsilon \{\pm \}`$, is $`\text{SL}(2,)`$โequivalent to a $`\text{SU}(2)`$โrepresentation $`\overline{\rho }_{\theta ,\epsilon }`$ which satisfies
$$\overline{\rho }_{\theta ,\epsilon }(\mu )=\text{diag}(e^{2\pi i\theta },e^{2\pi i\theta }),\overline{\rho }_{\theta ,\epsilon }(\lambda )=\text{diag}(L_\epsilon ,L_\epsilon ^1),$$
where $`\mu `$ and $`\lambda `$ are the elements of $`\pi `$ corresponding to the meridian and the preferred longitude of $`K`$, and
$$L_\pm =L_\pm (\theta )=\lambda _{11}(e^{2\pi i\theta },u_\pm )=1+e^{4\pi i\theta }2e^{4\pi i\theta }+e^{8\pi i\theta }+u_\pm \left(e^{4\pi i\theta }e^{4\pi i\theta }\right).$$
We note that $`\overline{\rho }_{\theta ,\epsilon }`$ and $`\overline{\rho }_{\theta ,\epsilon }`$ are $`\text{SU}(2)`$โequivalent. In particular, the space $`_{\text{nab}}`$ can be parametrized by the two arcs $`(e^{2\pi i\theta },u_+(\theta ))`$, $`\theta [1/6,1/3]`$, and $`(e^{2\pi i\theta },u_{}(\theta ))`$, $`\theta [1/6,1/3]`$. These two arcs only coincide at the endpoints, so topologically $`_{\text{nab}}`$ is a circle.
This proposition follows from \[28, Proposition 4\], see also \[18, Proposition 5.4\]. For a more geometric argument determining the topological type of $`_{\text{nab}}`$, see .
For $`\theta [1/3,1/6][1/6,1/3]`$ and $`\epsilon \{\pm \}`$, the representation $`\overline{\rho }_{\theta ,\epsilon }`$ extends to a representation of $`\pi _1(M_{p/q})`$ if and only if
$$\overline{\rho }_{\theta ,\epsilon }(\mu )^p\overline{\rho }_{\theta ,\epsilon }(\lambda )^q=1,$$
i.e. if and only if
$$e^{2\pi ip\theta }=L_\epsilon (\theta )^q.$$
(0.68)
From this (use e.g. (5.2)) we see that
###### Corollary 2
Let $`p/q`$ be arbitrary. The moduli space of irreducible flat $`\text{SU}(2)`$โconnections on $`M_{p/q}`$ is a finite set.
Let us end this section by finding the abelian $`\text{SU}(2)`$โrepresentations of $`\pi _1(M_{p/q})`$ (up to equivalence). Therefore, let $`\theta ]1/2,1/2]`$ and let $`\rho _\theta `$ be the assignment
$$\rho _\theta (\mu )=\text{diag}(e^{2\pi i\theta },e^{2\pi i\theta }).$$
By (0.57) this assignment extends to an abelian $`\text{SU}(2)`$โrepresentation of $`\pi `$ for any $`\theta ]1/2,1/2]`$ by letting $`\rho _\theta (y)=\rho _\theta (x)`$. Moreover, any abelian $`\text{SU}(2)`$โrepresentation of $`\pi `$ is $`\text{SU}(2)`$โequivalent to $`\rho _\theta `$ for some $`\theta ]1/2,1/2]`$. For any $`\theta ]1/2,1/2]`$ we have $`\rho _\theta (\lambda )=1`$, and $`\rho _\theta `$ extends to a representation of $`\pi _1(M_{p/q})`$ if and only if $`\rho _\theta (\mu )^p=1`$, i.e. if and only if $`p\theta `$. If $`A=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)`$ then $`A\rho _\theta (\mu )A^1=\rho _\theta (\mu )`$ so we can assume that $`\theta [0,1/2]`$. Note, moreover, that two matrices $`\text{diag}(e^{i\varphi },e^{i\varphi })`$ and $`\text{diag}(e^{i\psi },e^{i\psi })`$, $`\varphi ,\psi ]\pi ,\pi ]`$, are conjugate in $`\text{SU}(2)`$ if and only if $`\varphi =\psi `$ or $`\varphi =\psi `$. We conclude
###### Proposition 3
For $`p0`$ the set of conjugacy classes of abelian $`\text{SU}(2)`$โrepresen-tations of $`\pi _1(M_{p/q})`$ is given by
$$\left\{[\rho _{j/|p|}]\right|j=0,1,\mathrm{},\left[\frac{|p|}{2}\right]\},$$
where $`[|p|/2]`$ is the integer part of $`|p|/2`$. For $`p=0`$ the set of conjugacy classes of abelian $`\text{SU}(2)`$โrepresentations of $`\pi _1(M_{p/q})`$ is given by
$$\left\{[\rho _\theta ]\right|\theta [0,\frac{1}{2}]\},$$
so topologically this set is a closed interval.
### 5.2 ChernโSimons invariants
We begin by recalling formulas from for the ChernโSimons invariants of the flat $`\text{SU}(2)`$โconnections on $`M_{p/q}`$. The basic tool will be Theorem 6 below due to P. A. Kirk and E. P. Klassen.
Let $`M`$ be a closed oriented $`3`$โmanifold with a knot $`K`$ in its interior and let $`X`$ be the complement of a tubular neighborhood of $`K`$ in $`M`$. Moreover, let $`\mu `$ be a meridian of $`K`$ and $`\lambda `$ a longitude, both in $`X`$. Let $`G`$ be $`\text{SU}(2)`$ or $`\text{SL}(2,)`$. If $`\rho :\pi _1(X)G`$ is a representation, then $`\rho `$ extends to a $`G`$โrepresentation of $`\pi _1(M)`$ if and only if $`\rho (\mu )=1`$. Now assume that $`\rho _t:\pi _1(X)G`$ is a piecewise smooth path of representations, $`tI=[0,1]`$. Choose a piecewise smooth path $`g:IG`$ such that
$`g_t\rho _t(\mu )g_t^1`$ $`=`$ $`\text{diag}(e^{2\pi i\alpha (t)},e^{2\pi i\alpha (t)}),`$
$`g_t\rho _t(\lambda )g_t^1`$ $`=`$ $`\text{diag}(e^{2\pi i\beta (t)},e^{2\pi i\beta (t)})`$ (0.69)
for some piecewise smooth curves $`\alpha ,\beta `$. If $`G=\text{SU}(2)`$ this is always possible by \[18, Lemma 3.1\], and in that case $`\alpha `$ and $`\beta `$ are real-valued. If $`G=\text{SL}(2,)`$ the above is possible if the path $`\rho `$ avoids the parabolic representations (i.e. upper triangular with $`1`$s or $`1`$s on the diagonal), cf. \[18, Remark p. 354\]. (See the text in connection to (0.66) for the case of the figure $`8`$ knot.) In that case the curves $`\alpha `$ and $`\beta `$ are complex-valued. We then have
###### Theorem 6 ( \[18, Theorem 4.2\] )
Assume that $`\rho _0(\mu )=\rho _1(\mu )=1`$. Thinking of $`\rho _0`$ and $`\rho _1`$ as flat $`G`$โconnections on $`M`$, we have
$$\text{CS}(\rho _1)\text{CS}(\rho _0)=2_0^1\beta (t)\alpha ^{}(t)\text{d}t(mod),$$
where CS is the ChernโSimons functional associated to $`G`$.
We note that in case $`G=\text{SL}(2,)`$ the ChernโSimons functional takes values in $`/`$.
Next consider a knot $`K`$ in $`S^3`$ and let $`X`$ be the knot complement. Let $`p/q`$ be a rational number and let $`N_{p/q}`$ be the closed oriented $`3`$โmanifold obtained by $`p/q`$ surgery on $`S^3`$ along $`K`$. Let $`\mu `$ and $`\lambda `$ be classes in $`\pi _1(X)`$ represented by respectively a meridian and the preferred longitude of $`K`$. Choose integers $`c,d`$ such that $`pdqc=1`$. Let $`V`$ be a tubular neighborhood of $`K`$ considered as a subspace of $`N_{p/q}`$. We note that $`\mu ^{}=p\mu +q\lambda `$ and $`\lambda ^{}=c\mu +d\lambda `$ are represented by respectively a meridian of $`V`$ and a longitude of $`V`$. Assume that $`\rho _t:\pi _1(X)G`$, $`tI`$, is a piecewise smooth curve of representations from the trivial representation to a representation, which extends to a representation of $`\pi _1(N_{p/q})`$, i.e. $`\rho _1(\mu ^{})=1`$. Assume, moreover, that $`\rho _t`$ avoids the parabolic representations in case $`G=\text{SL}(2,)`$, and choose curves $`\alpha ,\beta `$ as in (5.2) with $`\alpha (0)=\beta (0)=0`$. By Theorem 6 we have
$`\text{CS}(\rho _1)`$ $`=`$ $`2{\displaystyle _0^1}(c\alpha (t)+d\beta (t))(p\alpha ^{}(t)+q\beta ^{}(t))\text{d}t`$
$`=`$ $`2{\displaystyle _0^1}\beta (t)\alpha ^{}(t)\text{d}tcp\alpha ^2(1)dq\beta ^2(1)2cq\alpha (1)\beta (1)(mod).`$
(We have corrected a sign error in \[18, Formula (\*) p. 361\].) We note that this expression is independent of the choice of $`c,d`$. The condition $`\rho _1(\mu ^{})=1`$ is equivalent to
$$p\alpha (1)+q\beta (1).$$
(0.71)
Now let $`K`$ be the figure $`8`$ knot and let $`\overline{\rho }_{\theta ,\epsilon }`$ be a $`\text{SU}(2)`$โrepresentation of $`\pi _1(M_{p/q})`$, i.e. (0.68) is satisfied. Following we determine a formula for $`\text{CS}(\overline{\rho }_{\theta ,\epsilon })`$. For later we will here pay special attention to the branches of the logarithm. By Proposition 2 we have
$`\text{Re}\left(L_{}(\theta )\right)`$ $`=`$ $`\text{Re}\left(L_+(\theta )\right)=2\mathrm{cos}^2(4\pi \theta )\mathrm{cos}(4\pi \theta )2,`$
$`\text{Im}\left(L_\pm (\theta )\right)`$ $`=`$ $`2\mathrm{sin}(4\pi \theta )\sqrt{\mathrm{cos}^2(4\pi \theta )\mathrm{cos}(4\pi \theta ){\displaystyle \frac{3}{4}}}`$ (0.72)
for all $`\theta [1/3,1/6][1/6,1/3]`$. From these identities we see that
$$L_\pm (1/3)=L_\pm (1/6)=1,L_\pm (1/4)=1,$$
(0.73)
and that $`\text{Im}(L_+)<0`$ and $`\text{Im}(L_{})>0`$ on $`]1/6,1/4[`$ with the opposite signs on $`]1/4,1/3[`$. We conclude that $`L_+(\theta )`$ and $`L_{}(\theta )`$ run through $`S^1`$ both beginning and ending in $`1`$, $`L_+`$ in the anti-clockwise and $`L_{}`$ in the clockwise direction, as $`\theta `$ runs through $`[1/6,1/3]`$. We can therefore use the principal logarithm Log to define continuous curves $`\beta _\pm :[1/6,1/3]`$ by
$$\beta _\pm (\theta )=\frac{1}{2\pi i}\text{Log}(L_\pm (\theta ))+f_\pm (\theta ),$$
(0.74)
where
$$f_+(\theta )=\{\begin{array}{cc}0,& \theta =\frac{1}{6},\hfill \\ 1,& \theta ]\frac{1}{6},\frac{1}{3}],\hfill \end{array}$$
(0.75)
and
$$f_{}(\theta )=\{\begin{array}{cc}0,& \theta [\frac{1}{6},\frac{1}{3}[,\hfill \\ 1,& \theta =\frac{1}{3}.\hfill \end{array}$$
(0.76)
By (0.68) and the definition of $`\beta _\pm `$ the representation $`\overline{\rho }_{\theta ,\epsilon }`$ extends to a representation of $`\pi _1(M_{p/q})`$ if and only if
$$p\theta +q\beta _\epsilon (\theta )$$
(0.77)
which is nothing but (0.71) (reparametrized). We note that $`\beta _\pm `$ are smooth on $`]\frac{1}{6},\frac{1}{3}[`$ but not in the end points $`1/6`$ and $`1/3`$. The terms $`f_\pm `$ have been chosen so that $`\beta _\pm (1/6)=1/2`$. This is needed for the proof of
###### Proposition 4 ( \[18, p. 362\] )
Let $`\theta [1/6,1/3]`$ and let $`\overline{\rho }_{\theta ,\pm }`$ be as in Proposition 2. If $`\overline{\rho }_{\theta ,\epsilon }`$ extends to a representation of $`\pi _1(M_{p/q})`$ for a $`\epsilon \{\pm 1\}`$ then
$$\text{CS}(\overline{\rho }_{\theta ,\epsilon })=\frac{1}{6}cp\theta ^2dq\beta _\epsilon ^2(\theta )2cq\theta \beta _\epsilon (\theta )2_{1/6}^\theta \beta _\epsilon (t)\text{d}t(mod),$$
where $`\beta _\pm `$ are the curves defined by (0.74).
Kirk and Klassen prove the above result by explicitly constructing a piecewise smooth path $`\rho :[0,1]\text{Hom}(\pi ,\text{SL}(2,))`$ from the trivial representation to $`\overline{\rho }_{\frac{1}{6},+}=\overline{\rho }_{\frac{1}{6},}`$ with piecewise smooth curves $`\alpha ,\beta :[0,1]`$ as in (5.2) satisfying $`\alpha (1)=\frac{1}{6}`$ and $`\beta (1)=\frac{1}{2}`$. Moreover, they use that $`_0^1\beta (t)\alpha ^{}(t)\text{d}t=\frac{1}{12}`$ and the fact that $`y\overline{\rho }_{y,\epsilon }`$, $`[1/6,\theta ]\text{Hom}(\pi ,\text{SU}(2))`$ is a path from $`\overline{\rho }_{\frac{1}{6},\epsilon }`$ to $`\overline{\rho }_{\theta ,\epsilon }`$ with associated functions $`\alpha (y)=y`$ and $`\beta (y)=\beta _\epsilon (y)`$. By the choice of the functions $`f_\pm `$, these $`\alpha `$โ and $`\beta `$โfunctions are continuations of the ones used for the path $`\rho `$ from the trivial representation to $`\overline{\rho }_{\frac{1}{6},\pm }`$.
Kirk and Klassen argue that $`_0^1\beta (t)\alpha ^{}(t)\text{d}t=\frac{1}{12}`$ using a comparison between a computer calculation and the ChernโSimons invariants of flat $`\text{SU}(2)`$โconnection on the Seifert manifold $`M_3`$. By following Kirk and Klassenโs argument \[18, pp. 361โ362\] it is actually not hard to give an explicit calculation of this integral. Let us give some details. The path $`\rho `$ from the trivial connection to $`\overline{\rho }_{\frac{1}{6},\pm }`$ consists of three parts, where $`\beta `$ is identically zero along the first two parts. We can therefore concentrate on the third part. Let $`a=(3+\sqrt{5})/2`$, and let $`\alpha :[0,1]`$ be a piecewise smooth curve from $`\frac{1}{2\pi i}\text{Log}(\sqrt{a})`$ to $`1/6`$. Let $`t(s)=e^{4\pi i\alpha (s)}`$ and choose a piecewise smooth solution $`u(s)`$ to $`\varphi (t(s),u(s))=0`$, where $`\varphi `$ is given by (0.64). Then $`s\rho _{(e^{2\pi i\alpha (s)},u(s))}=:\eta _s`$ is the third piece of our curve $`\rho `$ (reparametrized). Assuming $`\alpha (s)\frac{1}{2}`$ (so as to avoid the parabolic representations) we can diagonalize exactly as demonstrated after Corollary 1 and get
$$\stackrel{~}{\eta }_s(\mu )=\text{diag}(T,T^1),\stackrel{~}{\eta }_s(\lambda )=\text{diag}(\lambda _{11}(T,u),\lambda _{11}(T,u)^1),$$
where $`T=T(s)=e^{2\pi i\alpha (s)}`$, $`u=u(s)`$, $`\lambda _{11}`$ is as below Corollary 1, and where $`\stackrel{~}{\eta }_s(\gamma )=U^1\eta _s(\gamma )U`$ for $`\gamma \pi `$, where $`U=U(s)=\left(\begin{array}{cc}1& 1/(T^21)\\ 0& 1\end{array}\right)`$. This shows that $`\alpha (s)`$ indeed plays the role as the $`\alpha `$โcurve for our path $`\eta _s`$. The $`\beta `$โcurve should be a piecewise smooth curve $`\beta (s)`$ starting at zero such that
$$e^{2\pi i\beta (s)}=\lambda _{11}(e^{2\pi i\alpha (s)},u(s))$$
(0.78)
for $`s[0,1]`$. We note that $`\lambda _{11}(e^{2\pi i\alpha (1)},u(1))=1`$ so we must have $`\beta (1)\frac{1}{2}+`$. Since $`\overline{\rho }_{1/6,\pm }`$ extends to a $`\text{SU}(2)`$โrepresentation of $`\pi _1(M_3)`$ we get from (5.2) that
$$2_0^1\beta (s)\alpha ^{}(s)\text{d}s\frac{1}{12}+\frac{1}{3}\beta (1)(mod)$$
is a ChernโSimons value of a flat $`\text{SU}(2)`$โconnection on $`M_3`$, so $`_0^1\beta (s)\alpha ^{}(s)\text{d}s`$ is real.
We now only have to choose a nice $`\alpha `$โcurve. Let $`\delta `$ be a small positive parameter less than $`|\alpha (0)|<1/6`$, and let $`\alpha =\alpha _1+\alpha _2+\alpha _3`$, where $`\alpha _1`$ is the line segment $`[\alpha (0),i\delta ]`$, $`\alpha _2`$ is the part of the circle with centre zero and radius $`\delta `$ running from $`i\delta `$ to $`\delta `$, and $`\alpha _3=[\delta ,1/6]`$. Here the parameter $`\delta `$ is introduced to avoid passing through zero (so as to avoid parabolic representations). With this choice it is not hard to show that there is a unique continuous solution $`(u(s),\beta (s))`$ to $`\varphi (t(s),u(s))=0`$ and (0.78) with $`\beta (0)=0`$ and $`\beta (1)=\frac{1}{2}`$ and to show that $`_0^1\beta (s)\alpha ^{}(s)\text{d}s=\frac{1}{12}`$ for this solution. (Use that the integral is real and independent of $`\delta `$ and calculate its $`\delta 0_+`$ limit).
Proposition 4 gives a formula for the ChernโSimons invariants of the irreducible flat $`\text{SU}(2)`$โconnections on $`M_{p/q}`$. Let us also determine the ChernโSimons invariants of the reducible flat $`\text{SU}(2)`$โconnections on $`M_{p/q}`$. Therefore let $`\rho _{j/|p|}`$ be as in Proposition 3, $`p0`$, and let $`\alpha (t)=jt/|p|`$ and $`\beta (t)=0`$, $`t[0,1]`$, and get by (5.2) that
$$\text{CS}(\rho _{j/|p|})=cpj^2/p^2=cj^2/p(mod),$$
(0.79)
where $`c`$ is the inverse of $`qmodp`$.
In case $`p=0`$ the moduli space of flat reducible $`\text{SU}(2)`$โconnections on $`M_0`$ can be identified with a closed interval, cf. Proposition 3. Since the ChernโSimons functional is constant on each of the connected components of the moduli space of flat connections, we conclude that the ChernโSimons invariant of any of the reducible $`\text{SU}(2)`$โconnections is equal to the ChernโSimons invariant of the trivial connection, i.e. it is equal to zero. This of course also follows from Kirk and Klassenโs result. In fact, if $`\rho _\theta `$ denotes the abelian representation from Proposition 3, then we can put $`\alpha (t)=\theta t`$ and $`\beta (t)=0`$, $`t[0,1]`$, and get by (5.2) that
$$\text{CS}(\rho _\theta )=0(mod).$$
(0.80)
### 5.3 A comparison between ChernโSimons invariants and critical values of the phase functions $`\mathrm{\Phi }_n^{a,b}`$
The purpose of this section is to combine the ChernโSimons theory described in the previous two sections with the asymptotic analysis in Sec. 4.3. First we will describe a correspondence between the critical points of the phase functions $`\mathrm{\Phi }_n^{a,b}`$ in (0.36) and the nonabelian $`\text{SL}(2,)`$โrepresentations of $`\pi _1(M_{p/q})`$. Thereafter we will show that the set of ChernโSimons invariants of flat irreducible $`\text{SU}(2)`$โconnections on $`M_{p/q}`$ is a certain subset of the critical values of the functions $`\mathrm{\Phi }_n^{a,b}`$. Like in Sec. 4.3 we will most of the time work with the shifted phase functions $`\mathrm{\Psi }_n^{a,b}`$ in (0.39). We refer to Remark 5 below for a comparison between the phase functions $`\mathrm{\Phi }_n^{a,b}`$ and $`\mathrm{\Psi }_n^{a,b}`$. Recall the set $`\stackrel{~}{๐ฉ}`$ given by (0.65).
###### Theorem 7
The map $`(x,y)\rho _{(e^{\pi ix},e^{2\pi iy}1)}`$ gives a surjection $`\phi `$ from the set of critical points $`(x,y)`$ of the functions $`\mathrm{\Psi }_n^{a,b}`$, $`a,b,n`$, with $`x`$ onto the set of representations $`\rho _{(s,u)}:\pi \text{SL}(2,)`$, $`(s,u)\stackrel{~}{๐ฉ}`$, which extend to $`\text{SL}(2,)`$โrepresentations of $`\pi _1(M_{p/q})`$.
If $`(x,y)`$ is a critical point, let us say of $`\mathrm{\Psi }_n^{a,b}`$, then $`(x+2k,y+l)`$ is a critical point of $`\mathrm{\Psi }_{n+pk}^{a+l,b+2k}`$ for any $`k,l`$, so $`\phi ^1(\phi (x,y))=\{(x+2k,y+l)|k,l\}`$ if $`x`$.
* Proof. The last statement follows immediately from (4.3), so let us concentrate on the first part. Let $`\psi :^2`$ be given by
$$\psi (z,w)=(1zw)(wz)zw,$$
so $`\psi (v^2,w)=0`$ if and only if $`(v,w)`$ is a solution to (0.42). Then
$$\psi (z,u+1)=z\varphi (z,u)$$
for $`(z,u)^{}\times `$, where $`\varphi `$ is the function (0.64). It follows that $`(v,w)`$ is a solution to (0.42) if and only if either $`(v,w)=(0,0)`$ or $`(v,w1)\stackrel{~}{๐ฉ}`$.
Next assume that $`(v,u)\stackrel{~}{๐ฉ}`$ and put $`(z,w)=(v^2,u+1)`$. Since $`\psi (z,w)=0`$ and $`z0`$ we have that $`w0`$ and $`wz0`$. Therefore $`1zw=wz/(wz)`$ and $`(v,w)`$ is a solution to (0.41) if and only if
$`v^p`$ $`=`$ $`\left({\displaystyle \frac{(wz)^2}{wz}}\right)^q=\left(wz^12+zw^1\right)^q`$
$`=`$ $`\left(\lambda _{11}(v,w1)1z^2+zw^1+wz+z\right)^q.`$
But $`\psi (z,w)=0`$ and $`w0`$ implies that $`1z^2+zw^1+wz+z=0`$. Thus $`(v,w)`$ is a solution to (0.41) if and only if $`(s,u)=(v,u)`$ is a solution to (0.67).
The above shows together with the discussion around (0.67) that there is a one to one correspondence between common nonzero solutions $`(v,w)`$ to (0.41), (0.42) with $`v^21`$ and representations $`\rho _{(v,w1)}:\pi \text{SL}(2,)`$, $`(v,w1)\stackrel{~}{๐ฉ}`$, which extend to $`\text{SL}(2,)`$โrepresentations of $`\pi _1(M_{p/q})`$. By the remarks following (4.3) this proves the theorem.
Recall the set $`๐ฎ`$ in (0.5). Moreover, recall that if $`(x,y)`$ is a critical point of $`\mathrm{\Psi }_n^{a,b}`$, $`a,b,n`$, with $`(x,y)๐ฎ`$ then $`x`$ is automatically satisfied. By Proposition 2 and the above theorem we therefore have
###### Corollary 3
The surjection $`\phi `$ in Theorem 7 restricts to a surjection from the set of critical points $`(x,y)`$ of $`\mathrm{\Psi }_n^{a,b}`$, $`a,b,n`$, with $`(x,y)๐ฎ`$ onto the set of representations $`\rho _{(s,u)}:\pi _1(M_{p/q})\text{SL}(2,)`$, $`(s,u)\stackrel{~}{๐ฉ}`$, which are equivalent to $`\text{SU}(2)`$โrepresentations of $`\pi _1(M_{p/q})`$.
###### Theorem 8
Let $`(x,y)^2`$ such that $`\rho _{(e^{\pi ix},e^{2\pi iy}1)}`$ is equivalent to a nonabelian $`\text{SU}(2)`$โrepresentation of $`\pi _1(M_{p/q})`$ and choose in accordance with Theorem 7 integers $`a,b,n`$ such that $`(x,y)`$ is a critical point of $`\mathrm{\Psi }_n^{a,b}`$. If $`a^{},b^{},n^{}`$ is another such set of integers, then $`b^{}=b`$ and $`a^{}+n^{}/q=a+n/q`$. In fact we have
$`a+{\displaystyle \frac{n}{q}}`$ $`=`$ $`y+{\displaystyle \frac{p}{2q}}x+{\displaystyle \frac{i}{2\pi }}\left(\text{Log}\left(1e^{2\pi i(x+y)}\right)\text{Log}\left(1e^{2\pi i(xy)}\right)\right),`$
$`b`$ $`=`$ $`x+{\displaystyle \frac{i}{2\pi }}\left(\text{Log}\left(1e^{2\pi i(x+y)}\right)+\text{Log}\left(1e^{2\pi i(xy)}\right)\right).`$
Moreover, for any such set of integers $`a,b,n`$ we have
$$\text{CS}(\overline{\rho }_{\theta ,\epsilon })=\mathrm{\Psi }_n^{a,b}(x,y)(mod),$$
where $`\theta [1/3,1/6][1/6,1/3]`$ and $`\epsilon \{\pm \}`$ with $`e^{2\pi i\theta }=e^{\pi ix}`$ and $`1+u_\epsilon (\theta )=e^{2\pi iy}`$, and $`\overline{\rho }_{\theta ,\epsilon }`$ is a $`\text{SU}(2)`$โrepresentation of $`\pi _1(M_{p/q})`$ equivalent to $`\rho _{(e^{\pi ix},e^{2\pi iy}1)}`$ and given by Proposition 2.
###### Remark 5
i) In general we have $`\mathrm{\Phi }_n^{a,b}=\mathrm{\Psi }_n^{a+b,ab}`$, so there is a ono to one correspondence between the phase functions $`\mathrm{\Phi }_n^{a,b}`$ and the phase functions $`\mathrm{\Psi }_n^{a^{},b^{}}`$ with $`a^{}+b^{}`$ even. If $`(x,y)`$ is a critical point of $`\mathrm{\Psi }_n^{a^{},b^{}}`$ with $`a^{}+b^{}`$ odd, then $`(x,y)`$ is also a critical point of $`\mathrm{\Psi }_{nq}^{a^{}+1,b^{}}`$ by Theorem 8. In that way we see that any critical point of one of the functions $`\mathrm{\Psi }_n^{a^{},b^{}}`$ is also a critical point of one of the functions $`\mathrm{\Phi }_n^{a,b}`$.
ii) Assume that $`|p/q|>\sqrt{20}`$ and assume that $`(x_j,y_j)`$, $`j=1,2`$, are degenerate critical points of phase functions $`\mathrm{\Phi }_n^{a,b}`$ belonging to $`๐ฎ`$. Let $`v_j=e^{\pi ix_j}`$ and $`w_j=e^{2\pi iy_j}`$. By Proposition 1 we have $`\mathrm{cos}(2\pi x_1)=\mathrm{cos}(2\pi x_2)`$ and $`w_j=w_\epsilon (x_j)`$, where $`\epsilon =\text{sign}(p/q)`$, so $`w_1=w_2`$. Moreover, $`x_1,x_2[2/3,1/3][1/3,2/3]+2`$. We are interested in comparing the equivalence classes of the representations $`\rho _j=\rho _{(v_j,w_j1)}`$. For that purpose we can without loss of generality assume that $`x_1,x_2[1/3,2/3]`$. Then $`x_2=x_1`$ or $`x_2=1x_1`$. In the second case $`v_2=v_1`$. Note here that $`v_1=v_1^1`$ if and only if $`v_1\{\pm i\}`$. But in that case $`\mathrm{cos}(2\pi x_1)=1`$ and by (0.53) we get that $`p/q\{\pm 2/\sqrt{5}\}`$ which is impossible. Thus $`\rho _1`$ and $`\rho _2`$ are not equivalent. If $`p`$ is odd we must have $`x_2=x_1`$ since only one of the points $`(v_1,w_1)`$, $`(v_1,w_1)`$ can be a solution to (0.41). If $`p`$ is even both $`(v_1,w_1)`$ and $`(v_1,w_1)`$ are solutions to both of the equations (0.41) and (0.42). Thus Proposition 1 shows that the set of degenerate critical points in $`๐ฎ`$ of the phase functions $`\mathrm{\Phi }_n^{a,b}`$, $`a,b,n`$, is either empty or corresponds under the surjection in Theorem 7 to one point in $`_{p/q}`$ in case $`p`$ is odd and to two points in the moduli space in case $`p`$ is even.
The phase functions $`\mathrm{\Psi }_n^{0,0}`$, $`\mathrm{\Psi }_n^{1,1}`$, $`\mathrm{\Psi }_n^{1,1}`$ and $`\mathrm{\Psi }_n^{2,0}`$ are the ones entering Conjecture 3. By Theorem 8 the parameter $`b`$ is fixed by a critical point of $`\mathrm{\Psi }_n^{a,b}`$, while $`a`$ and $`n`$ can be varied as long as we keep fixed $`a+n/q`$. Because of this (see also Lemma 5 below) we only need to consider $`\mathrm{\Psi }_n^{0,0}`$, $`\mathrm{\Psi }_n^{0,1}`$, and $`\mathrm{\Psi }_n^{0,1}`$. By Theorem 8 (and also the proof of this theorem and Proposition 2) we get
###### Corollary 4
Let $`(x,y)^2`$ such that $`\rho _{(e^{\pi ix},e^{2\pi iy}1)}`$ is equivalent to a nonabelian $`\text{SU}(2)`$โrepresentation of $`\pi _1(M_{p/q})`$. Let $`(x^{},y^{})^2`$ such that $`(e^{\pi ix^{}},e^{2\pi iy^{}})=(e^{\pi ix},e^{2\pi iy})`$ and $`\text{Re}(x^{})]1,1]`$ and $`\text{Re}(y^{})]1/2,1/2]`$. Moreover, let
$`n`$ $`=`$ $`q\left(y^{}+{\displaystyle \frac{p}{2q}}x^{}+{\displaystyle \frac{i}{2\pi }}\left(\text{Log}\left(1e^{2\pi i(x+y)}\right)\text{Log}\left(1e^{2\pi i(xy)}\right)\right)\right),`$
$`b`$ $`=`$ $`x^{}+{\displaystyle \frac{i}{2\pi }}\left(\text{Log}\left(1e^{2\pi i(x+y)}\right)+\text{Log}\left(1e^{2\pi i(xy)}\right)\right).`$
Then $`\theta :=x^{}/2[1/3,1/6][1/6,1/3]`$ and
$$b=\{\begin{array}{cc}\hfill 1,& \theta [1/3,1/4]\hfill \\ \hfill 0,& \theta ]1/4,1/6][1/6,1/4]\hfill \\ \hfill 1,& \theta ]1/4,1/3].\hfill \end{array}$$
Moreover, $`n`$ and $`(x^{},y^{})`$ is a critical point of $`\mathrm{\Psi }_n^{0,b}`$ and
$$\text{CS}(\overline{\rho }_{\theta ,\epsilon })=\mathrm{\Psi }_n^{0,b}(x^{},y^{})(mod),$$
where $`\epsilon \{\pm \}`$ with $`1+u_\epsilon (\theta )=e^{2\pi iy}`$, and $`\overline{\rho }_{\theta ,\epsilon }`$ is a $`\text{SU}(2)`$โrepresentation of $`\pi _1(M_{p/q})`$ equivalent to $`\rho _{(e^{\pi ix},e^{2\pi iy}1)}`$ and given by Proposition 2.
We note that the first part of Theorem 8 is an immediate consequence of (4.3). To prove the second part we start by observing that if $`a,a^{},n,n^{},b`$ are integers such that $`a^{}+n^{}/q=a+n/q`$ then $`\mathrm{\Psi }_n^{a,b}(x,y)\mathrm{\Psi }_n^{}^{a^{},b}(x,y)`$ is an integer independent of $`(x,y)^2`$. The remaining part of the proof will consists of a series of lemmas. We start by
###### Lemma 5
Let $`a,b,n`$ and assume that $`(x,y)`$ is a critical point of $`\mathrm{\Psi }_n^{a,b}`$. Let $`l,k`$ and put $`a^{}=a+l`$, $`b^{}=b+2k`$, and $`n^{}=n+pk`$. Then $`(x+2k,y+l)`$ is a critical point of $`\mathrm{\Psi }_n^{}^{a^{},b^{}}`$ and
$$\mathrm{\Psi }_n^{}^{a^{},b^{}}(x+2k,y+l)\mathrm{\Psi }_n^{a,b}(x,y)0(mod).$$
* Proof. That $`(x+2k,y+l)`$ is a critical point of $`\mathrm{\Psi }_n^{}^{a^{},b^{}}`$ is an immediate consequence of (4.3) and was already observed in Theorem 7. Put $`x^{}=x+2k`$ and $`y^{}=y+l`$. Since $`e^{2\pi ix^{}}=e^{2\pi ix}`$ and $`e^{2\pi iy^{}}=e^{2\pi iy}`$ we get
$`\mathrm{\Psi }_n^{}^{a^{},b^{}}(x^{},y^{})\mathrm{\Psi }_n^{a,b}(x,y)`$ $`=`$ $`a^{}x^{}+b^{}y^{}x^{}y^{}+{\displaystyle \frac{n^{}}{q}}x^{}{\displaystyle \frac{p}{4q}}x_{}^{}{}_{}{}^{2}{\displaystyle \frac{d}{q}}n_{}^{}{}_{}{}^{2}`$
$`axby+xy{\displaystyle \frac{n}{q}}x+{\displaystyle \frac{p}{4q}}x^2+{\displaystyle \frac{d}{q}}n^2`$
$`=`$ $`2(a+l)k+2{\displaystyle \frac{n}{q}}(1pd)k+{\displaystyle \frac{p}{q}}(1pd)k^2.`$
Here $`pdcq=1`$ for an integer $`c`$ so
$$\mathrm{\Psi }_n^{}^{a^{},b^{}}(x^{},y^{})\mathrm{\Psi }_n^{a,b}(x,y)=2(a+lnc)kpck^2.$$
By Proposition 2 and the above lemma we are thus left with the following to prove: Let $`\theta [1/3,1/6][1/6,1/3]`$ and assume that $`\overline{\rho }_{\theta ,\epsilon }`$ extends to a representation of $`\pi _1(M_{p/q})`$ for $`\epsilon \{\pm \}`$. Let $`b=b(\theta )`$ and $`n=n(\theta )`$ be the unique integers such that $`(2\theta ,y)`$ is a stationary point of $`\mathrm{\Psi }_n^{0,b}`$, where $`y=\frac{1}{2\pi i}\text{Log}(1+u_\epsilon (\theta ))`$ (so $`\text{Re}(y)=\frac{1}{2}`$) (see around (0.43) and (0.44)). Then
$$\text{CS}(\overline{\rho }_{\theta ,\rho })=\mathrm{\Psi }_n^{0,b}(2\theta ,y)(mod).$$
(0.81)
In the course of the proof of this identity we will prove that
$$\text{CS}(\overline{\rho }_{\theta ,\epsilon })=\frac{1}{6}\frac{p}{q}\theta ^2+\frac{m}{q}2\theta \frac{d}{q}m^22_{1/6}^\theta \beta _\epsilon (t)\text{d}t(mod),$$
(0.82)
for $`\theta [1/6,1/3]`$, where $`m=p\theta +q\beta _\epsilon (\theta )`$ by (0.77). Unfortunately the proof of (0.81) is rather technical, but we have tried to emphasize the main ideas of the proof here and defer many technicalities to the appendices C and D.
The ChernโSimons invariants of $`\text{SU}(2)`$โconnections are real, so we begin by investigating to what extent a general phase function $`\mathrm{\Psi }_n^{a,b}`$, $`a,b,n`$, is real in its critical points. Assume that $`(x,y)`$ is such a critical point and write $`z=e^{2\pi ix}`$ and $`w=e^{2\pi iy}`$ as usual. From (4.3) we find that
$$a\text{Re}(y)=\frac{p}{2q}\text{Re}(x)\frac{n}{q}\frac{1}{2\pi }\text{Im}\left(\text{Log}(1zw)\text{Log}(1zw^1)\right),$$
and
$$b\text{Re}(x)=\frac{1}{2\pi }\text{Im}\left(\text{Log}(1zw)+\text{Log}(1zw^1)\right).$$
Moreover, we have $`\text{Im}(x)=\frac{1}{2\pi }\text{Log}|z|`$ and $`\text{Im}(y)=\frac{1}{2\pi }\text{Log}|w|`$. By (0.39) we therefore get
$`\text{Im}\left(\mathrm{\Psi }_n^{a,b}(x,y)\right)`$ $`=`$ $`(a\text{Re}(y))\text{Im}(x)+(b\text{Re}(x))\text{Im}(y){\displaystyle \frac{p}{2q}}\text{Im}(x)\text{Re}(x)`$
$`+{\displaystyle \frac{n}{q}}\text{Im}(x)+{\displaystyle \frac{1}{4\pi ^2}}\text{Im}\left(\text{Li}_2(zw)\text{Li}_2(zw^1)\right).`$
This together with the above expressions for $`a\text{Re}(y)`$ and $`b\text{Re}(x)`$ leads to the formula
$`\text{Im}\left(\mathrm{\Psi }_n^{a,b}(x,y)\right)`$ $`=`$ $`{\displaystyle \frac{1}{4\pi ^2}}(\text{Im}(\text{Log}(1zw)\text{Log}(1zw^1))\text{Log}|z|`$
$`+\text{Im}(\text{Log}(1zw)+\text{Log}(1zw^1))\text{Log}|w|`$
$`+\text{Im}(\text{Li}_2(zw)\text{Li}_2(zw^1))).`$
By introducing the BlochโWigner dilogarithm function
$$D(z)=\text{Im}\left(\text{Li}_2(z)\right)+\text{Arg}(1z)\text{Log}|z|$$
we obtain
###### Lemma 6
Let $`a,b,n`$ and let $`(x,y)`$ be a critical point of $`\mathrm{\Psi }_n^{a,b}`$. Then
$$\text{Im}\left(\mathrm{\Psi }_n^{a,b}(x,y)\right)=\frac{1}{4\pi ^2}\left(D\left(e^{2\pi i(x+y)}\right)D\left(e^{2\pi i(xy)}\right)\right).$$
We note that $`D`$ is analytic on $`\{0,1\}`$ and continuous on $``$. Moreover, $`D`$ satisfies the identities
$$D(z)+D(\overline{z})=0,D(z)+D(1/z)=0$$
(0.83)
for all $`z\{0\}`$. From this we have
###### Corollary 5
Let $`a,b,n`$ and let $`(x,y)`$ be a critical point of $`\mathrm{\Psi }_n^{a,b}`$ with $`e^{2\pi ix}S^1`$ and $`e^{2\pi iy}`$. Then
$$\text{Im}\left(\mathrm{\Psi }_n^{a,b}(x,y)\right)=0.$$
###### Remark 6
Let $`(v,w)`$ be a non-zero solution to (0.41) and (0.42). Then $`(\overline{v},\overline{w})`$ is also a solution to these two equations as already observed. Let $`z=v^2`$ and write $`(z,w)=(e^{2\pi ix},e^{2\pi iy})`$. Then $`(\overline{z},\overline{w})=(e^{2\pi i\overline{x}},e^{2\pi i\overline{y}})`$. If $`(x,y)`$ is a critical point of $`\mathrm{\Psi }_n^{a,b}`$, then $`(\overline{x},\overline{y})`$ is a critical point of $`\mathrm{\Psi }_n^{a,b}`$. By Lemma 6 and (0.83) we have
$$\text{Im}\left(\mathrm{\Psi }_n^{a,b}(x,y)\right)=\text{Im}\left(\mathrm{\Psi }_n^{a,b}(\overline{x},\overline{y})\right).$$
Hence, if this value is different from zero, then either $`\mathrm{exp}(2\pi ir\mathrm{\Psi }_n^{a,b}(\overline{x},\overline{y}))`$ or $`\mathrm{exp}(2\pi ir\mathrm{\Psi }_n^{a,b}(x,y))`$ grows exponentially. We note that by Conjecture 3 and Corollary 5 stationary points leading to such exponential growth do not contribute to the large $`r`$ asymptotics of $`\overline{\tau }_r(M_{p/q})`$, see also Theorem 4.
We now embark upon proving (0.81). We start by reducing to the case $`\theta [1/6,1/3]`$. By Proposition 2 we know that the representations $`\overline{\rho }_{\theta ,\epsilon }`$ and $`\overline{\rho }_{\theta ,\epsilon }`$ are $`\text{SU}(2)`$โequivalent, so in particular they have the same ChernโSimons invariant.
###### Lemma 7
Let the situation be as in (0.81). Then
$$\mathrm{\Psi }_{n(\theta )}^{0,b(\theta )}(2\theta ,y)\mathrm{\Psi }_{n(\theta )}^{0,b(\theta )}(2\theta ,y)0(mod).$$
To prove this and (0.81) we need the technical Lemma 8 below keeping track of branches of the logarithm in certain expressions. Let $`I=[1/3,1/6][1/6,1/3]`$ and put
$`Q_1^\pm (\theta )`$ $`=`$ $`1\left(1+u_\pm (\theta )\right)^1e^{4\pi i\theta },`$
$`Q_2^\pm (\theta )`$ $`=`$ $`1\left(1+u_\pm (\theta )\right)e^{4\pi i\theta }`$ (0.84)
for $`\theta I`$, where $`u_\pm `$ are defined in Proposition 2. We also put $`Q_3^\pm (\theta )=1+u_\pm (\theta )`$.
###### Lemma 8
We have
$$\text{Log}\left(Q_1^\pm (\theta )\right)+\text{Log}\left(Q_2^\pm (\theta )\right)=\text{Log}\left(Q_1^\pm (\theta )Q_2^\pm (\theta )\right)$$
(0.85)
and
$$Q_1^\pm (\theta )Q_2^\pm (\theta )=e^{4\pi i\theta }$$
(0.86)
for all $`\theta I`$. Moreover
$$\text{Log}\left(Q_1^\pm (\theta )\right)+\text{Log}\left(Q_3^\pm (\theta )\right)\text{Log}\left(Q_2^\pm (\theta )\right)=\text{Log}\left(\frac{Q_1^\pm (\theta )Q_3^\pm (\theta )}{Q_2^\pm (\theta )}\right)+e_\pm (\theta )2\pi i$$
(0.87)
for all $`\theta I`$, where
$$e_+(\theta )=\{\begin{array}{cc}1,& \theta ]1/6,1/4],\\ 0,& \theta ]1/4,1/3[,\end{array}$$
and $`e_{}(\theta )=1e_+(\theta )`$ for $`\theta ]1/6,1/3[`$, $`e_\pm (1/4)=e_\pm (1/4)`$, $`e_\pm (\theta )=1e_\pm (\theta )`$ for $`\theta I\{\pm 1/6,\pm 1/4,\pm 1/3\}`$, and $`e_\pm (\theta )=0`$ for $`\theta \{\pm 1/6,\pm 1/3\}`$. Finally
$$\frac{Q_1^\pm (\theta )Q_3^\pm (\theta )}{Q_2^\pm (\theta )}=L_\pm (\theta )$$
(0.88)
for all $`\theta I`$, where $`L_\pm `$ are given by Proposition 2.
The proof of this lemma is given in Appendix D.
* Proof of Lemma 7 By (0.43) and (0.44) we have
$$b(\theta )=2\theta \frac{1}{2\pi i}\left(\text{Log}\left(Q_1^\epsilon (\theta )\right)+\text{Log}\left(Q_2^\epsilon (\theta )\right)\right)$$
and
$$n(\theta )=p\theta +qy+q\frac{1}{2\pi i}\left(\text{Log}\left(Q_1^\epsilon (\theta )\right)\text{Log}\left(Q_2^\epsilon (\theta )\right)\right).$$
By Lemma 8
$$b(\theta )=2\theta \frac{1}{2\pi i}\text{Log}\left(e^{4\pi i\theta }\right).$$
(0.89)
Taking the real part of the expression for $`n(\theta )`$ we get
$$n(\theta )=p\theta +\frac{1}{2}q+\frac{q}{2\pi }\left(\text{Arg}\left(Q_1^\epsilon (\theta )\right)\text{Arg}\left(Q_2^\epsilon (\theta )\right)\right).$$
(0.90)
By (0.89) we get that
$$b(\theta )=b(\theta ),\theta [1/3,1/6][1/6,1/3]\{\pm 1/4\},$$
and
$$b(1/4)=1,b(1/4)=0.$$
By (0.90), (D.5), and (D.6) we get
$$n(\theta )=qn(\theta ),\theta [1/3,1/6][1/6,1/3]\{\pm 1/4\},$$
and
$$n(1/4)=n(1/4)\frac{p}{2}.$$
Let $`b=b(\theta )`$ and $`n=n(\theta )`$. Since $`\mathrm{\Psi }_{n(\theta )}^{0,b(\theta )}(2\theta ,y)`$ is real by Corollary 5 and since $`\text{Re}\left(\text{Li}_2(\overline{z})\right)=\text{Re}\left(\text{Li}_2(z)\right)`$ for $`z]1,\mathrm{}[`$ we find that
$`\mathrm{\Psi }_{n(\theta )}^{0,b(\theta )}(2\theta ,y)\mathrm{\Psi }_{n(\theta )}^{0,b(\theta )}(2\theta ,y)`$ $`=`$ $`(bb(\theta ))\text{Re}(y)4\theta \text{Re}(y)`$
$`+{\displaystyle \frac{(n+n(\theta )}{q}}2\theta {\displaystyle \frac{d(n^2n(\theta )^2)}{q}}.`$
(If $`\theta \{\pm 1/4\}`$, use that $`\text{Li}_2(e^{2\pi i(2\theta +y)})\text{Li}_2(e^{2\pi i(2\theta y)})`$ is the same for $`\theta =1/4`$ and $`\theta =1/4`$.) Assume first that $`\theta \pm 1/4`$. Then, since $`\text{Re}(y)=1/2`$, we get
$$\mathrm{\Psi }_{n(\theta )}^{0,b(\theta )}(2\theta ,y)\mathrm{\Psi }_{n(\theta )}^{0,b(\theta )}(2\theta ,y)=b+dq2dn.$$
Next assume that $`\theta =1/4`$. Then one finds
$$\mathrm{\Psi }_{n(\theta )}^{0,b(\theta )}(2\theta ,y)\mathrm{\Psi }_{n(\theta )}^{0,b(\theta )}(2\theta ,y)=\frac{n}{q}(1pd)\frac{p}{4q}(1pd).$$
But $`1pd=qc`$ for an integer $`c`$ so
$$\mathrm{\Psi }_{n(\theta )}^{0,b(\theta )}(2\theta ,y)\mathrm{\Psi }_{n(\theta )}^{0,b(\theta )}(2\theta ,y)=\left(\frac{p}{4}n\right)c.$$
By (0.77) $`p/4+q\beta _\epsilon (1/4)`$. But $`L_\pm (1/4)=1`$ so $`\beta _\epsilon (1/4)=f_\epsilon (1/4)`$, so $`p`$ is divisible by $`4`$.
By using that $`y=\frac{1}{2\pi i}\text{Log}\left(Q_3^\epsilon (\theta )\right)`$ together with Lemma 8 we obtain the alternative formula
$$n(\theta )=p\theta +qe_\epsilon (\theta )+\frac{q}{2\pi i}\text{Log}\left(L_\epsilon (\theta )\right).$$
This and (0.74) immediately leads to
$$n(\theta )=p\theta +q\beta _\epsilon (\theta )+q(e_\epsilon (\theta )f_\epsilon (\theta ))$$
(0.91)
for $`\theta [1/6,1/3]`$. In the proof of (0.81) we need certain symmetries for the functions $`L_\pm `$. First we note that
$$L_\pm \left(\frac{1}{2}\theta \right)=L_{}(\theta )$$
(0.92)
for $`\theta [1/6,1/3]`$ by (5.2). Second, by the next lemma and Proposition 2, we have
$$L_{}(\theta )=L_\pm (\theta )^1$$
(0.93)
for $`\theta [1/3,1/6][1/6,1/3]`$. (In particular, $`\overline{\rho }_{\theta ,\pm }`$ extends to a representation of $`\pi _1(M_{p/q})`$ if and only if $`\overline{\rho }_{\theta ,}`$ extends to a representation of $`\pi _1(M_{p/q})`$.)
###### Lemma 9
Let $`s^{}`$ and let $`u_\pm `$ be the two solutions to $`\varphi (s^2,u)=0`$. Then
$$(1+u_+)(1+u_{})=1,$$
and
$$\lambda _{11}(s,u_+)\lambda _{11}(s,u_{})=1.$$
* Proof. We have
$$1+u_\pm =\frac{1}{2}(s^2+s^21)\pm \frac{1}{2}\sqrt{s^4+s^42(s^2+s^2)1}$$
so
$$(1+u_+)(1+u_{})=\frac{1}{4}(s^2+s^21)^2\frac{1}{4}(s^4+s^42(s^2+s^2)1)=1.$$
Moreover,
$`\lambda _{11}(s,u_\pm )`$ $`=`$ $`{\displaystyle \frac{1}{2}}(s^4+s^4){\displaystyle \frac{1}{2}}(s^2+s^2)1`$
$`\pm {\displaystyle \frac{1}{2}}(s^2s^2)\sqrt{s^4+s^42(s^2+s^2)1}`$
so
$`\lambda _{11}(s,u_+)\lambda _{11}(s,u_{})`$ $`=`$ $`\left({\displaystyle \frac{1}{2}}(s^4+s^4){\displaystyle \frac{1}{2}}(s^2+s^2)1\right)^2`$
$`{\displaystyle \frac{1}{4}}(s^2s^2)^2\left(s^4+s^42(s^2+s^2)1\right).`$
By simple reductions one gets the result.
We are now ready to finalize the proof of (0.81) and thereby the proof of Theorem 8.
* Proof of (0.81) By Lemma 7 we can assume that $`\theta [1/6,1/3]`$. Write $`\theta _0`$ for $`\theta `$ in the following. Let us first observe that formula (0.82) is an immediate consequence of Proposition 4. In fact, by letting $`c,d`$ be integers as in Proposition 4 we get
$`2cq\theta _0\beta _\epsilon (\theta _0)=2cp\theta _0^22cm\theta _0,`$
$`dq\beta _\epsilon ^2(\theta _0)={\displaystyle \frac{d}{q}}\left(m^22pm\theta _0+p^2\theta _0^2\right),`$
and therefore
$$cp\theta _0^2dq\beta _\epsilon ^2(\theta _0)2cq\theta _0\beta _\epsilon (\theta _0)=\left(c\frac{dp}{q}\right)p\theta _0^22\left(c\frac{dp}{q}\right)m\theta _0\frac{d}{q}m^2.$$
But $`cdp/q=1/q`$, hence (0.82) follows. On the other hand we have
$$\mathrm{\Psi }_n^{0,b}(2\theta _0,y)=(b2\theta _0)y+\frac{n}{q}2\theta _0\frac{p}{q}\theta _0^2\frac{d}{q}n^2+\frac{1}{4\pi ^2}\left(\text{Li}_2(z_0w_\epsilon )\text{Li}_2(z_0w_\epsilon ^1)\right),$$
where $`z_0=e^{4\pi i\theta _0}`$ and $`w_\epsilon =1+u_\epsilon (\theta _0)`$. By (0.89) and (0.91) we have $`b2\theta _0=\frac{1}{2\pi i}\text{Log}(z_0)`$ and $`n=m+q\left(e_\epsilon (\theta _0)f_\epsilon (\theta _0)\right)`$ so
$`\mathrm{\Psi }_n^{0,b}(2\theta _0,y)`$ $`=`$ $`{\displaystyle \frac{p}{q}}\theta _0^2+{\displaystyle \frac{m}{q}}2\theta _0{\displaystyle \frac{d}{q}}m^2+2\left(e_\epsilon (\theta _0)f_\epsilon (\theta _0)\right)\theta _0l_\epsilon (\theta _0)`$
$`+{\displaystyle \frac{1}{4\pi ^2}}\left(\text{Log}(z_0)\text{Log}(w_\epsilon )+\text{Li}_2(z_0w_\epsilon )\text{Li}_2(z_0w_\epsilon ^1)\right),`$
where $`l_\epsilon (\theta _0)=2dm\left(e_\epsilon (\theta _0)f_\epsilon (\theta _0)\right)+dq\left(e_\epsilon (\theta _0)f_\epsilon (\theta _0)\right)^2`$. We therefore get that
$`\mathrm{\Psi }_n^{0,b}(2\theta _0,y)\text{CS}(\overline{\rho }_{\theta _0,\epsilon })`$ $`=`$ $`{\displaystyle \frac{1}{6}}+2\left(e_\epsilon (\theta _0)f_\epsilon (\theta _0)\right)\theta _0`$
$`+{\displaystyle \frac{1}{4\pi ^2}}\left(\text{Log}(z_0)\text{Log}(w_\epsilon )+\text{Li}_2(z_0w_\epsilon )\text{Li}_2(z_0w_\epsilon ^1)\right)`$
$`+2{\displaystyle _{1/6}^{\theta _0}}\beta _\epsilon (t)\text{d}t(mod).`$
For $`\theta _0[1/6,1/4]`$ we note that $`e_\epsilon (\theta _0)=f_\epsilon (\theta _0)`$. We will consider the special cases $`\theta _0\{1/6,1/4,1/3\}`$ first and then handle the other cases afterwards.
The cases $`\theta _0\{1/6,1/3\}`$. In these cases we have $`w_\epsilon =1`$ so
$`\mathrm{\Psi }_n^{0,b}(2\theta _0,y)\text{CS}(\overline{\rho }_{\theta _0,\epsilon })`$
$`={\displaystyle \frac{1}{6}}+{\displaystyle \frac{i}{4\pi }}\text{Log}(z_0)+2\left(e_\epsilon (\theta _0)f_\epsilon (\theta _0)\right)\theta _0+2{\displaystyle _{1/6}^{\theta _0}}\beta _\epsilon (t)\text{d}t(mod).`$
If $`\theta _0=1/6`$ we immmediately get that this is zero. If $`\theta _0=1/3`$ we get
$`\mathrm{\Psi }_n^{0,b}(2\theta _0,y)\text{CS}(\overline{\rho }_{\theta _0,\epsilon })`$ $`=`$ $`{\displaystyle \frac{1}{6}}+{\displaystyle \frac{i}{4\pi }}(4\pi i/32\pi i)+{\displaystyle \frac{2}{3}}\left(e_\epsilon (\theta _0)f_\epsilon (\theta _0)\right)`$
$`+2f_\epsilon (1/4)(1/31/6)`$
$`+{\displaystyle \frac{1}{\pi i}}{\displaystyle _{1/6}^{1/3}}\text{Log}(L_\epsilon (t))\text{d}t(mod).`$
But if $`\theta _0[1/4,1/3]`$ then
$$_{1/4}^{\theta _0}\text{Log}(L_\pm (\theta )\text{d}\theta =_{1/4}^{1/2\theta _0}\text{Log}(L_{}(t))\text{d}t=_{1/2\theta _0}^{1/4}\text{Log}(L_\pm (t))\text{d}t,$$
by (0.92) and (0.93) so
$$_{1/2\theta _0}^{\theta _0}\text{Log}(L_\pm (\theta )\text{d}\theta =0.$$
(0.94)
In particular,
$$\mathrm{\Psi }_n^{0,b}(2\theta _0,y)\text{CS}(\overline{\rho }_{\theta _0,\epsilon })=\frac{1}{3}\frac{2}{3}f_\epsilon (1/3)+\frac{1}{3}f_\epsilon (1/4)(mod),$$
where we also use that $`e_\pm (1/3)=0`$ by Lemma 8. By (0.75) and (0.76) this is zero.
The case $`\theta _0=1/4`$. In this case we have $`z_0=1`$ so
$`\mathrm{\Psi }_n^{0,b}(2\theta _0,y)\text{CS}(\overline{\rho }_{\theta _0,\epsilon })`$ $`=`$ $`{\displaystyle \frac{1}{6}}+{\displaystyle \frac{1}{4\pi ^2}}\left(i\pi \text{Log}(w_\epsilon )+\text{Li}_2(|w_\epsilon |)\text{Li}_2(|w_\epsilon |^1)\right)`$
$`+2{\displaystyle _{1/6}^{1/4}}\beta _\epsilon (t)\text{d}t(mod).`$
Since $`(v_0,w_\epsilon )=(e^{2\pi i\theta _0},w_\epsilon )`$ is a solution to (0.42) we have $`w_\epsilon \{(3\sqrt{5})/2,(3+\sqrt{5})/2\}`$, and by (5.3) we then conclude that $`w_\epsilon =\frac{3\sqrt{5}}{2}`$ if $`\epsilon =`$ and $`w_\epsilon =\frac{3+\sqrt{5}}{2}`$ if $`\epsilon =+`$. By (C.11) we then get
$`\mathrm{\Psi }_n^{0,b}(2\theta _0,y)\text{CS}(\overline{\rho }_{\theta _0,\epsilon })`$ $`=`$ $`{\displaystyle \frac{1}{6}}+{\displaystyle \frac{1}{4\pi ^2}}\left(\pi ^2\epsilon {\displaystyle \frac{\pi ^2}{5}}\right)+2{\displaystyle _{1/6}^{1/4}}\beta _\epsilon (t)\text{d}t`$
$`=`$ $`{\displaystyle \frac{1}{6}}+{\displaystyle \frac{1}{4}}\left(1\epsilon {\displaystyle \frac{1}{5}}\right)+{\displaystyle \frac{1}{6}}f_\epsilon (1/4)`$
$`+{\displaystyle \frac{1}{\pi }}{\displaystyle _{1/6}^{1/4}}\text{Arg}(L_\epsilon (t))\text{d}t(mod),`$
and this is zero by (0.75), (0.76), and (C.5).
The case $`\theta ]1/6,1/3[\{1/4\}`$. We have
$`\mathrm{\Psi }_n^{0,b}(2\theta _0,y)\text{CS}(\overline{\rho }_{\theta _0,\epsilon })`$ $`=`$ $`{\displaystyle \frac{1}{6}}+{\displaystyle \frac{1}{4\pi ^2}}\text{Log}(z_0)\text{Log}(w_\epsilon )+2\left(e_\epsilon (\theta _0)f_\epsilon (\theta _0)\right)\theta _0`$
$`+2{\displaystyle _{1/6}^{\theta _0}}\beta _\epsilon (t)\text{d}t+R(2\theta _0,y)(mod),`$
where
$$R(x,y)=\frac{1}{4\pi ^2}\left(\text{Li}_2(e^{2\pi i(x+y)})\text{Li}_2(e^{2\pi i(xy)})\right).$$
Let us write $`u`$ for $`u_\epsilon `$ in the following. By definition of the dilogarithm we have
$$4\pi ^2R(2\theta _0,y)=_0^{(1+u)^1e^{4\pi i\theta _0}}\frac{\text{Log}(1t)}{t}\text{d}t_0^{(1+u)e^{4\pi i\theta _0}}\frac{\text{Log}(1t)}{t}\text{d}t.$$
By (5.3) we have $`1+u_{}(\theta )11+u_+(\theta )<0`$ for all $`\theta [1/6,1/3]`$ and $`1+u_\pm (\theta )=1`$ if and only if $`\theta \{1/6,1/3\}`$. Let $`\theta _1=1/6`$ if $`\theta _0[1/6,1/4[`$ and $`\theta _1=1/3`$ if $`\theta _0]1/4,1/3[`$. We note that $`z\text{Log}(1z)/z`$ is analytic on $`[1,\mathrm{}[`$ so by Cauchyโs theorem
$`4\pi ^2R(2\theta _0,y)`$ $`=`$ $`{\displaystyle _0^{e^{4\pi i\theta _1}}}{\displaystyle \frac{\text{Log}(1t)}{t}}\text{d}t+{\displaystyle _{e^{4\pi i\theta _1}}^{(1+u)^1e^{4\pi i\theta _0}}}{\displaystyle \frac{\text{Log}(1t)}{t}}\text{d}t`$
$`{\displaystyle _0^{e^{4\pi i\theta _1}}}{\displaystyle \frac{\text{Log}(1t)}{t}}\text{d}t{\displaystyle _{e^{4\pi i\theta _1}}^{(1+u)e^{4\pi i\theta _0}}}{\displaystyle \frac{\text{Log}(1t)}{t}}\text{d}t`$
$`=`$ $`{\displaystyle _{e^{4\pi i\theta _1}}^{(1+u)^1e^{4\pi i\theta _0}}}{\displaystyle \frac{\text{Log}(1t)}{t}}\text{d}t{\displaystyle _{e^{4\pi i\theta _1}}^{(1+u)e^{4\pi i\theta _0}}}{\displaystyle \frac{\text{Log}(1t)}{t}}\text{d}t.`$
The curves $`\gamma _\pm (\theta )=(1+u(\theta ))^{\pm 1}e^{4\pi i\theta }`$ are smooth on $`]1/6,1/3[`$ so
$`4\pi ^2R(2\theta _0,y)`$ $`=`$ $`\underset{\eta 0_+}{lim}({\displaystyle _{\theta _1+\mu \eta }^{\theta _0}}\text{Log}(1\gamma _{}(\theta )){\displaystyle \frac{\gamma _{}^{}(\theta )}{\gamma _{}(\theta )}}\text{d}\theta `$
$`{\displaystyle _{\theta _1+\mu \eta }^{\theta _0}}\text{Log}(1\gamma _+(\theta )){\displaystyle \frac{\gamma _+^{}(\theta )}{\gamma _+(\theta )}}\text{d}\theta ),`$
where $`\mu =1`$ if $`\theta _1=1/6`$ and $`\mu =1`$ if $`\theta _1=1/3`$. (The parameter $`\eta `$ is necessary because $`u`$ is not differentiable in $`1/6`$ and $`1/3`$.) It follows that
$$4\pi ^2R(2\theta _0,y)=\underset{\eta 0_+}{lim}\left(4\pi iR_1(\theta _0,\eta )R_2(\theta _0,\eta )\right),$$
where
$`R_1(\theta _0,\eta )`$ $`=`$ $`{\displaystyle _{\theta _1+\mu \eta }^{\theta _0}}\text{Log}\left(Q_1^\epsilon (\theta )\right)\text{Log}\left(Q_2^\epsilon (\theta )\right)\text{d}\theta ,`$
$`R_2(\theta _0,\eta )`$ $`=`$ $`{\displaystyle _{\theta _1+\mu \eta }^{\theta _0}}\left\{\text{Log}\left(Q_1^\epsilon (\theta )\right)+\text{Log}\left(Q_2^\epsilon (\theta )\right)\right\}{\displaystyle \frac{u^{}(\theta )}{1+u(\theta )}}\text{d}\theta ,`$
where the functions $`Q_i^\pm `$ are defined above Lemma 8. By Lemma 8
$`R_2(\theta _0,\eta )`$ $`=`$ $`{\displaystyle _{\theta _1+\mu \eta }^{\theta _0}}\text{Log}\left(e^{4\pi i\theta }\right){\displaystyle \frac{u^{}}{1+u}}\text{d}\theta `$
$`=`$ $`4\pi i{\displaystyle _{\theta _1+\mu \eta }^{\theta _0}}\theta {\displaystyle \frac{u^{}(\theta )}{1+u(\theta )}}\text{d}\theta 2\pi ib(\theta _0){\displaystyle _{\theta _1+\mu \eta }^{\theta _0}}{\displaystyle \frac{u^{}(\theta )}{1+u(\theta )}}\text{d}\theta .`$
Since $`u^{}(\theta )/(1+u(\theta ))=\frac{\text{d}}{\text{d}\theta }\mathrm{log}(1+u(\theta ))`$ for any branch $`\mathrm{log}`$ of the logarithm defined on an open section of $`^{}`$ containing $`]\mathrm{},0[`$ we have
$$_{\theta _1+\mu \eta }^{\theta _0}\theta \frac{u^{}(\theta )}{1+u(\theta )}\text{d}\theta =\left[\theta \text{Log}(1+u(\theta ))\right]_{\theta _1+\mu \eta }^{\theta _0}_{\theta _1+\mu \eta }^{\theta _0}\text{Log}(1+u(\theta ))\text{d}\theta ,$$
and
$$_{\theta _1+\mu \eta }^{\theta _0}\frac{u^{}(\theta )}{1+u(\theta )}\text{d}\theta =\left[\text{Log}(1+u(\theta ))\right]_{\theta _1+\mu \eta }^{\theta _0}.$$
We therefore get
$`R(2\theta _0,y)`$ $`=`$ $`{\displaystyle \frac{1}{4\pi ^2}}\underset{\eta 0_+}{lim}\left(4\pi iR_1(\theta _0,\eta )R_2(\theta _0,\eta )\right)`$
$`=`$ $`{\displaystyle \frac{b(\theta _0)}{2}}\theta _1{\displaystyle \frac{1}{4\pi ^2}}\text{Log}(z_0)\text{Log}(w_\epsilon )`$
$`+{\displaystyle \frac{i}{\pi }}\underset{\eta 0_+}{lim}\left(R_1(\theta _0,\eta )+{\displaystyle _{\theta _1+\mu \eta }^{\theta _0}}\text{Log}\left(Q_3^\epsilon (\theta )\right)\text{d}\theta \right),`$
where we use that $`1+u_\epsilon (\theta )=Q_3^\epsilon (\theta )=w_\epsilon `$ and (0.89). By Lemma 8 we get
$$R_1(\theta _0,\eta )+_{\theta _1+\mu \eta }^{\theta _0}\text{Log}\left(Q_3^\epsilon (\theta )\right)\text{d}\theta =_{\theta _1+\mu \eta }^{\theta _0}\text{Log}\left(L_\epsilon (\theta )\right)+2\pi ie_\epsilon (\theta )\text{d}\theta ,$$
so
$$\underset{\eta 0_+}{lim}\left(R_1(\theta _0,\eta )+_{\theta _1+\mu \eta }^{\theta _0}\text{Log}\left(Q_3^\epsilon (\theta )\right)\text{d}\theta \right)=e_\epsilon (\theta _0)2\pi i(\theta _0\theta _1)+_{\theta _1}^{\theta _0}\text{Log}(L_\epsilon (\theta ))\text{d}\theta .$$
But then
$`\mathrm{\Psi }_n^{0,b}(2\theta _0,y)\text{CS}(\overline{\rho }_{\theta _0,\epsilon })`$ $`=`$ $`{\displaystyle \frac{1}{6}}+2e_\epsilon (\theta _0)\theta _12f_\epsilon (\theta _0)\theta _0+{\displaystyle \frac{b(\theta _0)}{2}}\theta _1`$
$`+2{\displaystyle _{1/6}^{\theta _0}}\beta _\epsilon (t)\text{d}t{\displaystyle \frac{1}{\pi i}}{\displaystyle _{\theta _1}^{\theta _0}}\text{Log}(L_\epsilon (\theta ))\text{d}\theta (mod).`$
Here
$$2_{1/6}^{\theta _0}\beta _\epsilon (t)\text{d}t=2f_\epsilon (\theta _0)\left(\theta _0\frac{1}{6}\right)+\frac{1}{\pi i}_{\theta _1}^{\theta _0}\text{Log}(L_\epsilon (\theta ))\text{d}\theta ,$$
so
$`\mathrm{\Psi }_n^{0,b}(2\theta _0,y)\text{CS}(\overline{\rho }_{\theta _0,\epsilon })`$ $`=`$ $`{\displaystyle \frac{1}{\pi i}}\left({\displaystyle _{1/6}^{\theta _0}}\text{Log}(L_\epsilon (\theta ))\text{d}\theta {\displaystyle _{\theta _1}^{\theta _0}}\text{Log}(L_\epsilon (\theta ))\text{d}\theta \right)`$
$`+{\displaystyle \frac{1}{6}}+2e_\epsilon (\theta _0)\theta _1{\displaystyle \frac{1}{3}}f_\epsilon (\theta _0)+{\displaystyle \frac{b(\theta _0)}{2}}\theta _1(mod).`$
The subcase $`\theta _0]1/6,1/4[`$. Here we have $`\theta _1=1/6`$ and $`b(\theta _0)=0`$ so the result follows by the fact that $`e_\epsilon (\theta )=f_\epsilon (\theta )`$ for $`\theta ]1/6,1/4[`$.
The subcase $`\theta _0]1/4,1/3[`$. In this case we have $`b(\theta _0)=1`$ and $`\theta _1=1/3`$ so
$`\mathrm{\Psi }_n^{0,b}(2\theta _0,y)\text{CS}(\overline{\rho }_{\theta _0,\epsilon })`$ $`=`$ $`{\displaystyle \frac{1}{\pi i}}\left({\displaystyle _{1/6}^{\theta _0}}\text{Log}(L_\epsilon (\theta ))\text{d}\theta {\displaystyle _{1/3}^{\theta _0}}\text{Log}(L_\epsilon (\theta ))\text{d}\theta \right)`$
$`+{\displaystyle \frac{1}{3}}+{\displaystyle \frac{2}{3}}e_\epsilon (\theta _0){\displaystyle \frac{1}{3}}f_\epsilon (\theta _0)(mod).`$
By (0.94) we get
$$_{1/3}^{\theta _0}\text{Log}(L_\pm (t))\text{d}t=_{1/3}^{\frac{1}{2}\theta _0}\text{Log}(L_\pm (t))\text{d}t=_{1/6}^{\theta _0}\text{Log}(L_\pm \left(\frac{1}{2}t\right))\text{d}t.$$
By (0.92) and (0.93) we then have
$$_{1/3}^{\theta _0}\text{Log}(L_\pm (t))\text{d}t=_{1/6}^{\theta _0}\text{Log}(L_\pm (t))\text{d}t,$$
so
$$\mathrm{\Psi }_n^{0,b}(2\theta _0,y)\text{CS}(\overline{\rho }_{\theta _0,\epsilon })=\frac{1}{3}+\frac{2}{3}e_\epsilon (\theta _0)\frac{1}{3}f_\epsilon (\theta _0)(mod),$$
and this is zero by Lemma 8, (0.75) and (0.76).
###### Remark 7
If $`u_\pm =u_\pm (v)`$ are the two solutions to $`\varphi (v^2,u)=0`$ for $`v^{}`$ fixed, then $`\lambda _{11}(v,u_{})=\lambda _{11}(v,u_+)^1`$ by Lemma 9. By the proof of Theorem 7 we therefore conclude that $`(v,1+u_\epsilon )`$, $`\epsilon \{\pm \}`$, is a solution to (0.41) and (0.42) if and only if $`(v,1+u_\epsilon )`$ is a solution to (0.42) and
$$v^p=\left(\frac{wv^2}{1v^2w}\right)^q.$$
(0.95)
If we work with the invariants $`\tau _r`$ instead of the invariants $`\overline{\tau }_r`$, then by (0.16) we have to change $`p/q`$ to $`p/q`$ everywhere in the above. The equation (0.42) will be the same but (0.41) will change to (0.95). If $`(v,u)\stackrel{~}{๐ฉ}`$ then we find as in the proof of Theorem 7 that $`(v,w)=(v,u+1)`$ is a solution to (0.95) if and only if $`(v,u)`$ is a solution to
$$v^p=\lambda _{11}(v,u)^q,$$
which is the โwrongโ equation. This is one of the main reasons for working with $`\overline{\tau }_r`$ instead of $`\tau _r`$. Another reason is that one would get a minus sign in (0.81).
Let us end this section by some futher results on the ChernโSimons invariants of flat irreducible $`\text{SU}(2)`$โconnections on $`M_{p/q}`$. By (0.93) $`\overline{\rho }_{\theta ,\epsilon }`$ extends to a $`\text{SU}(2)`$โrepresentation of $`\pi _1(M_{p/q})`$ if and only if $`\overline{\rho }_{\theta ,\epsilon }`$ extends to a $`\text{SU}(2)`$โrepresentation of $`\pi _1(M_{p/q})`$. In that case we find that
$$\text{CS}(\overline{\rho }_{\theta ,\epsilon })=\text{CS}(\overline{\rho }_{\theta ,\epsilon })$$
by using (0.77).
If $`p`$ is even then $`e^{2\pi ip(1/2\theta )}=e^{2\pi ip\theta }`$. Since $`L_\epsilon (1/2\theta )=L_\epsilon (\theta )^1`$ by (0.92) and (0.93) and since $`q`$ is odd it follows that $`\overline{\rho }_{\theta ,\epsilon }`$ extends to a $`\text{SU}(2)`$โrepresentation of $`\pi _1(M_{p/q})`$ if and only if $`\overline{\rho }_{1/2\theta ,\epsilon }`$ extends to such a representation, see (0.68). Assuming that $`\overline{\rho }_{\theta ,\epsilon }`$ extends to such a representation we get
$$\text{CS}(\overline{\rho }_{\frac{1}{2}\theta ,\epsilon })=\text{CS}(\overline{\rho }_{\theta ,\epsilon })\frac{1}{2}\frac{p}{2}(mod)$$
by using (0.94) and (0.77).
By (0.73) and (0.68) we find that $`\overline{\rho }_{1/6,}=\overline{\rho }_{1/6,+}`$ is a $`\text{SU}(2)`$โrepresentation of $`\pi _1(M_{p/q})`$ if and only if $`p=6m+3`$, $`m`$, and q is odd. In that case
$$\text{CS}(\overline{\rho }_{1/6,\epsilon })=\frac{cp}{36}\frac{dq}{4}\frac{1}{2}d(mod).$$
By (0.73) and (0.68) we also have that $`\overline{\rho }_{1/3,}=\overline{\rho }_{1/3,+}`$ is a $`\text{SU}(2)`$โrepresentation of $`\pi _1(M_{p/q})`$ if and only if $`p=6m+3`$, $`m`$, and q is even. In that case
$$\text{CS}(\overline{\rho }_{1/3,\epsilon })=\frac{1}{2}\frac{cp}{9}\frac{dq}{4}(mod).$$
Finally, by (0.73) and (0.68), $`\rho _{1/4,}`$ and $`\rho _{1/4,+}`$ are $`\text{SU}(2)`$โrepresentations of $`\pi _1(M_{p/q})`$ if and only if $`4`$ divides $`p`$. In that case we find that
$`\text{CS}(\overline{\rho }_{1/4,+})`$ $`=`$ $`{\displaystyle \frac{1}{5}}{\displaystyle \frac{cp}{16}}(mod),`$
$`\text{CS}(\overline{\rho }_{1/4,})`$ $`=`$ $`{\displaystyle \frac{1}{5}}{\displaystyle \frac{cp}{16}}(mod).`$
This follows by using (C.5).
We have in the following appendices collected material of a technical nature.
Appendix A Proofs of Lemma 1, Lemma 3 and (0.23)
* Proof of Lemma 1 Let $`a>0`$. Let $`\epsilon =1`$ if $`\text{Im}(\zeta )0`$ and let $`\epsilon =1`$ otherwise. Put $`\delta _a^{}=[a,\sqrt{1}\epsilon a]`$ and $`\delta _a^+=[\sqrt{1}\epsilon a,a]`$. (Here, as usual, $`[z_1,z_2]`$ denotes the line segment in $``$ beginning at $`z_1`$ and ending at $`z_2`$.) We have
$$\frac{S_\gamma (\zeta \gamma )}{S_\gamma (\zeta +\gamma )}=\mathrm{exp}\left(\frac{1}{2}_{C_R}\frac{e^{\zeta z}}{\mathrm{sinh}(\pi z)z}\text{d}z\right).$$
By an elementary argument one finds that the integrals $`_{\delta _a^\pm }\frac{e^{\zeta z}}{\mathrm{sinh}(\pi z)z}\text{d}z`$ converge to zero as $`a\mathrm{}`$. Therefore
$$_{C_R}\frac{e^{\zeta z}}{\mathrm{sinh}(\pi z)z}\text{d}z=\epsilon 2\pi \sqrt{1}\left(b_\epsilon +\underset{n=1}{\overset{\mathrm{}}{}}\text{Res}_{z=\epsilon \sqrt{1}n}\left\{\frac{e^{\zeta z}}{\mathrm{sinh}(\pi z)z}\right\}\right),$$
where $`b_1=0`$ and $`b_1=\text{Res}_{z=0}\left\{\frac{e^{\zeta z}}{\mathrm{sinh}(\pi z)z}\right\}=\frac{\zeta }{\pi }`$. For $`n\{0\}`$ we have
$$\text{Res}_{z=\sqrt{1}n}\left\{\frac{e^{\zeta z}}{\mathrm{sinh}(\pi z)z}\right\}=\frac{(1)^ne^{\sqrt{1}\zeta n}}{\pi \sqrt{1}n},$$
so
$$_{C_R}\frac{e^{\zeta z}}{\mathrm{sinh}(\pi z)z}\text{d}z=(1\epsilon )\sqrt{1}\zeta 2\text{Log}\left(1+e^{\epsilon \sqrt{1}\zeta }\right)$$
giving the result.
To prove the identity (0.23) we use the power series expansion
$$\text{Li}_2(z)=\underset{n=1}{\overset{\mathrm{}}{}}\frac{z^n}{n^2}$$
(A.1)
for the dilogarithm, valid for $`|z|1`$. In the course of the following proof we will establish the identity
$$\frac{\zeta ^2}{2}\frac{\pi ^2}{6}\text{Li}_2(e^{\sqrt{1}\zeta })=\text{Li}_2(e^{\sqrt{1}\zeta })$$
(A.2)
valid for $`\zeta =\pm \pi `$ and all $`\zeta `$ with $`|\text{Re}(\zeta )|<\pi `$.
* Proof of (0.23) Note first that the integral $`A_\gamma (\zeta ):=\frac{1}{4}_{C_R}\frac{e^{\zeta z}}{\mathrm{sinh}(\pi z)\gamma z^2}\text{d}z`$ is convergent for all $`\zeta `$ with $`|\text{Re}(\zeta )|\pi `$ since
$$_{\mathrm{}}^R\frac{e^{\zeta t}}{\mathrm{sinh}(\pi t)t^2}\text{d}t=_R^{\mathrm{}}\frac{e^{\zeta t}}{\mathrm{sinh}(\pi t)t^2}\text{d}t$$
and
$$\left|_R^{\mathrm{}}\frac{e^{\zeta t}}{\mathrm{sinh}(\pi t)t^2}\text{d}t\right|\frac{2}{1e^{2\pi R}}_R^{\mathrm{}}e^{(\text{Re}(\zeta )\pi )t}\frac{1}{t^2}\text{d}t.$$
Let $`b=\text{sign}(\text{Im}(\zeta ))`$, where $`\text{sign}(0)`$ can be put to both $`1`$ and $`1`$ in the following. Moreover, let $`h`$ be a positive parameter and let $`\delta _h^{}(t)=(1+ib)ht+ibh`$ for $`t[1,0]`$ and let $`\delta _h^+(t)=(1ib)ht+ibh`$ for $`t[0,1]`$. It is elementary to show that the integrals $`_{\delta _h^\pm }\frac{e^{\zeta z}}{\mathrm{sinh}(\pi z)z^2}\text{d}z`$ converge to zero as $`h`$ converges to infinity for $`|\text{Re}(\zeta )|\pi `$. By the residue theorem we conclude that
$$A_\gamma (\zeta )=\frac{1}{4\gamma }2\pi i\underset{n=1}{\overset{\mathrm{}}{}}\text{Res}_{z=in}\left\{\frac{e^{\zeta z}}{\mathrm{sinh}(\pi z)z^2}\right\}$$
for $`\text{Im}(\zeta )0`$ and $`|\text{Re}(\zeta )|\pi `$ and
$$A_\gamma (\zeta )=\frac{1}{4\gamma }2\pi i\underset{n=0}{\overset{\mathrm{}}{}}\text{Res}_{z=in}\left\{\frac{e^{\zeta z}}{\mathrm{sinh}(\pi z)z^2}\right\}$$
for $`\text{Im}(\zeta )0`$ and $`|\text{Re}(\zeta )|\pi `$. Using (A.1) this leads directly to
$$A_\gamma (\zeta )=\{\begin{array}{cc}\frac{1}{2\sqrt{1}\gamma }\text{Li}_2(e^{\sqrt{1}\zeta })\hfill & ,\text{Im}(\zeta )0,\hfill \\ \frac{1}{2\sqrt{1}\gamma }\left[\frac{\zeta ^2}{2}\frac{\pi ^2}{6}\text{Li}_2(e^{\sqrt{1}\zeta })\right]\hfill & ,\text{Im}(\zeta )0.\hfill \end{array}$$
for $`|\text{Re}(\zeta )|\pi `$. Left is to prove the identity (A.2). To this end, let
$$g(\zeta )=\frac{\zeta ^2}{2}\frac{\pi ^2}{6}\text{Li}_2(e^{i\zeta })\text{Li}_2(e^{i\zeta })$$
for $`\zeta \mathrm{\Omega }:=\{\zeta ||\text{Re}(\zeta )|<\pi \}`$. By (0.22) we have
$$g^{}(\zeta )=\zeta +i\left(\text{Log}(1+e^{i\zeta })\text{Log}(1+e^{i\zeta })\right)$$
and therefore $`e^{ig^{}(\zeta )}=1`$. Since $`\mathrm{\Omega }`$ is connected, $`g`$ is $`C^1`$ and $`g^{}(0)=0`$ we get that $`g^{}`$ is identically zero on $`\mathrm{\Omega }`$ so $`g`$ is constant on $`\mathrm{\Omega }`$. Now $`g(0)=\frac{\pi ^2}{6}2\text{Li}_2(1)`$ and
$$\text{Li}_2(1)=\underset{n=1}{\overset{\mathrm{}}{}}\frac{(1)^n}{n^2}=\frac{\pi ^2}{12}$$
so $`g(0)=0`$. Finally we note that $`g`$ is well-defined and continuous on $`\mathrm{\Omega }\{\pm \pi \}`$ so $`g(\pm \pi )=0`$ by continuity.
Note that the function $`g`$ in the above proof is a well-defined analytic function on $`W=\{\zeta \text{Re}(\zeta )\pi +2\pi \}`$ and that $`g`$ is continuous on $`W`$. As in the proof above we find that $`g`$ is constant on each connected component of $`W`$. Moreover, we can on each of these connected components choose a branch of the dilogarithm such that $`g`$ extends to a continuous (and hence a constant) function on the connected set $`W`$.
* Proof of Lemma 3 The function $`1/\mathrm{sinh}(w)`$ has a simple pole at $`w=0`$ with principal part $`1/w`$, i.e.
$$\varphi (w)=\frac{1}{\mathrm{sinh}(w)}\frac{1}{w}$$
is holomorphic in a neighborhood of zero, in fact on the open disk $`D(0,\pi )`$ with centre $`0`$ and radius $`\pi `$. Let $`aR`$, let $`C_{R;a}=[a,R]\mathrm{{\rm Y}}_R[a,R]`$, and write $`I_\gamma (\zeta )=J_\gamma (\zeta )+K_\gamma (\zeta )`$, where
$`4J_\gamma (\zeta )`$ $`=`$ $`{\displaystyle _{C_{R;a}}}{\displaystyle \frac{e^{\zeta z}}{\mathrm{sinh}(\pi z)z}}\varphi (\gamma z)\text{d}z,`$
$`4K_\gamma (\zeta )`$ $`=`$ $`{\displaystyle _{\mathrm{}}^a}{\displaystyle \frac{e^{\zeta z}}{\mathrm{sinh}(\pi z)z}}\varphi (\gamma z)\text{d}z+{\displaystyle _a^{\mathrm{}}}{\displaystyle \frac{e^{\zeta z}}{\mathrm{sinh}(\pi z)z}}\varphi (\gamma z)\text{d}z.`$
To estimate $`K_\gamma (\zeta )`$ we simply use that
$$|\varphi (\gamma t)|\frac{1}{\mathrm{sinh}(\gamma t)}+\frac{1}{\gamma t}\frac{2}{\gamma t}$$
for $`t>0`$ leading to the bound
$`|K_\gamma (\zeta )|`$ $``$ $`{\displaystyle \frac{1}{2\gamma }}{\displaystyle _a^{\mathrm{}}}{\displaystyle \frac{e^{\text{Re}(\zeta )t}+e^{\text{Re}(\zeta )t}}{\mathrm{sinh}(\pi t)}}{\displaystyle \frac{1}{t^2}}\text{d}t`$
$``$ $`{\displaystyle \frac{1}{\gamma (1e^{2\pi a})}}{\displaystyle _a^{\mathrm{}}}\left(e^{(\pi \text{Re}(\zeta ))t}+e^{(\pi +\text{Re}(\zeta ))t}\right){\displaystyle \frac{1}{t^2}}\text{d}t.`$
For $`|\text{Re}(\zeta )|\pi `$ we therefore get
$$|K_\gamma (\zeta )|\frac{2}{\gamma (1e^{2\pi a})}_a^{\mathrm{}}\frac{1}{t^2}\text{d}t=\frac{2}{a\gamma (1e^{2\pi a})}.$$
If $`|\text{Re}(\zeta )|<\pi `$ we find that
$`|K_\gamma (\zeta )|`$ $``$ $`{\displaystyle \frac{1}{\gamma a^2(1e^{2\pi a})}}{\displaystyle _a^{\mathrm{}}}(e^{(\pi \text{Re}(\zeta ))t}+e^{(\pi +\text{Re}(\zeta ))t})\text{d}t`$
$`=`$ $`{\displaystyle \frac{1}{\gamma a^2(1e^{2\pi a})}}\left({\displaystyle \frac{e^{(\pi \text{Re}(\zeta ))a}}{\pi \text{Re}(\zeta )}}+{\displaystyle \frac{e^{(\pi +\text{Re}(\zeta ))a}}{\pi +\text{Re}(\zeta )}}\right).`$
Next let us estimate $`J_\gamma (\zeta )`$. First we use the standard estimate
$$\left|_{\mathrm{{\rm Y}}_R}\frac{e^{\zeta z}}{\mathrm{sinh}(\pi z)z}\varphi (\gamma z)\text{d}z\right|\pi RM(\zeta ,R),$$
where $`M(\zeta ,R)=\mathrm{max}_{z\mathrm{{\rm Y}}_R}\left|\frac{e^{\zeta z}}{\mathrm{sinh}(\pi z)z}\varphi (\gamma z)\right|.`$ We have
$$\left|\varphi (w)\right|=\frac{\mathrm{sinh}(w)w}{w\mathrm{sinh}(w)}.$$
Here $`\mathrm{sinh}(w)w=w^3h(w)`$ and $`w\mathrm{sinh}(w)=w^2k(w)`$, where $`h`$ and $`k`$ are entire functions. Note that $`k`$ is different from zero on $`D(0,\pi )`$. Since $`\gamma ]0,1[`$ we get
$$M(\zeta ,R)=2\gamma L(R)N(\zeta ,R),$$
where $`L(R)=\mathrm{max}_{|z|R}|h(z)/k(z)|`$ and $`N(\zeta ,R)=\mathrm{max}_{z\mathrm{{\rm Y}}_R}\left|\frac{e^{\zeta z}}{e^{\pi z}e^{\pi z}}\right|.`$ We note that
$$N(\zeta ,R)Q_\pm (R)\underset{z\mathrm{{\rm Y}}_R}{\mathrm{max}}\left|e^{(\zeta \pm \pi )z}\right|,$$
where $`Q_\pm (R)=\mathrm{max}_{z\mathrm{{\rm Y}}_R}\frac{1}{\left|1e^{\pm 2\pi z}\right|}`$. Put $`Q(R)=Q_{}(R)+Q_+(R)`$ and get
$$N(\zeta ,R)Q(R)\underset{\mu =\pm 1}{\mathrm{min}}\left(\underset{z\mathrm{{\rm Y}}_R}{\mathrm{max}}\left|e^{(\zeta +\mu \pi )z}\right|\right).$$
Since $`\text{Re}(z)[R,R]`$ and $`\text{Im}(z)[0,R]`$ for $`z\mathrm{{\rm Y}}_R`$ we finally get
$$N(\zeta ,R)Q(R)e^{2\pi R}\left(1+e^{\text{Im}(\zeta )R}\right).$$
We have thus obtained the estimate
$$\left|\frac{1}{4}_{\mathrm{{\rm Y}}_R}\frac{e^{\zeta z}}{\mathrm{sinh}(\pi z)z}\varphi (\gamma z)\text{d}z\right|\gamma B\left(1+e^{\text{Im}(\zeta )R}\right),$$
where $`B=\frac{\pi }{2}RL(R)Q(R)e^{2\pi R}`$.
Finally we have to estimate $`_R^a\frac{e^{\zeta t}e^{\zeta t}}{\mathrm{sinh}(\pi t)t}\varphi (\gamma t)\text{d}t`$. First observe that
$$h(y)=\underset{n=0}{\overset{\mathrm{}}{}}\frac{y^{2n}}{(2n+3)!}$$
and
$$k(y)=\underset{n=0}{\overset{\mathrm{}}{}}\frac{y^{2n}}{(2n+1)!}$$
so $`h(y)k(y)/6`$ for $`y`$. Therefore
$`\left|{\displaystyle _R^a}{\displaystyle \frac{e^{\zeta t}e^{\zeta t}}{\mathrm{sinh}(\pi t)t}}\varphi (\gamma t)\text{d}t\right|`$ $``$ $`{\displaystyle \frac{\gamma }{3}}{\displaystyle _R^a}{\displaystyle \frac{e^{\text{Re}(\zeta )t}+e^{\text{Re}(\zeta )t}}{e^{\pi t}e^{\pi t}}}\text{d}t`$
$``$ $`{\displaystyle \frac{\gamma }{3(1e^{2\pi R})}}{\displaystyle _R^a}\left(e^{(\pi \text{Re}(\zeta ))t}+e^{(\pi +\text{Re}(\zeta ))t}\right)\text{d}t.`$
For $`|\text{Re}(\zeta )|\pi `$ we get
$$|_R^a\frac{e^{\zeta t}e^{\zeta t}}{\mathrm{sinh}(\pi t)t}\varphi (\gamma (t)\text{d}t|\frac{2a\gamma }{3(1e^{2\pi R})}.$$
If $`|\text{Re}(\zeta )<\pi `$ we get
$`|{\displaystyle _R^a}{\displaystyle \frac{e^{\zeta t}e^{\zeta t}}{\mathrm{sinh}(\pi t)t}}\varphi (\gamma (t)\text{d}t|{\displaystyle \frac{\gamma }{3(1e^{2\pi R})}}{\displaystyle _0^a}(e^{(\pi \text{Re}(\zeta ))t}+e^{(\pi +\text{Re}(\zeta ))t})\text{d}t`$
$`={\displaystyle \frac{\gamma }{3(1e^{2\pi R})}}\left({\displaystyle \frac{1e^{(\pi \text{Re}(\zeta ))a}}{\pi \text{Re}(\zeta )}}+{\displaystyle \frac{1e^{(\pi +\text{Re}(\zeta ))a}}{\pi +\text{Re}(\zeta )}}\right).`$
The lemma now follows by putting $`a=1/\gamma >1>R`$ and $`A=\frac{13}{12(1e^{2\pi R})}`$.
Appendix B Proofs of the estimates (0.28) and (0.30)
Let
$`J_\pm ^{^{}}(r,\epsilon )`$ $`=`$ $`{\displaystyle _{C_\pm (\epsilon )}}\left(\mathrm{exp}\left(I_\gamma (\pi 2\pi x)I_\gamma (\pi +2\pi x)\right)1\right)e^{r\mathrm{\Phi }(x)}\text{d}x,`$
$`J_\pm ^{^{\prime \prime }}(r,\epsilon )`$ $`=`$ $`{\displaystyle _{C_\pm (\epsilon )}}(\mathrm{tan}(\pi rx)\sqrt{1})g_r(x)\text{d}x,`$
where $`\mathrm{\Phi }`$ is given by (0.29). Note first that we are free to deform the contour $`C_\pm (\epsilon )`$ as long as we stay inside the domain of analyticity of the integrands. For the integrals $`J_\pm ^{^{}}(r,\epsilon )`$ we will deform $`C_\pm (\epsilon )`$ to $`[\epsilon ,1\epsilon ]`$. Since the integrands of the integrals $`J_\pm ^{^{\prime \prime }}(r,\epsilon )`$ are analytic on $`\mathrm{\Omega }_{\frac{1}{2}r}\{(m+1/2)/r|m=0,1,\mathrm{},r1\}`$, we can deform $`C_\pm (\epsilon )`$ to $`C_\pm (0)`$ in these integrals without changing their sum, i.e.
$$J_{}^{^{\prime \prime }}(r,\epsilon )+J_+^{^{\prime \prime }}(r,\epsilon )=J_{}^{^{\prime \prime }}(r)+J_+^{^{\prime \prime }}(r),$$
where $`J_\pm ^{^{\prime \prime }}(r)=J_\pm ^{^{\prime \prime }}(r,0)`$. In the following calculations we will need the identity
$$\text{Re}\left(\mathrm{\Phi }(x)\right)=2\pi \text{Re}(x)\text{Im}(x)\pi \text{Im}(x)\frac{1}{\pi }\text{Im}\left(\text{Li}_2\left(e^{2\pi \text{Im}(x)}e^{2\pi i\text{Re}(x)}\right)\right)$$
(B.1)
valid for $`x\mathrm{\Omega }_{\mathrm{}}`$ (see (0.20)). This identity is an immediate consequence of (A.2).
* Proof of (0.28) Let us first estimate $`J_+^{^{\prime \prime }}(r)`$. We partition $`C_+(0)`$ into the three pieces $`C_+^1=[\sqrt{1},0]`$, $`C_+^2=[1+\sqrt{1},\sqrt{1}]`$, and $`C_+^3=[1,1+\sqrt{1}]`$. Put
$$I_+^i(r)=_{C_+^i}(\mathrm{tan}(\pi rx)\sqrt{1})g_r(x)\text{d}x.$$
By (0.26) we immediately get
$$|I_+^1(r)|2_0^1e^{2\pi rt}|g_r(\sqrt{1}t)|\text{d}t.$$
To be able to use (0.23) we introduce the positive parameter $`\epsilon `$ again. In fact, we have by (0.23), Lemma 3 and Lebesgueโs dominated convergence theorem that
$$|I_+^1(r)|2\mathrm{exp}(2A+2B\pi /r)\underset{\epsilon 0_+}{lim}_0^1e^{2\pi rt}\left|e^{r\mathrm{\Phi }(\epsilon +\sqrt{1}t)}\right|\text{d}t.$$
By (B.1) we immediately get that
$$\left|e^{r\mathrm{\Phi }(\epsilon +\sqrt{1}t)}\right|=\mathrm{exp}\left(r\left[2\pi \epsilon t\pi t\frac{1}{\pi }\text{Im}\left(\text{Li}_2\left(e^{2\pi \sqrt{1}\epsilon }e^{2\pi t}\right)\right)\right]\right).$$
Now by continuity of $`(t,\epsilon )\text{Li}_2\left(e^{2\pi \sqrt{1}\epsilon }e^{2\pi t}\right)`$ on $`[0,1]\times [0,1]`$ we can remove the parameter $`\epsilon `$ again by using Lebusgueโs dominated convergence theorem once more. This gives us
$$\underset{\epsilon 0_+}{lim}_0^1e^{2\pi rt}\left|e^{r\mathrm{\Phi }(\epsilon +\sqrt{1}t)}\right|\text{d}t=_0^1e^{3\pi rt}\text{d}t$$
leading to the estimate
$$|I_+^1(r)|\frac{2}{3\pi r}\mathrm{exp}(2A+2B\pi /r)\left(1e^{3\pi r}\right).$$
Next we estimate $`I_+^2`$. By (0.26) we get
$$|I_+^2(r)|4e^{2\pi r}_0^1|g_r(\sqrt{1}+t)|\text{d}t.$$
Similarly to the analysis of $`I_+^1`$ we introduce the parameter $`\epsilon `$ and get
$$|I_+^2(r)|4\mathrm{exp}(2A+2B\pi /r)e^{2\pi r}\underset{\epsilon 0_+}{lim}_\epsilon ^{1\epsilon }|e^{r\mathrm{\Phi }(\sqrt{1}+t)}|\text{d}t.$$
Here
$$|e^{r\mathrm{\Phi }(\sqrt{1}+t)}|=e^{\pi r(12t)}\mathrm{exp}\left(\frac{r}{\pi }\text{Im}\left(\text{Li}_2(e^{2\pi }e^{2\pi \sqrt{1}t})\right)\right)$$
by (B.1), so by Lebesgueโs dominated convergence theorem we get
$$|I_+^2(r)|4\mathrm{exp}(2A+2B\pi /r)e^{3\pi r}_0^1e^{2\pi rt}\mathrm{exp}\left(\frac{r}{\pi }\text{Im}\left(\text{Li}_2(e^{2\pi }e^{2\pi \sqrt{1}t})\right)\right)\text{d}t.$$
By definition of the dilogarithm we have
$$\text{Im}\left(\text{Li}_2(e^{2\pi }e^{2\pi \sqrt{1}t})\right)=_0^1\frac{\text{Arg}\left(1se^{2\pi }e^{2\pi \sqrt{1}t}\right)}{s}\text{d}s,$$
which is non-negative for $`t[0,1/2]`$. For $`t[1/2,1]`$ we use that
$$\text{Im}\left(\text{Li}_2(e^{2\pi }e^{2\pi \sqrt{1}t})\right)=\text{Im}\left(\text{Li}_2(e^{2\pi \sqrt{1}t})\right)+_{1e^{2\pi }}^1\frac{\text{Arg}\left(1se^{2\pi \sqrt{1}t}\right)}{s}\text{d}s,$$
where the last integral is positive. The first term is bounded from below by $`_{n=1}^{\mathrm{}}\frac{1}{n^2}=\frac{\pi ^2}{6}`$ by (0.24). We therefore end up with
$`|I_+^2(r)|`$ $``$ $`4\mathrm{exp}(2A+2B\pi /r)e^{(31/6)\pi r}{\displaystyle _0^1}e^{2\pi rt}\text{d}t`$
$`=`$ $`{\displaystyle \frac{2}{\pi r}}\mathrm{exp}(2A+2B\pi /r)e^{(31/6)\pi r}\left(e^{2\pi r}1\right).`$
Finally, we estimate $`I_+^3`$. Similarly to the other cases we get
$$|I_+^3(r)|2\mathrm{exp}(2A+2B\pi /r)\underset{\epsilon 0_+}{lim}_0^1e^{2\pi rt}\left|e^{r\mathrm{\Phi }(1\epsilon +\sqrt{1}t)}\right|\text{d}t.$$
By (B.1) we have
$$\left|e^{r\mathrm{\Phi }(1\epsilon +\sqrt{1}t)}\right|=\mathrm{exp}\left(r\left[\pi (12\epsilon )t\frac{1}{\pi }\text{Im}\left(\text{Li}_2(e^{2\pi t}e^{2\pi \sqrt{1}\epsilon })\right)\right]\right)$$
leading to
$$|I_+^3(r)|2\mathrm{exp}(2A+2B\pi /r)_0^1e^{\pi rt}\text{d}t=\frac{2}{\pi r}\mathrm{exp}(2A+2B\pi /r)\left(1e^{\pi r}\right).$$
By letting $`C_{}^1=[\sqrt{1},0]`$, $`C_{}^2=[1\sqrt{1},\sqrt{1}]`$, and $`C_{}^3=[1,1\sqrt{1}]`$ and
$$I_{}^i(r)=_{C_{}^i}(\mathrm{tan}(\pi rx)+\sqrt{1})g_r(x)\text{d}x$$
we find (now by the use of (0.27)) an upper bound for $`|I_{}^i(r)|`$ identical with the upper bound for $`|I_+^i(r)|`$, $`i=1,2,3`$, with the exception that $`\mathrm{exp}(2A+2B\pi /r)`$ should be replaced by $`\mathrm{exp}(2A+4B\pi /r)`$ in these bounds. To conclude we have shown that there exists a constant $`K_1`$ independent of $`r`$ and $`\epsilon `$ such that
$$|J_+^{^{\prime \prime }}(r,\epsilon )+J_{}^{^{\prime \prime }}(r,\epsilon )|\underset{i=1}{\overset{3}{}}\left(|I_+^i|+|I_{}^i|\right)\frac{K_1}{r}$$
for all $`r_2`$.
* Proof of (0.30) By the remarks prior to the proof of (0.28) we have
$$J_\pm ^{^{}}(r,\epsilon )=_\epsilon ^{1\epsilon }\mathrm{exp}\left(r\mathrm{\Phi }(t)\right)h_\gamma (t)\text{d}t,$$
where
$$h_\gamma (t)=\mathrm{exp}\left(I_\gamma (\pi 2\pi t)I_\gamma (\pi +2\pi t)\right)1.$$
The integrand is continuous on $`[0,1]`$. Therefore
$$\left|J_\pm ^{^{}}(r,\epsilon )\right|_0^1\mathrm{exp}\left(r\text{Re}\left(\mathrm{\Phi }(t)\right)\right)\left|h_\gamma (t)\right|\text{d}t$$
for all $`\epsilon ]0,\frac{1}{4r}[`$. By (0.24) and (0.25) and the remarks prior to (0.25) we have
$$\text{Re}\left(\mathrm{\Phi }(t)\right)=\frac{1}{2\pi }\text{Im}\left(\text{Li}_2(e^{2\pi \sqrt{1}t})\text{Li}_2(e^{2\pi \sqrt{1}t})\right)=\frac{1}{\pi }\text{Cl}_2(2\pi t)\frac{1}{2\pi }\text{Vol}(K).$$
Therefore
$$|J_\pm ^{^{}}(r,\epsilon )|\mathrm{exp}\left(\frac{r}{2\pi }\text{Vol}(K)\right)_0^1\left|h_\gamma (t)\right|\text{d}t.$$
By definition we have
$$h_\gamma (t)=\underset{n=1}{\overset{\mathrm{}}{}}\frac{1}{n!}\left(I_\gamma (\pi 2\pi t)I_\gamma (\pi +2\pi t)\right)^n.$$
From Lemma 3 we get
$$\left|I_\gamma (\pi 2\pi t)I_\gamma (\pi +2\pi t)\right|\left(Cf(t)+D\right)\gamma $$
for $`t]0,1[`$, where $`C`$ and $`D`$ are positive constants independent of $`\gamma `$ and $`t`$, and
$$f(t)=\frac{1}{t}+\frac{1}{1t}.$$
Since $`f:]0,1[`$ is bigger than or equal to $`4`$ we can choose $`C`$ so big that
$$\left|I_\gamma (\pi 2\pi t)I_\gamma (\pi +2\pi t)\right|Cf(t)\gamma $$
for $`t]0,1[`$. From Lemma 3 we also have
$$\left|h_\gamma (t)\right|\mathrm{exp}\left(|I_\gamma (\pi 2\pi t)I_\gamma (\pi +2\pi t)|\right)\mathrm{exp}(4A+4B\pi /r)$$
for $`t[0,1]`$, where $`A`$ and $`B`$ are as in Lemma 3, so $`\frac{\pi }{r}\mathrm{exp}(4A+4B\pi /r)`$ is an upper bound for both of the integrals $`_0^\gamma |h_\gamma (t)|\text{d}t`$ and $`_{1\gamma }^1|h_\gamma (t)|\text{d}t`$ (for $`r4`$). Left is to evaluate (for $`r7`$)
$$_\gamma ^{1\gamma }|h_\gamma (t)|\text{d}t\underset{n=1}{\overset{\mathrm{}}{}}\frac{C^n\gamma ^n}{n!}_\gamma ^{1\gamma }f(t)^n\text{d}t.$$
By using that
$$f(t)^n=\underset{k=0}{\overset{n}{}}\left(\begin{array}{c}n\\ k\end{array}\right)\frac{1}{t^k}\frac{1}{(1t)^{nk}}$$
we get
$`{\displaystyle _\gamma ^{1\gamma }}f(t)^n\text{d}t`$ $`=`$ $`{\displaystyle \underset{k=0}{\overset{n}{}}}\left(\begin{array}{c}n\\ k\end{array}\right){\displaystyle _\gamma ^{1\gamma }}{\displaystyle \frac{1}{t^k}}{\displaystyle \frac{1}{(1t)^{nk}}}\text{d}t`$
$``$ $`{\displaystyle \underset{k=0}{\overset{n}{}}}\left(\begin{array}{c}n\\ k\end{array}\right)\left(2^{nk}{\displaystyle _\gamma ^{1/2}}{\displaystyle \frac{1}{t^k}}\text{d}t+2^k{\displaystyle _{1/2}^{1\gamma }}{\displaystyle \frac{1}{(1t)^{nk}}}\text{d}t\right)`$
$``$ $`2^n{\displaystyle \underset{k=0}{\overset{n}{}}}\left(\begin{array}{c}n\\ k\end{array}\right)\left({\displaystyle _\gamma ^{1/2}}{\displaystyle \frac{1}{t^k}}\text{d}t+{\displaystyle _\gamma ^{1/2}}{\displaystyle \frac{1}{t^{nk}}}\text{d}t\right)`$
$``$ $`2^{n+1}\left({\displaystyle \underset{k=0}{\overset{n}{}}}\left(\begin{array}{c}n\\ k\end{array}\right)\right){\displaystyle _\gamma ^{1/2}}{\displaystyle \frac{1}{t^n}}\text{d}t=2^{2n+1}{\displaystyle _\gamma ^{1/2}}{\displaystyle \frac{1}{t^n}}\text{d}t.`$
Here
$$_\gamma ^{1/2}\frac{1}{t}\text{d}t=\text{Log}(2)\text{Log}(\gamma )\text{Log}(r)$$
and
$$_\gamma ^{1/2}\frac{1}{t^n}\text{d}t=\frac{1}{n1}\left(\frac{1}{\gamma ^{n1}}2^{n1}\right)\frac{1}{n1}\frac{1}{\gamma ^{n1}}$$
for $`n2`$. Therefore
$`{\displaystyle _\gamma ^{1\gamma }}|h_\gamma (t)|\text{d}t`$ $``$ $`2\gamma \left(4C\text{Log}(r)+{\displaystyle \underset{n=2}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(4C)^n}{(n1)n!}}\right)`$
$``$ $`{\displaystyle \frac{2\pi }{r}}\left(4C\text{Log}(r)+\mathrm{exp}(4C)4C1\right).`$
We conclude that there exists a constant $`K_2`$ independent of $`r`$ and $`\epsilon `$ such that
$$|J_+^{^{}}(r,\epsilon )+J_{}^{^{}}(r,\epsilon )|K_2\frac{\text{Log}(r)}{r}\mathrm{exp}\left(\frac{r}{2\pi }\text{Vol}(K)\right),$$
for all $`r_2`$.
Appendix C The case $`M_0`$
The manifold $`M_0`$ is the mapping torus of a torus with monodromy matrix
$$\left(\begin{array}{cc}2& 1\\ 1& 1\end{array}\right),$$
see \[18, p. 366\]. The invariant $`\tau _r(M_0)`$ has been calculated by Jeffrey \[13, Theorem 4.1\]. This theorem gives the large $`r`$ asymptotics of the invariant as well. In fact, we have
$$\tau _r(M_0)=\frac{1}{2}\frac{1}{2\sqrt{5}}\frac{1}{\sqrt{5}}\left(\mathrm{exp}\left(2\pi \sqrt{1}r\left(\frac{1}{5}\right)\right)+\mathrm{exp}\left(2\pi \sqrt{1}r\frac{1}{5}\right)\right)$$
(C.1)
which at the same time can be taken as the large $`r`$ asymptotics of the invariant. Let us relate this result to our contour integral formula for the invariant $`\tau _r(M_0)=\overline{\tau }_r(M_0)`$. By Lemma 2
$$\tau _r(M_0)=\frac{ri}{4}_{C_r^1}\mathrm{cot}(\pi rx)\left(_{C_r^2}\mathrm{tan}(\pi ry)f_{0,r}(x,y)\text{d}y\right)\text{d}x,$$
where
$$f_{0,r}(x,y)=\mathrm{sin}(\pi x)e^{2\pi irxy}\frac{S_{\pi /r}(\pi +2\pi (xy))}{S_{\pi /r}(\pi +2\pi (x+y))}.$$
Following the discussion in Sec. 4.3 the relevant (shifted) phase functions to consider are given by
$$\mathrm{\Psi }_0^{a,b}(x,y)=ax+byxy+\frac{1}{4\pi ^2}\left(\text{Li}_2(e^{2\pi i(x+y)})Li_2(e^{2\pi i(xy)})\right),$$
where $`a,b`$ are certain integers. If we put $`\mathrm{\Psi }=\mathrm{\Psi }_0^{a,b}`$, then by (4.3) $`(x,y)`$ is a critical point of $`\mathrm{\Psi }`$ if and only if
$`0`$ $`=`$ $`2\pi i(ay)+\text{Log}(1zw)\text{Log}(1zw^1),`$
$`0`$ $`=`$ $`2\pi i(bx)+\text{Log}(1zw)+\text{Log}(1zw^1),`$ (C.2)
where $`z=e^{2\pi ix}`$ and $`w=e^{2\pi iy}`$ as usual, and this set of equations implies that
$`wz`$ $`=`$ $`1zw,`$
$`(1zw)(wz)`$ $`=`$ $`zw,`$ (C.3)
compare with (0.41) and (0.42). We note that the first of these equations is equivalent to
$$w1=z(1w)$$
so $`w=1`$ or $`z=1`$. For $`z=1`$ we get $`w^2+3w+1=0`$ so $`w=\frac{3\pm \sqrt{5}}{2}`$. For $`w=1`$ we find that $`z^23z+1=0`$ so $`z=\frac{3\pm \sqrt{5}}{2}`$. However, only the point $`(z,w)=((3\sqrt{5})/2,1)`$ satisfies that $`zw,zw^1[1,\mathrm{}[`$. For this point we find that $`y`$ and from (5.3) we get that $`y=a`$. But then
$$\mathrm{\Psi }(x,y)=ax+byxy+\frac{1}{4\pi ^2}\left(\text{Li}_2(z)\text{Li}_2(z)\right)=yx+baxy=ba.$$
Let $`\overline{\rho }_{\theta ,\epsilon }`$ be the nonabelian $`\text{SU}(2)`$โrepresentations of $`\pi `$ from Proposition 2, where $`\theta [1/3,1/6][1/6,1/3]`$ and $`\epsilon \{\pm \}`$. Here $`\overline{\rho }_{\theta ,\epsilon }`$ and $`\overline{\rho }_{\theta ,\epsilon }`$ are conjugate. By (0.68) $`\overline{\rho }_{\theta ,\epsilon }`$ extends to a representation of $`\pi _1(M_0)`$ if and only if
$$L_\epsilon (\theta )=1.$$
But this happens if and only if $`\theta =\pm 1/4`$ for both $`\epsilon =+`$ and $`\epsilon =`$, i.e. the set of conjugacy classes of nonabelian representations of $`\pi _1(M_0)`$ into $`\text{SU}(2)`$ is $`\{[\overline{\rho }_{1/4,}],[\overline{\rho }_{1/4,+}]\}`$.
By (0.80) the flat reducible $`\text{SU}(2)`$โconnections on $`M_0`$ all have a ChernโSimons invariant equal to zero, so we conclude that the image set of the $`\text{SU}(2)`$ ChernโSimons functional on $`M_0`$ has at most three elements. By \[18, Theorem 5.6 and precedent text\] we can therefore conclude that the set of ChernโSimons invariants of flat $`\text{SU}(2)`$โconnections on $`M_0`$ is
$$\{0mod,\frac{1}{5}mod,\frac{1}{5}mod\}.$$
(C.4)
In particular, the set of ChernโSimons invariants of $`[\overline{\rho }_{1/4,+}]`$ and $`[\overline{\rho }_{1/4,}]`$ is equal to $`\{\frac{1}{5}mod,\frac{1}{5}mod\}`$. We note that (C.1) and (C.4) prove the AEC for the invariants $`\tau _r(M_0)`$.
By (0.73) and (0.74) we have $`\beta _\pm (1/4)=f_\pm `$, where $`f_{}=0`$ and $`f_+=1`$. By Proposition 4 and (0.74) we find
$$\text{CS}(\overline{\rho }_{1/4,\pm })=\pm \frac{1}{6}\frac{1}{\pi }_{1/6}^{1/4}\text{Arg}(L_\pm (t))\text{d}t(mod),$$
where we use that $`L_\pm (t)S^1`$. By (0.93) we have
$$_{1/6}^{1/4}\text{Arg}(L_{}(t))\text{d}t=_{1/6}^{1/4}\text{Arg}(L_+(t))\text{d}t$$
so we finally get
$$\text{CS}(\overline{\rho }_{1/4,\pm })=\pm \left(\frac{1}{6}\frac{1}{\pi }_{1/6}^{1/4}\text{Arg}(L_+(t))\text{d}t\right)(mod).$$
Note that $`\text{Im}(L_+(t))<0`$ on $`]1/6,1/4[`$ so
$$\frac{1}{12}=\frac{1}{\pi }(\pi )(1/41/6)<\frac{1}{\pi }_{1/6}^{1/4}\text{Arg}(L_+(t))\text{d}t<0.$$
Therefore
$$1/6\frac{1}{\pi }_{1/6}^{1/4}\text{Arg}(L_+(t))\text{d}t]\frac{1}{6},\frac{1}{4}[.$$
Since this value $`mod`$ belongs to the set $`\{\pm 1/5mod\}`$ we conclude that it is equal to $`1/5`$ so
$$\frac{1}{\pi }_{1/6}^{1/4}\text{Arg}(L_{}(t))\text{d}t=\frac{1}{\pi }_{1/6}^{1/4}\text{Arg}(L_+(t))\text{d}t=\frac{1}{30}.$$
(C.5)
Let us finally calculate the value of $`\mathrm{\Psi }=\mathrm{\Psi }_0^{a,b}`$ in the critical points corresponding to the solutions $`(z,w)=((3+\sqrt{5})/2,1)`$ and $`(z,w)=(1,\frac{3\pm \sqrt{5}}{2})`$ to (5.3). Since the set of solutions to (5.3) is in one to one correspondence with the set of conjugacy classes of nonabelian $`\text{SL}(2,)`$โrepresentations of $`\pi _1(M_0)`$ by the proof of Theorem 7 and since the subset of solutions $`(z,w)`$ with $`zS^1`$ and $`w]\mathrm{},0[`$ corresponds to the set of conjugacy classes of $`\text{SU}(2)`$โrepresentations we see that the points $`(z,w)=(1,\frac{3\pm \sqrt{5}}{2})`$ correspond to nonabelian $`\text{SU}(2)`$โrepresentations of $`\pi _1(M_0)`$ while the points $`(z,w)=(\frac{3\pm \sqrt{5}}{2},1)`$ correspond to nonabelian $`\text{SL}(2,)`$โrepresentations of $`\pi _1(M_0)`$ which are not equivalent to $`\text{SU}(2)`$โrepresentations.
For $`(z,w)=(e^{2\pi ix},e^{2\pi iy})=((3+\sqrt{5})/2,1)`$ we find again that $`y=a`$ and then $`\mathrm{\Psi }(x,y)=ab`$ (independent of the choice of branch of the dilogarithm along $`]1,\mathrm{}[`$).
Finally, let $`(z,w)=(e^{2\pi ix},e^{2\pi iy})=(1,\frac{3\pm \sqrt{5}}{2})`$. The real values of the right-hand sides of (5.3) do not depend on $`a`$ and $`b`$. Taking the imaginary values of these equations we get
$`0`$ $`=`$ $`2\pi \left(a\text{Re}(y)\right)+\text{Im}\left(\text{Log}(1+w)\text{Log}(1+w^1)\right),`$
$`0`$ $`=`$ $`2\pi (bx)+\text{Im}\left(\text{Log}(1+w)+\text{Log}(1+w^1)\right).`$
The second of these equations is equivalent to
$$xb=\frac{1}{2\pi }\left(\text{Arg}(1+w)+\text{Arg}(1+w^1)\right)=\frac{1}{2}$$
for both $`w=\frac{3\pm \sqrt{5}}{2}`$, and the first is equivalent to
$$a\text{Re}(y)=\frac{1}{2\pi }\left(\text{Arg}(1+w^1)\text{Arg}(1+w)\right)=\{\begin{array}{cc}\frac{1}{2},& w=\frac{3+\sqrt{5}}{2},\\ \frac{1}{2},& w=\frac{3\sqrt{5}}{2}.\end{array}$$
We have
$$\mathrm{\Psi }(x,y)=ax+byxy+\frac{1}{4\pi ^2}\left(\text{Li}_2(|w|)\text{Li}_2(1/|w|)\right),$$
and since this is real by Corollary 5, we get
$$\mathrm{\Psi }(x,y)=ax+(bx)\text{Re}(y)+\frac{1}{4\pi ^2}\text{Re}\left(\text{Li}_2(|w|)\text{Li}_2(1/|w|)\right).$$
Here
$`ax+(bx)\text{Re}(y)`$ $`=`$ $`(bx)\left(\text{Re}(y)a\right)+ab=ab+{\displaystyle \frac{1}{2}}\left(a\text{Re}(y)\right)`$
$`=`$ $`ab+\{\begin{array}{cc}\frac{1}{4},& w=\frac{3+\sqrt{5}}{2},\\ \frac{1}{4},& w=\frac{3\sqrt{5}}{2}.\end{array}`$
For $`z[0,\mathrm{}[`$ we have
$$\text{Li}_2(z)+\text{Li}_2(1/z)=\frac{\pi ^2}{6}\frac{1}{2}\text{Log}^2(z).$$
This identity e.g. follows by differentiating the difference of the two sides in the identity, showing that this difference is constant on $`[0,\mathrm{}[`$, and then evaluating in $`z=1`$ using that $`\text{Li}_2(1)=\pi ^2/12`$. Therefore
$$\text{Li}_2(t)=\frac{\pi ^2}{6}\frac{1}{2}\text{Log}^2(t)\text{Li}_2(t^1)$$
(C.7)
for $`t>1`$ for a branch of $`\text{Li}_2`$ continuously extended across $`]1,\mathrm{}[`$. Let $`w_\pm =\frac{3\pm \sqrt{5}}{2}`$. We note that
$$w_+w_{}=1.$$
(C.8)
Moreover, by \[21, Formula (1.20) p. 7\],
$$\text{Li}_2\left(\frac{3\sqrt{5}}{2}\right)=\frac{\pi ^2}{15}\text{Log}^2\left(\frac{1+\sqrt{5}}{2}\right).$$
(C.9)
Note also that
$$\left(\frac{1+\sqrt{5}}{2}\right)^2=\frac{3+\sqrt{5}}{2}.$$
(C.10)
Assume that $`w=w_{}=\frac{3\sqrt{5}}{2}`$. Then, by (C.7),
$$\text{Li}_2(|w|)=\frac{\pi ^2}{6}\frac{1}{2}\text{Log}^2(|w|)\text{Li}_2(|w|^1),$$
where $`\text{Log}(|w|)=\text{Log}(|w|)+i\pi `$, so
$$\text{Li}_2(|w|)=\frac{\pi ^2}{2}\frac{\pi ^2}{6}\frac{1}{2}\text{Log}^2(|w|)\text{Li}_2(|w|^1)i\pi \text{Log}(|w|).$$
But then
$$\text{Li}_2(|w|)\text{Li}_2(|w|^1)=\frac{\pi ^2}{2}\frac{\pi ^2}{6}\frac{1}{2}\text{Log}^2(|w|)2\text{Li}_2(|w|^1)i\pi \text{Log}(|w|).$$
Here $`|w|^1=\frac{3\sqrt{5}}{2}`$ by (C.8) so by (C.9) and (C.10) we get
$$\text{Li}_2(|w|)\text{Li}_2(|w|^1)=\frac{\pi ^2}{2}\frac{\pi ^2}{6}\frac{2\pi ^2}{15}i\pi \text{Log}(|w|)=\frac{\pi ^2}{5}i\pi \text{Log}(|w|).$$
By (C.8) we conclude that
$$\text{Li}_2(|w_\pm |)\text{Li}_2(|w_\pm |^1)=\frac{\pi ^2}{5}i\pi \text{Log}(|w_\pm |).$$
(C.11)
For $`w=w_\pm `$ we therefore get
$$\mathrm{\Psi }(x,y)=ab\pm \frac{1}{4}\frac{1}{20}=ab\pm \frac{1}{5}.$$
Note here that we use the identities (C.5) and (C.11) in the proof of Theorem 8 to handle the cases $`\theta =\pm 1/4`$, so we can not here refer to this theorem.
We have above shown that the set of values $`mod`$ of $`\mathrm{\Psi }_0^{a,b}`$ in its critical points is identical with the set (C.4) of ChernโSimons invariants of flat $`\text{SU}(2)`$โconnections on $`M_0`$ for arbitrary $`a,b`$. The determinant of the Hessian $`H`$ of $`\mathrm{\Psi }_0^{a,b}`$ in a critical point $`(x,y)`$ is given by
$$det(H)=12\left(z+\frac{1}{z}\right)$$
according to (0.49). Thus all the critical points are non-degenerate. It is interesting to note that to obtain the leading order large $`r`$ asymptotics of $`\tau _r(M_0)`$ using the saddle point method we have to use critical points corresponding to $`\text{SL}(2,C)`$โrepresentations of $`\pi _1(M_0)`$ which are not equivalent to $`\text{SU}(2)`$โrepresentations in order to get the part of that leading asymptotics being associated to the reducible flat $`\text{SU}(2)`$โconnections on $`M_0`$. By Conjecture 3 this phenomenon should only occur for $`p/q=0`$. Recall here that the case $`p/q=0`$ is in any case special; only in that case the moduli space of reducible flat $`\text{SU}(2)`$โconnections on $`M_{p/q}`$ is not discrete.
For $`(z,w)=(1,\frac{3\pm \sqrt{5}}{2})`$ we see that $`det(H)=5`$. Moreover, $`z=e^{2\pi ix}=1`$ implies that $`x\frac{1}{2}+`$, so $`\mathrm{sin}(\pi x)=\pm 1`$. Thus we see that the right-hand side of the formula (0.11) gives the part of the leading order large $`r`$ asymptotics of $`\tau _r(M_0)`$ corresponding to the irreducible flat $`\text{SU}(2)`$โconnections on $`M_0`$ if we let $`m_{\overline{\rho }}=4`$ for both points in $`_0^{}`$.
The identity (C.11) was proved by using the explicit value of the dilogarithm in $`(3\sqrt{5})/2`$, cf. (C.9). We note that only very few explicit values of the dilogaritm are known, see \[21, Chap 1\].
Appendix D Proof of Lemma 8
Let us begin by showing the identitites (0.86) and (0.88). Let $`\theta I`$ and let $`u=u_\pm (\theta )`$, $`Q_i=Q_i^\pm (\theta )`$ and $`t=e^{4\pi i\theta }`$ and get
$`Q_1Q_2`$ $`=`$ $`\left(1(1+u)^1t\right)\left(1(1+u)t\right)=1t\left(1+u+(1+u)^1\right)+t^2`$
$`=`$ $`t(1+u)^1\left((1+u)^2(t+t^1)(u+1)+1\right)=t,`$
where the last equality follows by the fact that $`\varphi (t,u)=0`$, where $`\varphi `$ is given by (0.64).
To show (0.88) we observe that
$`{\displaystyle \frac{Q_1Q_3}{Q_2}}`$ $`=`$ $`{\displaystyle \frac{Q_3Q_1^2}{Q_1Q_2}}=t^1Q_3Q_1^2`$
$`=`$ $`t^1(1+u)\left(1t(1+u)^1\right)^2=t^1(1+u)+t(1+u)^12`$
by (0.86). Now $`\varphi (t,u)=0`$ implies that
$$t(1+u)^1=t^2+1tt(u+1)$$
leading to the identity
$$\frac{Q_1Q_3}{Q_2}=1+t^12t+t^2+u(t^1t)=L_\pm (\theta ).$$
To prove the identities (0.85) and (0.87) it is necessary to examine the arguments of the functions $`Q_i^\pm (\theta )`$, $`i=1,2`$. By Lemma 9 we note that
$$Q_1^\pm (\theta )=Q_2^{}(\theta )$$
(D.1)
for $`\theta I`$.
Assume first that $`\theta [1/6,1/3]`$. An elementary calculation shows that
$`1+u_+(\theta )[1,(\sqrt{5}3)/2],`$
$`1+u_{}(\theta )[(3+\sqrt{5})/2,1]`$ (D.2)
for $`\theta [1/6,1/3]`$ and $`1+u_\pm (\theta )=1`$ if and only if $`\theta \{1/6,1/3\}`$. In particular, $`1+u_\pm (\theta )`$ is negative. We therefore get
$$\text{Im}\left(Q_1\right)=|1+u|^1\mathrm{sin}(4\pi \theta )\{\begin{array}{cc}>0,& \theta [1/6,1/4[,\hfill \\ =0,& \theta =1/4,\hfill \\ <0,& \theta ]1/4,1/3].\hfill \end{array}$$
By (5.3) $`\text{Re}(Q_1)=1(1+u)^1\mathrm{cos}(4\pi \theta )`$. Since $`1+u`$ is negative we see that $`\text{Re}(Q_1)`$ have the same sign as $`\mathrm{cos}(4\pi \theta )(1+u)`$, where $`\text{sign}(0)=0`$ as usual. By (5.3) we conclude that
$$\text{Re}\left(Q_1^{}(\theta )\right)>0$$
for all $`\theta [1/6,1/3]`$. Let us next consider $`Q_1^+`$. First note that
$$\mathrm{cos}(4\pi \theta )u_+(\theta )1=\frac{1}{2}\sqrt{\mathrm{cos}^2(4\pi \theta )\mathrm{cos}(4\pi \theta )\frac{3}{4}}.$$
We therefore get that $`\text{Re}\left(Q_1^+(\theta )\right)`$ has the opposite sign as $`\mathrm{cos}^2(4\pi \theta )\mathrm{cos}(4\pi \theta )1`$. By the assumption on $`\theta `$ we have $`\mathrm{cos}(4\pi \theta )[1,1/2]`$, and we therefore get
$$\text{Re}\left(Q_1^+(\theta )\right)\{\begin{array}{cc}>0,& \theta [1/6,\theta _0[]1/2\theta _0,1/3],\hfill \\ =0,& \theta \{\theta _0,1/2\theta _0\},\hfill \\ <0,& \theta ]\theta _0,1/2\theta _0[,\hfill \end{array}$$
where $`\theta _0]1/6,1/4[`$ is the unique element such that $`\mathrm{cos}(4\pi \theta _0)=(1\sqrt{5})/2]1,1/2[`$ is the negative solution to $`t^2t1=0`$.
Let $`\psi _i^\pm (\theta )]\pi ,\pi ]`$ be the principal argument of $`Q_i^\pm (\theta )`$. Then the above analysis shows, also using (D.1), that
$$\psi _2^{}(\theta )=\psi _1^+(\theta )\{\begin{array}{cc}]0,\frac{\pi }{2}[,\hfill & \theta [\frac{1}{6},\theta _0[,\hfill \\ =\frac{\pi }{2},\hfill & \theta =\theta _0,\hfill \\ ]\frac{\pi }{2},\pi [,\hfill & \theta ]\theta _0,\frac{1}{4}[,\hfill \\ =\pi ,\hfill & \theta =\frac{1}{4},\hfill \\ ]\pi ,\frac{\pi }{2}[,\hfill & \theta ]\frac{1}{4},\frac{1}{2}\theta _0[,\hfill \\ =\frac{\pi }{2},\hfill & \theta =\frac{1}{2}\theta _0,\hfill \\ ]\frac{\pi }{2},0[,\hfill & \theta ]\frac{1}{2}\theta _0,\frac{1}{3}],\hfill \end{array}$$
(D.3)
and
$$\psi _2^+(\theta )=\psi _1^{}(\theta )\{\begin{array}{cc}]0,\frac{\pi }{2}[,\hfill & \theta [\frac{1}{6},\frac{1}{4}[,\hfill \\ =0,\hfill & \theta =\frac{1}{4},\hfill \\ ]\frac{\pi }{2},0[,\hfill & \theta ]\frac{1}{4},\frac{1}{3}],\hfill \end{array}$$
(D.4)
so
$$\psi _1^\pm (\theta )+\psi _2^\pm (\theta )\{\begin{array}{cc}]0,\pi [,\hfill & \theta [\frac{1}{6},\theta _0[,\hfill \\ ]\frac{\pi }{2},\frac{3\pi }{2}[,\hfill & \theta ]\theta _0,\frac{1}{4}[,\hfill \\ ]\frac{\pi }{2},\pi ],\hfill & \theta =\frac{1}{4},\hfill \\ ]\frac{3\pi }{2},\frac{\pi }{2}[,\hfill & \theta ]\frac{1}{4},\frac{1}{2}\theta _0[,\hfill \\ ]\pi ,0[,\hfill & \theta [\frac{1}{2}\theta _0,\frac{1}{3}].\hfill \end{array}$$
By (0.86) we have
$$\psi _1^\pm (\theta )+\psi _2^\pm (\theta )4\pi \theta +2\pi ,$$
so we conclude that
$$\psi _1^\pm (\theta )+\psi _2^\pm (\theta )=4\pi \theta $$
for all $`\theta [1/6,1/3]`$ proving (0.85) for these $`\theta `$.
Next assume that $`\theta [1/3,1/6]`$. First observe that
$$Q_i^\pm (\theta )=\overline{Q_i^\pm (\theta )},\theta I\{\pm 1/4\},$$
(D.5)
and
$$Q_i^\pm (1/4)=Q_i^\pm (1/4)$$
(D.6)
for $`i=1,2`$. By (D.6) we immediately get that (0.85) holds for $`\theta =1/4`$. For $`\theta [1/3,1/6]\{1/4\}`$, (0.85) follows by (D.5) and the fact that $`\text{Log}(\overline{p})=\overline{\text{Log}(p)}`$ for $`p]\mathrm{},0]`$.
Note that (0.87) is true if we choose $`e_\pm (\theta )`$ such that
$$\psi _1^\pm (\theta )+\pi \psi _2^\pm (\theta )e_\pm (\theta )2\pi ]\pi ,\pi ].$$
By (D.3) and (D.4) we have that $`\psi _1^\pm (\theta )\psi _2^\pm (\theta )]\pi ,\pi ]`$, and we conclude that we have to put $`e_\pm (\theta )=0`$ if and only if $`\psi _1^\pm (\theta )\psi _2^\pm (\theta )`$ and $`e_\pm (\theta )=1`$ elsewhere. By (D.3) and (D.4) we conclude that $`e_{}(1/4)=0`$ and $`e_+(1/4)=1`$.
Assume that $`\theta [1/6,1/3]\{1/4\}`$. Then $`\psi _1^\pm (\theta )`$ and $`\psi _2^\pm (\theta )`$ both belong to either $`]\pi ,0[`$ or to $`]0,\pi [`$. We use this fact together with the fact that $`\mathrm{cot}:]m\pi ,(m+1)\pi [`$ is strictly decreasing for any $`m`$. In fact, $`\mathrm{cot}(\psi _i)=\frac{\text{Re}(Q_i)}{\text{Im}(Q_i)}`$ so
$`\mathrm{cot}(\psi _1)`$ $`=`$ $`\mathrm{cot}(4\pi \theta )+{\displaystyle \frac{|1+u|}{\mathrm{sin}(4\pi \theta )}},`$
$`\mathrm{cot}(\psi _2)`$ $`=`$ $`\mathrm{cot}(4\pi \theta )+{\displaystyle \frac{|1+u|^1}{\mathrm{sin}(4\pi \theta )}}.`$
By this we find that
$$\text{sign}\left(\psi _1^\pm (\theta )\psi _2^\pm (\theta )\right)=\text{sign}\left(\mathrm{sin}(4\pi \theta )\left(1\left(1+u_\pm (\theta )\right)^2\right)\right).$$
Since $`(1+u_+(\theta ))^21`$ and $`(1+u_{}(\theta ))^21`$ with equalities if and only if $`\theta \{1/6,1/3\}`$, we get $`\psi _1^\pm (\theta )=\psi _2^\pm (\theta )`$ for $`\theta \{1/6,1/3\}`$ and
$$\psi _1^\epsilon (\theta )<\psi _2^\epsilon (\theta )$$
for $`\theta ]1/6,1/4[`$ and $`\epsilon =`$ or for $`\theta ]1/4,1/3[`$ and $`\epsilon =+`$. Moreover,
$$\psi _1^\epsilon (\theta )>\psi _2^\epsilon (\theta )$$
for $`\theta ]1/6,1/4[`$ and $`\epsilon =+`$ or for $`\theta ]1/4,1/3[`$ and $`\epsilon =`$. We therefore get that (0.87) is true if we put $`e_\pm (\theta )=0`$ for $`\theta \{1/6,1/3\}`$ and
$$e_+(\theta )=\{\begin{array}{cc}1,& \theta ]1/6,1/4[,\hfill \\ 0,& \theta ]1/4,1/3[,\hfill \end{array}$$
and let $`e_{}(\theta )=1e_+(\theta )`$ for $`\theta ]1/6,1/3[\{1/4\}`$.
Let us finally consider the case $`\theta [1/3,1/6]`$. First observe that
$$Q_3^\pm (\theta )=Q_3^\pm (\theta )$$
(D.7)
for all $`\theta I`$. By this and (D.6) we conclude that $`e_\pm (1/4)=e_\pm (1/4)`$. For $`\theta [1/3,1/6]\{1/4\}`$ we get by (D.5) and (D.7) that
$`\text{Log}\left(Q_1^\pm (\theta )\right)+\text{Log}\left(Q_3^\pm (\theta )\right)\text{Log}\left(Q_2^\pm (\theta )\right)=\text{Log}\left(Q_3^\pm (\theta )\right)\overline{\text{Log}\left(Q_3^\pm (\theta )\right)}`$
$`+\overline{\text{Log}\left(Q_1^\pm (\theta )\right)+\text{Log}\left(Q_3^\pm (\theta )\right)\text{Log}\left(Q_2^\pm (\theta )\right)}`$
$`=2i\text{Im}\left(\text{Log}\left(Q_3^\pm (\theta )\right)\right)+\overline{\left(\text{Log}\left({\displaystyle \frac{Q_1^\pm (\theta )Q_3^\pm (\theta )}{Q_2^\pm (\theta )}}\right)+e_\pm (\theta )2\pi i\right)}.`$
By using that $`Q_3`$ is negative and that $`L_\pm (\theta )]\mathrm{},0]`$ for $`\theta ]1/6,1/3[`$ we get for $`\theta ]1/3,1/6[\{1/4\}`$ that
$`\text{Log}\left(Q_1^\pm (\theta )\right)+\text{Log}\left(Q_3^\pm (\theta )\right)\text{Log}\left(Q_2^\pm (\theta )\right)`$ $`=`$ $`\text{Log}\left({\displaystyle \frac{Q_1^\pm (\theta )Q_3^\pm (\theta )}{Q_2^\pm (\theta )}}\right)`$
$`+\left(1e_\pm (\theta )\right)2\pi i,`$
so we conclude that $`e_\pm (\theta )=1e_\pm (\theta )`$ for these $`\theta `$. Finally we get for $`\theta \{1/3,1/6\}`$ that $`Q_1^\pm (\theta )=Q_2^\pm (\theta )`$ so (0.87) is satisfied for $`e_\pm (\theta )=0`$ for these $`\theta `$.
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# The Indefinite Logarithm, Logarithmic Units, and the Nature of Entropy
## 1 Introduction
The goal of this paper is to help clear up what is perceived to be a widespread confusion that can found in many popular sources (websites, popular books, etc.) regarding the proper mathematical status of a variety of physical quantities that are conventionally defined on logarithmic scales.
As an example of a logarithmic quantity about which much confusion still lingers, we focus on the quantity of thermodynamic entropy, and its close relationship to (and really, identity with) the concepts of entropy and information as defined within the context of information theory. Although many physicists and information theorists have understood quite well the mathematical reasons for the underlying unity between the entropy concepts in these two domains, many others still do not, and continue to believe that the physical and information-theoretic concepts of entropy are somehow fundamentally different from each other.
Some (although not all) of the confusion that we have seen expressed in this regard can be traced back to the historical accident that thermodynamic entropy is most often measured in natural logarithm units, while information-theoretic entropy is more frequently measured in units of the logarithm base 2 (i.e., in bits), or some multiple thereof. But of course, the choice of the logarithm base in the definition of entropy is completely inessential, and amounts merely to a choice of oneโs unit of measurement, which went without saying in Boltzmannโs era, and which Shannon himself pointed out in his seminal work on information theory.
Even further, the supposed distinction between the โphysicalโ nature of thermodynamic entropy (as measured in, say, Joules per Kelvin) and the allegedly more โmathematicalโ nature of information entropy (measured in bits) can also be seen as a totally artificial distinction, one resulting from nothing other than the fact that early thermodynamicists were not yet aware that physical entropy really *is* nothing other than a measurable manifestation of what is at root merely a purely mathematical, statistical quantity.
In fact, as we will review, any given unit of physical entropy can be *exactly* identified with a corresponding (purely abstract) *mathematical* unit, while still remaining consistent with all observed empirical data. The fundamental scientific principle of adopting the most parsimonious theory that explains the data (a.k.a. Ockhamโs razor) then *demands* that as good scientists we *must* indeed adopt this identification between the physical and mathematical domains, and take it seriously as holding the status of our best available model of reality, at least until empirical evidence to the contrary is found.
Although these issues are already quite well understood in certain circles, we nevertheless felt that, as a public service, it would be worthwhile to compose a short paper that elaborates on the mathematical foundations of these issues in some detail. Two fundamental mathematical concepts which I have to be found rather useful in explaining these kinds of issues are concepts that I refer to as the โindefinite logarithmโ and โlogarithmic units.โ The definition and discussion of these concepts will form the main mathematical core of this paper.
Although this material seems to be already essentially common (or intuitively obvious) knowledge among many of the leading researchers who deal every day with the physics of information, I have found in my experience that misunderstandings and confusion regarding these issues nevertheless still abound in other communities.
The reader should please note that, since this material seems to hold the status of being considered obvious or common knowledge in certain circles, this paper is by no means intended to claim any kind of intellectual priority on these issues. Rather, it is being written simply because the author is not presently aware of an accessible reference on this subject that explicitly explains these issues with a sufficient degree of pedagogical detail to satisfy general audiences.
The author welcomes comments and feedback from readers that may help point the author at seminal references or review articles in the mathematical literature that may elucidate these same issues, though quite possibly using different terminology.
## 2 The Indefinite Logarithm
We use the standard notation $`\mathrm{log}_ba`$ for the logarithm, base $`b`$, of $`a`$. Of course, this expression is well-defined for all real $`a,b>0`$, and even for all complex $`a,b0`$ as a multi-valued function. But, what if no particular base $`b`$ is selected?
Of course, as a matter of notational convenience in mathematical literature, $`\mathrm{log}a`$ is often defined to be simply a shorthand for the frequently-used natural logarithm, $`\mathrm{ln}a=\mathrm{log}_ea`$, or, in more everyday applied contexts, for the decimal-based $`\mathrm{log}_{10}a`$. But the topic of this paper is not situations such as these in which some definite base really exists but is merely left implicit by the notation. Rather, here we would like to discuss the concept of a โnewโ kind of logarithm function wherein *no specific base is implied at all*. We dub this the *indefinite logarithm*, and we will give it a formal definition in a moment.
Of course, without a specific base, the result of the logarithm cannot be an ordinary number, since any specific numeric result would imply some specific base that must have been used. Instead, we can decare the output of the indefinite logarithm to be a different (i.e., non-numeric) type of mathematical object representing the result of performing this more abstract operation.
The form of this new type of object can be rigorously defined using standard mathematical concepts. Of course, as with any type of mathematical object, there is an infinite variety of ways in which we could satisfactorily represent these new objects in terms of more standard mathematical objects. Here is the representation that we find most convenient for purposes of this paper:
###### Definition 2.1 (Indefinite logarithm).
For any given real number $`x>0`$, the *indefinite logarithm $`L`$ of $`x`$*, written $`L=[\mathrm{log}x]`$, is a special type of mathematical object called a โlogarithmic quantityโ object, which we define as follows:
$$L=[\mathrm{log}x]:\lambda (b>0).\mathrm{log}_bx=\{(b,y)|b>0,y=\mathrm{log}_bx\}.$$
(1)
Here, $`:`$ denotes โis defined as,โ and in the first expression after the $`:`$ we are using Churchโs lambda-calculus notation for functions (see, for example, ), which gives us a concise way of saying that the indefinite logarithm $`[\mathrm{log}x]`$ for any given $`x`$ is defined to be the unary function object $`L:`$ mapping real numbers $`b>0`$ to the logarithm $`y=\mathrm{log}_bx`$ of the (given constant) $`x`$ to the (variable, function argument) base $`b`$. Meanwhile, on the right, we are merely writing out the standard โgraph representationโ of this function object explicitly as the set of all ordered pairs $`(b,y)`$ consisting of a base $`b>0`$, followed by the ordinary (definite) logarithm $`y=\mathrm{log}_bx`$, of $`x`$ to the base $`b`$.
Clearly, the result of the indefinite logarithm, as defined above, does not select any preferred base, yet it contains โall of the informationโ about the logarithm of $`x`$ taken to all possible bases. So, in this sense, it is not any less descriptive than a definite logarithm.
Although the above definition is restricted to positive real numbers (since this all that we need for our subsequent discussions), it could easily be extended to non-zero complex numbers if desired.
We can define addition of indefinite logarithms by adding their corresponding $`y`$ values:
###### Definition 2.2 (Sum of indefinite logarithms).
Given two indefinite logarithm objects $`L_1`$ and $`L_2`$, their sum $`L=L_1+L_2`$ is defined by $`L(b):L_1(b)+L_2(b)`$ for all $`b>0`$. Or, stated a bit more formally using Lambda calculus, $`L:\lambda b.L_1(b)+L_2(b)`$.
We can similarly define negation of indefinite logarithms by simply negating their $`y`$ values:
###### Definition 2.3 (Negation of indefinite logarithms).
Given an indefinite logarithm object $`L`$, its negative $`L^{}=L`$ is defined by $`L^{}(b):L(b)`$ for all $`b>0`$. Or, $`L^{}:\lambda b.L(b)`$.
If we add any indefinite logarithm to its negation, or take the indefinite logarithm of 1, we get a unique indefinite-logarithm object called the *null* or 0 indefinite logarithm, which returns 0 for all bases:
###### Definition 2.4 (Null indefinite logarithm).
The indefinite logarithm of 1, i.e. $`L_0=[\mathrm{log}1]`$, is (by previous definitions) the function $`\lambda b.0`$ over reals $`b>0`$. This $`L_0`$ will be called the *null indefinite logarithm* and will sometimes be written $`[0]`$.
Of course, is the identity element for the addition operation on natural logarithms; that is, for any $`L`$, we have $`L+[0]=L`$.
Note that there is no corresponding concept of a unit indefinite logarithm, i.e., a multiplicative identity. That is, indefinite logarithms are inherently scale-free objects; that is, they are non-scalar quantities. Further, the space of indefinite logarithms does not even need to be considered to be closed under multiplication. A meaningful multiplication operation can be defined if a product of two indefinite logarithms is considered to be a distinct type (similarly to how the product of two lengths is an area), but we will not develop that here.
Consistently with all of the above definitions, and with the ordinary definition of multiplication as repeated addition, indefinite logarithms can also be multiplied by arbitrary (positive, negative, or zero) real numbers:
###### Definition 2.5 (Indefinite logarithms multiplied by scalars).
Given an indefinite logarithm quantity $`L`$ and real number $`r`$, define the product $`L^{}=rL=Lr`$ of $`L`$ times $`r`$ by $`L^{}(b):rL(b)`$. That is, $`L^{}=\lambda b.rL(b)`$.
Finally, solving the preceding expression for $`r`$ allows us to recognize and define the result of the ratio of two indefinite logarithms as being an ordinary number:
###### Definition 2.6 (Ratio of indefinite logarithms).
Given two indefinite logarithms $`L_1`$ and $`L_2`$, their ratio $`r=L_1/L_2`$ is defined as the real number $`r:L_1(b)/L_2(b)`$, where $`b`$ is any positive real number. (The value of $`r`$ does not depend on $`b`$.)
As an immediate consequence of the above definitions, the ratio of the indefinite logarithms of two numbers $`a`$ and $`c`$ is simply $`[\mathrm{log}a]/[\mathrm{log}c]=\mathrm{log}_ca`$, which (note) is the same as the ratio $`\mathrm{log}_ba/\mathrm{log}_bc`$ of the definite logarithms of $`a`$ and $`c`$ to any common base $`b`$. Thus, to emphasize, *the ratio of two logarithms is independent of what base we are working in, and thus remains well-defined even for indefinite logarithms*.
The above fact is important for our discussions in subsequent sections.
It is also worth noting that the indefinite logarithm shares all of the mathematically important properties of the ordinary logarithm (aside from not being a number), including the following useful identities:
* $`[\mathrm{log}xy]=[\mathrm{log}x]+[\mathrm{log}y]`$
* $`[\mathrm{log}x/y]=[\mathrm{log}x][\mathrm{log}y]`$
* $`[\mathrm{log}x^y]=y[\mathrm{log}x]`$.
Finally, it is useful to also define an *indefinite exponential* function, which maps a given indefinite logarithm object back to the unique real number of which it is the indefinite logarithm.
###### Definition 2.7 (Indefinite exponential).
For any indefinite logarithm object $`L=[\mathrm{log}x]`$, let *the indefinite exponential of $`L`$*, written $`[\mathrm{exp}L]`$, be given by simply $`[\mathrm{exp}L]=x`$.
Another way to define $`[\mathrm{exp}L]`$, which is helpful when we are not explicitly given the $`x`$ such that $`L=[\mathrm{log}x]`$, is simply to say that $`[\mathrm{exp}L]=b^{L(b)}`$, where $`b>0`$ is any positive real number; all such $`b`$ give the same value for $`[\mathrm{exp}L]`$.
## 3 Logarithmic Quantities
In the above, we occasionally referred to the indefinite logarithm objects as โquantitiesโ in order to anticipate what we will now discuss, which is that indefinite logarithmic quantities (which we defined as pure mathematical entities) behave formally in a way that is exactly analogous to how dimensional *physical* quantities (such as length, time, and mass) behave. Indeed, logarithmic quantities can be viewed as โnaturalโ dimensional quantities that exist independently of any particular models of physics.
Even further, later we will argue that logarithmic quantities can be understood as being the underlying essence behind certain quantities (in particular, thermodynamic entropy and physical information) that are frequently perceived as being โphysicalโ rather than mathematical in nature. We will also argue that the question of whether these quantities are โreallyโ physical or mathematical ones is an ill-posed one, being predicated on an entirely false dichotomy that has no real meaning.
First, what do we mean by a *quantity*, in general? (Regardless, for now, of whether it is supposed to be โmathematicalโ or โphysical.โ) For our purposes, a quantity is an object selected from a set having a structure similar to that of the real number system, but without any built-in unit, that is, with no pre-ordained object to be designated โ1.โ In abstract algebra terms, a set of quantities forms an (abstract) vector space over the reals, with a definite 0 quantity, an addition operation, a negation operation, and the ability to multiply by reals, but without a predefined unit quantity, and without necessarily any assigned meaning for a product of quantities.
A bit more generally and formally, we can define:
###### Definition 3.1 (Quantity spaces).
Given any field $`F`$ (in the standard abstract algebra sense of โfield,โ i.e., a commutative division ring with unity), a *quantity space Q over F* is simply a vector space over $`F`$, that is, an Abelian group of objects to be called *quantities*, which is closed under the additional operation of multiplication by the (scalar) elements of $`F`$.
Although quantity spaces, being vector spaces, may in general be many-dimensional, in this article we will primarily work with examples of quantity spaces that are only one-dimensional.
We can now define the concept of a *logarithmic (quantity) space*.
###### Definition 3.2 (Logarithmic spaces).
A *logarithmic space* (or *logarithmic quantity space*) is a quantity space $`Q`$ over $``$ in which each quantity $`qQ`$ is identified with an indefinite logarithm object $`L`$ (as defined in the preceding section), or with a series of indefinite logarithm objects, in the case of multidimensional logarithmic quantity spaces. The vector addition, negation, and scalar multiplication operations are identified with the corresponding operations on the indefinite logarithms. The null (0) vector comes from the null indefinite logarithm.
Now, a *logarithmic quantity* $`q`$ is simply a member of some logarithmic quantity space $`Q`$. By a *scalar* logarithmic quantity or *logarithmic scalar*, we mean a member of a one-dimensional logarithmic space. Members of $`n`$-dimensional logarithmic quantity spaces will be called *$`n`$-dimensional* logarithmic quantities.
## 4 Logarithmic Units
Quantities in general (and logarithmic quantities in particular) have the property that there is no natural, built-in unit quantity that is automatically provided by the quantity space itself. However, in any given quantity space $`Q`$, we can always choose some arbitrary $`uQ`$ to be designated as a provisional *unit quantity*, and then all quantities in $`Q`$ can be described in terms of scalar multiples of that unit. (For elements of multidimensional quantity spaces, a series of multiples is needed.) We may even have several different units $`u_1,u_2,\mathrm{}Q`$, and express quantities sometimes as multiples of $`u_1`$, sometimes as multiples of $`u_2`$, etc., and convert between expressions utilizing different units by multiplying them by appropriate conversion factors.
Of course, we are already familiar with these properties of quantities from their use in ordinary physics, in which (for example) spatial distances (and multi-dimensional displacement vectors) are considered to be quantities, rather than just pure numbers, and we can choose any number of units (meters, feet, etc.) for expressing them. Space itself (in the traditional continuum description) does not have any natural โunit length,โ only arbitrary units that we chose by convention.<sup>1</sup><sup>1</sup>1Emerging theories of quantum gravity suggest that the Planck length $`\mathrm{}_P=\sqrt{\mathrm{}G/c^3}`$ (or some small multiple of it) may play the role of a โnaturalโ unit length, in some sense which is not yet fully understood. Nevertheless, we are still free, if we wish, to treat lengths as quantities that can be represented in arbitrary units. Other examples of commonly used physical quantities include time, velocity, mass, and energy. (Of course, there are many others as well.)
Now, the primary observation of the previous section is that spaces of indefinite logarithm objects (logarithmic spaces) provide exactly the structure of quantity spaces; thus, we can represent all one-dimensional indefinite logarithm objects as scalar multiples of some arbitrarily chosen โunitโ indefinite logarithm object. Indefinite logarithm objects thus naturally have the same mathematical status, in this sense, as do physical quantities.
## 5 Logarithmic Scales
Logarithmic quantities and units, in one guise or another, are of course very widely used today, for quantifying a wide variety of concepts in different fields of study. Some examples include:
* Relative signal amplitudes or power levels, in physics and engineering.
* Earthquake strength (Richter scale) in seismology.
* Tonal intervals on a musical scale.
* Entropy (in, as we will see, both the thermodynamic and information-theory senses).
* Information, in the information theory sense.
What is lacking presently, however, is the ubiquitous understanding that all of these disparate types of quantities can be understood as dealing with what is fundamentally the same underlying system of logarithmic units, as we defined above. The various โdifferentโ logarithmic scales that are in use are really distinguished only by different choices of terminology for discussing logarithmic quantities, different sizes and names of the logarithmic units used for expressing them, and the application of these units in describing different domains of study.
To illustrate, let us now identify and name a variety of indefinite logarithm objects that are popularly used as units in which logarithmic quantities of interest are expressed in various fields. This list is ordered from the smallest logarithmic unit to the largest, emphasizing that logarithmic units (like numbers) are comparable across domains.
* $`\mathrm{cent}=[\mathrm{log}2]/1,200`$. In music theory, the *cent* is 1/100th of a minor second, or 1/1,200<sup>th</sup> of an octave.
* $`\mathrm{m2}=[\mathrm{log}2]/12`$. In music theory, the *minor second* m2 is 100 cents or 1/12<sup>th</sup> of an octave.
* $`\mathrm{M2}=[\mathrm{log}2]/6`$. In music theory, the *major second* M2 is 200 cents or 1/6<sup>th</sup> of an octave.
* $`\mathrm{dB}=0.1[\mathrm{log}10]`$. The *decibel*. This is the smallest logarithmic unit in widespread use outside music theory, usually for expressing the magnitude of the ratio between signal strengths.
* $`\mathrm{b}=[\mathrm{log}2]`$. In information theory, the binary digit or *bit*. This is the smallest non-null logarithmic unit with an integer argument. In music theory, the same logarithmic unit is called an *octave* P8.
* $`\mathrm{n}=[\mathrm{log}e]`$. The natural-log unit or *nat*. As we will explain in more detail later, this mathematical unit can be exactly identified with the physical unit $`k_\mathrm{B}`$ known as Boltzmannโs constant. When used to express a ratio of current or voltage levels, the nat is called a *Neper* or Np.
* $`\mathrm{Np}=2[\mathrm{log}e]`$. The magnitude of the *Neper* of a ratio of currents or voltages, when translated to a ratio of power levels. (It is doubled because the power is the square of the current or voltage.)
* $`\mathrm{o}=[\mathrm{log}8]=3\mathrm{b}`$. The octal digit, which could be abbreviated *oit* (in analogy with bit). It is equal to three bits. Used as an information unit in computer engineering.
* $`\mathrm{d}=[\mathrm{log}10]`$. The decimal digit, abbreviable as *dit*. In various contexts, this unit is also known as *Bel*, *power of ten*, *order of magnitude*, *Richter-scale point*, or *decade*.
* $`\mathrm{h}=[\mathrm{log}16]=4\mathrm{b}`$. In computer engineering, the hexadecimal digit is a unit of information, which might be called a *hit*, but in practice, it is called a *nibble* or *nybble*.
* $`\mathrm{B}=[\mathrm{log}256]=8\mathrm{b}=2\mathrm{h}`$. The usual definition of a *byte* in computer engineering; sometimes called an *octet* in network engineering.
* $`\mathrm{kcal}/\mathrm{mol}/\mathrm{K}503.6k_\mathrm{B}/\mathrm{molecule}[\mathrm{log}(4.9\times 10^{218})]/\mathrm{molecule}`$. In chemistry, the *kilocalorie per mole per degree Kelvin* is a common intensive unit of thermodynamic entropy, equivalent to about $`503.6k_\mathrm{B}`$ (or nats or Nepers) per molecule.
* $`\mathrm{kb}=1,000\mathrm{b}=[\mathrm{log}2^{1000}]`$. An information unit known as a *kilobit* in telecommunications.
* $`\mathrm{kb}=1,024\mathrm{b}=[\mathrm{log}2^{2^{10}}]`$. An information unit called a *kibibit*, also known as a *kilobit* in computer engineering.
* Multiplying the above definitions by 8 gives the standard definitions for the *kilobyte* unit of information, in the telecommunication and computer-engineering contexts respectively.
* Similarly for higher powers of 1,000 (or 1,024), with prefixes mega- (M), giga- (G), tera- (T), peta- (P), exa- (E), zetta- (Z), and yotta- (Y).
* $`\mathrm{J}/\mathrm{K}7.243\times 10^{22}k_\mathrm{B}11\mathrm{ZB}[\mathrm{log}10^{3.14558\times 10^{22}}]`$. In thermodynamics, the *Joule per Kelvin* is a common extensive unit of bulk thermodynamic entropy. Converted into information units, it is about 11 zettabytes, meaning $`11\times 1,024^7\mathrm{B}`$.
* $`\mathrm{kcal}/\mathrm{K}=4186.8\mathrm{J}/\mathrm{K}3.03\times 10^{26}k_\mathrm{B}45.2\mathrm{YB}`$. In chemistry, the *kilocalorie per degree Kelvin* is a common extensive unit of bulk thermodynamic entropy. In information units, it is about 45 yottabytes, meaning $`45\times 1,024^8\mathrm{B}`$.
Of course, one could systematically define and name still larger (or smaller) logarithmic units by applying larger (or smaller) order-of-magnitude prefixes to the above.
The point of this exercise is to emphasize that all of the supposed disparate logarithmic scales that are in use in these various fields are ultimately all just different views of the same fundamental logarithmic scale. The various quantities and units discussed on all of these logarithmic scales are all exactly comparable with each other (with the exception of intensive units such as kcal/mol/K, which are only comparable if specific quantity of material is chosen, e.g.1 molecule in the above).
Among the logarithmic quantities that are in widespread use, perhaps the quantity whose status as a logarithmic quantity is least widely appreciated in some circles is the quantity known as *thermodynamic entropy*. Reviewing why this โphysicalโ quantity is indeed, at root, truly just a logarithmic quantity is the subject of the next section.
## 6 Logarithmic Units and Entropy
The original definition of the quantity known as *entropy*, first introduced by Rudolph Clausius in the mid-1800s , was (in differential form)
$$\mathrm{d}S=\mathrm{d}Q/T$$
(2)
where $`\mathrm{d}Q`$ represents an infinitesimal increment of heat energy added to or removed from a system, and $`T`$ is the temperature of the system. Although this is only a differential definition, we can presume that a physical system has a property called its total entropy $`S`$, changes of which correspond to the increments $`\mathrm{d}S`$. (However, the original definition did not specify the base value of $`S`$ for any particular cases.)
Clausius observed that in any thermodynamic process, the total entropy (as he defined it) never decreased, since heat always moved spontaneously from higher-temperature systems to lower-temperature ones, and never vice-versa. He postulated that the principle of the non-decrease of entropy could be introduced as a fundamental law of physics (โsecond law of thermodynamicsโ), equivalent to the other (pre-existing) versions of the second law (impossibility of perpetual motion machines, etc.).
Now, prima facie, Clausiusโ entropy does not seem in any way to be a logarithmic quantity. But, with the subsequent development of statistical mechanics by Maxwell , Boltzmann (see ), and Gibbs in the late 1800s, the thermodynamic entropy came to be understood as really being a statistical quantity that is naturally defined on a logarithmic scale. The โBoltzmanโ form of the definition of entropy (which evolved gradually from the $`H`$ quantity originally defined by Boltzmann in ) was expressed as
$$S=k_\mathrm{B}\mathrm{ln}W,$$
(3)
where $`W`$ denoted the number of possible distinct microscopic ways of arranging the system (consistently with its macroscopic description), and $`k_\mathrm{B}`$ denoted a fundamental entropy unit first used by Planck (according to ) which came to be called *Boltzmannโs constant*, which had a value that (in conventional units of heat over temperature) was found to be equal to about $`1.38\times 10^{23}\mathrm{J}/\mathrm{K}`$.
Now, the traditional stance as to the status of this equation, which is maintained today by many of the more traditional-minded thermodynamicists, is that the entropy $`S`$ is fundamentally a โphysicalโ quantity, namely a ratio of heat to temperature, and Boltzmannโs equation (3) predicts what the value of this quantity will be as a multiple of the Boltzmannโs constant unit, where the multiplier is the pure number obtained from $`\mathrm{ln}W`$.
However, what is arguably the preferred (simpler and more modern) perspective on Boltzmannโs equation is that it is merely a way of rendering (in traditional units) the more elegant and fundamental relation
$$S=[\mathrm{log}W],$$
(4)
where now entropy is taken as being at root just an indefinite logarithmic quantity, in the abstract sense that we outlined in the previous sections.
The modern form (4) can be seen as being exactly equivalent to (3) if we simply declare that
$$k_\mathrm{B}=[\mathrm{log}e],$$
(5)
since $`[\mathrm{log}e]\mathrm{ln}W=[\mathrm{log}W]`$. Note, in particular, that the choice of using $`[\mathrm{log}e]`$ as the unit in the original equation (3) was a completely arbitrary one, and was merely a consequence of the choice of using the base-$`e`$ (โnaturalโ) logarithm in the formula. So, we could equally validly re-render eq. (3) in any of the following ways:
$`S`$ $`=`$ $`k_\mathrm{b}\mathrm{log}_2W`$ (6)
$`S`$ $`=`$ $`k_\mathrm{o}\mathrm{log}_8W`$ (7)
$`S`$ $`=`$ $`k_\mathrm{d}\mathrm{log}_{10}W,`$ (8)
where $`k_\mathrm{b}=[\mathrm{log}2]=k_\mathrm{B}\mathrm{ln}2`$, $`k_\mathrm{o}=[\mathrm{log}8]=k_\mathrm{B}\mathrm{ln}8`$, and $`k_\mathrm{d}=[\mathrm{log}10]=k_\mathrm{B}\mathrm{ln}10`$ are respectively binary, octal, and decimal entropy units, corresponding to the bases of the logarithms used. Of course, any of the other logarithmic units listed in the previous section (and a continuum of other units as well) could also have been used in Boltzmannโs relation, with a suitable choice of logarithm base.
Observe now that the relations $`k_\mathrm{B}=[\mathrm{log}e]`$ and $`k_\mathrm{B}1.38\times 10^{23}\mathrm{J}/\mathrm{K}`$ imply that $`1\mathrm{K}1.38\times 10^{23}\mathrm{J}/[\mathrm{log}e]`$, in other words, the Kelvin (or any temperature unit) is fundamentally just an expression of an amount of energy per logarithmic unit of some arbitrary size. Here it is expressed as energy per nat or Neper, where this unit quantifies the increase in the indefinite logarithm of the number of states when the state count is multiplied by $`e`$. Of course, we could equally well express the Kelvin in terms of logarithmic units of other sizes as well, for example, multiplying top and bottom by $`(\mathrm{ln}10)`$ gives $`1\mathrm{K}3.18\times 10^{23}\mathrm{J}/[\mathrm{log}10]`$, where we see we have now expressed the Kelvin in units of Joules per decade (Bel, order of magnitude, power of ten etc.) of increase in the number of states.
Indeed, the modern thermodynamic definition of temperature is indeed just
$$T=\mathrm{d}Q/\mathrm{d}S,$$
(9)
where $`\mathrm{d}Q`$ is the amount of heat that must be added to a system in order to increase its entropy $`S=[\mathrm{log}W]`$ by a small amount $`\mathrm{d}S`$, or in other words to increase its number of states by the multiplicative factor $`[\mathrm{exp}\mathrm{d}S]`$, where note we are using our indefinite exponential notation from earlier.
## 7 Logarithmic Units and Information
Just as with entropy, the amount of information content or information capacity $`I`$ of a system can be expressed very elegantly and generically as an indefinite logarithmic quantity, that is, as
$$I=[\mathrm{log}W]$$
(10)
where $`W`$ is again the number of ways of arranging the system, or a subsystem of it whose state can be controlled, e.g., for purposes of storing or communicating a message.
The only real distinction between entropy and information is a distinction of epistemological status.
Entropy is usually taken to refer to that part of the physical information that is *unknown*, or in other words is not included in the available overall description of a physical situation; this information is not considered part of the so-called โmacrostateโ of the system.
The word โinformation,โ on the other hand, is sometimes reserved to implicitly connote information that is (or could be) explicitly known, although this more restricted usage is becoming less common. More and more, physicists who use the word โinformationโ understand that physical information, in general, could have the status of being either known information, or unknown information (entropy). One word that has been proposed to indicate known information as opposed to entropy, but which is not yet very popular, is *extropy*, which was coined to serve as a complement to the word โentropy.โ
Of course, due to the (historically dominant) use of binary codes in our modern digital systems, information has been traditionally measured in $`[\mathrm{log}2]`$ units (bits), or multiples thereof (such as bytes), rather than in $`[\mathrm{log}e]`$ units (nats) or $`[\mathrm{log}10]`$ units (decades). However, as we have been emphasizing, this difference in the conventional choice of units does not at all imply that information is fundamentally a different kind of quantity from the logarithmic quantities that are used in other contexts.
In fact, we argue that the only distinction between information/entropy and other types of logarithmic quantities is that information is a logarithmic quantity derived from an absolute pure number that represents a โnumber of alternativesโ in some sense, while most other logarithmic quantities (e.g. octaves, decibels, Richter scale points) are derived from pure numbers representing ratios between physical quantities (pitch, signal power, Earthquake strength). But fundamentally, although the sources of the pure numbers $`x`$ in the two kinds of cases are different, this does not matter; the values of $`[\mathrm{log}x]`$ are always still fundamentally the very *same* type of mathematical object.
Thus, a bit of information *is*, mathematically, the very *same* kind of object as an octave of pitch. A Boltzmannโs constant unit of entropy is the same kind of object as a Naper of current ratio. A decimal-digit-sized quantity of information is the *same* as a Richter-scale point of relative earthquake strength. Fundamentally, the only import of the different names for these mathematical objects is to connote their use in describing different types of situations.
Just as the number 2 is still a 2 whether we are talking about two giraffes or two potatoes, likewise the indefinite-logarithm object $`[\mathrm{log}2]`$ is still a $`[\mathrm{log}2]`$ unit whether we are talking about a $`[\mathrm{log}2]`$ amount of information (called a bit) or a $`[\mathrm{log}2]`$ size of musical interval (called an octave). And a $`[\mathrm{log}e]`$ unit is still a $`[\mathrm{log}e]`$ unit, whether we are talking about a $`[\mathrm{log}e]`$ unit of thermodynamic entropy (called Boltzmannโs constant) or a $`[\mathrm{log}e]`$ unit of voltage ratios (called a Neper). A $`[\mathrm{log}10]`$ unit is still a $`[\mathrm{log}10]`$ unit, whether we are talking about a $`[\mathrm{log}10]`$ unit of signal power ratio (called a Bel) or a $`[\mathrm{log}10]`$ unit of information (called a decimal digit).
In other words, logarithmic units are logarithmic units, and logarithmic quantities (expressed by a real number times a logarithmic unit) are mathematically always the very same kind of entity, no matter the domain. The only differences are in the size of the standard units that are conventionally used in a given context, the names that we call them, and how we apply them.
## 8 Discussion
Today, thanks to Boltzmann and his followers, we know that a certain quantity that used to be thought of as โphysical,โ namely entropy, is really just a โmathematicalโ quantity, namely an indefinite logarithmic quantity derived from the number of states, or in other words a kind of information. Since indefinite logarithmic units are scale-free, there is no natural unit or โatomโ of entropy or information that we must use, only units that we choose rather arbitrarily, by convention or for mathematical or technological convenience, such as the nat (Boltzmannโs constant) or the bit.
In the future, it is possible that we might discover that other quantities that are currently thought of as โphysicalโ could at root turn out to really be logarithmic quantities as well. For example, Tommaso Toffoli has speculated that, just as entropy turned out to be equivalent to the logarithm of the number of possible states that a system could (statically) be in, perhaps *energy* could be shown in some way to really be equivalent a logarithm of the number of possible *computations* that a system could carry out dynamically at the microscale. As of this writing, this intriguing idea is still rather far from being substantiated, but the history of how our understanding of the quantity of entropy has evolved indeed makes Toffoliโs proposal seem like an idea worth exploring.
Besides entropy and (perhaps) energy, one wonders whether other kinds of physical quantities such as distances and times might potentially be shown to ultimately be logarithmic quantities as well. That this might be true is hinted at by the Bekenstein-Hawking formula for the entropy of a black hole, $`S=A/4`$, where $`A`$ is the holeโs event horizon area in Planck units, and $`S`$ is the entropy in nats. Thus, for example, we could assign the Planck unit of length to be the square root of a nat, $`\mathrm{}_\mathrm{P}=[\mathrm{log}e]^{1/2}`$, and then write $`A=[\mathrm{log}W^4]`$, where $`W`$ is the number of states of the black hole, and this would be consistent with the Bekenstein relation as well as the entropy relation $`S=[\mathrm{log}W]`$. However, in this line of thought, it remains obscure why the area should be the indefinite logarithm of the fourth power of the number of states, and why the length unit should have dimensions of a square root of a logarithmic unit. Still, this may be an interesting line to pursue further.
## 9 Conclusion
In this paper, we have reviewed a well-defined mathematical concept of an indefinite logarithm function in which no particular logarithm base is selected, and have shown that the entities returned by this function can be used as the basis for a system of mathematical quantities that is exactly analogous in its behavior to systems of dimensioned physical quantities. In fact, this mathematical system of indefinite logarithmic quantities exactly corresponds the physical quantity known as entropy, when the logarithms are applied to the number of distinguishable physical states that are consistent with a given abstract description of the system. This quantity (the indefinite logarithm of the number of states) is also called โinformationโ in a slightly broader context.
In other words, physical (thermodynamic) entropy really *is* nothing but (unknown) information in the physical state, and its quantity really *is* nothing other than the indefinite logarithm of the state count. Further, Boltzmannโs constant $`k_\mathrm{B}`$ is really nothing other than a representation (in conventional physical units of energy over temperature) of the specific (and arbitrarily chosen) abstract indefinite-logarithm unit $`[\mathrm{log}e]`$, which is known as the nat or the Neper in other contexts. And, thermodynamic temperature really is nothing but the energy per logarithmic unit, for small increments in the indefinite logarithm of the state count.
There are speculations that other quantities such as energy that we presently think of as being fundamentally โphysicalโ in nature (as opposed to mathematical) might (similarly to entropy) someday be revealed to be, at root, derived from logarithmic quantities of some sort.
Of course, if physics can someday be exactly described by mathematics, as most theoretical physicists believe (or at least hope), then ultimately, the entire distinction between mathematical and physical quantities becomes somewhat of an artificial and illusory one, since we cannot then rule out the possibility that our physical universe may really be nothing but a particular (very elaborate) mathematical structure, one in which we (and our thought processes) happen to be embedded. What is a โphysicalโ quantity then ultimately becomes only a question of which mathematical quantities happen to arise naturally within the context of the particular mathematical structures that make up our physical universe.
To conclude, although there are probably no substantive ideas in this paper that have not been said many times before, somewhere in the literature (though perhaps in different terms), and although many of these ideas would likely be considered self-evident to professional mathematicians, we nevertheless felt that many of these ideas lack exposure at present within certain communities, and that it would be worthwhile to present and explain them again, so as to facilitate the more widespread understanding of these issues. We hope that this paper serves that purpose, at least.
Note: The reference list below is still under construction. The author would appreciate receiving from readers suggested references to appropriate prior sources that discuss these or similar ideas, so that he can cite the sources in future versions of this paper, as well as in future papers on related topics.
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# Probing the density dependence of the symmetry potential at low and high densities
## Acknowledgments
Q. Li thanks the Alexander von Humboldt-Stiftung for a fellowship. This work is partly supported by the National Natural Science Foundation of China under Grant No. 10255030, by GSI, BMBF, DFG, and VolkswagenStiftung.
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# Differential equations satisfied by modular forms and K3 surfaces
## 1. Introduction
Lian and Yau studied arithmetic properties of mirror maps of pencils of certain $`K3`$ surfaces, and further, they considered mirror maps of certain families of CalabiโYau threefolds . Lian and Yau observed in a number of explicit examples a mysterious relationship (now the so-called mirror moonshine phenomenon) between mirror maps and the McKayโThompson series (Hauptmoduls of one variable associated to a genus zero congruence subgroup of $`SL_2()`$) arising from the Monster. Inspired by the work of Lian and Yau, VerrillโYui further computed more examples of mirror maps of one-parameter families of lattice polarized $`K3`$ surfaces with Picard number $`19`$. The outcome of VerrillโYuiโs calculations suggested that the mirror maps themselves are not always Hauptmoduls, but they are commensurable with Hauptmoduls (referred as the modularity of mirror maps). This fact was indeed established by Doran for $`M_n`$-lattice polarized $`K3`$ surfaces of Picard number $`19`$ with maximal unipotent monodromy (where $`M_n=U(E_8)^22n`$). More generally, Doran considered the commensurability of โmaximal $`n`$-dimensional families of rank $`20n`$ lattice polarized families of $`K3`$ surfaces, and he showed that all such families of $`K3`$ surfaces are commensurable to autormorphic forms.
The mirror maps were calculated via the PicardโFuchs differential equations of the $`K3`$ families in question. Therefore, the determination of the PicardโFuchs differential equations played the central role in their investigations.
In this paper, we will address the inverse problem of a kind. That is, instead of starting with families of $`K3`$ surfaces or families of CalabiโYau threefolds, we start with modular forms and functions of more than one variable associated to certain subgroups of $`SL_2()`$.
More specifically, the main focus our discussions in this paper are on modular forms and functions of two variables. Here is the precise definition.
###### Definition 1.1.
Let $``$ denote the upper half-plane $`\{\tau :\mathrm{}\tau >0\}`$, and let $`^{}=\{\mathrm{}\}`$. Let $`\mathrm{\Gamma }_1`$ and $`\mathrm{\Gamma }_2`$ be two subgroups of $`SL_2()`$ commensurable with $`SL_2()`$. We call a function $`F:^{}\times ^{}`$ of two variables a modular form (of two variables) of weight $`(k_1,k_2)`$ on $`\mathrm{\Gamma }_1\times \mathrm{\Gamma }_2`$ with character $`\chi `$ if $`F`$ is meromorphic on $`^{}\times ^{}`$ such that
$$F(\gamma _1\tau _1,\gamma _2\tau _2)=\chi (\gamma _1,\gamma _2)(c_1\tau _1+d_1)^{k_1}(c_2\tau _2+d_2)^{k_2}F(\tau _1,\tau _2)$$
for all
$$\gamma _1=\left(\begin{array}{cc}a_1& b_1\\ c_1& d_1\end{array}\right)\mathrm{\Gamma }_1,\gamma _2=\left(\begin{array}{cc}a_2& b_2\\ c_2& d_2\end{array}\right)\mathrm{\Gamma }_2.$$
If $`F`$ is a modular form (of two variables) of weight $`(0,0)`$ with trivial character, then we also call $`F`$ a modular function (of two variables) on $`\mathrm{\Gamma }_1\times \mathrm{\Gamma }_2`$.
Notation. We let $`q_1=e^{2\pi i\tau _1}`$ and $`q_2=e^{2\pi i\tau _2}`$. For a variable $`t`$ we let $`D_t`$ denote the the differential operator $`t\frac{}{t}`$.
###### Remark 1.1.
Stienstra and Zagier have introduced the notion of bi-modular forms (of two variables). Let $`\mathrm{\Gamma }\text{SL}_2()`$, and let $`\tau _1,\tau _2`$. Let $`k_1,k_2`$ be integers. A two-variable meromorphic function $`F:\times `$ is called a bi-modular form of weight $`(k_1,k_2)`$ on $`\mathrm{\Gamma }`$ if for any $`\gamma =\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)\mathrm{\Gamma }`$, it satisfies the transformation formula:
$$F(\gamma \tau _1,\gamma \tau _2)=(c\tau _1+d)^{k_1}(c\tau _2+d)^{k_2}F(\tau _1,\tau _2).$$
For instance,
$$F(\tau _1,\tau _2)=\tau _1\tau _2$$
is a bi-modular form for $`\text{SL}_2()`$ of weight $`(1,1)`$. Another typical example is
$$F(\tau _1,\tau _2)=E_2(\tau _1)\frac{1}{\tau _1\tau _2},$$
which is a bi-modular form of weight $`(2,0)`$ for $`\text{SL}_2()`$.
For bi-modular forms of StienstraโZagier, the fundamental domain $`\times /\mathrm{\Gamma }`$ is not of finite volume. On the other hand, for our modular forms (of two variables), the fundamental domain is $`/\mathrm{\Gamma }_1\times /\mathrm{\Gamma }_2`$, which is always of finite volume.
We should emphasize that the two notions of two variable modular forms (namely, our modular forms and bi-modular forms of Stienstra and Zagier) are indeed different. Also we mention that our modular forms are not a special case of Hilbert modular forms.
The problems that we will consider here are formulated as follows : Given a modular form $`F`$ (of two variables), determine a differential equation it satisfies, and construct a family of $`K3`$ surfaces (or degenerations of a family of CalabiโYau threefolds at some limit points) having the determined differential equation as its PicardโFuchs differential equation. This kind of problem may be called a geometric realization problem.
In fact, a similar problem was already considered by Lian and Yau in their papers . They discussed the so-called โmodular relationsโ involving power series solutions to second and third order differential equations of Fuchsian type (e.g., hypergeometric differential equations $`{}_{2}{}^{}F_{1}^{},_3F_2`$) and modular forms of weight $`4`$ using mirror symmetry. More recently, van Enckevort and van Straten considered the following geometric realization problem: Starting with a certain forth order differential equation whose monodromy representation can be calculated, find a one-parameter families of CalabiโYau threefolds (if it exists), whose associated PicardโFuchs differential equation is the given one. Also a recent article of Doran and Morgan addressed the geometric realization question in the context of an old question of Griffiths: When does an integral variation of Hodge structure come from geometry?. A rigorous answer was presented for one-parameter families of CalabiโYau threefolds with $`h^{2,1}=1`$ with generalized PicardโFuch differential eqations, relating mirror symmetry and integral variations of Hodge structure.
In this paper, we will focus our discussion on modular forms (of two variables) of weight $`(1,1)`$. We will determine the differential equations satisfied by modular forms (of two variables) of weight $`(1,1)`$ associated to $`\mathrm{\Gamma }_1\times \mathrm{\Gamma }_2`$ where $`\mathrm{\Gamma }_i`$ are genus zero subgroups of $`SL_2()`$ of the form $`\mathrm{\Gamma }_0(N)`$ and $`\mathrm{\Gamma }_0(N)^{}`$. Then the existence and the construction of particular modular forms of weight $`(1,1)`$ are discussed, using solutions of some hypergeometric differential equations. Moreover, we determine the differential equations they satisfy. Further, several examples of modular forms (of two variables) and their differential equations are discussed aiming to realize these differential equations as the PicardโFuchs differential equations of some families of $`K3`$ surfaces (or degenerations of families of CalabiโYau threefolds) with large Picard numbers $`19,18,17`$ and $`16`$.
It should be pointed out that our paper and our results have non-empty intersections with the results of Lian and Yau . Indeed, our approach rediscovers some of the examples of Lian and Yau.
Our contributions may be summarized as follows. From geometric point of veiw, we give examples of two-parameter families of $`K3`$ surfaces which after pull-back along a morphism from $`(t_1,t_2)`$-space to $`(x,y)`$-space decouple as a direct product of two one-parameter families of elliptic curves. From function theoretic point of veiw, we give examples of non-trivial substitions transforming certain (two-variables) GKZ hypergeometic hypergeometric functions into a product of two (one-variable) GKZ hypergeometric fucntions. Finally, from moduli point of view, we give examples of moduli spaces for $`K3`$ surfaces with extra structure and show that these moduli spaces are quotients of $`\times `$.
## 2. Differential equations satisfied by modular forms (of two variables)
We will now determine differential equations satisfied by modular forms (of two variables) of weight $`(1,1)`$ on $`\mathrm{\Gamma }_1\times \mathrm{\Gamma }_2`$.
###### Theorem 2.1.
Let $`F(\tau _1,\tau _2)`$ be a modular form (of two variables) of weight $`(1,1)`$, and let $`x(\tau _1,\tau _2)`$ and $`y(\tau _1,\tau _2)`$ be non-constant modular functions (of two variables) on $`\mathrm{\Gamma }_1\times \mathrm{\Gamma }_2`$, where $`\mathrm{\Gamma }_i`$ ($`i=1,2`$) are subgroups of $`SL_2()`$ commensurable with $`SL_2()`$. Then $`F`$, as a function of $`x`$ and $`y`$, satisfy a system of partial differential equations
(2.1)
$$\begin{array}{c}\hfill D_x^2F+a_0D_xD_yF+a_1D_xF+a_2D_yF+a_3F=0,\\ \hfill D_y^2F+b_0D_xD_yF+b_1D_xF+b_2D_yF+b_3F=0,\end{array}$$
where $`a_i`$ and $`b_i`$ are algebraic functions of $`x`$ and $`y`$, and can be expressed explicitly as follows. Suppose that, for each function $`t`$ among $`F`$, $`x`$, and $`y`$, we let
$$G_{t,1}=\frac{D_{q_1}t}{t}=\frac{1}{2\pi i}\frac{dt}{td\tau _1},G_{t,2}=\frac{D_{q_2}t}{t}=\frac{1}{2\pi i}\frac{dt}{td\tau _2}.$$
Then we have
$$a_0=\frac{2G_{y,1}G_{y,2}}{G_{x,1}G_{y,2}+G_{y,1}G_{x,2}},b_0=\frac{2G_{x,1}G_{x,2}}{G_{x,1}G_{y,2}+G_{y,1}G_{x,2}},$$
$$a_1=\frac{G_{y,2}^2(D_{q_1}G_{x,1}2G_{F,1}G_{x,1})G_{y,1}^2(D_{q_2}G_{x,2}2G_{F,2}G_{x,2})}{G_{x,1}^2G_{y,2}^2G_{y,1}^2G_{x,2}^2},$$
$$b_1=\frac{G_{x,2}^2(D_{q_1}G_{x,1}2G_{F,1}G_{x,1})+G_{x,1}^2(D_{q_2}G_{x,2}2G_{F,2}G_{x,2})}{G_{x,1}^2G_{y,2}^2G_{y,1}^2G_{x,2}^2},$$
$$a_2=\frac{G_{y,2}^2(D_{q_1}G_{y,1}2G_{F,1}G_{y,1})G_{y,1}^2(D_{q_2}G_{y,2}2G_{F,2}G_{y,2})}{G_{x,1}^2G_{y,2}^2G_{y,1}^2G_{x,2}^2},$$
$$b_2=\frac{G_{x,2}^2(D_{q_1}G_{y,1}2G_{F,1}G_{y,1})+G_{x,1}^2(D_{q_2}G_{y,2}2G_{F,2}G_{y,2})}{G_{x,1}^2G_{y,2}^2G_{y,1}^2G_{x,2}^2},$$
$$a_3=\frac{G_{y,2}^2(D_{q_1}G_{F,1}G_{F,1}^2)G_{y,1}^2(D_{q_2}G_{F,2}G_{F,2}^2)}{G_{x,1}^2G_{y,2}^2G_{y,1}^2G_{x,2}^2},$$
and
$$b_3=\frac{G_{x,2}^2(D_{q_1}G_{F,1}G_{F,1}^2)+G_{x,1}^2(D_{q_2}G_{F,2}G_{F,2}^2)}{G_{x,1}^2G_{y,2}^2G_{y,1}^2G_{x,2}^2}.$$
In order to prove Theorem 2.1, we first need the following lemma, which is an analogue of the classical Ramanujanโs differential equations
$$D_qE_2=\frac{E_2^2E_4}{12}=24\underset{n}{}\frac{n^2q^n}{(1q^n)^2},$$
$$D_qE_4=\frac{E_2E_4E_6}{3}=240\underset{n}{}\frac{n^4q^n}{(1q^n)^2},$$
$$D_qE_6=\frac{E_2E_6E_4^2}{2}=\underset{n}{}\frac{n^6q^n}{(1q^n)^2}$$
where
(2.2)
$$E_k=1\frac{2k}{B_k}\underset{n}{}\frac{n^{k1}q^n}{1q^n}$$
are the Eisenstein series of weight $`k`$ on $`SL_2()`$, where $`B_k`$ denotes the $`k`$-th Bernoulli number, e.g., $`B_2=\frac{1}{6},B_4=\frac{1}{30}`$ and $`B_6=\frac{1}{42}`$.
###### Lemma 2.2.
We retain the notations of Theorem 2.1. Then
(a) $`G_{x,1}`$ and $`G_{y,1}`$ are modular forms (of two variables) of weight $`(2,0)`$,
(b) $`G_{x,2}`$ and $`G_{y,2}`$ are modular forms (of two variables) of weight $`(0,2)`$,
(c) $`D_{q_1}G_{x,1}2G_{F,1}G_{x,1}`$, $`D_{q_1}G_{y,1}2G_{F,1}G_{y,1}`$ and $`D_{q_1}G_{F,1}G_{F,1}^2`$ are modular forms (of two variables) of weight $`(4,0)`$, and
(d) $`D_{q_2}G_{x,2}2G_{F,2}G_{x,2}`$, $`D_{q_2}G_{y,2}2G_{F,2}G_{y,2}`$ and $`D_{q_2}G_{F,2}G_{F,2}^2`$ are modular forms (of two variables) of weight $`(0,4)`$.
###### Proof.
We shall prove (a) and (c); the proof of (b) and (d) is similar.
By assumption, $`x`$ is a modular function (of two variables) on $`\mathrm{\Gamma }_1\times \mathrm{\Gamma }_2`$. That is, for all $`\gamma _1=\left(\begin{array}{cc}a_1& b_1\\ c_1& d_1\end{array}\right)\mathrm{\Gamma }_1`$ and all $`\gamma _2=\left(\begin{array}{cc}a_2& b_2\\ c_2& d_2\end{array}\right)\mathrm{\Gamma }_2`$, one has
$$x(\gamma _1\tau _1,\gamma _2\tau _2)=x(\tau _1,\tau _2)$$
Taking the logarithmic derivatives of the above equation with respect to $`\tau _1`$, we obtain
$$\frac{1}{(c_1\tau _1+d_1)^2}\frac{\dot{x}}{x}(\gamma _1\tau _1,\tau _2)=\frac{\dot{x}}{x}(\tau _1,\tau _2),$$
or
(2.3)
$$G_{x,1}(\gamma _1\tau _1,\gamma _2\tau _2)=(c_1\tau _1+d_1)^2G_{x,1}(\tau _1,\tau _2),$$
where we let $`\dot{x}`$ denote the derivative of the two-variable function $`x`$ with respect to the first variable. This shows that $`G_{x,1}`$ is a modular form of weight $`(2,0)`$ on $`\mathrm{\Gamma }_1\times \mathrm{\Gamma }_2`$ with the trivial character. The proof for the case $`G_{y,1}`$ is similar.
Likewise, taking the logarithmetic derivatives of the equation
$$F(\gamma _1\tau _1,\gamma _2\tau _2)=\chi (\gamma _1,\gamma _2)(c_1\tau _1+d_1)(c_2\tau _2+d_2)F(\tau _1,\tau _2)$$
with respect to $`\tau _1`$, we obtain
$$\frac{1}{(c_1\tau _1+d_1)^2}\frac{\dot{F}}{F}(\gamma _1\tau _1,\gamma _2\tau _2)=\frac{c_1}{(c_1\tau _1+d_1)}+\frac{\dot{F}}{F}(\tau _1,\tau _2),$$
or, equivalently
(2.4)
$$G_{F,1}(\gamma _1\tau _1,\gamma _2\tau _2)=\frac{c_1(c_1\tau _1+d_1)}{2\pi i}+(c_1\tau _1+d_1)^2G_{F,1}(\tau _1,\tau _2).$$
Now, differentiating (2.3) with respect to $`\tau _1`$ again, we obtain
$$\frac{\dot{G}_{x,1}}{(c_1\tau _1+d_1)^2}(\gamma _1\tau _1,\gamma _2\tau _2)=2c_1(c_1\tau _1+d_1)G_{x,1}(\tau _1,\tau _2)+(c_1\tau _1+d_1)^2\dot{G}_{x,1}(\tau _1,\tau _2),$$
or
$$D_{q_1}G_{x,1}(\gamma _1\tau _1,\gamma _2\tau _2)=\frac{c_1(c_1\tau _1+d_1)^3}{\pi i}G_{x,1}(\tau _1,\tau _2)+(c_1\tau _1+d_1)^4D_{q_1}G_{x,1}(\tau _1,\tau _2).$$
On the other hand, we also have, by (2.3) and (2.4),
$$G_{F,1}G_{x,1}(\gamma _1\tau _1,\gamma _2\tau _2)=\frac{c_1(c_1\tau _1+d_1)^3}{2\pi i}G_{x,1}(\tau _1,\tau _2)+(c_1\tau _1+d_1)^4G_{F,1}G_{x,1}(\tau _1,\tau _2).$$
From these two equations we see that $`D_{q_1}G_{x,1}2G_{F,1}G_{x,1}`$ is a modular form (of two variables) of weight $`(4,0)`$ with the trivial character.
Finally, differentiating (2.4) with respect to $`\tau _1`$ and multiplying by $`(c_1\tau _1+d_1)^2`$ we have
$$\begin{array}{cc}\hfill D_{q_1}G_{F,1}(\gamma _1\tau _1,\gamma _2\tau _2)& =\frac{c_1^2(c_1\tau _1+d_1)^2}{(2\pi i)^2}+\frac{c_1(c_1\tau _1+d_1)^3}{\pi i}G_{F,1}(\tau _1,\tau _2)\hfill \\ & +(c_1\tau _1+d_1)^4D_{q_1}G_{F,1}(\tau _1,\tau _2).\hfill \end{array}$$
Combining this with the square of (2.4) we see that $`D_{q_1}G_{F,1}G_{F,1}^2`$ is a modular form of weight $`(4,0)`$ on $`\mathrm{\Gamma }_1\times \mathrm{\Gamma }_2`$. This completes the proof of the lemma. โ
###### Proof of Theorem 2.1.
In light of Lemma 2.2, the functions $`a_k`$, $`b_k`$ are all modular functions on $`\mathrm{\Gamma }_1\times \mathrm{\Gamma }_2`$, and thus can be expressed as algebraic functions of $`x`$ and $`y`$. Therefore, it suffices to verify (2.1) as formal identities. By the chain rule we have
$$\left(\begin{array}{c}D_{q_1}F\\ D_{q_2}F\end{array}\right)=\left(\begin{array}{cc}x^1D_{q_1}x& y^1D_{q_1}y\\ x^1D_{q_2}x& y^1D_{q_2}y\end{array}\right)\left(\begin{array}{c}D_xF\\ D_yF\end{array}\right).$$
It follows that
$$\left(\begin{array}{c}D_xF\\ D_yF\end{array}\right)=\frac{F}{G_{x,1}G_{y,2}G_{x,2}G_{y,1}}\left(\begin{array}{cc}G_{y,2}& G_{y,1}\\ G_{x,2}& G_{x,1}\end{array}\right)\left(\begin{array}{c}G_{F,1}\\ G_{F,2}\end{array}\right)$$
Writing
$$\mathrm{\Delta }=G_{x,1}G_{y,2}G_{x,2}G_{y,1},$$
and
$$\mathrm{\Delta }_x=G_{F,1}G_{y,2}G_{F,2}G_{y,1},\mathrm{\Delta }_y=G_{x,2}G_{F,1}+G_{x,1}G_{F,2},$$
we have
(2.5)
$$D_xF=F\frac{\mathrm{\Delta }_x}{\mathrm{\Delta }},D_yF=F\frac{\mathrm{\Delta }_y}{\mathrm{\Delta }}.$$
Applying the same procedure on $`D_xF`$ again, we obtain
$$\begin{array}{cc}\hfill \left(\begin{array}{c}D_x^2F\\ D_yD_xF\end{array}\right)& =\frac{1}{\mathrm{\Delta }}\left(\begin{array}{cc}G_{y,2}& G_{y,1}\\ G_{x,2}& G_{x,1}\end{array}\right)\left(\begin{array}{c}D_{q_1}(F\mathrm{\Delta }_x/\mathrm{\Delta })\\ D_{q_2}(F\mathrm{\Delta }_x/\mathrm{\Delta })\end{array}\right)\hfill \\ & =\frac{F}{\mathrm{\Delta }}\left(\begin{array}{cc}G_{y,2}& G_{y,1}\\ G_{x,2}& G_{x,1}\end{array}\right)\left\{\frac{\mathrm{\Delta }_x}{\mathrm{\Delta }}\left(\begin{array}{c}G_{F,1}\\ G_{F,2}\end{array}\right)+\left(\begin{array}{c}D_{q_1}(\mathrm{\Delta }_x/\mathrm{\Delta })\\ D_{q_2}(\mathrm{\Delta }_x/\mathrm{\Delta })\end{array}\right)\right\}.\hfill \end{array}$$
That is,
(2.6)
$$D_x^2F=F\frac{\mathrm{\Delta }_x^2}{\mathrm{\Delta }^2}+\frac{F}{\mathrm{\Delta }}\left(G_{y,2}D_{q_1}\frac{\mathrm{\Delta }_x}{\mathrm{\Delta }}G_{y,1}D_{q_2}\frac{\mathrm{\Delta }_x}{\mathrm{\Delta }}\right)$$
and
(2.7)
$$D_yD_xF=F\frac{\mathrm{\Delta }_x\mathrm{\Delta }_y}{\mathrm{\Delta }^2}+\frac{F}{\mathrm{\Delta }}\left(G_{x,2}D_{q_1}\frac{\mathrm{\Delta }_x}{\mathrm{\Delta }}+G_{x,1}D_{q_2}\frac{\mathrm{\Delta }_x}{\mathrm{\Delta }}\right).$$
We then substitute (2.5), (2.6), and (2.7) into (2.1) and find that (2.1) indeed holds. (The details are tedious, but straightforward calculations. We omit the details here.) โ
## 3. Modular forms (of two variables) associated to solutions of hypergeometric differential equations
Here we will construct modular forms (of two variables) of weight $`(1,1)`$ using solutions of some hypergeometric differential equations. Our main result of this section is the following theorem.
###### Theorem 3.1.
Let $`0<a<1`$ be a positive real number. Let $`f(t)=_2F_1(a,a;1;t)`$ be a solution of the hypergeometric differential equation
(3.1)
$$t(1t)f^{\prime \prime }+[1(1+2a)t]f^{}a^2f=0.$$
Let
$$F(t_1,t_2)=f(t_1)f(t_2)(1t_1)^a(1t_2)^a,$$
$$x=\frac{t_1+t_2}{(t_11)(t_21)},y=\frac{t_1t_2}{(t_1+t_2)^2}.$$
Then $`F`$ is a modular form of weight $`(1,1)`$ for $`\mathrm{\Gamma }_1\times \mathrm{\Gamma }_2`$, provided that $`t_1`$ and $`t_2`$ are modular functions (of one variable) for $`\mathrm{\Gamma }_1`$ and $`\mathrm{\Gamma }_2`$, respectively. Furthermore, $`F`$, as a function of $`x`$ and $`y`$, is a solution of the partial differential equations
(3.2)
$$D_x(D_x2D_y)F+x(D_x+a)(D_x+1a)F=0,$$
and
(3.3)
$$D_y^2Fy(2D_yD_x+1)(2D_yD_x)F=0,$$
where $`D_x=/x`$ and $`D_y=/y`$.
###### Remark 3.1.
Theorem 2.1 of Lian and Yau is essentially the same as our Theorem 3.1, though the formulation and proof are different.
The condition that $`t_1,t_2`$ are modular functions (of one variable) for $`\mathrm{\Gamma }_1`$ and $`\mathrm{\Gamma }_2`$ are used to draw the conclustion that $`F`$ is a modular form (of two variables) for $`\mathrm{\Gamma }_1\times \mathrm{\Gamma }_2`$. However, the modular property of $`t_1,t_2`$ is irrelevant to derive 3.2 and 3.3 from 3.1.
We will present our proof of Theorem 3.1 now. For this, we need one more ingredient, namely, the Schwarzian derivatives.
###### Lemma 3.2.
Let $`f(t)`$ and $`f_1(t)`$ be two linearly independent solutions of a differential equation
$$f^{\prime \prime }+p_1f^{}+p_2f=0.$$
Set $`\tau :=f_1(t)/f(t)`$. Then the associated Schwarzian differential equation
$$2Q\left(\frac{dt}{d\tau }\right)^2+\{t,\tau \}=0,$$
where $`\{t,\tau \}`$ is the Schwarzian derivative
$$\{t,\tau \}=\frac{dt^3/d\tau ^3}{dt/d\tau }\frac{3}{2}\left(\frac{dt^2/d\tau ^2}{dt/d\tau }\right)^2,$$
satisfies
$$Q=\frac{4p_22p_1^{}p_1^2}{4}.$$
###### Proof.
This is standard, and proof can be found, for instance, in Lian and Yau . โ
###### Proof of Theorem 3.1.
Let $`f_1`$ be another solution of (3.1) linearly independent of $`f`$, and set $`\tau =f_1/f`$. Then a classical identity asserts that
$$f^2=c\mathrm{exp}\left\{^t\frac{1(1+2a)u}{u(1u)}๐u\right\}\frac{dt}{d\tau }=\frac{cdt/d\tau }{t(1t)^{2a}},$$
where $`c`$ is a constant depending on the choice of $`f_1`$. Thus, letting
$$q_1=e^{2\pi if_1(t_1)/f(t_1)}\text{and}q_2=e^{2\pi if_1(t_2)/f(t_2)},$$
the function $`F`$, with a suitable choice of $`f_1`$, is in fact
$$F(t_1,t_2)=\left(\frac{D_{q_1}t_1D_{q_2}t_2}{t_1t_2}\right)^{1/2}.$$
We now apply the differential identities in (2.1), which hold for arbitrary $`F`$, $`x`$, and $`y`$. We have
$$G_{x,1}:=\frac{D_{q_1}x}{x}=\frac{(1+t_2)D_{q_1}t_1}{(t_1+t_2)(1t_1)},G_{x,2}:=\frac{D_{q_2}x}{x}=\frac{(1+t_1)D_{q_2}t_2}{(t_1+t_2)(1t_2)},$$
$$G_{y,1}:=\frac{D_{q_1}y}{y}=\frac{(t_2t_1)D_{q_1}t_1}{t_1(t_1+t_2)},G_{y,2}:=\frac{D_{q_2}y}{y}=\frac{(t_1t_2)D_{q_2}t_2}{t_2(t_1+t_2)},$$
$$G_{F,1}:=\frac{D_{q_1}F}{F}=\frac{t_1D_{q_1}^2t_1(D_{q_1}t_1)^2}{2t_1D_{q_1}t_1},G_{F,2}:=\frac{D_{q_2}F}{F}=\frac{t_2D_{q_2}^2t_2(D_{q_2}t_2)^2}{2t_2D_{q_2}t_2}.$$
It follows that
$$a_0:=\frac{2G_{y,1}G_{y,2}}{G_{x,1}G_{y,2}+G_{y,1}G_{x,2}}=\frac{2(t_11)(t_21)}{t_1t_2+1}=\frac{2}{1+x},$$
$$b_0:=\frac{2G_{x,1}G_{x,2}}{G_{x,1}G_{y,2}+G_{y,1}G_{x,2}}=\frac{2t_1t_2(t_1+1)(t_2+1)}{(t_1t_2)^2(t_1t_2+1)}=\frac{2y(1+2x)}{(1+x)(14y)},$$
$$\begin{array}{cc}\hfill a_1:& =\frac{G_{y,2}^2(D_{q_1}G_{x,1}2G_{F,1}G_{x,1})G_{y,1}^2(D_{q_2}G_{x,2}2G_{F,2}G_{x,2})}{G_{x,1}^2G_{y,2}^2G_{y,1}^2G_{x,2}^2}\hfill \\ & =\frac{t_1+t_2}{t_1t_2+1}=\frac{x}{1+x},\hfill \end{array}$$
$$\begin{array}{cc}\hfill b_1:& =\frac{G_{x,2}^2(D_{q_1}G_{x,1}2G_{F,1}G_{x,1})+G_{x,1}^2(D_{q_2}G_{x,2}2G_{F,2}G_{x,2})}{G_{x,1}^2G_{y,2}^2G_{y,1}^2G_{x,2}^2}\hfill \\ & =\frac{t_1t_2(t_1+1)(t_2+1)}{(t_1t_2)^2(t_1t_2+1)}=\frac{y(1+2x)}{(1+x)(14y)},\hfill \end{array}$$
$$a_2:=\frac{G_{y,2}^2(D_{q_1}G_{y,1}2G_{F,1}G_{y,1})G_{y,1}^2(D_{q_2}G_{y,2}2G_{F,2}G_{y,2})}{G_{x,1}^2G_{y,2}^2G_{y,1}^2G_{x,2}^2}=0,$$
$$\begin{array}{cc}\hfill b_2:& =\frac{G_{x,2}^2(D_{q_1}G_{y,1}2G_{F,1}G_{y,1})+G_{x,1}^2(D_{q_2}G_{y,2}2G_{F,2}G_{y,2})}{G_{x,1}^2G_{y,2}^2G_{y,1}^2G_{x,2}^2}\hfill \\ & =\frac{2t_1t_2}{(t_1t_2)^2}=\frac{2y}{14y}.\hfill \end{array}$$
Moreover, we have
$$\begin{array}{cc}\hfill a_3:& =\frac{G_{y,2}^2(D_{q_1}G_{F,1}G_{F,1}^2)G_{y,1}^2(D_{q_2}G_{F,2}G_{F,2}^2)}{G_{x,1}^2G_{y,2}^2G_{y,1}^2G_{x,2}^2}\hfill \\ & =\frac{(t_11)(t_21)(t_1+t_2)\left\{t_1^2\dot{t}_2^4(2\dot{t}_1\stackrel{\dot{}\dot{}\dot{}}{t}_13\ddot{t}_1^2)t_2^2\dot{t}_1^4(2\dot{t}_2\stackrel{\dot{}\dot{}\dot{}}{t}_23\ddot{t}_2^2)\right\}}{4(t_1t_2)(t_1^2t_2^21)\dot{t}_1^4\dot{t}_2^4},\hfill \end{array}$$
where, for brevity, we let $`\dot{t}_j`$, $`\ddot{t}_j`$, $`\stackrel{\dot{}\dot{}\dot{}}{t}_j`$ denote the derivatives $`D_{q_j}t_j`$, $`D_{q_j}^2t_j`$, and $`D_{q_j}^3t_j`$, respectively. To express $`a_3`$ in terms of $`x`$ and $`y`$, we note that, by Lemma 3.2,
$$\begin{array}{cc}\hfill 2\dot{t}_j\stackrel{\dot{}\dot{}\dot{}}{t}_j3\ddot{t}_j^2& =\dot{t}_j^4\left(\frac{4a^2}{t_j(1t_j)}2\frac{d}{dt_j}\frac{1(1+2a)t_j}{t_j(1t_j)}\frac{(1(1+2a)t_j)^2}{t_j^2(1t_j)^2}\right)\hfill \\ & =\frac{(t_j1)^2+4a(1a)t_j}{t_j^2(t_j1)^2}\dot{t}_j^4.\hfill \end{array}$$
It follows that
$$a_3=a(1a)\frac{t_1+t_2}{t_1t_2+1}=\frac{a(1a)x}{1+x}.$$
Likewise, we have
$$\begin{array}{cc}\hfill b_3:& =\frac{G_{x,2}^2(D_{q_1}G_{F,1}G_{F,1}^2)+G_{x,1}^2(D_{q_2}G_{F,2}G_{F,2}^2)}{G_{x,1}^2G_{y,2}^2G_{y,1}^2G_{x,2}^2}\hfill \\ & =a(1a)\frac{t_1t_2(t_1+t_2)}{(t_1t_2)^2(t_1t_2+1)}=\frac{a(1a)xy}{(1+x)(14y)}.\hfill \end{array}$$
Then, by (2.1), the function $`F`$, as a function of $`x`$ and $`y`$, satisfies
(3.4)
$$D_x^2F\frac{2}{1+x}D_xD_yF+\frac{x}{1+x}D_xF+\frac{a(1a)x}{1+x}F=0$$
and
(3.5)
$$\begin{array}{cc}& D_y^2F+\frac{2y(1+2x)}{(1+x)(14y)}D_xD_yF+\frac{y(1+2x)}{(1+x)(14y)}D_xF\hfill \\ & \frac{2y}{14y}D_yF+\frac{a(1a)xy}{(1+x)(14y)}F=0.\hfill \end{array}$$
Finally, we can deduce the claimed differential equations by taking (3.4) times $`(1+x)`$ and (3.5) times $`(14y)`$ minus (3.4) times $`y`$, respectively. โ
## 4. Examples
###### Example 4.1.
Let $`j`$ be the elliptic modular $`j`$-function, and let $`E_4(\tau )=1+240_n\frac{n^3q^n}{1q^n}`$, $`q=e^{2\pi i\tau }`$, be the Eisenstein series of weight $`4`$ on $`SL_2()`$. Set
$$x=2\frac{1/j(\tau _1)+1/j(\tau _2)1728/(j(\tau _1)j(\tau _2))}{1+\sqrt{(11728/j(\tau _1))(11728/j(\tau _2))}},y=\frac{1}{j(\tau _1)j(\tau _2)x^2},$$
and
$$F=(E_4(\tau _1)E_4(\tau _2))^{1/4}.$$
Then $`F`$ satisfies the system of partial differential equations:
$$(1432x)D_x^2F2D_xD_yF432xD_x60xF=0,$$
$$(14y)D_y^2F+4yD_xD_yFyD_x^2FyD_xF2yD_yF=0.$$
We should remark that the functions $`x`$ and $`y`$ are modular functions (of two variables) for $`\mathrm{\Gamma }_1\times \mathrm{\Gamma }_2`$ where $`\mathrm{\Gamma }_1=\mathrm{\Gamma }_2`$ is a subgroup of $`SL_2()`$ of index $`2`$. On the other hand, in the sense of Stienstra-Zagier, $`x`$ and $`y`$ are bi-modular functions for the group $`SL_2()`$ (cf. Remark 1.1).
We have noticed that this system of differential equation belongs to a general class of partial differential equations which involve solutions of hypergeometric hypergeometric differential equations discussed in Theorem 3.1.
Here we will prove the assertion of Example 4.1 using Theorem 3.1.
###### Proof of Example 4.1.
We first make a change of variable $`x\overline{x}/432`$. For convenience, we shall denote the new variable $`\overline{x}`$ by $`x`$ again. Thus, we are required to show that the functions
$$x=864\frac{1/j(\tau _1)+1/j(\tau _2)1728/(j(\tau _1)j(\tau _2))}{1+\sqrt{(11728/j(\tau _1))(11728/j(\tau _2))}},y=\frac{432^2}{j(\tau _1)j(\tau _2)x^2},$$
and $`F=(E_4(\tau _1)E_4(\tau _2))^{1/4}`$ satisfy
$$(1+x)D_x^2F2D_xD_yF+xD_x+\frac{5}{36}xF=0,$$
and
$$(14y)D_y^2F+4yD_xD_yFyD_x^2FyD_xF2yD_yF=0.$$
For brevity, we let $`j_1`$ denote $`j(\tau _1)`$ and $`j_2`$ denote $`j(\tau _2)`$. We now observe that the function $`x`$ can be alternatively expressed as
$$\begin{array}{cc}\hfill x& =864\frac{1/j_1+1/j_21728/(j_1j_2)}{1(11728/j_1)(11728/j_2)}\left(1\sqrt{(11728/j_1)(11728/j_2)}\right)\hfill \\ & =\frac{1}{2}\left(\sqrt{(11728/j_1)(11728/j_2)}1\right).\hfill \end{array}$$
Setting
$$t_1=\frac{\sqrt{11728/j_1}1}{\sqrt{11728/j_1}+1},t_2=\frac{\sqrt{11728/j_2}1}{\sqrt{11728/j_2}+1},$$
we have
$$x=\frac{t_1+t_2}{(t_11)(t_21)}.$$
Moreover, the functions $`j_k`$, written in terms of $`t_k`$, are $`j_k=432(t_k1)^2/t_k`$ for $`k=1,2`$. It follows that
$$y=\frac{432^2}{j_1j_2x^2}=\frac{t_1t_2}{(t_1+t_2)^2}.$$
In view of Theorem 3.1, setting
$$t=\frac{\sqrt{11728/j(\tau )}1}{\sqrt{11728/j(\tau )}+1}$$
it remains to show that the function $`f(t)=E_4(\tau )^{1/4}(1t)^{1/6}`$ is a solution of the hypergeometric differential equation
$$t(1t)f^{\prime \prime }+(1+4t/3)f^{}\frac{1}{36}f=0,$$
or equivalently, that
$$\frac{E_4(\tau )^{1/4}}{(1t)^{1/6}}=_2F_1(1/6,1/6;1;t).$$
This, however, follows from the classical identity
$$E_4(\tau )^{1/4}=_2F_1(\frac{1}{12},\frac{5}{12};1;\frac{1728}{j(\tau )})$$
and Kummerโs transformation formula
$$\begin{array}{cc}& \left(\frac{1+\sqrt{1z}}{2}\right)_2^{2a}F_1(a,b;a+b+\frac{1}{2};z)\hfill \\ & =_2F_1(2a,ab+\frac{1}{2};a+b+\frac{1}{2};\frac{\sqrt{1z}1}{\sqrt{1z}+1}).\hfill \end{array}$$
This completes the proof of Example 4.1. โ
###### Remark 4.1.
The functions $`x`$ and $`y`$ in Example 4.1 (up to constant multiple) have also appeared in the paper of Lian and Yau , Corollary 1.2, as the mirror map of the family of $`K3`$ surfaces defined by degree $`12`$ hypersurfaces in the weighted projective space $`^3[1,1,4,6]`$. Further, this $`K3`$ family is derived from the square of a family of elliptic curves in the weighted projective space $`^2[1,2,3]`$. (The geometry behind this phenomenon is the so-called ShoidaโInose structures, which has been studied in detail by Long for one-parameter families of $`K3`$ surfaces, and their PicardโFuchs differential equations.) Lian and Yau proved that the mirror map of the $`K3`$ family can be given in terms of the elliptic $`j`$-function, and indeed, by the functions $`x`$ and $`y`$ (up to constant multiple). We will discuss more examples of families of $`K3`$ surfaces, their PicardโFuchs differential equations and mirror maps in the section 6.
Along the same vein, we obtain more examples of modular forms of weight $`(1,1)`$ and modular functions on $`\mathrm{\Gamma }_0(N)\times \mathrm{\Gamma }_0(N)`$ for $`N=2,3,4`$.
###### Theorem 4.1.
We retain the notations of Theorem 3.1. Then the solutions of the differential equations (3.2) and (3.3) for the cases $`a=1/2,1/3,1/4,1/6`$ can be expressed in terms of modular forms and modular functions on $`\mathrm{\Gamma }_0(N)\times \mathrm{\Gamma }_0(N)`$ for some $`N`$.
(a) For $`a=1/2`$, they are given by
$$F(\tau _1,\tau _2)=\theta _4(\tau _1)^2\theta _4(\tau _2)^2,t=\theta _2(\tau )^4/\theta _3(\tau )^4,$$
which are modular on $`\mathrm{\Gamma }_0(4)\times \mathrm{\Gamma }_0(4)`$.
(b) For $`a=1/3`$, they are
$$F(\tau _1,\tau _2)=\frac{1}{2}(3E_2(3\tau _1)E_2(\tau _1))^{1/2}(3E_2(3\tau _2)E_2(\tau _2))^{1/2},t=27\frac{\eta (3\tau )^{12}}{\eta (\tau )^{12}},$$
which are modular on $`\mathrm{\Gamma }_0(3)\times \mathrm{\Gamma }_0(3)`$.
(c) For $`a=1/4`$, they are
$$F(\tau _1,\tau _2)=(2E_2(2\tau _1)E_2(\tau _1))^{1/2}(2E_2(2\tau _2)E_2(\tau _2))^{1/2},t=64\frac{\eta (2\tau )^{24}}{\eta (\tau )^{24}},$$
which are modular are $`\mathrm{\Gamma }_0(2)\times \mathrm{\Gamma }_0(2)`$.
(d) For $`a=1/6`$, they are given as in Example 4.1.
Here
$$\eta (\tau )=q^{1/24}\underset{n}{}(1q^n),q=e^{2\pi i\tau }$$
is the Dedekind eta-function, and
$$\theta _2(\tau )=q^{1/4}\underset{n}{}q^{n(n+1)},\theta _3(\tau )=\underset{n}{}q^{n^2},\theta _4(\tau )=\underset{n}{}(1)^nq^{n^2}$$
are theta-series.
###### Lemma 4.2.
Let $`\mathrm{\Gamma }`$ be a subgroup of $`SL_2()`$ commensurable with $`SL_2()`$. Let $`f(\tau )`$ be a modular form (of one variable) of weight $`1`$, and $`t(\tau )`$ be a non-constant modular function (of one variable) on $`\mathrm{\Gamma }`$. Then, setting
$$G_t=\frac{D_qt}{t},G_f=\frac{D_qf}{f},$$
we have
$$D_t^2f+\frac{D_qG_t2G_fG_t}{G_t^2}D_tf\frac{D_qG_fG_f^2}{G_t^2}f=0.$$
###### Proof of Theorem 4.1.
To prove part (a) we use the well-known identities
$$\theta _3^2=_2F_1(\frac{1}{2},\frac{1}{2};1;\frac{\theta _2^4}{\theta _3^4})$$
(see for a proof using Lemma 4.2) and
$$\theta _3^4=\theta _2^4+\theta _4^4.$$
Applying Theorem 3.1 and observing that
$$\theta _3^2\left(1\frac{\theta _2^4}{\theta _3^4}\right)^{1/2}=\theta _3^2\frac{\theta _4^2}{\theta _3^2}=\theta _4^2,$$
we thus obtain the claimed differential equation.
For parts (b), we need to show that the function
$$f(\tau )=\frac{(3E_2(3\tau )E_2(\tau ))^{1/2}}{(1t)^{1/3}}$$
satisfies
$$t(1t)\frac{d^2}{dt^2}f+(15t/3)\frac{d}{dt}f\frac{1}{9}f=0,$$
or, equivalently,
(4.1)
$$(1t)D_t^2f\frac{2}{3}tD_tf\frac{1}{9}tf=0.$$
Let $`G_t`$ and $`G_f`$ be defined as in Lemma 4.2. For convenience we also let $`g=(3E_2(3\tau )E_2(\tau ))/2`$. We have
$$G_t=\frac{1}{2}(3E_2(3\tau )E_2(\tau ))=g$$
and
$$G_f=\frac{D_qg}{2g}\frac{1}{3(1t)}D_qt=\frac{D_qg}{2g}+\frac{t}{3(1t)}g.$$
It follows that
$$\frac{D_qG_t2G_fG_t}{G_t^2}=g^2\left(D_qg2\left(\frac{D_qg}{2g}+\frac{t}{3(1t)}g\right)g\right)=\frac{2t}{3(1t)}.$$
Moreover, we can show that $`(D_qG_fG_f^2)/G_t^2`$ is equal to $`t/(9(1t))`$ by comparing enough Fourier coefficients. This establishes (4.1) and hence part (b).
The proof of part (c) is similar, and we shall skip the details here. โ
## 5. More examples
We may also consider groups like $`\mathrm{\Gamma }_0(N)^{}\times \mathrm{\Gamma }_0(N)^{}`$ where $`\mathrm{\Gamma }_0(N)^{}`$ denotes the group generated by $`\mathrm{\Gamma }_0(N)`$ and the AtkinโLehner involution $`w_N=\left(\begin{array}{cc}0& 1\\ N& 0\end{array}\right)`$ for some $`N`$. (Note that $`\mathrm{\Gamma }_0(N)^{}`$ is contained in the normalizer of $`\mathrm{\Gamma }_0(N)`$ in $`SL_2()`$.) Also the entire list of $`N`$ giving rise to genus zero groups $`\mathrm{\Gamma }_0(N)^{}`$ is known (cf.), and we will be interested in some of those genuz zero groups. We can determine differential equations satisfied by modular forms (of two variables) of weight $`(1,1)`$ on $`\mathrm{\Gamma }_0(N)^{}\times \mathrm{\Gamma }_0(N)^{}`$ for some $`N`$ (giving rise to genus zero subgroups $`\mathrm{\Gamma }_0(N)^{}`$).
We first prove a generalization of Theorem 3.1.
###### Theorem 5.1.
Let $`0<a,b<1`$ be positive real numbers. Let $`f(t)=_2F_1(a,b;1;t)`$ be a solution of the hypergeometric differential equation
(5.1)
$$t(1t)f^{\prime \prime }+[1(1+a+b)t]f^{}abf=0.$$
Set
$$F(t_1,t_2)=f(t_1)f(t_2)(1t_1)^{(a+b)/2}(1t_2)^{(a+b)/2},$$
$$x=t_1+t_22,y=(1t_1)(1t_2).$$
Then $`F`$, as a function of $`x`$ and $`y`$, satisfies
(5.2)
$$D_x^2F+2D_xD_yF\frac{1}{x+y+1}D_xF+\frac{x}{x+y+1}D_yF+\frac{(2abab)x}{2(x+y+1)}F=0$$
and
(5.3)
$$\begin{array}{cc}& D_y^2F+\frac{2y}{x^2}D_xD_yF+\frac{y^2}{x^2(x+y+1)}D_xF+\frac{yxx^2}{x(x+y+1)}D_yF\hfill \\ & \frac{(a+b)(a+b2)(x^2+x)+(ab)^2xy(4ab2a2b)y}{4x(x+y+1)}F=0.\hfill \end{array}$$
###### Proof.
The proof is very similar to that of Theorem 3.1. Let $`f_1`$ be another solution of the hypergeometric differential equation (5.1), and set $`\tau :=f_1/f`$. We find
$$f^2=c\mathrm{exp}\left\{^t\frac{1(1+a+b)u}{u(1u)}๐u\right\}\frac{dt}{d\tau }=\frac{cdt/d\tau }{t(1t)^{a+b}}$$
for some constant $`c`$ depending on the choice of $`f_1`$. Thus, setting
$$q_1=e^{2\pi if_1(t_1)/f(t_1)}\text{and}q_2=e^{2\pi if_1(t_2)/f(t_2)},$$
we have
$$F(t_1,t_2)=c^{}\left(\frac{D_{q_1}t_1D_{q_2}t_2}{t_1t_2}\right)^{1/2}$$
for some constant $`c^{}`$. We now apply the differential identities (2.1). We have, for $`j=1,2`$,
$$G_{x,j}:=\frac{D_{q_j}x}{x}=\frac{D_{q_j}t_j}{t_1+t_22},G_{y,j}:=\frac{D_{q_j}y}{y}=\frac{D_{q_j}t_j}{1t_j},$$
and
$$G_{F,j}:=\frac{D_{q_j}F}{F}=\frac{t_jD_{q_j}^2t_j(D_{q_j}t_j)^2}{2t_jD_{q_j}t_j}.$$
It follows that the coefficients in (2.1) are
$$a_0=2,b_0=\frac{2(1t_1)(1t_2)}{(t_1+t_22)^2}=\frac{2y}{x^2},$$
$$a_1=\frac{1}{t_1t_2}=\frac{1}{x+y+1},b_1=\frac{(1t_1)^2(1t_2)^2}{t_1t_2(t_1+t_22)^2}=\frac{y^2}{x^2(x+y+1)},$$
$$a_2=\frac{t_1+t_22}{t_1t_2}=\frac{x}{x+y+1},$$
$$b_2=\frac{t_1^2+t_1t_2+t_2^22t_12t_2+1}{t_1t_2(t_1+t_22)}=\frac{yxx^2}{x(x+y+1)}.$$
Moreover, we have
$$\begin{array}{cc}\hfill a_3& =\left\{\frac{(1t_1)^2(2\dot{t}_1\stackrel{\dot{}\dot{}\dot{}}{t}_13\ddot{t}_1^2)}{4(t_1t_2)\dot{t}_1^4}+\frac{(1t_2)^2(2\dot{t}_2\stackrel{\dot{}\dot{}\dot{}}{t}_23\ddot{t}_2^2)}{4(t_1t_2)\dot{t}_2^4}\frac{2t_1t_2t_1t_2}{4t_1^2t_2^2}\right\}\hfill \\ & \times (t_1+t_22),\hfill \end{array}$$
where we, as before, employ the notations $`\dot{t}_j`$, $`\ddot{t}_j`$, $`\stackrel{\dot{}\dot{}\dot{}}{t}_j`$ for the derivatives $`D_{q_j}t_j`$, $`D_{q_j}^2t_j`$, and $`D_{q_j}^3t_j`$, respectively. Now, by Lemma 3.2, we have
$$2\dot{t}_j\stackrel{\dot{}\dot{}\dot{}}{t}_j3\ddot{t}_j^2=\dot{t}_j^4\frac{(ab)^2t_j^2(1t_j)^2+(4ab2a2b)t_j}{t_j^2(1t_j)^2}.$$
It follows that
$$a_3=\frac{(2abab)(t_1+t_22)}{2t_1t_2}=\frac{(2abab)x}{x+y+1}.$$
A similar calculation shows that
$$b_3=\frac{(a+b)(a+b2)(x^2+x)+(ab)^2xy(4ab2a2b)y}{4x(x+y+1)}.$$
This proves the claimed result. โ
###### Remark 5.1.
It should be pointed out that the first identity in our proof of Theorem 5.1 is equivalent to the formula in Proposition 4.4 of Lian and Yau .
We now obtain new examples of modular forms of weight $`(1,1)`$ on $`\mathrm{\Gamma }_0(N)^{}\times \mathrm{\Gamma }_0(N)^{}`$ for some $`N`$.
###### Theorem 5.2.
When the pairs of numbers $`(a,b)`$ in Theorem 5.1 are given by $`(1/12,5/12)`$, $`(1/12,7/12)`$, $`(1/8,3/8)`$, $`(1/8,5/8)`$, $`(1/6,1/3)`$, $`(1/6,2/3)`$, $`(1/4,1/4)`$ and $`(1/4,3/4)`$, the solutions $`F(t_1,t_2)`$ of the differential equations (5.2) and (5.3) are modular forms of weight $`(1,1)`$ on $`\mathrm{\Gamma }_0(N)^{}\times \mathrm{\Gamma }_0(N)^{}`$ with $`N=1,1,2,2,3,3,4,4`$, respectively.
###### Proof.
We shall prove only the cases $`(a,b)=(1/6,1/3)`$ and $`(1/6,2/3)`$; the other cases can be proved in the same manner.
Let
$$s(\tau )=27\frac{\eta (3\tau )^{12}}{\eta (\tau )^{12}},E_2(\tau )=124\underset{n=1}{\overset{\mathrm{}}{}}\frac{nq^n}{1q^n}.$$
From the proof of Part (b) of Theorem 4.1 we know that
$$f(\tau )=\frac{(3E_2(3\tau )E_2(\tau ))^{1/2}}{(1s)^{1/3}},$$
as a function of $`s`$, is equal to $`\sqrt{2}_2F_1(1/3,1/3;1;s)`$. Now, applying the quadratic transformation formula
$${}_{2}{}^{}F_{1}^{}(\alpha ,\beta ;\alpha \beta +1;x)=(1x)_2^\alpha F_1(\frac{\alpha }{2},\frac{1+\alpha }{2}\beta ;\alpha \beta +1;\frac{4x}{(1x)^2})$$
for hypergeometric functions (see, for example \[1, Theorem 3.1.1\]) with $`\alpha =\beta =1/3`$, we obtain
$$(3E_2(3\tau )E_2(\tau ))^{1/2}=\sqrt{2}_2F_1(\frac{1}{6},\frac{1}{3};1;\frac{4s}{(1s)^2}).$$
Observing that the action of the Atkin-Lehner involution $`w_3`$ sends $`s`$ to $`1/s`$, we find that the function $`s/(1s)^2`$ is modular on $`\mathrm{\Gamma }_0(3)^{}`$. This proves that $`F(t_1,t_2)`$ is a modular form of weight $`(1,1)`$ for $`\mathrm{\Gamma }_0(3)^{}\times \mathrm{\Gamma }_0(3)^{}`$ in the case $`(a,b)=(1/6,1/3)`$.
Furthermore, an application of another hypergeometric function identity
$${}_{2}{}^{}F_{1}^{}(\alpha ,\beta ;\gamma ;x)=(1x)_2^\alpha F_1(\alpha ,\gamma \beta ;\gamma ;\frac{x}{x1})$$
yields
$$(3E_2(3\tau )E_2(\tau ))^{1/2}=\sqrt{2}\left(\frac{1s}{1+s}\right)_2^{1/3}F_1(\frac{1}{6},\frac{2}{3};1;\frac{4s}{(1+s)^2}).$$
This corresponds to the case $`(a,b)=(1/6,2/3)`$. Again, the function $`4s/(1+s)^2`$ is modular on $`\mathrm{\Gamma }_0(3)^{}`$. This implies that $`F(t_1,t_2)`$ is a modular form of weight $`(1,1)`$ for $`\mathrm{\Gamma }_0(3)^{}\times \mathrm{\Gamma }_0(3)^{}`$ for the case $`(a,b)=(1/6,2/3)`$. โ
###### Remark 5.2.
For the remaining pairs $`(a,b)`$ in Theorem 5.2, we simply list the exact expressions of $`F(t_1,t_2)`$ in terms of modular forms as proofs are similar.
For $`(a,b)=(1/12,5/12)`$ and $`(1/12,7/12)`$, they are
$$\left(\frac{E_6(\tau _1)E_6(\tau _2)}{E_4(\tau _1)E_4(\tau _2)}\right)^{1/2},and\left(\frac{E_8(\tau _1)E_8(\tau _2)}{E_6(\tau _1)E_6(\tau _2)}\right)^{1/2},$$
respectively, where $`E_k`$ are the Eisenstein series in (2.2).
For $`(a,b)=(1/8,3/8)`$ and $`(1/8,5/8)`$, they are
$$\underset{j=1}{\overset{2}{}}\left(\frac{1+s_j}{1s_j}(2E_2(2\tau _j)E_2(\tau _j))\right)^{1/2},and\underset{j=1}{\overset{2}{}}\left(\frac{1s_j}{1+s_j}(2E_2(2\tau _j)E_2(\tau _j))\right)^{1/2},$$
respectively, where $`s_j=64\eta (\tau _j)^{24}/\eta (\tau _j)^{24}`$.
For $`(a,b)=(1/6,1/3)`$ and $`(1/6,2/3)`$, they are
$$\underset{j=1}{\overset{2}{}}\left(\frac{1+s_j}{1s_j}(3E_2(3\tau _j)E_2(\tau _j))\right)^{1/2},and\underset{j=1}{\overset{2}{}}\left(\frac{1s_j}{1+s_j}(3E_2(3\tau _j)E_2(\tau _j))\right)^{1/2},$$
respectively, where $`s_j=27\eta (3\tau _j)^{12}/\eta (\tau _j)`$.
For $`(a,b)=(1/4,1/4)`$ and $`(1/4,3/4)`$, they are
$$\underset{j=1}{\overset{2}{}}\left(2E_2(2\tau _j)E_2(\tau _j)\right)^{1/2},and\underset{j=1}{\overset{2}{}}\left(2E_2(2\tau _j)E_2(\tau _j)\right)^{1/2}\frac{1s_j}{1+s_j},$$
respectively, where $`s_j=\theta _2(\tau _j)^4/\theta _3(\tau _j)^4`$.
## 6. PicardโFuchs differential equations of Familes of $`K3`$ surfaces : Part I
One of the motivations of our investigation is to understand the mirror maps of families of $`K3`$ surfaces with large Picard nubmers, e.g., $`19,18,17`$ or $`16`$. Some examples of such families of $`K3`$ surfaces were discussed in LianโYau , HosonoโLianโYau and also in Verrill-Yui . Some of $`K3`$ families occured considering degenerations of CalabiโYau families.
Our goal here is to construct families of $`K3`$ surfaces whose PicardโFuchs differential equations are given by the differential equations satisfied by modular forms (of two variables) we constructed in the earlier sections. In this section, we will look into the families of $`K3`$ surfaces appeared in Lian and Yau .
Let $`S`$ be a $`K3`$ surface. We recall some general theory about $`K3`$ surfaces which are relevant to our discussion. We know that
$$H^2(S,)(E_8)^2U^3$$
where $`U`$ is the hyperbolic plane $`\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)`$ and $`E_8`$ is the even unimodular negative definite lattice of rank $`8`$. The Picard group of $`S`$, $`\text{Pic}(S)`$, is the group of linear equivalence classes of Cartier divisors on $`S`$. Then $`\text{Pic}(S)`$ injects to $`H^2(X,)`$, and the image of $`\text{Pic}(S)`$ is the algebraic cycles in $`H^2(S,)`$. As $`Pic(S)`$ is torsion-free, it may be regarded as a lattice in $`H^2(S,)`$, called the Picard lattice, and its rank is denoted by $`\rho (S)`$.
According to ArnoldโDolgachev , two $`K3`$ surfaces form a mirror pair $`(S,\widehat{S})`$ if
$$\text{Pic}(S)_{H^2(S,)}^{}=\text{Pic}(\widehat{S})U\text{as lattices}$$
In terms of ranks, a mirror pair $`(S,\widehat{S})`$ is related by the identity:
$$22\rho (S)=\rho (\widehat{S})+2\rho (S)+\rho (\widehat{S})=20.$$
###### Example 6.1.
We will be interested in mirror pairs of $`K3`$ surfaces $`(S,\widehat{S})`$ whose Picard lattices are of the form
$$Pic(S)=UandPic(\widehat{S})=U_2(E_8)^2.$$
We go back to our Example 4.1, and discuss geometry behind that example. Associated to this example, there is a family of $`K3`$ surfaces in the weighted projective $`3`$-space $`^3[1,1,4,6]`$ with weight $`(q_1,q_2,q_3,q_4)=(1,1,4,6)`$. There is a mirror pair of $`K3`$ surfaces $`(S,\widehat{S})`$. Here we know (cf. Belcastro ) that
$$\text{Pic}(S)=U\text{so that }\rho (S)=2,$$
and that $`S`$ has a mirror partner $`\widehat{S}`$ whose Picard lattice is given by
$$\text{Pic}(\widehat{S})=U(E_8)^2\text{ so that }\rho (\widehat{S})=18.$$
The mirror $`K3`$ family can be defined by a hypersurface in the orbifold ambient space $`^3[1,1,4,6]/G`$ of degree $`12`$. Here $`G`$ is the discrete group of symmetry and can be given explicitly by $`G=(/3)\times (/2)=g_1\times g_2`$ where $`g_1,g_2`$ are generatoers whose actions are given by:
$$\begin{array}{ccc}g_1:(Y_1,Y_2,Y_3,Y_4)& & (\zeta _3Y_1,Y_2,\zeta _3^1Y_3,Y_4)\\ g_2:(Y_1,Y_2,Y_3,Y_4)& & (Y_1,Y_2,Y_3,Y_4)\end{array}$$
(Here $`\zeta _3=e^{2\pi i/3}`$.) The $`G`$-invariant monomials are
$$Y_1^{12},Y_2^{12},Y_3^3,Y_4^2,Y_1^6Y_2^6,Y_1Y_2Y_3Y_4.$$
The matrix of exponents is the following $`6\times 5`$ matrix
$$\left(\begin{array}{ccccc}12& 0& 0& 0& 1\\ 0& 12& 0& 0& 1\\ 0& 0& 3& 0& 1\\ 0& 0& 0& 2& 1\\ 6& 6& 0& 0& 1\\ 1& 1& 1& 1& 1\end{array}\right)$$
whose rank is $`2`$. Therefore we may conclude that the typical $`G`$-invariant polynomials is in $`2`$-parameters, and $`\widehat{S}`$ can be defined by the following $`2`$-parameter family of hypersurfaces of degree $`12`$
$$Y_1^{12}+Y_2^{12}+Y_3^3+Y_4^2+\lambda Y_1Y_2Y_3Y_4+\varphi Y_1^6Y_2^6=0$$
in $`^3[1,1,4,6]/G`$ with parameters $`\lambda `$ and $`\varphi `$.
How do we conmpute the PicardโFuchs differential equation of this $`K3`$ family?
Several physics articles are devoted to this question. For instance, KlemmโLercheโMayr , HosonoโKlemmโTheisenโYau , Lian and Yau determined the PicardโFuchs differential equation of the CalabiโYau family using the GKZ hypergeometric system. Also it was noticed (cf. , ) that the PicardโFuchs system of this family of $`K3`$ surfaces can be realized as the degeneration of the PicardโFuchs systems of the CalabiโYau family. The family of CalabiโYau threefolds is a degree $`24`$ hypersurfaces in $`^4[1,1,2,8,12]`$ with $`h^{1,1}=3`$. The defining equation for this family is given by
$$Z_1^{24}+Z_2^{24}+Z_3^{12}+Z_4^3+Z_5^212\psi _0Z_1Z_2Z_3Z_4Z_5$$
$$6\psi _1(Z_1Z_2Z_3)^6\psi _2(Z_1Z_2)^{12}=0.$$
Its PicardโFuchs system is given by
$$\begin{array}{ccc}L_1& =& \mathrm{\Theta }_x(\mathrm{\Theta }_x2\mathrm{\Theta }_z)12x(6\mathrm{\Theta }_x+5)(6\mathrm{\Theta }_x+1)\\ L_2& =& \mathrm{\Theta }_y^2y(2\mathrm{\Theta }_y\mathrm{\Theta }_z+1)(2\mathrm{\Theta }_y\mathrm{\Theta }_z)\\ L_3& =& \mathrm{\Theta }_z(\mathrm{\Theta }_z2\mathrm{\Theta }_y)z(2\mathrm{\Theta }_z\mathrm{\Theta }_x+1)(2\mathrm{\Theta }_z\mathrm{\Theta }_x)\end{array}$$
where
$$x=\frac{2\psi _1}{1728^2\psi _0^6},y=\frac{1}{\psi _2^2}\text{and}z=\frac{\psi _2}{4\psi _1^2}$$
are deformation coordinates.
Now the intersection of this CalabiโYau hypersurface with the hyperplane $`Z_2tZ_1=0`$ gives rise to a family of $`K3`$ surfaces
$$b_0Y_1Y_2Y_3Y_4+b_1Y_1^{12}+b_2Y_2^{12}+b_3Y_3^3+b_4Y_4^2+b_5Y_1^6Y_2^6=0$$
in $`^3[1,1,4,6]`$ of degree $`12`$. Taking $`(b_0,b_1,b_2,b_3,b_4,b_5)=(\lambda ,1,1,1,1,\varphi )`$ we obtain the $`2`$-parameter family of $`K3`$ surfaces described above. The PicardโFuchs system of this $`K3`$ family is obtained by taking the limit $`y=0`$ in the PicardโFuchs system for the CalabiโYau family:
$$\begin{array}{ccc}L_1& =& \mathrm{\Theta }_x(\mathrm{\Theta }_x2\mathrm{\Theta }_z)12x(6\mathrm{\Theta }_x+5)(6\mathrm{\Theta }_x+1)\\ L_3& =& \mathrm{\Theta }_z^2z(2\mathrm{\Theta }_z\mathrm{\Theta }_x+1)(2\mathrm{\Theta }_z\mathrm{\Theta }_x)\end{array}$$
Further, if we intersect this $`K3`$ family with the hyperplane $`Y_2sY_1=0`$, we obtain a family of elliptic curves:
$$c_0W_1W_2W_3+c_1W_1^6+c_2W_2^3+c_3W_3^2=0$$
in $`^2[1,2,3]`$, whose PicardโFuchs equation is given by
$$L=\mathrm{\Theta }_x^212x(6\mathrm{\Theta }_x+5)(6\mathrm{\Theta }_x+1).$$
Here we describe a relation of the PicardโFuchs system of the above family of $`K3`$ surfaces to the differential equation discussed in Example 4.1.
###### Remark 6.1.
We note that, in view of our proof of Example 4.1, the process of setting $`z=0`$ in the above PicardโFuchs system $`\{L_1,L_3\}`$ is equivalent to setting $`t_1=0`$ or $`t_2=0`$ in $`x`$ and $`y`$ in Example 4.1. Our Theorem 3.1 then implies that $`F(t)=(1t)_2^{1/6}F_1(1/6,1/6;1;t)`$ satisfies
$$(1+x)D_x^2F+xD_xF+\frac{5}{36}xF=0$$
with $`x=t/(1t)`$, or equivalently, (making a change of variable $`xx`$)
$$x(1x)F^{\prime \prime }+(12x)F^{}\frac{5}{36}F=0$$
with $`x=t/(t1)`$. That is,
$$(1t)_2^{1/6}F_1(\frac{1}{6},\frac{1}{6};1;t)=_2F_1(\frac{1}{6},\frac{5}{6};1;\frac{t}{t1}).$$
This is the special case of the hypergeometric series identity
$$(1t)_2^aF_1(a,b;c;t)=_2F_1(a,cb;c;\frac{t}{t1}).$$
We will discuss more examples of PicardโFuchs systems of CalabiโYau threefolds and $`K3`$ surfaces, which have already been considered by several people. For instance, the articles , , and obtained the PicardโFuchs operators for CalabiโYau hypersurfaces with $`h^{1,1}3`$. The next two examples consider CalabiโYau hypersurfaces with $`h^{1,1}>3`$, and the paper of Lian and Yau addressed the question of determining the PicardโFuchs system of the families of $`K3`$ surfaces $`^3[1,1,2,2]`$ of degree $`6`$ and $`^3[1,1,2,4]`$ of degree $`8`$. Their results are that
(1) there is an elliptic fibration on these $`K3`$ surfaces, and the PicardโFuchs systems of the $`K3`$ families can be derived from the PicardโFuchs system of the elliptic pencils, and that
(2) the solutions of the PicardโFuchs systems for the $`K3`$ families are given by โsquaresโ of those for the elliptic families.
The system of partial differential equations considered by Lian and Yau is
$$\begin{array}{ccc}L_1& =& \mathrm{\Theta }_x(\mathrm{\Theta }_x2\mathrm{\Theta }_z)\lambda x(\mathrm{\Theta }_x+\frac{1}{2}+\nu )(\mathrm{\Theta }_x+\frac{1}{2}\nu )\\ L_2& =& \mathrm{\Theta }_z^2z(2\mathrm{\Theta }_z\mathrm{\Theta }_x+1)(2\mathrm{\Theta }_z\mathrm{\Theta }_x)\end{array}$$
and an ordinary differential equations
$$L=\mathrm{\Theta }_x^2\lambda x(\mathrm{\Theta }_x+\frac{1}{2}+\nu )(\mathrm{\Theta }_x+\frac{1}{2}\nu )$$
where $`\mathrm{\Theta }_x=x\frac{}{x}`$, etc.) and $`\lambda ,\nu `$ are complex numbers.
Also they noted that the $`K3`$ families correspond, respecitvely, to the families of CalabiโYau threefolds $`^4[1,1,2,4,4]`$ of degree $`12`$ and $`^4[1,1,2,4,8]`$ of degree $`16`$. However, the PicardโFuchs systems for the CalabiโYau families are not explicitly determined.
###### Example 6.2.
We now consider a family of $`K3`$ surfaces $`^3[1,1,2,4]`$. of degree $`8`$. This $`K3`$ family is realized as the degeneration of the family of CalabiโYau hypersurfaces $`^4[1,1,2,4,8]`$ of degree $`16`$ and $`h^{1,1}=4`$. The most generic defining equation for this family is given by
$$a_0Z_1Z_2Z_3Z_4Z_5+a_1Z_1^{16}+a_2Z_2^{16}+a_3Z_3^8+a_4Z_4^4+a_5Z_5^2+a_6Z_3^2Z_4Z_5+a_7Z_1^8Z_2^8=0$$
Again the intersection with the hyperplane $`Z_2tZ_1=0`$ gives rise to a family of $`K3`$ surfaces $`^3[1,1,2,4]`$:
$$Y_1^8+Y_2^8+Y_3^4+Y_4^2+\lambda Y_1Y_2Y_3Y_4+\varphi Y_1^4Y_2^4=0$$
Let $`S`$ denote this family of $`K3`$ surfaces. Then
$$Pic(S)=M_{(1,1),(1,1),0}\text{with }\rho (S)=3.$$
The mirror family $`\widehat{S}`$ exists and its Picard lattice is
$$Pic(\widehat{S})=E_8D_7U\text{with }\rho (\widehat{S})=17.$$
The Picard lattices are determined by Belcastro . The intersection of this family of $`K3`$ surfaces with the hyperplane $`Y_2sY_2=0`$ gives rise to the pencil of elliptic curves
$$c_0W_1W_2W_3+c_1W_1^4+c_2W_2^4+c_3W_3^2=0$$
in $`^2[1,1,2]`$ of degree $`4`$. This means that this family of $`K3`$ surfaces has the elliptic fibration with section.
Now translate this โinductiveโ structure to the PicardโFuchs systems. The PicardโFuchs system for the $`K3`$ family is given by
$$\begin{array}{ccc}L_1& =& \mathrm{\Theta }_x(\mathrm{\Theta }_x2\mathrm{\Theta }_z)64x(\mathrm{\Theta }_x+\frac{1}{2}+\frac{1}{4})(\mathrm{\Theta }_x+\frac{1}{2}\frac{1}{4})\\ L_2& =& \mathrm{\Theta }_z^2z(2\mathrm{\Theta }_z\mathrm{\Theta }_x+1)(2\mathrm{\Theta }_z\mathrm{\Theta }_x)\end{array}$$
and the PicardโFuchs defferential equation of the elliptic family is given by
$$L=\mathrm{\Theta }_x^264x(\mathrm{\Theta }_x+\frac{1}{2}+\frac{1}{4})(\mathrm{\Theta }_x+\frac{1}{2}\frac{1}{4})$$
The same remark as Remark 6.1 is valid for the PicardโFuchs system $`\{L_1,L_3\}`$ which corresponds to Theorem 4.1 (b) with $`a=1/3`$.
###### Example 6.3.
We consider a family of $`K3`$ surfaces $`^3[1,1,2,2]`$ of degree $`6`$. This $`K3`$ family is realized as the degeneration of the family of CalabiโYau hypersurfaces $`^4[1,1,2,4,4]`$ of degree $`12`$ and $`h^{1,1}=5`$:
$$a_0Z_1Z_2Z_3Z_4Z_5+a_1Z_1^{12}+a_2Z_2^{12}+a_3Z_3^6+a+4Z_4^3+a_5Z_5^3+a_6Z_1^6Z_2^6=0.$$
The intersection of this CalabiโYau hypersurface with the hyperplane $`Z_2tZ_1=0`$ gives rise to the family of $`K3`$ hypersurfaces $`^3[1,1,2,4]`$:
$$Y_1^6+Y_2^6+Y_3^3+Y_4^3+\lambda Y_1Y_2Y_3Y_4+\varphi Y_1^3Y_2^3=0.$$
Let $`S`$ denote this family of $`K3`$ surfaces. Then
$$\text{Pic}(S)=M_{(1,1,1),(1,1,1),0}\text{with }\rho (S)=4.$$
There is a mirror family of $`K3`$ surfaces, $`\widehat{S}`$ with
$$\text{Pic}(\widehat{S})=E_8D_4A_2U\text{with }\rho (\widehat{S})=16.$$
The Picard lattices are determined by Belcastro .
The intersection of this $`K3`$ family with the hyperplane $`Y_2sY_1=0`$ gives rise to the family of elliptic curves
$$c_0W_1W_2W_3+c_1W_1^3+c_2W_2^3+c_3W_3^3=0$$
in $`^2[1,1,1]`$ of degree $`3`$.
The PicardโFuchs system of this $`K3`$ family is
$$\begin{array}{ccc}L_1& =& \mathrm{\Theta }_x(\mathrm{\Theta }_x2\mathrm{\Theta }_z)27x(\mathrm{\Theta }_x+\frac{1}{2}+\frac{1}{6})(\mathrm{\Theta }_x+\frac{1}{2}\frac{1}{6})\\ L_2& =& \mathrm{\Theta }_z^2z(2\mathrm{\Theta }_z\mathrm{\Theta }_x+1)(2\mathrm{\Theta }_z\mathrm{\Theta }_x)\end{array}$$
and the PicardโFuchs differential equation for the elliptic family is given by
$$L=\mathrm{\Theta }_x^227x(\mathrm{\Theta }_x+\frac{1}{2}+\frac{1}{6})(\mathrm{\Theta }_x+\frac{1}{2}\frac{1}{6})$$
We note that the same remark is valid for the PicardโFuchs system $`\{L_1,L_3\}`$ corresponding to $`a=1/4`$ in Theorem 4.1(c).
We will summarize the above discussions for the families of $`K3`$ surfaces in the following form.
###### Proposition 6.1.
The PicardโFuchs systems of families of $`K3`$ surfaces obtained by Lian and Yau can be reconstructed starting from the modular forms (of two variables) and then finding the differential equations satisfied by them. In other words, the differential equations satisfied by the modular forms (of two variables) are realized as the PicardโFuchs differential equations of the families of $`K3`$ surfaces, establishing, in a sense, the โmodularityโ of the $`K3`$ families.
## 7. PicardโFuchs differential equations of families of $`K3`$ surfaces: Part II
The purpose of this section is to study (one-parameter) families of $`K3`$ surfaces (some of which are realized as degenerations of some families of CalabiโYau threefolds), whose mirror maps are expressed in terms of Hauptmodules for genus zero subgroups of the form $`\mathrm{\Gamma }_0(N)^{}`$, aiming to identify their PicardโFuchs systems with differential equations assocaited to some to modular forms (of two variable) (e.g., in Theorem 5.1).
Dolgachev has discussed several examples of families of $`M_N`$-polarized $`K3`$ surfaces corresponding to $`\mathrm{\Gamma }_0(N)^{}`$ for small values of $`N`$, e.g., $`N=1,2`$ and $`3`$.
Lian and Yau have given examples of families of $`K3`$ surfaces and their PicardโFuchs differential equqtions of order $`3`$. The modular groups are genus zero subgroups of the form $`\mathrm{\Gamma }_0(N)^{}`$ where $`N`$ ranging from $`1`$ to $`30`$. Here we try to analyze their examples and their method in relation to our results in the section 5.
###### Example 7.1.
We start with the hypergeometric equation:
$$t(1t)f^{\prime \prime }+[1(1+a+b)t]f^{}abf=0$$
in Theorem 5.1. Take $`a=b=\frac{1}{4}`$ and consider a one-parameter deformation of this equation of the form:
$$t(1t)f^{\prime \prime }+(1\frac{3}{2}t)f^{}\frac{1}{16}(14\nu ^2)f=0$$
with a deformation parameter $`\nu `$. This has a unique solution $`f_0(t)`$ near $`t=0`$ with $`f_0(0)=1`$, and a solution $`f_1(t)`$ with $`f_1(t)=f_0(t)\text{log}t+O(t)`$. The inverse $`t(q)`$ of the power series $`q=\text{exp}(\frac{f_1(t)}{f_0(t)})=t+O(t^2)`$ defines an invertible holomorphic function in a disc, and $`t(q)`$ is the so-called mirror map. Put
$$x(q)=\frac{1}{\lambda }t(\lambda q)\text{for a given }\lambda .$$
One of the main results of Lian and Yau is that for any complex numbers $`\lambda ,\nu `$ with $`\lambda 0`$, there is a power series identity:
$${}_{3}{}^{}F_{2}^{}(\frac{1}{2},\frac{1}{2}+\nu ,\frac{1}{2}\nu ;1,1;\lambda x(q))^2=\frac{x^{\mathrm{\hspace{0.17em}2}}}{x^2(1\lambda x)}$$
in the common domain of definitions of both sides. As before, $`x^{}(q)=D_qx(q)`$.
For instance, take $`(\lambda ,\nu )=(2^63^3,\frac{1}{3}),(2^8,\frac{1}{4}),(2^23^3,\frac{1}{6})`$ and $`(2^6,0)`$, then these relations are given below. The mirror maps in these examples are expressed in terms of Hauptmodules of genus zero modular groups of the form $`\mathrm{\Gamma }_0(N)^{}`$ ($`\mathrm{\Gamma }_0(1)^{}=\mathrm{\Gamma }`$).
$$\begin{array}{ccccc}\text{Label}& \text{Modular Relation}& & & \text{Modular Group}\\ I& :\left(_{n=0}^{\mathrm{}}\frac{(6n)!}{(3n)!(n!)^3}\frac{1}{j(\tau )^n}\right)^2& =& E_4(q)& \mathrm{\Gamma }\\ II& :\left(_{n=0}^{\mathrm{}}\frac{(4n)!}{(n!)^4}x_2(\tau )^n\right)^2& =& \frac{x_2^{\mathrm{\hspace{0.17em}2}}}{x^2(1256x)}& \mathrm{\Gamma }_0(2)^{}\\ III& :\left(_{n=0}^{\mathrm{}}\frac{(2n)!(3n)!}{(n!)^5}x_3(\tau )^n\right)^2& =& \frac{x_3^{\mathrm{\hspace{0.17em}2}}}{x_3^2(1108x_3)}& \mathrm{\Gamma }_0(3)^{}\\ IV& :\left(_{n=0}^{\mathrm{}}\frac{(2n)!^3}{(n!)^6}x_4(\tau )^n\right)^2& =& \frac{x_4^{\mathrm{\hspace{0.17em}2}}}{x_4^2(164x_4)}& \mathrm{\Gamma }_0(4)^{}.\end{array}$$
Here $`j(\tau ),x_2(\tau ),x_3(\tau )`$ and $`x_4(\tau )`$ are Hauptmodules for the genus zero subgroups $`\mathrm{\Gamma },\mathrm{\Gamma }_0(2)^{},\mathrm{\Gamma }_0(3)^{}`$ and $`\mathrm{\Gamma }_0(4)^{}`$, respectively. Observe that in each modular relation, the right hand side is a modular form of weight $`4`$ on the corresponding genus zero subgroup.
We know that $`{}_{3}{}^{}F_{2}^{}(\frac{1}{2},\frac{1}{2}+\nu ,\frac{1}{2}\nu ;1,1;\lambda x)`$ is a unique solution with the leading term $`1+O(x)`$ to the differential operator
$$L=\mathrm{\Theta }_x^3\lambda x(\mathrm{\Theta }_x+\frac{1}{2})(\mathrm{\Theta }_x+\frac{1}{2}+\nu )(\mathrm{\Theta }_x+\frac{1}{2}\nu ).$$
In these examples, this differential operator is identified with the PicardโFuchs differential operator for a one-parameter family of $`K3`$ surfaces, which are obtained by degenerating CalabiโYau families. (Cf. Lian and Yau , Klemm, Lercher and Myer .)
$$\begin{array}{cccc}& \text{CY family}& K3\text{ family}& \text{PF Operator}\\ I& X(1,1,2,2,2)[8]& X(1,1,1,3)[6]& \mathrm{\Theta }^38x(6\mathrm{\Theta }+5)(6\mathrm{\Theta }+3)(6\mathrm{\Theta }+1)\\ II& X(1,1,2,2,6)[12]& X(1,1,1,1)[4]& \mathrm{\Theta }^34x(4\mathrm{\Theta }+3)(4\mathrm{\Theta }+2)(4\mathrm{\Theta }+1)\\ III& X(1,1,2,2,2,2)[6,4]& X(1,1,1,1,1)[3,2]& \mathrm{\Theta }^36x(2\mathrm{\Theta }+1)(3\mathrm{\Theta }+2)(3\mathrm{\Theta }+1)\\ IV& X(1,1,2,2,2,2,2)[4,4,4]& X(1,1,1,1,1,1)[2,2,2]& \mathrm{\Theta }^38x(2\mathrm{\Theta }+1)^3\end{array}$$
The $`K3`$ families I and II have already been discussed in LianโYau (see also VerrillโYui ) in relation to mirror maps. The Picard group of I (resp. II) is given by
$$(E_8)^2U_2<4>\text{(resp. }(E_8)^2U_2<2>).$$
The CalabiโYau family III can be realized as a complete intersection of the two hypersurfaces:
$$\begin{array}{c}Y_1^6+Y_2^6+Y_3^3+Y_4^3+Y_5^3+Y_6^3=0\\ Y_1^4+Y_2^4+Y_3^2+Y_4^2+Y_5^2+Y_6^2=0\end{array}$$
This CalabiโYau family has $`h^{1,2}=68`$ and $`h^{1,1}=2`$. The $`K3`$ family is realized as the fiber space by setting
$$Y_1=Z_1^{1/2},Y_2=\lambda Z_1^{1/2},\text{and}Y_i=Z_i\text{for }i=3,\mathrm{},6$$
where $`\lambda ^1`$ is a parameter. That is, we obtain a family of complete intersection $`K3`$ surfaces $`X(1,1,1,1,1)[3,2]`$:
$$\begin{array}{c}(1+\lambda ^6)Z_1^3+Z_3^3+Z_4^3+Z_5^3+Z_6^3=0\\ (1+\lambda ^4)Z_1^2+Z_3^2+Z_4^2+Z_5^2+Z_6^2=0\end{array}$$
Question: What is the Picard group of this K3 family?
In the similar manner, the CalabiโYau family IV can be realized as a complete intersection of the three hypersurfaces:
$$\begin{array}{c}Y_1^4+Y_2^4+Y_3^2+Y_4^2+Y_5^2+Y_6^2+Y_7^2=0\\ Z_1^4+Z_2^4+Z_3^2+Z_4^2+Z_5^2+Z_6^2+Z_7^2=0\\ W_1^4+W_2^4+W_3^2+W_4^2+W_5^2+W_6^2+W_7^2=0\end{array}$$
The $`K3`$ family is realized as the fiber space by setting
$$Y_1=Y_1^{\frac{1}{2}},Y_2=\lambda Y_1^{\frac{1}{2}}\text{and}Y_i=Y_i^{}\text{for }i=3,\mathrm{},7$$
and similarly for $`Z_1,Z_2`$ and $`W_1,W_2`$ where $`\lambda ^1`$ is a parameter.
This gives rise to the $`K3`$ family $`X(1,1,1,1,1,1)[2,2,2]`$:
$$\begin{array}{c}(1+\lambda ^4)Y_1^2+Y_3^2+Y_4^2+Y_5^2+Y_6^2+Y_7^{}=0\\ (1+\lambda ^4)Z_1^2+Z_3^2+Z_4^2+Z_5^2+Z_6^2+Z_7^2=0\\ (1+\lambda ^4)W_1^2+W_3^2+W_4^2+W_5^2+W_6^2+W_7^2=0\end{array}$$
Question: What is the Picard group of this K3 family?
Here is the summary:
(1) One starts with a Hauptmodule $`x(=x(q))`$ for a genus zero subgroup $`\mathrm{\Gamma }_0(N)^{}`$;
(2) then there associate a modular form $`\frac{x^{\mathrm{\hspace{0.17em}2}}}{xr(x)}`$ of weight $`4`$,
(3) and a power series solution $`\omega _0(x)`$ of an order three differential operator;
(4) this differential operator coincides with the PicardโFuchs differential operator of a one-parameter family of $`K3`$ surfaces in weighted projective spaces.
Lian and Yau further considered generalizations of the above phenomenon, constructing many more examples. Given a genus zero subgroup of the form $`\mathrm{\Gamma }_0(N)^{}`$ and a Hauptmodul $`x(q)`$, constract (by taking a Schwarzian derivative) a modular form $`E`$ of weight $`4`$ of the form $`\frac{x^{\mathrm{\hspace{0.17em}2}}}{xr(x)}`$ and a differential operator $`L`$ whose monodromy has maximal unipotency at $`x=0`$, such that $`LE^{1/2}=0`$. Further, identify $`L`$ as the PicardโFuchs differential operator of a family of $`K3`$ surfaces. Let $`\omega _0(x)`$ denotes the fundamental period of this manifold. Then it should be subject to the modular relation
$$\omega _0(x)^2=\frac{x^{\mathrm{\hspace{0.17em}2}}}{xr(x)}$$
How do we associate modular forms of weight $`(1,1)`$ corresponding to the groups $`\mathrm{\Gamma }_0(N)^{}\times \mathrm{\Gamma }_0(N)^{}`$ in this situation?
Taking the square root of both sides of the modular relation, we obtain that $`\omega _0(x)^{1/2}`$ is a modular form (of one variable) of weight $`1`$ for the group $`\mathrm{\Gamma }_0(N)^{}`$. Then taking $`\omega _0(q_1)\omega _0(q_2)`$, we see that this is a modular form for $`\mathrm{\Gamma }_0(N)^{}\times \mathrm{\Gamma }_0(N)^{}`$ of weight $`(1,1)`$. Then this modular form (of two variables) satisfies a differential equation, which may be identified with the PicardโFuchs differential equation of the K3 family considered above. We summarize the above discussion in the following proposition.
###### Proposition 7.1.
The examples IโIV above are related to our Theorem 5.2. Indeed, the connection is established by the identity
$${}_{2}{}^{}F_{1}^{}(a,b;a+b+\frac{1}{2};z)^2=_3F_2(2a,a+b,2b;a+b+\frac{1}{2},2a+2b;z).$$
More explicitly, the examples IโIV correspond to the cases $`(1/12,5/12)`$, $`(1/8,3/8)`$, $`(1/6,1/3)`$, and $`(1/4,1/4)`$, respectively.
Note that the generalized hypergeometric series $`{}_{3}{}^{}F_{2}^{}(\alpha _1,\alpha _2,\alpha _3;1,1;z)`$ satisfies the differential equation of the form:
$$[\mathrm{\Theta }_z^3\lambda z(\mathrm{\Theta }_z+\alpha _1)(\mathrm{\Theta }_z+\alpha _2)(\mathrm{\Theta }_z+\alpha _3)]f=0$$
for some $`\alpha _1,\alpha _2,\alpha _3`$ and $`\lambda ,0`$.
A natural question we may ask now is: Is is possible to construct families of $`K3`$ surfaces corresponding to Theorem 5.2 from this observation?
When the order $`3`$ differential equation of this form becomes the symmetric square of an order $`2`$ differential equation, and if the order $`2`$ differential equation is realized as the PicardโFuchs differential equation of a family of elliptic curves, we may be able to construct a family of $`K3`$ surfaces using the method of Long , especially when the Picard number of the $`K3`$ family in question is $`19`$ or $`20`$. In fact, RodriguezโVillegas has discussed $`4`$ families of $`K3`$ surfaces which fall into this class.
However, at the moment, we do not know if there are readily available methods for constructing $`K3`$ families starting from differential equations.
###### Remark 7.1.
If we consider the order $`4`$ generalized hypergeomtric series, there are $`14`$ differential equations are of the form
$$[\mathrm{\Theta }_z^4\lambda z(\mathrm{\Theta }_z+\alpha _1)(\mathrm{\Theta }_z+\alpha _2)(\mathrm{\Theta }_z+\alpha _3)(\mathrm{\Theta }_z+\alpha _4)]f=0$$
for some $`\alpha _i`$ and $`\lambda ,0`$. These $`14`$ differential operators can be found in AlmkvistโZudilin . As Doran and Morgan explained, only 13 of the 14 such operators are known to be realizable as the PicardโFuch differential operator for a family of smooth CalabiโYau threefolds with $`h^{2,1}=1`$. For the $`13`$ cases, Klemm and Theisen (see also Villegas ) found the corresponding families of CalabiโYau threefolds in weighted projective spaces. The missing case is $`(\alpha _1,\alpha _2,\alpha _3,\alpha _4)=(1/12,5/12,7/12,11/12)`$. For more thorough discussions on this topic, the reader should consult the article of Doran and Morgan , where classify integral monodromy representations.
## 8. Generalizations and open problems
###### Problem 1.
We have determined differential equations satisfied by modular forms (of two variables) of weight $`(1,1)`$. The arguments can be generalized to modular forms (of two variables) of any weight $`(k_1,k_2)`$, using the result of Yang . However, differerntial equations satisfied by them are getting too big to display.
###### Problem 2.
A natural generalization is to consider modular forms of three (or more than three) variables $`F(\tau _1,\tau _2,\tau _3)`$ of weight $`(k_1,k_2,k_3)`$ on $`\mathrm{\Gamma }_1\times \mathrm{\Gamma }_2\times \mathrm{\Gamma }_3`$.
Examples of this kind should correspond to PicardโFuchs differential equations of families of CalabiโYau threefolds, or PicardโFuchs differential equations of degenerate families of CalabiโYau fourfolds.
## Acknowledgments
The collaboration for this work started when Y. Yang visited N. Yui at Queenโs University, Kingston Canada, in August 2004. The work was completed at Tsuda College, Tokyo Japan 2005, and revisions were incorporated at Tsuda College again in the summer of 2006. Both authors were visiting professors at that institution in June 2005, and August 2006. We thank the hospitality of Tsuda College.
We thank Jan Stienstra for his interest and helpful comments and suggestions, and Chuck Doran and Don Zagier for their comments on the earlier version(s) of this paper.
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# The maximum size of ๐ฟ-functions
## 1. Introduction and statement of results
A fundamental problem in analytic number theory is to calculate the maximum size of $`L`$-functions in the critical strip. For example, the importance of the Lindelรถf Hypothesis, which is a consequence of the Riemann Hypothesis, is that it provides at least a crude estimate for the maximum in the case of the Riemann zeta-function. In this paper we use a variety of methods to conjecture the true rate of growth.
Consider first the Riemann zeta-function, which is a prototypical $`L`$-function. The Lindelรถf Hypothesis asserts that for every $`\epsilon >0`$, $`\zeta (\frac{1}{2}+\mathrm{i}t)=O(t^\epsilon )`$ (here we assume $`t`$ is positive). Under the Riemann Hypothesis, one can show that
$$\zeta (\frac{1}{2}+\mathrm{i}t)=O\left(\mathrm{exp}\left(C\frac{\mathrm{log}t}{\mathrm{log}\mathrm{log}t}\right)\right)$$
(1.1)
for some constant $`C`$ (see Theorem 14.14A of , for example). Several results of the form
$$|\zeta (\frac{1}{2}+\mathrm{i}t)|=\mathrm{\Omega }\left(\mathrm{exp}\left(C^{}\sqrt{\frac{\mathrm{log}t}{\mathrm{log}\mathrm{log}t}}\right)\right)$$
(1.2)
have also been established. Assuming the Riemann Hypothesis (RH), Montgomery showed $`C^{}1/20`$. Balasubramanian and Ramachandra improved the constant $`C^{}`$ and removed the assumption of RH. Soundararajan further improved the estimate to $`C^{}1`$, and he has also obtained similar results for the central values of a family of $`L`$-functions. In fact, Soundararajanโs calculations show that the proportion of $`t`$ for which $`\zeta (\frac{1}{2}+\mathrm{i}t)`$ is this big is quite large, suggesting that it may get bigger still. Numerical calculations of Kotnik indicate that $`C^{}>2`$ and perhaps $`C^{}`$ can be much larger.
We are interested in finding out which of equations (1.1) or (1.2) is closer to the truth. This is part of a class of problems that has recently come to be known as the โ1 or 2?โ question, where one has an $`O`$-result and an $`\mathrm{\Omega }`$-result which, suitably interpreted, differ by a factor of $`2`$. In the case here, the unknown factor is the power of $`\mathrm{log}t`$ in the exponential. The calculations in this paper support the view that โ1โ is the correct answer in this case, and we make the following conjecture:
###### Conjecture A.
$$\underset{t[0,T]}{\mathrm{max}}|\zeta (\frac{1}{2}+\mathrm{i}t)|=\mathrm{exp}\left((1+o(1))\sqrt{\frac{1}{2}\mathrm{log}T\mathrm{log}\mathrm{log}T}\right).$$
(1.3)
Similar arguments to those presented for $`\zeta (\frac{1}{2}+\mathrm{i}t)`$ work for $`S(t)`$, the error term in the number of zeros of the zeta-function up to height $`t`$, and lead to
###### Conjecture B.
$$\underset{t\mathrm{}}{lim\; sup}\frac{S(t)}{\sqrt{\mathrm{log}t\mathrm{log}\mathrm{log}t}}=\frac{1}{\pi \sqrt{2}}.$$
(1.4)
Much recent progress in understanding analytic properties of $`L`$-functions has come from the idea of a โfamilyโ of $`L`$-functions with an associated symmetry type. The idea is that to a collection of $`L`$-functions, with appropriate natural conditions, one can associate a classical compact group: unitary, symplectic, or orthogonal. One expects the analytic properties of the $`L`$-functions to be largely governed only by the symmetry type. Here we apply this philosophy to conjecture the maximal size of the critical values of $`L`$-functions.
A family $``$ of $`L`$-functions is partially ordered by the โconductorโ $`c(F)`$ for $`F`$. Our calculations assume that $`\mathrm{\#}\{F:c(F)<D\}D`$. Straightforward modifications can handle the case in which the family grows like $`D^A`$ for any $`A>0`$.
For a more detailed discussion of families of $`L`$-functions, see . (However, note that introduces a refined notion of โconductorโ, which, asymptotically, is the logarithm of the โusualโ conductor used here.)
###### Conjecture C.
Suppose $``$ is a family of $`L`$-functions and, for $`F`$, let $`c(F)`$ denote the conductor of $`F`$. With $`B=1/2`$ for unitary families and $`B=1`$ for symplectic and orthogonal families, we have
$$\underset{\begin{array}{c}F\\ c(F)D\end{array}}{\mathrm{max}}|F(\frac{1}{2})|=\mathrm{exp}\left((1+o(1))\sqrt{B\mathrm{log}D\mathrm{log}\mathrm{log}D}\right).$$
(1.5)
The implied constant depends only on $``$.
For example, for the symplectic family of real primitive Dirichlet $`L`$-functions, $`L(s,\chi _d)`$, where $`\chi _d=\left(\genfrac{}{}{0pt}{}{d}{}\right)`$, we conjecture that
$$\underset{|d|D}{\mathrm{max}}|L(\frac{1}{2},\chi _d)|=\mathrm{exp}\left((1+o(1))\sqrt{\mathrm{log}D\mathrm{log}\mathrm{log}D}\right).$$
(1.6)
Similarly, for the orthogonal family of Dirichlet series associated with holomorphic cusp forms, $`L(s,f)`$, where $`fS_k(\mathrm{\Gamma }_0(N))`$, the conductor is $`kN`$, so we conjecture that
$$\underset{\begin{array}{c}fS_k(\mathrm{\Gamma }_0(N))\\ kND\end{array}}{\mathrm{max}}|L(\frac{k}{2},f)|=\mathrm{exp}\left((1+o(1))\sqrt{\mathrm{log}D\mathrm{log}\mathrm{log}D}\right).$$
(1.7)
Note that Conjecture C contains Conjecture A, because any primitive $`L`$-function, $`L(s)`$, has associated with it the unitary family
$$_L:=\{F_y(s):=L(s+iy)|y\}$$
(1.8)
with conductor $`c(F_y)|y|`$.
Our conjectures suggest that on the critical line the answer to the โ1 or 2โ question is โ1โ. Work of Montgomery and Vaughan and Granville and Soundararajan has suggested that the answer also is โ1โ on the $`1`$-line. Thus in both cases, the maximum value the $`L`$-function attains appears to be closer to the $`\mathrm{\Omega }`$-result than to the $`O`$-result.
In Section 2 we use a rigorous approximation to the zeta-function due to Gonek, Hughes, and Keating to justify Conjecture A. This approximation represents $`\zeta (s)`$ as a product over primes times a product over zeros. We use characteristic polynomials of random unitary matrices to model the product over zeros, and a separate probabilistic model due to Granville and Soundararajan for the product over primes. The approximation to the zeta-function has a parameter which controls the relative contribution of the primes and the zeros. We show that the predicted maximal order of the zeta-function is the same independent of the choice of parameter. That is, whether we use only the primes, or only the zeros, or some combination of the two, we obtain Conjecture A.
In Section 3 we use recent conjectures for moments of $`L`$-functions to give an alternative justification for Conjecture A. Our approach also provides new limits on the range of validity of those conjectured moments.
In Section 4 we modify the treatment in Section 2 to obtain Conjecture B.
In Section 5 we describe how to extend our approach to obtain Conjecture C. We also describe some other approaches to obtaining these conjectures and then indicate possible arguments against the conjectures.
Finally, in Appendix A we prove a theorem about random matrix polynomials that is used in Section 2.
We thank Andrew Granville and Soundararajan for allowing us to incorporate their work on extreme values using the primes, and we thank Andrew Booker, Brian Conrey, Hugh Montgomery and Doug Ulmer for helpful conversations.
## 2. A probabilistic model for the zeta-function
Gonek, Hughes, and Keating have proved that if $`s=\sigma +\mathrm{i}t`$, with $`0\sigma 1`$ and $`|t|2`$, then for $`X>2`$ and $`K`$ any positive integer,
$$\zeta (s)=P_X(s)Z_X(s)\left(1+O\left(\frac{X^{2\sigma +K}}{(|t|\mathrm{log}X)^K}\right)+O(X^\sigma \mathrm{log}X)\right),$$
(2.1)
where
$$P_X(s):=\mathrm{exp}\left(\underset{nX}{}\frac{\mathrm{\Lambda }(n)}{n^s\mathrm{log}n}\right),$$
(2.2)
$`\mathrm{\Lambda }(n)`$ is the von-Mangoldt function, and
$$Z_X(s):=\mathrm{exp}\left(\underset{\rho }{}U((s\rho )\mathrm{log}X)\right).$$
(2.3)
Here the $`\rho `$ are non-trivial zeros of $`\zeta (s)`$ and $`U(z)=_0^{\mathrm{}}u(x)E_1(z\mathrm{log}x)dx`$, where $`E_1(z)=_z^{\mathrm{}}\frac{e^w}{w}dw`$ is the exponential integral and $`u`$ is any smooth function supported in $`[e^{11/X},e]`$.
The parameter $`X`$ controls the relative influence of the primes and the zeros. If $`X`$ is large, there are many primes in $`P_X(s)`$, and only the zeros very close to $`s`$ effect the product in $`Z_X(s)`$, while if $`X`$ is small, the zeros further away from $`s`$ make a contribution to $`Z_X(s)`$, but the number of primes in $`P_X(s)`$ is diminished. When $`X`$ is not too large, we expect $`Z_X`$ and $`P_X`$ to behave somewhat independently, and Gonek, Hughes, and Keating give evidence of this. In the remainder of this section we describe probabilistic models for $`P_X`$ and $`Z_X`$ which, assuming independence, will give Conjecture A. In Section 2.1 we describe our model for the large values of $`|Z_X|`$, establish some new results on the size of characteristic polynomials of random unitary matrices, and justify Conjecture A by choosing $`X`$ small. In Section 2.2 we describe Granville and Soundararajanโs model for $`P_X`$ and justify Conjecture A by choosing $`X`$ large. Then, in Section 2.3 we combine $`Z_X`$ and $`P_X`$, showing that intermediate values of $`X`$ also lead to Conjecture A.
### 2.1. A random matrix model for large values of $`Z_X(\frac{1}{2}+\mathrm{i}t)`$
Here we study the characteristic polynomial
$$\mathrm{\Lambda }_U(\theta )=det\left(IUe^{\mathrm{i}\theta }\right)=\underset{n=1}{\overset{N}{}}\left(1e^{\mathrm{i}(\theta _n\theta )}\right)$$
(2.4)
of a random unitary matrix $`UU(N)`$ chosen uniformly with respect to Haar measure. The characteristic polynomial $`\mathrm{\Lambda }_U(\theta )`$ was first developed as a model for the Riemann zeta-function by Keating and Snaith . In , building on , it is argued that for $`tT`$, $`Z_X(\frac{1}{2}+\mathrm{i}t)`$, given by (2.3), can be modeled by $`\mathrm{\Lambda }_U(\theta )`$, where $`UU(N)`$ with
$$N=\left[\frac{\mathrm{log}T}{e^\gamma \mathrm{log}X}\right].$$
(2.5)
We will prove a result about the value distribution and maximal size of $`|\mathrm{\Lambda }_U(\theta )|`$ in Appendix A, and use it to conjecture the distribution of large values of $`|Z_X|`$.
The largest value of $`|\mathrm{\Lambda }_U(\theta )|`$ is $`2^N`$, and values near this occur when $`U`$ is in a small neighborhood of scalar multiples of the identity matrix. If $`X=e^{o(\mathrm{log}\mathrm{log}T)}`$, this violates the known bound on $`|\zeta (\frac{1}{2}+\mathrm{i}t)|`$, so our model for the large values of $`|Z_X(\frac{1}{2}+\mathrm{i}t)|`$ must do something more subtle than just take the maximum of $`|\mathrm{\Lambda }_U|`$ over all $`UU(N)`$.
If $`T`$ and $`X`$ are thought of as fixed, then matrices of size $`N=[\frac{\mathrm{log}T}{e^\gamma \mathrm{log}X}]`$ should model the zeta-function as long as $`T/X^{e^\gamma }<t<T`$. If $`X>2`$, say, then up to constants there are $`T\mathrm{log}T`$ zeros in this interval. Therefore, in order to have the same number of eigenvalues, one needs
$$M=[T\mathrm{log}T/N]\mathrm{exp}\left(e^\gamma N\mathrm{log}X\right)\mathrm{log}X$$
(2.6)
matrices. Thus, one plausible guess for the maximum value of $`|Z_X(\frac{1}{2}+\mathrm{i}t)|`$ for $`0<t<T`$ is $`K=K(M,N)`$, where $`N`$ and $`M`$ are given in (2.5) and (2.6), respectively, and $`K`$ is the smallest possible function of $`M`$ and $`N`$ such that
$$\left\{\underset{1jM}{\mathrm{max}}\underset{\theta }{\mathrm{max}}|\mathrm{\Lambda }_{U_j}(\theta )|K\right\}1$$
(2.7)
as $`N\mathrm{}`$. Such a $`K`$ is found in Theorem 2.1.
We have glossed over some issues here, but we argue that they are not essential. First, we have claimed that matrices of size $`N`$ model $`Z_X(\frac{1}{2}+\mathrm{i}t)`$ for $`T/X^{e^\gamma }<t<T`$, whereas we want $`0<t<T`$. However, if $`X\mathrm{}`$, then $`[T/X^{e^\gamma },T]`$ will cover almost all of $`[0,T]`$, and so should capture the maximum. Secondly, we have been slightly cavalier in dropping the condition that $`N`$ should be an integer, which will have an effect on the number of matrices, $`M`$, we maximize over. However, as we will see below, the answer depends only on the logarithmic size of $`M`$, so this is not a serious problem. Finally, we remark that the placement of $`e^\gamma `$ in our definitions of $`N`$ and $`M`$ is actually irrelevant: our heuristics are sufficiently robust that increasing $`N`$ by any fixed constant and decreasing $`M`$ correspondingly leads to the same conjectured maximum. We include the $`e^\gamma `$ factor to be consistent with , where the precise choice of $`N`$ does matter.
We now find an explicit $`K`$ satisfying (2.7).
###### Theorem 2.1.
Fix $`\delta >0`$. Let $`M=\mathrm{exp}(N^\beta )`$, with $`\delta <\beta <2\delta `$, and set
$$K_\epsilon (N)=\mathrm{exp}\left(\left(\sqrt{1\frac{1}{2}\beta }+\epsilon \right)\sqrt{\mathrm{log}M}\sqrt{\mathrm{log}N}\right).$$
(2.8)
If $`U_1,\mathrm{},U_M`$ are chosen independently from $`U(N)`$, then as $`N\mathrm{}`$,
$$\left\{\underset{1jM}{\mathrm{max}}\underset{\theta }{\mathrm{max}}|\mathrm{\Lambda }_{U_j}(\theta )|K_\epsilon (N)\right\}1$$
(2.9)
for all $`\epsilon >0`$ and for no $`\epsilon <0`$.
###### Proof.
Note that by the independence of the $`U_j`$,
$$\left\{\underset{1jM}{\mathrm{max}}\underset{\theta }{\mathrm{max}}|\mathrm{\Lambda }_{U_j}(\theta )|K_\epsilon (N)\right\}=\left\{\underset{\theta }{\mathrm{max}}|\mathrm{\Lambda }_U(\theta )|K_\epsilon (N)\right\}^M$$
(2.10)
and, for this to tend to $`1`$ as $`N\mathrm{}`$, we must have
$$M\mathrm{log}\left(1\left\{\underset{\theta }{\mathrm{max}}|\mathrm{\Lambda }_U(\theta )|>K_\epsilon (N)\right\}\right)0.$$
(2.11)
Thus, the proof of the theorem (and all similar ones in this paper) requires knowledge of the tails of the distribution, and this is given by Lemma A.1. If $`M=\mathrm{exp}(N^\beta )`$ and $`K=K_\epsilon (N)`$ is as in (2.8) with $`0<\beta <2`$ and $`\epsilon >\sqrt{1\beta /2}`$, then by Lemma A.1 one easily finds that
$`M\mathrm{log}\left(1\left\{\underset{\theta }{\mathrm{max}}|\mathrm{\Lambda }_U(\theta )|>K_\epsilon (N)\right\}\right)`$ $`=\mathrm{exp}(N^\beta )\mathrm{exp}\left({\displaystyle \frac{\mathrm{log}^2K}{\mathrm{log}N\mathrm{log}\mathrm{log}K}}(1+o(1))\right)`$ (2.12)
$`0`$
as $`N\mathrm{}`$ for all $`\epsilon >0`$, but for no $`\epsilon <0`$. โ
To summarize, we use the characteristic polynomials $`\mathrm{\Lambda }_U(\theta )`$ of random unitary matrices $`UU(N)`$ to model $`Z_X(\frac{1}{2}+\mathrm{i}t)`$. To model the large values of $`|Z_X(\frac{1}{2}+\mathrm{i}t)|`$ for $`t[0,T]`$ we choose $`N`$ as in (2.5) with $`\mathrm{log}X<(\mathrm{log}N)^A`$ for some $`A`$, and we take about
$$M=N^c\mathrm{exp}\left(e^\gamma N\mathrm{log}X\right)$$
(2.13)
different matrices. Here $`c0`$ is fixed, and we include it to allay concerns that choosing too few matrices may miss some large values. With these values for $`M`$ and $`N`$, it follows that $`\beta `$ in Theorem 2.1 is $`1`$, and this leads to the following conjecture:
###### Conjecture D.
If $`2<X<\mathrm{log}^AT`$, then
$$\underset{t[0,T]}{\mathrm{max}}|Z_X(\frac{1}{2}+\mathrm{i}t)|=\mathrm{exp}\left((1+o(1))\sqrt{\frac{1}{2}\mathrm{log}T\mathrm{log}\mathrm{log}T}\right)$$
(2.14)
as $`T\mathrm{}`$.
We can now complete our argument for Conjecture A.
###### Justification of Conjecture A: .
By the prime number theorem and (2.2), we see that
$$|P_X(\frac{1}{2}+\mathrm{i}t)|\mathrm{exp}\left(\underset{nX}{}\frac{\mathrm{\Lambda }(n)}{\sqrt{n}\mathrm{log}n}\right)=O\left(\mathrm{exp}\left(3\frac{\sqrt{X}}{\mathrm{log}X}\right)\right).$$
(2.15)
Thus, if $`X=O(\mathrm{log}T)`$ and $`T/X^{e^\gamma }<t<T`$, then
$$P_X(\frac{1}{2}+\mathrm{i}t)=O\left(\mathrm{exp}\left(C\frac{\sqrt{\mathrm{log}t}}{\mathrm{log}\mathrm{log}t}\right)\right).$$
(2.16)
Combining this, (2.1), and Conjecture D, we obtain Conjecture A. โ
The argument above essentially splits the critical line into blocks of size $`1`$, maximizes over each block, and then finds the maximum of the maxima. However, one might instead wish to sample the critical line at many evenly spaced points. If they are not too sparse, then a value close to the global maximum in $`[0,T]`$ will be found. The following lemma explains why this is the case.
###### Lemma 2.2.
Suppose $`|\zeta (\frac{1}{2}+\mathrm{i}t_0)|=m_T:=\mathrm{max}_{t[0,T]}|\zeta (\frac{1}{2}+\mathrm{i}t)|`$. There is an absolute constant $`A>0`$ such that if $`|tt_0|<A/\mathrm{log}T`$, then $`|\zeta (\frac{1}{2}+\mathrm{i}t)|>\frac{1}{2}m_T`$.
###### Proof.
We can estimate the size of the derivative of the zeta-function near $`\frac{1}{2}+\mathrm{i}t_0`$ using Cauchyโs theorem. If we integrate around a circle of size $`1/\mathrm{log}T`$ and use the functional equation, we find that there exists an absolute constant $`c_1`$ such that if $`|s(\frac{1}{2}+\mathrm{i}t_0)|<c_1/\mathrm{log}T`$, then $`\zeta ^{}(s)m_T\mathrm{log}T`$. This gives the lemma. โ
As a random matrix model for $`|Z_X(\frac{1}{2}+\mathrm{i}t)|`$ when it is sampled at evenly spaced points, one might consider the largest value of $`K=K(N,X)`$ such that
$$\left\{\underset{1jN^c\mathrm{exp}(N\mathrm{log}X)}{\mathrm{max}}|\mathrm{\Lambda }_{U_j}(0)|K\right\}1.$$
(2.17)
The following theorem determines $`K`$ explicitly as a function of $`N`$ and $`X`$ and shows that such sampling is sufficient to capture the large values.
###### Theorem 2.3.
Let $`M=N^ce^{N\mathrm{log}X}`$, where $`2<X<N`$ and $`c>0`$ is fixed. If
$$K_\epsilon (N)=\mathrm{exp}\left((\frac{1}{\sqrt{2}}+\epsilon )\sqrt{N\mathrm{log}N\mathrm{log}X}\right),$$
(2.18)
and $`U_1,\mathrm{},U_M`$ are chosen independently from $`U(N)`$, then as $`N\mathrm{}`$,
$$\left\{\underset{1jM}{\mathrm{max}}|\mathrm{\Lambda }_{U_j}(0)|K_\epsilon (N)\right\}1$$
(2.19)
for all $`\epsilon >0`$ and for no $`\epsilon <0`$.
###### Proof.
As in the proof of Theorem 2.1, since the $`U_j`$ are independent we have
$$\mathrm{log}\left\{\underset{1jM}{\mathrm{max}}|\mathrm{\Lambda }_{U_j}(0)|K_\epsilon (N)\right\}=M\mathrm{log}\left(1\left\{|\mathrm{\Lambda }_U(0)|>K_\epsilon (N)\right\}\right).$$
(2.20)
Theorem 3.5 of asserts that if $`\delta >0`$ is fixed and $`\mathrm{exp}(N^\delta )K\mathrm{exp}(N^{1\delta })`$, then
$$\left\{|\mathrm{\Lambda }_U(0)|>K\right\}=\mathrm{exp}\left(\frac{\mathrm{log}^2K}{\mathrm{log}N\mathrm{log}\mathrm{log}K}(1+o_\delta (1))\right)$$
(2.21)
as $`N\mathrm{}`$. Hence, if $`M=N^c\mathrm{exp}(N\mathrm{log}X)`$ and $`K_\epsilon (N)`$ is given by (2.18), then the left-hand side of (2.20) tends to zero for all $`\epsilon >0`$, but for no $`\epsilon <0`$. โ
Note that the statements of Theorems 2.1 and 2.3 are almost identical and, in particular, one can capture the largest values of $`|\mathrm{\Lambda }_U|`$, hence $`|Z_X|`$, just by sampling at individual points; it is not necessary to find the maxima of the individual polynomials. This is significant for our modeling of the prime contribution $`P_X`$, for in that case we are only able to sample at individual points, and there is nothing comparable to a sequence of polynomials over which we can maximize individually.
### 2.2. Probabilistic model for large values of $`P_X`$
The material in this section was provided to us by Granville and Soundararajan.
First note that
$`\mathrm{log}P_X(\frac{1}{2}+\mathrm{i}t)`$ $`=`$ $`{\displaystyle \underset{pX}{}}{\displaystyle \frac{1}{p^{\frac{1}{2}+\mathrm{i}t}}}+O\left({\displaystyle \underset{p\sqrt{X}}{}}{\displaystyle \frac{1}{p}}\right)`$ (2.22)
$`=`$ $`{\displaystyle \underset{pX}{}}{\displaystyle \frac{1}{p^{\frac{1}{2}+\mathrm{i}t}}}+O(\mathrm{log}\mathrm{log}X).`$ (2.23)
Hence
$`P_X(\frac{1}{2}+\mathrm{i}t)`$ $`=`$ $`\mathrm{exp}\left({\displaystyle \underset{pX}{}}{\displaystyle \frac{1}{p^{\frac{1}{2}+\mathrm{i}t}}}\right)\times \mathrm{exp}(O(\mathrm{log}\mathrm{log}X))`$ (2.24)
$`=`$ $`\mathrm{exp}\left(P_X^{}(\frac{1}{2}+\mathrm{i}t)\right)\times \mathrm{exp}(O(\mathrm{log}\mathrm{log}X)),`$ (2.25)
say. We will see that this approximation is adequate as long as $`X=\mathrm{exp}(o(\sqrt{\mathrm{log}T\mathrm{log}\mathrm{log}T}))`$.
The method is based on treating the $`p^{\mathrm{i}t}`$ as independent random variables. The large values of $`P_X^{}(\frac{1}{2}+\mathrm{i}t)`$ can then be obtained from the following lemma, the proof of which involves calculating the moments of the distribution.
###### Lemma 2.4.
Let $`\{z_j\}`$ be a sequence of independent random variables distributed uniformly on the unit circle and let $`\{a_j\}`$ be a sequence of bounded real numbers such that for all $`n3`$,
$$\frac{1}{V_J^{\frac{n}{2}}}\underset{1jJ}{}a_j^n0$$
(2.26)
as $`J\mathrm{}`$, where
$$V_J:=\underset{1jJ}{}a_j^2.$$
(2.27)
Then, as $`J\mathrm{}`$, the distribution of
$$Y_J:=๐ข\underset{1jJ}{}a_jz_j$$
(2.28)
tends to a Gaussian with mean 0 and variance $`\frac{1}{2}V_J`$.
Applying the lemma to $`P_X^{}(\frac{1}{2}+\mathrm{i}t)`$ with $`z_j=p_j^{\mathrm{i}t}`$, where $`p_j`$ is the $`j`$th prime, and $`a_j=1/\sqrt{p_j}`$, we see that $`V_J\mathrm{log}\mathrm{log}X`$. For $`X=\mathrm{exp}(\sqrt{\mathrm{log}T})`$ we model the maximum of $`|P_X^{}(\frac{1}{2}+\mathrm{i}t)|`$ by independently choosing $`T\mathrm{log}^cT`$ values of $`t`$. By (2.24) this yields
$$\underset{t[0,T]}{\mathrm{max}}|P_X(\frac{1}{2}+\mathrm{i}t)|=\mathrm{exp}\left((1+o(1))\sqrt{\frac{1}{2}\mathrm{log}T\mathrm{log}\mathrm{log}T}\right).$$
(2.29)
If $`X=\mathrm{exp}(\sqrt{\mathrm{log}T})`$, the method of Section 2.1 predicts that
$$\underset{t[0,T]}{\mathrm{max}}|Z_X(\frac{1}{2}+\mathrm{i}t)|=O\left(\mathrm{exp}\left(\sqrt{\mathrm{log}T}\right)\right).$$
(2.30)
These estimates together with (2.1) give an alternative justification of Conjecture A.
### 2.3. Combining $`Z_X`$ and $`P_X`$
If $`X=\mathrm{exp}(\mathrm{log}^\alpha T)`$ with $`0<\alpha <\frac{1}{2}`$, then the largest values of $`|Z_X|`$ and $`|P_X|`$ are approximately the same size and both will contribute to the largest values of $`|\zeta (\frac{1}{2}+\mathrm{i}t)|`$. Specifically, applying Theorem 2.1 with $`N=\mathrm{log}T/\mathrm{log}X`$ (so that $`Z_X`$ is modeled by $`\mathrm{\Lambda }_U`$) and $`M=T\mathrm{log}X`$ (so that we sample enough characteristic polynomials to cover the critical line between $`t=0`$ and $`t=T`$), the previous analysis using characteristic polynomials predicts that $`|Z_X(\frac{1}{2}+\mathrm{i}t)|`$ gets as large as
$$\mathrm{exp}\left(\frac{1}{\sqrt{2}}\sqrt{(12\alpha )\mathrm{log}T\mathrm{log}\mathrm{log}T}\right),$$
(2.31)
and $`|P_X(\frac{1}{2}+\mathrm{i}t)|`$ gets as large as
$$\mathrm{exp}\left(\sqrt{\alpha \mathrm{log}T\mathrm{log}\mathrm{log}T}\right).$$
(2.32)
The product of these two quantities is larger than our conjectured maximum of $`|\zeta (\frac{1}{2}+\mathrm{i}t)|`$, as it should be, because we do not expect $`|Z_X|`$ and $`|P_X|`$ to attain their maximum values simultaneously. Instead of multiplying the maxima, we must find the distribution of the large values of the product $`|Z_XP_X|`$ in order to check that our method is consistent throughout the range $`0<\alpha <\frac{1}{2}`$. This is a calculation involving the tails of the distributions of $`Z_X`$ and $`P_X`$.
In Section 2.2 we saw that, conjecturally, if $`X=\mathrm{exp}(\mathrm{log}^\alpha T)`$, then the tail of $`\mathrm{log}P_X`$ has the same distribution as a Gaussian with mean 0 and variance $`\sigma _P^2=\frac{1}{2}\alpha \mathrm{log}\mathrm{log}T`$. For $`Z_X`$, Lemma A.1 states that if $`\delta >0`$ is fixed and $`\delta \lambda 1\delta `$, then
$$\left\{\underset{\theta }{\mathrm{max}}|\mathrm{\Lambda }_U(\theta )|\mathrm{exp}(N^\lambda )\right\}=\mathrm{exp}\left(\frac{N^{2\lambda }}{(1\lambda )\mathrm{log}N}(1+o(1))\right).$$
(2.33)
Since $`\mathrm{\Lambda }_U`$ models $`Z_X`$, the lemma provides the tail of the distribution of $`\mathrm{log}|Z_X(\frac{1}{2}+\mathrm{i}t)|`$. Assuming that $`Z_X`$ and $`P_X`$ are essentially statistically independent, we should expect the distribution of $`\mathrm{log}|Z_X|+\mathrm{log}|P_X|`$ to be the convolution of the two distributions. Hence, for the tail we convolve the tails of the two distributions. Thus, for large $`K`$ we expect that
$`{\displaystyle \frac{1}{T}}\mathrm{meas}\left\{0<t<T:\mathrm{log}|P_X(\frac{1}{2}+\mathrm{i}t)|+\mathrm{log}|Z_X(\frac{1}{2}+\mathrm{i}t)|\mathrm{log}K\right\}`$ (2.34)
$`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}\mathrm{exp}\left(\left({\displaystyle \frac{(\mathrm{log}Kx)^2}{\alpha \mathrm{log}\mathrm{log}T}}+{\displaystyle \frac{x^2}{(1\alpha )\mathrm{log}\mathrm{log}T\mathrm{log}x}}\right)(1+o(1))\right)dx`$ (2.35)
$`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}\mathrm{exp}\left(f_K(x)(1+o(1))\right)dx,`$ (2.36)
say. By the saddle point method, this equals
$$\mathrm{exp}\left(f_K(x_0)(1+o(1))\right),$$
(2.37)
where $`x_0`$ is such that $`f_K^{}(x_0)=0`$. If we write $`K=\mathrm{exp}\left(d\sqrt{\mathrm{log}T\mathrm{log}\mathrm{log}T}\right)`$, then solving
$$\begin{array}{c}0=f_K^{}(x_0)=\frac{2\left(x_0d\sqrt{\mathrm{log}T\mathrm{log}\mathrm{log}T}\right)}{\alpha \mathrm{log}\mathrm{log}T}+\frac{2x_0}{(1\alpha )\mathrm{log}\mathrm{log}T\mathrm{log}x_0}\hfill \\ \hfill +\frac{x_0}{((1\alpha )\mathrm{log}\mathrm{log}T\mathrm{log}x_0)^2}\end{array}$$
(2.38)
yields the solution $`x_0d(12\alpha )\sqrt{\mathrm{log}T\mathrm{log}\mathrm{log}T}`$ (which is justified so long as $`0<\alpha <1/2`$). Thus,
$$f_K(x_0)=(2+o(1))d^2\mathrm{log}T,$$
(2.39)
and this leads to the following conjecture:
###### Conjecture E.
For $`d>0`$ fixed and $`T\mathrm{}`$, we have
$$\frac{1}{T}\mathrm{meas}\left\{0<t<T:|\zeta (\frac{1}{2}+\mathrm{i}t)|>\mathrm{exp}(d\sqrt{\mathrm{log}T\mathrm{log}\mathrm{log}T})\right\}=\mathrm{exp}\left(2d^2\mathrm{log}T(1+o(1))\right).$$
(2.40)
By Lemma 2.2, $`|\zeta (\frac{1}{2}+\mathrm{i}t)|`$ is close to its maximum value over a window of size $`C/\mathrm{log}T`$, so we wish to find the smallest $`d`$ such that
$$\mathrm{meas}\left\{0<t<T:|\zeta (\frac{1}{2}+\mathrm{i}t)|\mathrm{exp}(d\sqrt{\mathrm{log}T\mathrm{log}\mathrm{log}T})\right\}\frac{1}{\mathrm{log}T},$$
(2.41)
that is, such that $`T\mathrm{log}T\mathrm{exp}(2d^2(1+o(1))\mathrm{log}T)1`$. This happens if $`d=\sqrt{\frac{1}{2}}+\epsilon `$ for any $`\epsilon >0`$, but for no $`\epsilon <0`$. Once more this gives Conjecture A, and this time in a way that is independent of the choice of $`X=\mathrm{exp}(\mathrm{log}^\alpha T)`$ for $`0<\alpha <1/2`$.
## 3. Bounds based on conjectures for moments
In this section we obtain Conjecture A by using conjectures for moments of the zeta-function. Our method also leads to limits on the possible uniformity of the conjectured moments.
Our approach here is based on the work of Conrey and Gonek . Let
$$m_T:=\underset{t[0,T]}{\mathrm{max}}|\zeta (\frac{1}{2}+\mathrm{i}t)|,$$
(3.1)
and note that we have the trivial inequality
$$m_T^{2k}\left(\frac{1}{T}_0^T|\zeta (\frac{1}{2}+\mathrm{i}t)|^{2k}dt\right).$$
(3.2)
It follows that estimates for the right-hand side of the inequality imply lower bounds for the maximum size of the zeta-function.
Keating and Snaith used random matrix theory to conjecture that if $`k>1/2`$ is fixed, then as $`T\mathrm{}`$,
$$\frac{1}{T}_0^T|\zeta (\frac{1}{2}+\mathrm{i}t)|^{2k}dt\frac{G^2(k+1)}{G(2k+1)}a(k)\mathrm{log}^{k^2}T,$$
(3.3)
where $`G`$ is the Barnes $`G`$-function, and
$$a(k)=\underset{\begin{array}{c}p\\ \text{prime}\end{array}}{}\left(1\frac{1}{p}\right)^{k^2}\underset{m=0}{\overset{\mathrm{}}{}}\left(\frac{\mathrm{\Gamma }(m+k)}{m!\mathrm{\Gamma }(k)}\right)^2p^m.$$
(3.4)
This conjecture is for $`k`$ fixed, but we would like to let $`k\mathrm{}`$, because the $`2k`$th root of the right-hand side of (3.2) then actually tends to $`m_T`$. Thus, we would like to know how large $`k`$ can be as a function of $`T`$.
Conrey and Gonek showed that if formula (3.3) holds for $`k=\sqrt{\mathrm{log}T/\mathrm{log}\mathrm{log}T}`$ then
$$m_T\mathrm{exp}\left(C_1\sqrt{\frac{\mathrm{log}T}{\mathrm{log}\mathrm{log}T}}\right),$$
(3.5)
and if it holds for $`k`$ as large as $`\mathrm{log}T/\mathrm{log}\mathrm{log}T`$ then
$$m_T\mathrm{exp}\left(C_2\frac{\mathrm{log}T}{\mathrm{log}\mathrm{log}T}\right),$$
(3.6)
where $`C_1`$ and $`C_2`$ are given explicitly. Hughes gave a convexity argument to show that formula (3.3) must fail before $`k=\mathrm{log}T/\mathrm{log}\mathrm{log}T`$. However, using the last point at which convexity holds for (3.3), one still obtains (3.6), but with a smaller constant $`C_2`$.
In all of these cases one only requires a lower bound for the right-hand side of (3.3). Our approach here is to use the mean value formula (3.3) to obtain upper bounds instead of lower bounds for $`m_T`$. As a consequence, we also obtain restrictions on the possible range of validity of (3.3) for $`k`$ growing with $`T`$. Specifically we prove the following:
###### Theorem 3.1.
Formula (3.3) does not hold for
$$k(2\sqrt{2}+\epsilon )\sqrt{\frac{\mathrm{log}T}{\mathrm{log}\mathrm{log}T}}$$
(3.7)
for any fixed $`\epsilon >0`$.
Our method allows us to get upper bounds for $`m_T`$ and, in particular, we obtain
###### Theorem 3.2.
If formula (3.3) holds for $`k=\mathrm{log}^\delta T`$ for some $`\delta <\frac{1}{2}`$, then
$$m_T\mathrm{exp}\left(\mathrm{log}^{1\delta }T\right).$$
(3.8)
Moreover, if formula (3.3) holds for $`k=\sqrt{2\mathrm{log}T/\mathrm{log}\mathrm{log}T}`$, then
$$m_T\mathrm{exp}\left(\sqrt{\frac{1}{2}\mathrm{log}T\mathrm{log}\mathrm{log}T}+O\left(\frac{\sqrt{\mathrm{log}T}\mathrm{log}\mathrm{log}\mathrm{log}T}{\sqrt{\mathrm{log}\mathrm{log}T}}\right)\right),$$
(3.9)
and if formula (3.3) holds for $`k=\sqrt{8\mathrm{log}T/\mathrm{log}\mathrm{log}T}`$, then
$$m_T\mathrm{exp}\left(\sqrt{\frac{1}{2}\mathrm{log}T\mathrm{log}\mathrm{log}T}+O\left(\frac{\sqrt{\mathrm{log}T}\mathrm{log}\mathrm{log}\mathrm{log}T}{\sqrt{\mathrm{log}\mathrm{log}T}}\right)\right).$$
(3.10)
Theorem 3.2 says that if formula (3.3) holds true until $`k=\sqrt{8\mathrm{log}T/\mathrm{log}\mathrm{log}T}`$ (after which we know it must fail), then we have
$$m_T=\mathrm{exp}\left(\sqrt{\frac{1}{2}\mathrm{log}T\mathrm{log}\mathrm{log}T}\left(1+O\left(\frac{\mathrm{log}\mathrm{log}\mathrm{log}T}{\mathrm{log}\mathrm{log}T}\right)\right)\right)$$
(3.11)
Note that this implies the conjecture found in the previous section. This is not surprising because, as we will show, the arithmetic factor in the conjectured moments is smaller than the other factors. Thus, this bound is coming just from the random matrix model.
If the true order of the zeta-function is larger than the bound in (3.11), then one would like to know where our calculation fails. Since formula (3.3) is only the leading order term in the asymptotic expansion for the $`2k`$th moment of the zeta-function, it is possible that the lower order terms dominate when $`k=\mathrm{log}^\delta T`$. However, this is unlikely. In the random matrix case, the $`2k`$th moment is given by
$`๐ผ\left\{\right|\mathrm{\Lambda }_U(\theta |^{2k}\}`$ $`={\displaystyle \frac{G^2(k+1)}{G(2k+1)}}{\displaystyle \frac{G(N+1)G(N+2k+1)}{G^2(N+k+1)}}`$ (3.12)
$`={\displaystyle \frac{G^2(k+1)}{G(2k+1)}}\mathrm{exp}\left(k^2\mathrm{log}N+{\displaystyle \frac{k^3}{N}}{\displaystyle \frac{7k^4}{12N^2}}+{\displaystyle \frac{k^5}{2N^3}}+\mathrm{}\right),`$ (3.13)
and one sees that the first term dominates even for $`k`$ as large as $`N^\delta T`$ provided that $`\delta <1`$.
In the zeta-function case, the complete main term of the $`2k`$th moment has been conjectured (see , Conjecture 1.5.1). One can check that the contribution from the primes is bounded by $`\mathrm{exp}(ck^2)`$, which is insufficient to affect the estimate for $`m_T`$. Thus, if our conjecture for the growth of $`|\zeta (\frac{1}{2}+\mathrm{i}t)|`$ is incorrect, then the main term in the mean value must take a new form for $`k=\mathrm{log}^\delta T`$ for some $`\delta <\frac{1}{2}`$. If (1.1) is the true maximal size, then by equation (3.21) the conjectured mean value can only hold for $`k\mathrm{log}\mathrm{log}T`$.
Our main tool for finding upper bounds is the following lemma.
###### Lemma 3.3.
For all positive real $`k`$ we have
$$m_T^{2k}2^{2k}\mathrm{log}T_0^T|\zeta (\frac{1}{2}+\mathrm{i}t)|^{2k}dt,$$
(3.14)
where the implied constant is absolute.
###### Proof.
Suppose that $`|\zeta (\frac{1}{2}+\mathrm{i}t_0)|=m_T`$, where $`0<t_0<T`$. By Lemma 2.2 there is an absolute constant $`A>0`$ such that if $`|tt_0|A/\mathrm{log}T`$ then $`|\zeta (\frac{1}{2}+\mathrm{i}t)|\frac{1}{2}m_T`$. This gives
$`{\displaystyle _0^T}|\zeta (\frac{1}{2}+\mathrm{i}t)|^{2k}dt`$ $``$ $`{\displaystyle _{t_0A/\mathrm{log}T}^{t_0+A/\mathrm{log}T}}|\zeta (\frac{1}{2}+\mathrm{i}t)|^{2k}dt`$ (3.15)
$``$ $`{\displaystyle \frac{A}{\mathrm{log}T}}\left({\displaystyle \frac{m_T}{2}}\right)^{2k},`$ (3.16)
as claimed. โ
###### Proof of Theorems 3.1 and 3.2.
By Lemma 3.3 and (3.2) there exists an absolute constant $`C`$ such that
$$\left(\frac{1}{T}_0^T|\zeta (\frac{1}{2}+\mathrm{i}t)|^{2k}dt\right)^{1/2k}m_T2(CT\mathrm{log}T)^{1/2\mathrm{}}\left(\frac{1}{T}_0^T|\zeta (\frac{1}{2}+\mathrm{i}t)|^2\mathrm{}dt\right)^{1/2\mathrm{}}.$$
(3.17)
We use these inequalities to prove Theorem 3.2 first.
The asymptotic expansion for the Barnes $`G`$-function (see Barnes ) gives
$$\frac{G(1+k)^2}{G(1+2k)}=\mathrm{exp}\left(k^2\left(\mathrm{log}k+\frac{3}{2}2\mathrm{log}2\right)\frac{1}{12}\mathrm{log}k+\frac{1}{12}\mathrm{log}2+\zeta ^{}(1)+O\left(\frac{1}{k}\right)\right)$$
(3.18)
for $`k1`$. Furthermore, Conrey and Gonek have shown that
$$\mathrm{log}a(k)k^2\mathrm{log}(2e^\gamma \mathrm{log}k)+o(k^2)\text{ for }k\mathrm{}.$$
(3.19)
Thus, if (3.3) holds, then
$$\mathrm{log}(\frac{1}{T}_0^T|\zeta (\frac{1}{2}+\mathrm{i}t)|^2\mathrm{}\mathrm{d}t)^{1/2\mathrm{}}=\frac{1}{2}\mathrm{}\mathrm{log}\mathrm{log}T\frac{1}{2}\mathrm{}\mathrm{log}\mathrm{}+O(\mathrm{}\mathrm{log}\mathrm{log}\mathrm{}).$$
(3.20)
It therefore follows from (3.17) that
$$\mathrm{log}m_T\frac{\mathrm{log}T+\mathrm{log}\mathrm{log}T}{2\mathrm{}}+\frac{1}{2}\mathrm{}\mathrm{log}\mathrm{log}T\frac{1}{2}\mathrm{}\mathrm{log}\mathrm{}+O(\mathrm{}\mathrm{log}\mathrm{log}\mathrm{}).$$
(3.21)
Setting $`\mathrm{}=\mathrm{log}^\delta T`$, we find that
$$m_T\mathrm{exp}\left(\frac{1}{2}\mathrm{log}^{1\delta }T+O(\mathrm{log}^\delta T\mathrm{log}\mathrm{log}T)\right),$$
(3.22)
which gives the first inequality in Theorem 3.2. Setting $`\mathrm{}=c\sqrt{\frac{\mathrm{log}T}{\mathrm{log}\mathrm{log}T}}`$, we find that
$$m_T\mathrm{exp}\left(\left(\frac{1}{2c}+\frac{c}{4}\right)\sqrt{\mathrm{log}T\mathrm{log}\mathrm{log}T}+O\left(\frac{\sqrt{\mathrm{log}T}\mathrm{log}\mathrm{log}\mathrm{log}T}{\sqrt{\mathrm{log}\mathrm{log}T}}\right)\right),$$
(3.23)
which is minimized by taking $`c=\sqrt{2}`$. This gives the second estimate in Theorem 3.2.
By (3.17) and (3.20) we also have
$$m_T\mathrm{exp}\left(\frac{1}{2}k\mathrm{log}\mathrm{log}T\frac{1}{2}k\mathrm{log}k+O(k\mathrm{log}\mathrm{log}k)\right).$$
(3.24)
If we set $`k=c\sqrt{\frac{\mathrm{log}T}{\mathrm{log}\mathrm{log}T}}`$, we find that
$$m_T\mathrm{exp}\left(\frac{c}{4}\sqrt{\mathrm{log}T\mathrm{log}\mathrm{log}T}+O\left(\frac{\sqrt{\mathrm{log}T}\mathrm{log}\mathrm{log}\mathrm{log}T}{\sqrt{\mathrm{log}\mathrm{log}T}}\right)\right).$$
(3.25)
Choosing $`c=2\sqrt{2}`$, we obtain the third inequality in Theorem 3.2. If $`c>2\sqrt{2}`$, this contradicts (3.23) and thereby establishes Theorem 3.1. โ
## 4. Bounds for $`S(t)`$
Recall that $`S(t)`$ is the error term in the counting function for the number of non-trivial zeros of the zeta-function with imaginary part less than $`t`$. It may also be expressed as $`S(t)=\frac{1}{\pi }๐ช\mathrm{log}\zeta (\frac{1}{2}+\mathrm{i}t)`$ (see Titchmarsh ). Since $`\zeta (\frac{1}{2}+\mathrm{i}t)`$ is essentially $`P_X(\frac{1}{2}+\mathrm{i}t)Z_X(\frac{1}{2}+\mathrm{i}t)`$ and $`Z_X`$ is modeled by $`\mathrm{\Lambda }_U(\theta )`$, one would expect that if $`X`$ is sufficiently small, so that the contribution from $`P_X`$ is negligible, then $`\frac{1}{\pi }๐ช\mathrm{log}\zeta (\frac{1}{2}+\mathrm{i}t)`$ too can be modeled by random matrix theory, in particular, by
$$\frac{1}{\pi }๐ช\mathrm{log}\mathrm{\Lambda }_U(0).$$
(4.1)
Evidence for this is presented in . This is the basis for our Conjecture B.
###### Theorem 4.1.
Set
$$K_\epsilon (N)=\left(\frac{1}{\sqrt{2}}+\epsilon \right)\sqrt{N\mathrm{log}N}$$
(4.2)
and $`M=N^ce^N`$, where $`c>0`$ is fixed. If $`U_1,\mathrm{},U_M`$ are chosen independently from $`U(N)`$, then as $`N\mathrm{}`$,
$$\left\{\underset{1jM}{\mathrm{max}}๐ช\mathrm{log}\mathrm{\Lambda }_{U_j}(0)K_\epsilon (N)\right\}1$$
(4.3)
for all $`\epsilon >0`$ and for no $`\epsilon <0`$.
###### Proof.
This follows along the same lines as the proof of Theorem 2.3. Since we are making independent choices,
$$\left\{\underset{1jM}{\mathrm{max}}๐ช\mathrm{log}\mathrm{\Lambda }_{U_j}(0)K\right\}=\left\{๐ช\mathrm{log}\mathrm{\Lambda }_U(0)K\right\}^M.$$
(4.4)
For this to tend to $`1`$, we need
$$M\mathrm{log}\left\{๐ช\mathrm{log}\mathrm{\Lambda }_U(0)K\right\}0.$$
(4.5)
By Theorem 3.6 of Hughes, Keating and OโConnell , if $`K=N^\lambda `$, where $`\delta <\lambda <1\delta `$ and $`\delta >0`$ is fixed, then
$$\{๐ช\mathrm{log}\mathrm{\Lambda }_U(0)K)\}=1\mathrm{exp}(\frac{K^2}{\mathrm{log}N\mathrm{log}K}(1+o_\delta (1))).$$
(4.6)
One can easily check that if $`K_\epsilon (N)`$ is given by (4.2), then (4.5) holds for all $`\epsilon >0`$, but for no $`\epsilon <0`$. โ
Conjecture B now follows in the same manner as the justification of Conjecture A; that is, by controlling the prime contribution from $`๐ช\mathrm{log}P_X(\frac{1}{2}+\mathrm{i}t)`$.
## 5. Other families and other arguments
### 5.1. Other families: symplectic and orthogonal
The analogue of Gonek, Hughes, and Keatingโs approximation to the zeta-function has not yet been extended to the case of other $`L`$-functions near the critical point. However, it is believed that the characteristic polynomials of symplectic (or orthogonal) matrices model the central value of $`L`$-functions taken from a symplectic (or orthogonal) family of $`L`$-functions . Thus, the methods developed in section 2.1 can be applied. Moreover, we can still estimate the maximal size of critical values by using a partial Euler product and modifying the method of Granville and Soundararajan. Finally, we can also apply the method involving mean values.
We give as an example finding the large values of the characteristic polynomials of the symplectic group at the critical point. The orthogonal family is treated in an almost identical way. The characteristic polynomial of an $`N\times N`$ symplectic matrix ($`N`$ must be even) with eigenvalues $`e^{\pm \mathrm{i}\theta _n}`$ is
$$Z(U,0)=\underset{j=1}{\overset{N/2}{}}(1e^{\mathrm{i}\theta _n})(1e^{\mathrm{i}\theta _n}).$$
(5.1)
Keating and Snaith calculated the moment generating function and found that
$$๐ผ_{Sp(N)}\left\{Z(U,0)^s\right\}=2^{Ns}\underset{j=1}{\overset{N/2}{}}\frac{\mathrm{\Gamma }(N/2+j+1)\mathrm{\Gamma }(s+j+1/2)}{\mathrm{\Gamma }(j+1/2)\mathrm{\Gamma }(s+N/2+j+1)}.$$
(5.2)
A long but straightforward calculation using Stirlingโs asymptotic series for the gamma function shows that if $`\delta >0`$ is fixed and $`\delta <\lambda <1\delta `$, then for $`A(N)=N^\lambda `$, $`B(N)=\frac{N^{2\lambda }}{(1\lambda )\mathrm{log}N}`$, and fixed $`s0`$, we have
$$\underset{N\mathrm{}}{lim}\frac{1}{B}\mathrm{log}๐ผ_{Sp(N)}\left\{Z(U,0)^{sB/A}\right\}=\frac{1}{2}s^2.$$
(5.3)
From this, large deviation theory (for example, see ) allows us to deduce that if $`\mathrm{exp}(N^\delta )K\mathrm{exp}(N^{1\delta })`$, then
$$_{Sp(N)}\left\{Z(U,0)>K\right\}=\mathrm{exp}\left(\frac{\mathrm{log}^2K}{2\mathrm{log}N2\mathrm{log}\mathrm{log}K}(1+o_\delta (1))\right)$$
(5.4)
as $`N\mathrm{}`$. Comparing this with (2.21), the analogous statement for the unitary group, we note the extra factor of $`2`$ in the denominator. This difference explains why the constant $`B`$ in Conjecture C equals $`1`$ rather than $`1/2`$.
We now see, by methods identical to those of the previous section, that if $`M=N^ce^N`$ for any fixed $`c0`$, and if $`K_\epsilon =\mathrm{exp}\left((1+\epsilon )\sqrt{N\mathrm{log}N}\right)`$, then if $`U_1,\mathrm{},U_M`$ are chosen independently from $`Sp(N)`$,
$$_{Sp(N)}\left\{\underset{1jM}{\mathrm{max}}Z(U_j,0)K_\epsilon (N)\right\}1$$
(5.5)
as $`N\mathrm{}`$ for all $`\epsilon >0`$ and for no $`\epsilon <0`$.
Consider for instance the family of all quadratic Dirichlet $`L`$-functions $`L(s,\chi _d)`$. For characters with modulus around $`D`$, random matrix theory suggests (see Keating and Snaith ) that $`N=\mathrm{log}D`$ is the correct identification between the size of the matrix and the conductor (though note that in $`N`$ is *half* the size of the symplectic matrix). Furthermore, it is well known that there are asymptotically $`6D/\pi ^2`$ primitive discriminants less than $`D`$. Thus, we conjecture that
$$\underset{\begin{array}{c}|d|D\\ \chi _d\text{ real}\end{array}}{\mathrm{max}}|L(\frac{1}{2},\chi _d)|=\mathrm{exp}\left((1+o(1))\sqrt{\mathrm{log}D\mathrm{log}\mathrm{log}D}\right).$$
(5.6)
Similarly, the moment generating function has been calculated for the orthogonal case (see ) and, if $`N`$ is even, we have
$$๐ผ_{SO(N)}\left\{Z(U,0)^s\right\}=2^{Ns}\underset{j=1}{\overset{N/2}{}}\frac{\mathrm{\Gamma }(N/2+j1)\mathrm{\Gamma }(s+j1/2)}{\mathrm{\Gamma }(j1/2)\mathrm{\Gamma }(s+N/2+j1)}.$$
(5.7)
Equations (5.3) and (5.4) apply to the orthogonal case without change, so by the same reasoning as previously, if $`M=N^ce^N`$ for any fixed $`c0`$, and $`K_\epsilon =\mathrm{exp}\left((1+\epsilon )\sqrt{N\mathrm{log}N}\right)`$, then
$$_{SO(N)}\left\{\underset{1jM}{\mathrm{max}}Z(U_j,0)K_\epsilon (N)\right\}1$$
(5.8)
as $`N\mathrm{}`$ for all $`\epsilon >0`$, but for no $`\epsilon <0`$.
Next we consider how to adapt the Granville-Soundararajan argument involving the product over primes to the symplectic case. We require the following lemma.
###### Lemma 5.1.
Let $`\{x_j\}`$ be a sequence of independent *real* random variables with mean 0 and variance 1, and let $`\{a_j\}`$ be a bounded sequence of real numbers such that for all $`n3`$,
$$\frac{1}{V_J^{\frac{n}{2}}}\underset{1jJ}{}a_j^n0$$
(5.9)
as $`J\mathrm{}`$, where
$$V_J:=\underset{1jJ}{}a_j^2.$$
(5.10)
Then as $`J\mathrm{}`$, the distribution of
$$Y_J:=\underset{1jJ}{}a_jx_j$$
(5.11)
tends to a Gaussian with mean 0 and variance $`V_J`$.
Just as in the treatment involving characteristic polynomials, the fact that the variance is $`V_j`$ for these families instead of $`V_j/2`$ leads to the constant $`B=1`$ in Conjecture C, instead of $`B=1/2`$ for the unitary family dealt with previously.
### 5.2. Other arguments, for
Although the conjectures in this paper are based on very recent work, Hugh Montgomery has pointed out to us that a similar conjecture can be obtained by viewing $`\mathrm{log}|\zeta (\frac{1}{2}+\mathrm{i}t)|`$ as a Gaussian distributed random variable with variance $`C\mathrm{log}\mathrm{log}T`$, where one estimates $`m_T`$ by sampling $`T^A`$ times.
Soundararajan suggests a different way to use the moments of the zeta-function to conjecture an upper bound. The proof of Lemma 3.3 showed that if
$$\frac{1}{T}\mathrm{meas}\left\{t[0,T]:|\zeta (\frac{1}{2}+\mathrm{i}t)|\tau \right\}\frac{1}{T},$$
(5.12)
then $`m_T2\tau `$. For when $`|\zeta (\frac{1}{2}+\mathrm{i}t)|`$ is very large, it must remain large over an interval of size $`c/\mathrm{log}T`$. Now
$$\tau ^{2k}\mathrm{meas}\left\{t[0,T]:|\zeta (\frac{1}{2}+\mathrm{i}t)|\tau \right\}_0^T|\zeta (\frac{1}{2}+\mathrm{i}t)|^{2k}dt,$$
(5.13)
so if equation (3.3) holds, then we have
$$\frac{1}{T}\mathrm{meas}\left\{t[0,T]:|\zeta (\frac{1}{2}+\mathrm{i}t)|\tau \right\}\tau ^{2k}\frac{G^2(k+1)}{G(2k+1)}(\mathrm{log}T)^{k^2}.$$
(5.14)
The right-hand side is less than $`c/(T\mathrm{log}T)`$ (which means there is only one place where the maximum occurs) when
$$\tau \mathrm{exp}\left(\frac{\mathrm{log}T}{2k}+\frac{k}{2}\mathrm{log}\mathrm{log}T\frac{k}{2}\mathrm{log}k\right).$$
(5.15)
The minimum of this is $`\mathrm{exp}\left(\sqrt{\frac{1}{2}\mathrm{log}T\mathrm{log}\mathrm{log}T}\right)`$ and it occurs when $`k=\sqrt{2\mathrm{log}T/\mathrm{log}\mathrm{log}T}`$.
### 5.3. Other arguments, against
We now discuss potential arguments against the conjectures in this paper. One possibility, so fundamental that it cannot be addressed, is that the large values of an $`L`$-function may be so rare that these statistical models cannot detect them. Indeed, since these are problems in number theory, there may be number-theoretic constructions of large values which contradict our conjectures. The two examples below, due to Brian Conrey, suggest the kinds of things we have in mind.
The first argument invokes an analogy with the divisor function $`d(n)=_{d|n}1`$ and the related function $`\omega (n)=_{p|n}1`$, where here the sum is over prime divisors of $`n`$. Since $`\omega (n)`$ is $`\mathrm{log}\mathrm{log}n`$ on average and has a Gaussian distribution, the relation $`d(n)=2^{\omega (n)}`$ for square-free $`n`$ might lead one to conjecture that for such $`n`$, $`d(n)`$ is bounded by
$$\mathrm{exp}(c\sqrt{\mathrm{log}n/\mathrm{log}\mathrm{log}n}).$$
(5.16)
However, we know how to construct large values of $`d(n)`$, and it is easy to see that for $`n`$ squarefree, $`d(n)`$ can get as large as
$$\mathrm{exp}(c\mathrm{log}n/\mathrm{log}\mathrm{log}n).$$
(5.17)
The second argument concerns the Fourier coefficients $`a_n`$ of cusp forms. For integer weight cusp forms, rescaled so that for $`p`$ prime we have $`|a_p|2`$, the coefficients $`a_n`$ can get as large as
$$\mathrm{exp}(c\mathrm{log}n/\mathrm{log}\mathrm{log}n).$$
(5.18)
In other words, they can get about as large as $`d(n)`$. The question is: can the coefficients of half-integral weight forms also get this large? If they can, then our conjecture on the maximal size of the critical values of a symplectic family of $`L`$-functions is incorrect. For if $`fS_k(\mathrm{\Gamma }_0(N))`$, then $`L_f(\frac{1}{2},\chi _d)=c_d^2/\sqrt{d}`$, where $`c_d`$ is a Fourier coefficient of the half-integral weight form associated with $`f`$ by the Shimura correspondence.
Our methods cannot be directly applied to produce a version of Conjecture B for symplectic and orthogonal families. But if one were to believe that for a family $``$ of $`L`$-functions with $`c(F)`$ denoting the conductor of $`F`$,
$$\underset{c(F)\mathrm{}}{lim\; sup}\frac{๐ช\mathrm{log}F(\frac{1}{2})}{\sqrt{\mathrm{log}c(F)\mathrm{log}\mathrm{log}c(F)}}=\sqrt{B}$$
(5.19)
where $`B=1/2`$ for unitary families and $`B=1`$ for symplectic and orthogonal families, then since the rank of an elliptic curve is related to the order of vanishing of its associated $`L`$-function, this could lead to new information about large ranks. That is, if (5.19) is true, then it suggests that for rational elliptic curves we have
$$\underset{c_E\mathrm{}}{lim\; sup}\frac{\mathrm{rank}(E)}{\sqrt{\mathrm{log}c_E\mathrm{log}\mathrm{log}c_E}}=1,$$
(5.20)
where $`c_E`$ is the conductor of $`E`$. Note that this is smaller than the ranks of elliptic curves found by Ulmer in the function field case.
## Appendix A The tail of the distribution of $`\mathrm{max}_\theta |\mathrm{\Lambda }_U(\theta )|`$
Here we prove the random matrix polynomial result used in Section 2.
###### Lemma A.1.
If $`\delta >0`$ is fixed and $`\delta \lambda 1\delta `$, then
$$\left\{\underset{\theta }{\mathrm{max}}|\mathrm{\Lambda }_U(\theta )|\mathrm{exp}(N^\lambda )\right\}=\mathrm{exp}\left(\frac{N^{2\lambda }}{(1\lambda )\mathrm{log}N}(1+o(1))\right).$$
(A.1)
###### Proof of Lemma A.1.
Bernsteinโs inequality for polynomials implies that for any matrix $`U`$,
$$\underset{\theta }{\mathrm{max}}|\mathrm{\Lambda }_U^{}(\theta )|N\underset{\theta }{\mathrm{max}}|\mathrm{\Lambda }_U(\theta )|.$$
(A.2)
Thus, if $`\varphi `$ is a point at which the maximum of $`|\mathrm{\Lambda }_U(\theta )|`$ occurs, and if we indicate the maximum by $`m_U`$, then for $`|\theta \varphi |1/N`$,
$$|\mathrm{\Lambda }_U(\theta )|m_U|\theta \varphi |Nm_U.$$
(A.3)
It follows that
$`{\displaystyle _0^{2\pi }}|\mathrm{\Lambda }_U(\theta )|^{2k}d\theta `$ $`m_U^{2k}{\displaystyle _{1/N}^{1/N}}\left(1|x|N\right)^{2k}dx`$ (A.4)
$`=m_U^{2k}{\displaystyle \frac{2}{N}}{\displaystyle \frac{1}{2k+1}}.`$ (A.5)
Combining this with the trivial lower bound for $`m_U`$, we find that
$$\frac{1}{2\pi }_0^{2\pi }|\mathrm{\Lambda }_U(\theta )|^{2k}d\theta m_U^{2k}\frac{2k+1}{2}N_0^{2\pi }|\mathrm{\Lambda }_U(\theta )|^{2k}d\theta .$$
(A.6)
This bound holds for any matrix. We now average over all $`N\times N`$ unitary matrices with respect to Haar measure. That is, we calculate the expectation $`๐ผ_N`$ of $`|\mathrm{\Lambda }_U(\theta )|^{2k}`$. Set
$$๐ผ_N\left\{|\mathrm{\Lambda }_U(\theta )|^{2k}\right\}=M_N(2k).$$
(A.7)
Keating and Snaith have shown that
$$M_N(2k)=\frac{G^2(k+1)}{G(2k+1)}\frac{G(1+N)G(1+N+2k)}{G^2(1+N+k)},$$
(A.8)
where $`G`$ is the Barnes $`G`$-function. Note that this is independent of $`\theta `$. Therefore, by (A.6)
$$M_N(2k)๐ผ\left\{m_U^{2k}\right\}\pi (2k+1)NM_N(2k).$$
(A.9)
Hughes, Keating and OโConnell have shown that if $`A(N)=N^\lambda `$ with $`\delta <\lambda <1\delta `$ and $`\delta >0`$ fixed, and if
$$B(N)=\frac{N^{2\lambda }}{(1\lambda )\mathrm{log}N},$$
(A.10)
then for $`s0`$,
$$\underset{N\mathrm{}}{lim}\frac{1}{B(N)}\mathrm{log}M_N\left(\frac{sB(N)}{A(N)}\right)=\frac{1}{4}s^2.$$
(A.11)
Since
$`{\displaystyle \frac{1}{B(N)}}\mathrm{log}M_N(sB(N)/A(N))`$ $`{\displaystyle \frac{1}{B(N)}}\mathrm{log}๐ผ\left\{m_U^{sB(N)/A(N)}\right\}`$ (A.12)
$``$ $`{\displaystyle \frac{1}{B(N)}}\mathrm{log}M_N(sB(N)/A(N))+O\left({\displaystyle \frac{(\mathrm{log}N)^2}{N^{2\lambda }}}\right),`$ (A.13)
we conclude that for $`s0`$,
$$\underset{N\mathrm{}}{lim}\frac{1}{B(N)}\mathrm{log}๐ผ\left\{\mathrm{exp}\left(\frac{sB(N)\mathrm{log}(\mathrm{max}_\theta |\mathrm{\Lambda }_U(\theta )|)}{A(N)}\right)\right\}=\frac{1}{4}s^2.$$
(A.14)
From this, large deviation theory (see, for example, ) allows us to deduce that
$$\underset{N\mathrm{}}{lim}\frac{1}{B(N)}\mathrm{log}\left\{\mathrm{log}\underset{\theta }{\mathrm{max}}|\mathrm{\Lambda }_U(\theta )|A(N)\right\}=1.$$
(A.15)
Inserting $`A(N)=N^\lambda `$ and $`B(N)`$ from (A.10), we obtain the statement in the lemma. โ
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# Can a dust dominated universe have accelerated expansion?
## 1 Introduction
Kolb and coworkers argued recently that growing perturbations on a scale larger than the Hubble length may lead to accelerated expansion as observed from the center of the perturbation. This was then followed up by Wiltshire and Carter et al. , who considered a dust dominated universe model where the observed universe is an underdense bubble in an Einstein-de Sitter universe. They calculated the luminosity-redshift relationship for their model and found very small deviations from the standard $`\mathrm{\Lambda }`$CDM model. Thus, they were able to obtain a good fit to the Hubble diagram of type Ia supernovae.
However, critical papers have appeared. We will now proceed to give a brief overview over some of these papers and the arguments therein. By considering a universe model where sub-Hubble perturbations are absent, Flanagan argues that the contributions from the super-Hubble perturbations to the value of the deceleration parameter are so small that they cannot be responsible for the acceleration of the universe. Geshnizjani et al. argue that to second order in spatial gradients the super-Hubble perturbations only amount to a renormalization of local spatial curvature, and thus cannot account for the negative deceleration parameter. Making an expansion to Newtonian order in potential and velocity, but taking into account fully non-linear density inhomogeneities, Siegel and Fry obtain similar results, and conclude that inhomogeneity contributions cannot mimic the effects of dark energy or induce an accelerated expansion.
One of the strongest and most convincing critiques of the conclusions of Kolb et al. is that of Hirata and Seljak . They argue from the Raychaudhuri equation that in a dust dominated universe there must be a non-vanishing vorticity in order to obtain a negative deceleration parameter. Then they showed that the perturbations considered in have vanishing vorticity and hence cannot lead to accelerated expansion. In yet another critique of the work of Kolb et al., Rรคsรคnen also claims that their conclusions are ruled out. The basis for this conclusion was an analysis of the model proposed in by applying the Buchert formalism for backreaction . This conclusion is also supported in works by Giovannini , where the author performs a general analysis of the class of models considered in . He finds that dust-dominated models cannot have accelerated expansion, and hence, the conclusion of Kolb et al. cannot be correct.
Several of the arguments against accelerated expansion induced by inhomogeneities in matter are based upon perturbation calculations. However, as noted in , one way of evading these arguments is through non-perturbative effects. This is what we will investigate in the present work. We want to know if Einsteinโs field equations permit inhomogeneities to change the positive deceleration parameter of a dust-dominated homogeneous universe model to a negative value.
Unfortunately, a general analysis incorporating all possible inhomogeneous models is not feasible, so we will instead concentrate on one specific class of such models, namely the dust-dominated, spherically symmetric models. Moffat suggested in that it might be possible to find solutions with accelerated expansion. Considering the field equations for these models, an expression for the local deceleration parameter was derived, and it was claimed that this might indeed allow for negative values. Later, in , this claim was moderated considerably. Here it was claimed that the *volume averaged* deceleration parameter, and not the local deceleration parameter, could be negative. No explicit solution of the field equation were actually found in these analyses. Instead, the conclusions were reached by investigating the field equations directly.
In this work, we will consider the same class of spherically symmetric, inhomogeneous universe models and present a general solution to the field equations. We then show that no such model can have an accelerated expansion in the sense of a negative local deceleration parameter.
The structure of this paper is as follows. In Sect. 2 we present the field equations for a general spherically symmetric, inhomogeneous universe model. In Sect. 3 we discuss the properties of these models and whether they can have accelerated expansion. Finally, in Sect. 4 we give a brief summary of our work and present our conclusion.
## 2 Spherically symmetric, inhomogeneous universe models
The spherically symmetric, inhomogeneous universe models are described by the Lemaรฎtre-Tolman-Bondi (LTB) space-time . The line element can be written as
$$ds^2=dt^2+X^2(r,t)dr^2+R^2(r,t)d\mathrm{\Omega }^2,$$
(1)
where $`X(r,t)`$ and $`R(r,t)`$ are general function to be determined by the field equations and boundary and initial conditions. They can be thought of as generalized, position dependent scale factors in the radial and transverse directions. The Einstein field equations are
$$G_{\mu \nu }=R_{\mu \nu }\frac{1}{2}g_{\mu \nu }R=\kappa T_{\mu \nu },$$
(2)
where $`\kappa =8\pi G`$.
We assume the universe to be occupied by an ideal fluid with energy density $`\rho (r,t)`$ and an isotropic pressure $`p(r,t)`$. In the comoving coordinates defined in Eq. (1), the energy-momentum tensor of the fluid can be written as $`T_{\mu \nu }=\text{diag}(\rho ,p,p,p)`$. A relation between $`X`$ and $`R`$ can be found by solving the $`t`$-$`r`$ component of the field equations, for which the right-hand side vanishes. The relation one finds is
$$X(r,t)=\frac{R^{}(r,t)}{f(r)},$$
(3)
where $`f(r)`$ is an arbitrary function of $`r`$ only. Throughout this paper, we will use a $`{}_{}{}^{}=d/dr`$ to denote differentiation with respect to $`r`$ and $`\dot{}=d/dt`$ for differentiation with respect to $`t`$.
Following Moffat , we define two โHubble parametersโ
$$H_{}=\frac{\dot{R}}{R}\text{and}H_r=\frac{\dot{R}^{}}{R^{}},$$
(4)
which are a measure of the local expansion rate in the perpendicular and radial directions respectively, and the deceleration parameter
$$q_{}=\frac{1}{H_{}^2}\frac{\ddot{R}}{R}.$$
(5)
Using these definitions, the time-time and space-space components of the field equations take the form
$$H_{}^2+2H_rH_{}\frac{\beta }{R^2}\frac{\beta ^{}}{RR^{}}=\kappa \rho $$
(6)
and
$$6H_{}^2q_{}+2H_{}^22\frac{\beta }{R^2}2H_rH_{}+\frac{\beta ^{}}{RR^{}}=\kappa (\rho +3p).$$
(7)
where $`\beta f^21`$. These are the generalization of the Friedmann equations in a homogeneous universe to a spherically symmetric, inhomogeneous universe. Note that the function $`\beta (r)`$ is often written as $`\beta (r)=2E(r)r^2`$, where $`E(r)`$ determines the local curvature radius.
## 3 Expansion rate of dust-dominated LTB models
From now on we will assume the cosmic fluid to be pressure-less matter. Adding Eqs. (6) and (7), we then obtain the following expression for the deceleration parameter
$$q_{}=\frac{1}{2}\frac{\beta }{2\dot{R}^2}.$$
(8)
The condition for accelerated expansion now takes the form
$$\beta >\dot{R}^2>0\text{ or }f^2>1+\dot{R}^2.$$
(9)
Substituting the definition for the deceleration parameter in Eq. (5) into Eq. (8), we arrive at the expression
$$2R\ddot{R}+\dot{R}^2=\beta .$$
(10)
A single integration of this equation yields
$$R\dot{R}^2=\beta R+\alpha \text{ or }H_{}^2=\frac{\beta }{R^2}+\frac{\alpha }{R^3},$$
(11)
where the function $`\alpha (r)`$ enters as an integration โconstantโ and depends only on the radial coordinate $`r`$. Comparing this equation with the ordinary Friedmann equation for homogeneous models, we see that the dynamical effects of $`\beta `$ and $`\alpha `$ are similar to those of curvature and dust, respectively. Therefore, $`\alpha (r)`$ is regarded as the gravitational mass function, and one often chooses $`\alpha (r)r^3`$. Note, however, that $`\alpha `$ does not appear in the energy-momentum tensor.
From Eqs. (9) and (11), we find that $`\alpha `$ and $`\beta `$ must satisfy the inequality
$$\beta R<\alpha (r)<0$$
(12)
in order for the deceleration parameter to be negative. Thus, it appears that the LTB models seem to allow for accelerated expansion even for dust-dominated universe models, as long as $`\alpha `$ and $`\beta `$ are chosen appropriately. The dynamical effect of $`\alpha `$ would then correspond to that of dust with negative energy density in a homogeneous universe model. It should be noted that the inequality above forbids accelerated expansion in a โbig bangโ model where the scale factor has the initial value $`R(0,r)=0`$, which implies that $`\alpha (r)>0`$. However, this initial condition may not be physically realistic. The universe may have started with a finite scale factor, or maybe has collapsed and reached a finite minimum radius before expanding again. In such models accelerated expansion does not seem to be forbidden. However, as we will show shortly, such solutions are unphysical. First, we find a general solution of the field equation and show how such an accelerated solution behaves.
We introduce a conformal time $`\eta `$ defined via the differential relation $`\beta ^{1/2}dt=Rd\eta `$. This allow us to integrate Eq. (11) to give a parametric solution in terms of $`\eta `$. The solutions can be grouped into three different classes according to the sign of the local curvature $`\beta `$. Choosing $`t_0=\eta _0=0`$, the three classes are
$`R`$ $`=`$ $`{\displaystyle \frac{\alpha }{2\beta }}(\mathrm{cosh}\eta 1)+R_0\left[\mathrm{cosh}\eta +\sqrt{{\displaystyle \frac{\alpha +\beta R_0}{\beta R_0}}}\mathrm{sinh}\eta \right]`$ (13)
$`\sqrt{\beta }t`$ $`=`$ $`{\displaystyle \frac{\alpha }{2\beta }}(\mathrm{sinh}\eta \eta )+R_0\left[\mathrm{sinh}\eta +\sqrt{{\displaystyle \frac{\alpha +\beta R_0}{\beta R_0}}}\left(\mathrm{cosh}\eta 1\right)\right]\text{}\beta >0`$
for positive $`\beta `$,
$`R`$ $`=`$ $`{\displaystyle \frac{\alpha }{2|\beta |}}(1\mathrm{cos}\eta )+R_0\left[\mathrm{cos}\eta +\sqrt{{\displaystyle \frac{\alpha +\beta R_0}{|\beta |R_0}}}\mathrm{sin}\eta \right]`$ (14)
$`\sqrt{|\beta |}t`$ $`=`$ $`{\displaystyle \frac{\alpha }{2|\beta |}}(\eta \mathrm{sin}\eta )+R_0\left[\mathrm{sin}\eta +\sqrt{{\displaystyle \frac{\alpha +\beta R_0}{|\beta |R_0}}}(1\mathrm{cos}\eta )\right]\text{}\beta <0`$
for negative $`\beta `$, and finally
$$R=\left(R_0^{3/2}+\frac{3}{2}\sqrt{\alpha }t\right)^{2/3}$$
(15)
for vanishing $`\beta `$. We allow the initial size of the universe to be non-zero, hence the non-vanishing value for $`R_0R(t=0,r)`$. Eqs. (13)-(15) with $`R_0=0`$ represent the usual form of the LTB solution of the field equations. The scale factor $`R`$ as a function of time is shown for some typical parameters in Fig. 1. The top curve represents a solution with accelerating expansion.
In order for a dust-dominated solution in the LTB space-time to undergo an epoch of accelerated expansion, the inequality in Eq. (12) has to be satisfied. Can a physically realistic solution satisfy this? Differentiating Eq. (11) and comparing the result with Eq. (6) we obtain a relation between the density distribution and $`\alpha `$ for this class of models,
$$\kappa \rho =\frac{\alpha ^{}}{R^2R^{}}=3\frac{\alpha ^{}}{V^{}},$$
(16)
where $`V=R^3`$ is a comoving volume. A physically realistic model must have positive density everywhere, and hence, it must have $`\alpha ^{}>0`$. Furthermore, the angular part of the line element is $`R^2d\mathrm{\Omega }^2`$ where $`d\mathrm{\Omega }`$ is a solid angle element. At the origin, $`r=0`$, this part of the line element must vanish, and thus $`R(0,t)=0`$. From Eq. (11) we then get that $`\alpha (0)=0`$. Since $`\alpha ^{}(r)>0`$, it follows that $`\alpha (r)0`$ for all $`r`$. But, as we showed above, accelerated expansion is only possible for models with $`\alpha (r)<0`$. Hence, the dust dominated LTB universe models must have decelerated expansion.
We shall now arrive at the same conclusion by an alternative approach related to the Raychaudhuri equation, which describes the flow of a cosmological fluid in space-time. It can be written as
$$_\mu a^\mu =\dot{\theta }+\frac{\theta ^2}{3}+2(\sigma ^2\omega ^2)+R_{\mu \nu }u^\mu u^\nu ,$$
(17)
where $`a^\mu u_{;\alpha }^\mu u^\alpha `$ is the four-acceleration of the fluid, $`\theta u_{;\mu }^\mu `$ is the expansion rate of the fluid and $`\sigma `$ and $`\omega `$ are the shear and vorticity scalars, respectively. The latter two quantities are defined as
$$\sigma ^2=\frac{1}{2}\sigma _{\mu \nu }\sigma ^{\mu \nu }\text{and}\omega ^2=\frac{1}{2}\omega _{\mu \nu }\omega ^{\mu \nu }$$
(18)
where
$$\sigma _{\mu \nu }=u_{(\mu ;\nu )}+a_{(\mu }u_{\nu )}\frac{1}{3}\theta (g_{\mu \nu }+u_\mu u_\nu )$$
(19)
and
$$\omega _{\mu \nu }=u_{[\mu ;\nu ]}+a_{[\mu }u_{\nu ]}.$$
(20)
In a comoving reference frame the four-velocity of the fluid can be written as $`u^\mu =[1,0,0,0]`$. The divergence of the four-acceleration will then vanish, and the Raychaudhuri equation can be simplified to
$$0=\dot{\theta }+\frac{\theta ^2}{3}+2(\sigma ^2\omega ^2)+R_{00}.$$
(21)
Analogously to the Hubble parameter in FRW models, we can define an effective local Hubble parameter from the expansion rate of the fluid:
$$H=\frac{1}{3}u_{;\mu }^\mu .$$
(22)
This allow us to define an effective deceleration parameter in terms of $`H`$ and $`\dot{H}`$:
$$q\frac{\dot{H}+H^2}{H^2}=13\dot{\theta }/\theta ^2.$$
(23)
Using the relation $`R_{00}=\frac{\kappa }{2}(\rho +3p)`$, Eq. (21) can be written
$$\theta ^2q=6(\sigma ^2\omega ^2)+\frac{3}{2}\kappa (\rho +3p).$$
(24)
While Eq. (24) is valid for any metric, let us now specialize to the case of an LTB model. In it was claimed that the vorticity scalar is non-vanishing for the models considered in the paper, which may permit negative values of q. We find, however, that the vorticity tensor (20) vanishes for a comoving fluid in the metric (1). This incorporates also the models considered in . For this metric the shear can be written as
$$\sigma ^2=\frac{1}{3}(H_rH_{})^2.$$
(25)
Finally, assuming the fluid to be pressure-less, Eq. (24) reduces to
$$\theta ^2q=2(H_rH_{})^2+\frac{3}{2}\kappa \rho .$$
(26)
This equation shows that the dust dominated LTB-models have $`q0`$, meaning that they cannot have accelerated expansion.
Although we have now shown that the LTB models cannot have accelerated expansion, this doesnโt necessarily mean that they cannot explain the data used to infer this expansion. The first data used to conclude that the expansion seems to be accelerating were the measurements of the luminosities of supernovae of type Ia (SNIa) . Since then newer data have appeared that strengthen this claim even further . However, accelerated expansion follows logically from these data only if one assumes that the universe is homogeneous. The added freedom of having a position dependent expansion in LTB models allows one to explain the data without the need for the expansion to accelerate locally. The explanation would then be that the expansion rate is highest at $`r=0`$ and decreases with distance from the center, since the oldest supernovae are also farthest away. In this case a decelerating universe would seem to be accelerating. Cรฉlรฉrier shows in how such inhomogeneities can mimic the effects of dark energy in supernova observations, when interpreted within the framework of FRW models. More specifically, Iguchi et al. show that models with $`\alpha (r)>0`$ can reproduce the luminosity distance-redshift relationship of a $`\mathrm{\Lambda }`$CDM model up to $`z1`$, but not for higher redshifts. However, when compared to actual SNIa observations, it is easy to find models that fit the data even better than the best-fit $`\mathrm{\Lambda }`$CDM model . In fact, as we show in , these models might also be compatible with the CMB power spectrum, even though they do not contain any form of dark energy.
Finally, we wish to stress the point that in this work we have shown that the LTB models cannot have an expansion that is *locally* accelerated. This does not exclude the possibility that there can be a so-called *volume averaged acceleration*, where a scale factor defined via the physical volume of a comoving region is found to have a positive double time derivative. Recent papers discussing the possibilities of this effect explaining the apparent accelerated expansion are Refs. . At the present time there does not appear to exist a consensus among cosmologists on this issue.
## 4 Conclusion
Recently it was claimed by Kolb et al. that inhomogeneous perturbations of FRW universe models can result in accelerated cosmic expansion, and thereby eliminate the need to postulate the existence of the mysterious dark energy. However, several papers have appeared since criticizing this work heavily. These attribute the effects claimed by Kolb et al. to incorrect handling of second order terms in the perturbative expansion. When dealt with correctly, one would then assume that these terms do not lead to accelerated expansion. However, this leaves open the possibility that the full, non-perturbative solutions of the Einstein equations for inhomogeneous models might exhibit accelerated expansion.
In this work we investigated a special class of inhomogeneous models which can be solved analytically, the so-called Lemaรฎtre-Tolman-Bondi (LTB) models. These are the spherically symmetric, inhomogeneous models. We found a general solution to a dust-dominated LTB model, and showed that there exist solutions that appear to have accelerated expansion. However, these solutions turned out to be unphysical. Finally, it was shown that any physically realistic solution of Einsteinโs equations must have decelerated expansion if it contains only matter.
## Acknowledgments
MA acknowledges support from the Norwegian Research Council through the project 159637/V30 โ โShedding Light on Dark Energyโ.
## References
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# Irreversibility for All Bound Entangled States
## Appendix A Appendix
Definition. For any continuous function $`f`$ of state on Hilbert space $`C^d`$ one defines a mixed convex roof $`\widehat{f}`$
$$\widehat{f}(\rho )=inf\underset{i}{}p_if(\rho _i),$$
(9)
where the infimum is taken over all finite decompositions $`_ip_i\rho _i=\rho `$.
Proposition 3. The infimum is attained, and the optimal ensemble can be chosen to have $`d^2+1`$ elements.
Proof. We use standard techniques from information theory Csiszar (see also Uhlmann ). First, let us show that for any finite decomposition $`\rho =_{i=1}^np_i\rho _i`$, we can provide a decomposition $`\rho =_{i=1}^{d^2+1}q_i\sigma _i`$ with $`d^2+1`$ elements, such that
$$\underset{i=1}{\overset{n}{}}p_if(\rho _i)=\underset{i=1}{\overset{d^2+1}{}}q_if(\sigma _i).$$
(10)
To this, consider convex hull $`๐`$ of the set $`\{(\rho _i,f(\rho _i))\}_{i=1}^n`$. The point $`x=(_{i=1}^np_i\rho _i,_{i=1}^np_if(\rho _i))`$ belongs $`๐`$. The set $`๐`$ is a compact convex set, actually a polyhedron, in $`d^2`$-dimensional real affine space (this comes from the fact that states belongs to the real $`d^2`$ dimensional space of Hermitian operators and have unit trace). The set of extremal points is included in the set $`\{\rho _i,f(\rho _i)\}_{i=1}^n`$. Then from Caratheodory theorem it follows that $`x`$ can be written as a convex combination of at most $`d^2+1`$ extremal points, i.e. $`x=_{i_j}q_{i_j}(\rho _{i_j},f(\rho _{i_j}))`$ where $`j=1,\mathrm{}d^2+1`$. Writing $`q_{i_j}=q_j`$, $`\rho _{i_j}=\sigma _j`$ we get $`_{i=1}^np_i\rho _i=_{j=1}^{d^2+1}q_j\sigma _j`$ and $`_{i=1}^np_if(\rho _i)=_{j=1}^{d^2+1}q_jf(\sigma _j)`$. Thus we have found a decomposition that has $`d^2+1`$ elements, and returns the same value of average, so that the infimum can be taken solely over such decompositions. Then from continuity of the function and compactness of the set of states it follows that the infimum is attained.
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# Quadratic response theory for spin-orbit coupling in semiconductor heterostructures
## I Introduction
The Rashba Hamiltonian Rashba (1960) is the prototype of a class of effective-mass Hamiltonians describing spin-orbit coupling in semiconductors. Winkler (2003) These models have been under intensive study in the past several years due to theoretical and experimental advances in spin-related phenomena such as the intrinsic spin Hall effect Murakami et al. (2003); Murakami et al. (2004a); Sinova et al. (2004); Culcer et al. (2004); Shen et al. (2004, 2005); Ma et al. (2004); Murakami et al. (2004b); Schliemann and Loss (2005); Bernevig and Zhang (2005, 2004); Zhang and Yang (2005); Hu (2005); Kato et al. (2004); Wunderlich et al. (2005) and the spin galvanic and circular photogalvanic effects. Ganichev et al. (2001); Ganichev et al. (2002a, b); Golub (2003); Ganichev et al. (2003); Belโkov et al. (2003); Ganichev and Prettl (2003); Ganichev et al. (2004) Such effects are generated by spin-orbit coupling terms in the conduction or valence bands of clean nonmagnetic semiconductors. Although a variety of different two- and three-dimensional semiconductor systems are under investigation, one of the most widely studied is a heterostructure between semiconductors with the zinc-blende structure, in which an external electric field can be used to tune the relative contributions from the Rashba and Dresselhaus spin-splitting terms. Nitta et al. (1997); Schliemann et al. (2003); Ganichev et al. (2004)
In a two-dimensional effective-mass model, the Rashba coupling has no coordinate dependence. But in a three-dimensional theory, it is usually separated into (1) a contribution proportional to the macroscopic electric field generated by gate voltages, dopants, and free carriers; and (2) $`\delta `$ functions representing the contribution from the rapid change in potential at a heterojunction. Leibler (1977); Vasโko (1979); Bastard (1981, 1982); Lassnig (1985); Lommer et al. (1985); Gerchikov and Subashiev (1992); Foreman (1993); Pfeffer and Zawadzki (1995); de Andrada e Silva et al. (1997); Schรคpers et al. (1998); Winkler (2000, 2003) However, the assumption that a heterojunction can be represented by a short-range $`\delta `$ potential has never been justified from first principles. In a self-consistent theory with electron-electron interactions, there are in general long-range Coulomb multipole potentials that are not well localized at the interface. These long-range potentials contribute spin-dependent terms to the Hamiltonian. Thus, it is important to establish the conditions under which such terms will appear in the effective-mass Hamiltonian for a heterojunction.
This paper examines the long-range terms in the self-energy of a quasiparticle for the case of a lattice-matched semiconductor heterostructure. The potential energy of the ions is described using norm-conserving pseudopotentials. Bachelet et al. (1982); Goedecker et al. (1996); Hartwigsen et al. (1998) The heterostructure pseudopotential is treated as a perturbation of a bulk reference crystal, with the self-energy calculated to second order in the perturbation using quadratic response theory. Hu and Zaremba (1988); Pitarke et al. (1995); Bergara et al. (1997, 1999); Nazarov et al. (2002); Nagao et al. (2003) This approach is well justified numerically, since the linear response alone has been shown to give excellent predictions for the valence band offset in a variety of material systems, including isovalent and heterovalent heterostructures. Resta et al. (1989); Baroni et al. (1988, 1989); Peressi et al. (1990, 1991); Colombo et al. (1991); Baroni et al. (1992); Peressi and Baroni (1994); Montanari et al. (1996)
The approach used here follows closely earlier work by Sham Sham (1966) on the theory of shallow impurity states in bulk semiconductors. Shamโs work is generalized to include nonlocal spin-dependent potentials and terms of higher order in the crystal momentum. Even for local spin-independent potentials, the present work includes terms neglected in Shamโs analysis, such as dipole potentials in the quadratic response.
This paper is limited to a study of the electron self-energy in the limit of small crystal momentum. The derivation of an effective-mass Hamiltonian from these results is presented in the following paper. Foreman (2005)
As a workable approximation scheme, the calculation of the linear response is carried out to terms two orders in $`q`$ higher than the lowest nonvanishing term, where $`๐ช`$ is the crystal momentum transfer of the perturbing potential. The quadratic response is calculated to the same order in $`q`$ as the lowest nonvanishing term in the linear response. (See the following paper Foreman (2005) for further discussion of this approximation scheme.) Three classes of heterostructure perturbations are considered:
(I) Heterovalent perturbations with nonzero charge. In this case the perturbing potential includes a monopole of $`O(q^2)`$, and the analysis is performed to an accuracy of $`O(q^0)`$ in the linear response and $`O(q^2)`$ in the quadratic response.
(II) Isovalent perturbations for which the linear response has a nonzero dipole moment. In this case the leading term in the linear response is $`O(q^1)`$, so the linear response is evaluated to $`O(q)`$ and the quadratic response is evaluated to $`O(q^1)`$.
(III) Isovalent perturbations for which the linear response has no dipole moment. In this case the leading term in the linear response is $`O(q^0)`$, so the linear response is evaluated to $`O(q^2)`$ and the quadratic response is evaluated to $`O(q^0)`$.
The perturbations that make up a given heterostructure are generally a mixture of classes I, II, and III. The simplest situation is that of an isovalent heterostructure made up of semiconductors with the zinc-blende structure, such as GaAs/AlAs or InAs/GaSb. In this case, every ionic perturbation is an isovalent perturbation from class III.
Most theoretical and experimental studies of the Rashba spin-splitting Hamiltonian have dealt with this type of heterostructure. This case is therefore studied in greatest detail here, by working out the explicit form of the self-energy from crystal symmetry. The results show that in this case the Rashba Hamiltonian contains only short-range terms (to within the accuracy of the stated approximation scheme). However, there are long-range spin-dependent terms that are not of the Rashba form.
In a heterovalent system such as Ge/GaAs, the ionic perturbations are from class I. However, since macroscopic accumulations of charge are energetically unfavorable, real heterostructures tend to be macroscopically neutral. Baroni et al. (1989); Harrison et al. (1978) Such a nominally heterovalent class I problem can therefore often be reduced to an isovalent class II or III problem by grouping the ions together in neutral clusters and treating these clusters as the basic unit. This approach is discussed further in Appendix A.
In wurtzite heterostructures such as GaN/AlN, the ionic perturbations are from class II, since the site symmetry of atoms in the wurtzite structure (space group $`C_{6v}^4`$) permits a dipole moment. These dipole terms produce spontaneous polarization along the hexagonal $`c`$ axis in bulk wurtzite crystals, leading to macroscopic interface charge at heterojunctions. Vanderbilt and King-Smith (1993); Bernardini et al. (1997a, b); Bernardini and Fiorentini (1998); Fiorentini et al. (1999); Bechstedt et al. (2000); Bernardini et al. (2001); Zoroddu et al. (2001) Such charge produces macroscopic electric fields that generate different piezoelectric strain fields in different materials. The present theory, which is restricted to lattice-matched heterostructures, is therefore not generally applicable to wurtzite systems (except in the unrealistic special case Vanderbilt and King-Smith (1993) where the interface polarization charge is exactly cancelled by an external interface charge). However, the results derived here provide a first step towards a more general theory dealing with lattice-mismatched heterostructures.
The paper begins in Sec. II by establishing the basic definitions and notation for the Green function and self-energy that are used throughout the paper. The finite-temperature formalism is used (both for generality and because it facilitates the derivation of Ward identities), although the main interest of this paper is the limit of an insulator at zero temperature. In Sec. III, the self-energy is expanded in powers of the perturbing potential using vertex functions. A set of Ward identities is derived for the vertex functions at finite and zero temperature. Section IV presents general expressions for the nonlinear density response, including Ward identities for the static polarization. The bare ionic perturbations are screened in Sec. V, where the proper vertex functions and proper polarization are introduced.
A detailed analysis of the small wave vector properties of the linear screened potential is carried out in Sec. VI for the special case of a local spin-independent perturbation. The quadratic response for the same case is considered in Sec. VII, and the linear and quadratic contributions to the self-energy are derived in Sec. VIII. In Sec. IX it is shown that in the norm-conserving pseudopotential formalism, the contributions from the nonlocal spin-dependent part of the perturbing potential merely renormalize the contributions from the local part of the perturbation. The main results of the paper are discussed and summarized in Sec. X.
## II Green function and self-energy
This section establishes the notation, basic definitions, and symmetry properties of the Green function and self-energy used in subsequent sections of the paper. The starting point is the definition of the one-particle thermal Green function Abrikosov et al. (1975); Fetter and Walecka (2003); Negele and Orland (1998)
$$G_{ss^{}}(๐ฑ,\tau ;๐ฑ^{},\tau ^{})=T_\tau [\widehat{\psi }_s(๐ฑ,\tau )\widehat{\psi }_s^{}^{}(๐ฑ^{},\tau ^{})],$$
(1)
where $`s=\pm \frac{1}{2}`$ labels the $`z`$ component of the spin, $`\tau `$ is the imaginary time, $`T_\tau `$ is the time ordering operator, and $`\widehat{\psi }_s(๐ฑ,\tau )=e^{\widehat{K}\tau }\widehat{\psi }_s(๐ฑ)e^{\widehat{K}\tau }`$ and $`\widehat{\psi }_s^{}(๐ฑ,\tau )=e^{\widehat{K}\tau }\widehat{\psi }_s^{}(๐ฑ)e^{\widehat{K}\tau }`$ are field operators in the Heisenberg picture. The angular brackets denote a thermal average
$$\widehat{O}=e^{\beta \mathrm{\Omega }}\mathrm{Tr}(e^{\beta \widehat{K}}\widehat{O}),$$
(2)
where $`\beta =1/k_BT`$ is the inverse temperature, $`\mathrm{Tr}`$ denotes a trace over the many-particle Fock space, $`\widehat{K}=\widehat{H}\mu \widehat{N}`$ is the grand Hamiltonian (with $`\mu `$ the chemical potential and $`\widehat{N}`$ the number operator), and $`e^{\beta \mathrm{\Omega }}=\mathrm{Tr}(e^{\beta \widehat{K}})`$. The many-particle Hamiltonian is defined by
$$\begin{array}{c}\widehat{H}=\underset{s,s^{}}{}\widehat{\psi }_s^{}(๐ฑ)h_{ss^{}}(๐ฑ,๐ฑ^{})\widehat{\psi }_s^{}(๐ฑ^{})d^3xd^3x^{}\hfill \\ \hfill +\frac{1}{2}\underset{s,s^{}}{}\frac{\widehat{\psi }_s^{}(๐ฑ)\widehat{\psi }_s^{}^{}(๐ฑ^{})\widehat{\psi }_s^{}(๐ฑ^{})\widehat{\psi }_s(๐ฑ)}{|๐ฑ๐ฑ^{}|}d^3xd^3x^{},\end{array}$$
(3)
where $`h=h^{}`$ is the Hamiltonian of a single noninteracting particle, and Hartree atomic units are used.
Since $`\widehat{K}`$ is time independent, $`G`$ has the form $`G_{ss^{}}(๐ฑ,\tau ;๐ฑ^{},\tau ^{})=G_{ss^{}}(๐ฑ,๐ฑ^{},\tau \tau ^{})`$, with $`G_{ss^{}}(๐ฑ,๐ฑ^{},\tau \beta )=G_{ss^{}}(๐ฑ,๐ฑ^{},\tau )`$ for $`0<\tau <\beta `$. This permits the Fourier series representation (for $`\beta <\tau <\beta `$)
$`G(\tau )`$ $`={\displaystyle \frac{1}{\beta }}{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}G(\zeta _n)e^{i\zeta _n\tau },`$ (4a)
$`G(\zeta _n)`$ $`={\displaystyle _0^\beta }G(\tau )e^{i\zeta _n\tau }๐\tau ,`$ (4b)
where $`\zeta _n=(2n+1)\pi /\beta `$, and $`G(\tau )`$ denotes a single-particle operator whose matrix elements in the $`|๐ฑ,s`$ basis are $`G_{ss^{}}(๐ฑ,๐ฑ^{},\tau )`$. A continuous Green function $`G(\omega )`$ may then be defined by analytic continuation of $`G(\zeta _n)`$ from the discrete frequencies $`\omega =\mu +i\zeta _n`$.
The Green function $`G(\omega )`$ satisfies Dysonโs equation
$$[\omega h\mathrm{\Sigma }(\omega )]G(\omega )=G(\omega )[\omega h\mathrm{\Sigma }(\omega )]=1,$$
(5)
which is an implicit definition for the self-energy operator
$$\mathrm{\Sigma }(\omega )=\omega hG^1(\omega ).$$
(6)
A formal solution to Eq. (5) can be constructed by solving the nonhermitian eigenvalue equations
$`[h+\mathrm{\Sigma }(\omega )]|\psi _n(\omega )`$ $`=E_n(\omega )|\psi _n(\omega ),`$ (7a)
$`[h+\mathrm{\Sigma }^{}(\omega )]|\chi _n(\omega )`$ $`=E_n^{}(\omega )|\chi _n(\omega ),`$ (7b)
which are also referred to as Dysonโs equations. It is usually assumed that the solutions to (7) form a complete biorthonormal Morse and Feshbach (1953) set with the properties
$$\chi _n(\omega )|\psi _n^{}(\omega )=\delta _{nn^{}},$$
(8a)
$$\underset{n}{}|\psi _n(\omega )\chi _n(\omega )|=1,$$
(8b)
although this is difficult to prove in general. Friedman (1956) If Eqs. (8) are valid, then $`G(\omega )`$ is given by Layzer (1963); Morse and Feshbach (1953)
$$G(\omega )=\underset{n}{}\frac{|\psi _n(\omega )\chi _n(\omega )|}{\omega E_n(\omega )},$$
(9)
which satisfies the Dyson equation (5) by construction.
Symmetries of the many-particle Hamiltonian $`\widehat{H}`$ under time reversal and space group operations imply corresponding symmetries of the one-particle operators $`G`$ and $`\mathrm{\Sigma }`$. Derivations of the most useful symmetry relations are presented in Appendix B.
## III Vertex functions
In this section, a perturbative approach to the Dyson equation (7a) is developed by using vertex functions to expand the self-energy in powers of the perturbing potential.
### III.1 Definitions
The single-particle Hamiltonian $`h`$ is chosen here to have the form
$$h_{ss^{}}(๐ฑ,๐ฑ^{})=\frac{1}{2}^2\delta (๐ฑ๐ฑ^{})\delta _{ss^{}}+v_{ss^{}}^{\text{ext}}(๐ฑ,๐ฑ^{}),$$
(10)
where the fixed external potential $`v^{\text{ext}}`$ is a norm-conserving ionic pseudopotential, Bachelet et al. (1982); Goedecker et al. (1996); Hartwigsen et al. (1998) which accounts for both spin-orbit coupling Hybertsen and Louie (1986); Surh et al. (1991); Hemstreet et al. (1993); Theurich and Hill (2001) and scalar relativistic effects. The Dyson equation (7a) is therefore
$$\begin{array}{c}\frac{1}{2}^2\psi _s(๐ฑ,\omega )+\underset{s^{}}{}V_{ss^{}}(๐ฑ,๐ฑ^{},\omega )\psi _s^{}(๐ฑ^{},\omega )d^3x^{}\hfill \\ \hfill =E(\omega )\psi _s(๐ฑ,\omega ),\end{array}$$
(11)
in which $`\psi `$ is a spinor wave function, and the total potential energy $`V`$ is
$$V_{ss^{}}(๐ฑ,๐ฑ^{},\omega )=v_{ss^{}}^{\text{ext}}(๐ฑ,๐ฑ^{})+\mathrm{\Sigma }_{ss^{}}(๐ฑ,๐ฑ^{},\omega ).$$
(12)
The external pseudopotential can be separated as $`v^{\text{ext}}=v^{(0)}+v`$, where $`v^{(0)}`$ is the potential of some periodic reference crystal, and $`v`$ is a nonperiodic perturbation associated with a heterostructure or an impurity. It is assumed that the total potential (12) can be represented as a power series in the perturbation $`v`$:
$$V_{s_1s_2}(๐ฑ_1,\tau _1;๐ฑ_2,\tau _2)V(12)=\underset{\nu =0}{\overset{\mathrm{}}{}}V^{(\nu )}(12),$$
(13)
where $`V^{(0)}`$ is the potential (12) when $`v^{\text{ext}}=v^{(0)}`$, and the numerical arguments on the right-hand side are shorthand for the space, spin, and time coordinates $`(1)=(๐ฑ_1,s_1,\tau _1)`$. Although the upper limit of the formal expansion (13) is written as $`\nu =\mathrm{}`$, this may well be an asymptotic series, and in practice only a finite number of terms are retained.
The linear and quadratic terms of (13) are
$$\begin{array}{cc}\hfill V^{(1)}(12)& =\mathrm{\Gamma }^{(1)}(1243)v(34),\hfill \\ \hfill V^{(2)}(12)& =\frac{1}{2}\mathrm{\Gamma }^{(2)}(124365)v(34)v(56),\hfill \end{array}$$
(14)
where $`\mathrm{\Gamma }^{(\nu )}`$ is called the vertex function of order $`\nu `$. Here a summation or integration of repeated coordinates is assumed, and the labels are ordered as the trace of a matrix product. The perturbation $`v`$ is taken to be an instantaneous static potential of the form
$$v(34)=v_{s_3s_4}(๐ฑ_3,๐ฑ_4)\delta (\tau _3\tau _4).$$
(15)
The vertex functions are by definition functional derivatives of $`V`$ with respect to $`v`$:
$$\begin{array}{cc}\hfill \mathrm{\Gamma }^{(1)}(1243)& =\frac{\delta V(12)}{\delta v(34)},\hfill \\ \hfill \mathrm{\Gamma }^{(2)}(124365)& =\frac{\delta ^2V(12)}{\delta v(34)\delta v(56)}\hfill \\ & =\frac{\delta \mathrm{\Gamma }^{(1)}(1243)}{\delta v(56)},\hfill \end{array}$$
(16)
which may also be expressed as
$$\begin{array}{cc}\hfill \mathrm{\Gamma }^{(1)}(1243)& =\delta (13)\delta (24)+\frac{\delta \mathrm{\Sigma }(12)}{\delta v(34)},\hfill \\ \hfill \mathrm{\Gamma }^{(2)}(124365)& =\frac{\delta ^2\mathrm{\Sigma }(12)}{\delta v(34)\delta v(56)},\hfill \end{array}$$
(17)
in which $`\delta (12)=\delta _{s_1s_2}\delta (๐ฑ_1๐ฑ_2)\delta (\tau _1\tau _2)`$. Note that upon application of the Fourier transforms defined in Appendix C, the above equations hold equally well in momentum and frequency space.
It is convenient at this point to carry out the time integrals in Eq. (14). This eliminates the variables $`\tau _3,\tau _4,\mathrm{}`$ from $`\mathrm{\Gamma }^{(\nu )}`$ and $`v`$, and reduces the time dependence of the $`\delta (13)\delta (24)`$ term in Eq. (17) to $`\delta (\tau _1\tau _2)`$. It is assumed below that this has been done.
### III.2 Ward identities
The vertex functions satisfy various Ward identities Noziรจres (1964); Sham and Kohn (1966); Jones and March (1973) for certain limiting values of their arguments. One set of these can be derived by varying the chemical potential $`\mu `$ by a small amount $`\delta \mu `$. Since $`\widehat{K}=\widehat{H}\mu \widehat{N}`$, this is equivalent to varying $`v`$ by $`\delta v(12)=\delta \mu \delta (12)`$. For this special case, Eqs. (16) and (17) reduce to
$$\begin{array}{cc}\hfill \mathrm{\Gamma }^{(1)}(1233)& =\delta (12)\frac{\delta \mathrm{\Sigma }(12)}{\delta \mu },\hfill \\ \hfill \mathrm{\Gamma }^{(2)}(124355)& =\frac{\delta \mathrm{\Gamma }^{(1)}(1243)}{\delta \mu },\hfill \end{array}$$
(18)
which involve a trace over the input variables of $`\mathrm{\Gamma }`$. If $`\mathrm{\Sigma }`$ is analytically continued as a function of $`\omega =\mu +i\zeta _n`$, the variation with respect to $`\mu `$ may be written as
$$\frac{\delta }{\delta \mu }\frac{}{\omega }+\frac{}{\mu }.$$
(19)
Now for the special case of an insulator at $`T=0`$, the chemical potential can have any value in the range $`\mu _{N1}<\mu <\mu _N`$, where $`\mu _N`$ is the minimum energy needed to add one particle to the ground state of an $`N`$-particle system, and the energy gap is $`E_\text{g}=\mu _N\mu _{N1}`$. \[For small temperatures $`\mu `$ approaches the well-defined limit $`\frac{1}{2}(\mu _N+\mu _{N1})`$, Ashcroft and Mermin (1976) but at exactly $`T=0`$ it becomes ill defined.\] Since $`\mu `$ can vary arbitrarily within the gap for a system with finite $`E_\text{g}`$, the Ward identities for the insulator reduce to
$$\begin{array}{cc}\hfill \mathrm{\Gamma }^{(1)}(1233)& =\delta (12)\frac{\mathrm{\Sigma }(12)}{\omega },\hfill \\ \hfill \mathrm{\Gamma }^{(2)}(124355)& =\frac{\mathrm{\Gamma }^{(1)}(1243)}{\omega },\hfill \\ \hfill \mathrm{\Gamma }^{(2)}(123344)& =\frac{^2\mathrm{\Sigma }(12)}{\omega ^2},\hfill \end{array}$$
(20)
which generalize and extend the results derived for spin-independent local potentials in Refs. Sham, 1966 and Sham and Kohn, 1966.
## IV Nonlinear density response
### IV.1 Definitions
In this section, perturbation theory (see Appendix D) is used to evaluate the electron density of the system with Hamiltonian $`\widehat{H}=\widehat{H}_0+\widehat{H}_1`$, in which $`\widehat{H}_0`$ is the Hamiltonian of the reference crystal and $`\widehat{H}_1`$ is the perturbation due to $`v`$:
$$\widehat{H}_1=\mathrm{tr}(\widehat{\rho }v)=\underset{s,s^{}}{}\widehat{\rho }_{s^{}s}(๐ฑ^{},๐ฑ)v_{ss^{}}(๐ฑ,๐ฑ^{})d^3xd^3x^{},$$
(21)
where $`\widehat{\rho }`$ is the density operator
$$\widehat{\rho }_{s^{}s}(๐ฑ^{},๐ฑ)=\widehat{\psi }_s^{}(๐ฑ)\widehat{\psi }_s^{}(๐ฑ^{}).$$
(22)
The mean nonlocal electron density in the perturbed system is defined as
$$n_{ss^{}}(๐ฑ,๐ฑ^{})=\widehat{\rho }_{ss^{}}(๐ฑ,๐ฑ^{},\tau ),$$
(23)
which is independent of $`\tau `$. If $`n`$ is evaluated using the perturbation theory formula (139), one obtains a power series in $`v`$:
$$n_{ss^{}}(๐ฑ,๐ฑ^{})=\underset{\nu =0}{\overset{\mathrm{}}{}}n_{ss^{}}^{(\nu )}(๐ฑ,๐ฑ^{}),$$
(24)
in which
$$n_{ss^{}}^{(0)}(๐ฑ,๐ฑ^{})=\widehat{\rho }_{ss^{}}(๐ฑ,๐ฑ^{},\tau )_0$$
(25)
is the density of the reference crystal. (The notation $`\widehat{O}_0`$ refers to a thermal average with respect to the reference crystal; see Appendix D.) The terms of order $`\nu >0`$ are given by
$$n^{(\nu )}(00^{})=\frac{1}{\nu !}\mathrm{\Pi }^{(\nu )}(00^{},1^{}1,\mathrm{},\nu ^{}\nu )v(11^{})\mathrm{}v(\nu \nu ^{}),$$
(26)
where $`\mathrm{\Pi }^{(\nu )}`$ is the $`\nu ^{\text{th}}`$-order static polarization (or density correlation function), which is defined in Eq. (141). Here and below, the numerical arguments of $`n^{(\nu )}`$, $`\mathrm{\Pi }^{(\nu )}`$, and $`v`$ are time-independent quantities of the form $`(0)=(๐ฑ_0,s_0)`$.
### IV.2 Ward identities
The linear and quadratic density response are given explicitly by
$$\begin{array}{cc}\hfill n^{(1)}(12)& =\mathrm{\Pi }^{(1)}(1243)v(34),\hfill \\ \hfill n^{(2)}(12)& =\frac{1}{2}\mathrm{\Pi }^{(2)}(124365)v(34)v(56),\hfill \end{array}$$
(27)
which shows that $`\mathrm{\Pi }`$ may be defined as a functional derivative of $`n`$ with respect to $`v`$:
$$\begin{array}{cc}\hfill \mathrm{\Pi }^{(1)}(1243)& =\frac{\delta n(12)}{\delta v(34)},\hfill \\ \hfill \mathrm{\Pi }^{(2)}(124365)& =\frac{\delta ^2n(12)}{\delta v(34)\delta v(56)}\hfill \\ & =\frac{\delta \mathrm{\Pi }^{(1)}(1243)}{\delta v(56)}.\hfill \end{array}$$
(28)
For the special case $`\delta v(12)=\delta \mu \delta (12)`$ representing a variation in chemical potential of $`\delta \mu `$ (see Sec. III.2), these expressions give
$$\begin{array}{cc}\hfill \mathrm{\Pi }^{(1)}(1233)& =\frac{n(12)}{\mu },\hfill \\ \hfill \mathrm{\Pi }^{(2)}(124355)& =\frac{\mathrm{\Pi }^{(1)}(1243)}{\mu },\hfill \end{array}$$
(29)
which are the Ward identities for the static polarization. For an insulator at $`T=0`$, $`\mu `$ is indefinite, and these reduce to
$$\mathrm{\Pi }^{(1)}(1233)=0,\mathrm{\Pi }^{(2)}(124355)=0.$$
(30)
## V Screening
### V.1 Potential
In this section, the concept of screening is used to extract the long-range Coulomb interaction terms from the polarization and vertex functions. The first-order screened potential $`\phi `$ is defined by adding the Coulomb potential generated by $`n^{(1)}`$ to $`v`$:
$$\phi (12)=v(12)+u(1243)\mathrm{\Pi }^{(1)}(3465)v(56).$$
(31)
Here $`u`$ represents the Coulomb interaction, which is spin-independent and local in coordinate space at both the input and output:
$$u(1243)=\delta _{s_1s_2}\delta _{s_3s_4}\delta (๐ฑ_1๐ฑ_2)\delta (๐ฑ_3๐ฑ_4)u(๐ฑ_1๐ฑ_3).$$
(32)
In momentum space (see Appendix C) this has the form
$$u(1243)=\delta _{s_1s_2}\delta _{s_3s_4}\delta _{๐ค_1๐ค_2,๐ค_3๐ค_4}u(๐ค_1๐ค_2),$$
(33)
in which $`u(๐ค)=v_c(๐ค)/\mathrm{\Omega }`$, where $`\mathrm{\Omega }`$ is the crystal volume and
$$v_c(๐ค)=\{\begin{array}{cc}4\pi /k^2\hfill & \text{if }k0\text{,}\hfill \\ 0\hfill & \text{if }k=0\text{.}\hfill \end{array}$$
(34)
Another way of writing the screened potential is
$$\phi (12)=ฯต^1(1243)v(34),$$
(35)
in which the inverse static electronic dielectric matrix is
$$ฯต^1(1243)=\delta (13)\delta (24)+u(1265)\mathrm{\Pi }^{(1)}(5643).$$
(36)
The dielectric matrix $`ฯต`$ satisfies
$$ฯต(1243)ฯต^1(3465)=ฯต^1(1243)ฯต(3465)=\delta (15)\delta (26),$$
(37)
and is given explicitly below in Eq. (47). For an insulator, the Ward identity (30) yields
$$ฯต^1(1233)=\delta (12).$$
(38)
The second-order potential $`\phi ^{(2)}`$ is just the Coulomb potential generated by $`n^{(2)}`$:
$$\phi ^{(2)}(12)=u(1243)n^{(2)}(34).$$
(39)
By translation symmetry, $`\mathrm{\Pi }^{(1)}(1243)=0`$ unless $`๐ค_1๐ค_2=๐ค_3๐ค_4+๐`$, where $`๐`$ is a reciprocal lattice vector of the reference crystal. Equation (31) may therefore be written as
$$\begin{array}{c}\phi _{ss^{}}(๐ค,๐ค^{})=v_{ss^{}}(๐ค,๐ค^{})+\delta _{ss^{}}v_c(๐ช)\underset{๐ค^{\prime \prime }}{}\underset{๐}{}\hfill \\ \hfill \times \mathrm{\Pi }_{\lambda ^{}\lambda }(๐ช;๐ค^{\prime \prime },๐ค^{\prime \prime }+๐ช+๐)v_{\lambda \lambda ^{}}(๐ค^{\prime \prime }+๐ช+๐,๐ค^{\prime \prime }),\end{array}$$
(40)
where $`๐ช๐ค๐ค^{}`$ and
$$\begin{array}{c}\mathrm{\Pi }_{\lambda ^{}\lambda }(๐ช;๐ค,๐ค+๐ช+๐)\hfill \\ \hfill =\frac{1}{\mathrm{\Omega }}\underset{๐ค_2}{}\mathrm{\Pi }_{\alpha \alpha ,\lambda ^{}\lambda }^{(1)}(๐ช+๐ค_2,๐ค_2;๐ค,๐ค+๐ช+๐).\end{array}$$
(41)
This simplified form of $`\mathrm{\Pi }`$ is introduced because the Coulomb potential depends only on the local spin-independent density $`n(๐ฑ)n_{\alpha \alpha }(๐ฑ,๐ฑ)`$.
### V.2 Vertex functions
Given Eqs. (35) and (39), one can rewrite the total potentials (14) as
$$\begin{array}{c}\begin{array}{cc}\hfill V^{(1)}(12)& =\stackrel{~}{\mathrm{\Gamma }}^{(1)}(1243)\phi (34),\hfill \\ \hfill V^{(2)}(12)& =\frac{1}{2}\stackrel{~}{\mathrm{\Gamma }}^{(2)}(124365)\phi (34)\phi (56)\hfill \end{array}\hfill \\ \hfill +\stackrel{~}{\mathrm{\Gamma }}^{(1)}(1243)\phi ^{(2)}(34),\end{array}$$
(42)
in which the proper vertex function $`\stackrel{~}{\mathrm{\Gamma }}^{(\nu )}`$ is defined in perturbation theory as the sum of all $`\nu ^{\text{th}}`$-order vertex diagrams that cannot be separated into two disconnected parts by cutting one Coulomb interaction line or one electron propagator. (An alternative definition would be as the $`\nu ^{\text{th}}`$ functional derivative of $`V`$ with respect to $`\phi `$.Hedin (1965); Hedin and Lundqvist (1969)) From the above results, $`\mathrm{\Gamma }`$ and $`\stackrel{~}{\mathrm{\Gamma }}`$ are related by
$$\begin{array}{c}\begin{array}{cc}\hfill \mathrm{\Gamma }^{(1)}(1243)& =\stackrel{~}{\mathrm{\Gamma }}^{(1)}(1265)ฯต^1(5643),\hfill \\ \hfill \mathrm{\Gamma }^{(2)}(124365)& =\stackrel{~}{\mathrm{\Gamma }}^{(2)}(128709)ฯต^1(7843)ฯต^1(9065)\hfill \end{array}\hfill \\ \hfill +\stackrel{~}{\mathrm{\Gamma }}^{(1)}(1287)u(8709)\mathrm{\Pi }^{(2)}(904365).\end{array}$$
(43)
In an insulator, the Ward identities (30) and (38) yield
$$\begin{array}{cc}\hfill \mathrm{\Gamma }^{(1)}(1233)& =\stackrel{~}{\mathrm{\Gamma }}^{(1)}(1244),\hfill \\ \hfill \mathrm{\Gamma }^{(2)}(124355)& =\stackrel{~}{\mathrm{\Gamma }}^{(2)}(128799)ฯต^1(7843),\hfill \\ \hfill \mathrm{\Gamma }^{(2)}(123344)& =\stackrel{~}{\mathrm{\Gamma }}^{(2)}(125566).\hfill \end{array}$$
(44)
Hence, for insulators, the Ward identities (20) are valid for both $`\mathrm{\Gamma }`$ and $`\stackrel{~}{\mathrm{\Gamma }}`$.
### V.3 Polarization
In a similar fashion, one can define the proper polarization $`\stackrel{~}{\mathrm{\Pi }}`$ as the sum of all static polarization diagrams that cannot be split by cutting a Coulomb line. Thus
$$\mathrm{\Pi }^{(1)}(1243)=\stackrel{~}{\mathrm{\Pi }}^{(1)}(1265)ฯต^1(5643),$$
(45)
which has the form of a Dyson equation: Jones and March (1973)
$$\begin{array}{cc}\hfill \mathrm{\Pi }^{(1)}(1243)& =\stackrel{~}{\mathrm{\Pi }}^{(1)}(1243)+\stackrel{~}{\mathrm{\Pi }}^{(1)}(1265)u(5687)\mathrm{\Pi }^{(1)}(7843)\hfill \\ & =\stackrel{~}{\mathrm{\Pi }}^{(1)}(1243)+\mathrm{\Pi }^{(1)}(1265)u(5687)\stackrel{~}{\mathrm{\Pi }}^{(1)}(7843).\hfill \end{array}$$
(46)
This can be used to verify that the dielectric matrix
$$ฯต(1243)=\delta (13)\delta (24)u(1265)\stackrel{~}{\mathrm{\Pi }}^{(1)}(5643)$$
(47)
is indeed the inverse of Eq. (36). Equations (45) and (46) can also be written as
$$\mathrm{\Pi }^{(1)}(1243)=ฯต^1(6512)\stackrel{~}{\mathrm{\Pi }}^{(1)}(5643),$$
(48)
in which the symmetry property (144) was used. The total and proper quadratic polarizations are likewise related by (see Fig. 5 of Ref. Sham, 1966)
$$\begin{array}{c}\mathrm{\Pi }^{(2)}(124365)=ฯต^1(8712)\stackrel{~}{\mathrm{\Pi }}^{(2)}(78092^{}1^{})ฯต^1(9043)\hfill \\ \hfill \times ฯต^1(1^{}2^{}65).\end{array}$$
(49)
In an insulator, Eqs. (38) and (45) give $`\stackrel{~}{\mathrm{\Pi }}^{(1)}(1233)=\mathrm{\Pi }^{(1)}(1244)=0`$, which implies that $`ฯต(1233)=\delta (12)`$. The inverses of Eqs. (45) and (49), i.e.,
$$\stackrel{~}{\mathrm{\Pi }}^{(1)}(1243)=\mathrm{\Pi }^{(1)}(1265)ฯต(5643),$$
(50)
$$\stackrel{~}{\mathrm{\Pi }}^{(2)}(124365)=ฯต(8712)\mathrm{\Pi }^{(2)}(78092^{}1^{})ฯต(9043)ฯต(1^{}2^{}65),$$
(51)
then show that the insulator Ward identities (30) are valid for both $`\mathrm{\Pi }`$ and $`\stackrel{~}{\mathrm{\Pi }}`$.
Since the Coulomb interaction depends only on the reduced polarization matrix (41), Eq. (46) can be reduced to
$$\begin{array}{c}\mathrm{\Pi }_{ss^{}}(๐ช+๐;๐ค,๐ค+๐ช+๐^{})=\stackrel{~}{\mathrm{\Pi }}_{ss^{}}(๐ช+๐;๐ค,๐ค+๐ช+๐^{})\hfill \\ \hfill +\underset{๐^{\prime \prime }}{}\stackrel{~}{\mathrm{\Pi }}(๐ช+๐,๐ช+๐^{\prime \prime })v_c(๐ช+๐^{\prime \prime })\\ \hfill \times \mathrm{\Pi }_{ss^{}}(๐ช+๐^{\prime \prime };๐ค,๐ค+๐ช+๐^{}),\end{array}$$
(52)
in which a scalar version of $`\mathrm{\Pi }`$ is defined by
$$\mathrm{\Pi }(๐ช,๐ช+๐)=\underset{๐ค}{}\mathrm{\Pi }_{\lambda \lambda }(๐ช;๐ค,๐ค+๐ช+๐).$$
(53)
Now $`v_c(๐ช+๐^{\prime \prime })`$ is nonsingular in the limit $`q0`$ when $`๐^{\prime \prime }\mathrm{๐}`$, so it is convenient to regroup the series expansion of Eq. (52) so as to isolate the terms $`v_c(๐ช)`$: Ambegaokar and Kohn (1960); Jon (a)
$$\begin{array}{c}\mathrm{\Pi }_{ss^{}}(๐ช+๐;๐ค,๐ค+๐ช+๐^{})=P_{ss^{}}(๐ช+๐;๐ค,๐ค+๐ช+๐^{})\hfill \\ \hfill +P(๐ช+๐,๐ช)v_c(๐ช)\mathrm{\Pi }_{ss^{}}(๐ช;๐ค,๐ค+๐ช+๐^{}).\end{array}$$
(54)
Here $`P`$ is the sum of all polarization diagrams that cannot be separated by cutting a Coulomb line labeled with $`๐ช`$ (although they may be split by cutting lines labeled $`๐ช+๐^{\prime \prime }`$ with $`๐^{\prime \prime }\mathrm{๐}`$). This will be called the regular polarization; it is related to the proper polarization $`\stackrel{~}{\mathrm{\Pi }}`$ by Eq. (52) with $`\mathrm{\Pi }P`$ and $`๐^{\prime \prime }\mathrm{๐}`$.
Both $`\stackrel{~}{\mathrm{\Pi }}`$ and $`P`$ are well behaved in the limit $`q0`$, but $`P`$ is more convenient for analysis of the small-$`q`$ behavior of $`\mathrm{\Pi }`$ because, unlike the case for $`\stackrel{~}{\mathrm{\Pi }}`$, it does not require the inversion of matrices (see Ref. Onida et al., 2002 for further discussion and an alternative approach). From the relationship between $`P`$ and $`\stackrel{~}{\mathrm{\Pi }}`$, it is apparent that the insulator Ward identities (30) for $`\mathrm{\Pi }`$ and $`\stackrel{~}{\mathrm{\Pi }}`$ hold for $`P`$ as well.
## VI Linear response to a local perturbation
In this section the properties of the screened potential $`\phi `$ are examined in greater detail for the case of a local perturbing potential, Sham (1966) which by definition has the form
$$\begin{array}{cc}\hfill v(๐ฑ,๐ฑ^{})& =\delta (๐ฑ๐ฑ^{})v(๐ฑ),\hfill \\ \hfill v(๐ค,๐ค^{})& =v(๐ค๐ค^{}).\hfill \end{array}$$
(55)
Here the spin indices were omitted because a hermitian, time-reversal invariant, local potential must be a spin scalar (see Sec. IX). Contributions from the nonlocal part of the perturbation are considered in Sec. IX.
### VI.1 Screened potential
With this simplification, all of the polarization matrices can be reduced to the scalar form (53), and the screened potential (40) simplifies to Sham (1966)
$$\phi (๐ช)=v(๐ช)+v_c(๐ช)\underset{๐}{}\mathrm{\Pi }(๐ช,๐ช+๐)v(๐ช+๐).$$
(56)
Likewise, the local version of Eq. (54) is
$$\begin{array}{c}\mathrm{\Pi }(๐ช+๐,๐ช+๐^{})=P(๐ช+๐,๐ช+๐^{})\hfill \\ \hfill +P(๐ช+๐,๐ช)v_c(๐ช)\mathrm{\Pi }(๐ช,๐ช+๐^{}).\end{array}$$
(57)
It is convenient at this point to define a macroscopic static electronic dielectric function Jon (b)
$$\begin{array}{cc}\hfill ฯต(๐ค)& =1v_c(๐ค)P(๐ค,๐ค)\hfill \\ & =1/[1+v_c(๐ค)\mathrm{\Pi }(๐ค,๐ค)],\hfill \end{array}$$
(58)
which may be used to express $`\mathrm{\Pi }`$ as a function of the regular polarization $`P`$:
$`\mathrm{\Pi }(๐ช,๐ช+๐)`$ $`=ฯต^1(๐ช)P(๐ช,๐ช+๐),`$ (59a)
$`\mathrm{\Pi }(๐ช+๐,๐ช)`$ $`=ฯต^1(๐ช)P(๐ช+๐,๐ช),`$ (59b)
$`\mathrm{\Pi }(๐ช+๐,๐ช+๐^{})`$ $`=P(๐ช+๐,๐ช+๐^{})+P(๐ช+๐,๐ช)`$
$`\times ฯต^1(๐ช)v_c(๐ช)P(๐ช,๐ช+๐^{}).`$ (59c)
Here all of the nonanalytic behavior at small $`q`$ is contained in the factors $`v_c(๐ช)`$ and $`ฯต^1(๐ช)`$.
Upon substituting (59) into (56), one obtains the screened potentials
$$\phi (๐ช)=\frac{v(๐ช)}{ฯต(๐ช)}+\frac{v_c(๐ช)}{ฯต(๐ช)}\underset{๐\mathrm{๐}}{}P(๐ช,๐ช+๐)v(๐ช+๐),$$
(60)
$$\begin{array}{c}\phi (๐ช+๐)=v(๐ช+๐)+v_c(๐ช+๐)P(๐ช+๐,๐ช)\phi (๐ช)\hfill \\ \hfill +v_c(๐ช+๐)\underset{๐^{}\mathrm{๐}}{}P(๐ช+๐,๐ช+๐^{})v(๐ช+๐^{}),\end{array}$$
(61)
where both expressions are valid for arbitrary $`๐ช`$ and $`๐`$, but the latter is more useful for investigating the behavior of $`\phi `$ in the neighborhood of a nonzero reciprocal lattice vector. For small $`๐ช`$, the first term in (60) is the macroscopic screening that occurs even for slowly varying potentials \[with $`v(๐ค)=0`$ for $`๐ค`$ outside the first Brillouin zone\], while the second term is a local-field correction Louie et al. (1975); Ortiz et al. (1990) arising from the microscopic inhomogeneity of the reference crystal.
### VI.2 Power series expansions
The next step is to establish the small-$`q`$ properties of $`P`$. Since the only singular Coulomb terms in $`P(๐ช+๐,๐ช+๐^{})`$ are the factors $`v_c(๐ช+๐^{\prime \prime })`$ with $`๐^{\prime \prime }\mathrm{๐}`$, the regular polarization $`P(๐ช+๐,๐ช+๐^{})`$ is analytic for $`q<G_{\text{min}}`$, where $`G_{\text{min}}`$ is the magnitude of the smallest nonzero reciprocal lattice vector. In addition, the symmetry property (144) gives $`P(๐ฑ,๐ฑ^{})=P(๐ฑ^{},๐ฑ)`$ or
$$P(๐ค,๐ค^{})=P(๐ค^{},๐ค),$$
(62)
while the Ward identity (30) for an insulator implies that
$$\underset{q0}{lim}P(๐ช,๐ช+๐)=0.$$
(63)
From these results, we see that the matrix $`P_{\mathrm{๐๐}^{}}(๐ช)P(๐ช+๐,๐ช+๐^{})`$ has the Taylor series expansion
$`P_{\mathrm{๐๐}}(๐ช)`$ $`=P_{\mathrm{๐๐}}^{[2]}+P_{\mathrm{๐๐}}^{[4]}+P_{\mathrm{๐๐}}^{[6]}+O(q^8),`$ (64a)
$`P_{\mathrm{๐}๐}(๐ช)`$ $`=P_{\mathrm{๐}๐}^{[1]}+P_{\mathrm{๐}๐}^{[2]}+P_{\mathrm{๐}๐}^{[3]}+P_{\mathrm{๐}๐}^{[4]}+O(q^5),`$ (64b)
$`P_{\mathrm{๐๐}^{}}(๐ช)`$ $`=P_{\mathrm{๐๐}^{}}^{[0]}+P_{\mathrm{๐๐}^{}}^{[1]}+P_{\mathrm{๐๐}^{}}^{[2]}+O(q^3),`$ (64c)
with $`P_{\mathrm{๐๐}}(๐ช)=P_{\mathrm{๐},๐}(๐ช)`$. Here $`P_{\mathrm{๐๐}^{}}^{[l]}`$ denotes a general polynomial of order $`l`$ in the Cartesian components of $`๐ช`$; for example,
$$P_{\mathrm{๐๐}}^{[2]}=P_{\mathrm{๐๐}}^{[2]}(๐ช)=q_\alpha q_\beta P_{\mathrm{๐๐}}^{\alpha \beta },$$
(65)
in which $`P_{\mathrm{๐๐}}^{\alpha \beta }`$ is a constant, and a sum over the Cartesian components $`\alpha `$ and $`\beta `$ is implicit.
In the limit of small $`q`$, the dielectric function of an insulator therefore tends toward a finite but direction-dependent limit:
$$\underset{q0}{lim}ฯต(๐ช)ฯต(\widehat{๐ช})=14\pi P_{\mathrm{๐๐}}^{\alpha \beta }\frac{q_\alpha q_\beta }{q^2},$$
(66)
in which $`\widehat{๐ช}=๐ช/q`$. This behavior contributes nonanalytic terms in the small-$`q`$ expansion of $`ฯต^1(๐ช)`$:
$$\begin{array}{c}\frac{1}{ฯต(๐ช)}=\frac{1}{ฯต(\widehat{๐ช})}[1+w_c(๐ช)(P_{\mathrm{๐๐}}^{[4]}+P_{\mathrm{๐๐}}^{[6]})+w_c^2(๐ช)(P_{\mathrm{๐๐}}^{[4]})^2]\hfill \\ \hfill +O(q^6),\end{array}$$
(67)
in which $`w_c(๐ช)=v_c(๐ช)/ฯต(\widehat{๐ช})`$.
### VI.3 Pseudopotential
To proceed further it is necessary to make some assumptions about the perturbing pseudopotential $`v(๐ค)`$. This is taken to be a superposition of spherically symmetric ionic perturbations, each of which has the Gaussian form used in Refs. Bachelet et al., 1982; Goedecker et al., 1996; Hartwigsen et al., 1998. The pseudopotential can therefore be written as
$$v(๐ค)=v_{\text{an}}(๐ค)+v_c(๐ค)\rho (๐ค),$$
(68)
in which $`v_{\text{an}}(๐ค)`$ and $`\rho (๐ค)`$ are entire analytic functions of $`๐ค`$, and $`\rho (๐ค)`$ represents a portion of the pseudocharge density of the perturbation \[the other portion being given by $`k^2v_{\text{an}}(๐ค)/4\pi `$\].
For a single ion, these functions have the form of a Gaussian times a polynomial in $`k^2`$: Goedecker et al. (1996); Hartwigsen et al. (1998)
$$\begin{array}{cc}\hfill \rho (๐ค)& =\frac{Z_v}{\mathrm{\Omega }}[1+\frac{1}{2}(kr_0)^2]e^{(kr_0)^2/2},\hfill \\ \hfill v_{\text{an}}(๐ค)& =\frac{1}{\mathrm{\Omega }}[2\pi Z_vr_0^2+g(k^2r_0^2)]e^{(kr_0)^2/2},\hfill \end{array}$$
(69)
where $`Z_v`$ is the charge of the ion, $`r_0`$ is a core radius parameter, and $`g(x)`$ is a cubic polynomial given in Refs. Goedecker et al., 1996 and Hartwigsen et al., 1998. Here $`\rho (๐ค)`$ has been defined in such a way that its Taylor series contains no term proportional to $`k^2`$:
$$\rho (๐ค)=\rho _0+\rho _4k^4+\rho _6k^6+\mathrm{}.$$
(70)
Therefore, the only term in $`v_c(๐ค)\rho (๐ค)`$ that does not vanish in the limit $`k0`$ is the divergent term $`4\pi Z_v/\mathrm{\Omega }k^2`$. This term has been eliminated Yin and Cohen (1982) at $`k=0`$ \[by the definition (34) of $`v_c(๐ค)`$\] because the ion is assumed to be accompanied by $`Z_v`$ electrons, so that the crystal remains neutral after the perturbation.
If the perturbation contains more than one ion (e.g., the quasiatoms defined in Appendix A), $`\rho (๐ค)`$ is just a general Taylor series
$$\rho (๐ค)=\rho _0+\rho ^{[1]}+\rho ^{[2]}+\mathrm{},$$
(71)
although the linear response can always be treated as a superposition of individual ions.
For $`๐\mathrm{๐}`$, $`v(๐ช+๐)`$ is analytic for $`q<G`$, with the Taylor series
$$v(๐ช+๐)=v_๐^{[0]}+v_๐^{[1]}+v_๐^{[2]}+v_๐^{[3]}+O(q^4),$$
(72)
in which $`v_๐^{[0]}=v(๐)`$. A similar expansion is valid for $`v_c(๐ช+๐)`$.
### VI.4 Effective macroscopic density
It is now convenient to rewrite Eq. (60) in a form modeled after the familiar expressions for screening in a homogeneous system: not (a)
$$\phi (๐ช)=v_{\text{an}}(๐ช)+\frac{v_c(๐ช)\overline{n}(๐ช)}{ฯต(๐ช)},$$
(73)
in which
$$\overline{n}(๐ช)=\rho (๐ช)+P_{\mathrm{๐๐}}(๐ช)v_{\text{an}}(๐ช)+\underset{๐\mathrm{๐}}{}P_{\mathrm{๐}๐}(๐ช)v(๐ช+๐)$$
(74)
is an effective macroscopic electron density, which is analytic for $`q<G_{\text{min}}`$. This includes the (partial) bare charge $`\rho (๐ช)`$, the macroscopic charge induced by $`v_{\text{an}}(๐ช)`$, and the local-field corrections. Note that since $`\phi (๐ช)=v(๐ช)+v_c(๐ช)n^{(1)}(๐ช)`$, one can also write $`\overline{n}(๐ช)=[\rho (๐ช)+n^{(1)}(๐ช)]ฯต(๐ช)`$. Also note that $`\overline{n}_0=\rho _0`$, because the second term in (74) is $`O(q^2)`$ and the last term is $`O(q)`$. The leading contributions to $`\phi (๐ช)`$ for small $`q`$ are therefore
$$\begin{array}{c}\phi (๐ช)=\rho _0[w_c(๐ช)+w_c^2(๐ช)(P_{\mathrm{๐๐}}^{[4]}+P_{\mathrm{๐๐}}^{[6]})+w_c^3(๐ช)(P_{\mathrm{๐๐}}^{[4]})^2]\hfill \\ \hfill +w_c(๐ช)(\overline{n}^{[1]}+\overline{n}^{[2]}+\overline{n}^{[3]}+\overline{n}^{[4]})\\ \hfill +w_c^2(๐ช)P_{\mathrm{๐๐}}^{[4]}(\overline{n}^{[1]}+\overline{n}^{[2]})+v_{\text{an}}(๐ช)+O(q^3),\end{array}$$
(75)
in which $`v_{\text{an}}(๐ช)`$ is to be replaced by its Taylor series expansion.
The first set of terms in (75) is proportional to $`\rho _0=Z_v/\mathrm{\Omega }`$. These terms are just the power series expansion for $`\rho _0v_c(๐ช)/ฯต(๐ช)`$. Such terms are present in general, but they vanish for isovalent perturbations (e.g., Al substituting for Ga in GaAs).
The remaining nonanalytic terms depend on $`\overline{n}^{[l]}`$, where $`l1`$. The symmetry of $`\overline{n}^{[l]}`$ may differ from that of $`\rho ^{[l]}`$. For example, the term $`w_c(๐ช)n^{[1]}(๐ช)`$ would contribute a dipole field if it were present, but $`\rho ^{[1]}`$ vanishes for a spherically symmetric atom. However, since the symmetry of $`P`$ is the same as that of the reference crystal, the symmetry of $`\overline{n}`$ is just the site symmetry at the position of the ionic perturbation (i.e., the maximal common subgroup of the reference crystal space group and the full rotation group at the given atomic site). Therefore, for crystals of sufficiently low symmetry (e.g., wurtzite), $`\overline{n}^{[1]}`$ may contribute a dipole field even though $`\rho ^{[1]}`$ does not.
### VI.5 Special cases
The general expression (75) is quite cumbersome and is unlikely to be used in its entirety for any particular material system. In many cases one would only be interested in retaining terms that are two orders in $`q`$ higher than the lowest nonvanishing term. Thus, for heterovalent perturbations with $`\rho _00`$ (i.e., class I of the Introduction), Eq. (75) could be simplified to
$$\begin{array}{c}\phi (๐ช)=v_{\text{an}}(\mathrm{๐})+w_c(๐ช)(\rho _0+\overline{n}^{[1]}+\overline{n}^{[2]})\hfill \\ \hfill +\rho _0w_c^2(๐ช)P_{\mathrm{๐๐}}^{[4]}+O(q),\end{array}$$
(76)
which contains monopole, dipole, and quadrupole terms, plus a correction to the monopole term describing the wave vector dependence of the dielectric function. For isovalent perturbations in crystals with atomic site symmetry that supports a dipole moment (class II), a suitable approximation would be
$$\begin{array}{c}\phi (๐ช)=v_{\text{an}}(๐ช)+w_c(๐ช)(\overline{n}^{[1]}+\overline{n}^{[2]}+\overline{n}^{[3]})\hfill \\ \hfill +w_c^2(๐ช)P_{\mathrm{๐๐}}^{[4]}\overline{n}^{[1]}+O(q^2),\end{array}$$
(77)
which includes additional octopole terms. Finally, for isovalent perturbations in crystals with site symmetry that does not support a dipole moment (class III), one would use
$$\begin{array}{c}\phi (๐ช)=v_{\text{an}}(๐ช)+w_c(๐ช)(\overline{n}^{[2]}+\overline{n}^{[3]}+\overline{n}^{[4]})\hfill \\ \hfill +w_c^2(๐ช)P_{\mathrm{๐๐}}^{[4]}\overline{n}^{[2]}+O(q^3).\end{array}$$
(78)
As an explicit example, consider the case of isovalent substitutions in a crystal with the zinc-blende or diamond structure (space group $`T_d^2`$ or $`O_h^7`$), both of which have site symmetry $`T_d`$ at the atomic sites. In this case $`\rho _0=0`$, and the only invariants of order $`q^4`$ or lower are 1, $`q^2`$, $`q_xq_yq_z`$, $`q^4`$, and $`q_x^4+q_y^4+q_z^4`$. Hence, the quadratic terms are isotropic ($`P_{\mathrm{๐๐}}^{\alpha \beta }=P_2\delta _{\alpha \beta }`$, $`\overline{n}^{[2]}=\overline{n}_2q^2`$) and the long-wavelength dielectric function reduces to a constant \[$`ฯต(\widehat{๐ช})=ฯต=14\pi P_2`$\]. The leading contributions to $`\phi (๐ช)`$ may be written as
$$\begin{array}{c}\phi (๐ช)=\frac{4\pi \overline{n}_2}{ฯต}(1\delta _{\mathrm{๐ช๐}})+C_1+C_2q^2+C_3\frac{q_xq_yq_z}{q^2}\hfill \\ \hfill +C_4\frac{q_x^4+q_y^4+q_z^4}{q^2}+O(q^3),\end{array}$$
(79)
where $`C_i`$ is a constant. The terms $`C_3`$ and $`C_4`$ represent octopole and hexadecapole moments, respectively.
### VI.6 Nonzero reciprocal lattice vectors
Turning now to $`\phi (๐ช+๐)`$, Eq. (61) can be written in the condensed notation
$$\phi (๐ช+๐)=R_{\mathrm{๐๐}}(๐ช)\phi (๐ช)+\xi _๐(๐ช),$$
(80)
in which $`R_{\mathrm{๐๐}}(๐ช)`$ and $`\xi _๐(๐ช)`$ are analytic for $`q<G_{\text{min}}`$:
$$R_{\mathrm{๐๐}}(๐ช)=\{\begin{array}{cc}1\hfill & \text{if }๐=\mathrm{๐},\hfill \\ v_c(๐ช+๐)P_{\mathrm{๐๐}}(๐ช)\hfill & \text{if }๐\mathrm{๐},\hfill \end{array}$$
(81)
$$\begin{array}{c}\xi _๐(๐ช)=(1\delta _{\mathrm{๐๐}})[v(๐ช+๐)+v_c(๐ช+๐)\hfill \\ \hfill \times \underset{๐^{}\mathrm{๐}}{}P_{\mathrm{๐๐}^{}}(๐ช)v(๐ช+๐^{})].\end{array}$$
(82)
For $`๐\mathrm{๐}`$, Eqs. (64b) and (72) show that $`R_{\mathrm{๐๐}}(๐ช)`$ is of order $`q`$ or higher. The explicit form of the Taylor series for $`R_{\mathrm{๐๐}}(๐ช)`$ is determined by finding the invariants of the group of the wave vector $`๐`$ in the reference crystal, where different $`๐`$ vectors are treated as inequivalent. Thus, for general $`๐`$, the linear term is nonvanishing. However, the leading term in the Taylor series for $`\xi _๐(๐ช)`$ (with $`๐\mathrm{๐}`$) is a constant.
Hence, for $`๐\mathrm{๐}`$, the nonanalytic terms in $`\phi (๐ช+๐)`$ are at least one order in $`q`$ higher than those in $`\phi (๐ช)`$. In the zinc-blende example discussed above, one has
$$\phi (๐ช+๐)=C_3R_{\mathrm{๐๐}}^\alpha \frac{q_\alpha q_xq_yq_z}{q^2}+\chi _๐(๐ช)+O(q^3),$$
(83)
in which $`\chi _๐(๐ช)`$ is analytic, and $`R_{\mathrm{๐๐}}^\alpha `$ is the linear coefficient in the Taylor series for $`R_{\mathrm{๐๐}}(๐ช)`$. The nonanalytic term in (83) is a hexadecapole moment that is invariant with respect to the group of $`๐`$.
## VII Quadratic response to a local perturbation
To calculate the quadratic density $`n^{(2)}(๐ค)`$, it is helpful to begin by considering the following partial density obtained from the local version of Eqs. (27), (35), and (49):
$$\begin{array}{c}\stackrel{~}{n}^{(2)}(๐ค)=\frac{1}{2}\underset{๐ค_1}{}^{}\underset{๐_1๐_2}{}\stackrel{~}{\mathrm{\Pi }}^{(2)}(๐ค,๐ค_1+๐_1,๐ค๐ค_1+๐_2)\hfill \\ \hfill \times \phi (๐ค_1+๐_1)\phi (๐ค๐ค_1+๐_2),\end{array}$$
(84)
where the summation on $`๐ค_1`$ is limited to the first Brillouin zone of the reference crystal. Here the proper polarization $`\stackrel{~}{\mathrm{\Pi }}^{(2)}(๐ค,๐ค_1,๐ค_2)`$ vanishes unless $`๐ค=๐ค_1+๐ค_2+๐`$, where $`๐`$ is any reciprocal lattice vector. It satisfies the symmetry relations (144), the reduced form of which is
$$\stackrel{~}{\mathrm{\Pi }}^{(2)}(๐ค,๐ค_1,๐ค_2)=\stackrel{~}{\mathrm{\Pi }}^{(2)}(๐ค,๐ค_2,๐ค_1)=\stackrel{~}{\mathrm{\Pi }}^{(2)}(๐ค_1,๐ค,๐ค_2).$$
(85)
It also satisfies the Ward identity (30):
$$\underset{k0}{lim}\stackrel{~}{\mathrm{\Pi }}^{(2)}(๐ค,๐ค_1,๐ค๐ค_1+๐_2)=0.$$
(86)
With these constraints, the Taylor series expansion of the polarization matrix $`\stackrel{~}{\mathrm{\Pi }}_{\mathrm{๐๐}_1๐_2}^{(2)}(๐ค,๐ค_1,๐ค_2)=\stackrel{~}{\mathrm{\Pi }}^{(2)}(๐ค+๐,๐ค_1+๐_1,๐ค_2+๐_2)`$ has a form similar to that given for $`P`$ in Eq. (64):
$$\begin{array}{cc}\hfill \stackrel{~}{\mathrm{\Pi }}_{\mathrm{๐๐๐}}^{(2)}(๐ค,๐ค_1,๐ค_2)& =k_\alpha k_{1\beta }k_{2\gamma }\stackrel{~}{\mathrm{\Pi }}_{\mathrm{๐๐๐}}^{\alpha \beta \gamma }+O(k^4),\hfill \\ \hfill \stackrel{~}{\mathrm{\Pi }}_{\mathrm{๐๐๐}}^{(2)}(๐ค,๐ค_1,๐ค_2)& =k_{1\beta }k_{2\gamma }\stackrel{~}{\mathrm{\Pi }}_{\mathrm{๐๐๐}}^{\beta \gamma }+O(k^3),\hfill \\ \hfill \stackrel{~}{\mathrm{\Pi }}_{\mathrm{๐}๐_1๐_2}^{(2)}(๐ค,๐ค_1,๐ค_2)& =k_\alpha \stackrel{~}{\mathrm{\Pi }}_{\mathrm{๐}๐_1๐_2}^\alpha +O(k^2),\hfill \\ \hfill \stackrel{~}{\mathrm{\Pi }}_{\mathrm{๐๐}_1๐_2}^{(2)}(๐ค,๐ค_1,๐ค_2)& =\stackrel{~}{\mathrm{\Pi }}_{\mathrm{๐๐}_1๐_2}+O(k),\hfill \end{array}$$
(87)
where the order of the leading term is equal to the number of $`๐`$ vectors that are zero. Here $`O(k^n)`$ denotes a term of order $`k^pk_1^qk_2^r`$, where $`n=p+q+r`$.
The partial quadratic density (84) generates a Coulomb potential $`v^{(2)}(๐ช)=v_c(๐ช)\stackrel{~}{n}^{(2)}(๐ช)`$, which is then screened to produce $`\phi ^{(2)}(๐ช)`$ of Eq. (39). This potential is calculated by replacing $`v(๐ช)`$ with $`v^{(2)}(๐ช)`$ in Eqs. (60) and (61). The result may be written as
$$\phi ^{(2)}(๐ช)=\frac{v_c(๐ช)\overline{n}^{(2)}(๐ช)}{ฯต(๐ช)},$$
(88)
where $`\overline{n}^{(2)}(๐ช)=n^{(2)}(๐ช)ฯต(๐ช)`$ is an effective โexternalโ density
$$\overline{n}^{(2)}(๐ช)=\underset{๐}{}R_{\mathrm{๐}๐}(๐ช)\stackrel{~}{n}^{(2)}(๐ช+๐),$$
(89)
in which
$$R_{\mathrm{๐}๐}(๐ช)=\{\begin{array}{cc}1\hfill & \text{if }๐=\mathrm{๐},\hfill \\ P_{\mathrm{๐}๐}(๐ช)v_c(๐ช+๐)\hfill & \text{if }๐\mathrm{๐}.\hfill \end{array}$$
(90)
In Eq. (89), $`\stackrel{~}{n}^{(2)}(๐ช+๐)`$ is given by (84), where $`\phi (๐ค+๐)`$ can be expressed in terms of $`\phi (๐ค)`$ using Eq. (80). The resulting expression for (89) can be written as $`\overline{n}^{(2)}(๐ช)=\overline{n}_A^{(2)}(๐ช)+\overline{n}_B^{(2)}(๐ช)+\overline{n}_C^{(2)}(๐ช)`$, in which
$$\begin{array}{cc}\hfill \overline{n}_A^{(2)}(๐ช)& =\frac{1}{2}\underset{๐ค}{}^{}A(๐ช,๐ค,๐ช๐ค)\phi (๐ค)\phi (๐ช๐ค),\hfill \\ \hfill \overline{n}_B^{(2)}(๐ช)& =\frac{1}{2}\underset{๐ค}{}^{}[B(๐ช,๐ค,๐ช๐ค)\phi (๐ช๐ค)\hfill \\ & +B(๐ช,๐ช๐ค,๐ค)\phi (๐ค)],\hfill \\ \hfill \overline{n}_C^{(2)}(๐ช)& =\frac{1}{2}\underset{๐ค}{}^{}C(๐ช,๐ค,๐ช๐ค).\hfill \end{array}$$
(91)
Here the functions $`A`$, $`B`$, and $`C`$, which are defined in Appendix F, have the Taylor series expansions
$$\begin{array}{cc}\hfill A(๐ค,๐ค_1,๐ค_2)& =k_\alpha k_{1\beta }k_{2\gamma }A_{\alpha \beta \gamma }+O(k^4),\hfill \\ \hfill B(๐ค,๐ค_1,๐ค_2)& =k_\alpha k_{2\gamma }B_{\alpha \gamma }+O(k^3),\hfill \\ \hfill C(๐ค,๐ค_1,๐ค_2)& =k_\alpha C_\alpha +O(k^2).\hfill \end{array}$$
(92)
For $`๐ช`$ values inside the first Brillouin zone, the functions $`A`$, $`B`$, and $`C`$ in (91) are analytic for all $`๐ค`$ values included in the summation, but $`\phi (๐ค)`$ is nonanalytic at $`๐ค=\mathrm{๐}`$. Therefore it is possible that the $`\phi `$ terms in (91) may produce nonanalytic behavior in $`\overline{n}^{(2)}(๐ช)`$. $`\overline{n}_C^{(2)}(๐ช)`$ is obviously analytic in $`๐ช`$, as is the second term in $`\overline{n}_B^{(2)}(๐ช)`$. The first term in $`\overline{n}_B^{(2)}(๐ช)`$ is as well, since a small variation $`\delta ๐ช`$ can be eliminated from $`\phi (๐ช๐ค)`$ with an equal variation $`\delta ๐ค=\delta ๐ช`$. This slightly shifts the zone boundary in the summation, but $`\phi (๐ช๐ค)`$ is analytic at the zone boundary, so $`\overline{n}_B^{(2)}(๐ช)`$ is analytic for small $`๐ช`$.
However, for $`\overline{n}_A^{(2)}(๐ช)`$ this argument is no longer valid. The singularities in $`\phi (๐ค)`$ and $`\phi (๐ช๐ค)`$ merge when $`๐ช=\mathrm{๐}`$, producing nonanalytic behavior in $`\overline{n}_A^{(2)}(๐ช)`$ at this point. The contribution from the nonanalytic part of $`\overline{n}_A^{(2)}(๐ช)`$ is examined in Appendix G, where it is shown to be negligible under all three approximation schemes defined in the Introduction. Therefore, only the analytic part of $`\overline{n}^{(2)}(๐ช)`$ is retained here.
The leading contributions to the quadratic screened potential are therefore
$$\phi ^{(2)}(๐ช)=w_c(๐ช)(q_\alpha \overline{n}_\alpha ^{(2)}+q_\alpha q_\beta \overline{n}_{\alpha \beta }^{(2)})+O(q),$$
(93)
where $`\overline{n}_\alpha ^{(2)}`$ and $`\overline{n}_{\alpha \beta }^{(2)}`$ are Taylor series coefficients for the analytic part of $`\overline{n}^{(2)}(๐ช)`$. The absence of a constant term in the power series for $`\overline{n}^{(2)}(๐ช)`$ is a consequence of the Ward identity (86). Equation (93) is used in its full form only for isovalent class III perturbations. For class II, the $`\overline{n}_{\alpha \beta }^{(2)}`$ term is negligible, while for class I, the entire contribution from $`\phi ^{(2)}(๐ช)`$ is negligible. not (b)
In the vicinity of a nonzero reciprocal lattice vector, $`\phi ^{(2)}(๐ช+๐)`$ can be written in a form similar to (80):
$$\phi ^{(2)}(๐ช+๐)=R_{\mathrm{๐๐}}(๐ช)\phi ^{(2)}(๐ช)+\xi _๐^{(2)}(๐ช),$$
(94)
in which $`\xi _๐^{(2)}(๐ช)`$ is given by Eq. (82) with $`v(๐ค)`$ replaced by $`v^{(2)}(๐ค)=v_c(๐ค)\stackrel{~}{n}^{(2)}(๐ค)`$. Using the same type of analysis as before, one finds that the nonanalytic part of $`\xi _๐^{(2)}(๐ช)`$ is $`O(q)`$ for class I, $`O(q^3)`$ for class II, and $`O(q^5)`$ for class III. Therefore, the limit $`\xi _๐^{(2)}(\mathrm{๐})`$ is well defined, and the leading terms in $`\phi ^{(2)}(๐ช+๐)`$ are given by
$$\phi ^{(2)}(๐ช+๐)=w_c(๐ช)(q_\alpha q_\beta \overline{n}_\alpha ^{(2)}R_{\mathrm{๐๐}}^\beta )+\xi _๐^{(2)}(\mathrm{๐})+O(q).$$
(95)
Since the leading terms here are $`O(q^0)`$, this contribution is negligible for classes I and II.
Note that the quadratic response for a heterostructure cannot be written as a superposition of spherically symmetric atomic perturbations; one must also include diatomic perturbations with axial symmetry $`C_\mathrm{}v`$ (for a heteronuclear diatomic molecule) or $`D_\mathrm{}h`$ (for a homonuclear diatomic molecule). Tinkham (1964) The symmetry of $`\overline{n}^{(2)}(๐ช)`$ is determined by the maximal common subgroup of the reference crystal space group and the molecular point group. For example, for a perturbation at neighboring atomic sites in zinc-blende, the symmetry of $`\overline{n}^{(2)}(๐ช)`$ is $`C_{3v}`$, which supports a nonvanishing dipole moment $`\overline{n}_\alpha ^{(2)}`$.
In general, $`\overline{n}_\alpha ^{(2)}`$ is nonvanishing for any heteronuclear perturbation, because $`C_\mathrm{}v`$ itself permits the existence of a dipole. Such dipoles therefore always appear in heterostructures involving more than one type of atomic perturbation (e.g., InAs/GaSb). (This property of the nonlinear response was deduced from numerical calculations of band offsets in Ref. Dandrea et al., 1992.) Furthermore, the quadrupole term $`\overline{n}_{\alpha \beta }^{(2)}\frac{1}{3}\overline{n}_{\lambda \lambda }^{(2)}\delta _{\alpha \beta }`$ is nonvanishing for any diatomic perturbation, since isotropy requires cubic symmetry.
## VIII Self-energy for a local perturbation
### VIII.1 Linear terms
The above results may now be used to calculate the self-energy $`\mathrm{\Sigma }`$ and the total potential $`V`$ defined in Eqs. (12) and (42). The total linear potential (42) is given for the case of a local perturbation by
$$V_{ss^{}}^{(1)}(๐ค+๐,๐ค^{}+๐^{};\omega )=\underset{๐^{\prime \prime }}{}\stackrel{~}{\mathrm{\Gamma }}_{ss^{}}^{(1)}(๐ค+๐,๐ค^{}+๐^{};๐ช+๐^{\prime \prime };\omega )\phi (๐ช+๐^{\prime \prime }),$$
(96)
where $`๐ช=๐ค๐ค^{}`$. Here $`\phi (๐ช+๐^{\prime \prime })`$ can be expressed in terms of $`\phi (๐ช)`$ using Eq. (80); this yields
$$V_{ss^{}}^{(1)}(๐ค+๐,๐ค^{}+๐^{};\omega )=\mathrm{\Lambda }_{ss^{}}(๐ค,๐ค^{};๐,๐^{};\omega )\phi (๐ช)+W_{ss^{}}^{(1)}(๐ค,๐ค^{};๐,๐^{};\omega ),$$
(97)
in which
$$W_{ss^{}}^{(1)}(๐ค,๐ค^{};๐,๐^{};\omega )=\underset{๐^{\prime \prime }\mathrm{๐}}{}\stackrel{~}{\mathrm{\Gamma }}_{ss^{}}^{(1)}(๐ค+๐,๐ค^{}+๐^{};๐ช+๐^{\prime \prime };\omega )\xi _{๐^{\prime \prime }}(๐ช)$$
(98)
is an analytic function of $`๐ค`$ and $`๐ค^{}`$ (and therefore also of $`๐ช`$). The nonanalytic terms are all contained in the screened potential $`\phi (๐ช)=v_{\text{an}}(๐ช)+v_c(๐ช)\overline{n}(๐ช)/ฯต(๐ช)`$, which is multiplied by the effective macroscopic vertex function
$$\mathrm{\Lambda }_{ss^{}}(๐ค,๐ค^{};๐,๐^{};\omega )=\underset{๐^{\prime \prime }}{}\stackrel{~}{\mathrm{\Gamma }}_{ss^{}}^{(1)}(๐ค+๐,๐ค^{}+๐^{};๐ช+๐^{\prime \prime };\omega )R_{๐^{\prime \prime }\mathrm{๐}}(๐ช).$$
(99)
Since $`\mathrm{\Lambda }_{ss^{}}`$ is an analytic function of $`๐ค`$ and $`๐ค^{}`$, it can be expanded in a Taylor series \[treating $`๐ช=๐ค๐ค^{}`$ and $`๐=\frac{1}{2}(๐ค+๐ค^{})`$ as the independent variables\], with the result
$$\begin{array}{c}\mathrm{\Lambda }_{ss^{}}(๐ค,๐ค^{};๐,๐^{};\omega )=[\delta _{ss^{}}\delta _{\mathrm{๐๐}^{}}\mathrm{\Sigma }_{ss^{}}^{(0)}(๐,๐^{};\omega )/\omega ]+q_\alpha \mathrm{\Lambda }_{ss^{}}^{(\alpha |)}(๐,๐^{};\omega )+Q_\alpha \mathrm{\Lambda }_{ss^{}}^{(|\alpha )}(๐,๐^{};\omega )\hfill \\ \hfill +q_\alpha q_\beta \mathrm{\Lambda }_{ss^{}}^{(\alpha \beta |)}(๐,๐^{};\omega )+Q_\alpha Q_\beta \mathrm{\Lambda }_{ss^{}}^{(|\alpha \beta )}(๐,๐^{};\omega )+q_\alpha Q_\beta \mathrm{\Lambda }_{ss^{}}^{(\alpha |\beta )}(๐,๐^{};\omega )+O(q^3).\end{array}$$
(100)
Here the lowest-order term is determined by the Ward identity (20), and the Taylor series coefficients such as $`\mathrm{\Lambda }_{ss^{}}^{(\alpha |\beta )}`$ are given in Appendix H. The analytic potential (98) can be expanded in the same way:
$$\begin{array}{c}W_{ss^{}}^{(1)}(๐ค,๐ค^{};๐,๐^{};\omega )=W_{ss^{}}^{(1)}(\mathrm{๐},\mathrm{๐};๐,๐^{};\omega )+q_\alpha W_{ss^{}}^{(\alpha |)}(๐,๐^{};\omega )+Q_\alpha W_{ss^{}}^{(|\alpha )}(๐,๐^{};\omega )\hfill \\ \hfill +q_\alpha q_\beta W_{ss^{}}^{(\alpha \beta |)}(๐,๐^{};\omega )+Q_\alpha Q_\beta W_{ss^{}}^{(|\alpha \beta )}(๐,๐^{};\omega )+q_\alpha Q_\beta W_{ss^{}}^{(\alpha |\beta )}(๐,๐^{};\omega )+O(q^3).\end{array}$$
(101)
In the expansion (100), a term such as $`q_\alpha `$ appears in the total potential (97) as a multiplicative factor in front of the screened potential $`\phi (๐ช)`$. In coordinate space, this term therefore takes the gradient of $`\phi (๐ฑ)`$, generating the $`\alpha `$ component of the macroscopic electric field produced by the perturbation $`v(๐ฑ)`$. Likewise, the term $`Q_\alpha `$ has the form of a (symmetrized) crystal momentum operator that acts upon the envelope functions in an effective-mass theory. The Taylor series (101) for the analytic potential (98) is interpreted in the same way, except that these terms produce only short-range localized potentials because they are analytic functions of $`๐ช`$.
The various terms in Eq. (100) therefore give rise to various long-range spin-dependent potentials whose particular form is determined by the symmetry of the coefficients $`\mathrm{\Lambda }_{ss^{}}^{(\alpha |)}`$, etc. The specific term that generates the long-range Rashba effect is $`\mathrm{\Lambda }_{ss^{}}^{(\alpha |\beta )}`$, since this term is linear in the electric field $`q_\alpha \phi (๐ช)`$ and linear in the crystal momentum $`Q_\beta `$. The usual short-range part of the Rashba coupling is generated by the analogous term $`W_{ss^{}}^{(\alpha |\beta )}`$ in Eq. (101).
The complete expression for $`V^{(1)}`$ is obtained by inserting the expansion (100) for $`\mathrm{\Lambda }`$ and one of the three expansions (76), (77), or (78) for $`\phi (๐ช)`$ into Eq. (97). For the specific example of isovalent perturbations in zinc-blende materials treated in Eq. (79), one finds
$$\begin{array}{c}V_{ss^{}}^{(1)}(๐ค+๐,๐ค^{}+๐^{};\omega )=\frac{1\delta _{\mathrm{๐ช๐}}}{q^2}\{[\delta _{ss^{}}\delta _{\mathrm{๐๐}^{}}\mathrm{\Sigma }_{ss^{}}^{(0)}(๐,๐^{};\omega )/\omega ][4\pi \overline{n}_2q^2/ฯต+C_3q_xq_yq_z+C_4(q_x^4+q_y^4+q_z^4)]\hfill \\ \hfill +C_3q_xq_yq_z[q_\alpha \mathrm{\Lambda }_{ss^{}}^{(\alpha |)}(๐,๐^{};\omega )+Q_\alpha \mathrm{\Lambda }_{ss^{}}^{(|\alpha )}(๐,๐^{};\omega )]\}+\text{analytic terms}+O(q^3),\end{array}$$
(102)
where the analytic terms include $`W^{(1)}`$ and contributions from the analytic part of $`\phi (๐ช)`$. From this result it can be seen that the Rashba effect in isovalent zinc-blende materials does not include any long-range terms (to within the accuracy of the present approximation scheme), since $`\mathrm{\Lambda }_{ss^{}}^{(\alpha |\beta )}`$ contributes only to $`O(q^3)`$. However, there are other long-range spin-splitting terms of $`O(q^2)`$ or lower, and the Rashba effect does contribute nonnegligible long-range terms for perturbations in classes I and II.
### VIII.2 Quadratic terms
Turning now to the quadratic response, the two contributions to $`V^{(2)}`$ in Eq. (42) will be denoted $`V^{(2\text{a})}`$ and $`V^{(2\text{b})}`$, respectively. The first of these is given by
$$V_{ss^{}}^{(2\text{a})}(๐ค+๐,๐ค^{}+๐^{};\omega )=\frac{1}{2}\underset{๐ค_1}{}^{}\underset{๐_1}{}\underset{๐_2}{}\stackrel{~}{\mathrm{\Gamma }}_{ss^{}}^{(2)}(๐ค+๐,๐ค^{}+๐^{};๐ค_1+๐_1,๐ช๐ค_1+๐_2;\omega )\phi (๐ค_1+๐_1)\phi (๐ช๐ค_1+๐_2).$$
(103)
Upon inserting Eq. (80) for $`\phi (๐ค+๐)`$ into the right-hand side, one obtains an expression for $`V^{(2\text{a})}`$ very similar to that found in Eqs. (91) and (147) for the effective quadratic density $`\overline{n}^{(2)}(๐ช)`$. Just as before, there are both analytic and nonanalytic contributions. The nonanalytic contributions can be evaluated using the method outlined in Appendix G; the results show that the nonanalytic terms in $`V^{(2\text{a})}`$ are $`O(q^1)`$ for class I, $`O(q)`$ for class II, and $`O(q^3)`$ for class III. (An explicit expression for the leading nonanalytic term in class I was given previously by Sham. Sham (1966); Sak (1968)) Therefore, the nonanalytic contributions are negligible in all three cases, and the leading term in $`V^{(2\text{a})}`$ is just a constant:
$$V_{ss^{}}^{(2\text{a})}(๐ค+๐,๐ค^{}+๐^{};\omega )=V_{ss^{}}^{(2\text{a})}(๐,๐^{};\omega )+O(q).$$
(104)
This term is negligible under the approximation schemes for classes I and II.
Finally, the second contribution to $`V^{(2)}`$ in Eq. (42) is given by
$$V_{ss^{}}^{(2\text{b})}(๐ค+๐,๐ค^{}+๐^{};\omega )=\underset{๐^{\prime \prime }}{}\stackrel{~}{\mathrm{\Gamma }}_{ss^{}}^{(1)}(๐ค+๐,๐ค^{}+๐^{};๐ช+๐^{\prime \prime };\omega )\phi ^{(2)}(๐ช+๐^{\prime \prime }),$$
(105)
where $`\phi ^{(2)}(๐ช+๐^{\prime \prime })`$ was given previously in Eq. (94). Inserting this expression into Eq. (105), one obtains
$$V_{ss^{}}^{(2\text{b})}(๐ค+๐,๐ค^{}+๐^{};\omega )=\mathrm{\Lambda }_{ss^{}}(๐ค,๐ค^{};๐,๐^{};\omega )\phi ^{(2)}(๐ช)+W_{ss^{}}^{(2\text{b})}(๐ค,๐ค^{};๐,๐^{};\omega ),$$
(106)
in which $`\mathrm{\Lambda }`$ was defined in Eq. (99), and
$$W_{ss^{}}^{(2\text{b})}(๐ค,๐ค^{};๐,๐^{};\omega )=\underset{๐^{\prime \prime }\mathrm{๐}}{}\stackrel{~}{\mathrm{\Gamma }}_{ss^{}}^{(1)}(๐ค+๐,๐ค^{}+๐^{};๐ช+๐^{\prime \prime };\omega )\xi _{๐^{\prime \prime }}^{(2)}(๐ช).$$
(107)
Unlike the case for Eq. (98), this is not an analytic function of $`๐ช`$. However, as discussed below Eq. (94), the nonanalytic portion is $`O(q)`$ and therefore vanishes at $`q=0`$. An explicit expression for $`V^{(2\text{b})}`$ can now be obtained by inserting the expansion (93) for $`\phi ^{(2)}(๐ช)`$ into Eq. (106):
$$\begin{array}{c}V_{ss^{}}^{(2\text{b})}(๐ค+๐,๐ค^{}+๐^{};\omega )=q_\alpha w_c(๐ช)\{\overline{n}_\alpha ^{(2)}[\delta _{ss^{}}\delta _{\mathrm{๐๐}^{}}\mathrm{\Sigma }_{ss^{}}^{(0)}(๐,๐^{};\omega )/\omega +q_\beta \mathrm{\Lambda }_{ss^{}}^{(\beta |)}(๐,๐^{};\omega )+Q_\beta \mathrm{\Lambda }_{ss^{}}^{(|\beta )}(๐,๐^{};\omega )]\}\hfill \\ \hfill +q_\alpha q_\beta w_c(๐ช)\{\overline{n}_{\alpha \beta }^{(2)}[\delta _{ss^{}}\delta _{\mathrm{๐๐}^{}}\mathrm{\Sigma }_{ss^{}}^{(0)}(๐,๐^{};\omega )/\omega ]\}+W_{ss^{}}^{(2\text{b})}(\mathrm{๐},\mathrm{๐};๐,๐^{};\omega )+O(q).\end{array}$$
(108)
The first and second terms are dipole and quadrupole potentials, respectively, while the last term is just a constant. In class II, only the leading $`O(q^1)`$ dipole term is retained; in class I, the entire expression (108) is neglected.
## IX Nonlocal perturbations
An arbitrary nonlocal potential can be separated into local and nonlocal parts (although this separation is not unique):
$$v_{ss^{}}(๐ฑ,๐ฑ^{})=v_{ss^{}}^{\text{loc}}(๐ฑ,๐ฑ^{})+v_{ss^{}}^{\text{nl}}(๐ฑ,๐ฑ^{}).$$
(109)
Here the local part $`v_{ss^{}}^{\text{loc}}(๐ฑ,๐ฑ^{})`$ has the form of Eq. (55) and is treated according to the methods developed above. This section considers the changes in the preceding expressions that may be necessary in the case of nonlocal perturbations, particularly those involving spin-orbit coupling.
### IX.1 Analytic form
A general nonlocal potential can be written as
$$v_{ss^{}}(๐ฑ,๐ฑ^{})=\delta _{ss^{}}v_0(๐ฑ,๐ฑ^{})+๐_{ss^{}}๐ฏ(๐ฑ,๐ฑ^{}),$$
(110)
where $`v_0`$ is a scalar relativistic potential, $`๐`$ is the Pauli matrix, and $`๐ฏ`$ is a pseudovector (similar to orbital angular momentum) that accounts for spin-orbit coupling. If $`v`$ is hermitian and time-reversal invariant, then $`v_0`$ is real and symmetric, while $`๐ฏ(๐ฑ,๐ฑ^{})=๐ฏ(๐ฑ^{},๐ฑ)`$ is imaginary and antisymmetric. Thus $`๐ฏ`$ can have no local component, and a local time-reversal-invariant potential must be a spin scalar:
$$v_{ss^{}}^{\text{loc}}(๐ฑ,๐ฑ^{})=\delta _{ss^{}}\delta (๐ฑ๐ฑ^{})v_{\text{loc}}(๐ฑ).$$
(111)
In the norm-conserving pseudopotential formalism, the nonlocal part of the ionic pseudopotential $`v_{ss^{}}^{\text{nl}}(๐ฑ,๐ฑ^{})`$ is confined to a small region near the nucleus, typically either having the form of a polynomial times a Gaussian Bachelet et al. (1982); Goedecker et al. (1996); Hartwigsen et al. (1998) or vanishing absolutely outside a core region of radius $`r_c`$. Kerker (1980); Troullier and Martins (1991) As a result, $`v_{ss^{}}^{\text{nl}}(๐ค,๐ค^{})`$ is an entire analytic function of $`๐ค`$ and $`๐ค^{}`$.
It is important to note that this analytic form relies upon a physical approximation. In an all-electron calculation where the pseudopotential approximation is not used, the spin-orbit coupling does in general include a contribution from the long-range Coulomb part of the ionic potential. The choice of an analytic pseudopotential $`v_{ss^{}}^{\text{nl}}(๐ค,๐ค^{})`$ is therefore an approximation, in which the spin-orbit coupling is assumed to be dominated by the contribution from the ionic core. Conventional norm-conserving pseudopotentials incorporate all relativistic corrections of order $`Z^2\alpha ^2`$ (where $`Z`$ is the atomic number and $`\alpha `$ is the fine-structure constant), but neglect various terms of order $`\alpha ^2`$, Kleinman (1980); Bachelet and Schlรผter (1982) including the spin-orbit coupling from the long-range (but slowly varying) Coulomb potential outside the core region. This approximation is used in all that follows. It greatly simplifies the analysis of spin-dependent perturbations, as shown below.
### IX.2 Screening
Consider now the description of screening for a spin-dependent perturbation. The relationship between the total polarization $`\mathrm{\Pi }`$ and the regular polarization $`P`$ was given above in Eq. (54) for a general nonlocal potential. Setting $`๐=\mathrm{๐}`$ in this equation gives
$$\mathrm{\Pi }_{ss^{}}(๐ช;๐ค,๐ค+๐ช+๐^{})=ฯต^1(๐ช)P_{ss^{}}(๐ช;๐ค,๐ค+๐ช+๐^{}),$$
(112)
where $`ฯต(๐ช)`$ is the same scalar dielectric function defined above in Eq. (58). Substituting this result into Eq. (54) then yields
$$\begin{array}{c}\mathrm{\Pi }_{ss^{}}(๐ช+๐;๐ค,๐ค+๐ช+๐^{})=P_{ss^{}}(๐ช+๐;๐ค,๐ค+๐ช+๐^{})\hfill \\ \hfill +P(๐ช+๐,๐ช)\frac{v_c(๐ช)}{ฯต(๐ช)}P_{ss^{}}(๐ช;๐ค,๐ค+๐ช+๐^{}).\end{array}$$
(113)
Equations (112) and (113) replace the scalar equations (59) derived previously.
If the perturbation is now separated into local and nonlocal parts, the linear screened potential (40) can be written in a form similar to (73):
$$\phi _{ss^{}}(๐ค,๐ค^{})=v_{ss^{}}^{\text{an}}(๐ค,๐ค^{})+\delta _{ss^{}}\frac{v_c(๐ช)[\overline{n}(๐ช)+n_{\text{nl}}(๐ช)]}{ฯต(๐ช)},$$
(114)
in which $`v_{ss^{}}^{\text{an}}(๐ค,๐ค^{})=v_{ss^{}}^{\text{nl}}(๐ค,๐ค^{})+\delta _{ss^{}}v_{\text{an}}(๐ช)`$, $`๐ช=๐ค๐ค^{}`$, $`\overline{n}(๐ช)`$ is the effective density (74) for the local potential, and
$$\begin{array}{c}n_{\text{nl}}(๐ช)=\underset{๐ค^{\prime \prime }}{}\underset{๐}{}P_{\lambda ^{}\lambda }(๐ช;๐ค^{\prime \prime },๐ค^{\prime \prime }+๐ช+๐)\hfill \\ \hfill \times v_{\lambda \lambda ^{}}^{\text{nl}}(๐ค^{\prime \prime }+๐ช+๐,๐ค^{\prime \prime })\end{array}$$
(115)
is a correction to $`\overline{n}(๐ช)`$ from the nonlocal part of the perturbation. Since $`n_{\text{nl}}(๐ช)=O(q)`$ is an analytic function with the same site symmetry as $`\overline{n}(๐ช)`$, this term does not produce any qualitative changes in $`\phi `$; it merely renormalizes $`\overline{n}(๐ช)`$. The only qualitatively new contribution to $`\phi `$ is the spin-dependent term $`v_{ss^{}}^{\text{nl}}(๐ค,๐ค^{})`$ itself, which is analytic.
For wave vectors in the vicinity of a nonzero reciprocal lattice vector, it is convenient to write Eq. (40) in the following alternative form:
$$\phi _{ss^{}}(๐ค+๐,๐ค^{})=R_{\mathrm{๐๐}}(๐ช)\phi _{ss^{}}(๐ค,๐ค^{})+\xi _{ss^{}}^๐(๐ค,๐ค^{}).$$
(116)
Here $`R_{\mathrm{๐๐}}(๐ช)`$ was defined in Eq. (81), and
$$\begin{array}{c}\xi _{ss^{}}^๐(๐ค,๐ค^{})=\delta _{ss^{}}[\xi _๐(๐ช)+(1\delta _{\mathrm{๐๐}})v_c(๐ช+๐)n_{\text{nl}}(๐ช+๐)]\hfill \\ \hfill +v_{ss^{}}^{\text{nl}}(๐ค+๐,๐ค^{})R_{\mathrm{๐๐}}(๐ช)v_{ss^{}}^{\text{nl}}(๐ค,๐ค^{})\end{array}$$
(117)
is a generalization of the function $`\xi _๐(๐ช)`$ defined in Eq. (82). This is an analytic function of $`๐ค`$ and $`๐ค^{}`$ for $`q<G_{\text{min}}`$.
If Eqs. (114) and (116) are inserted into the nonlocal version of Eq. (96) \[i.e., Eq. (42)\], it is apparent that the nonlocal part of the perturbation produces no qualitative change in the total linear potential $`V^{(1)}`$. The only change is a simple renormalization of the analytic and nonanalytic terms in $`V^{(1)}`$.
The same conclusion also holds for the quadratic potential $`V^{(2)}`$. Thus, the correct qualitative form of $`V^{(1)}`$ and $`V^{(2)}`$ can be derived by ignoring the nonlocal (and spin-dependent) part of the perturbing potential, and including spin only in the vertex function $`\mathrm{\Lambda }`$ and the analytic parts of $`V^{(1)}`$ and $`V^{(2)}`$. This is precisely the approach used in Secs. VI, VII, and VIII.
The key to obtaining this simple result is the fact that $`v_{ss^{}}^{\text{nl}}(๐ค,๐ค^{})`$ is analytic. As shown above, this relies upon the approximation of neglecting spin-orbit coupling outside the atomic cores. Such an approximation would also be possible (and even desirable for its simplicity) in an all-electron calculation where the core electrons are treated explicitly.
## X Summary and conclusions
This paper has presented an analysis of the self-energy of an electron in a lattice-matched semiconductor heterostructure for small values of the crystal momentum. A general theory of nonlinear response for nonlocal spin-dependent perturbations was developed in terms of vertex functions and the static polarization, and applied to the case of quadratic response in a periodic insulator at zero temperature. A set of Ward identities was established for nonlocal spin-dependent potentials. The heterostructure perturbation was separated into a local spin-independent part and a nonlocal spin-dependent part, and the contributions from these were analyzed separately. Due to the neglect of spin-orbit coupling outside the atomic cores, the nonlocal part of the potential is analytic in momentum space. As a result, the nonlocal part of the perturbation merely renormalizes the contributions from the local part.
The main results of the paper are presented in Eqs. (97), (104), and (106). The total linear potential (97) has the form $`V^{(1)}=\mathrm{\Lambda }\phi +W^{(1)}`$, in which all of the nonanalytic contributions come from the screened scalar potential $`\phi `$. This has a form $`\phi (๐ช)=v_{\text{an}}(๐ช)+v_c(๐ช)\overline{n}(๐ช)/ฯต(๐ช)`$ similar to that for screening in a homogeneous system, except that the effective density $`\overline{n}(๐ช)`$ has the site symmetry of the perturbation and the macroscopic dielectric function $`ฯต(๐ช)`$ has the symmetry of the reference crystal.
Spin-dependent contributions come from the analytic part $`W^{(1)}`$ and the vertex function $`\mathrm{\Lambda }`$. The vertex function can be expanded in a Taylor series (100), in which $`q_\alpha `$ takes the gradient (in coordinate space) of $`\phi (๐ช)`$, while $`Q_\alpha `$ is a crystal momentum operator in effective-mass theory. The generalized Rashba effect comes from the term linear in $`q_\alpha `$ and $`Q_\beta `$, but there are other spin-splitting contributions from the lower-order terms as well. A more detailed analysis of the various terms is given in the following paper on effective-mass theory. Foreman (2005)
The total quadratic potential of Eqs. (104) and (106) has a similar form $`V^{(2)}=\mathrm{\Lambda }\phi ^{(2)}+W^{(2)}`$, in which $`W^{(2)}`$ is analytic to within the accuracy of the approximation scheme defined in the Introduction. The quadratic screened potential is $`\phi ^{(2)}(๐ช)=v_c(๐ช)\overline{n}^{(2)}(๐ช)/ฯต(๐ช)`$, where the effective external density $`\overline{n}^{(2)}(๐ช)`$ has the site symmetry of a diatomic perturbation in the reference crystal. Due to the Ward identities for an insulator, the leading term in the power series expansion of $`\overline{n}^{(2)}(๐ช)`$ is a dipole term. The Rashba term in the quadratic potential is always negligible under the approximation scheme used here.
The results derived in this paper are used to develop a first-principles effective-mass theory in the following paper. Foreman (2005) The present results are of crucial importance in establishing clearly defined limitations on the validity of this theory. Most previous formulations of effective-mass theory have been based on non-self-consistent empirical pseudopotentials, for which the possibility of long-range Coulomb interactions is not even considered. However, as shown here, long-range potentials arising from nonanalytic terms in the screening potentialโand even the charge density itselfโmust be considered in general.
The omission of such terms is partially justified (to a certain order of approximation) in isovalent zinc-blende systems, where high crystal symmetry eliminates the contributions from dipole and quadrupole terms in the linear response. Baroni et al. (1989) However, it is not fully justified even in zinc-blende, since the leading octopole terms are of a lower order than the position dependence of the effective mass, which is often included in heterostructure effective-mass calculations. The following paper Foreman (2005) accounts for all terms of the same order as the position dependence of the effective mass, including the octopole and hexadecapole potentials derived here in Eqs. (79) and (102).
A pioneering paper by Sham on effective-mass theory for shallow impurity states Sham (1966) dealt with many of the same issues (for local spin-independent potentials), but at a lower order of approximation. In particular, Sham considered only the lowest-order terms in cubic crystals. At this level of approximation, the total polarization can be treated as analytic \[see Eq. (4.8) of Ref. Sham, 1966\], whereas the present Eq. (59) shows that this is no longer true for terms of higher order (such as those needed for the analysis in Ref. Foreman, 2005) or crystals of lower symmetry. The present work provides a systematic framework for extending Shamโs analysis to crystals of general symmetry and terms of arbitrary order.
The value of establishing such a framework is demonstrated by the result (93) derived here for the leading dipole term in the quadratic density response. Sham has stated that the quadratic density response contains no dipole terms, not (c) but the justification for this statement is not clear because no details of his calculation were given. However, numerical evidence to the contrary was subsequently provided by Dandrea, Duke, and Zunger in a first-principles study of band offsets in InAs/GaSb superlattices. Dandrea et al. (1992) They deduced that the calculated difference between the macroscopic interface dipoles for InSb and GaAs interfaces must be a nonlinear effect (because such differences do not exist in linear response theory Baroni et al. (1989) in cubic crystals), but did not inquire further into its origin.
To the authorโs knowledge, the present derivation provides the first direct explanation for their result, and the first demonstration that dipole terms are a general feature of the quadratic density response. The magnitude of such dipoles is smallโcontributing 50 to 100 meV to the band offset of typical no-common-atom heterojunctionsDandrea et al. (1992); Seidel et al. (1997)โbut they play an important role in explaining the experimentally observed asymmetry of band offsets in such systems. Seidel et al. (1997)
As a final note, it is worth drawing attention to a fundamental property of the nonlinear response of insulators that apparently is not widely known. For example, in Refs. Baroni et al., 1989 and Peressi et al., 1990, Baroni et al. have pointed out that โwithin linear response theory, the electronic charge induced by a given perturbation is proportional to the charge of the perturbing potential,โ Peressi et al. (1990) which implies that โwithin linear response theory, isovalent substitutional impurities carry no net charge.โ Baroni et al. (1989) Although the restriction to linear response theory is necessary in general, the results derived here (in Sec. VII) demonstrate that the quadratic density response of an insulator to a charged perturbation also carries no net charge (i.e., it vanishes in the limit of small wave vectors). Indeed, upon replacing Eqs. (30) and (49) with their higher-order generalizations, one finds that this statement remains true for the nonlinear density response (26) of arbitrary order.
This result stems from the Ward identities (30) for the total static polarization and proper polarization (see Sec. V.3) of an insulator at zero temperature. As a consequence of these identities, one can therefore state that in an insulator, the total electronic charge induced by a given perturbation is exactly linearly proportional to the charge of the perturbing potential. Of course, this statement assumes that the system remains insulating over the full range of the perturbation (from zero to full strength); otherwise, the perturbation theory used here is no longer valid.
###### Acknowledgements.
This work was supported by Hong Kong UGC Grant HIA03/04.SC02.
## Appendix A Reducing heterovalent perturbations to isovalent perturbations
In a heterovalent system such as Ge/GaAs, Harrison et al. (1978) the ionic perturbations are from class I. However, it is often possible to reduce such problems to an equivalent class II or III problem, because accumulations of macroscopic charge are energetically unfavorable; therefore, the interfaces grown in real heterojunctions tend to be macroscopically neutral.
For Ge/GaAs, an ideal (110) heterojunction is already neutral, but an ideal (001) or (111) interface would have a large macroscopic interface charge, leading to a large compensating interface free-carrier density. Harrison et al. (1978) Since this is not observed experimentally, the atoms in a real interface are believed to be arranged in one or more layers of mixed composition, such that the net macroscopic interface charge is zero. Baroni et al. (1989); Harrison et al. (1978) (This is similar to the concept of surface reconstruction, but the interface layers differ from the bulk only in chemical composition, not structure.) In such a system it is possible to replace the heterovalent ionic perturbations with a set of equivalent isovalent perturbations, simply by grouping the atoms together in clusters.
The first step is to define quasiatomic building blocks using a modified version of Evjenโs technique. Evjen (1932) Let $`\mathrm{\Omega }_0(๐ซ)`$ be a Wigner-Seitz cell (of the reference crystal) that is centered on position $`๐ซ`$, and let $`๐ฉ`$ be the number of atoms in any primitive cell of the reference crystal. For a given atom $`a`$ at position $`๐ซ_a`$ in the heterostructure, the quasiatomic potential $`\overline{v}_a`$ is defined in terms of the ionic potentials $`v_a^{}`$ for all atoms $`a^{}`$ via
$$\overline{v}_a=\underset{a^{}}{}w_a^{}(a)v_a^{}.$$
(118)
Here $`w_a^{}(a)`$ is a weight factor, defined as $`w_a^{}(a)=1/๐ฉ`$ if atom $`a^{}`$ lies inside $`\mathrm{\Omega }_0(๐ซ_a)`$, $`w_a^{}(a)=0`$ if $`a^{}`$ lies outside $`\mathrm{\Omega }_0(๐ซ_a)`$, and $`w_a^{}(a)=1/m๐ฉ`$ if $`a^{}`$ lies on the surface of $`\mathrm{\Omega }_0(๐ซ_a)`$ \[where $`m`$ is the number of cells $`\mathrm{\Omega }_0(๐ซ_a+๐)`$ that share atom $`a^{}`$, with $`๐`$ any Bravais lattice vector of the reference crystal\]. In a bulk crystal, these quasiatoms are neutral objects with the site symmetry of atom $`a`$ in the reference crystal. Therefore, in a heterostructure, the quasiatoms carry a charge only near the heterojunctions.
For Ge/GaAs, each quasiatomic building block is constructed from $`\frac{1}{2}`$ of the potential for a given ion plus $`\frac{1}{8}`$ of the potential for each of its four nearest neighbors. In a bulk zinc-blende crystal, these quasiatoms have $`T_d`$ symmetry and possess no charge, no dipole moment, and no quadrupole moment. Therefore, in a heterostructure, the quasiatoms have monopole, dipole, and quadrupole moments only near the heterojunctions.
It is assumed here that the ions in the mixed-composition interface layers form a periodic array, so that a two-dimensional superlattice translation symmetry exists in the directions parallel to the junction plane. In this case, one can define a three-dimensional โslab-adaptedโ Heine (1963); Kleinman (1981) unit cell of quasiatoms that is large enough to contain 100% of the ions in the mixed-composition layers. This unit cell has no net charge (since the interface is assumed to be macroscopically neutral), and for some choices of compositional mixing, it may also have no dipole moment. Harrison et al. (1978)
Thus, if one treats these slab-adapted unit cells as the fundamental perturbations, this type of heterovalent class I perturbation can be replaced by an equivalent neutral perturbation from class II or class III. In general the interface cells do have a dipole moment, so the interface perturbations are class II, while the bulk perturbations are class III. However, an interface dipole term of $`O(q^1)`$ is physically equivalent to a bulk quadrupole term of $`O(q^0)`$. Therefore, the approximation scheme defined in the Introduction yields results of the same overall accuracy for both the class II interface and class III bulk perturbations in Ge/GaAs.
## Appendix B Symmetry properties
This appendix considers some symmetry properties of $`G`$ and $`\mathrm{\Sigma }`$. Time-reversal symmetry is developed from the properties of zero- and one-particle states. The vacuum state $`|0`$ is defined to be time-reversal invariant:
$$\widehat{\mathrm{\Theta }}|0=|0,$$
(119)
where $`\widehat{\mathrm{\Theta }}`$ is the antiunitary time-reversal operator. The phase of $`\widehat{\mathrm{\Theta }}`$ may be partially defined by letting the single-particle basis states $`|๐ฑ,s=\widehat{\psi }_s^{}(๐ฑ)|0`$ satisfy
$$\widehat{\mathrm{\Theta }}|๐ฑ,s=(1)^{s1/2}|๐ฑ,s.$$
(120)
This relation is consistent with the operator equation Merzbacher (1998)
$$\widehat{\mathrm{\Theta }}\widehat{\psi }_s^{}(๐ฑ)\widehat{\mathrm{\Theta }}^{}=(1)^{s1/2}\widehat{\psi }_s^{}(๐ฑ).$$
(121)
One may therefore define $`\widehat{\mathrm{\Theta }}`$ over the entire many-particle Fock space by Eqs. (119) and (121).
The next step is to use the identity Sakurai (1994)
$$\beta |\widehat{A}|\alpha =\stackrel{~}{\beta }|\widehat{\mathrm{\Theta }}\widehat{A}\widehat{\mathrm{\Theta }}^{}|\stackrel{~}{\alpha }^{}=\stackrel{~}{\alpha }|\widehat{\mathrm{\Theta }}\widehat{A}^{}\widehat{\mathrm{\Theta }}^{}|\stackrel{~}{\beta },$$
(122)
in which $`\widehat{A}`$ is a linear operator and $`|\stackrel{~}{\alpha }=\widehat{\mathrm{\Theta }}|\alpha `$. If the many-particle Hamiltonian $`\widehat{H}`$ is time-reversal invariant (i.e., $`[\widehat{\mathrm{\Theta }},\widehat{H}]=0`$), one has
$$\widehat{\mathrm{\Theta }}[\widehat{\psi }_s(๐ฑ,\tau )]^{}\widehat{\mathrm{\Theta }}^{}=(1)^{s1/2}\widehat{\psi }_s^{}(๐ฑ,\tau ),$$
(123)
which holds for complex $`\tau `$. From this and Eq. (122) one immediately obtains
$$G_{ss^{}}(๐ฑ,\tau ;๐ฑ^{},\tau ^{})=(1)^{ss^{}}G_{s^{},s}(๐ฑ^{},\tau ^{};๐ฑ,\tau ),$$
(124)
which is the generalization of an ordinary Green-function โreciprocityโ relation Mor to the interacting-particle case. (A similar expression was given in Ref. Noziรจres, 1964, but with the sign term omitted.) Since the change of time variables in (124) does not alter $`\tau \tau ^{}`$, the Fourier transform of (124) is just
$$G_{ss^{}}(๐ฑ,๐ฑ^{},\omega )=(1)^{ss^{}}G_{s^{},s}(๐ฑ^{},๐ฑ,\omega ).$$
(125)
Now $`\widehat{H}`$ is time-reversal invariant if and only if $`h`$ is, which from (120) and (122) implies that
$$h_{ss^{}}(๐ฑ,๐ฑ^{})=(1)^{ss^{}}h_{s^{},s}(๐ฑ^{},๐ฑ).$$
(126)
Equation (6) then shows that $`\mathrm{\Sigma }`$ has the same time-reversal properties as $`G`$:
$$\mathrm{\Sigma }_{ss^{}}(๐ฑ,๐ฑ^{},\omega )=(1)^{ss^{}}\mathrm{\Sigma }_{s^{},s}(๐ฑ^{},๐ฑ,\omega ).$$
(127)
$`G`$ and $`\mathrm{\Sigma }`$ may also satisfy other conditions derived from linear symmetries of $`\widehat{H}`$. Consider the linear many-particle operator $`\widehat{Q}`$ defined by
$$\widehat{Q}=\underset{s,s^{}}{}\widehat{\psi }_s^{}(๐ฑ)q_{ss^{}}(๐ฑ,๐ฑ^{})\widehat{\psi }_s^{}(๐ฑ^{})d^3xd^3x^{},$$
(128)
in which $`q`$ is a linear single-particle operator. $`\widehat{Q}`$ obeys the commutation relations
$`[\widehat{\psi }_s(๐ฑ),\widehat{Q}]`$ $`={\displaystyle \underset{s^{}}{}}{\displaystyle q_{ss^{}}(๐ฑ,๐ฑ^{})\widehat{\psi }_s^{}(๐ฑ^{})d^3x^{}},`$ (129a)
$`[\widehat{\psi }_s^{}(๐ฑ),\widehat{Q}]`$ $`={\displaystyle \underset{s^{}}{}}{\displaystyle \widehat{\psi }_s^{}^{}(๐ฑ^{})q_{s^{}s}(๐ฑ^{},๐ฑ)d^3x^{}}.`$ (129b)
The commutator of any two such operators is another operator with the same form:
$$[\widehat{Q}_1,\widehat{Q}_2]=\widehat{Q}_3;q_3[q_1,q_2].$$
(130)
Now suppose that the Hamiltonian has the symmetry $`[\widehat{K},\widehat{Q}]=[\widehat{H},\widehat{Q}]=0`$. From (130), this is possible only if $`[h,q]=0`$. One can then use Eqs. (1) and (129) and the cyclic property of the trace to show that
$$[G,q]=0,$$
(131)
which further implies that $`[\mathrm{\Sigma },q]=0`$.
This result is applied to lattice translations throughout this paper, and to other space group operations in the following paper. Foreman (2005) Also of interest in the present paper is the spin operator $`๐ฌ_{ss^{}}(๐ฑ,๐ฑ^{})=\frac{1}{2}๐_{ss^{}}\delta (๐ฑ๐ฑ^{})`$, where $`๐`$ is the Pauli spin matrix vector. If $`h`$ is independent of spin (i.e., $`[h,๐ฌ]=0`$), then $`[G,๐ฌ]=0`$, and $`G`$ and $`\mathrm{\Sigma }`$ have the scalar form
$$\mathrm{\Sigma }_{ss^{}}(๐ฑ,๐ฑ^{},\omega )=\delta _{ss^{}}\mathrm{\Sigma }(๐ฑ,๐ฑ^{},\omega ).$$
(132)
However, if $`h`$ includes spin-orbit coupling (which is the case studied here), $`\mathrm{\Sigma }`$ is nondiagonal.
## Appendix C Fourier transforms
The Fourier transforms of the potential with respect to momentum and frequency are defined by
$$V(๐ค,๐ค^{})=\frac{1}{\mathrm{\Omega }}_\mathrm{\Omega }_\mathrm{\Omega }e^{i๐ค๐ฑ}V(๐ฑ,๐ฑ^{})e^{i๐ค^{}๐ฑ^{}}d^3xd^3x^{},$$
(133)
$$V(\zeta _n,\zeta _n^{})=\frac{1}{\beta }_0^\beta _0^\beta e^{i\zeta _n\tau }V(\tau ,\tau ^{})e^{i\zeta _n^{}\tau ^{}}๐\tau ๐\tau ^{}.$$
(134)
Since $`V(\tau ,\tau ^{})=V(\tau \tau ^{})`$, the latter integral is always diagonal in $`n`$:
$$V(\zeta _n,\zeta _n^{})=V(\zeta _n)\delta _{nn^{}}.$$
(135)
For many-variable quantities such as the vertex function $`\mathrm{\Gamma }^{(\nu )}`$, the Fourier integrals for $`\mathrm{\Gamma }^{(\nu )}(๐ค,๐ค^{};๐ช,๐ช^{};\mathrm{})`$ have the same form as (133) for each pair of $`(๐ค,๐ค^{})`$ variables.
For a function of the form $`f(๐ฑ)=f(r)Y_l^m(\widehat{๐ฑ})`$, where $`Y_l^m`$ is a spherical harmonic, the Fourier transform is of the form $`f(๐ค)=f(k)Y_l^m(\widehat{๐ค})`$, where for $`k>0`$,
$$f(k)=\frac{4\pi }{\mathrm{\Omega }}(i)^l_0^{\mathrm{}}r^2f(r)j_l(kr)๐r,$$
(136)
in which $`j_l(kr)`$ is a spherical Bessel function. For the special case $`f(r)=r^n`$, Eq. (6.561.14) of Ref. Gradshteyn and Ryzhik, 1994 gives
$$f(k)=\frac{4\pi }{\mathrm{\Omega }}(i)^l\frac{\sqrt{\pi }\mathrm{\Gamma }[\frac{1}{2}(ln+3)]}{2^{n1}\mathrm{\Gamma }[\frac{1}{2}(l+n)]}k^{n3},$$
(137)
in which $`\mathrm{\Gamma }(z)=_0^{\mathrm{}}e^tt^{z1}๐t`$. Equation (137) is valid for $`k>0`$, $`n>1`$, and $`l>n3`$. Gradshteyn and Ryzhik (1994)
## Appendix D Perturbation theory
The starting point for the perturbation theory used in Sec. IV is the standard formula Abrikosov et al. (1975); Fetter and Walecka (2003)
$$\widehat{A}_H(\tau )=\frac{T_\tau [\widehat{A}_I(\tau )\widehat{๐ฐ}]_0}{\widehat{๐ฐ}_0},$$
(138)
where $`\widehat{A}_H(\tau )`$ is a Heisenberg picture operator, $`\widehat{A}_I(\tau )`$ is the same operator in the interaction picture, $`\widehat{O}`$ denotes the thermal average (2) with respect to $`\widehat{K}`$, $`\widehat{O}_0`$ is a thermal average with respect to $`\widehat{K}_0=\widehat{H}_0\mu \widehat{N}`$, and $`\widehat{๐ฐ}=T_\tau \{\mathrm{exp}[_0^\beta \widehat{H}_1(\tau )๐\tau ]\}`$. If $`\widehat{๐ฐ}`$ is expanded in a power series, terms of equal order in the numerator and denominator can be grouped together as
$$\begin{array}{c}\widehat{A}_H(\tau )=\widehat{A}_I(\tau )_0+\underset{\nu =1}{\overset{\mathrm{}}{}}\frac{(1)^\nu }{\nu !}_0^\beta ๐\tau _1\mathrm{}_0^\beta ๐\tau _\nu \hfill \\ \hfill \times T_\tau [\mathrm{\Delta }\widehat{H}_1(\tau _1)\mathrm{}\mathrm{\Delta }\widehat{H}_1(\tau _\nu )\mathrm{\Delta }\widehat{A}_I(\tau )]_0,\end{array}$$
(139)
where $`\mathrm{\Delta }\widehat{A}_{I,H}(\tau )=\widehat{A}_{I,H}(\tau )\widehat{A}_I(\tau )_0`$. A more compact expression for (139) is
$$\mathrm{\Delta }\widehat{A}_H(\tau )=T_\tau [\mathrm{\Delta }\widehat{A}_I(\tau )\widehat{๐ฒ}]_0,$$
(140)
where $`\widehat{๐ฒ}=T_\tau \{\mathrm{exp}[_0^\beta \mathrm{\Delta }\widehat{H}_1(\tau )๐\tau ]\}`$.
## Appendix E Polarization
The static polarization (26) is defined by
$$\begin{array}{c}\mathrm{\Pi }^{(\nu )}(00^{},11^{},\mathrm{},\nu \nu ^{})\hfill \\ \hfill =_0^\beta ๐\tau _1\mathrm{}_0^\beta ๐\tau _\nu D^{(\nu )}(00^{},11^{},\mathrm{},\nu \nu ^{}),\end{array}$$
(141)
where $`D`$ is the dynamic polarization
$$\begin{array}{c}D^{(\nu )}(00^{},11^{},\mathrm{},\nu \nu ^{})\hfill \\ \hfill =(1)^\nu T_\tau [\mathrm{\Delta }\widehat{\rho }(00^{})\mathrm{\Delta }\widehat{\rho }(11^{})\mathrm{}\mathrm{\Delta }\widehat{\rho }(\nu \nu ^{})]_0.\end{array}$$
(142)
Here a superfluous second time variable has been added (for notational convenience) to the interaction picture operators according to the definition
$$\widehat{\rho }_{ss^{}}(๐ฑ,\tau ;๐ฑ^{},\tau ^{})\widehat{\rho }_{ss^{}}(๐ฑ,๐ฑ^{},\tau \tau ^{}),$$
(143)
where $`\tau ^{}0`$. Since $`D^{(\nu )}`$ is periodic (with period $`\beta `$) in all time variables $`\tau _\lambda `$, but depends on time only via the intervals $`\tau _\lambda \tau _0`$ (for $`\lambda =1,2,\mathrm{},\nu `$), it follows that $`\mathrm{\Pi }^{(\nu )}`$ is independent of time.
By definition, $`D`$ is symmetric with respect to interchange of any pair of operators $`\mathrm{\Delta }\widehat{\rho }`$; thus
$$\begin{array}{c}D^{(\nu )}(\mathrm{},ii^{},\mathrm{},jj^{},\mathrm{})\hfill \\ \hfill =D^{(\nu )}(\mathrm{},jj^{},\mathrm{},ii^{},\mathrm{}).\end{array}$$
(144)
Another constraint on $`D`$ can be derived from time-reversal symmetry using the methods of Appendix B:
$$D^{(\nu )}(00^{},11^{},\mathrm{},\nu \nu ^{})=(1)^sD^{(\nu )}(\overline{0}^{}\overline{0},\overline{1}^{}\overline{1},\mathrm{},\overline{\nu }^{}\overline{\nu }).$$
(145)
Here $`(\overline{\lambda })=(๐ฑ_\lambda ,s_\lambda ,\tau _\lambda )`$ and
$$s=\underset{\lambda =0}{\overset{\nu }{}}(s_\lambda s_\lambda ^{}).$$
(146)
These symmetries are valid for the static polarization $`\mathrm{\Pi }`$ as well (with the time variables omitted).
## Appendix F Functions $`A`$, $`B`$, and $`C`$
The functions $`A`$, $`B`$, and $`C`$ introduced in Eq. (91) are defined by
$$\begin{array}{cc}\hfill A(๐ค,๐ค_1,๐ค_2)& =\underset{๐}{}\underset{๐_1}{}\underset{๐_2}{}R_{\mathrm{๐}๐}(๐ค)\stackrel{~}{\mathrm{\Pi }}^{(2)}(๐ค+๐,๐ค_1+๐_1,๐ค_2+๐_2)R_{๐_1\mathrm{๐}}(๐ค_1)R_{๐_2\mathrm{๐}}(๐ค_2),\hfill \\ \hfill B(๐ค,๐ค_1,๐ค_2)& =\underset{๐}{}\underset{๐_1}{}\underset{๐_2}{}R_{\mathrm{๐}๐}(๐ค)\stackrel{~}{\mathrm{\Pi }}^{(2)}(๐ค+๐,๐ค_1+๐_1,๐ค_2+๐_2)\xi _{๐_1}(๐ค_1)R_{๐_2\mathrm{๐}}(๐ค_2),\hfill \\ \hfill C(๐ค,๐ค_1,๐ค_2)& =\underset{๐}{}\underset{๐_1}{}\underset{๐_2}{}R_{\mathrm{๐}๐}(๐ค)\stackrel{~}{\mathrm{\Pi }}^{(2)}(๐ค+๐,๐ค_1+๐_1,๐ค_2+๐_2)\xi _{๐_1}(๐ค_1)\xi _{๐_2}(๐ค_2).\hfill \end{array}$$
(147)
## Appendix G Nonanalytic terms in $`\overline{n}_A^{(2)}(๐ช)`$
To leading order, the term $`\overline{n}_A^{(2)}(๐ช)`$ in Eq. (91) is
$$\overline{n}_A^{(2)}(๐ช)=\frac{1}{2}A_{\alpha \beta \gamma }q_\alpha \underset{๐ค}{}^{}k_\beta (q_\gamma k_\gamma )\phi (๐ค)\phi (๐ช๐ค)+O(q^3),$$
(148)
in which $`A_{\alpha \beta \gamma }`$ is the Taylor series coefficient (92). The only contribution to Eq. (148) that is of order $`q^2`$ comes from the monopole terms in $`\phi (๐ช)`$. For any perturbation comprising a finite number of atoms, the expansion (75) begins as $`\phi (๐ช)=Z_vw_c(๐ช)/\mathrm{\Omega }+O(q^1)`$, where $`Z_v`$ is the net valence charge of the ionic perturbation.
If this lowest-order term is considered, the value of Eq. (148) can be estimated by extending the summation to all values of $`๐ค`$. For a cubic crystal with scalar $`ฯต`$, the convolution can be performed using the Fourier transforms in Appendix C; the result is
$$\overline{n}_A^{(2)}(๐ช)=\frac{\pi ^2Z_v^2}{8\mathrm{\Omega }ฯต^2}A_{\alpha \beta \gamma }\frac{q_\alpha q_\beta q_\gamma }{q}+O(q^3).$$
(149)
This shows explicitly that the leading nonanalytic term in $`\overline{n}_A^{(2)}(๐ช)`$ is $`O(q^2)`$.
However, this term exists only for the heterovalent perturbations of class I in the Introduction. For the isovalent perturbations ($`Z_v=0`$) of classes II and III, there is no $`O(q^2)`$ term in $`\overline{n}_A^{(2)}(๐ช)`$. A similar analysis shows that for class II, $`\overline{n}_A^{(2)}(๐ช)=O(q^4)`$, while for class III, $`\overline{n}_A^{(2)}(๐ช)=O(q^6)`$. Therefore, according to the approximation scheme adopted in this paper, $`\overline{n}_A^{(2)}(๐ช)`$ is negligible in all three cases.
## Appendix H Vertex function Taylor series
Since the proper vertex function $`\stackrel{~}{\mathrm{\Gamma }}^{(1)}`$ is an analytic function of $`๐ค`$ and $`๐ค^{}`$, it can be expanded in a Taylor series in the variables $`๐ช=๐ค๐ค^{}`$ and $`๐=\frac{1}{2}(๐ค+๐ค^{})`$. This yields an expression similar to Eq. (100):
$$\begin{array}{c}\stackrel{~}{\mathrm{\Gamma }}_{ss^{}}^{(1)}(๐ค+๐,๐ค^{}+๐^{};๐ช+๐^{\prime \prime };\omega )=\stackrel{~}{\mathrm{\Gamma }}_{ss^{}}^{(1)}(๐,๐^{};๐^{\prime \prime };\omega )+q_\alpha \stackrel{~}{\mathrm{\Gamma }}_{ss^{}}^{(\alpha |)}(๐,๐^{};๐^{\prime \prime };\omega )+Q_\alpha \stackrel{~}{\mathrm{\Gamma }}_{ss^{}}^{(|\alpha )}(๐,๐^{};๐^{\prime \prime };\omega )\hfill \\ \hfill +q_\alpha q_\beta \stackrel{~}{\mathrm{\Gamma }}_{ss^{}}^{(\alpha \beta |)}(๐,๐^{};๐^{\prime \prime };\omega )+Q_\alpha Q_\beta \stackrel{~}{\mathrm{\Gamma }}_{ss^{}}^{(|\alpha \beta )}(๐,๐^{};๐^{\prime \prime };\omega )+q_\alpha Q_\beta \stackrel{~}{\mathrm{\Gamma }}_{ss^{}}^{(\alpha |\beta )}(๐,๐^{};๐^{\prime \prime };\omega )+O(q^3).\end{array}$$
(150)
The Taylor series coefficients in Eq. (100) for the effective vertex function $`\mathrm{\Lambda }`$ are therefore given by
$$\begin{array}{cc}\hfill \mathrm{\Lambda }_{ss^{}}^{(\alpha |)}(๐,๐^{};\omega )& =\stackrel{~}{\mathrm{\Gamma }}_{ss^{}}^{(\alpha |)}(๐,๐^{};\mathrm{๐};\omega )+\underset{๐^{\prime \prime }\mathrm{๐}}{}\stackrel{~}{\mathrm{\Gamma }}_{ss^{}}^{(1)}(๐,๐^{};๐^{\prime \prime };\omega )R_{๐^{\prime \prime }\mathrm{๐}}^\alpha ,\hfill \\ \hfill \mathrm{\Lambda }_{ss^{}}^{(|\alpha )}(๐,๐^{};\omega )& =\stackrel{~}{\mathrm{\Gamma }}_{ss^{}}^{(|\alpha )}(๐,๐^{};\mathrm{๐};\omega ),\hfill \\ \hfill \mathrm{\Lambda }_{ss^{}}^{(\alpha \beta |)}(๐,๐^{};\omega )& =\stackrel{~}{\mathrm{\Gamma }}_{ss^{}}^{(\alpha \beta |)}(๐,๐^{};\mathrm{๐};\omega )+\underset{๐^{\prime \prime }\mathrm{๐}}{}[\stackrel{~}{\mathrm{\Gamma }}_{ss^{}}^{(\alpha |)}(๐,๐^{};๐^{\prime \prime };\omega )R_{๐^{\prime \prime }\mathrm{๐}}^\beta +\stackrel{~}{\mathrm{\Gamma }}_{ss^{}}^{(1)}(๐,๐^{};๐^{\prime \prime };\omega )R_{๐^{\prime \prime }\mathrm{๐}}^{\alpha \beta }],\hfill \\ \hfill \mathrm{\Lambda }_{ss^{}}^{(|\alpha \beta )}(๐,๐^{};\omega )& =\stackrel{~}{\mathrm{\Gamma }}_{ss^{}}^{(|\alpha \beta )}(๐,๐^{};\mathrm{๐};\omega ),\hfill \\ \hfill \mathrm{\Lambda }_{ss^{}}^{(\alpha |\beta )}(๐,๐^{};\omega )& =\stackrel{~}{\mathrm{\Gamma }}_{ss^{}}^{(\alpha |\beta )}(๐,๐^{};\mathrm{๐};\omega )+\underset{๐^{\prime \prime }\mathrm{๐}}{}\stackrel{~}{\mathrm{\Gamma }}_{ss^{}}^{(|\beta )}(๐,๐^{};๐^{\prime \prime };\omega )R_{๐^{\prime \prime }\mathrm{๐}}^\alpha ,\hfill \end{array}$$
(151)
in which $`R_{\mathrm{๐๐}}^\alpha `$ and $`R_{\mathrm{๐๐}}^{\alpha \beta }`$ are the Taylor series coefficients for $`R_{\mathrm{๐๐}}(๐ช)`$. In these expressions, some special values of the coefficients for the case $`๐^{\prime \prime }=\mathrm{๐}`$ are given by the Ward identity (20):
$$\begin{array}{cc}\hfill \stackrel{~}{\mathrm{\Gamma }}_{ss^{}}^{(1)}(๐,๐^{};\mathrm{๐};\omega )& =\delta _{ss^{}}\delta _{\mathrm{๐๐}^{}}\frac{\mathrm{\Sigma }_{ss^{}}^{(0)}(๐,๐^{};\omega )}{\omega },\hfill \\ \hfill \stackrel{~}{\mathrm{\Gamma }}_{ss^{}}^{(|\alpha )}(๐,๐^{};\mathrm{๐};\omega )& =\frac{}{k_\alpha }\frac{\mathrm{\Sigma }_{ss^{}}^{(0)}(๐ค+๐,๐ค+๐^{};\omega )}{\omega }|_{๐ค=\mathrm{๐}},\hfill \\ \hfill \stackrel{~}{\mathrm{\Gamma }}_{ss^{}}^{(|\alpha \beta )}(๐,๐^{};\mathrm{๐};\omega )& =\frac{1}{2}\frac{^2}{k_\alpha k_\beta }\frac{\mathrm{\Sigma }_{ss^{}}^{(0)}(๐ค+๐,๐ค+๐^{};\omega )}{\omega }|_{๐ค=\mathrm{๐}}.\hfill \end{array}$$
(152)
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# Young, Jupiter-Mass Objects in Ophiuchus
## 1. Introduction
Free-floating objects with masses comparable to the masses of the most substantial extrasolar planets have been difficult to find and even more difficult to confirm. Several groups have reported sources with masses below 10 M<sub>J</sub> (Martรญn et al. 2004; Lucas et al. 2003). The limited wavelength range of the photometry available for these objects makes the inital source identifications uncertain. All of the candidate objects lie at distances $``$450 pc, where their extreme faintness makes them difficult to confirm spectroscopically. In fact, the only spectrospically โconfirmedโ extremely low mass source from one of these samples (Zapatero Osorio et al. 2002) may be an older foreground object with higher mass (Burgasser et al. 2004). It is well worth continuing to search to search for a sample of young objects with extremely low masses, both because of the clues they provide about star and planet formation and because they could serve as a testbed for ideas about the structure and early evolution of massive planets.
## 2. Experimental Design
The problem of finding Jupiter-mass objects in nearby molecular clouds is more one of identification than detection. Models for very young stellar objects predict that even extremely low mass bodies are within the reach of direct observations in the near and mid-IR (Baraffe et al. 2003; Burrows et al. 2001).
The Spitzer Legacy Program โFrom Molecular Cores to Planet-Forming Disksโ (c2d Evans et al. 2003) provides mid-IR fluxes (in the \[3.6 and \[4.5\] micron bands) for objects with theoretical masses down to 2 M<sub>J</sub>. Figure 1 illustrates, however, that using IRAC colors alone, one cannot distinguish young, low-mass objects from background stars and galaxies. Our survey of part of the Ophiuchus cloud in I,J,H, and Ks using MOSAIC II and ISPI on the Blanco 4m telescope has 10$`\sigma `$ limits of I=23.5, J=20, H=19, and Ks=18.5. These limits allow us (based on theoretical isochrones) to detect 2 M<sub>J</sub>, 10<sup>6</sup> year old objects even in the presence of modest extinction (A$`{}_{\mathrm{V}}{}^{}<`$10). The fluxes between 0.8 and 3.5 $`\mu `$m, where our survey should be complete for all 2 M<sub>J</sub> sources and where most of the flux from these sources emerges, provide us with a way to build enough color-color and color-magnitude spaces to break the degeneracy between our target population and the myriad of contaminants.
## 3. A Sample of Candidate Young, Jupiter-Mass Objects
In our first-round analysis of $``$0.5 sq degrees in Ophiuchus, we start with 19,000 objects detected at $`>5\sigma `$ in all 5 bands used for our cuts: I,J,H,Ks, and \[3.6\]. In our current pass through the sample, we use a set of empirical criteria based on the nominal colors and magnitudes of a 10 M<sub>J</sub>, 10<sup>6</sup> year old object (Chabrier et al. 2000). We eliminate all sources with J$`<15.09`$ and I-J$`<2.94`$, thereby removing sources that are too bright either because they are foreground objects, luminous background objects, or galaxies with blue I-J colors. Of the remaining 6,000 sources that are faint in J and red in I-J, most are reddened background M stars. We look at this reduced sample in the IJH and IJK color planes and deredden all sources back to the theoretical main sequence for 10<sup>6</sup> year old objects. At this point, only 50 objects have I-J$`>2.94`$. Since brown dwarfs get monotonically redder in K-L as spectral types get later (Golimowski et al. 2004), we cut the sample further by requiring that dereddened K-\[3.6\]$`>0.46`$, leaving us with 37 total candidate 10<sup>6</sup> year old, 1 to 10 M<sub>J</sub> objects. Figure 2 shows the observed colors of one of our candidate objects. Though the colors agree quite well with model predictions for a young 2 M<sub>J</sub> object, higher-mass late M and early L type field brown dwarfs have similar colors (Patten et al. 2004).
Ultimately, we need spectroscopy to confirm that our candidate objects have low gravities, and therefore have low masses. The shape of broad H<sub>2</sub>O and CH<sub>4</sub> absorption bands in near-IR spectra are sensitive to gravity (Lucas et al. 2001), while the relative strengths of the bands can provide spectral types (Geballe et al. 2002). Once we have spectroscopically confirmed objects with masses of 1-10 M<sub>J</sub>, the observed colors of these objects can add confidence to the low-mass nature of other candidate objects in our survey. We will also adjust the models based on the observed spectra and colors of our confirmed objects and tighten our selection criteria for future passes through our data.
Spectroscopy is not the only way to gain confidence in our selection techniques or to produce a subsample with higher reliability. Excess emission in the IRAC bands has been reported around a spectroscopically confirmed 15 M<sub>J</sub> object in Chamaeleon (Luhman et al. 2005). We have recently examined our sample of candidates for evidence of excess mid-IR emission from circum-object disks. Most of our candidates are too faint to be detected in IRAC bands 3 or 4 of the c2d survey, even if they have excess emission from a disk. Among the 37 candidates in our sample, 5 show excess emission in \[5.8\] and/or \[8.0\] compared to the IRAC colors of field brown dwarfs with comparable near-IR colors (Patten et al. 2004). 2 of the 5 sources showing mid-IR excess are also dectected at 24 $`\mu `$m in the c2d MIPS data. Preliminary modeling of the SEDs of our candidates detected at 24 $`\mu `$m indicates that these objects, with possible masses as low as 5 M<sub>J</sub>, could have flared circum-object disks.
### Acknowledgments.
The authors wish to thank Neal Evans and the Cores to Disks Legacy team for providing IRAC data for our regions. We also thank Giovanni Fazio and Brian Patten for providing their IRAC results for field brown dwarfs.
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# Propagation of Avalanches in Mn12-acetate: Magnetic Deflagration
## Abstract
Local time-resolved measurements of fast reversal of the magnetization of single crystals of Mn<sub>12</sub>-acetate indicate that the magnetization avalanche spreads as a narrow interface that propagates through the crystal at a constant velocity that is roughly two orders of magnitude smaller than the speed of sound. We argue that this phenomenon is closely analogous to the propagation of a flame front (deflagration) through a flammable chemical substance.
Mn<sub>12</sub>-acetate (hereafter Mn<sub>12</sub>-ac) is a prototypical molecular magnet composed of magnetic molecules,
\[Mn<sub>12</sub>O<sub>12</sub>(CH<sub>3</sub>COO)<sub>16</sub>(H<sub>2</sub>O)<sub>4</sub>\]$``$2CH<sub>3</sub>COOH$``$4H<sub>2</sub>O,
with cores consisting of twelve Mn atoms strongly coupled by exchange to form superparamagnetic clusters of spin $`S=10`$ at low temperatures Sessoli . Arranged in a centered tetragonal lattice, the spin of the Mn<sub>12</sub> clusters is subject to strong magnetic anisotropy along the symmetry axis (the c-axis of the crystal). Below the blocking temperature of $`3.5`$K, the crystal exhibits remarkable staircase magnetic hysteresis due to resonant quantum spin tunneling between energy levels on opposite sides of the anisotropy barrier corresponding to different spin projections, as illustrated in Fig. 1(a).Friedman This and other interesting properties of Mn<sub>12</sub>-ac have been intensively studied in the last decade (see Refs. review1, ; review2, ; friedmanreview, ; Barco, for reviews).
It has been known for some time Paulsen that Mn<sub>12</sub>-ac crystals exhibit an abrupt reversal of their magnetic moment under certain conditions. This phenomenon, also observed in other molecular magnets, has been attributed to a thermal runaway (avalanche) in which the initial relaxation of the magnetization toward the direction of the field results in the release of heat that further accelerates the magnetic relaxation. Direct measurements of the heat emitted by Mn<sub>12</sub>-ac crystals,Fominaya as well as measurements of the magnetization reversal in pulsed magnetic fields,avalanches have confirmed the thermal nature of the avalanches. More recently, the electromagnetic signal associated with avalanches was detected Tejada ; Tejada2 and it was argued that if the radiation is of thermal origin it would indicate a significant increase in the temperature of the crystal. This has not been confirmed by direct bulk measurements of the temperature using a thermometer. Evidence has been obtained Bal that the avalanche may not be a uniform process throughout the sample. No clear understanding of the avalanche process has emerged to date.
In this Letter we report local time-resolved measurements of fast magnetization reversal (avalanches) in mm-size single crystals of Mn<sub>12</sub>-ac. We show that a magnetic avalanche takes the form of a thin interface between regions of opposite magnetization which propagates throughout the crystal with a constant field-dependent speed ranging from $`1`$ to $`15`$ m/s. We demonstrate that this phenomenon is closely analogous to the propagation of a flame front (deflagration) through a flammable chemical substance.
Microscopic arrays of Hall bars were used to measure the magnetization of three single crystals of Mn<sub>12</sub>-ac with dimensions: sample 1, $`0.29\times 0.29\times 0.64`$ mm<sup>3</sup>; sample 2, $`0.28\times 0.28\times 1.44`$ mm<sup>3</sup>; sample 3, $`0.24\times 0.24\times 1.02`$ mm<sup>3</sup>. Eleven Hall bars of dimensions $`10\times 10`$ $`\mu `$m<sup>2</sup> with $`30`$ $`\mu `$m intervals were used for sample 1, and $`30\times 30`$ $`\mu `$m<sup>2</sup> with $`130`$ $`\mu `$m intervals for samples 2 and 3.
Using an excitation current of $`2`$ $`\mu `$A, the Hall bar signal was amplified by a factor of $`1200`$, and detected and recorded by several digital scopes and a data acquisition card. In order to ensure proper synchronization, one channel of each scope was anchored to the same signal. The Hall sensor and amplifier introduced combined delays of up to $`3`$ $`\mu `$s. A magnetic field was applied in the z-direction along the crystal easy axis (see Fig. 1), lowering (raising) the energy of the states corresponding to spin projections along (opposite to) the field direction. The Hall sensors were aligned to detect the magnetic induction of the sample in the $`x`$-direction. $`B_x`$ is proportional to the spatial derivative of $`M_z`$ in the region near the sensor. For a uniform magnetization in the z-direction, $`B_x`$ derives from the gradient at the sample ends, which is proportional to $`M_z`$ itself. During an avalanche, there is a large contribution to $`B_x`$ from the local region corresponding to the avalanche front, where $`M_z/z`$ is large.
The samples were immersed in liquid <sup>3</sup>He; most of the data were obtained at the base temperature of $`250`$ mK. The few points measured at $`400`$ and $`650`$ mK were found to lie on the same curve within the scatter of the data, indicating that the temperature dependence is weak. A longitudinal magnetic field (parallel to the easy axis) was swept back and forth through the hysteresis loop to $`\pm 6`$ T until an avalanche was triggered. As reported in an earlier paper MMMsuzuki , avalanches occur in a stochastic way at $`0.25`$ K both at resonant magnetic fields (where energy levels on opposite side of the barrier match, see Fig. 1) and away from resonance. Avalanches were also found for sample 2 for zero field-cooled conditions, where the sample starts from zero magnetization (instead of full saturation).
Figure 2 shows an avalanche for sample 1. Steps due to quantum tunneling of the magnetization were observed, with a magnetization that was almost uniform throughout the sample, until an avalanche occured, as shown in the inset. During the avalanche the Hall bar recorded a large peak in $`B_x`$, signaling the abrupt onset of a highly non-uniform magnetization.
For a field sweep rate of $`10`$ mT/s and temperature $`0.25`$ K, Fig. 3 shows an avalanche triggered at $`4`$ T and recorded for sample 1 by seven of the eleven sensors placed in sequential positions near the center of the sample. The avalanche was triggered above the top-most sensor and traveled downward (see Fig. 1). $`B_x`$ displays the largest peak at the center due to the finite size of the sample. The inset shows the sensor position as a function of the time at which the sensor registered the peak amplitude. The slope of the straight line drawn through these points yields a constant velocity of $`12`$ m/s for this avalanche.
Figure 4 summarizes the data obtained for the velocity of propagation of avalanches recorded in different longitudinal magnetic fields for all three samples. For avalanches starting from full magnetization, the data for samples 2 and 3 lie on approximately the same curve. The velocity decreases with decreasing longitudinal magnetic field and goes to zero at about $`0.6`$ T, below which no avalanches can occur. Smaller velocities are obtained for avalanches in sample 2 when starting from the zero-field-cooled condition. Avalanches for sample 1 were obtained only at relatively high magnetic fields in the vicinity of $`4`$ T; the velocities for this sample range in value and do not appear to be consistent with data for the other two samples.
Interestingly, as shown in Fig. 5, an approximate collapse is obtained for all the data when plotted as a function of $`g\mu _BHS(\mathrm{\Delta }M/M_{sat})`$, the energy per molecule released during an avalanche. Thus, avalanches require the release of a threshold energy, above which they propagate with a speed that appears to be a linear function of the energy for the range investigated in these experiments.
Our observations cannot be attributed to magnetization reversal associated with domain wall motion, since there is no long-range order in our system. Some insight can be obtained by noting that, from a thermodynamic point of view, a crystal of Mn<sub>12</sub> molecules placed in a magnetic field opposite to the magnetic moment is equivalent to a metastable (flammable) chemical substance. In our case, the role of the chemical energy stored in a molecule is played by the difference in the Zeeman energy, $`\mathrm{\Delta }E=2g\mu _BHS`$, for states of the Mn<sub>12</sub>-ac molecule that correspond to $`๐`$ parallel and antiparallel to $`๐`$; here $`g=1.94`$ is the gyromagnetic factor and $`\mu _B`$ is the Bohr magneton. For Mn<sub>12</sub>-ac in a field of a few Tesla, $`\mathrm{\Delta }E`$ is below $`0.01`$eV, as is the energy barrier, $`U(H)`$, between spin-up and spin-down states due to the magnetic anisotropy. Thus, for the avalanches in Mn<sub>12</sub>-ac, $`\mathrm{\Delta }E`$ and $`U`$ are two orders of magnitude smaller than typical energies of chemical reactions. However, our temperature range is also more than two orders of magnitude below room temperature, making the analogy rather close.
A well-known mechanism for the release of energy by a metastable chemical substance is combustion or slow burning, technically referred to as deflagration.LL It occurs as a flame of finite width, $`\delta `$, propagates at a constant speed, $`v`$, small compared to the speed of sound. The parameter $`\delta `$ is determined by the distance, $`\delta \sqrt{\kappa \tau }`$ through which the heat diffuses during the time of the โchemical reactionโ $`\tau `$. In our case
$$\tau (H)=\tau _0\mathrm{exp}\left[\frac{U(H)}{k_BT_f}\right],$$
(1)
where $`\tau _010^7`$s is the attempt time Friedman and $`T_f`$ is the temperature of the flame. The dynamics of the flame are governed by the thermal diffusivity, $`\kappa `$, which obeys:
$$\frac{T}{t}=\kappa ^2T.$$
(2)
For $`\kappa `$ independent of $`T`$, substituting $`T=T(xvt)`$ at $`x>vt`$, one obtains $`T=T_f\mathrm{exp}[v(xvt)/\kappa ]`$ in front of the interface, which yields $`v\delta =\kappa `$. An interface thickness that is at most the distance between sensors, $`\delta 30\mu `$m, and the experimentally measured velocities of the order of $`115`$ m/s, yield an upper bound on $`\kappa `$ in the range $`10^5`$m<sup>2</sup>/s to $`10^4`$m<sup>2</sup>/s, consistent with heat pulse experiments.javier
Combining $`v\delta =\kappa `$ with $`\delta \sqrt{\kappa \tau }`$, one obtains:
$$v\frac{\delta }{\tau }\sqrt{\kappa /\tau }=\left(\frac{\kappa }{\tau _0}\right)^{1/2}\mathrm{exp}\left[\frac{U(H)}{2k_BT_f}\right].$$
(3)
The strongest dependence of $`v`$ on $`H`$ derives from the exponential, which contains the known dependencefriedman2 of the energy barrier, $`U(H)`$, on the magnetic field. Assuming that the temperature of the flame is proportional to the released magnetic energy density, $`T_f=CH\mathrm{\Delta }M`$, it is possible to fit all the data with a single value of the proportionality constant $`C`$, as shown in Fig. 4. The flame temperature obtained from these fits ranges from $`8.5`$ K for an avalanche at $`1`$ T to $`26`$ K for an avalanche triggered at $`4`$ T. We note that the prefactor obtained from the fit is consistent with the values of $`\kappa `$ and $`\tau _0`$ discussed earlier.note Here we have used the simplest model of deflagration, a widely studied phenomenon that is known to be quite complex.Combustion This crude model captures the overall behavior, and yields parameters that are quite reasonable in size. A more complete theory is needed to account for the apparent data collapse of Fig, 5.
Recent bolometer measurements of the radiation generated by a magnetic avalanche Tejada ; Tejada2 gave puzzling results that can be understood within our model. One enigma was that the sample temperature measured by a thermometer directly following an avalanche was lower ($`<6`$K) than the temperature registered by the bolometer if one assumed thermal radiation. A second puzzle was that the reversal of the magnetization during the avalanche occured on a much shorter time scale than the cooling of the sample following the avalanche.Tejada2 We suggest that the radiation observed during avalanches is generated by the narrow hot interface (flame) that propagates through the crystal. The temperature of the bulk of the crystal (including the โashโ left behind the interface) is always significantly lower than the temperature of the interface itself. The time of the magnetization reversal is determined by the time $`t=l/v`$ needed for the interface to sweep the sample of length $`l`$. In our case $`t0.1`$ms, while the time needed for the โashโ to reach equilibrium with the thermal bath can be much longer.
The strongest evidence that our observations are due to deflagration is the presence of a well defined propagating front requiring a threshold energy traveling at a subsonic velocity. The deflagration mechanism provides the condition needed for the avalanche to occur. This condition is the same as the condition needed to sustain the propagation of a flame through a chemical substance. It is well-known that deflagration of a flammable gas will not occur in a pipe of diameter, $`d`$, small compared to the width of the flame, $`\delta `$. In our case $`\delta `$ must be small compared to the diameter of the crystal. If this condition is not satisfied, the heat generated by the magnetization reversal diffuses mostly through the walls of the sample and cannot sustain the propagation of the interface. This explains why avalanches only occur in larger crystals with sufficiently large magnetization opposite to the direction of the field. The latter condition coincides with the condition of โflammabilityโ Combustion needed to provide sufficient heating (that is, the large $`T_f`$) required for $`\delta <d`$. It is interesting to note in this connection that the few avalanches that were recorded around $`1`$ T did not result in full reversal of the magnetization. At these low fields the conditions for ignition are marginally satisfied, and the โflameโ is extinguished before the process of magnetization reversal has been completed. In addition to available magnetic energy (flammability), the conditions for ignition may also depend on the shape and quality of the crystal, which may account for differences observed for different samples.
Slow burning at a subsonic speed (deflagration) is governed by the linear process of thermal conductivity. In addition to deflagration, unstable chemical substances also exhibit detonation, which can be caused by instability of the flame or by direct initiation other than through deflagration.Combustion The initial stage of the detonation corresponds to a non-linear supersonic shock wave.LL ; Combustion Theory and experimental studies of advanced stages of detonation are lacking. Based on the close analogy between unstable chemical substances and molecular magnets, the latter may well exhibit โmagnetic detonationโ under the right conditions.
In conclusion, we have demonstrated that avalanches in the magnetization reversal of sufficiently large crystals of magnetic molecules are very similar to flame propagation (deflagration) through a metastable chemical substance. The analogy between the two systems derives from the magnetic bi-stability of molecular nanomagnets. Our observation of โmagnetic deflagrationโ offers a potentially important new way to investigate the phenomenon of flame propagation (and, possibly detonation). In contrast to deflagration in flammable chemical substances, the analogous process of โmagnetic deflagrationโ in molecular nanomagnets is non-destructive, reversible, and much easier to control.
We are grateful to D. Graybill for participation in various aspects of this work, and to K. M. Mertes, J. R. Friedman, M. Bal, and J. Tejada for valuable discussions. Work at City College was supported by NSF grant DMR-0451605. Support was provided for EMC by NSF Grant No. EIA-0310517, and for GC by NSF Grant No. CHE-0071334. E. Z. acknowledges the support of the Israeli Science Foundation Center of Excellence.
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# Thermodynamics of a Fermi liquid in a magnetic field
## Abstract
We extend previous calculations of the non-analytic terms in the spin susceptibility $`\chi _s(T)`$ and the specific heat $`C(T)`$ to systems in a magnetic field. Without a field, $`\chi _s(T)`$ and $`C(T)/T`$ are linear in $`T`$ in $`2D`$, while in $`3D`$, $`\chi _s(T)T^2`$ and $`C(T)/TT^2\mathrm{log}T`$. We show that in a magnetic field, the linear in $`T`$ terms in 2D become scaling functions of $`\mu _BH/T`$. We present explicit expressions for these functions and show that at high fields, $`\mu _BHT`$, $`\chi _s(T,H)`$ scales as $`|H|`$. We also show that in 3D, $`\chi _s(T,H)`$ becomes non-analytic in a field and at high fields scales as $`H^2\mathrm{log}|H|`$.
Landau Fermi liquid theory landau provides the basis for our present understanding of correlated electronic systems. The theory predicts that, in any Fermi liquid, the spin susceptibility $`\chi _s(T)`$ and the specific heat coefficient $`\gamma (T)=C(T)/T`$ tend to a constant at $`T0`$ landau ; AGD . Later, Landau theory has been extended to include the leading temperature dependence of $`\chi _s(T)`$ and $`\gamma (T)`$ which turn out to be non-analytic in dimensions $`D3`$ doniach ; amit ; pethick ; belitz ; millis ; bedell ; baranov ; andrey1 ; dassarma ; andrey2 ; aleiner ; galitski . Like the zero-temperature terms, the thermal corrections come from fermions in the immediate vicinity of the Fermi surface. In 2D systems, both $`\chi _s(T)`$ and $`\gamma (T)`$ are linear in $`T`$ bedell ; baranov ; millis ; andrey1 and the coefficients are expressed in terms of charge and spin components of the scattering amplitude at the scattering angle $`\theta =\pi `$ andrey2 ; aleiner . In 3D, $`\chi _s(T)`$ is quadratic in $`T`$, i.e., is analytic amit ; pethick ; belitz , while $`\gamma (T)`$ is non-analytic and scales as $`T^2\mathrm{log}T`$ doniach ; amit ; pethick .
In this communication, we extend previous works to systems in a magnetic field $`H`$. We consider $`S=1/2`$ charge-less fermions (like $`{}_{}{}^{3}He`$ atoms) for which the magnetic field adds spin-dependent Zeeman term $`\pm \mu _BH`$ to the fermionic dispersion. We show that, in the presence of a field, $`\mathrm{\Delta }\chi _s(T,H)=\chi _s(T,H)\chi _s(0,0)`$ and $`\mathrm{\Delta }\gamma (T,H)=\gamma (T,H)\gamma (0,0)`$ become scaling functions of $`\mu _BH/T`$: $`\mathrm{\Delta }\chi _s(T,H)=\mathrm{\Delta }\chi _s(T,0)f_\chi (\mu _BH/T)`$, $`\mathrm{\Delta }\gamma (T,H)=\mathrm{\Delta }\gamma (T,0)f_\gamma (\mu _BH/T)`$. We present the expressions for these functions to second order in the interaction potential $`U`$. For 2D systems, we show that at $`\mu _BHT`$ (but still, $`\mu _BHE_F`$), $`\mathrm{\Delta }\chi _s(T,H)`$ scales as $`|H|`$ and weakly depends on $`T`$. In the same field range, $`\delta \gamma (T,H)`$ is still linear in $`T`$, but the prefactor is different from that at $`H=0`$. For 3D systems, we show that $`\mathrm{\Delta }\chi _s(T,H)`$ becomes non-analytic at a non-zero $`H`$. The non-analytic term in $`\mathrm{\Delta }\chi _s(T,H)`$ scales as $`H^2\mathrm{log}[max(\mu _BH,T)/E_F]`$. The specific heat coefficient $`\gamma (H,T)`$ in a field still scales as $`T^2\mathrm{log}T`$, but, like in 2D, the prefactor changes between $`H=0`$ and $`\mu _BHT`$.
The analysis of the behavior of $`\mathrm{\Delta }\chi _s(T,H)`$ and $`\mathrm{\Delta }\gamma (T,H)`$ in a magnetic field may be useful for experimental verifications of the non-analytic behavior of thermodynamic variables. It is more straightforward to analyze the dependence on the magnetic field rather than the dependence on the temperature. In particular, recent measurements of the temperature dependence of the spin susceptibility in Si inversion layers R1 didnโt yield conclusive results on whether the $`T`$ dependence of $`\chi _s(T)`$ is indeed linear, as some temperature dependence inevitably comes from spins on the substrate. We propose to measure the field dependence of the spin susceptibility at a given $`T`$ and use our scaling functions to fit the data.
The point of departure for our calculations is the Luttiger-Ward expression for the thermodynamic potential. To simplify the presentation, we assume that the interaction potential $`U(q)`$ is independent on $`q`$. We restore the momentum dependence of $`U(q)`$ in the final formulas. To second-order in $`U`$, the thermodynamic potential is given by
$$\mathrm{\Phi }=\mathrm{\Phi }_0\frac{U^2}{2}T\underset{n}{}_q\frac{d^dq}{(2\pi )^d}\mathrm{\Pi }^{}(\stackrel{}{q},\mathrm{\Omega }_n,T)\mathrm{\Pi }^{}(\stackrel{}{q},\mathrm{\Omega }_n,T),$$
(1)
where $`\mathrm{\Phi }_0`$ is the thermodynamic potential for free fermions, and $`\mathrm{\Pi }^{}(\stackrel{}{q},\mathrm{\Omega }_n,T)`$ and $`\mathrm{\Pi }^{}(\stackrel{}{q},\mathrm{\Omega }_n,T)`$ are the particle-hole bubbles composed of fermions with spin up or spin down, respectively.
Previous studies of the spin susceptibility and the specific heat in a zero magnetic field established that the non-analytic temperature behavior of $`\mathrm{\Delta }\chi _s(T)`$ and $`\mathrm{\Delta }\gamma (T)`$ originates from the non-analyticity of the polarization operator either near $`q=0`$ (Landau damping) doniach ; amit ; pethick ; andrey1 or near $`q=2k_F`$ (a dynamic Kohn anomaly) belitz ; baranov ; millis ; andrey1 ; aleiner ; galitski . The $`2k_F`$ non-analyticity contributes to the spin susceptibility and the specific heat, while the $`q=0`$ non-analyticity only contributes to the non-analyticity in the specific heat. This can be easily understood as the non-analytic term in the zero field spin susceptibility $`\mathrm{\Delta }\chi _s(T)`$ describes a singular response to an infinitesimally small magnetic field. A magnetic field splits Fermi momentum $`k_F`$ into $`k_F^{}`$ and $`k_F^{}`$. The small $`q`$ form of $`\mathrm{\Pi }^{}(\stackrel{}{q},\mathrm{\Omega }_n,T)`$ and $`\mathrm{\Pi }^{}(\stackrel{}{q},\mathrm{\Omega }_n,T)`$ is unaffected by this splitting, up to terms of order $`(\mu _BH/E_F)^2`$, hence the response to the infinitesimal field must be analytic in $`T`$. At the same time, singular $`2k_F`$ contribution to $`\mathrm{\Phi }(T)`$ at zero field originates from the fact that the two polarization operators in Eq. (1) are non-analytic at the same $`q=2k_F`$. In a field the singularities in spin-up and spin-down polarization operators occur at different $`q=2k_F^{}`$ and $`q=2k_F^{}`$. Accordingly, a magnetic field regularizes $`2k_F`$ non-analyticity in the thermodynamic potential, but for a price that the linear response to the field, i.e. the spin susceptibility $`\mathrm{\Delta }\chi (T,H=0)`$, becomes non-analytic.
Our goal is to analyze the forms of the susceptibility and the specific heat at a finite $`H`$, i.e., beyond the linear response theory. We consider the fields for which $`\mu _BH`$ is comparable to $`T`$, but still $`\mu _BHE_F`$. For these fields, the non-analytic contribution to $`\mathrm{\Phi }`$ from small $`q`$ are unaffected by the field . However the $`2k_F`$ contribution is field dependent and evolves at $`\mu _BHT`$.
The calculation of $`\mathrm{\Delta }\mathrm{\Phi }=\mathrm{\Phi }\mathrm{\Phi }_0`$ is somewhat tricky. In principle, all one has to do is to evaluate particle-hole bubbles for fermions with up and down spins at a finite $`T`$, substitute the results into Eq. (1), integrate over momentum $`q`$ and sum over Matsubara frequencies $`\mathrm{\Omega }_n`$. In practice, however, this computation is easy to perform only for small $`q`$ part as for $`qk_F`$, the non-analytic part of the polarization bubble is associated with the Landau damping, which does not depend on $`T`$, apart from regular $`(T/E_F)^2`$ corrections. Accordingly, one can safely use the known analytical forms of $`\mathrm{\Pi }(q,\mathrm{\Omega })`$ at $`T=0`$. For $`q`$ near $`2k_F`$, non-analytic terms in $`\mathrm{\Pi }^{}(\stackrel{}{q},\mathrm{\Omega }_n,T)`$ and $`\mathrm{\Pi }^{}(\stackrel{}{q},\mathrm{\Omega }_n,T)`$ contain scaling functions of $`T/\omega `$, which are only available in integral forms millis . This substantially complicates direct calculation of the $`2k_F`$ term. There exists, however, a way to compute the $`2k_F`$ term, which avoids dealing with the $`2k_F`$ polarization bubbles at a finite $`T`$. This method explores the fact that only the non-analytic parts $`\mathrm{\Pi }^{}(\stackrel{}{q},\mathrm{\Omega }_n,T)`$ and $`\mathrm{\Pi }^{}(\stackrel{}{q},\mathrm{\Omega }_n,T)`$ for $`q`$ near $`2k_F`$ contribute the non-analyticity in the thermodynamic potential. Earlier works have demonstrated that the $`2k_F`$ non-analyticity in $`\mathrm{\Pi }(\stackrel{}{q},\mathrm{\Omega }_n,T)`$ comes from fermions in the particle-hole bubble with momenta near $`\pm \stackrel{}{q}/2`$ belitz ; andrey2 . This implies that, out of four fermions in the second order, two-bubble diagram for the thermodynamic potential in Fig. 1, two fermions with opposite spins have momenta near $`\stackrel{}{q}/2`$, while the other two fermions have momenta near $`\stackrel{}{q}/2`$. Then the $`2k_F`$ part of the $`\mathrm{\Delta }\mathrm{\Phi }`$ can be re-written as the integral over small $`\stackrel{}{k}`$ and small $`\stackrel{}{p}`$ of
$`\mathrm{\Delta }\mathrm{\Phi }_{2k_F}={\displaystyle \frac{U^2}{2}}{\displaystyle \underset{\omega _m,\omega _m^{},\omega _m^{\prime \prime }}{}}{\displaystyle d^2qd^2kd^2p}`$
$`G^{}(\stackrel{}{q}/2+\stackrel{}{k},\omega _m+\omega _m^{})G^{}(\stackrel{}{q}/2,\omega _m^{})\times `$
$`G^{}(\stackrel{}{q}/2+\stackrel{}{k},\omega _m+\omega _m^{\prime \prime })G^{}(\stackrel{}{q}/2+\stackrel{}{p},\omega _m^{\prime \prime })`$ (2)
or, equivalently, as
$$\mathrm{\Delta }\mathrm{\Phi }_{2k_F}=\frac{U^2}{2}\underset{n}{}\frac{d^dq^{}}{(2\pi )^d}\left[\mathrm{\Pi }^{}(\stackrel{}{q^{}},\mathrm{\Omega }_n,T)\right]^2,$$
(3)
where the integration is confined to small $`\stackrel{}{q^{}}=\stackrel{}{k}\stackrel{}{p}`$. In other words, the non-analytic $`2k_F`$ contribution to the thermodynamic potential can be re-expressed in terms of the particle-hole bubble for fermions with opposite spins and a small momentum transfer. The non-analytic term in $`\mathrm{\Pi }`$ at small $`\stackrel{}{q^{}}`$ does not depend on temperature (apart from irrelevant corrections), hence $`\mathrm{\Pi }^{}(\stackrel{}{q^{}},\mathrm{\Omega }_n,T)`$ can be safely approximated by $`\mathrm{\Pi }^{}(\stackrel{}{q^{}},\mathrm{\Omega }_n,0)`$. At the same time, the polarization bubble $`\mathrm{\Pi }^{}(\stackrel{}{q},\mathrm{\Omega }_n,0)`$ strongly depends on the magnetic field (contrary to $`\mathrm{\Pi }^{}(\stackrel{}{q},\mathrm{\Omega }_n,0)`$), and this gives rise to the scaling dependence on $`\mu _BH/T`$.
Combining the $`q=0`$ and $`2k_F`$ contributions, we obtain for the thermodynamic potential
$`\mathrm{\Delta }\mathrm{\Phi }={\displaystyle \frac{U^2}{2}}{\displaystyle \underset{n}{}}{\displaystyle \frac{d^dq}{(2\pi )^d}}`$
$`\left[\left(\mathrm{\Pi }^{}(\stackrel{}{q},\mathrm{\Omega }_n,0)\right)^2+\mathrm{\Pi }^{}(\stackrel{}{q},\mathrm{\Omega }_n,0)\mathrm{\Pi }^{}(\stackrel{}{q},\mathrm{\Omega }_n,0)\right],`$ (4)
where the integration involves only small $`q`$.
We next proceed separately with $`2D`$ and $`3D`$ cases. In 2D we have
$`\mathrm{\Pi }^{}(\stackrel{}{q},\mathrm{\Omega }_n,T)={\displaystyle \frac{m}{2\pi }}{\displaystyle \frac{|\mathrm{\Omega }_n|}{\sqrt{(\mathrm{\Omega }_ni\delta \mu )^2+(\upsilon _Fq)^2}}}+\mathrm{}`$
$`\mathrm{\Pi }^{}(\stackrel{}{q},\mathrm{\Omega }_n,T)={\displaystyle \frac{m}{2\pi }}{\displaystyle \frac{|\mathrm{\Omega }_n|}{\sqrt{\mathrm{\Omega }_n^2+(\upsilon _Fq)^2}}}+\mathrm{},`$ (5)
where dots stand for analytic terms, expandable in powers of $`\mathrm{\Omega }_n^2`$ or $`q^2`$, and $`\delta \mu =\mu _{}\mu _{}=2\mu _BH`$. Substituting Eq. (5) into Eq. (4) and integrating explicitly over momentum $`q`$ we obtain
$$\mathrm{\Delta }\mathrm{\Phi }=\left(\frac{m}{2\pi }\right)^2\frac{U^2T}{8\pi v_F^2}\underset{n}{}\mathrm{\Omega }_{n}^{}{}_{}{}^{2}\mathrm{log}\left[\frac{(\mathrm{\Omega }_n2i\mu _BH)^2\mathrm{\Omega }_n^2}{E_F^4}\right].$$
(6)
Differentiating Eq. (6) with respect to $`H`$, we obtain
$$\mathrm{\Delta }M=\frac{\mathrm{\Delta }\mathrm{\Phi }}{H}=\frac{\mu _{B}^{}{}_{}{}^{4}m^4U^2H^3}{\pi ^3k_F^2}T\underset{n}{}\frac{1}{\mathrm{\Omega }_{n}^{}{}_{}{}^{2}+(2\mu _BH)^2}$$
The sum over Matsubara frequencies can be easily evaluated and yields
$`\mathrm{\Delta }M`$ $`=`$ $`{\displaystyle \frac{\mu _Bm^4U^2}{4\pi ^3k_F^2}}T^2\left[\left({\displaystyle \frac{\mu _BH}{T}}\right)^2\mathrm{coth}({\displaystyle \frac{\mu _BH}{T}})\right]`$ (7)
$`=`$ $`\mu _BAT^2x^2\left[1+2n_B(2x)\right],`$
where
$$A=\frac{m^4U^2}{4\pi ^3k_F^2},x=\frac{\mu _BH}{T}.$$
(8)
We see from Eq. (7) that $`\mathrm{\Delta }M`$ increases in a field by two reasons. First, the field leads to a finite magnetization at $`T=0`$, and second, a finite field populates the system with spin waves precessing at the energy $`\mu _BH`$. Differentiating (7) again over $`H`$, we obtain the spin susceptibility in the form
$$\mathrm{\Delta }\chi (T,H)=\chi (T,H)\chi (0,0)=\mu _B^2ATf_\chi (x),$$
(9)
where
$$f_\chi (x)=\frac{x}{\mathrm{sinh}^2(x)}\left[\mathrm{sinh}(2x)x\right].$$
(10)
For vanishing $`H`$, i.e., at $`x0`$, $`f_\chi (0)=1`$, and
$$\mathrm{\Delta }\chi (T,H)=\chi (T,0)\chi (0,0)=\mu _B^2AT.$$
(11)
This coincides with the earlier result andrey1 . In the opposite limit of large $`x`$, $`f_\chi (x)2x`$, and
$$\mathrm{\Delta }\chi (T,H)=2\mu _B^2AT|x|=2\mu _B^3A|H|.$$
(12)
We see that at high fields, the spin susceptibility scales as $`|H|`$, i.e., is non-analytic in $`H`$.
In Fig. 2 we plot the susceptibility as a function of temperature at a given $`H`$, and as a function of the magnetic field at a given $`T`$. Note that at a finite $`H`$, the Bose term in Eq. (8) gives rise to a negative derivative of $`\mathrm{\Delta }\chi /T`$. This in turn gives rise to a shallow minimum in the temperature dependence of $`\mathrm{\Delta }\chi (T,H)`$.
The specific heat $`\mathrm{\Delta }\gamma (T,H)=\gamma (T,H)\gamma (0,0)`$ is obtained by differentiating Eq. (6) twice over $`T`$. At $`H=0`$, $`\mathrm{\Delta }\gamma (T,H)=6AT\zeta (3)`$ andrey1 ; andrey2 ; aleiner . At a finite $`H`$,
$`\mathrm{\Delta }\gamma (T,H)=6AT\zeta (3)+2AT{\displaystyle _0^{\mu _BH/T}}{\displaystyle \frac{dxx^3}{\mathrm{sinh}^3x}}\times `$
$`\left(x\mathrm{cosh}x\mathrm{sinh}x\right)=ATf_\gamma (x),`$ (13)
where
$`f_\gamma (x)=3\left(Li_3(e^{2x})+2xLi_2(e^{2x})2x^2\mathrm{log}(1e^{2x})\right)`$
$`+6\zeta (3)2x^3+4x^3\mathrm{coth}xx^3{\displaystyle \frac{1}{\mathrm{sinh}^2x}}\left(\mathrm{sinh}2xx\right)`$
and $`Li`$ are polylogarithmic functions. At $`x1`$, $`f_\gamma (x)6\zeta (3)\frac{x^4}{6}`$ and
$$\mathrm{\Delta }\gamma (\frac{\mu _BH}{T}1)AT\left(6\zeta (3)\frac{1}{6}\left(\frac{\mu _BH}{T}\right)^4\right).$$
In the opposite limit of $`x1`$, $`f_\gamma (x)=3\zeta (3)+4x^4e^{2x}`$, and
$$\mathrm{\Delta }\gamma (\frac{\mu _BH}{T}1)AT\left(3\zeta (3)+4\left(\frac{\mu _BH}{T}\right)^4e^{2\mu _BH/T}\right).$$
We see that in both limits the temperature dependence of the specific heat is linear in $`T`$, but the prefactor changes by a factor of $`2`$ between small and high fields. This result could be anticipated as a high magnetic field eliminates the non-analyticity in the polarization bubble $`\mathrm{\Pi }^{}(\stackrel{}{q},\mathrm{\Omega }_n,)`$, such that only the second term in Eq. (4) contributes to the $`T`$ term in $`\mathrm{\Delta }\gamma (T,H)`$.
The extension of the above results to an arbitrary $`U(q)`$ is straightforward. For the susceptibility, the prefactor in Eq. (9), contains $`U(2k_F)`$ instead of $`U`$belitz ; millis ; andrey1 . For the specific heat coefficient, we have, instead of (13)
$`\mathrm{\Delta }\gamma (T,H)={\displaystyle \frac{3\zeta (3)m^4}{2\pi ^3k_F^2}}T\times [(U(0){\displaystyle \frac{1}{2}}U(2k_F))^2`$
$`+{\displaystyle \frac{U^2(2k_F)}{4}}(1+2{\displaystyle \frac{f_\gamma (x)3\zeta (3)}{3\zeta (3)}})].`$ (14)
The combinations $`U(0)1/2U(2k_F)`$ and $`1/2U(2k_F)`$ are charge and spin components of the scattering amplitude $`A(\pi )`$, respectively. At large $`x`$, $`f_\gamma (x)3\zeta (3)`$, and the last term in the r.h.s. of (14) vanishes. This obviously implies that at a large field, only the charge the longitudinal spin components of the scattering amplitude contribute to $`\mathrm{\Delta }\gamma (T)`$.
We next consider the $`3D`$ case. The polarization operators at small $`q`$ are given by AGD :
$`\mathrm{\Pi }^{}(\stackrel{}{q},\mathrm{\Omega }_n,0)={\displaystyle \frac{mk_F}{2\pi ^2}}{\displaystyle \frac{\mathrm{\Omega }_n}{v_Fq}}\mathrm{arctan}{\displaystyle \frac{(\mathrm{\Omega }_ni\delta \mu )}{v_Fq}}+\mathrm{}`$
$`\mathrm{\Pi }^{}(\stackrel{}{q},\mathrm{\Omega }_n,T)={\displaystyle \frac{mk_F}{2\pi ^2}}{\displaystyle \frac{\mathrm{\Omega }_n}{v_Fq}}\mathrm{arctan}{\displaystyle \frac{\mathrm{\Omega }_n}{v_Fq}}+\mathrm{}`$ (15)
As before, dots stand for analytic terms, expandable in powers of $`\mathrm{\Omega }_n^2`$ or $`q^2`$, and $`\delta \mu =\mu _{}\mu _{}=2\mu _BH`$. Differentiating the thermodynamic potential, Eq. (4), with respect to $`H`$, we obtain
$$\mathrm{\Delta }M=\frac{\mathrm{\Delta }\mathrm{\Phi }}{H}=\frac{\mu _B(mUk_F)^2}{4\pi ^5v_F^3}T\underset{n}{}\mathrm{\Omega }_{n}^{}{}_{}{}^{2}\mathrm{arctan}\frac{2\mu _BH}{|\mathrm{\Omega }_n|}.$$
Differentiating further with respect to $`H`$, we obtain
$`\mathrm{\Delta }\chi _s(T,H)={\displaystyle \frac{\mu _B(mUk_F)^2}{2\pi ^5v_F^3}}`$
$`\left[T{\displaystyle \underset{n=1}{\overset{M}{}}}\mathrm{\Omega }_n+4(\mu _BH)^2T{\displaystyle \underset{n=1}{\overset{M}{}}}{\displaystyle \frac{\mathrm{\Omega }_n}{\mathrm{\Omega }_n^2+4\mu _B^2H^2}}\right],`$ (16)
where $`ME_F/T`$ is the upper cutoff in the summation over frequency. The first term in the r.h.s of Eq. (16) is the susceptibility at zero field. By power counting, one might expect the spin susceptibility $`\chi _s(T)`$ in 3D to scale as $`T^2\mathrm{log}T`$. However, the Matsubara sum $`T_{n=1}^M\mathrm{\Omega }_n`$ only contains a T-independent term, of order $`E_F^2`$, and a term $`(1/6)\pi T^2`$. This last term is universal, but it is analytic in $`T`$. As a result, $`\mathrm{\Delta }\chi _s(T,0)T^2`$ is analytic and essentially irrelevant as the analytic in $`T`$ contributions to $`\chi _s(T)`$ are already present in the Lindhard function for free fermions. The absence of the non-analytic temperature correction to the spin susceptibility in $`3D`$ was first noticed in Ref. belitz , (see also pethick ). The second term in the r.h.s. of (16) is the extra contribution in a finite field. Evaluating the Matsubara sum we find that this contribution scales as $`H^2\mathrm{log}\{max(T,\mu _BH)\}`$. We see therefore that in a finite magnetic field, $`\chi _s(T)`$ does indeed become non-analytic. Casting $`\mathrm{\Delta }\chi _s(T,H)`$ into the scaling form, we obtain
$$\mathrm{\Delta }\chi _s(T,H)=\chi _0\left(\frac{mUk_F}{2\pi ^2}\right)^2\left(\frac{T}{E_F}\right)^2g\left(\frac{\mu _BH}{T}\right),$$
(17)
where $`\chi _0=\mu _B^2k_F^3/(2\pi ^2E_F)`$ is Pauli susceptibility, and to a logarithmic accuracy,
$$g(x)=x^2\mathrm{log}\left[\frac{E_F}{Tmax\{x,1\}}\right].$$
(18)
The $`H^2\mathrm{log}H`$ dependence of $`\chi (H)`$ was earlier reported by Misawa misawa . However, his prefactor is different from the one we obtained.
Differentiating the thermodynamic potential twice over $`T`$, we also obtain field dependence of the specific heat coefficient. The field dependence in 3D parallels the one for 2D systems. Namely, at zero field,
$`\mathrm{\Delta }\gamma (T,0)={\displaystyle \frac{3}{20}}{\displaystyle \frac{(mk_F)^2}{\pi ^2}}`$ (19)
$`\times \left[(U(0){\displaystyle \frac{1}{2}}U(2k_F))^2+{\displaystyle \frac{3}{4}}U^2(2k_F)\right]\left({\displaystyle \frac{T}{E_F}}\right)^2\mathrm{ln}{\displaystyle \frac{E_F}{T}}.`$
In a finite field, the charge part is not affected, while in the spin part, the logarithmic factor $`3\mathrm{log}\frac{E_F}{T}`$ is replaced by $`\mathrm{log}\frac{E_F}{T}+2\mathrm{log}\frac{E_F}{max\{T,\mu _BH\}}`$. As a result, at $`\mu _BHT`$, $`\mathrm{\Delta }\gamma (T,H)`$ still behaves as $`T^2\mathrm{log}T`$, but the prefactor gets smaller.
To summarize, in this paper we analyzed non-analytic terms in the magnetization, the spin susceptibility and the specific heat of 2D and 3D Fermi liquids, placed into an external magnetic field $`\mu _BHE_F`$. We obtained the non-analytic terms in the forms of scaling functions of $`\mu _BH/T`$. We found that at $`\mu _BHT`$, the spin susceptibility scales as $`|H|`$ in 2D and as $`H^2\mathrm{log}|H|`$ in 3D. The specific heat in a field preserves the same temperature dependence as in the absence of a field, but the prefactor changes between small and large $`\mu _BH/T`$.
We thank D. Belitz, P. Fulde, D. Maslov and T. Vojta for useful discussions. The research has been supported by NSF DMR 0240238 (A. Ch.), and the Visitors Program of the Max Planck Society (J. B.)
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# The Brightest Point X-Ray Sources in Elliptical Galaxies and the Mass Spectrum of Accreting Black Holes
## 1. INTRODUCTION
Chandra has revolutionized the study of point X-ray sources in the nearby Universe. The majority of these are interpreted to be X-ray binaries (XRBs; for a general review on Chandra performance see Weisskopf et al. 2003, for extragalactic X-ray binaries see, e.g., Kim & Fabbiano, 2004; Jordรกn et al., 2004). Elliptical galaxies out to the Virgo cluster have now been studied and have surprised us with the large number (typically $`100`$ per galaxy) of detectable point X-ray sources down to X-ray luminosities of about $`10^{37}`$ erg s<sup>-1</sup>. Two main population characteristics have attracted considerable attention so far: (i) the X-ray luminosity function (hereafter XLF) that may or may not exhibit a break at about $`45\times 10^{38}`$ erg s<sup>-1</sup> (for a recent update see Kim & Fabbiano, 2004); (ii) the high fraction of sources coincident with identified globular clusters (GCs) in ellipticals.
The shape of the XLF has been debated since the first observations of ellipticals were reported. Sarazin et al. (2000) identified a shape that required two power laws with a โbreakโ or a โkneeโ at $`\mathrm{\hspace{0.17em}3.2}\times \mathrm{\hspace{0.17em}10}^{38}`$ ergs s<sup>-1</sup>. Kim & Fabbiano (2003) argued that the break may result from biases affecting the detection threshold of the data. In the following few years longer exposures became possible and more and more ellipticals were added in the observed sample with low enough sensitivity (see, e.g. Gilfanov, 2004). The current situation is probably best summarized in Kim & Fabbiano (2004). They analyzed a large sample of elliptical galaxies with varying sizes of point source samples, and they concluded that: although XLFs of individual galaxies do not require a broken-power-law fit, the combined sample of sources from all the galaxies considered shows a statistically significant requirement for two power laws and a break at $`5\pm 1.6\times 10^{38}`$ ergs s<sup>-1</sup>. They found the best-fit slope of the lower end of the differential XLF to be $`\alpha _\mathrm{d}=1.8\pm 0.2`$ and the best-fit slope of the upper end to be $`\alpha _\mathrm{d}=2.8\pm 0.6`$. It is important to note that the break location is consistent with the Eddington luminosity for a $`1.9\pm 0.6`$ M neutron star (NS) accreting helium-rich material (for hydrogen rich donor this value is as large as $`3.2\pm 1M_{}`$). In what follows we consider the results of the Kim & Fabbiano (2004) study as representing our current observational understanding of the XLF in ellipticals. We address the question of the interpretation of this understanding and what it implies about the properties of the sources contributing to the observed XLFs.
Large fractions (20% โ 70%) of the point sources in ellipticals have been reported to be associated with globular clusters (see, e.g., Sarazin et al., 2003, and references therein). These high fractions have led to the suggestion that all point X-ray sources seen currently in ellipticals have been formed through stellar interactions and that sources that are not associated with GCs have originated in GCs and have either been ejected or the parent GCs have been destroyed by the galaxian tidal forces. As much interesting as this suggestion is, it raises the question: why would the field stellar population of ellipticals not lead to XRB formation as it has occurred in the Milky Way, for example? One could speculate that the field population is just too old and the XRBs that were formed at some point have completed their X-ray emitting life. However, such a speculation is inconsistent with the expectation that XRBs with low mass donors can live for several Gyrs as the donors lose mass and the binaries enter a transient phase. Such systems would become detectable as bright X-ray sources during the disk outbursts (Piro & Bildsten, 2002). Moreover, it has recently been pointed out that the high rate of source coincidence with GCs actually appears consistent with some of the sources having been formed in the field (Juett, 2005), although the result may sensitively depend on the definition of GC concentration in ellipticals. Most importantly for the present study, of the bright sources above the XLF break only a very small fraction is associated with GCs (e.g., only one source in the Virgo cluster sample as reported by Jordรกn et al., 2004).
In a recent Letter Bildsten & Deloye (2004, hereafter DB) have suggested an explanation that couples the two population characteristics: the XLF shape and the source coincidence with GCs. They identify the point X-ray sources as ultra-compact binaries (UCBs) that form predominantly in GCs. They are neutron stars accreting from low-mass He or C/O white dwarfs and they contribute to the ellipticals XLF early in their lifetime when they are still bright. Their association with GCs is important in replenishing the population through tidal interactions and allowing a significant number of sources in this bright phase, even though there is no ongoing star formation in ellipticals nor in GCs. DB estimate the rate of ultracompact binary formation to be consistent with the number sources observed and conclude that the model XLF slope is in agreement with the slope below the break as derived by Kim & Fabbiano (2004). The association of the XLF break with the NS Eddington limit for He-rich accretion is also consistent with this interpretation. However, the upper end of the XLF (at luminosities in excess of the break location) are not easy to interpret. Deloye & Bildsten suggest that some of the NS ultra-compact sources can reach super-Eddington luminosities. However luminosities in excess of $`10^{39}`$ erg s<sup>-1</sup> are very difficult to explain with NS accretors. Therefore the origin of the upper-end slope is not naturally connected to NS UCBs formed in GCs.
In this paper we address the question of the upper-end XLF slope and its origin. We consider the previously made suggestion (Sarazin et al., 2000) that the XLF above the break at $`5\times 10^{38}`$ erg s<sup>-1</sup> is populated by XRBs with black hole (BH) accretors. Given that GCs are not expected to harbor a significant number of BH-XRBs (see Kalogera et al., 2004, and references therein), we suggest that the vast majority of these BH-XRBs are part of the galactic-field stellar population in ellipticals. As we will show most of donors in these binaries are of low-enough mass that the XRBs are expected to be transient and therefore they populate the XLF only during disk outbursts when they typically emit at the Eddington luminosity for their BH mass. We further suggest that the slope of the upper XLF is a footprint of the BH mass spectrum in the BH XRBs under consideration. We present analytical derivations that demonstrate this link and we develop a method that allows us to infer the underlying BH mass spectrum consistent with the current upper-end XLF slope (Kim & Fabbiano, 2004). We also show that given the current observations it is possible to constrain the strength of magnetic braking acting in these XRBs, the type of BH donors, as well as the transient duty cycle to some extent. We also examine the quantitative robustness of our results against variations of some basic assumptions. This analysis is presented in ยง 2 and 3. We conclude with a discussion of our results and possible connections to population synthesis calculations (ยง 4).
## 2. Black Hole X-Ray Binaries in Ellipticals
We consider XRBs that could possibly populate the part of the observed XLF above the reported break at $`46\times \mathrm{\hspace{0.17em}10}^{38}`$ erg s<sup>-1</sup>, and therefore we focus on BH accretors (masses in excess of $`23`$ M). Given the current estimates for the ages of stellar populations in ellipticals (in their majority 8 to 12 Gyr, although some estimates are slightly shorter than 5 Gyr; see Ryden et al., 2001; Temi et al., 2005), we expect that donor masses are lower than $`11.5`$ M. Given these mass ratios and the properties of similar observed systems in the Milky Way (i.e., soft X-ray transients), these BH XRBs are expected to be transient X-ray sources (see McClintock & Remillard 2005 for a review of BH X-ray binaries in the Milky Way), where mass transfer is driven by the Roche-lobe filling donor. Given the above upper limit on the donor mass for ellipticals, we expect that there will be three different types of low-mass donors: (i) Main Sequence (MS) stars, (ii) Evolved or Red Giant Branch (RG) stars, and (iii) White Dwarf (WD) donors. Each of these sub-populations of BH XRBs will have different typical mass-transfer rates and binary property distributions, and therefore we examine them separately in our analysis that follows.
### 2.1. Transient X-Ray Sources
We adopt the current understanding for the origin of transient behavior in XRBs (for a recent review, see King, 2005). To identify transient systems in our modeling we consider the typical mass-transfer (MT) rate associated with each type of XRB donor ($`\dot{M}_i`$) and compare it to the critical MT rates for transient behavior ($`\dot{M}_{\mathrm{crit}}`$): if $`\dot{M}_i<\dot{M}_{\mathrm{crit}}`$, then the accretion disk is expected to be thermally unstable and the binary system is assumed to be a transient X-ray source. The value of this critical MT rate for the disk instability is not precisely known and its functional dependence on disk and binary properties are subject to uncertainties associated with our current theoretical understanding of the disk instability. However, both the qualitative concept of the existence of a critical rate for the instability to set in and its quantitative estimates by recent studies appear to be in good agreement with the behavior of Galactic X-ray transients. Therefore, we adopt the current understanding and quantitative estimates. More specifically, for hydrogen-rich donors, we adopt the $`\dot{M}_{\mathrm{crit}}`$ value derived by Dubus et al. (1999), and for helium or carbon-oxygen donors, we adopt the value derived by Menou et al. (2002). The effects of quantitative deviations from the adopted expressions are discussed in what follows.
To account for the contribution of transient sources in the XLF among any persistent sources, assumptions about the XRB luminosity at outburst and the transient duty cycle need to be made.
When a XRB is identified as transient, we assume that during the disk outburst the X-ray luminosity $`L_\mathrm{X}`$ is equal to the Eddington luminosity $`L_{\mathrm{Edd}}`$ associated to the BH accretor:
$$L_{\mathrm{Edd}}=\frac{4\pi cGM_{\mathrm{BH}}}{\kappa }=5\times 10^{37}\frac{m_{\mathrm{BH}}}{\kappa }\mathrm{ergs}\mathrm{s}^1,$$
(1)
where $`m_{\mathrm{BH}}`$ is the accretor mass in $`M_{}`$ and $`\kappa `$ is the opacity of the accreting material in cm<sup>2</sup> g<sup>-1</sup>. We adopt electron scattering opacities equal to 0.32 and 0.19 for hydrogen and helium or carbon-oxygen rich material, respectively.
We note that of the 15 confirmed transient BH XRBs in our Galaxy, 3 appear to reach possibly super-Eddington luminosities (McClintock & Remillard, 2005) at outburst (although distance estimate uncertainties cannot be ignored). Two of them, V4641 Sgr and 4U 1543-47, have early-type donors. Such donors are not present in elliptical galaxies with population ages of $`510`$ Gyr. The third one, GRS 1915+105, has a low-mass giant donor and an orbital period of 33 days. However its X-ray luminosity at outburst just barely exceeds its $`L_{\mathrm{Edd}}`$, by 40% only. Given distance uncertainties associated with such an estimate, we conclude that we can neglect the possibility of super-Eddington luminosities during outburst in our XLF modeling. On the other hand outburst peak luminosities cover a significant range at sub-Eddington values. During primary<sup>1</sup><sup>1</sup>1We use the term โprimaryโ to distinguish from โfollow-upโ outbursts that are occasionally observed very soon after primary ones with peak luminosities orders of magnitude below the Eddington limit (Remillard & McClintock 2005, private communication). Such small outbursts do not reflect an extremely short duty cycle and do not contribute to the high-end XLF of interest here. outbursts peak $`L_X`$ values can be lower than $`L_{\mathrm{Edd}}`$ by factors of a few (McClintock & Remillard 2005, private communication). As part of our analysis we examine the effect of such variations on the methodology and conclusions presented here (see ยง 3.1.1). From an observationally point of view it has been shown (Zezas et al. 2004 and Zezas 2005, private communication) that variability in X-ray fluxes (and hence luminosities) by factors of a few (typical among accreting sources and detected with Chandra observations at different epochs) do not alter the XLF slopes as measured for nearby galaxies within the current errors. Therefore observationally the reported XLF slopes appear to be robust. Consequently we can use them to learn about the underlying XRB population with considerable confidence.
At present there are no strong constraints on the duty cycles either from observations or from theoretical considerations. Among known Galactic X-ray transients, typical duty cycles of a few % is favored for hydrogen donors (Tanaka & Shibazaki, 1996). To our knowledge, there are no data on duty cycles for transients with a WD companion. In what follows we investigate how plausible duty cycle assumptions affect the upper-end XLF shape. In particular, we consider two specific cases: one of constant duty cycle equal to $`\eta =0.01`$; another of a variable (dependent on MT rates) duty cycle equal to
$$\eta =0.1\left(\frac{\dot{M}_\mathrm{d}}{\dot{M}_{\mathrm{crit}}}\right)^\delta $$
(2)
where $`\delta =1`$ is assumed. The first of these two cases corresponds to the standard assumption of a constant duty cycle often made in the literature. The second case is motivated primarily by our plan to examine how one example form of a MT-dependent duty cycle affects our analysis and results. Admittedly the specific choice of the dependence on $`\dot{M}_{\mathrm{crit}}`$ shown above is not solidly motivated, given all the uncertainties of the outburst mechanism. However it implies a correlation of the duty cycle with how strong a transient the system is: the further away from the critical MT rate, the smaller the duty cycle. We stress that in most of our analysis we adopt this form with $`\delta =1`$ as a plausible example and throughout the paper we contrast the results to those obtained with a constant duty cycle. Furthermore in ยง 3.1.1 we examine the sensitivity of our results for main-sequence donors on the choice of the duty-cycle dependence on MT extensively: we adopt $`\delta >1`$ in eq.(2) and we also introduce yet one other example of a MT-dependent duty cycle:
$$\eta =0.1\left(\frac{\dot{M}_\mathrm{d}}{\dot{M}_{\mathrm{EDD}}}\right)^\delta $$
(3)
for which we examine various values of $`\delta `$. Once again there is no solid theoretical motivation for this latter functional choice. However, it provides us with a better understanding of how sensitive our results are to the details of the possible duty-cycle dependence on MT properties.
### 2.2. Main Sequence Donors
Mass transfer in BH XRBs with hydrogen-rich, low-mass MS donors is expected to be driven by angular momentum losses due to magnetic braking (MB) and gravitational radiation (GR). In the case of conservative mass transfer (Verbunt, 1993) the angular momentum loss rate are connected to the MT rate as follows:
$$\frac{\dot{J}_{\mathrm{gr}}}{J_{\mathrm{orb}}}+\frac{\dot{J}_{\mathrm{mb}}}{J_{\mathrm{orb}}}=\frac{\dot{M}_\mathrm{d}}{M_\mathrm{d}}\left(\frac{5}{6}+\frac{n}{2}\frac{M_\mathrm{d}}{M_{\mathrm{BH}}}\right),$$
(4)
where $`n`$ is the radius-mass exponent for the donor. For a low-mass MS star:
$$r_\mathrm{d}m_\mathrm{d},$$
(5)
where $`r_\mathrm{d}=R_\mathrm{d}/R_{}`$ and $`m_\mathrm{d}=M_\mathrm{d}/M_{}`$ are the donor stellar radius and mass in solar units. Therefore $`ndlnR_\mathrm{d}/dlnM_\mathrm{d}`$ is equal to $`1`$ for MS donors.
According to general relativity the rate of angular momentum loss due to GR is given by:
$`{\displaystyle \frac{\dot{J}_{\mathrm{gr}}}{J_{\mathrm{orb}}}}={\displaystyle \frac{32G^3}{5c^5}}{\displaystyle \frac{M_{\mathrm{BH}}M_\mathrm{d}(M_{\mathrm{BH}}+M_\mathrm{d})}{A^4}}`$
$`=2.6\times 10^{17}{\displaystyle \frac{m_{\mathrm{BH}}m_\mathrm{d}(m_{\mathrm{BH}}+m_\mathrm{d})}{a^4}}\mathrm{s}^1,`$ (6)
where $`a`$ is the orbital semi-major axis in units of solar radius. For a mass ratio $`qM_d/M_{\mathrm{BH}}<0.8`$ (Paczyลski, 1971) and using the mass-radius relation eq. (5):
$$a=\frac{1}{0.46}m_\mathrm{d}^{2/3}(m_{\mathrm{BH}}+m_\mathrm{d})^{1/3}.$$
(7)
We consider two derivations of the angular momentum loss rate due to magnetic braking: (i) the Skumanich law based on the empirical relation for slowly rotating stars adopted from (Rappaport et al., 1983) (RVJ), and (ii) the revised law based on X-ray observations of faster rotating dwarfs adopted from (Ivanova & Taam, 2003) (IT):
$`\dot{J}_{\mathrm{mb}}^{\mathrm{RVJ}}`$ $`=3.8\times 10^{30}M_\mathrm{d}R_{}^4(R_\mathrm{d}/R_{})^2\mathrm{\Omega }^3\mathrm{dyn}\mathrm{cm}`$ (8)
$`\dot{J}_{\mathrm{mb}}^{\mathrm{IT}}`$ $`=6\times 10^{30}(R_\mathrm{d}/R_{})^4\left({\displaystyle \frac{\mathrm{\Omega }^{1.3}\mathrm{\Omega }_\mathrm{x}^{1.7}}{\mathrm{\Omega }_{}^3}}\right)\mathrm{dyn}\mathrm{cm},`$ (9)
where $`\mathrm{\Omega }`$ \[s<sup>-1</sup>\] is the stellar angular velocity which is equal to the binary orbital velocity assuming the star is in full synchronization with the binary orbit, $`\mathrm{\Omega }_{}=5\times 10^6`$ s<sup>-1</sup> is the Sunโs angular velocity, and $`\mathrm{\Omega }_\mathrm{x}=10\mathrm{\Omega }_{}`$. Using Keplerโs law the above are re-written as:
$$\frac{\dot{J}_{\mathrm{mb}}^{\mathrm{RVJ}}}{J_{\mathrm{orb}}}=7.2\times 10^{15}\frac{m_\mathrm{d}^2(m_{\mathrm{BH}}+m_\mathrm{d})^2}{m_{\mathrm{BH}}a^5}\mathrm{s}^1,$$
(10)
$$\frac{\dot{J}_{\mathrm{mb}}^{\mathrm{IT}}}{J_{\mathrm{orb}}}=2.7\times 10^{17}\frac{m_\mathrm{d}^3(m_{\mathrm{BH}}+m_\mathrm{d})^{1.15}}{m_{\mathrm{BH}}a^{2.45}}\mathrm{s}^1.$$
(11)
In XRBs with BH masses $`3M_{}`$ and MS donor masses $`1.0M_{}`$, it is $`J_{\mathrm{gr}}J_{\mathrm{mb}}^{\mathrm{RVJ}}`$ and $`J_{\mathrm{gr}}J_{\mathrm{mb}}^{\mathrm{IT}}`$. The lifetime of the MS-BH XRBs is much longer in the latter case.
As mentioned earlier in this study we adopt the derivation of the critical MT rate below which the accretion disk becomes unstable for irradiated disks presented by Dubus et al. 1999:
$$\dot{M}_{\mathrm{crit}}=1.5\times 10^{15}m_{\mathrm{BH}}^{0.4}\left(\frac{R_{\mathrm{disk}}}{10^{10}\mathrm{cm}}\right)^{2.1}\mathrm{g}\mathrm{s}^1.$$
(12)
Here $`R_{\mathrm{disk}}`$ is the radius of the accretion disk. We note that the exact value of this critical rate is subject to uncertainties associated with our limited understanding of the disk instability, but we adopt the above expression as indicative of the process and we continue with our analysis.
From eq.(4), (10) and (12) it can be shown numerically, that for the RVJ MB law and for low-mass donors, there is a BH mass $`M_{\mathrm{PT}}`$ of $`5M_{}`$ that separates BH-MS systems into persistent ($`M_{\mathrm{BH}}<M_{\mathrm{PT}}`$), and transient ($`M_{\mathrm{BH}}>M_{\mathrm{PT}}`$). From more detailed binary evolutionary calculations using the stellar evolution and MT code described in (Ivanova & Taam, 2004), we find that this boundary is about 10 $`M_{}`$ (see Fig. 1 for details). For BHs less massive than this critical mass, the XRBs are persistent as long as the donor masses are higher than about 0.3 $`M_{}`$. In these persistent sources the MT rates driven by the RVJ type of MB turn out to be $`0.010.25`$ of the black holesโs Eddington rate. As a result, the persistent X-ray luminosity for these systems is $`10^{38}`$ erg s<sup>-1</sup>, i.e., below the bright $`L_X`$ range we consider here. Therefore we conclude that the persistent BH-MS binaries driven by the RVJ type of MB cannot contribute significantly to the upper-end XLFs of ellipticals.
Next we examine whether the transient phases associated with BH-MS binaries and the RVJ MB law are important when they reach X-ray luminosities comparable to the Eddington limit. For $`M_{\mathrm{BH}}>10`$$`M_{}`$, the outburst luminosity is expected to be in excess of $`1.5\times 10^{39}`$ erg s<sup>-1</sup>. However, this lower limit is comparable to the highest luminosity seen currently in XLFs of ellipticals (Kim & Fabbiano, 2004), and therefore these systems cannot contribute significantly to the observed XLFs. The last possibility is outbursts from transient BH-MS with $`M_{\mathrm{BH}}<10`$$`M_{}`$ and donors less massive than $`0.3`$ M. Such low mass donors are out of thermal equilibrium and significantly expanded ($`3\times `$) compared to an undisturbed MS star of the same mass. In the case of the MT dependent duty-cycle $`\eta `$ is about a few %. We note, however, that applicability of MB for these stars is very questionable, as it is generally accepted that MB does not operate in fully convective stars, which are found to be less massive than $`0.35`$ M for undisturbed stars. In principle, however, โerodedโ, out-of-thermal-equilibrium low-mass MS donors like the ones in XRBs do not necessarily become fully convective at the same critical mass as stars with no prior MT evolution. For this reason we have used detailed MT calculations (Fig. 1) with a stellar-evolution code to examine this behavior further. We find that in the donor-mass range $`0.150.3M_{}`$ the radiative core is extremely small, even for these โerodedโ stars, and therefore applying angular momentum loss due to MB is not reasonable. Instead MB activity and hence mass transfer is expected to be interrupted until eventually GR drives Roche-lobe overflow much later.
Based on the above line of arguments we conclude that BH-MS binaries evolving according to the RVJ MB law are not expected to contribute significantly to the high-end XLFs of ellipticals.
In the case of the IT MB prescription, BH-MS systems are transient for all BHs masses $`M_{\mathrm{BH}}>3M_{}`$ and for all low-mass MS donors. The reason is that the IT MB is weaker; consequently the donors are mildly out of thermal equilibrium and the mass transfer rates are lower compared to the RVJ MB case. Using again detailed MT evolutionary simulations we find that $`\dot{M}/\dot{M}_{\mathrm{crit}}0.25\pm 0.15`$ (see Fig. 2). Therefore, in the case of the MT dependent duty-cycle $`\eta `$ is again about a few %. In both cases (RVJ MB and IT MB), the value of the duty cycle is consistent with observations for BHs of different masses, though the transiency occurs at very different donor masses.
We conclude that, regardless of the specific MB law, it is rather unlikely that persistent sources with a BH accretor and a hydrogen-rich, low-mass MS donor populate at any significant fraction the upper-end XLF of ellipticals; only transient BH-MS sources driven by the IT MB law can populate this X-ray luminosity range.
### 2.3. Red Giant Donors
For orbital periods more than about a day, MT occurs when the low-mass donor is a subgiant or a giant. The driving force is the nuclear expansion of the donor, and a simple analytic prescription for the MT is (Webbink, Rappaport, Savonije 1983; Ritter 1999; see also King 2005):
$$\dot{M}\mathrm{rg}=3.4\times 10^{15}a^{1.4}\frac{m_\mathrm{d}^{1.47}}{(m_{\mathrm{BH}}+m_\mathrm{d})^{0.465}}\mathrm{g}\mathrm{s}^1.$$
(13)
It has been shown (by King et al., 1997; King, 2000) that such wider XRBs are transient, regardless of the BH mass (original derivations were based on a somewhat different expression for the critical MT rate for transient behavior, but still quite similar to eq. .
### 2.4. White Dwarf Donors
A typical WD mass-radius relation is (see, e.g., Rappaport et al., 1987):
$$r_\mathrm{d}=0.0128m_\mathrm{d}^{1/3}$$
(14)
Using an approximation for the Roche Lobe radius (from Paczyลski, 1971) and assuming that the mass of WD is much smaller than a BH mass, we can show that
$$a=0.0278m_\mathrm{d}^{1/3}\left(\frac{m_\mathrm{d}}{m_{\mathrm{BH}}}\right)^{1/3}$$
(15)
We consider again conservative mass transfer (4) but without any MB losses, adopting $`n=1/3`$ and assuming that $`m_\mathrm{d}m_{\mathrm{BH}}`$. Then
$$\dot{M}_\mathrm{d}\frac{32G^3}{3.3c^5}\frac{M_{\mathrm{BH}}^2M_\mathrm{d}^2}{A^4}=7.9\times 10^{16}\frac{m_{\mathrm{BH}}^2m_\mathrm{d}^2}{a^4}\mathrm{g}\mathrm{s}^1$$
(16)
We substitute here eq. (15) and then have
$$\dot{m}_\mathrm{d}=2\times 10^3m_\mathrm{d}^{4\frac{2}{3}}m_{\mathrm{BH}}^{\frac{2}{3}}M_{}\mathrm{yr}^1$$
(17)
In what follows we adopt the critical MT rate for He-rich donors from Menou et al. (2002), but we note that the expression is subject to quantitative uncertainties associated with the current understanding of the disk instability:
$$\dot{M}_{\mathrm{crit}}=5.9\times 10^{16}m_{\mathrm{BH}}^{0.87}\left(\frac{R_{\mathrm{disk}}}{10^{10}}\right)^{2.62}\mathrm{g}\mathrm{s}^1$$
(18)
For a large mass ratio $`q_{\mathrm{BH}}=m_{\mathrm{BH}}/m_\mathrm{d}`$ the Roche lobe of the accretor is $`r_{\mathrm{RL}}0.7a`$ and
$$r_{\mathrm{disk}}2/3r_{\mathrm{RL}}=0.013m_\mathrm{d}^{2/3}m_{\mathrm{BH}}^{1/3}$$
(19)
$$\dot{m}_{\mathrm{crit}}=1.7\times 10^{12}m_\mathrm{d}^{1.74}M_{}\mathrm{yr}^1$$
(20)
BH-WD binaries will be transient if
$$\frac{\dot{m}_\mathrm{d}}{\dot{m}_{\mathrm{crit}}}=1.2\times 10^9m_\mathrm{d}^{6.4}m_{\mathrm{BH}}^{\frac{2}{3}}1$$
(21)
Therefore the maximum donor mass that leads to transient behavior in BH-WD binaries is:
$$m_{\mathrm{tr}}=0.038m_{\mathrm{BH}}^{0.1}$$
(22)
The time interval in Gyr needed for the WD donor mass to evolve from $`m_d^{11/3}(T_1)`$ to $`m_d^{11/3}(T_2)`$ is (using eq. 17):
$`T_2T_1`$ $`=`$ $`{\displaystyle \frac{3}{22}}\times 10^6m_{\mathrm{BH}}^{2/3}\times `$ (23)
$`\left(m_d(T_2)^{11/3}m_d(T_1)^{11/3}\right)`$
Here we assume that the mass of the BH is constant, since $`m_{\mathrm{BH}}m_d`$. Consequently and using eq. (22) we can find that the time a BH-WD system spends in the persistent state is $`t_{\mathrm{pers}}\mathrm{\hspace{0.17em}20}\times \mathrm{\hspace{0.17em}10}^6m_{\mathrm{BH}}^{0.3}`$ yr, i.e., it very weakly depends on the accretor mass (the dependence on the initial donor mass is negligible, below that 1%). We note that through this persistent phase the MT rate will be comparable or higher to the Eddington limit only for a very short time, $`t_{\mathrm{Edd}}\mathrm{\hspace{0.17em}2}\times \mathrm{\hspace{0.17em}10}^6m_{\mathrm{BH}}^{0.9}`$ yr <sup>2</sup><sup>2</sup>2Although MT is non-conservative during the super-Eddington accretion, and eq. (4) formally should not be applied, this result is well consistent with the detailed calculations that take into account non-conservative MT.
The evolution of BH-WD systems with C/O WD companions is rather similar. The critical MT rate (using Menou et al. 2002) is
$$\dot{m}_{\mathrm{crit}}=9.4\times 10^{13}m_\mathrm{d}^{1.47}M_{}\mathrm{yr}^1$$
(24)
and
$$m_{\mathrm{tr}}=0.03m_{\mathrm{BH}}^{0.1}.$$
(25)
The time that a BH-WD system with a C/O rich donor spends in the persistent state is also not very long, $`t_{\mathrm{pers}}\mathrm{\hspace{0.17em}50}\times \mathrm{\hspace{0.17em}10}^6m_{\mathrm{BH}}^{0.3}`$ yr.
We conclude that BH-WD binaries that contribute to the current upper-end XLFs of ellipticals are expected to be transient sources<sup>3</sup><sup>3</sup>3This would not be true if BH-WD binaries continuously formed, but this is not possible in the galactic field of ellipticals and is not even expected in globular clusters, since BHs tend to dynamically separate from the rest of the cluster and eject one another (Kulkarni et al., 1993; Sigurdsson & Hernquist, 1993; Watters et al., 2000).
## 3. Mass Spectrum Weighting Factor and Transient Duty Cycle
In the previous section we have shown that the upper-end XLF of ellipticals is dominated by transient BH XRBs possibly with a variety of donors: MS (for the case of the IT MB prescription) and RG stars, and WD donors with masses lower than $`0.035`$ M. All these systems contribute to the XLF only when in outburst, when their $`L_XL_{\mathrm{Edd}}M_{\mathrm{BH}}`$. Consequently the slope of the upper-end XLF can serve as a footprint of the BH mass distribution of accretors in the contributing BH XRBs. These contributing BH XRBs are just a sub-set (those in outburst) of the true population of BH XRBs in ellipticals determined by the duty cycle of BH transients binaries. For the general case of a transient duty cycle that is dependent on the BH accretor mass (and possibly other quantities), the differential XLF $`n(L)_{\mathrm{obs}}`$ and the underlying BH mass distribution in XRBs $`n(m)_{\mathrm{BH}}`$ are connected by:
$$n(L_X)_{\mathrm{obs}}=n(m_{\mathrm{BH}})\times W(m_{\mathrm{BH}}),$$
(26)
where $`W(m_{\mathrm{BH}})`$ is a weighting factor related to the dependence of the transient duty cycle on $`m_{\mathrm{BH}}`$.
The observed slope of the differential upper-end XLF is $`\alpha _\mathrm{d}=2.8\pm 0.6`$: $`n(L_X)_{\mathrm{obs}}L_X^{\alpha _\mathrm{d}}`$ (the slope of the cumulative upper-end XLF reported by Kim & Fabbiano 2004 is $`\alpha _\mathrm{c}=1.8\pm 0.6`$). Assuming that $`n(m_{\mathrm{BH}})m_{\mathrm{BH}}^\beta `$ and $`W(m_{\mathrm{BH}})m_{\mathrm{BH}}^\gamma `$, the slope characterizing the underlying BH mass distribution in XRBs is:
$$\beta =\alpha _\mathrm{d}\gamma .$$
(27)
For the standard assumption of a constant duty cycle, $`\beta =\alpha _\mathrm{d}=2.8\pm 0.6`$. In the following subsections we derive $`W(m_{\mathrm{BH}})`$ and $`\gamma `$ for all three types of donors in one example case of a duty cycle dependent on the binary and MT properties (see, e.g., eq.2).
### 3.1. Main Sequence Donors
In what follows we estimate the typical duty cycle for BH-MS transients averaged over the possible distribution of donor masses. In the case of the IT MB perscription, the angular momentum loss rate due to GR is comparable or even more important than MB for all BH masses above 3 $`M_{}`$ and donor masses $`1.01.2M_{}`$. The MT timescale during the transient phase is longer than the donorโs thermal timescale when the donor is $`0.25`$ M. For this range the donor is in thermal equilibrium and the approximation for the mass-radius dependence eq. (5) can be used. Using eq. (4), (6) and (12), and the fitting formula for the Roche lobe radius from Eggleton (1983), we find:
$$\frac{\dot{m}}{\dot{m}_{\mathrm{crit}}}0.054\frac{m_{\mathrm{BH}}^{0.4}(q_{\mathrm{BH}}^{2/3}\mathrm{log}(1+q_{\mathrm{BH}}^{1/3})+0.6)^{2.1}}{m_\mathrm{d}^2(1+q_{\mathrm{BH}})(4/31/q_{\mathrm{BH}})},$$
(28)
where $`q_{\mathrm{BH}}=m_{\mathrm{BH}}/m_\mathrm{d}`$ is the mass ratio. For $`q_{\mathrm{BH}}1`$ we have
$$\frac{\dot{m}}{\dot{m}_{\mathrm{crit}}}0.014m_{\mathrm{BH}}^{0.4}m_\mathrm{d}^2$$
(29)
At donor masses smaller than $`0.25`$ M, the donor is out of the thermal equlibirum and its radius is about twice bigger than predicted by eq. (5). We then find:
$$\frac{\dot{m}}{\dot{m}_{\mathrm{crit}}}0.0004m_{\mathrm{BH}}^{0.4}m_\mathrm{d}^2$$
(30)
We note that the split into the two expressions above is a rough, but useful approximation.
As discussed in ยง 2.2, for IT MB, a BH-MS system is transient throughout the MT phase. The MT rates are well below the Eddington limit for the BH mass and therefore we assumed that MT is fully conservative; i.e., $`M_\mathrm{d}+M_{\mathrm{BH}}=M_{\mathrm{tot}}`$ is constant with time. In what follows we assume a flat current mass distribution for donors ($`N/m_\mathrm{d}=const`$). We are guided in this choice by results from binary population synthesis calculations (with the StarTrack code; Belczynski et al. 2002 and 2005; Belczynski 2005, private communication). We integrate eq. 29 for $`m_\mathrm{d}`$ from $`\mathrm{\hspace{0.17em}0.25}\mathrm{to}\mathrm{\hspace{0.17em}1}`$ M and using eq. (2), we find that at present the probability that a system with a BH accretor of $`m_{\mathrm{BH}}`$ is in outburst and therefore contributes to the upper-end XLF is:
$$W(m_{\mathrm{BH}})=\frac{_{0.25}^{m_{\mathrm{TO}}}\eta \frac{N}{m_\mathrm{d}}๐m_\mathrm{d}}{_{0.25}^{m_{\mathrm{TO}}}\frac{N}{m_\mathrm{d}}๐m_\mathrm{d}}0.05m_{\mathrm{BH}}^{0.4},$$
(31)
where $`m_{\mathrm{TO}}`$ is the turn-off MS mass for the elliptical galaxy in solar units. We note that this result is valid only for large mass ratios $`q_{\mathrm{BH}}>>1`$. The contribution of BH-MS system when donors have masses $`0.25M_{}`$ (systems where the donor is out of the thermal equilibrium) is less significant.
It is also important to note here that for MS donors the factor $`W`$ does not appear to depend on time (i.e., the age of the elliptical galaxy). Such a time dependence would enter in relation to the value of the maximum donor mass (turn-off mass for the host galaxy). However we find that $`W(m_{\mathrm{BH}})`$ is a very weak function of $`m_{\mathrm{TO}}`$, and therefore it is not sensitive to the elliptical age.
#### 3.1.1 Monte Carlo Simulations
In principle, prolonged mass accretion onto the BHs can affect their mass spectrum. Since this effect cannot be included analytically, we have examined it quantitatively using simple Monte Carlo simulations. We set up the simulations assuming a flat BH-MS birth (MT onset) rate and a flat mass distribution for donors at the onset of the MT phase, without any restrictions on the mass ratio $`q_{\mathrm{BH}}`$. Donor masses at birth were varied in from 0.1 $`M_{}`$ to the current $`m_{\mathrm{TO}}`$ at an elliptical age of $`10`$ Gyr, assumed as a standard value. Each BH-MS binary was evolved to the current elliptical age using equations shown in the ยง 2.2. For the MT evolution we took into account both IT MB and GR, and for the Roche lobe radius we adopted the approximation by Eggleton (1983). A binary is removed from the MT population if the donor mass falls below $`0.05`$ M. If the MT timescale is longer than the thermal timescale of the donor, the donor radius evolution is simply proportional to the mass lost due to MT. On the other hand, if the MT timescale is shorter than the donorโs thermal timescale, the donor is out of the thermal equilibrium. In this case we modify the evolution of the donor radius using a prescription that is in acceptable agreement with our results from detailed MT calculations with the stellar evolution: $`\delta r\delta m\sqrt{\dot{m}_{\mathrm{TH}}/\dot{m}}`$, where $`\dot{m}_{\mathrm{TH}}=m_\mathrm{d}/t_{\mathrm{TH}}`$ is the MT rate driven on the donorโs thermal timescale $`t_{\mathrm{TH}}`$. Evolution of the transient systems follows the adopted duty cycle (eq. with $`\delta =1`$).
Based on the results of our Monte Carlo simulations we find that: (i) due to accretion the BH mass spectrum slope increases by about $`0.2`$, i.e., $`\beta =\beta _0+0.2`$, where $`\beta _0`$ is the BH mass slope at MT onset; (ii) the slope of the BH mass spectrum at the beginning of mass transfer best reproduces the observations with $`\beta _0=2.3\pm 0.6`$ (see Fig. 3). We also find that the relation between $`\beta `$ and $`\beta _0`$ is not sensitive to the age of the elliptical (as long as it is in the range of a few to several Gyr).
Next we examine how the results in this section are affected by plausible variations of a number of possibly oversimplifying assumptions made so far. Although it is possible to re-derive the analytical expressions with different assumptions, it requires repeating essentially the same analysis multiple times, something inappropriate for presentation here. Instead it is more efficient to examine these effects using the Monte Carlo simulations for XLF slopes.
For the tests that follow we adopt $`\beta =2.3`$ and examine how the cumulative XLF slopes $`\alpha _c`$ are affected.
We first examine the effect of random variations by a factor of 2 of outburst peak X-ray luminosities of individual transient systems. The results are shown in Fig. 4 (top panel): the solid line is our standard model where peak $`L_X`$ are set equal to $`L_{\mathrm{Edd}}`$ (same as the solid curve in Fig. 3 top panel) and the dash-dotted line is the result we obtain when random $`L_X`$ variations are introduced. They are essentially indistinguishable, and therefore the derivation of a BH mass spectrum slope from the observed XLF slope appears highly robust. This is consistent with the findings of observational studies of such variations (Zezas et al. 2004).
Next we examine the effect of $`\delta `$ values in the case of the MT-dependent duty cycle in eq. 2 other than 0 or 1. A set of curves for $`\delta =`$0.5, 1, and 2 are shown in Fig. 4 (top panel). Once again we conclude that the variations are essentially negligible to such quantitative changes. Instead the slope behavior seems to be dominated by the qualitative character of $`\eta `$ in eq. 2: the stronger the transient character, the smaller the duty cycle.
Last we examine the effect of changing the functional dependence of the duty cycle normalizing it to the Eddington MT rate instead of the critical rate for transient behavior (eq. ). The results are shown in Fig. 4 (bottom panel) for $`\delta =`$0.5, 1, and 2 and they are compared to the our standard case of $`\eta `$ (from eq. with $`\delta =1`$). This is the only case where significant variations are evident. More specifically, the resultant XLFs are steeper (absolute $`\alpha _c`$ values are higher) and they are more dependent on the choice of $`\delta `$. Given that the results are shown for a fixed $`\beta `$ value, it means that, for a specific observed XLF slope, this different form of the MT-dependent duty cycle would lead to the derivation of a flatter BH mass spectrum compared to that derived using eq. (3). We conclude that better understanding of the transient duty cycle and its dependence on binary or MT properties is important for obtaining reliable quantitative results in the future.
### 3.2. Red Giant Donors
In principle we should repeat the above estimate for BH-RG transients. However, from eq. (12), (13) and (19) we obtain:
$$\frac{\dot{m}_{\mathrm{RG}}}{\dot{m}_{\mathrm{crit}}}0.2a^{0.7}m_\mathrm{d}^{1.5},$$
(32)
and it is evident that, in this case, the MT-dependent duty cycle $`\eta `$ (see eq. 2) is independent of the BH mass. Therefore, for RG donors, $`\gamma =0`$. Instead it is significantly dependent on the RG donor mass. Since in ellipticals the typical mass for RG donors is about the same as the mass of the turn-off of MS stars, the duty cycle for RG donors depends on the turn-off mass, and hence on the age T \[Gyr\] of the elliptical. For a solar metallicity and stars of $`1`$ M$`{}_{}{}^{}M1.5`$ M, the approximate dependence of the turn-off mass with the age is $`m_\mathrm{d}m_{\mathrm{TO}}2T^{1/3}`$ (here we used the simplified evolutionary code from Hurley et al., 2000). Using eq. (2), and integrating eq. (32) over binary separations (similar to our integrations for MS donors; see eq. 31), we find:
$$W(T)0.03T^{0.5},$$
(33)
Here we assume a distribution of orbital separations for BH-RG binaries before MT starts that is flat in the logarithm. The reason for this choice is that this is appropriate for the distribution of zero-age binaries and the shape is actually preserved through wind mass loss, common-envelope evolution, and asymmetric explosions with small kicks appropriate for black holes (Kalogera & Webbink 1998).
It is interesting to note that for RG donors $`W`$ is dependent on the galaxy age whereas for MS donors the dependence on the BH mass dominates.
### 3.3. White Dwarf Donors
During the transient stage, the WD mass can be written as a function of time $`T`$ in Gyr (using eq. 23):
$`m_\mathrm{d}(T)=0.0134T^{3/11}m_{\mathrm{BH}}^{2/11}`$ (34)
We can then calculate the probability for a BH-WD system to contribute to the upper-end XLF at an elliptical age $`T`$. We assume that (i) all accreting BH-WD systems were formed within a short interval of elliptical ages $`T_{\mathrm{start}}`$ to $`T_{\mathrm{fin}}`$ (in Gyrs) several Gyrs ago when star formation was still occurring in the elliptical; and (ii) $`TT_{\mathrm{fin}}>t_{\mathrm{pers}}`$,i.e., the binary is a transient at time $`T`$. The latter assumption is well justified given the short duration of the persistent phase (see ยง 2.4). We further adopt a constant BH-WD formation rate between $`T_{\mathrm{start}}`$ and $`T_{\mathrm{fin}}`$, i.e., $`\frac{N}{t}=const`$. The probability then is expressed by the duty-cycle weighting factor at $`T`$ for a given BH mass:
$`W(T;m_{\mathrm{BH}})`$ $`=`$ $`{\displaystyle \frac{_{m_{\mathrm{d1}}}^{m_{\mathrm{d2}}}\eta \frac{N}{m_\mathrm{d}}๐m_\mathrm{d}}{_{m_{\mathrm{d1}}}^{m_{\mathrm{d2}}}\frac{N}{m_\mathrm{d}}๐m_\mathrm{d}}}`$ (35)
$`=`$ $`0.1{\displaystyle \frac{_{m_{\mathrm{d1}}}^{m_{\mathrm{d2}}}\frac{\dot{m}_\mathrm{d}}{\dot{m}_{\mathrm{crit}}}\frac{N}{t}\frac{t}{m_\mathrm{d}}๐m_\mathrm{d}}{_{m_{\mathrm{d1}}}^{m_{\mathrm{d2}}}\frac{N}{t}\frac{t}{m_\mathrm{d}}๐m_\mathrm{d}}}`$
$`=`$ $`0.1{\displaystyle \frac{_{m_{\mathrm{d1}}}^{m_{\mathrm{d2}}}\dot{m}_{\mathrm{crit}}^1๐m_\mathrm{d}}{_{m_{\mathrm{d1}}}^{m_{\mathrm{d2}}}\dot{m}_\mathrm{d}^1๐m_\mathrm{d}}}`$
Here $`m_{\mathrm{d1}}=m_\mathrm{d}(TT_{\mathrm{start}};m_{\mathrm{BH}})=m_\mathrm{d}(t_1;m_{\mathrm{BH}})`$ and $`m_{\mathrm{d2}}=m_\mathrm{d}(TT_{\mathrm{fin}};m_{\mathrm{BH}})=m_\mathrm{d}(t_2,m_{\mathrm{BH}})`$; $`m_{\mathrm{d2}}>m_{\mathrm{d1}}`$. Then, using eq. (17), (20) and (34), we obtain:
$`W(T;m_{\mathrm{BH}})`$ $`=`$ $`1.6\times 10^4m_{\mathrm{BH}}^{0.5}t_1^{7/4}`$ (36)
$`\times {\displaystyle \frac{1(t_2/t_1)^{3/4}}{1(t_2/t_1)}}`$
Therefore for WD donors $`\gamma =0.5`$ and $`W`$ depends on both the BH mass and the galaxy age.
## 4. Accreting Black Hole Mass Spectrum
So far we have derived the dependence of the duty-cycle weighting factor $`W`$ (eq.26) on the accreting BH mass and the age of the host elliptical galaxy, for the different types of BH donors. In order to make progress and develop a method for deriving constraints on the slope $`\beta `$ of the underlying accreting BH mass spectrum we need to examine which of the possible donor populations dominate the observed upper-end XLF under what conditions. The answer to this question requires large-scale population synthesis models that are not part of the scope of this paper. However, as is shown below, we can use a number of different arguments and pieces of evidence to derive tentative constraints. The primary purpose of our analysis is not to derive an unambiguous constraint at present, but instead to develop a methodology for how to derive the most reliable constraints, given the current uncertainties associated with our current understanding of these X-ray binaries.
For RG donors we find that in the case of a MT-dependent duty cycle $`\eta `$ is about an order of magnitude smaller than for MS donors: by comparing eq. (31) and eq. (33), using an elliptical-galaxy age of at least a few Gyr (more typical is 12 Gyr) and a BH mass of at least $`3`$ M the difference with the $`\eta `$ value for MS donors is a factor of 15. We also note that among known BH X-ray transients in our Galaxy, the ratio of transient BH-RG systems to transient BH-MS systems is about 1:2 (see Table 4.1 in McClintock & Remillard, 2005). Consequently we conclude that BH-RG transients cannot be important contributors to the upper-end XLF of elliptical galaxies for the example case of the MT-dependent duty cycle as defined by $`\eta `$. Only in the case of the constant, MT-independent (and therefore donor and BH-mass independent) duty cycle we expect BH-RG transients to be a significant population of the observed upper-end XLF.
Let us consider the case when the number of transient BH-RG systems exceeds the number of transient BH-MS system in a way that the contributions of the two populations become comparable at some age of the elliptical. We also assume that $`\beta _0`$ should be the same for both populations. In this case, the resultant combined XLF will be flatter than the XLF provided by only BH-MS contributors. Secondly, as the contribution of BH-RG system decreases with elliptical age (see eq. 33), the XLF becomes steeper, evolving towards the slope characteristic for BH-MS binaries. It is possible that this is the kind of behavior that we observe in XLFs of ellipticals (see Fig. 3): we note that younger ellipticals appear to have flatter XLFs, although uncertainties are significant.
For WD donors we find $`W(m)t_1^{7/4}`$ (see eq. 36), implying that the probability of each BH-WD transient contributing to the observed XLF decreases with the age of the elliptical galaxy (similar to the case of RG donors, but unlike the case of MS donors). Let us consider an elliptical where the formation of BH-WD systems has ended at least a few Gyr ago. We also consider that BH-MS systems are transient and have $`W(M)`$ according eq. (31). In order for BH-WD binaries to contribute significantly to the observed XLF they must form at a rate such that more than $`8,000`$ BH-WDs form for each BH-MS. This ratio is calculated adopting a value for the age of the elliptical of $`3.5`$ Gyr (among the lowest reported in the literature) and for a choice of BH masses in binaries with WD and MS donors, so that their Eddington X-ray luminosities are comparable, and therefore they contribute to the same X-ray luminosity bin (3 M for WD and 5 M for MS donors). For more typical, older ellipticals with ages closer to 10 Gyr the required ratio becomes even higher than $`8,000`$. According to binary population synthesis models for the Milky Way published so far, the number of formed BH-WD LMXBs exceeds the number of BH-MS LMXBs by at most a factor of 100 (Hurley et al., 2002). Furthermore, the lifetime of BH-MS binaries is of order 1 Gyr (or a few Gyr; see also Fig. 1), whereas the lifetime BH-WD binaries is longer, but cannot exceed the age of the elliptical galaxy ($`10`$ Gyr). We conclude that the number of BH-WD LMXBs could be at most about a factor of 1000 higher than the number of BH-MS LMXBs, but this ratio is still below what is required for BH-WD to become an important contributor. So, if the ratio of BH-MS binaries to BH-RG binaries in ellipticals is similar to that in the Milky Way, BH-WD XRBs will not be a significant contributor to the XLF. Based again on the discrepancy between the duty cycles for WD and RG donors (smaller for WDs by a factor of $`800`$), expect that BH-RG transients dominate over BH-WD transients too.
For the case of an example MT-dependent duty cycle (expressed by $`\eta `$ in eq. 2) we conclude that: (i) if IT MB describes the angular momentum loss best, then only BH-MS transients significantly contribute to the XLFs of elliptical galaxies; consequently $`\beta =2.5\pm 0.6`$; (ii) if instead RVJ MB is a better prescription, then the XLF is dominated by BH-RG binaries and $`\beta =2.8\pm 0.6`$.
For the case of a constant duty cycle independent of the donor type, it is clear that the XRB type with the highest formation rate should dominate the XLF. According to formation rates calculated by Hurley et al. (2002), BH-WD binaries form at a rate about 100 times higher than BH-MS and BH-RG binaries. Consequently, WDs would be expected to dominate the transient population and this is certainly not true for the Milky Way. Therefore we conclude that the assumption of a constant, MT-independent duty cycle is most probably not realistic.
Overall, we conclude that MS or RG donors dominate, depending on whether the IT or RVJ MB prescription is more realistic. Consequently, the slope of the accreting BH mass spectrum is $`\beta =2.5\pm 0.6`$ ($`\beta _0=2.3\pm 0.6`$) or $`\beta =2.8\pm 0.6`$, respectively. These quantitative results are of course dependent on the adopted example form of the MT-dependent duty cycle (eq. ). As shown in ยง 3.1.1 a possible different form (e.g., eq. ) could lead to somewhat flatter values for $`\beta `$.
Next we consider the fact that the upper-end XLF of ellipticals is not a perfect power-law up to arbitrarily high $`L_X`$ values; instead there is a usually smooth cut-off behavior that limits the maximum $`L_X`$ observed at $`2\times 10^{39}`$ erg s<sup>-1</sup>. In Fig. 3 we show the cumulative XLF associated with a model population of BH-MS binaries with a BH mass spectrum with a differential slope of $`\beta =2.3`$ and with an imposed upper limit of $`15`$ M on the maximum BH mass present in the XRB population. We obtain a model XLF that behaves very similarly to observed XLFs (dash-dotted line). Clearly this is just to show the importance of the qualitative effect of a BH mass cut-off on the cumulative XLF.
## 5. DISCUSSION
We consider the upper-end XLF of ellipticals (above the reported break at $`46\times 10^{38}`$ erg s<sup>-1</sup>) and suggest that it is populated by BH X-ray transients at outburst emitting approximately at the Eddington limit. We argue that the upper-end XLF slope is a footprint of the underlying accreting BH mass spectrum modified by a weighting function related to the transient duty cycle. We show that this weighting factor is generally dependent on the BH mass and/or the age of the host galaxy and the derived power-law dependence is different for each of the possible BH donor types: MS, RG, and WD. Our predicted dominance of X-ray transients at outburst contributing to the upper-end XLF could possibly be tested by future high-resolution X-ray observations designed to achieve long-term monitoring probably at time scales of years or longer. Unfortunately, given the uncertainties in the theory of the thermal disk instability, it is not possible to make any predictions about the expected duration of these outbursts.
Based on our analysis and prior population synthesis results we conclude that a constant transient duty cycle independent of the donor type can be excluded. Instead a duty cycle dependent on the binary and MT properties seems to be required. Given the uncertainties associated with transient duty cycles at present, we adopt a couple of different formulations of such a dependence (eqs. ), as reasonable examples, which in no way exhaust the possibilities. In the specific case of a duty cycle that depends on the ratio of the binary mass transfer rate to the critical rate for transient behavior (see eq. 2) we conclude find that the BH X-ray transients forming the upper-end XLF in ellipticals have a dominant donor type and an accreting BH mass spectrum slope $`\beta `$ that depend on the strength of MB angular momentum loss: (i) for the IT MB prescription, only BH-MS transients significantly contribute to the upper-end XLF and $`\beta =2.5\pm 0.6`$ ($`\beta _0=2.3\pm 0.6`$ ); (ii) for the RVJ MB prescription, the XLF is dominated by BH-RG binaries and $`\beta =2.8\pm 0.6`$. We note that these quantitative results do depend on our conclusions about which donor-type population dominates based on currently published population synthesis models (Hurley et al., 2002) and on available observations of BH X-ray systems in our Galaxy. If, e.g., the relative fraction of BH-RG transients in ellipticals is larger than the observed relative fraction in our Galaxy, we expect that BH-RG binaries contribution will lead to a time-dependence of XLF slopes, where younger ellipticals will have a slope predicted for BH-RG binaries, and older ellipticals a steeper slop predicted for BH-MS binaries.
The primary goals of this study are to present (i) the line of arguments that connects the upper-end XLF of ellipticals to BH XRBs formed in the galactic field and (ii) the methodology for how to extract information about the accreting BH mass spectrum from the observed XLF slopes. We have further obtained quantitative results on the BH mass spectrum slope under certain reasonable assumptions, some of which (e.g., the functional form of the MT-dependent duty cycle) represent mere examples. A careful examination of the robustness of these quantitative results has been for the case of MS donors. It has been found that the derived slopes are robust against (i) random variations by factors of a few of the outburst peak luminosity of individual sources, and against (ii) variations of the possible duty-cycle dependence on the critical MT rate for outburst behavior. However, completely difference duty-cycle dependencies cannot be excluded. An improved understanding of this issue would be required to derive reliable quantitative conclusions about the value of the BH mass spectrum slope in transient XRBs in ellipticals.
We expect that our analysis and methods can be used to reveal more information about the formation of BH XRBs in elliptical galaxies. More specifically they could eventually be used to constrain the physical connection between massive stars in XRB progenitors and the resultant BH masses. Current simulations assume either an artificially constant mass for BHs formed (usually at $`10`$ M), or a constant mass fraction of the progenitor leading to the remnant objects, or a remnant mass relation consistent with core-collapse simulations. Constraints on the accreting BH mass spectrum as those discussed here could contribute to our understanding of core collapse, and the connection of BH masses to their progenitor masses.
We thank K. Belczynski, J. McClintock, and R. Remillard for useful discussions and the anonymous referee for suggestions that greatly improved the manuscript and motivated us to perform a number of tests. This work is partially supported by a NASA Chandra Theory Award to N. Ivanova and a Packard Foundation Fellowship in Science and Engineering to V. Kalogera. VK also acknowledges the hospitality of the Aspen Center for Physics where part of this work was completed.
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# Recent Experimental Tests of Special Relativity
## 1 Introduction
One hundred years after Einsteinโs first paper Einstein1905 special relativity is still standing up to all experimental tests and verifications. Over the last century a large number of such tests have provided what is certainly one of the most solid experimental bases of any present fundamental theory of physics. As a consequence special relativity is today underpinning all of present day physics, ranging from the standard model of particle physics (including nuclear and atomic physics) to general relativity and astronomy. That fact continues to push experimentalists to search for new experiments, or improve on previous ones, in order to uncover a possible violation of special relativity, as that would most certainly lead the way to a new conception of physics and of the universe surrounding us. Additional incentive for such tests comes from unification theories (e.g. string theories, loop quantum gravity), some of which KostoSam ; Damour1 ; Gambini suggest a violation of special relativity at some, a priori unknown, level. Given the strong theoretical motivation for such theories, but the lack of experimental data that would allow a more rigorous selection among the candidate theories and the parameter space of each class of such theories, any experimental results that could aid the theoretical efforts are certainly welcome.
The fundamental hypothesis of special relativity is what Einstein termed the โprinciple of relativityโ Einstein1905 , or in more modern terms Local Lorentz Invariance (LLI) Will . Loosely stated, LLI postulates that the outcome of any local test experiment is independent of the velocity of the freely falling apparatus. LLI can be viewed as a constituent part of the Einstein Equivalence Principle which is fundamental to general relativity and all metric theories of gravitation Will . The experiments presented in this paper test some aspect of LLI, as characterized in Lorentz violating theoretical frameworks like the ones briefly described in section 2.
We review and present two of our recent and ongoing experiments Wolf2003 ; WolfGRG ; Wolf2004 that test different aspects of LLI, analyzing and describing their outcome in two theoretical frameworks, the kinematical test theory of Robertson, Mansouri and Sexl (RMS) Robertson ; MaS and the Lorentz violating extension of the standard model (SME) Kosto1 . These experiments, a Michelson-Morley and Kennedy-Thorndike test (section 3), and an ongoing atomic clock test in the SME matter sector (section 4), are among the most precise LLI tests at present.
The vast majority of modern experiments that test LLI rely essentially on the stability of atomic clocks and macroscopic resonators, therefore improvements in oscillator technology have gone hand in hand with improved tests of LLI. The experiments presented here are no exception. All of them employ clocks and resonators developed and used primarily for other purposes (national and international time scales, frequency calibration, etc.) but adapted for tests of LLI.
## 2 Theoretical frameworks
Numerous test theories that allow the modeling and interpretation of experiments that test LLI have been developed. Kinematical frameworks Robertson ; MaS postulate a simple parametrisation of the Lorentz transformations with experiments setting limits on the deviation of those parameters from their special relativistic values. A more fundamental approach is offered by theories that parametrise the coupling between gravitational and non-gravitational fields (TH$`ฯต\mu `$ LightLee ; Will ; Blanchet or $`\chi `$g Ni formalisms) which allow the comparison of experiments that test different aspects of the EEP. Finally, formalisms motivated by unification theories Damour1 ; Damour2 ; Kosto1 have the advantage of opening the way to experimental investigations in the domain of the unification of gravity with the other fundamental forces of nature. In this work we restrict ourselves to two theoretical frameworks, the kinematical framework developed by Robertson, Mansouri and Sexl (RMS) and the more recent standard model extension (SME) of Kostelckรฝ and co-workers.
By construction, kinematical frameworks do not allow for any dynamical effects on the measurement apparatus. This implies that in all inertial frames two clocks of different nature (e.g. based on different atomic species) run at the same relative rate, and two length standards made of different materials keep their relative lengths. Coordinates are defined by the clocks and length standards, and only the transformations between those coordinate systems are modified. In general this leads to observable effects on light propagation in moving frames but, by definition, to no observable effects on clocks and length standards. In particular, no attempt is made at explaining the underlying physics (e.g. modified Maxwell and/or Dirac equations) that could lead to Lorentz violating light propagation but leave e.g. atomic energy levels unchanged. On the other hand dynamical frameworks (e.g. the TH$`ฯต\mu `$ formalism or the SME) in general use a modified general Lagrangian that leads to modified Maxwell and Dirac equations and hence to Lorentz violating light propagation and atomic properties, which is why they are considered more fundamental and more complete than the kinematical frameworks. Furthermore, as shown in KM , the SME is kept sufficiently general to, in fact, encompass the kinematical frameworks and some other dynamical frameworks (in particular the TH$`ฯต\mu `$ formalism) as special cases, although there are no simple and direct relationships between the respective parameters.
### 2.1 The Robertson, Mansouri & Sexl framework
Kinematical frameworks for the description of Lorentz violation have been pioneered by Robertson Robertson and further refined by Mansouri and Sexl MaS and others. Fundamentally the different versions of these frameworks are equivalent, and relations between their parameters are readily obtained. As mentioned above these frameworks postulate generalized transformations between a preferred frame candidate $`\mathrm{\Sigma }(T,๐)`$ and a moving frame $`S(t,๐ฑ)`$ where it is assumed that in both frames coordinates are realized by identical standards. The transformations of MaS (in differential form) for the case where the velocity of $`S`$ as measured in $`\mathrm{\Sigma }`$ is along the positive X-axis, and assuming Einstein synchronization in $`S`$ (in all of the following the choice of synchronization convention plays no role) are
$$dT=\frac{1}{a}\left(dt+\frac{vdx}{c^2}\right);dX=\frac{dx}{b}+\frac{v}{a}\left(dt+\frac{vdx}{c^2}\right);dY=\frac{dy}{d};dZ=\frac{dz}{d}$$
(1)
with $`c`$ the velocity of light in vacuum in $`\mathrm{\Sigma }`$, and $`๐ฏ`$ the velocity of $`S`$ in $`\mathrm{\Sigma }`$. In special relativity $`\alpha _{\mathrm{MS}}=1/2;\beta _{\mathrm{MS}}=1/2;\delta _{\mathrm{MS}}=0`$ and (1) reduces to the usual Lorentz transformations. Generally, the best candidate for $`\mathrm{\Sigma }`$ is taken to be the frame of the cosmic microwave background (CMB) Fixsen ; Lubin with the velocity of the solar system in that frame taken as $`v_{}377`$ km/s, decl. $`6.4^{}`$, $`RA11.2`$h.
Michelson-Morley type experiments MM determine the coefficient $`P_{MM}=(1/2\beta _{\mathrm{MS}}+\delta _{\mathrm{MS}})`$ of the direction dependent term. For many years the most stringent limit on that parameter was $`|P_{MM}|5\times 10^9`$ determined over 23 years ago in an outstanding experiment Brillet . Our experiment WolfGRG confirms that result with roughly equivalent uncertainty $`(2.2\times 10^9)`$. Recently an improvement to $`|P_{MM}|1.5\times 10^9`$ has been reported Muller . Kennedy-Thorndike experiments KT measure the coefficient $`P_{KT}=(\beta _{\mathrm{MS}}\alpha _{\mathrm{MS}}1)`$ of the velocity dependent term. The most stringent limit Schiller on $`|P_{KT}|`$ has been recently improved from Hils by a factor 3 to $`|P_{KT}|2.1\times 10^5`$. Our experiment WolfGRG improves this result by a factor of 70 to $`|P_{KT}|3.0\times 10^7`$. Finally Ives-Stilwell experiments IS measure $`\alpha _{\mathrm{MS}}`$. The most stringent result comes from the recent experiment of Saathoff which improves by a factor 4 our 1997 results WP , limiting $`|\alpha _{\mathrm{MS}}+1/2|`$ to $`2.2\times 10^7`$. The three types of experiments taken together then completely characterize any deviation from Lorentz invariance in this particular test theory, with present limits summarized in table 1 (but note that table 1 does not include new limits reported in these proceedings).
### 2.2 The Standard Model Extension
The general Lorentz violating Standard Model Extension (SME) was developed relatively recently by Kosteleckรฝ and co-workers Kosto1 , motivated initially by possible Lorentz violating phenomenological effects of string theory KostoSam . It consists of a parametrised version of the standard model Lagrangian that includes all Lorentz violating terms that can be formed from known fields, and includes (in its most recent version KostoGrav ) gravity.
The fundamental theory of the SME as applied to electrodynamics is laid out in KM and summarized below. We use that approach to model the MM and KT experiments in section 3.2. For the discussion of the atomic clock experiment of section 4 the SME matter sector is relevant. Its application to atomic physics, and in particular atomic clock experiments, is laid out in KL ; Bluhm and summarized below.
Generally, the SME characterizes a potential Lorentz violation using a number of parameters that are all zero in standard (non Lorentz violating) physics. These parameters are frame dependent and consequently vary as a function of the coordinate system chosen to analyze a given experiment. In principle they may be constant and non-zero in any frame (e.g. the lab frame). However, any non-zero values are expected to arise from Planck-scale effects in the early Universe. Therefore they should be constant in a cosmological frame (e.g. the one defined by the CMB radiation) or any frame that moves with a constant velocity and shows no rotation with respect to the cosmological one. Consequently the conventionally chosen frame to analyze and compare experiments in the SME is a sun-centered, non-rotating frame as defined in KM . The general procedure is to calculate the SME perturbation of the experimental observable in the lab frame (or cavity frame, or atom frame) and then to transform the lab frame SME parameters to the conventional sun-centered frame. This transformation will introduce a time variation of the frequency related to the movement of the lab with respect to the sun-centered frame (typically introducing time variations of sidereal and semi-sidereal periods for an Earth fixed experiment).
#### SME photon sector
The photon sector of the SME is described by a Lagrangian that takes the form
$$=\frac{1}{4}F_{\mu \nu }F^{\mu \nu }+\frac{1}{2}(k_{AF})^\kappa ฯต_{\kappa \lambda \mu \nu }A^\lambda F^{\mu \nu }\frac{1}{4}(k_F)_{\kappa \lambda \mu \nu }F^{\kappa \lambda }F^{\mu \nu }$$
(2)
where $`F_{\mu \nu }_\mu A_\nu _\nu A_\mu `$. The first term is the usual Maxwell part while the second and third represent Lorentz violating contributions that depend on the parameters $`k_{AF}`$ and $`k_F`$. For most analysis the $`k_{AF}`$ parameter is set to 0 for theoretical reasons (c.f. KM ), which is also well supported experimentally. The remaining dimensionless tensor $`(k_F)_{\kappa \lambda \mu \nu }`$ has a total of 19 independent components that need to be determined by experiment. Retaining only this term leads to Maxwell equations that take the familiar form but with $`๐`$ and $`๐`$ fields defined by a general matrix equation
$$\left(\begin{array}{c}๐\\ ๐\end{array}\right)=\left(\begin{array}{c}ฯต_0(\stackrel{~}{ฯต_r}+\kappa _{DE})\\ \sqrt{\frac{ฯต_0}{\mu _0}}\kappa _{HE}\end{array}\begin{array}{c}\sqrt{\frac{ฯต_0}{\mu _0}}\kappa _{DB}\\ \mu _0^1(\stackrel{~}{\mu _r}^1+\kappa _{HB})\end{array}\right)\left(\begin{array}{c}๐\\ ๐\end{array}\right)$$
(3)
where the $`\kappa `$ are $`3\times 3`$ matrices whose components are particular combinations of the $`k_F`$ tensor (c.f. equation (5) of KM ). If we suppose the medium of interest has general magnetic or dielectric properties, then $`\stackrel{~}{ฯต_r}`$ and $`\stackrel{~}{\mu _r}`$ are also 3 x 3 matrices. In vacuum $`\stackrel{~}{ฯต_r}`$ and $`\stackrel{~}{\mu _r}`$ are identity matrices. Equation (3) indicates a useful analogy between the SME in vacuum and standard Maxwell equations in homogeneous anisotropic media.
For the analysis of different experiments it turns out to be useful to introduce further combinations of the $`\kappa `$ matrices defined by:
$`(\stackrel{~}{\kappa }_{e+})^{jk}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(\kappa _{DE}+\kappa _{HB})^{jk},`$
$`(\stackrel{~}{\kappa }_e)^{jk}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(\kappa _{DE}\kappa _{HB})^{jk}{\displaystyle \frac{1}{3}}\delta ^{jk}(\kappa _{DE})^{ll},`$
$`(\stackrel{~}{\kappa }_{o+})^{jk}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(\kappa _{DB}+\kappa _{HE})^{jk},`$
$`(\stackrel{~}{\kappa }_o)^{jk}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(\kappa _{DB}\kappa _{HE})^{jk},`$
$`\stackrel{~}{\kappa }_{tr}`$ $`=`$ $`{\displaystyle \frac{1}{3}}(\kappa _{DE})^{ll}.`$ (4)
The first four of these equations define traceless $`3\times 3`$ matrices, while the last defines a single coefficient. All $`\stackrel{~}{\kappa }`$ matrices are symmetric except $`\stackrel{~}{\kappa }_{o+}`$ which is antisymmetric. These characteristics leave a total of 19 independent coefficients of the $`\stackrel{~}{\kappa }`$.
In general experimental results are quoted and compared using the $`\stackrel{~}{\kappa }`$ parameters rather than the original $`k_F`$ tensor components. The 10 independent components of the $`\stackrel{~}{\kappa }_{e+}`$ and $`\stackrel{~}{\kappa }_o`$ tensors, have been determined to $`2\times 10^{32}`$ by astrophysical tests KM . Of the 9 remaining independent components, 4 components of $`\stackrel{~}{\kappa }_e`$ and the 3 components of $`\stackrel{~}{\kappa }_{o+}`$ have been bounded by the resonator experiments reported here and in Muller ; Wolf2004 to parts in $`10^{15}`$ and $`10^{11}`$ respectively, with our results improving by up to a factor 10 on the best previous ones (c.f. Tab.7). The scalar $`\stackrel{~}{\kappa }_{tr}`$ has been bounded recently by our SME analysis Tobar2004 of the experiment of Saathoff to parts in $`10^5`$. In Tobar2004 we also propose several interferometer and resonator experiments that could improve the limit on $`\stackrel{~}{\kappa }_{tr}`$ to parts in $`10^{11}`$ and the limits on $`\stackrel{~}{\kappa }_{o+}`$ to parts in $`10^{15}`$. Finally, the remaining component $`\stackrel{~}{\kappa }_e^{ZZ}`$ is undetermined at present as it is not accessible to Earth fixed experiments. However, it should be accessible to experiments that are rotating in the laboratory, like the ones reported elsewhere in these proceedings, which should yield the first limits on that parameter and thereby complete the coverage of the parameter space in the SME photon sector. Present limits are summarized in Tab.2 (not including new limits reported in these proceedings).
#### SME matter sector
In the matter sector, the SME modifies the Lagrangian of a spin 1/2 fermion KostoCol ; KL via a number of parameterized Lorentz violating terms. When applied to atomic physics, this leads to a perturbation of the standard model Hamiltonian parametrised by 40 parameters for each fundamental particle (proton, neutron, electron), which in turn leads to a shift of the atomic energy levels and atomic transition frequencies (see KL ; Bluhm for details). Quite generally, the energy level shifts can be expressed in the form
$$\mathrm{\Delta }E=\widehat{m}_F(E_d^e+E_d^p+E_d^n)+\stackrel{~}{m}_F(E_q^e+E_q^p+E_q^n)$$
(5)
where $`E_d`$ and $`E_q`$ are energies given below, the superscripts $`e,p,n`$ stand for electron, proton and neutron and $`\widehat{m}_F`$ and $`\stackrel{~}{m}_F`$ are defined as
$$\widehat{m}_F:=\frac{m_F}{F},\stackrel{~}{m}_F:=\frac{3m_F^2F(F+1)}{3F^2F(F+1)}.$$
(6)
In general $`\mathrm{\Delta }E`$ of (5) will be time varying as the energies $`E_d^w,E_q^w`$ ($`w`$ stands for $`e,p,n`$) depend on the orientation of the angular momentum of $`w`$ with respect to the fixed stars (best approximation to the frame in which symmetry breaking took place in the early universe). Of particular interest will be (see section 4) Zeeman sublevels with $`m_F0`$ in which case the orientation of the quantization axis (quantization magnetic field) is relevant, so one can expect variations of $`\mathrm{\Delta }E`$ at sidereal and semi-sidereal frequencies due to the precession of the quantization axis with the rotation of the Earth.
The energies in (5) are KL
$`E_d^w`$ $`=`$ $`\beta _w\stackrel{~}{b}_3^w+\delta _w\stackrel{~}{d}_3^w+\kappa _w\stackrel{~}{g}_d^w`$
$`E_q^w`$ $`=`$ $`\gamma _w\stackrel{~}{c}_q^w+\lambda _w\stackrel{~}{g}_q^w.`$ (7)
In (7) the tilde quantities have the dimensions of energy and represent laboratory frame combinations of the SME parameters that need to be determined by experiment. They are time varying at sidereal and semi-sidereal frequencies as they are obtained by transforming the constant sun-centered-frame parameters to the laboratory frame. The other coefficients in (7) are constant and depend on the nuclear and electronic structure of the atom. Explicit expressions can be found in KL , with their values calculated for certain atoms and states (including the <sup>133</sup>Cs atom of interest to our experiment) in Bluhm .
Our experiment (see section 4) is sensitive to $`\stackrel{~}{c}_q^p`$. When transforming to the sun-centered-frame this parameter is a time varying combination of 8 constant SME parameters ($`\stackrel{~}{c}_Q`$, $`\stackrel{~}{c}_{}`$, $`\stackrel{~}{c}_X`$, $`\stackrel{~}{c}_Y`$, $`\stackrel{~}{c}_Z`$, $`\stackrel{~}{c}_{TX}`$, $`\stackrel{~}{c}_{TY}`$, $`\stackrel{~}{c}_{TZ}`$), which are generally used KL ; Bluhm to state and compare experimental results (see Tab. 3). In some publications Muller2005 ; Lane2005 the results are stated in terms of dimensionless sun-frame parameters related to the $`\stackrel{~}{c}`$ parameters by (c.f. Bluhm Appendix B)
$`\stackrel{~}{c}_Q`$ $`=`$ $`mc^2(c_{XX}+c_{YY}2c_{ZZ})`$
$`\stackrel{~}{c}_{}`$ $`=`$ $`mc^2(c_{XX}c_{YY})`$ (8)
$`\stackrel{~}{c}_J`$ $`=`$ $`mc^2|ฯต_{JKL}|c_{KL}`$
$`\stackrel{~}{c}_{TJ}`$ $`=`$ $`mc^2(c_{TJ}+c_{JT})`$
where m is the mass of the particle ($`m_n`$, $`m_p`$, or $`m_e`$), indices $`J,K,L`$ run over sun-frame spatial coordinates $`X,Y,Z`$ and the totally antisymmetric tensor $`ฯต_{JKL}`$ is defined with $`ฯต_{XYZ}=+1`$.
Existing bounds on the 40 parameters for each particle (n,p,e) come from clock comparison and magnetometer experiments using different atomic species (Bluhm and references therein, Cane ), from resonator experiments (including our experiments described in section 3 as analyzed recently by Mรผller) Muller2005 ; Wolf2004 ; Muller , and from analysis of Ives-Stilwell (Doppler-shift) experiments Lane2005 ; Saathoff . They are summarized in Tab. 3 below. The expected results of our present experiment (see section 4) are given in Tab. 3 in brackets. They correspond to first measurements of some parameters and an improvement by 11 and 14 orders of magnitude on others.
## 3 Michelson-Morley and Kennedy-Thorndike tests
In this section we review the results Wolf2003 ; WolfGRG ; Wolf2004 of our experiment that compares the frequencies of a cryogenic sapphire oscillator (CSO) and a hydrogen maser atomic clock. Both devices operate at microwave frequencies and are run and compared continuously for timekeeping purposes at the Paris observatory. We use that data to carry out Michelson-Morley and Kennedy-Thorndike experiments, searching for a dependence of the difference frequency on the orientation and/or the velocity of the CSO with respect to a prefered frame candidate.
The heart of the experiment is a monolithic sapphire crystal of cylindrical shape, about 5 cm diameter and 3 cm height. The resonance frequency is determined by exciting a so called Whispering Gallery mode, corresponding to a standing wave set up around the perimeter of the cylinder (see fig. 1 and WolfGRG for a detailed description). In our case the excited mode is a TE mode at 11.932 GHz, with dominant radial electric and vertical magnetic fields corresponding to propagation (Poynting) vectors in both directions around the circumference. The CSO is an active system oscillating at the resonant frequency (i.e. a classical loop oscillator which amplifies and re-injects the โnaturalโ resonator signal). Additionally the signal is locked to the resonance using the Pound-Drever technique (modulation at $``$ 80 kHz). The incident power is stabilized in the cryogenic environment and the spurious AM modulation is minimized using a servo loop. To minimize temperature sensitivity the resonator is heated (inside the 4 K environment) and stabilized to the temperature turning point ($``$ 6 K) of the resonator frequency which arises due to paramagnetic impurities in the sapphire. Under these conditions the loaded quality factor of the resonator is slightly below $`10^9`$. The resonator is kept permanently at cryogenic temperatures, with helium refills taking place about every 20 - 25 days.
The CSO is compared to a commercial (Datum Inc.) active hydrogen maser whose frequency is also regularly compared to caesium and rubidium atomic fountain clocks in the laboratory Bize . The CSO resonant frequency at 11.932 GHz is compared to the 100 MHz output of the hydrogen maser. The maser signal is multiplied up to 12 GHz of which the CSO signal is subtracted. The remaining $``$ 67 MHz signal is mixed to a synthesizer signal at the same frequency and the low frequency beat at $``$ 64 Hz is counted, giving access to the frequency difference between the maser and the CSO. The instability of the comparison chain has been measured at $`2\times 10^{14}\tau ^1`$, with long term instabilities dominated by temperature variations, but not exceeding $`10^{16}`$.
### 3.1 Results in the RMS framework
In the RMS framework our experiment sets the most stringent limit for Kennedy-Thorndike experiments (improving by a factor 70 over previous results) and is among the most precise Michelson-Morley tests (see table 1). Those results were reported in Wolf2003 ; WolfGRG and are summarized here.
In the RMS framework the frequency of a resonator in the lab frame $`S`$ is proportional to $`t_{c}^{}{}_{}{}^{1}`$ where $`t_c`$ is the return travel time of a light signal in the resonator. Setting $`c^2dT^2=dX^2+dY^2+dZ^2`$ in the preferred frame $`\mathrm{\Sigma }`$, and transforming according to (1) we find the coordinate travel time of a light signal in $`S`$:
$$dt=\frac{dl}{c}\left(1\left(\beta _{\mathrm{MS}}\alpha _{\mathrm{MS}}1\right)\frac{v^2}{c^2}\left(\frac{1}{2}\beta _{\mathrm{MS}}+\delta _{\mathrm{MS}}\right)\mathrm{sin}^2\theta \frac{v^2}{c^2}\right)+๐ช(4)$$
(9)
where $`dl=\sqrt{dx^2+dy^2+dz^2}`$ and $`\theta `$ is the angle between the direction of light propagation and the velocity v of $`S`$ in $`\mathrm{\Sigma }`$.
Calculating $`t_c`$ from (9) the relative frequency difference between the sapphire oscillator and the hydrogen maser (which, by definition, realizes coordinate time in $`S`$ masercom ) is
$$\frac{\mathrm{\Delta }\nu (t)}{\nu _0}=P_{KT}\frac{v(t)^2}{c^2}+P_{MM}\frac{v(t)^2}{c^2}\frac{1}{2\pi }_0^{2\pi }\mathrm{sin}^2\theta (t,\phi )๐\phi +๐ช(3)$$
(10)
where $`\nu _0`$ is the unperturbed frequency, $`v(t)`$ is the (time dependent) speed of the lab in $`\mathrm{\Sigma }`$, and $`\phi `$ is the azimuthal angle of the light signal in the plane of the cylinder. The periodic time dependence of $`v`$ and $`\theta `$ due to the rotation and orbital motion of the Earth with respect to the CMB frame allow us to set limits on the two parameters in (10) by fitting the periodic terms of appropriate frequency and phase (see Mike for calculations of similar effects for several types of oscillator modes). Given the limited durations of our data sets ($``$ 16 days) the dominant periodic terms arise from the Earthโs rotation, so retaining only those we have $`๐ฏ(t)=๐ฎ+\omega \times ๐`$ with $`๐ฎ`$ the velocity of the solar system with respect to the CMB, $`\omega `$ the angular velocity of the Earth, and $`๐`$ the geocentric position of the lab. We then find after some calculation.
$$\begin{array}{cc}\mathrm{\Delta }\nu /\nu _0& =P_{KT}(H\mathrm{sin}\lambda )\hfill \\ & +P_{MM}(A\mathrm{cos}\lambda +B\mathrm{cos}(2\lambda )+C\mathrm{sin}\lambda +D\mathrm{sin}\lambda \mathrm{cos}\lambda +E\mathrm{sin}\lambda \mathrm{cos}(2\lambda ))\hfill \end{array}$$
(11)
where $`\lambda =\omega t+\varphi `$, and A-E and $`\varphi `$ are constants depending on the latitude and longitude of the lab $`(48.7^{}`$N and $`2.33^{}`$E for Paris). Numerically $`H2.6\times 10^9`$, $`A8.8\times 10^8`$, $`B1.8\times 10^7`$, C-E of order $`10^9`$. We note that in (11) the dominant time variations of the two combinations of parameters are in quadrature and at twice the frequency which indicates that they should decorelate well in the data analysis allowing a simultaneous determination of the two (as confirmed by the correlation coefficients given below). Fitting this simplified model to our data we obtain results that differ by less than 10% from the results presented below that were obtained using the complete model ((10) including the orbital motion of the Earth).
For the RMS analysis we use 13 data sets in total spanning Sept. 2002 to Aug. 2003, of differing lengths (5 to 16 days, 140 days in total). The sampling time for all data sets was $`100`$ s except two data sets with $`\tau _0=12`$ s. To make the data more manageable we first average all points to $`\tau _0=2500`$ s. For the data analysis we simultaneously fit (using weighted least squares, WLS, c.f. WolfGRG ) an offset and a rate (natural frequency drift, typically $`1.7\times 10^{18}`$ s<sup>-1</sup>) per data set and the two parameters of the model (10). In the model (10) we take into account the rotation of the Earth and the Earthโs orbital motion, the latter contributing little as any constant or linear terms over the durations of the individual data sets are absorbed by the fitted offsets and rates.
Figure 2 shows the resulting values of the two parameters ($`P_{KT}`$ and $`P_{MM}`$) for each individual data set. A global WLS fit of the two parameters and the 13 offsets and drifts yields $`P_{MM}=(1.2\pm 1.9)\times 10^9`$ and $`P_{KT}=(1.6\pm 2.3)\times 10^7`$ ($`1\sigma `$ uncertainties), with the correlation coefficient between the two parameters less than 0.01 and all other correlation coefficients $`<0.06`$. The distribution of the 13 individual values around the ones obtained from the global fit is well compatible with a normal distribution ($`\chi ^2`$ = 10.7 and $`\chi ^2`$ = 14.6 for $`P_{MM}`$ and $`P_{KT}`$ respectively).
Systematic effects at diurnal or semi-diurnal frequencies with the appropriate phase could mask a putative sidereal signal. The statistical uncertainties of $`P_{MM}`$ and $`P_{KT}`$ obtained from the WLS fit above correspond to sidereal and semi-sidereal terms (from (11)) of $`7\times 10^{16}`$ and $`4\times 10^{16}`$ respectively so any systematic effects exceeding these limits need to be taken into account in the final uncertainty. We expect the main contributions to such effects to arise from temperature, pressure and magnetic field variations that would affect the hydrogen maser, the CSO and the associated electronics, and from tilt variations of the CSO which are known to affect its frequency (see section 3.2 for a detailed discussion). Our final uncertainties (the error bars in Fig. 2) are the quadratic sums of the statistical uncertainties from the WLS adjustment for each data set and the systematic uncertainties calculated for each data set from (11). For the global adjustment we average the systematic uncertainties from the individual data sets obtaining $`\pm 1.2\times 10^9`$ on $`P_{MM}`$ and $`\pm 1.9\times 10^7`$ on $`P_{KT}`$.
In the RMS framework, our experiment simultaneously constrains two combinations of the three parameters of the Mansouri and Sexl test theory (previously measured individually by Michelson-Morley and Kennedy-Thorndike experiments). We obtain $`\delta _{\mathrm{MS}}\beta _{\mathrm{MS}}+1/2=1.2(1.9)(1.2)\times 10^9`$ which is of the same order as the best previous results Muller ; Brillet , and $`\beta _{\mathrm{MS}}\alpha _{\mathrm{MS}}1=1.6(2.3)(1.9)\times 10^7`$ which improves the best previous limit Schiller by a factor of 70 (the first bracket indicates the $`1\sigma `$ uncertainty from statistics the second from systematic effects). We note that our value on $`\delta _{\mathrm{MS}}\beta _{\mathrm{MS}}+1/2`$ is compatible with the slightly significant recent result of Muller who obtained $`\delta _{\mathrm{MS}}\beta _{\mathrm{MS}}+1/2=(2.2\pm 1.5)\times 10^9`$.
As a result of our experiment the Lorentz transformations are confirmed in the RMS framework (c.f. Tab. 1) with an overall uncertainty of $`3\times 10^7`$ limited by our determination of $`\beta _{\mathrm{MS}}\alpha _{\mathrm{MS}}1`$ and the recent limit Saathoff of $`2.2\times 10^7`$ on the determination of $`\alpha _{\mathrm{MS}}`$. The latter is likely to improve in the coming years by experiments such as ACES (Atomic Clock Ensemble in Space ACES ) that will compare ground clocks to clocks on the international space station aiming at a $`10^8`$ measurement of $`\alpha _{\mathrm{MS}}`$.
### 3.2 Results in the SME
In the SME our experiment sets the presently most stringent limits on a number of photon sector parameters, improving previous results Muller by up to an order of magnitude. These results were first published in Wolf2004 and are reproduced here.
The SME perturbed frequency of a resonator can be calculated from equation (3) in the form (c.f. KM )
$`{\displaystyle \frac{\mathrm{\Delta }\nu }{\nu _0}}=`$ $``$ $`{\displaystyle \frac{1}{U}}{\displaystyle _V}d^3x(ฯต_0๐_{\mathrm{๐}}^{}{}_{}{}^{}\kappa _{DE}๐_\mathrm{๐}\mu _0^1๐_{\mathrm{๐}}^{}{}_{}{}^{}\kappa _{HB}๐_\mathrm{๐}`$
$`+`$ $`2\mathrm{R}\mathrm{e}(\sqrt{{\displaystyle \frac{ฯต_0}{\mu _0}}}๐_{\mathrm{๐}}^{}{}_{}{}^{}\kappa _{DB}๐_\mathrm{๐}))`$
where $`๐_\mathrm{๐},๐_\mathrm{๐},๐_\mathrm{๐},๐_\mathrm{๐}`$ are the unperturbed (standard Maxwell) fields and $`U=_Vd^3x(๐_\mathrm{๐}๐_{\mathrm{๐}}^{}{}_{}{}^{}+๐_\mathrm{๐}๐_{\mathrm{๐}}^{}{}_{}{}^{})`$. Note that, as shown in KL , the frequency of the H-maser is not affected to first order (because it operates on $`m_F=0`$ states) and Muller2 shows that the perturbation of the frequency due to the modification of the sapphire crystal structure (and hence the cavity size) is negligible with respect to the direct perturbation of the e-m fields.
The resonator is placed in the lab with its symmetry axis along the vertical. Applying (3.2) in the lab frame (z-axis vertical upwards, x-axis pointing south), with the fields calculated using a finite element technique as described in WolfGRG , we obtain an expression for the frequency variation of the resonator
$`{\displaystyle \frac{\mathrm{\Delta }\nu }{\nu _0}}`$ $`=`$ $`(_{DE})_{lab}^{xx}\left((\kappa _{DE})_{lab}^{xx}+(\kappa _{DE})_{lab}^{yy}\right)+(_{DE})_{lab}^{zz}(\kappa _{DE})_{lab}^{zz}`$ (13)
$`+`$ $`(_{HB})_{lab}^{xx}\left((\kappa _{HB})_{lab}^{xx}+(\kappa _{HB})_{lab}^{yy}\right)+(_{HB})_{lab}^{zz}(\kappa _{HB})_{lab}^{zz}`$
with the $`_{lab}`$ components given in Tab. 4. To obtain the values in Tab. 4 we take into account the fields inside the resonator (c.f. WolfGRG ) and outside ($`2\%`$ of the energy).
The last step is to transform the $`\kappa `$ tensors in (13) to the conventional sun-centered frame using the explicit transformations provided in KM , and to express the result in terms of the $`\stackrel{~}{\kappa }`$ tensors of (4). We obtain
$$\frac{\nu \nu _0}{\nu _0}=\underset{i}{}C_i\mathrm{cos}(\omega _iT_{}+\phi _i)+S_i\mathrm{sin}(\omega _iT_{}+\phi _i)$$
(14)
where $`\nu _0`$ is the unperturbed frequency difference, the sum is over the six frequencies $`\omega _i`$ of Tab.5, the coefficients $`C_i`$ and $`S_i`$ are functions of the Lorentz violating tensors $`\stackrel{~}{\kappa }_e`$ and $`\stackrel{~}{\kappa }_{o+}`$ (see Tab.5), $`T_{}=0`$ on December 17, 2001, 18:05:16 UTC, $`\phi _\omega _{}=\phi _{2\omega _{}}=0`$ and $`\phi _{(\omega _{}\pm \mathrm{\Omega }_{})}=\phi _{(2\omega _{}\pm \mathrm{\Omega }_{})}=\pm 4.682`$ rad. To obtain the relations of Tab.5 between $`C_i`$, $`S_i`$ and the SME parameters we have assumed zero values for the 10 independent components of the $`\stackrel{~}{\kappa }_{e+}`$ and $`\stackrel{~}{\kappa }_o`$ tensors, as those have been determined to $`2\times 10^{32}`$ by astrophysical tests KM .
To determine all 7 SME parameters appearing in Tab.5 one requires over a year of data in order to be able to decorrelate the annual sidebands from the sidereal and twice sidereal frequencies. To do so we have extended the data to 20 data sets in total, spanning Sept. 2002 to Jan. 2004, of differing lengths (5 to 20 days, 222 days in total). The sampling time for all data sets was $`100`$ s.
For the statistical analysis we first average the data to 2500 s sampling time and then simultaneously fit the 20 rates and offsets and the 12 parameters $`C_i`$ and $`S_i`$ of (14) to the complete data using two statistical methods, weighted least squares (WLS), which allows one to account for non-white noise processes (cf. Wolf2003 ), and individual periods (IP) as used in Muller . The two methods give similar results for the parameters (within the uncertainties) but differ in the estimated uncertainties (the IP uncertainties are a factor $`1.2`$ larger). Because IP discards a significant amount of data (about 10% in our case) we consider WLS the more realistic method and retain those results as the statistical uncertainties shown in Tab.5. We note that we now have sufficient data to decorrelate all 12 parameters ($`C_i`$, $`S_i`$) i.e. the WLS correlation coefficients between any two parameters or between any parameter and the fitted offsets and rates are all less than 0.20.
To investigate the distributions of our results we fit the coefficients $`C_i`$ and $`S_i`$ to each one of the 20 data sets individually with the results at the sidereal and semi-sidereal frequencies $`\omega _{}`$ and $`2\omega _{}`$ shown in Fig.3. If a genuine effect at those frequencies was present we would expect correlated phases of the individual points in Fig.3, but this does not seem to be supported by the data. A distribution of the phases may result from an effect at a neighboring frequency, in particular the diurnal and semi-diurnal frequencies $`\omega _{}\mathrm{\Omega }_{}`$ and $`2(\omega _{}\mathrm{\Omega }_{})`$ at which we would expect systematic effects to play an important role. Fig. 4 shows the amplitudes $`A_\omega =\sqrt{C_\omega ^2+S_\omega ^2}`$ resulting from least squares fits for a range of frequencies, $`\omega `$, around the frequencies of interest. We note that the fitted amplitudes at $`\omega _{}\mathrm{\Omega }_{}`$ and $`2(\omega _{}\mathrm{\Omega }_{})`$ are substantially smaller than those at $`\omega _{}`$ and $`2\omega _{}`$ and therefore unlikely to contribute to the distribution of the points in Fig.3.
Systematic effects at the frequencies $`\omega _i`$ could mask a putative Lorentz violating signal in our experiment and need to be investigated in order to be able to confirm such a signal or to exclude it within realistic limits. We have extensively studied all systematic effects arising from environmental factors that might affect our experiment. The resulting estimated contributions at the two central frequencies $`\omega _{}`$, $`2\omega _{}`$ and at the diurnal frequency $`\omega _{}\mathrm{\Omega }_{}`$ are summarized in Tab.6. The contributions at $`\omega _{}+\mathrm{\Omega }_{}`$ and $`2\omega _{}\pm \mathrm{\Omega }_{}`$ are not shown as they are identical to those at $`\omega _{}`$ and $`2\omega _{}`$ respectively.
We have compared the Hydrogen-maser (HM) used as our frequency reference to our highly stable and accurate Cs fountain clocks (FO2 and FOM). For example, the amplitudes at $`\omega _{}`$ and $`2\omega _{}`$ of the HM-FOM relative frequency difference over June-July 2003 were $`A_\omega _{}=(4.8\pm 4.7)\times 10^{16}`$ and $`A_{2\omega _{}}=(4.3\pm 4.7)\times 10^{16}`$. This indicates that any environmental effects on the HM at those frequencies should be below 5 parts in $`10^{16}`$ in amplitude. This is in good agreement with studies on similar HMs carried out in Parker that limited environmental effects to $`<`$ 3 to 4 parts in $`10^{16}`$.
To estimate the tilt sensitivity we have intentionally tilted the oscillator by $``$ 5 mrad off its average position which led to relative frequency variations of $`3\times 10^{13}`$ from which we deduce a tilt sensitivity of $`6\times 10^{17}\mu `$rad<sup>-1</sup>. This is in good agreement with similar measurements in ChangTh that obtained sensitivities of $`4\times 10^{17}\mu `$rad<sup>-1</sup>. Measured tilt variations in the lab at diurnal and semi-diurnal periods show amplitudes of 4.6 $`\mu `$rad and 1.6 $`\mu `$rad respectively which leads to frequency variations that do not exceed $`3\times 10^{16}`$ and $`1\times 10^{16}`$ respectively.
From the measurements of tilt sensitivity one can deduce the sensitivity to gravity variations (cf. ChangTh ), which in our case lead to a sensitivity of $`3\times 10^{10}g^1`$. Tidal gravity variations can reach $`10^7g`$ from which we obtain a maximum effect of $`3\times 10^{17}`$, one order of magnitude below the effect from tilt variations.
Variations of the ambient magnetic field in our lab. are dominated by the passage of the Paris Metro, showing a strong periodicity (โquietโ periods from 1 am to 5 am). The corresponding diurnal and semi-diurnal amplitudes are $`1.7\times 10^4`$ G and $`3.4\times 10^4`$ G respectively for the vertical field component and about 10 times less for the horizontal one. To determine the magnetic sensitivity of the CSO we have applied a sinusoidal vertical field of 0.1 G amplitude with a 200 s period. Comparing the CSO frequency to the FO2 Cs-fountain we see a clear sinusoidal signal (S/N $`>2`$) at the same period with an amplitude of $`7.2\times 10^{16}`$, which leads to a sensitivity of $`7\times 10^{15}`$ G<sup>-1</sup>. Assuming a linear dependence (there is no magnetic shielding that could lead to non-linear effects) we obtain effects of only a few parts in $`10^{18}`$.
Late 2002 we implemented an active temperature stabilization inside an isolated volume ($`15\mathrm{m}^3`$) that includes the CSO and all the associated electronics. The temperature is measured continously in two fixed locations (behind the electronics rack and on top of the dewar). For the best data sets the measured temperature variations do not exceed 0.02/0.01 K in amplitude for the diurnal and semi-diurnal components. A least squares fit to all our temperature data (taken simultaneously with our frequency measurements) yields amplitudes of $`A_\omega _{}=0.020`$ K and $`A_{2\omega _{}}=0.018`$ K with similar values at the other frequencies $`\omega _i`$ of interest, including the diurnal one ($`A_{\omega _{}\mathrm{\Omega }_{}}=0.022`$ K). Inducing a strong sinusoidal temperature variation ($`0.5`$ K amplitude at 12 h period) leads to no clearly visible effect on the CSO frequency. Taking the noise level around the 12 h period as the maximum effect we obtain a sensitivity of $`<4\times 10^{15}`$ per K. Using this estimate we obtain effects of $`<1\times 10^{16}`$ at all frequencies $`\omega _i`$.
Finally we have investigated the sensitivity of the CSO to atmospheric pressure variations. To do so we control the pressure inside the dewar using a variable valve mounted on the He-gas exhaust. During normal operation the valve is open and the CSO operates at ambient atmospheric pressure. For the sensitivity determination we have induced a sinusoidal pressure variation ($`14`$ mbar amplitude at 12 h period), which resulted in a clearly visible effect on the CSO frequency corresponding to a sensitivity of $`6.5\times 10^{16}`$ mbar<sup>-1</sup>. We have checked that the sensitivity is not significantly affected when changing the amplitude of the induced pressure variation by a factor 3. A least squares fit to atmospheric pressure data (taken simultaneously with our frequency measurements) yields amplitudes of $`A_\omega _{}=0.045`$ mbar and $`A_{2\omega _{}}=0.054`$ mbar with similar values at the other frequencies $`\omega _i`$ of interest, except the diurnal one for which $`A_{\omega _{}\mathrm{\Omega }_{}}=0.36`$ mbar. The resulting effects on the CSO frequency are given in Tab.6.
Our final results for the 7 components of $`\stackrel{~}{\kappa }_e`$ and $`\stackrel{~}{\kappa }_{o+}`$ are obtained from a least squares fit to the 12 measured coefficients of Tab.5. They are summarized and compared to the results of Muller in Tab.7.
We note that our results for $`\stackrel{~}{\kappa }_e^{XY}`$ and $`\stackrel{~}{\kappa }_e^{XZ}`$ are significant at about $`2\sigma `$, while those of Muller are significant at about the same level for $`(\stackrel{~}{\kappa }_e^{XX}\stackrel{~}{\kappa }_e^{YY})`$. The two experiments give compatible results for $`\stackrel{~}{\kappa }_e^{XZ}`$ (within the $`1\sigma `$ uncertainties) but not for the other two parameters, so the measured values of those are unlikely to come from a common source. Another indication for a non-genuine effect comes from figures 3 and 4, as we would expect any genuine effect to show an approximately coherent phase for the individual data sets in figure 3 and to display more prominent peaks in figure 4.
In conclusion, we have not seen any Lorentz violating effects in the general framework of the SME, and set limits on 7 parameters of the SME photon sector (cf. Tab. 7) which are up to an order of magnitude more stringent than those obtained from previous experiments Muller . Two of the parameters are significant (at $`2\sigma `$). We believe that this is most likely a statistical coincidence or a neglected systematic effect. To verify this, our experiment is continuing and new, more precise experiments are under way Mike .
## 4 Atomic clock test of Lorentz invariance in the SME matter sector
For this experiment we use one of the laser cooled fountain clocks operated at the Paris observatory, the <sup>133</sup>Cs and <sup>87</sup>Rb double fountain FO2 BizeJPB . We run it in Cs mode on the $`|F=4|F=3`$ hyperfine transition of the $`6S_{1/2}`$ ground state. Both hyperfine states are degenerate, with Zeeman substates $`m_F=[4,4]`$ and $`m_F=[3,3]`$ respectively. The clock transition used in routine operation is $`|F=4,m_F=0|F=3,m_F=0`$ at 9.2 GHz, which is magnetic field independent to first order. The first order magnetic field dependent Zeeman transitions ($`|F=4,m_F=i|F=3,m_F=i`$ with $`i=\pm 1,\pm 2,\pm 3`$) are used regularly for measurement and characterization of the magnetic field, necessary to correct the second order Zeeman effect of the clock transition. In routine operation the clock transition frequency stability of FO2 is $`1.6\times 10^{14}\tau ^{1/2}`$, and its accuracy $`7\times 10^{16}`$ BizeJPB ; Marion , the best performance of any clock at present.
In the presence of Lorentz violation the SME frequency shift of a $`Cs`$ $`|F=4,m_F|F=3,m_F`$ transition, arising from the energy level shifts described in section 2.2, has been calculated explicitly in Bluhm . It can be written in the form
$`\mathrm{}(\delta \omega _{SME})`$ $`=`$ $`s_1^p\left(\beta _p\stackrel{~}{b}_3^p\delta _p\stackrel{~}{d}_3^p+\kappa _p\stackrel{~}{g}_d^p\right)+s_2^p\left(\gamma _p\stackrel{~}{c}_q^p\lambda _p\stackrel{~}{g}_q^p\right)`$ (15)
$`+`$ $`s_1^e\left(\beta _e\stackrel{~}{b}_3^e\delta _e\stackrel{~}{d}_3^e+\kappa _e\stackrel{~}{g}_d^e\right)`$
where the tilde quantities are the SME matter sector parameters described in section 2.2. The quantities $`\beta _w,\delta _w,\kappa _w,\gamma _w,\lambda _w`$ depend on the nuclear and electronic structure, and are given in table II of Bluhm . The $`s`$ coefficients result from the application of the Wigner-Eckhart theorem and are also given in Bluhm . All coefficients entering equation (15) are summarized in table 8.
From equation (15) and table 8 we notice that all $`m_F0`$ Zeeman transitions are sensitive to a violation of Lorentz symmetry, but not the $`m_F=0`$ clock transition. So in principle a direct measurement of one of the Zeeman transitions with respect to the clock transition (used as the reference) can yield a test of Lorentz invariance. The sensitive axis of the experiment is defined by the direction of the quantization magnetic field used to separate the Zeeman substates (vertical in the case of FO2), hence the rotation of the earth provides a modulation of the Lorentz violating signal at sidereal and semi-sidereal frequencies, which could be searched for in the data.
However, in such a direct measurement the first order Zeeman shift of the $`m_F0`$ transition would be the dominant error source and largely degrade the sensitivity of the experiment. The complete frequency shift of a Cs hyperfine Zeeman transition is VanAud
$$\delta \omega =\delta \omega _{SME}+m_FK_Z^{(1)}B+\left(1\frac{m_F^2}{16}\right)K_Z^{(2)}B^2+\mathrm{\Delta }$$
(16)
where $`\delta \omega _{SME}`$ is the SME frequency shift given by (15), $`B`$ is the magnetic field seen by the atom, $`K_Z^{(1)}=44.035`$ rad s<sup>-1</sup> nT<sup>-1</sup> is the first order Zeeman coefficient, $`K_Z^{(2)}=2685.75`$ rad s<sup>-1</sup> T<sup>-2</sup> is the second order coefficient, and $`\mathrm{\Delta }`$ is the shift due to other systematic effects. In (16) the diurnal and semi-diurnal variations of $`B`$ would mimic a putative Lorentz violating signal appearing in the sidereal and semi-sidereal variations of $`\delta \omega _{SME}`$ and render such a measurement very uncertain.
A somewhat cleverer strategy is to take advantage of the linear dependence on $`m_F`$ of the first order Zeeman shift but quadratic dependence on $`m_F`$ of one of the SME terms (the $`s_2^p`$ term in (15)). That implies that when measuring โsimultaneouslyโ the $`m_F=3`$, $`m_F=3`$, and $`m_F=0`$ transitions and forming the observable $`(\omega _{+3}+\omega _32\omega _0)`$ one should obtain a quantity that is independent of the first order Zeeman shift, but still shows a deviation from zero and a sidereal and semi-sidereal modulation in the presence of Lorentz violation. Using (15) and (16) this observable is
$$(\omega _{+3}+\omega _32\omega _0)=\frac{1}{7}K_p\stackrel{~}{c}_q^p+K_{Z(obs)}^{(2)}B^2+\mathrm{\Delta }_{(obs)}$$
(17)
where $`K_{Z(obs)}^{(2)}`$ and $`\mathrm{\Delta }_{obs}`$ are now the second order Zeeman coefficient and correction from other systematic effects for the complete observable.
The first term of (17) characterizes a possible Lorentz violation in the SME and is time varying when transforming the lab frame parameter $`\stackrel{~}{c}_q^p`$ to the conventional sun-centered frame. The general form of that transformation yields Bluhm
$$\stackrel{~}{c}_q^p=\stackrel{~}{B}+\stackrel{~}{C}_\omega _{}\mathrm{cos}(\omega _{}t)+\stackrel{~}{S}_\omega _{}\mathrm{sin}(\omega _{}t)+\stackrel{~}{C}_{2\omega _{}}\mathrm{cos}(2\omega _{}t)+\stackrel{~}{S}_{2\omega _{}}\mathrm{sin}(2\omega _{}t)$$
(18)
where $`\omega _{}`$ is the frequency of rotation of the Earth. The coefficients $`\stackrel{~}{B}`$, $`\stackrel{~}{C}_\omega _{}`$, $`\stackrel{~}{S}_\omega _{}`$, $`\stackrel{~}{C}_{2\omega _{}}`$, $`\stackrel{~}{S}_{2\omega _{}}`$ are functions of the 8 constant sun frame SME parameters $`\stackrel{~}{c}_X^p`$, $`\stackrel{~}{c}_Y^p`$, $`\stackrel{~}{c}_Z^p`$, $`\stackrel{~}{c}_Q^p`$, $`\stackrel{~}{c}_{}^p`$, $`\stackrel{~}{c}_{TX}^p`$, $`\stackrel{~}{c}_{TY}^p`$, $`\stackrel{~}{c}_{TZ}^p`$ (see Bluhm for details) with the three $`\stackrel{~}{c}_{TJ}^p`$ components suppressed by a factor $`v_R/c10^6`$ related to the velocity $`v_R`$ of the lab due to the rotation of the Earth.
The observable we use (equation (17)) should be independent of any long term ($`>`$ few seconds) variations of the first order Zeeman effect and therefore any sidereal or semi-sidereal variation of the observable would be the result of Lorentz violation, if it exceeds the measurement noise and the limits imposed by other systematic effects (see below).
The FO2 setup is sketched in Fig.5. Cs atoms effusing from an oven are slowed using a counter propagating laser beam and captured in a lin $``$ lin optical molasses. Atoms are cooled by six laser beams supplied by preadjusted fiber couplers precisely attached to the vacuum tank and aligned along the axes of a 3 dimensional coordinate system, where the (111) direction is vertical. Compared to typical clock operation BizeJPB , the number of atoms loaded in the optical molasses has been reduced to $`2\times 10^7`$ atoms captured in 30 ms.
Atoms are launched upwards at 3.94 m.s<sup>-1</sup> by using a moving optical molasses and cooled to $`1\mu `$K in the moving frame by adiabatically decreasing the laser intensity and increasing the laser detuning. Atoms are then selected by means of a microwave excitation in the selection cavity performed in a bias magnetic field of $`20`$ $`\mu `$T, and of a push laser beam. Any of the $`|F=3,m_F`$ states can be prepared with a high degree of purity (few $`10^3`$) by tuning the selection microwave frequency. 52 cm above the capture zone, a cylindrical copper cavity (TE<sub>011</sub> mode) is used to probe the $`|F=3,m_F|F=4,m_F`$ hyperfine transition at 9.2 GHz. The Ramsey interrogation method is performed by letting the atomic cloud interact with the microwave field a first time on the way up and a second time on the way down. After the interrogation, the populations $`N_{F=4}`$ and $`N_{F=3}`$ of the two hyperfine levels are measured by laser induced fluorescence, leading to a determination of the transition probability $`P=N_{F=3}/(N_{F=3}+N_{F=4})`$ which is insensitive to atom number fluctuations. One complete fountain cycle from capture to detection lasts 1045 ms in the present experiment. From the transition probability, measured on both sides of the central Ramsey fringe, we compute an error signal to lock the microwave interrogation frequency to the atomic transition using a digital servo loop. The frequency corrections are applied to a computer controlled high resolution DDS synthesizer in the microwave generator. These corrections are used to measure the atomic transition frequency with respect to the local reference signal used to synthesize the microwave frequency.
The homogeneity and the stability of the magnetic field in the interrogation region is a crucial point for the experiment. A magnetic field of $`200`$ nT is produced by a main solenoid (length 815 mm, diameter 220 mm) and a set of 4 compensation coils. These coils are surrounded by a first layer of 3 cylindrical magnetic shields. A second layer is composed of 2 magnetic shields surrounding the entire experiment (optical molasses and detection zone included). Between the two layers, the magnetic field fluctuations are sensed with a flux-gate magnetometer and stabilized by acting on 4 hexagonal coils. The magnetic field in the interrogation region is probed using the $`|F=3,m_F=1|F=4,m_F=1`$ atomic transition with a sensitivity of $`7.0084`$ Hz.nT<sup>-1</sup>. Measurements of the transition frequency as a function of the launch height show a peak to peak spatial variation of $`230`$ pT over a range of 320 mm above the interrogation cavity. Measurements of the same transition as a function of time at the launch height of 791 mm show a magnetic field instability near 2 pT at $`\tau =`$1 s as indicated in figure 6. The long term behavior exhibits residual variations of the magnetic field induced by temperature fluctuations which could cause variations of the current flowing through solenoid, of the solenoid geometry, of residual thermoelectric currents, of the magnetic shield permeability, etc.
The experimental sequence is tailored to circumvent the limitation that the long term magnetic field fluctuations could cause. First $`|F=3,m_F=3`$ atoms are selected and the $`|F=3,m_F=3|F=4,m_F=3`$ transition is probed at half maximum on the red side of the resonance (0.528 Hz below the resonance center). The next fountain cycle, $`|F=3,m_F=+3`$ atoms are selected and the $`|F=3,m_F=+3|F=4,m_F=+3`$ transition is also probed at half maximum on the red side of the resonance. The third fountain cycle, $`|F=3,m_F=3`$ atoms are selected and the $`|F=3,m_F=3|F=4,m_F=3`$ transition is probed at half maximum on the blue side of the resonance (0.528 Hz above the resonance center). The fourth fountain cycle, $`|F=3,m_F=+3`$ atoms are selected and the $`|F=3,m_F=+3|F=4,m_F=+3`$ transition is probed on the blue side of the resonance. This 4180 ms long sequence is repeated so as to implement two interleaved digital servo loops finding the line centers of both the $`|F=3,m_F=3|F=4,m_F=3`$ and the $`|F=3,m_F=+3|F=4,m_F=+3`$ transitions. With this method, magnetic field fluctuations over timescales longer than 4 s are filtered in the comparison between the two transition frequencies. Every 400 fountain cycles, the above sequence is interrupted and the regular clock transition $`|F=3,m_F=0|F=4,m_F=0`$ is measured for 10 s allowing for an absolute calibration of the local frequency reference with a suitable statistical uncertainty. The overall statistical uncertainty of the experiment is dominated by the short term ($`\tau 4`$ s) magnetic field fluctuations (fig. 6).
We have taken data implementing the experimental sequence described above over a period of 21 days starting on march 30, 2005. The complete raw data (no post-treatment) is shown in figure 7, each point representing a $``$432 s measurement sequence of $`\omega _{+3}+\omega _32\omega _0`$ as described above. Figure 8 shows the frequency stability of the last continuous stretch of data ($``$10 days). We note the essentially white noise behavior of the data on figure 8, indicating that the experimental sequence successfully rejects all long term variations of the magnetic field or of other perturbing effects.
According to equation (16) the frequency of the observable should be the sum of the putative Lorentz violating signal and of the second order Zeeman and other possible systematic corrections. Figure 7 shows a clear offset of the data from zero, which, using a least squares fit, is found to be $`(5.5\pm 0.1)`$ mHz with a very slight linear drift of $`(1.8\pm 1.0)\times 10^7`$ mHz s<sup>-1</sup>.
For our magnetic field of 202.65 nT the second order Zeeman correction of the $`\omega _{+3}+\omega _32\omega _0`$ observable is $`2.0`$ mHz. This only partly explains the offset observed in the data. The remaining part is most likely due to the differential influence of the magnetic field on the $`m_F=\pm 3`$ transitions, resulting from slightly different trajectories of the atoms in the different $`m_F`$ states and magnetic field inhomogeneities (residual first order Zeeman effect). Such differences in the trajectories could be due to differences in the trapping and/or launching of the atoms, related to the slightly different response of the Zeeman substates to the trapping fields. To check this hypothesis we have looked at the time of flight (TOF) of the atoms as a function of $`m_F`$. An offset of $`150\mu `$s between the $`m_F=+3`$ and $`m_F=3`$ TOF is observed. We are presently studying this effect in more detail (Monte Carlo simulations using the magnetic field map, tests with $`m_F=\pm 1`$ and $`m_F=\pm 2`$ states, longer term observation of the TOF difference and its variation, etc.) in order to be able to completely characterize its influence on the offset in figure 7, and its variation at sidereal and semi-sidereal frequencies.
In this paper we provide, as a preliminary results, only the values and statistical uncertainties of the coefficients $`C_\omega _{}`$, $`S_\omega _{}`$, $`C_{2\omega _{}}`$, and $`S_{2\omega _{}}`$ obtained from a model of the form
$`{\displaystyle \frac{1}{2\pi }}(\omega _{+3}+\omega _32\omega _0)=At+B`$ $`+`$ $`C_\omega _{}\mathrm{cos}(\omega _{}t)+S_\omega _{}\mathrm{sin}(\omega _{}t)`$
$`+`$ $`C_{2\omega _{}}\mathrm{cos}(2\omega _{}t)+S_{2\omega _{}}\mathrm{sin}(2\omega _{}t),`$
and the corresponding order of magnitude limits we expect for the $`\stackrel{~}{c}^p`$ parameters (cf. equations (17), (18)) of the SME.
Figure 9 shows the amplitudes $`A_\omega =\sqrt{C_\omega ^2+S_\omega ^2}`$ of least squares fits for a range of frequencies including the two frequencies of interest. We note no particularly significant peak at any frequency, and even less so at the frequencies of interest. A least squares fit at those frequencies yields the results shown in table 9. The correlation coefficients between any two of the four parameters in table 9 do not exceed 0.07.
From equations (17), (18) and table I of Bluhm we deduce orders of magnitude for the limits on the $`\stackrel{~}{c}^p`$ parameters of the SME (see table 3). We expect to obtain limits on two combinations of the five parameters $`\stackrel{~}{c}_X^p`$, $`\stackrel{~}{c}_Y^p`$, $`\stackrel{~}{c}_Z^p`$, $`\stackrel{~}{c}_Q^p`$, $`\stackrel{~}{c}_{}^p`$ at a level of $`10^{25}`$ GeV, and two combinations of the three parameters $`\stackrel{~}{c}_{TX}^p`$, $`\stackrel{~}{c}_{TY}^p`$, $`\stackrel{~}{c}_{TZ}^p`$ at a level of $`10^{19}`$ GeV.
In summary, we have carried out an experiment using Zeeman transitions in a cold atom <sup>133</sup>Cs fountain clock to test Lorentz invariance in the framework of the matter sector of the SME. In this paper we give a detailed description of the experiment and the theoretical model, we show our data and statistics, and we discuss our still ongoing investigation of systematic effects. Pending the outcome of that investigation and a more detailed theoretical analysis of our experimental results (explicit transformation of $`\stackrel{~}{c}_q^p`$ for our case), we provide only first estimates of the limits that our experiment can set on linear combinations of 8 SME matter sector parameters for the proton. These limits would correspond to first ever measurements of some parameters, and improvements by 11 and 14 orders of magnitude on others. A complete analysis (including systematics) of our experiment with final results for the SME parameters and their uncertainties will be the subject of a near future publication.
## 5 Conclusion
One hundred years after the publication of Einsteinโs original paper Einstein1905 special relativity, and its fundamental postulate of Lorentz invariance (LLI) are still as โhealthyโ as in their first years, in spite of theoretical work (unification theories) that hint towards a violation of LLI, and tremendous experimental efforts to find such a violation. Our experiments over the last years have provided some of the most stringent tests of LLI WP ; Wolf2003 ; WolfGRG ; Wolf2004 , but have nonetheless only joined the growing number of experiments in scientific history that measure zero deviation from LLI, albeit with an ever decreasing uncertainty. In spite of that, experimental tests of LLI are continuing along two lines: decrease of the uncertainties (see for example the contributions on rotating Michelson-Morley experiments in this volume) on one hand, and new types of experiments, e.g. the atomic clock test reported here, on the other.
In this paper we have presented a review of our recent Michelson-Morley and Kennedy-Thorndike experiment (section 3), and reported first results of our ongoing experiment that tests Lorentz invariance in the matter sector using a cold Cs atomic fountain clock (section 4). We have briefly described the two theoretical frameworks used to model and analyze our experiments (the Robertson-Mansouri-Sexl (RMS) framework and the standard model extension (SME)), and derived experimental limits on a number of parameters of those frameworks. When compared to other experiments those limits are the most stringent at present for several parameters (see tables 1, 2, 3).
The next generation of Michelson-Morley experiments are based on similar technology as our experiment (section 3) or the equivalent approach at optical frequencies Muller , but take advantage of active rotation of the experiment (see the corresponding contributions in this volume). Rotation of the experiment (typically at about 0.1 Hz) allows much faster data integration and places the signal modulation frequency close to the optimum where resonators are the most stable. It is expected that such experiments will lead to order(s) of magnitude improvements on orientation dependent parameters in the theoretical frameworks, but they present no advantage for only velocity dependent parameters. For example, in the RMS rotating experiments are likely to provide new, more stringent limits for the Michelson-Morley parameter ($`P_{MM}=1/2\beta _{\mathrm{MS}}+\delta _{\mathrm{MS}}`$) but no improvements on the Kennedy-Thorndike one ($`P_{KT}=\beta _{\mathrm{MS}}\alpha _{\mathrm{MS}}1`$). So we expect our (and other) present limits on $`P_{MM}`$ to be significantly improved, but we see no obvious way of improving on our present limit on $`P_{KT}`$ in the near future.
Several improvements of our clock test of LLI in the SME matter sector (section 4) are possible. For example, using the unique capability of our double fountain (FO2) to run on both, Cs and Rb, we expect to be able to use Cs as the SME sensitive species and Rb (which is less sensitive to the SME Bluhm ) as the magnetic field probe. In that way we should be able to perform magnetic field independent measurements that could improve on our present results, and allow access to other SME parameters that we are insensitive to with our present set up. Also, rotation of the experiment could provide a method for faster modulation of the signal but is unpractical in an Earth bound laboratory. However, space missions with onboard atomic clocks are well suited for such a test. In particular the European ACES (Atomic Clock Ensemble in Space) mission ACES , scheduled for flight on the international space station (ISS) in 2009, seems very promising in this respect. It will include a laser cooled Cs clock (PHARAO) with expected performance at least equivalent to our FO2, but with the orientation of its quantization field axis modulated at a 90 min period (ISS orbital period) rather than 24 hr as in our case. This should allow for much faster data integration and significant improvement on the limits presented here.
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# AGN outflows in a Cosmological Context
## 1 INTRODUCTION
Outflows from active galactic nuclei (AGN) potentially play a very significant role in the evolution of large-scale structure. Such outflows consist of hot, tenuous out-flowing gas detected in absorption, and the powerful radio jets detected in some quasars. Blue-shifted absorption lines (relative to systematic velocity of the AGN), such as Ly$`\alpha `$ and the C iv and N v doublets, indicate gas moving away from the central source at high velocities (Ulrich 1988; Crenshaw et al. 1999). Broad absorption lines (BAL) indicating even higher velocities, are detected in luminous AGN (e.g., Hewett & Foltz 2003). Radio-loud quasars (RLQ) also carry substantial amounts of energy into the intergalactic medium (IGM) via collimated jets of relativistic plasma (e.g., Begelman, Blandford, & Rees 1984). AGN outflows may play a role in distributing magnetic fields into the IGM (e.g., Furlanetto & Loeb 2001), and they can impact the evolution of their host galaxies by, for example, regulating the growth of supermassive black holes (Wyithe & Loeb 2003). BAL outflows and RLQ are also possible mechanisms for heating the intra-cluster medium (e.g., Valageas & Silk 1999; Nath & Roychowdhury 2002). The magnitude of the influence of outflows on the IGM depends on a few key properties of AGN, such as the relationship between kinetic and bolometric luminosity, that still need to be constrained by observation.
In RLQ, the kinetic luminosity of the jet is thought to be correlated with the bolometric luminosity (Willott et al. 1999). Direct estimates of the kinetic luminosity in BAL outflows can be made if such quantities as the covering fraction of the outflows, the outflow velocity, the radius of the outflow, and the column density are known. The velocities are accurately obtained from the absorption lines, and the covering fractions can be estimated with statistical arguments (Weymann 1997). The radius of the outflow can be obtained through observations of the photoionizing flux of the central source (Krolik 1999). Estimates of the column densities can be made by studying the UV and X-ray absorption lines (Gallagher et al. 1999, 2001). Uncertainties in the above quantities translate to uncertainties in the kinetic luminosity, and therefore into uncertainties in the energy of AGN outflows. More energetic outflows can fill a larger volume, having a greater impact on large-scale structure.
In this paper, we model the growth of AGN outflows using energy conservation arguments and simple approximations about their geometry. Through numerical simulations we model the distribution of these outflows and hence we can estimate the degree to which they affect the universe on global scales, or the fractional volume occupied by AGN outflows. In ยง2, we describe some of the details of the cosmological simulation and the luminosity function that we have combined to obtain an AGN distribution. In ยง3 we explain the assumptions surrounding our physical model of the growth of heated bubbles around AGN as well as our selection of kinetic luminosity. In ยง4, we show the advantage of using a simulated density distribution over using simple statistical methods for the distribution of AGN. In ยง5 we examine the effects of simulation size and resolution, AGN lifetime, and AGN bias on the cumulative effect of outflows, and we conclude in ยง6. For the rest of this paper, we assume $`\mathrm{\Omega }_m=.27`$, $`\mathrm{\Omega }_\mathrm{\Lambda }=0.73`$, and $`\mathrm{\Omega }_b=.04`$ with $`\mathrm{\Omega }_bh^2=0.02`$, consistent with the WMAP data.
## 2 SIMULATING THE AGN ENVIRONMENT AND DISTRIBUTION
In order to understand the influence of AGN outflows on a global scale, it it useful to model the outflows in the context of large-scale structures. We assume here that AGN trace high density regions, and so use a gas density distribution to bias AGN in our study. Other studies have, for example, approximated the distribution of AGN with statistical formulae. Tegmark, Silk, & Evrard (1993) and Furlanetto & Loeb (2001; hereafter, F & L) both use Poisson distributions to model the spatial distribution of galaxies for simplicity, in order to obtain filling fractions of supernova-driven winds and of AGN outflows respectively. We combine a z-dependent luminosity function with a simulated gas density distribution, allowing us to consider the AGN in their appropriate environments rather than homogeneously distributing them throughout the universe. In ยง4, we compare a Poisson distribution of AGN outflows with our model.
### 2.1 Cosmological Density Distribution
In order to model the evolution of the gas density distribution in the universe, we use a standard Particle-Mesh code to simulate the distribution of the dark matter, and we assume that on the scales we are considering, the gas distribution follows that of the dark matter, which is supported by cosmological gas dynamics simulations (Gnedin 2000; Chiu & Ostriker 2000; Miller & Ostriker 2001; Somerville 2002; Tassis et al. 2003; Benson & Madau 2003; Susa & Umemura 2004; Shapiro, Iliev, & Raga 2004; Mo & Mao 2004). We also assume that the temperature of the cosmic gas is constant at 15,000 K. This assumption is, clearly, an oversimplification, since the temperature of the cosmic gas is known to evolve with time (Ricotti, Gnedin, & Shull 2000; Schaye et al. 2000; McDonald et al. 2001; Kim, Cristiani, & DโOdorico 2002; Theuns et al. 2002; Hui & Haiman 2003) and vary in space. However, since the typical sizes of AGN-driven bubbles are significantly larger than the scales over which the temperature of cosmic gas changes, this approximation is sufficient for our purposes.
### 2.2 Luminosity Function
The above simulation produces a gas density distribution at each time step, from $`z19`$ to $`z=0`$. In each step, we use a luminosity function to populate the simulation with AGN over a range of luminosities. In order to place AGN into our simulation, we must implement an AGN luminosity function that applies to high redshifts. The scarcity of high-z galaxy detections in surveys makes luminosity function predictions difficult, although the ability to make detections is improving. We use a model by Schirber & Bullock (2003) for the QSO luminosity function. The model satisfies all existing constraints from optical surveys, such as the Sloan Digital Sky Survey (SDSS) data and the Great Observatories Origins Deep Survey (GOODS) data for $`z>3`$ (Fan et al. 2001a, 2001b; Cristiani et al. 2004) and Two Degree Field (2dF) data for $`z<2.3`$ (Boyle et al. 2000). It should be noted that optical surveys perhaps underestimate the faint-end, low-z luminosity function. Hard X-ray detections of AGN yield a higher number of faint AGN at low-redshift than in previous optical surveys (Barger et al. 2005). The model that we use for now parameterizes the following luminosity function:
$$\varphi (L_B,z)=\frac{\varphi _{}/L_{}}{(L_B/L_{})^{\gamma _f}+(L_B/L_{})^{\gamma _b}}.$$
(1)
In the above equation, $`\varphi _{}`$ is the average comoving number density of AGN, $`L_{}`$ and $`L_B`$ are the break luminosity and AGN B-band luminosities respectively (given in units of $`L_{\mathrm{},B}`$ where we have followed S & B in using $`2.11\times 10^{33}\text{ ergs s}^1`$ for the B-band luminosity of the sun), and $`\gamma _f`$ and $`\gamma _b`$ are the faint and bright-end slopes respectively. As in S & B, we interpolate the SDSS and 2dF fits for $`2.3<z<3`$. The parameterization is given in Table 1, and the luminosity function for several different redshifts is plotted in Figure 1. Note that for $`z>3`$, $`L_{}`$ and $`\varphi _{}`$ depend on the weighted emissivity of AGN (accounting for contributions to the ionizing rate from other sources, such as stars). The weighted emissivity is given by $`\mathrm{log}_{10}\widehat{\epsilon }^Q=0.245z+0.596`$, a simple power law fit to the values given in S & B Table 1. We have chosen a minimum AGN luminosity of $`10^8L_{\mathrm{},B}`$ as in S & B, who argue that this represents the faintest Seyfert galaxies found. We extrapolate the high-z parameterization of Table 1 to $`z19`$, which introduces a significant uncertainty in the abundance of AGN at high redshift. However, as we show below, AGN at $`z<3`$ dominate, so this uncertainty affects our results insignificantly. The exact redshift range for each run depends on the size of the simulation box since larger boxes can sample an earlier, rarer population of AGN.
It should be noted that we used S & Bโs โModel Aโ to determine this particular luminosity function for $`z3`$. The model uses the emissivity of the ionizing background, assuming a stellar as well as an AGN contribution, to constrain the AGN luminosity function in combination with the SDSS data. The model implements the ionizing rates of McDonald & Miralda-Escudรฉ (2001) and assumes an escape fraction of $`0.16`$ (including the stellar contribution to the ionizing background). The model neglects optically obscured AGN on the basis that they do not make a significant contribution to the ionizing background. This means that the luminosity function we use could be excluding some AGN that host outflows. Of the models described by S & B, we have chosen โModel Aโ because it is most consistent (of the models studied by S & B) with the faint end luminosity function (for $`z>4`$), from GOODS (Cristiani et al. 2004).
We use the luminosity function to calculate the number of AGN at a given redshift and luminosity in each simulation box, and we then place each one into a random location, but with a bias toward high density regions. In order to distribute AGN so that they correspond to regions of high density, we calculate, for each cell, the probability of hosting an AGN given by
$$P(i,j,k)=\frac{\rho _m^\alpha (i,j,k)\mathrm{\Delta }V}{\underset{i,j,k=0}{\overset{N1}{}}(\rho _m^\alpha \mathrm{\Delta }V)}.$$
(2)
Above, $`\rho _m`$ is the matter density (in units of average baryon density) in each cell specified by coordinates $`i`$,$`j`$,$`k`$; $`N`$ is the size of simulation box; $`\mathrm{\Delta }V`$ is the comoving volume of each cell; and $`\alpha `$ is a bias parameter ensuring that regions of higher density are more likely to host AGN. The above probability function is independent of the characteristics of individual AGN (such as luminosity). Analysis of existing observations suggests that AGN are more biased toward high density regions at increasing redshifts (e.g., Croom et al. 2004; Porciani, Magliocchetti, & Norberg 2004). In ยง5.3, we will examine the effects of different bias parameters on our results. In order to place the AGN in specific cells, we randomly choose a cell and compare the value of $`P(i,j,k)`$ in that cell to a randomly generated number between zero and one. If $`P(i,j,k)`$ is greater than the random number, an AGN goes into the cell. Otherwise, the process repeats until the AGN has been placed in the simulation. This method has the effect of biasing AGN toward regions of higher density, while fully allowing for Poisson noise in their distribution.
## 3 AGN OUTFLOWS
The following sections detail the distribution and kinematics of AGN outflows in our model. Outflows are not observed in all AGN, and so ยง3.1 deals with our method of selecting AGN to host outflows. In ยง3.2, we describe our assumptions about the expansion and evolution of individual outflows into the IGM, and in ยง3.3, we discuss the effects of kinetic luminosity on the outflows.
### 3.1 Distribution of Outflows
In our model we assume that the AGN produce outflows that are responsible for distributing tenuous, hot gas into the IGM. The outflows may also be a mechanism for depositing metals and magnetic fields into the IGM. However, outflows are only associated with a fraction of AGN, which we must reflect in our model. Detection of blue-shifted broad absorption lines in an AGNโs spectrum indicates out-flowing gas from the nucleus. Until fairly recently, observations indicated that these broad absorption lines were limited to radio-quiet quasars (Stocke et al. 1992), and only seen in $`10\%`$ of them. However, recent detections of BAL outflows in RLQ (e.g., Brotherton et al. 1998; Menou et al. 2001) suggest that BAL outflows do not necessarily follow this radio dichotomy. There is, however, evidence for some luminosity dependence in the occurrence of BAL outflows. Crenshaw et al. (1999) examined the UV spectra of several Type I Seyfert galaxies obtained with the Hubble Space Telescope, and found intrinsic, narrow absorption lines in more than half of their sample. The absorption spectra of these low-luminosity objects show similarities to the BAL features in high-luminosity objects, suggesting a possible relationship between the two. On the higher luminosity end, Hewett & Foltz (2003) apply a magnitude correction to the Large Bright Quasar Survey (LBQS; Hewett, Foltz, & Chaffee 1995, 2001) and find a higher percentage of BAL quasars than previously determined. Magnitude and flux-limited samples can exclude BAL quasars in which the absorbing gas reduces the spectral energy distribution of the quasars in the wavelength range of selection. After applying a correction for this effect, Hewett & Foltz find that $`20\%`$ of the AGN in the sample host outflows.
In our model, we interpolate the above low and high-luminosity limits for the fraction of AGN hosting outflows, $`f_{out}`$, and combine this fraction with the luminosity function of ยง2.2. Our formulation is as follows:
$$f_{out}=\{\begin{array}{cc}.5\text{for }\mathrm{log}L_B<10,\hfill & \\ 1.50.1\mathrm{log}L_B\text{for }10\mathrm{log}L_B13,\hfill & \\ .2\text{for }\mathrm{log}L_B>13,\hfill & \end{array}$$
(3)
where the lower and higher-end luminosities are consistent with the fractions determined above. It is not well known whether the above statistics also apply to the population of X-ray detected AGN without optical counterparts. However, because X-ray surveys detect more low-luminosity AGN than optical surveys, it is possible that there are more AGN containing outflows than we here assume. Additionally, it should be noted that because of the range of covering fractions of AGN outflows, the fraction of AGN hosting outflows could be larger than what is observed, even in optical surveys, due to orientation effects (Morris 1988; Weymann et al. 1991; Hamann, Korista, & Morris 1993). Therefore, our assumptions should provide a conservative lower limit to the number of AGN hosting outflows.
### 3.2 Evolution of Outflows
In order to understand the degree to which active galaxies affect the IGM, it is important to have an accurate model of the expansion of hot, ionized gas into the IGM. The environment of the outflow must be considered in order to model the expansion. Using the density profile described in ยง2, we can, to a degree, reproduce the environment surrounding each of the individual active galaxies that we placed in our simulation in accordance with the luminosity function (eq. 1) and $`f_{out}`$ (eq. 3).
The lifetime of the active galactic nucleus is short compared to the expansion time of the bubbles, which lasts over the duration of our simulation. Therefore, we consider the active phase as a brief energy injection, followed by an adiabatic blast wave gathering up material in the IGM, analogous to the adiabatic phase of a supernova remnant. As in Scannapieco & Oh (2004), we have used the familiar Sedov-Taylor blast wave model to determine the size of the bubbles in our analysis. Their Equation 10, which we adopt for convenience of units, follows:
$$R_s=(1.7\text{ Mpc})\left(\frac{E_k}{10^{60}\text{ ergs}}\right)^{1/5}(1+\delta _m)^{1/5}(1+z)^{3/5}\left(\frac{t_{age}}{10^9\text{ yr}}\right)^{2/5}.$$
(4)
In the above equation, $`E_k`$ is the kinetic energy injected by the AGN, $`\delta _m`$ is the overdensity, and $`t_{age}`$ is the time since the active phase began. The overdensity is an average over all cells within the radius of the bubble, obtained using an iterative technique. The above Sedov-Taylor solution, with a constant density, is clearly not valid as the outflow escapes its host galaxy and travels first through a region with a steeply-falling off density profile, e.g. the NFW profile (Navarro, Frenk, & White 1997). However, our simulation does not resolve this stage of the expansion. The virial radius of a typical $`L_{}`$ galaxy is well below the resolution of our simulation. Without modeling the expansion of the outflows inside this region, there is uncertainty in the speed with which the outflows escape their host galaxies. However, this uncertainty is hidden by the parameter $`E_k`$, which describes the actual energy input from the AGN.
We assume that the bubbles expand according to the above equation until they reach a pressure equilibrium with their environment. If pressure equilibrium is reached before the energy injection has stopped, the growth of the bubbles is determined by the surrounding pressure and the injected energy, and the size is given by
$$R_P=(3.24\times 10^{25}\text{ Mpc})\left(\frac{3E_k}{4\pi P}\right)^{1/3}.$$
(5)
The pressure of the surrounding medium is given by $`P=(1+\delta _m)\overline{n}_bk_bT`$, where the quantity $`(1+\delta _m)\overline{n}_b`$ is the average gas density inside the bubble as determined above, and $`T=1.5\times 10^4\text{ K}`$ (the average temperature of the IGM). The kinetic energy in Equations 4 and 5 is given by $`E_k=\epsilon _{kB}L_Bt_{age}`$ where $`t_{age}`$ is the age of the AGN during the active phase, or the lifetime of the AGN once the active phase has ended. We assume a constant lifetime of $`10^8\text{ yrs}`$ for all AGN until ยง5.2 in which we will examine other lifetimes. The parameter $`\epsilon _{kB}`$ is given by the ratio of the kinetic luminosity of the outflow to the AGN B-band luminosity ($`L_k/L_B`$). Each AGN injects kinetic energy into the IGM at a rate ($`L_k`$) correlated with the AGNโs luminosity. We assume that the AGN are accreting at roughly their Eddington rates, so that $`L_{edd}L_{bol}`$. As in F & L, we adopt $`L_{bol}10L_B`$, consistent with Elvis et al. (1994). We describe our choice of the kinetic fraction, $`L_k/L_{bol}`$ or $`\epsilon _k`$, in the next section.
After the energy injection phase, bubbles in pressure equilibrium no longer expand as a result of the energy injection of the AGN. Any subsequent evolution of the bubbles in our simulation is determined by the Hubble expansion and the evolution of the density distribution in the surrounding environment. Overlap of the bubbles is not likely to affect the expansion significantly. When an expanding bubble in pressure equilibrium overlaps with the interior of another bubble, encountering densities much lower than those typical in the IGM, the expansion speed does not change, as there is no longer a pressure gradient.
Our model does not include radiative cooling, as the cooling times for these bubbles are typically much longer than the timescales we consider. However, for bubbles located in cluster environments (higher densities and temperatures, etc.), radiative cooling, as well as other physical processes, may become important. Accurately modeling outflow physics in these complex environments requires the use of hydrodynamical simulations. Separate work is currently underway to understand the impact of AGN outflows in these environments (e.g., Brรผggen & Kaiser 2002; Ruszkowski, Brรผggen, & Begelman 2004).
Figure 2 shows the evolution of the bubble size for two typical AGN (each with luminosities of $`10^9L_{\mathrm{},B}`$) residing in different environments within the simulation. The figure shows that the AGN residing in the lower density environment produces a larger bubble ($`4.8h^1\text{ Mpc}`$) than the AGN in the higher density environment ($`2.2h^1\text{ Mpc}`$).
In both F & L and Nath and Roychowdhury (2002), AGN outflows are treated as collimated jets that spread out into a cocoon after reaching pressure equilibrium at the end of the AGNโs active phase. Furthermore, both BAL AGN and RLQ are treated similarly, justified by the small covering fraction of BAL outflows, averaging at around $`10\%`$ (Weymann 1997). We likewise adopt the practice of treating the two different objects similarly, noting that there is additional incentive in the likelihood of overlap between RLQ and BAL AGN. Furthermore, because the time scales we consider are significantly longer than the AGN injection phase, we do not model collimated jets, but rather approximate the outflows as bubbles expanding into the IGM with spherical symmetry. However, as demonstrated by Figure 2 of F & L, if the energy injection is modeled assuming spherical symmetry during the active phase, rather than with collimated jets, the result is only a slightly smaller final comoving bubble size. The entirely spherical case produces a smaller cocoon because the surface area of the bubble causes it to decelerate sooner, whereas a collimated jet makes its way through the IGM more easily.
### 3.3 Kinetic Luminosity and Filling Fraction
Estimates of the kinetic fraction, $`\epsilon _k`$, are subject to a number of observational uncertainties surrounding BAL outflows. Observational constraints depend on quantities that are difficult to determine, such as the distance of outflows from the central source, and the covering fraction of the outflows (e.g., De Kool et al. 2001). We use our model of AGN outflows to calculate the filling fraction of outflows as a function of redshift, $`F(z)`$, for a box of length $`128h^1\text{ Mpc}`$ (with $`0.5h^1\text{ Mpc}`$ cells), assuming that all of the AGN are active for $`10^8\text{ yrs}`$, and that they follow a constant bias, $`\alpha =2`$. We then treat the kinetic fraction as a free parameter of our model. We start with $`\epsilon _{kB}=1`$, or a kinetic fraction $`\epsilon _k=0.1`$, as in F & L. Nath & Roychowdhury (2002) argue that $`\epsilon _k=0.1`$ is probably an upper limit for BAL outflows, assuming that the covering fraction of BAL outflows is $`10\%`$. Figure 3 shows the fractional volume filled with AGN outflows for decreasing kinetic fraction. We find that using a kinetic fraction of $`10\%`$, the entire simulation box is filled with outflows by $`z2`$. Similarly, Scannepieco & Oh (2004) find that a kinetic fraction of $`10\%`$ overestimates feedback effects in their model of AGN outflows.
Knowledge of the filling fraction of the Ly$`\alpha `$ forest at various redshifts can provide further constraints, if we assume that the AGN outflows consist of hot, tenuous gas that cannot occupy the same volume of space as Ly$`\alpha `$ absorbing regions. At low-z, simple conclusions drawn from observations of the Ly$`\alpha `$ forest constrain the fractional volume of voids to between $`70`$ and $`99.6\%`$ of the total volume, providing an upper limit to the filling fraction of AGN outflows (e.g., Penton, Stocke, & Shull 2004; Davรฉ et al. 1999). The vertical bar near $`z=0`$ in Figure 3 shows the constrained filling fraction of voids. An $`\epsilon _k`$ of less than $`10\%`$ produces filling factors that are consistent with the above values, but imposing more precise constraints from Ly$`\alpha `$ forest observations is difficult because of the range of column densities under consideration. For somewhat higher redshifts ($`1.7<z<3.8`$), Duncan, Ostriker, & Bajtlik (1989) have studied voids in the Ly$`\alpha `$ forest and determined that voids with sizes between $`10`$ and $`70h^1\text{ Mpc}`$ occupy $`<20\%`$ of the volume. The horizontal bar in Figure 3 shows this upper limit for the filling fraction of AGN outflows over the appropriate redshift range. We choose $`\epsilon _{kB}=.1`$ ($`\epsilon _k=1\%`$) as our fiducial value, as it seems to match observational and theoretical constraints more closely than the higher values.
## 4 COMPARISON WITH A POISSON DISTRIBUTION
As previously mentioned, it is possible to calculate the filling fraction of AGN outflows analytically by assuming that AGN are distributed according to Poisson statistics. In this section, we demonstrate the advantage of including a realistically evolving density distribution in our model.
We calculate the filling fraction of outflows using a Poisson distribution of sources and compare with our method of biasing AGN toward the high density regions within a cosmological simulation (see ยง2). We first calculate the porosity of AGN outflows at each redshift, given by:
$$Q(z)=\frac{4\pi }{3}_z^{\mathrm{}}\frac{dz^{}}{\tau _{AGN}}\frac{dt^{}}{dz^{}}_{L_{min}}^{L_{max}}R^3\varphi (L_B,z)f_{out}๐L_B,$$
(6)
where the above quantities are calculated in physical units. We then calculate the filling fraction assuming a Poisson distribution:
$$f(z)=1e^{Q(z)}.$$
(7)
In the above calculations, we have determined the volume of each outflow under pressure equilibrium conditions, and under the assumption that the energy injection is instantaneous. The pressure is determined by the average particle density at each redshift and a constant temperature of $`T=1.5\times 10^4\text{ K}`$. We compare the above filling fraction with that produced by our model, as described in the previous section, for a simulated volume of $`128^3h^3\text{ Mpc}^3`$ (with $`0.5h^1\text{ Mpc}`$ cell resolution) in Figure 4. The Figure shows that the Poisson distribution of sources produces a higher filling fraction than our model. The simulation provides a realistic environment for each AGN. The outflows of AGN residing in regions of higher density do not fill as large a volume because the IGM exerts more pressure on the bubbles than in a uniform, average density distribution. Also, the AGN themselves are not distributed over as large a volume of space, as they tend toward higher density regions. Therefore, their outflows do not fill as large a fraction of the simulation as if they were uniformly distributed.
## 5 RESULTS
We have seen in the previous two sections the general result of our model on the filling fraction of AGN outflows. In this section, we will examine the effects of varying simulation box size and resolution, AGN lifetime, and the distribution of AGN as determined by AGN bias. We will first conduct convergence studies, to determine the effects of box size and resolution on the rest of our analysis.
### 5.1 Convergence Studies
We ran our analysis on simulations of differing box sizes in order to choose the optimal box size for the rest of our studies. Larger boxes are more inclusive, but also much more computationally expensive, and so convergence is desirable. For each box size, we determine the fractional volume heated by AGN as a function of redshift using the model for bubble growth described in ยง3. We have evolved boxes of length $`64h^1\text{ Mpc}`$, $`128h^1\text{ Mpc}`$, and $`256h^1\text{ Mpc}`$ (each with $`1h^1\text{ Mpc}`$ cells) on a side. Figure 5 shows the filling fraction for each box size. The filling fraction does not vary dramatically between the two larger box sizes (for this reason, we did not complete the run for the largest, most expensive box size). It appears that we have reached convergence for the $`128h^1\text{ Mpc}`$ box. We therefore choose the $`128h^1\text{ Mpc}`$ box size for all of our parameter studies.
In addition to box size studies, we also examined the effects of simulation resolution on the AGN heated fractional volume. In simulations with finer resolution, individual cells can reach significantly higher densities (see Fig. 6), which directly affects the sizes of the bubbles in our simulations. The result is a different fractional volume affected by AGN depending on resolution. High resolution represents a more accurate picture of the simulated volume, but like large box size, is more computationally expensive because of the larger number of cells. Therefore, we increase the resolution in our simulation box until the filling fraction no longer depends on resolution, after which, there is no need for finer resolution. We studied the $`64h^1\text{ Mpc}`$ length box with resolutions of $`1h^1\text{ Mpc}`$ cells, $`0.5h^1\text{ Mpc}`$ cells, and finally $`0.25h^1\text{ Mpc}`$cells. We reach convergence in the filling fraction very quickly, as shown in Figure 7. The results for the $`0.5h^1\text{ Mpc}`$ cells and the $`0.25h^1\text{ Mpc}`$ cells are very similar, and so we choose the faster computation, $`0.5h^1\text{ Mpc}`$ cells, for our fiducial resolution.
### 5.2 AGN Lifetime
We have chosen for our fiducial model a constant lifetime for all AGN, $`\tau _{AGN}=10^8\text{ yrs}`$. Yu & Tremaine (2002) examine the dependence of lifetime on black hole mass and their results show modest variation in lifetime ($`30300\text{ Myr}`$) over the range of black hole masses of interest here. They determine AGN lifetime from a combination of the luminosity function and black hole number density, and their results are comparable to the associated Salpeter times, further evidence that black holes accrete most of their mass during their active phases. Our choice of $`\tau _{AGN}=10^8\text{ yrs}`$ is consistent with their results. We test the effect of using higher and lower lifetimes as well. We apply a constant lifetime of $`10^7\text{ yrs}`$ as our lower limit. According to S & B, $`10^7\text{ yrs}`$ is an approximate lower limit for AGN lifetime, obtained from arguments similar to those of Yu & Tremaine:
$$f_{on}(z)=\frac{\tau _{AGN}}{t_{Hubble}}\frac{\mathrm{\Phi }(>L_B,z)}{n(>M_{BH},z=0)}.$$
(8)
In the above equation (eq. 22 of S & B), the AGN lifetime is determined from the fraction of galaxies having active nuclei at a given redshift ($`f_{on}`$). This fraction can be determined by the ratio of the comoving number density of AGN (where $`\mathrm{\Phi }`$ is the number density of AGN with luminosities greater than $`L_B`$ at redshift $`z`$) to the black hole number density ($`n`$). S & B assume that the black hole number density only increases with time, so that eq. 8, with the number density evaluated at $`z=0`$, will give a minimum lifetime. For our upper limit, we use $`\tau _{AGN}=10^9\text{ yrs}`$, consistent with Croom et al. (2004) who determine an upper limit on AGN lifetime from arguments about the growth of dark matter halo mass. We do not examine redshift or mass dependence of $`\tau _{AGN}`$, because the variation is not significant over the range of lifetimes we consider.
We find that shorter lifetimes result in a lower filling fraction. In order to remain consistent with the luminosity function in the case of shorter lifetimes, many more AGN will become active at later redshifts than in the longer lifetime case. As AGN are born later in the simulation, they sit in regions of higher density than AGN born earlier, as the density distribution evolves with redshift. In this analysis, we have not included the possibility of recurrent activity associated with the same nucleus because the duty cycles of AGN are so low that it would introduce a small effect. Figure 8 shows $`F(z)`$ for lifetimes of $`10^7`$, $`10^8`$, and $`10^9\text{ yrs}`$. In order to test that our interpretation of the effect of evolving density distribution is correct, we calculate $`F(z)`$ for the three values of AGN lifetime under the assumption that the density distribution is uniform throughout the universe. Figure 9 shows the same range of lifetimes for the uniform density case. As expected, here we do not see the trend with lifetime because the density dependence has been removed. Therefore, because of the evolving density distribution, a shorter AGN lifetime will result in a lower filling fraction at $`z=0`$.
### 5.3 AGN Bias
Observations show that AGN are biased toward regions of high density, with this bias increasing toward higher redshifts. A simple, quantitative implementation of this observed bias is the relation $`n_{AGN}\rho _{m}^{}{}_{}{}^{\alpha }`$, where $`n_{AGN}`$ is the number density of AGN in units of average number density of AGN, and $`\alpha `$ is the linear bias parameter. Therefore, $`\alpha `$ is a measure of the correlation between matter density and AGN density. We examine simple constant bias models here as well as one redshift dependent model. We do not, however, examine the possibility of luminosity dependent bias here, nor do we distinguish between the bias of radio-loud and radio-quiet QSOs despite evidence that radio-loud sources are more strongly clustered. Therefore, our bias parameter, $`\alpha `$, should be considered as an average bias of the outflow-producing population of AGN.
Our fiducial model was run with a constant bias, $`\alpha =2`$, consistent with the average of many bias models. Because bias increases with redshift, we also ran a model with an increased constant bias of $`\alpha =3`$. With a larger bias, AGN will be more tightly clustered to higher density regions, preventing their outflows from growing as large, and resulting in more overlap between bubbles. The result is a lower filling fraction, as shown in Figure 10.
Since the value of the bias parameter is not constant, but more likely to be (at the very least) redshift dependent, we have also tested the following redshift dependent model from Croom et al. (2004):
$$\alpha (z)=0.53+0.289(1+z)^2.$$
(9)
The above is a simple model derived from 2dF data (up to $`z<2.48`$) combined with WMAP and 2dF cosmology. Rather than extrapolate this model over our entire redshift range, for $`z3`$ we use a constant bias of $`\alpha =5.154`$, determined by evaluating the above equation at $`z=3`$. The effects of redshift dependent bias are seen in Figure 10. At higher redshifts, AGN are more strongly biased toward high density regions than in either of our constant bias models, and so the filling fraction is significantly lower. At lower redshifts ($`z1`$), AGN are even less biased toward regions of high density than in our fiducial model, and so $`F(z)`$ increases approaching $`z=0`$.
## 6 DISCUSSION AND CONCLUSIONS
We have examined the filling fraction of AGN outflows in the context of large-scale cosmological simulations, and considered the influence of various observationally constrained parameters on the result. We find that the kinetic fraction of outflows need not be very high ($`10\%`$) for AGN outflows to fill the entire IGM by $`z2`$. Observations of gaps in the Ly$`\alpha `$ forest provide possible constraints on the filling fraction, that can in turn be used to place constraints on the kinetic luminosity. In our study we have used a luminosity function consistent with optical surveys to distribute outflows throughout our simulation, however, future studies will have to consider X-ray surveys, which predict more faint luminosity AGN.
Our model employs several simple approximations, but is nonetheless instructive. We have made the assumption that AGN outflows are spherical bubbles, propagating into the AGN adiabatically for a short while, until reaching pressure equilibrium with their environments. Spherical symmetry is not likely to remain intact once the outflow reaches far enough into the IGM. The outflows will expand away from large-scale structures, such as filaments, and into less dense regions. However, we can assume spherical symmetry as an average geometry, and the effects on the volume filling fraction are not likely to be very large. Our assumption that radiative cooling can be ignored is justified by the timescales involved. If we were to consider more complex environments in our model, such as those of clusters, we would need to include a great deal more physics (including cooling) requiring hydrodynamical simulations. While the actual outflow physics and geometry are likely to be more complex than we assume, our assumptions provide a good first-order approximation of the filling fraction of AGN.
The parameters we have studied, AGN lifetime, and AGN bias, contribute significantly to the AGN filling fraction as expected. We have assumed a particular evolution of the gas density distribution, in which the gas density profile is relatively uniform at high-z, and forms high density filaments at low-z. We have tested upper and lower limits for AGN lifetime, and find that shorter AGN lifetimes result in a lower filling fraction than longer lived AGN due to the evolving gas density distribution. We have examined the effects of different AGN biases on filling fraction as well. Larger bias results in an ultimately lower filling fraction. However, bias likely depends on factors such as redshift and luminosity, and the dependence of the evolution of the filling factor on bias is likely to be more complicated than in our simple scenario. Changes both in the bias, and in the AGN lifetime affect the filling fraction of outflows because of the importance of the underlying density distribution and its importance for determining AGN distribution and environments.
We are very grateful to Mitch Begelman for his contributions to this paper. Additionally, we thank Nahum Arav, Jack Gabel, Mateausz Ruszkowski, Mike Shull, and John Stocke for helpful discussions. This work was supported by NSF grant AST-0134373 and by the National Computational Science Alliance under grant AST-020018N, and utilized the SGI Origin 2000 array and IBM P690 array at the National Center for Supercomputing Applications.
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# X-ray Diagnostics of Grain Depletion in Matter Accreting onto T Tauri Stars
## 1 Introduction
X-rays, at least in principle, present a powerful method of studying directly the accretion processes of star formation. In the case of classical T Tauri stars (CTTS), the magnetospheric accretion paradigm posits that material is funnelled by magnetic fields onto the star from a truncated disk (e.g. Uchida & Shibata, 1984; Bertout et al., 1988), rather than through direct surface interaction with the disk itself. Matter impacts the stellar surface at free-fall velocities of up to a few hundred km s<sup>-1</sup>. The accreting material is then expected to form a shock at the stellar surface. The resulting shocked plasma temperature is of order a few million K, and it will therefore radiate predominantly in X-rays. For the ranges of mass accretion rates inferred for CTTS of $`10^6`$-$`10^{10}M_{}`$ yr<sup>-1</sup> (e.g. Hartigan et al., 1995; Gullbring et al., 1998; Johns-Krull et al., 2000), the X-ray luminosities resulting from accretion should be $`10^{31}`$ erg cm<sup>-2</sup> s<sup>-1</sup>โeasily sufficient to be observed in nearby associations of T Tauri stars and regions of star formation.
The above scenario of copious accretion-driven X-rays from CTTS contrasts sharply with what is observed. While strong X-ray emission appears to be a ubiquitous characteristic of CTTS (e.g. Feigelson & Montmerle, 1999), the source plasma temperatures are an order of magnitude higher than can be produced in accretion shocks and heating can instead be attributed to magnetic processes analogous to coronal activity in late-type main sequence and more evolved stars. Accretion activity is largely revealed by strong UV-optical continuum emission, suggesting that accretion shocks are either formed too deep in the stellar atmosphere to be observed (Drake, 2005), or that infall velocities are insufficient to attain X-ray temperatures in the shock. Nevertheless, two stars now stand out as probable exceptions and examples of objects whose high resolution X-ray spectra appear to be produced, at least in part, by accretion: TW Hya (Kastner et al., 2002; Stelzer & Schmitt, 2004) and BP Tau (Schmitt et al., 2005).
TW Hya has a dominant plasma temperature of $`3\times 10^6`$ K, as expected from an accretion shock resulting from free-fall of gas from a truncated disk. Based on the density-sensitive He-like Ne and O lines, both TW Hya and BP Tau seem to be characterised by plasma with high electron densities $`n_e10^{11}`$-$`10^{13}`$ cm<sup>-3</sup> at temperatures of $`3\times 10^6`$ K (Kastner et al., 2002; Stelzer & Schmitt, 2004), in contrast to all other single and active binary stars studied in the surveys of Testa et al. (2004) and Ness et al. (2004), where He-like O lines indicate $`n_e10^{10}`$ cm<sup>-3</sup>. The X-ray spectra of both stars exhibit extremely weak lines of Mg, Si and Fe, and instead are dominated by O and Ne. While this pattern is reminiscent of that seen in very active RS CVn-type binaries (e.g. Huenemoerder et al., 2001), Stelzer & Schmitt (2004) echoed the earlier suggestion of Herczeg et al. (2002) that the metal depletion is instead a signature of the accretion of grain-depleted gas. If the latter is indeed the case, X-rays then provide a unique means of investigating the matter in the inner accretion disk that is in the process of accreting.
In this Letter, we draw on the recent findings of Drake & Testa (2005), based on Chandra High Energy Transmission Grating (HETG) spectra of nearby stars, that the Ne/O abundance ratio is remarkably constant in stellar coronae. We use this result to investigate Ne/O in TW Hya and BP Tau, and show that this ratio provides a diagnostic of the depletion of O in grains in the circumstellar disk, and in particular in the very inner disk from which accreting material derives.
## 2 Observations and Analysis
We use spectral line fluxes for the H-like and He-like resonance lines of Ne and O obtained from XMM-Newton observations of TW Hya (acquired 2001 July) and BP Tau (acquired 2004 August) by Stelzer & Schmitt (2004) and Schmitt et al. (2005), respectively; the reader is referred to these works for further details regarding the observations. To provide an additional TW Hya association comparison for TW Hya itself, we also analysed the same lines for the CTTS TWA 5 studied recently by Argiroffi et al. (2005) based on XMM-Newton spectra.
TW Hya was observed earlier by the Chandra HETG and ACIS-S, on 2000 July 18 at UT10:19 (see Kastner et al., 2002, for further details). Here, we have analysed the products of standard CIAO 3.0 processing, and in particular the Medium Energy Grating spectrum. The measurement of Ne and O line spectral line fluxes and all calculations were performed using the IDL Package for INTeractive Analysis of Line Emission (PINTofALE)<sup>1</sup><sup>1</sup>1PINTofALE is freely available from http://hea-www.harvard.edu/PINTofALE/ following the methods described by Testa et al. (2004). The line fluxes of interest for this study are listed in Table 1.
The conversion of Ne and O fluxes into the Ne/O abundance ratio by number, $`A_{Ne}/A_O`$, is described by Drake & Testa (2005). Briefly, an abundance ratio in an optically-thin, collision-dominated plasma can be derived by using lines whose contribution functions $`G_{ji}(T)`$โthe product of the parent ion population and the line emissivityโhave very similar temperature dependence, simply from the ratio of observed line fluxes (corrected for attenuation by interstellar extinction), $`F_O`$ and $`F_{Ne}`$:
$$\frac{A_{Ne}}{A_O}=\overline{\left(\frac{G_O}{G_{Ne}}\right)}\frac{F_{Ne}}{F_O}$$
(1)
For Ne and O, such a ratio of $`G_{ji}(T)`$ functions can be constructed from the O VIII, Ne IX and Ne X resonance lines $`1s2p^1P_11s^2{}_{}{}^{1}S_{0}^{}`$ and $`2p^2P_{3/2,1/2}1s^2S_{1/2}`$, combined as follows
$$\overline{\left(\frac{G_O}{G_{Ne}}\right)}=\frac{G_{OVIII}}{G_{NeIX}+0.15G_{NeX}}.$$
(2)
Drake & Testa (2005) found $`\overline{G_O/G_{Ne}}=1.2\pm 0.1`$ for active stars. Ne/O abundance ratios derived for BP Tau, TW Hya and TWA 5 using this method are also listed in Table 1, together with the intervening absorbing columns used to correct the observed line fluxes for attenuation. These Ne/O ratios are compared with those presented by Drake & Testa (2005) for a sample of 21 post-T Tauri stars, including single main-sequence stars, giants, and tidally-interacting binaries in Figure 1.
## 3 Discussion
Based on the Ne/O ratios found for post-T Tauri stars, Drake & Testa (2005) drew two conclusions: (i) Ne/O is essentially constant in stellar coronae, and in full-disk observations is not susceptible to fractionation effects that often appear to characterise other elements with lower first ionisation potentials (FIPs) (e.g. Drake, 2003); (ii) the Ne/O abundance ratio in stars is significantly larger than current assessments of the solar ratio (also illustrated in Figure 1), but is in-line with inference from solar oscillations (Antia & Basu, 2005; Bahcall et al., 2005). The constancy of the coronal Ne/O ratio then allows us to use these elements for diagnostics purposes. Firstly, we can determine whether or not the suspected accretion shock has an Ne/O ratio consistent with that of the underlying star, as represented by the constant coronal Ne/O ratio illustrated in Figure 1. A ratio significantly different to the coronal value would provide further evidence that these lines are not formed in a โnormalโ coronal plasma. Secondly, departures from the coronal Ne/O ratio can then be used to infer compositional peculiarities in the accreting gas.
We conclude from Figure 1 that the Ne/O ratio of $`A_{Ne}/A_O1.0`$ we find in TW Hyaโin good agreement with the value obtained earlier by Kastner et al. (2002) based on a differential emission measure analysisโis significantly higher than that of more evolved stars (Figure 1). In contrast, Ne/O in BP Tau is perfectly consistent with that found for the rest of the sample of Drake & Testa (2005). This latter result also supports the interpretation of the Ne/O abundance in stellar coronae as representative of the ambient local cosmos: under the assumption that the Ne and O lines in BP Tau are formed from accreting material, as the anomalous O VII line ratios analysed by Schmitt et al. (2005) suggest, the ratio seen in BP Tau indicates that its circumstellar material shares the same Ne/O ratio.
We can confirm that the Ne/O ratio found for TW Hya does not simply a reflect an anomalous local environmental composition by comparison with the ratio seen in TWA 5, a CTTS from the same association as TW Hya. The X-ray spectrum of TWA 5 has been analysed recently by Argiroffi et al. (2005): unlike TW Hya and BP Tau, it does not show the high plasma density signature of an accretion shock, and its Ne and O lines can be interpreted as being formed in magnetically-heated coronal plasma. Using the Argiroffi et al. (2005) O and Ne line fluxes (Table 1) we obtain $`A_{Ne}/A_O=0.52\pm 0.09`$โin good agreement with their result based on a differential emission measure analysis and in excellent agreement with that of the rest of our star sample, except for TW Hya. Similarly, Kastner et al. (2004) find $`A_{Ne}/A_O0.5`$ for the TW Hya association multiple, non-accreting, weak-lined T Tauri star system HD 98800 based on Chandra HETG spectra.
Herczeg et al. (2002) suggested that the lack of Si in the UV spectrum of the accretion shock of TW Hya was a signature of the accretion of grain-depleted material. Stelzer & Schmitt (2004) invoked their suggestion to explain the metal-poor X-ray spectrum. However, since tidally-interacting binaries have also been shown to exhibit similar metal paucity (see, e.g. reviews by Drake, 2003; Audard, 2005), it is not immediately clear that the composition of the hot plasma in TW Hya is significantly different. The Ne/O ratio provides this evidence, and supports the conjecture that TW Hya is accreting grain-depleted gas. Unlike comparisons of Ne with metals such as Mg, Si and Fe, the Ne/O diagnostic appears to be robust to the effects of compositional fractionation seen in coronal plasma. Moreover, the Ne and O lines are formed at the temperature of $`3\times 10^6`$ K expected of the accretion shocks on T Tauri stars, whereas dominant lines of Fe XVII and higher charge states, Mg and Si are all formed at higher temperatures.
In the context of shock-produced versus coronal X-rays, it is worth noting that the observed departure from a โnormalโ chemical composition cannot be mimiced by a corona whose plasma is fed by accretion. The convective turnover time for CTTS is of order a few hundred days (Gilliland, 1986; Kim & Demarque, 1996) and so accreting plasma, and any abundance peculiarity in these last accumulating fractions, is rapidly subsumed by the underlying star that would feed the corona.
We interpret the different Ne/O ratios in BP Tau and TW Hya as arising from different depletions of O in the accreting gas. The comparison of BP Tau with post-T Tauri stars indicates that it is accreting material which is essentially undepleted. In the case of TW Hya, the O depletion amounts to a factor of $`2`$ or more. Why should BP Tau and TW Hya exhibit distinctly different O depletion in their accreting plasma? We propose that this reflects the different ages and evolutionary states of these stars.
There is some controversy concerning the age of BP Tau, though this is entirely the result of its uncertain Hipparcos parallax. The Hipparcos data indicate a distance of only 42-70 pc ($`\pm 1\sigma `$ range)โwell in front of the Taurus cloud at 140 pc in which it is generally thought to reside (Wichmann et al., 1998; Favata et al., 1998; Bertout et al., 1999). Based on this distance, Favata et al. (1998) estimated BP Tau to be as old as 35 Myr through comparison with evolutionary tracks. However, both Wichmann et al. (1998) and Bertout et al. (1999) refute the Hipparcos distance based on its low statistical significance. Bertout et al. (1999) find that the astrometric fit to the Hipparcos data for the distance of the Taurus cloud of 140 pc, instead of 42-70 pc, represents only a $`2\sigma `$ deviation, and point out that BP Tau is rather faint for the nominal distance to be reliable. They also show the Hipparcos distance of 32-52 pc for for DF Tauโthe other Taurus star found by Favata et al. (1998) to be of anomalous proximityโto be highly uncertain owing to binarity and photometric variability. Allowing for a somewhat smaller distance than 140 pc, but not as extreme as the Hipparcos value, Simon et al. (2000) estimate a range of 2-10 Myr. However, all other observational evidence points to BP Tau lying in the Taurus cloud, including radial velocity and proper motion (Hartmann & Stauffer, 1989; Bertout et al., 1999), extinction of $`A_V0.5`$ (e.g. Gullbring et al., 1996a; Muzerolle et al., 2003, and references therein) and strong interstellar Na I absorption (Gullbring et al., 1996b). Gullbring et al. (1998) estimate the age to be 0.6 Myr.
In contrast, the distance and age of TW Hya are much more secure: with an age of $`10`$ Myr (e.g. Zuckerman & Song, 2004), TW Hya appears significantly older than BP Tau. TW Hya is in fact one of the oldest CTTS known to be still accreting. More importantly, near infra-red measurements also show the inner disks of the two stars to be quite different. Veiling measurements for BP Tau show a very large excess of 0.6 at $`2.2\mu `$m (Muzerolle et al., 2003), and $`KL=0.57\pm 0.23`$ (Kenyon & Hartmann, 1995)โtypical of the accreting stars in Taurus, 95% of which have $`KL`$ colours in the range 0.35-1.2 (Meyer et al., 1997). Muzerolle et al. (2003) find that BP Tau has an optically thick disk whose dust truncation radius is essentially at the corotation radius from which magnetospheric accretion is thought to take placeโmaterial from both gas and dust will therefore accrete and this material will not be depleted.
Instead, TW Hya has very little near-infrared excess shortward of 10$`\mu `$m, with $`KL=0.25`$ (Sitko et al., 2000; Calvet et al., 2002; Uchida et al., 2004). This is more typical of $`KL`$ for the nonaccreting stars in Taurus that are in the range $`0.05`$-0.25 for similar spectral types (K7-M2) (Meyer et al., 1997). Indeed, the inner disk of TW Hya appears to be almost completely cleared (e.g. Calvet et al., 2002; Rettig et al., 2004). Consequently there is very little material in the form of small dust particles available to accrete and replenish the gain-depleted gas. Depletion is also in evidence in the unexpectedly weak or absent Si lines in UV spectra (Valenti et al., 2000; Herczeg et al., 2002), and in low Al and Si abundances in the gas jet of TW Hya inferred by Lamzin et al. (2004). This depletion is in qualitative accord with the infrared detection of excess emission at 8-13$`\mu `$m attributable to silicates, and the spectral energy distribution that indicates that grain formation and coagulation into larger particles is well advanced (Weinberger et al., 2002; Sitko et al., 2000; Calvet et al., 2002). Larger grains would be expected to settle in the midplane of the circumstellar disk (Chiang et al., 2001; DโAlessio et al., 1999). Calvet et al. (2002) speculate that the truncation of the outer disk at $`4`$ AU inferred from the negligible near-infrared excess could be a developing gap caused by a growing protoplanet. In this scenario, the depletion of metals and oxygen in the X-ray spectrum of the shocked accreting gas would be the result of retention of these elements in bodies of sufficient size that they are not significantly affected by the drag of inwardly migrating gas that feeds the accretion.
## 4 Conclusions
We have investigated the Ne/O abundance ratios of the accreting gas of the T Tauri stars TW Hya and BP Tau using high resolution X-ray spectra of H-like and He-like resonance lines, and have compared these to the same ratio for the coronal plasma of sample of magnetically-active main-sequence and evolved field stars. In the case of BP Tau, we find an Ne/O abundance ratio in good agreement with that of the coronal plasma of field stars. In contrast, the Ne/O ratio for TW Hya is higher than for the rest of the sample by a factor of two.
We interpret these results in terms of a relative depletion of O in the material accreting from the inner disk of TW Hya as compared with BP Tau. This can be attributed to the different, partially age-related, evolutionary states of the disks of these stars. In the BP Tau disk, dust is still present near the disk corotation radius and can be ionized and accreted, re-releasing elements depleted onto grains. In the more evolved TW Hya disk, evidence points to ongoing coagulation of grains into much larger bodies, and possibly planets, that can resist the drag of inward-migrating gas, and accreting gas is consequently depleted of grain-forming elements.
These results demonstrate the utility of high resolution X-ray spectra of CTTS for probing directly the chemical composition of the accreting gas and the state of evolution of the very inner protoplanetary disk.
We thank the NASA AISRP for providing financial assistance for the development of the PINTofALE package. JJD was supported by NASA contract NAS8-39073 to the Chandra X-ray Center during the course of this research. PT was supported by Chandra award number G03-4005A issued by CXC, and SAO contract SV3-73016 to MIT for support of CXC, which is operated by SAO for and on behalf of NASA under contracts NAS8-39073 and NAS8-03060. <sup>d</sup><sup>d</sup>footnotemark:
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# Introduction
## Introduction
The goal of these lectures is to convince the reader to construct the supersymmetric mechanics within the Hamiltonian framework, or, at least, to combine the superfield approach with the existing methods of Hamiltonian mechanics. The standard approach to construct the supersymmetric mechanics with more than two supercharges is the Lagrangian superfield approach. Surely, superfield formalism is a quite powerful method for the construction of supersymmetric theories. However, all superfield formalisms, being developed รก priori for field theory, are convenient for the construction of the field-theoretical models, which are covariant with respect to space-time coordinate transformations. However, the supermultiplets (i.e. the basic ingredients of superfield formalisms) do not respect the transformations mixing field variables. On the other hand, in supersymmetric mechanics these variables appear as spatial coordinates. In other words, the superfield approach, being applied to supersymmetric mechanics, provides us with a local construction of mechanical models. Moreover, the obtained models need to be re-formulated in the Hamiltonian framework, for the subsequent quantization. In addition, many of the numerous methods and statements in the Hamiltonian formalism could be easily extended to supersymmetric systems and applied there. Independently from the specific preferences, the โHamiltonian viewโ of the existing models of supersymmetric mechanics, which were built within the superfield approach, could establish unexpected links between different supermultiplets and models. Finally, the superfield methods seem to be too general in the context of simple mechanical systems.
For this reason, we tried to present some elements of Hamiltonian formalism, which do not usually appear in the standard textbooks on classical mechanics, but appear to be useful in the context of supersymmetric mechanics. We pay much attention to the procedure of Hamiltonian reduction, having in mind that it could be used for the construction of the lower-dimensional supersymmetric models from the existing higher-dimensional ones. Also, we devote a special attention to the Hopf maps and Kรคhler spaces, which are typical structures in supersymmetric systems. Indeed, to extend the number of supersymmetries (without extension of the fermionic degrees of freedom) we usually equip the configuration/phase space with complex structures and restrict them to be Kรคhler, hyper-Kรคhler, quaternionic and so on, often via a choice of the appropriate supermultiplets related to the real, complex, quaternionic structures. We illustrated these matters by examples of Hamiltonian reductions related with Hopf maps, having in mind that themy could be straightforwardly applied to supersymmetric systems. Also, we included some less known material related with Hopf fibrations. It concerns the generalization of the oscillator to spheres, complex projective spaces, and quaternionic projective spaces, as well as the reduction of the oscillator systems to Coulomb ones.
Most of the presented constructions are developed only for the zero and first Hopf maps. We tried to present them in the way, which will clearly show, how to extend them to the second Hopf map and the quaternionic case.
The last two sections are devoted to the super-Hamiltonian formalism. We present the superextensions of the Hamiltonian constructions, underlying the specific โsuperโ-properties, and present some examples. Then we provide the list of supersymmetric mechanics constructed within the Hamiltonian approach. Also in this case, we tried to arrange the material in such a way, as to make clear the relation of these constructions to complex structures and their possible extension to quaternionic ones.
The main references to the generic facts about Hamiltonian mechanics are the excellent textbooks , and on the supergeometry there exist the monographs . There are numerous reviews on supersymmetric mechanics. In our opinion the best introduction to the subject is given in refs. .
## 1 Hamiltonian formalism
In this Section we present some basic facts about the Hamiltonian formalism, which could be straightforwardly extended to the super-Hamiltonian systems.
We restrict ourselves to considering Hamiltonian systems with nondegenerate Poisson brackets. These brackets are defined, locally, by the expressions
$$\{f,g\}=\frac{f}{x^i}\omega ^{ij}(x)\frac{g}{x^j},det\omega ^{ij}0,$$
(1.1)
where
$`\{f,g\}=\{g,f\},\omega ^{ij}=\omega ^{ji}`$ (1.2)
$`\{\{f,g\},h\}+\mathrm{cycl}.\mathrm{perm}(f,g,h)=0,\omega _{,n}^{ij}\omega ^{nk}+\mathrm{cycl}.\mathrm{perm}(i,j,k)=0.`$ (1.3)
The Eq. (1.2) is known as a โantisymmetricity conditionโ, and the Eq.(1.3) is called Jacobi identity. Owing to the nondegeneracy of the matrix $`\omega ^{ij}`$, one can construct the nondegenerate two-form, which is closed due to Jacobi identity
$$\omega =\frac{1}{2}\omega _{ij}dx^idx^j:d\omega =0\omega _{ij,k}+\mathrm{cycl}.\mathrm{perm}(i,j,k)=0.$$
(1.4)
The manifold $`M`$ equipped with such a form, is called symplectic manifold, and denoted by $`(M,\omega )`$. It is clear that $`M`$ is an even-dimensional manifold, $`dimM=2N`$.
The Hamiltonian system is defined by the triple $`(M,\omega ,H)`$, where $`H(x)`$ is a scalar function called Hamiltonian.
The Hamitonian equations of motion yield the vector field preserving the symplectic form $`\omega `$
$$\frac{dx^i}{dt}=\{H,x^i\}=V_H^i:_{V_H}\omega =0.$$
(1.5)
Here $`_๐`$ denotes the Lie derivative along vector field $`๐`$.
Vice versa, any vector field, preserving the symplectic structure, is locally a Hamiltonian one. The easiest way to see it is to use homotopy formula
$$ฤฑ_๐d\omega +dฤฑ_๐\omega =_๐\omega dฤฑ_๐\omega =0.$$
(1.6)
Hence, $`ฤฑ_๐\omega `$ is a closed one-form and could be locally presented as follows: $`ฤฑ_๐\omega =dH(x)`$. The local function $`H(x)`$ is precisely the Hamiltonian, generating the vector field $`๐`$. The transformations preserving the symplectic structure are called canonical transformations.
Any symplectic structure could be locally presented in the form (Darboux theorem)
$$\omega _{can}=\underset{a=1}{\overset{N}{}}dp_adq^a,$$
(1.7)
where $`(p_a,q^a)`$ are the local coordinates of the symplectic manifold.
The vector field $`๐`$ defines a symmetry of the Hamiltonian system, if it preserves both the Hamiltonian $``$ and the symplectic form $`\omega `$: $`_{๐\omega }=0`$, $`๐=0`$. Hence,
$$๐=\{๐ฅ,\},\{๐ฅ,\}=0.$$
(1.8)
The $`2N`$-dimensional Hamiltonian system is called an integrable system, when it has $`N`$ functionally independent constants of motion being in involution (Liouville theorem),
$$\{๐ฅ_a,๐ฅ_b\}=0,\{,๐ฅ_b\}=0,=๐ฅ_1,a,b=1,\mathrm{},N.$$
(1.9)
When the constants of motion are noncommutative, the integrability of the system needs more than $`N`$ constants of motion. If
$$\{๐ฅ_\mu ,๐ฅ_\nu \}=f_{\mu \nu }(๐ฅ),\mathrm{corank}f_{\mu \nu }=K_0,\mu ,\nu =1,\mathrm{},KK_0,$$
(1.10)
then the system is integrable, if $`2N=K+K_0`$ . The system with $`K+K_02N`$ constants of motion is sometimes called a superintegrable system.
The cotangent bundle $`T^{}M_0`$ of any manifold $`M_0`$ (parameterized by local coordinates $`q^i`$) could be equipped with the canonical symplectic structure (1.7).
The dynamics of a free particle moving on $`M_0`$ is given by the Hamiltonian system
$$\left(T^{}M_0,\omega _{can},_0=\frac{1}{2}g^{ab}(q)p_ap_b\right),$$
(1.11)
where $`g^{ab}g_{bc}=\delta _c^a`$, and $`g_{ab}dq^adq^b`$ is a metric on $`M_0`$.
The interaction with a potential field could be incorporated in this system by the appropriate change of Hamiltonian,
$$_0=\frac{1}{2}g^{ab}(q)p_ap_b+U(q),$$
(1.12)
where $`U(q)`$ is a scalar function called potential. Hence, the corresponding Hamiltonian system is given by the triplet $`(T^{}M_0,\omega _{can},)`$.
In contrast to the potential field, the interaction with a magnetic field requires a change of symplectic structure. Instead of the canonical symplectic structure $`\omega _{can}`$, we have to choose
$$\omega _F=\omega _{can}+F,F=\frac{1}{2}F_{ab}(q)dq^adq^b,dF=0$$
(1.13)
where $`F_{ab}`$ are components of the magnetic field strength.
Hence, the resulting system is given by the triplet $`(T^{}M_0,\omega _F,)`$. Indeed, taking into account that the two-form $`F`$ is locally exact, $`F=dA`$, $`A=A_a(q)dq^a`$, we could pass to the canonical coordinates $`(\pi _a=p_a+A_a,q^a)`$. In these coordinates the Hamiltonian system assumes the conventional form
$$\left(T^{}M_0,\omega _{can}=d\pi _adq^a,=\frac{1}{2}g^{ab}(\pi _aA_a)(\pi _bA_b)+U(q)\right).$$
Let us also remind, that in the three-dimensional case the magnetic field could be identified with vector, whereas in the two-dimensional case it could be identified with (pseudo)scalar.
The generic Hamiltonian system could be described by the following (phase space) action
$$๐ฎ=๐t\left(๐_i(x)\dot{x}^i(x)\right),$$
(1.14)
where $`๐=๐_idx^i`$ is a symplectic one-form: $`d๐=\omega `$. Indeed, varying the action, we get the equations
$$\delta ๐ฎ=0,\dot{x}^i\omega _{ij}(x)=\frac{H}{x^i},\omega _{ij}=\frac{A_i}{x^j}\frac{A_j}{x^i}.$$
Though $`๐`$ is defined up to closed (locally exact) one-form, $`๐๐+df(x)`$, this arbitrariness has no impact in the equations of motion. It change the Lagrangian on the total derivative $`f_{,i}\dot{x}^i=df(x)/dt`$.
As an example, let us consider the particle in a magnetic field. The symplectic one-form corresponding to the symplectic structure (1.13), could be chosen in the form $`๐=(p_a+A_a)dq^a`$, $`d๐=\omega _F`$. Hence, the action (1.14) reads
$$๐ฎ=๐t\left((p_a+A_a)\dot{q}^a\frac{1}{2}g^{ab}(q)p_ap_bU(q)\right).$$
(1.15)
Varying this action by $`p`$, we get, on the extrema, the conventional second-order action for the system in a magnetic field
$$๐ฎ_0=๐t\left(\frac{1}{2}g_{ab}\dot{q}^a\dot{q}^b+A_a\dot{q}^aU(q)\right).$$
(1.16)
The presented manipulations are nothing but the Legendre transformation from the Hamiltonian formalism to the Lagrangian one.
### Particle in the Dirac monopole field
Let us consider the special case of a system on three-dimensional space moving in the magnetic field of a Dirac monopole. Its symplectic structure is given by the expression
$$\omega _D=dp_adq^a+s\frac{q^a}{2|q|^3}ฯต_{abc}dq^bdq^c.$$
(1.17)
The corresponding Poisson brackets are given by the relations
$$\{p_a,q^b\}=\delta _a^b,\{q^a,q^b\}=0,\{p_a,p_b\}=sฯต_{abc}\frac{q^c}{|q|^3}.$$
(1.18)
It is clear that the monopole field does not break the rotational invariance of the system. The vector fields generating $`SO(3)`$ rotations are given by the expressions
$$๐_a=ฯต_{abc}q^b\frac{}{q^c}ฯต_{abc}p_b\frac{}{p_c},[๐_a,๐_b]=ฯต_{abc}๐_c.$$
(1.19)
The corresponding Hamiltonian generators could be easily found as well
$$ฤฑ_{๐_a}\omega _D=d๐ฅ_a,\{๐ฅ_a,๐ฅ_b\}=ฯต_{abc}๐ฅ_c,$$
where
$$๐ฅ_a=ฯต_{abc}q^bp_c+s\frac{q^a}{|q|},J_aq^a=s|q|.$$
(1.20)
Now, let us consider the system given by the symplectic structure (1.17), and by the $`so(3)`$-invariant Hamiltonian
$$=\frac{p_ap_a}{2g}+U(|q|),\{๐ฅ_a,\}=0,$$
(1.21)
where $`g(|q|)dq^adq^a`$ is $`so(3)`$-invariant metric on $`M_0`$. In order to find the trajectories of the system, it is convenient to direct the $`q^3`$ axis along the vector $`๐=(๐ฅ_1,๐ฅ_2,๐ฅ_3)`$, i.e. to assume that $`๐=J_3J`$. Upon this choice of the coordinate system one has
$$\frac{q^3}{|q|}=\frac{s}{J}.$$
(1.22)
Then, we introduce the angle
$$\varphi =\mathrm{arctan}\frac{q^1}{q^2},\frac{d\varphi }{dt}=\frac{J^2s^2}{Jg|q|^2},$$
(1.23)
and get, after obvious manipulations
$$=\frac{J^2s^2}{2g|q|^2}+\frac{(J^2s^2)^2}{2Jg|q|^2}\left(\frac{d|q|}{d\varphi }\right)^2+U(|q|).$$
(1.24)
Here $``$ denotes the energy of the system.
From the expression (1.24) we find,
$$\varphi =\left(J\frac{s^2}{J}\right)\frac{d|q|}{\sqrt{2g|q|^2(U)J^2+s^2}}.$$
(1.25)
It is seen that, upon the replacement
$$U(q)U(q)+\frac{s^2}{2g|q|^2},$$
(1.26)
we shall eliminate in (1.25) the dependence on $`s`$, i.e. on a monopole field. The only impact of the monopole field on the trajectory will be the shift of the orbital plane given by (1.22).
Let us summarize our considerations. Let us consider the $`so(3)`$-invariant three-dimensional system
$$\omega _{can}=d๐ฉd๐ช,=\frac{๐ฉ^2}{2g}+U(|๐ช|),\{๐_0,\}=0,๐_0=๐ฉ\times ๐ช.$$
(1.27)
Then, replacing it by the following one:
$$\omega _{can}+s\frac{๐ช(d๐ช\times d๐ช)}{2|๐ช|^3},=\frac{1}{2g}\left(๐ฉ^2+\frac{s^2}{|๐ช|^2}\right)+U(|๐ช|),๐=๐_0+s\frac{๐ช}{|๐ช|},$$
(1.28)
we shall preserve the form of the orbit of the initial system, but shift it along $`๐`$ in accordance with (1.22).
One can expect that, when the initial system has a symmetry, additional with respect to the rotational one, the latter system will also inherit it. For the Coulomb system, $`U=\gamma /|๐ช|`$, this is indeed a case. The modified system (which is known as a MIC-Kepler system) possesses the hidden symmetry given by the analog of the Runge-Lenz vector, which is completely similar to the Runge-Lenz vector of the Kepler system .
### Kรคhler manifolds
One of the most important classes of symplectic manifolds is that of Kรคhler manifolds. The Hermitean manifold $`(M,g_{a\overline{b}}dz^adz^b)`$ is called Kรคhler manifold, if the imaginary part of the Hermitean structure is a symplectic two-form (see, e.g. ):
$$\omega =ig_{a\overline{b}}dz^ad\overline{z}^b:d\omega =0,detg_{a\overline{b}}0.$$
(1.29)
The Poisson brackets associated with this symplectic structure read
$$\{f,g\}_0=i\frac{f}{\overline{z}^a}g^{\overline{a}b}\frac{g}{z^b}i\frac{g}{z^b}g^{\overline{a}b}\frac{f}{\overline{z}^a},\mathrm{where}g^{\overline{a}b}g_{b\overline{c}}=\delta _{\overline{c}}^{\overline{a}}.$$
(1.30)
ยฟFrom the closeness of (1.29) it immediately follows, that the Kรคhler metric can be locally represented in the form
$$g_{a\overline{b}}dz^ad\overline{z}^b=\frac{^2K}{z^a\overline{z}^b}dz^ad\overline{z}^b,$$
(1.31)
where $`K(z,\overline{z})`$ is some real function called the Kรคhler potential. The Kรคhler potential is defined modulo holomorphic and antiholomorphic functions
$$K(z,\overline{z})K(z,\overline{z})+U(z)+\overline{U}(\overline{z}).$$
(1.32)
The local expressions for the differential-geometric objects on Kรคhler manifolds are also very simple. For example, the non-zero components of the metric connections (Cristoffel symbols) look as follows:
$$\mathrm{\Gamma }_{bc}^a=g^{\overline{n}a}g_{b\overline{n},c},\mathrm{\Gamma }_{\overline{b}\overline{c}}^{\overline{a}}=\overline{\mathrm{\Gamma }}_{bc}^a,$$
(1.33)
while the non-zero components of the curvature tensor read
$$R_{bc\overline{d}}^a=(\mathrm{\Gamma }_{bc}^a)_{,d},R_{\overline{b}\overline{c}d}^{\overline{a}}=\overline{R}_{bc\overline{d}}^a.$$
(1.34)
The isometries of Kรคhler manifolds are given by the holomorphic Hamiltonian vector fields
$$๐_\mu =V_\mu ^a(z)\frac{}{z^a}+\overline{V}_\mu ^{\overline{a}}(\overline{z})\frac{}{\overline{z}^a},,๐_\mu =\{๐_\mu ,\}_0,$$
(1.35)
where $`๐_\mu `$ is a real function, $`๐_\mu =\overline{๐}_\mu `$, called Killing potential. One has
$$[๐_\mu ,๐_\nu ]=C_{\mu \nu }^\lambda ๐_\lambda ,\{๐_\mu ,๐_\nu \}_0=C_{\mu \nu }^\lambda ๐_\lambda +\mathrm{const},$$
and
$$\frac{^2๐_\mu }{z^az^b}\mathrm{\Gamma }_{ab}^c\frac{๐_\mu }{z^c}=0.$$
The dynamics of a particle moving on the Kรคhler manifold in the presence of a constant magnetic field is described by the Hamiltonian system
$$\mathrm{\Omega }_B=dz^ad\pi _a+d\overline{z}^ad\overline{\pi }_a+iBg_{a\overline{b}}dz^ad\overline{z}^b,_0=g^{a\overline{b}}\pi _a\overline{\pi }_b$$
(1.36)
The isometries of a Kรคhler structure define the Noether constants of motion
$$๐ฅ_\mu J_\mu +B๐_\mu =V_\mu ^a\pi _a+\overline{V}_\mu ^{\overline{a}}\overline{\pi }_{\overline{a}}+B๐_\mu :\{\begin{array}{c}\{_0,J_\mu \}=0,\\ \{J_\mu ,J_\nu \}=C_{\mu \nu }^\lambda J_\lambda .\end{array}$$
(1.37)
One can easily check that the vector fields generated by $`๐ฅ_\mu `$ are independent of $`B`$
$$๐=V^a(z)\frac{}{z^a}V_{,b}^a\pi _a\frac{}{\pi _a}+\overline{V}^a(\overline{z})\frac{}{\overline{z}^a}\overline{V}_{,\overline{b}}^a\overline{\pi }_a\frac{}{\overline{\pi }_a}.$$
(1.38)
Hence, the inclusion of a constant magnetic field preserves the whole symmetry algebra of a free particle moving on a Kรคhler manifold.
### Complex projective space
The most known nontrivial example of a Kรคhler manifold is the complex projective space $`\text{I}\mathrm{CP}^\mathrm{N}`$. It is defined as a space of complex lines in $`\text{I}\mathrm{C}^{\mathrm{N}+1}`$: $`u^{\stackrel{~}{a}}\lambda u^{\stackrel{~}{a}}`$, where $`u^{\stackrel{~}{a}}`$, $`\stackrel{~}{a}=0,1,\mathrm{},N`$ are the Euclidean coordinates of $`\text{I}\mathrm{C}^{\mathrm{N}+1}`$, and $`\lambda \text{I}\mathrm{C}\{0\}`$. Equivalently, the complex projective space is the coset space $`\text{I}\mathrm{CP}^\mathrm{N}=\mathrm{SU}(\mathrm{N}+1)/\mathrm{U}(\mathrm{N})`$.
The complex projective space $`\text{I}\mathrm{CP}^\mathrm{N}`$ could be covered by $`N+1`$ charts marked by the indices $`\stackrel{~}{a}=0,a`$. The zero chart could be parameterized by the functions (coordinates) $`z_{(0)}^a=u^a/u^0`$, $`a=1,\mathrm{}N`$; the first chart by $`z_{(1)}^a=z^a/z^1`$, $`a=0,2,3\mathrm{},N`$, and so on.
Hence, the transition function from the $`\stackrel{~}{b}`$-th chart to the $`\stackrel{~}{c}`$-th one has the form
$$z_{(\stackrel{~}{c})}^{\stackrel{~}{a}}=\frac{z_{(\stackrel{~}{b})}^{\stackrel{~}{a}}}{z_{(\stackrel{~}{b})}^{\stackrel{~}{c}}},\mathrm{where}z_{(\stackrel{~}{a})}^{\stackrel{~}{a}}=1.$$
(1.39)
One can equip the $`\text{I}\mathrm{CP}^\mathrm{N}`$ by the Kรคhler metric, which is known under the name of Fubini-Study metric
$$g_{a\overline{b}}dz^adz^b=\frac{dzd\overline{z}}{1+z\overline{z}}\frac{(\overline{z}dz)(zd\overline{z})}{(1+z\overline{z})}.$$
(1.40)
Its Kรคhler potential is given by the expression
$$K=\mathrm{log}(1+z\overline{z}).$$
(1.41)
Indeed, it is seen that upon transformation from one chart to the other, given by (1.39), this potential changes by holomorphic and anti-holomorphic functions, i.e. the Fubini-Study metric is globally defined on $`\text{I}\mathrm{CP}^\mathrm{N}`$.
The Poisson brackets on $`\text{I}\mathrm{CP}^\mathrm{N}`$ are defined by the following relations:
$$\{z^a,\overline{z}^b\}=(1+z\overline{z})(\delta ^{a\overline{b}}+z^a\overline{z}^b),\{z^a,z^b\}=\{\overline{z}^a,\overline{z}^b\}=0.$$
(1.42)
It is easy to see that $`\text{I}\mathrm{CP}^\mathrm{N}`$ is a constant curvature space, with the symmetry algebra $`su(N+1)`$. This algebra is defined by the Killing potentials
$$๐_{\overline{a}b}=\frac{z^a\overline{z}^bN\delta _{\overline{a}b}}{1+z\overline{z}},๐_a^{}=\frac{z^a}{1+z\overline{z}},๐_a^+=\frac{\overline{z}^a}{1+z\overline{z}}.$$
(1.43)
The manifold $`\text{I}\mathrm{CP}^1`$ (complex projective plane) is isomorphic to the two-dimensional sphere $`S^2`$. Indeed, it is covered by the two charts, with the transition function $`z1/z`$. The symmetry algebra of $`\text{I}\mathrm{CP}^1`$ is $`su(2)=so(3)`$
$$\{x^i,x^j\}=ฯต^{ijk}x^k,i,j,k=1,2,3$$
(1.44)
where the Killing potentials $`x^i`$ look as follows:
$$x^1+ix^2=\frac{2z}{1+z\overline{z}},x^3=\frac{1z\overline{z}}{1+z\overline{z}}.$$
(1.45)
It is seen that these Killing potentials satisfy the condition
$$x^ix^i=1,$$
i.e. $`x^i`$ defines the sphere $`S^2`$ in the three-dimensional ambient space $`\mathrm{IR}^3`$. It is straightforwardly checked that $`z`$ are the coordinates of the sphere in the stereographic projection on $`\mathrm{IR}^2=\text{I}\mathrm{C}`$. The real part of the Fubini-Study structure gives the linear element of $`S^2`$, and the imaginary part coincides with the volume element of $`S^2`$.
On the other hand, these expressions give the embedding of the $`S^2`$ in $`S^3`$ (with ambient coordinates $`u^1,u^2`$) defining the so-called first Hopf map $`S^3/S^1=S^2`$. Below we shall describe this map in more detail.
### Hopf maps
The Hopf maps (or Hopf fibrations) are the fibrations of the sphere over a sphere,
$$S^{2p1}/S^{p1}=S^p,p=1,2,4,8.$$
(1.46)
These fibrations reflect the existence of real ($`p=1`$), complex ($`p=2`$), quaternionic ($`p=4`$) and octonionic ($`p=8`$) numbers.
We are interested in the so-called zero-th, first and second Hopf maps:
$`S^1/S^0=S^1(\mathrm{zero}\mathrm{Hopf}\mathrm{map})`$
$`S^3/S^1=S^2(\mathrm{first}\mathrm{Hopf}\mathrm{map})`$
$`S^7/S^3=S^4(\mathrm{second}\mathrm{Hopf}\mathrm{map})`$ .
Let us describe the Hopf maps in explicit terms. For this purpose, we consider the functions $`๐ฑ(u,\overline{u}),x_0(u,\overline{u})`$
$$๐ฑ=2u_1\overline{u}_2,x_{p+1}=u_1\overline{u}_1u_2\overline{u}_2,$$
(1.48)
where $`u_1,u_2`$, could be real, complex or quaternionic numbers. So, one can consider them as a coordinates of the $`2p`$-dimensional space $`\mathrm{IR}^{2\mathrm{p}}`$, where $`p=1`$ when $`u_{1,2}`$ are real numbers; $`p=2`$ when $`u_{1,2}`$ are complex numbers; $`p=4`$ when $`u_{1,2}`$ are quaternionic numbers; $`p=8`$ when $`u_{1,2}`$ are octonionic ones.
In all cases $`x_{p+1}`$ is a real number, while $`๐ฑ`$ is, respectively, a real number ($`p=1`$), complex number ($`p=2`$), quaternion($`p=4`$), or octonion ($`p=8`$). Hence, $`(x_0,๐ฑ)`$ parameterize the $`(p+1)`$-dimensional space $`\mathrm{IR}^{\mathrm{p}+1}`$.
The functions $`๐ฑ,x_{p+1}`$ remain invariant under transformations
$$u_a๐ u_a,\mathrm{where}g\overline{g}=1.$$
(1.49)
Hence
$`g=\pm 1\mathrm{for}p=1`$ (1.50)
$`g=\lambda _1+i\lambda _2,\lambda _1^2+\lambda _2^2=1\mathrm{for}p=2`$ (1.51)
$`g=\lambda _1+i\lambda _2+j\lambda _3+k\lambda _4,\lambda _1^2+\mathrm{}+\lambda _4^2=1\mathrm{for}p=4.`$ (1.52)
and similarly for the octonionic case $`p=8`$.
So, $`g`$ parameterizes the spheres $`S^{p1}`$ of unit radius. Notice that $`S^1,S^3,S^7`$ are the only parallelizable spheres. We shall also use the following isomorphisms between these spheres and groups: $`S^0=Z_2`$, $`S^1=U(1)`$, $`S^3=SU(2)`$.
We get that (1.48) defines the fibrations
$$\mathrm{IR}^2/\mathrm{S}^0=\mathrm{IR}^2,\mathrm{IR}^4/\mathrm{S}^1=\mathrm{IR}^3,\mathrm{IR}^8/\mathrm{S}^3=\mathrm{IR}^5,\mathrm{IR}^{16}/\mathrm{S}^7=\mathrm{IR}^9.$$
(1.53)
One could immediately check that the following equation holds:
$$๐ฑ\overline{๐ฑ}+x_{p+1}^2=(u_1\overline{u}_1+u_2\overline{u}_2)^2.$$
(1.54)
Thus, defining the $`(2p1)`$\- dimensional sphere in $`\mathrm{IR}^{2\mathrm{p}}`$ of the radius $`r_0`$: $`u_a\overline{u}_a=r_0`$, we will get the $`p`$-dimensional sphere in $`\mathrm{IR}^{\mathrm{p}+1}`$ with radius $`R_0=r_0^2`$
$$u_1\overline{u}_1+u_2\overline{u}_2=r_0^2๐ฑ\overline{๐ฑ}+x_0^2=r_0^4.$$
(1.55)
So, we arrive at the Hopf maps given by (1). The last, fourth Hopf map, $`S^{15}/S^7=S^8`$, corresponding to $`p=8`$, is related to octonions in the same manner.
For our purposes it is convenient to describe the the expressions (1.48) in a less unified way. For the zero Hopf map it is convenient to consider the initial and resulting ambient spaces $`\mathrm{IR}^2`$ as complex spaces $`\text{I}\mathrm{C}`$, parameterized by the single complex coordinates $`w`$ and $`z`$. In this case the map (1.48) could be represented in the form
$$w=z^2,$$
(1.56)
which is known as a Bohlin (or Levi-Civita) transformation relating the Kepler problem with the circular oscillator.
For the first and second Hopf maps it is convenient to represent the transformation (1.48) in the following form:
$$๐ฑ=u๐ธ\overline{u}.$$
(1.57)
Here, for the first Hopf map $`๐ฑ=(x^1,x^2,x^3)`$ parameterizes $`\mathrm{IR}^3`$, and $`u_1,u_2`$ parameterize $`\text{I}\mathrm{C}^2`$, and $`๐ธ=(\sigma ^1,\sigma ^2,\sigma ^3)`$ are Pauli matrices. This transformation is also known under the name of Kustaanheimo-Stiefel transformation. For the second Hopf map $`๐ฑ=(x^1,\mathrm{},x^5)`$ parameterizes $`\mathrm{IR}^5`$, and $`u_1,\mathrm{},u_4`$ parameterize $`\text{I}\mathrm{C}^4=\mathrm{IHI}^2`$, and $`๐ธ=(\gamma ^1,\mathrm{},\gamma ^4,\gamma ^5=\gamma ^1\gamma ^2\gamma ^3\gamma ^4)`$, where $`\gamma ^1,\mathrm{},\gamma ^4`$ are Euclidean four-dimensional gamma-matrices. The latter transformation is sometimes called Hurwitz transformation, or โgeneralized Kustaanheimo-Stiefelโ transformation.
## 2 Hamiltonian reduction
A Hamiltonian system which has a constant(s) of motion, can be reduced to a lower-dimensional one. The corresponding procedure is called Hamiltonian reduction. Let us explain the meaning of this procedure in the simplest case of the Hamiltonian reduction by a single constant of motion.
Let $`(\omega ,)`$ be a given $`2N`$-dimensional Hamiltonian system, with the phase space (local) coordinates $`x^A`$, and let $`๐ฅ`$ be its constant of motion, $`\{,๐ฅ\}=0`$. We go from the local coordinates $`x^A`$ to another set of coordinates, $`(,y^i,u)`$, where $`y^i=y^i(x)`$ are $`2N2`$ independent functions, which commute with $`๐ฅ`$,
$$\{y^i,๐ฅ\}=0,i=1,\mathrm{},2N2.$$
(2.58)
In this case the latter coordinate, $`u=u(x)`$, necessarily has a non-zero Poisson bracket with $`๐ฅ`$ (because the Poisson brackets are nondegenerate):
$$\{u(x),๐ฅ\}0.$$
(2.59)
Then, we immediately get that in these coordinates the Hamiltonian is independent of $`u`$
$$\{๐ฅ(,y,u),\}=\frac{}{u}\{u,๐ฅ\}0,=(๐ฅ,y).$$
(2.60)
On the other hand, from the Jacobi identity we get
$$\{\{y^i,y^j\},๐ฅ\}=\frac{\{y^i,y^j\}}{u}\{u,๐ฅ\}=0\{y^i,y^j\}=\omega ^{ij}(y,๐ฅ).$$
(2.61)
Since $`๐ฅ`$ is a constant of motion, we can fix its value
$$๐ฅ=c,$$
(2.62)
and describe the system in terms of the local coordinates $`y^i`$ only
$$(\omega (x),(x))\left(\omega _{red}(y,c)=\omega _{ij}(y,c)dy^idy^j,_{red}=(y,c)\right).$$
(2.63)
Hence, we reduced the initial $`2N`$-dimensional Hamiltonian system to a $`(2N2)`$-dimensional one.
Geometrically, the Hamiltonian reduction by $`๐ฅ`$ means that we fix the $`(2N1)`$\- dimensional level surface $`M_c`$ by the Eq.(2.62), and then factorize it by the action of a vector field $`\{๐ฅ,\}`$, which is tangent to $`M_c`$. The resulting space $`_0=M_c/\{๐ฅ,\}`$ is a phase space of the reduced system.
The Hamiltonian reduction by the $`K`$ commuting constants of motion $`๐ฅ`$, $`\{๐ฅ_\alpha ,๐ฅ_\beta \}=0`$ is completely similar to the above procedure. It reduces the $`2N`$ dimensional Hamiltonian system to a $`2(NK)`$ dimensional one.
When the constants of motion do not commute with each other, the reduction procedure is a bit more complicated.
Let the initial Hamiltonian system have $`K`$ constants of motion,
$$\{๐ฅ_\alpha ,\}=0,\{๐ฅ_\alpha ,๐ฅ_\beta \}=\omega _{\alpha \beta }(๐ฅ),\mathrm{corank}\omega _{\alpha \beta }|_{๐ฅ_\alpha =c_\alpha }=K_0.$$
(2.64)
Hence, one could choose the $`K_0`$ functions, which commute with the whole set of the constants of motion
$$\stackrel{~}{๐ฅ}_{\stackrel{~}{\alpha }}(๐ฅ):\{\stackrel{~}{๐ฅ}_{\stackrel{~}{\alpha }},๐ฅ_\beta \}|_{๐ฅ=c}=0,\stackrel{~}{\alpha }=1,\mathrm{}K_0.$$
(2.65)
The vector fields $`\{\stackrel{~}{๐ฅ}_{\stackrel{~}{\alpha }},\}`$ are tangent to the level surface
$$M_c:๐ฅ_\alpha =c_\alpha \mathrm{dim}M_c=2NK.$$
(2.66)
Factorizing $`M_c`$ by the action of the commuting vector fields $`\{\stackrel{~}{๐ฅ}_{\stackrel{~}{\alpha }},\}`$, we arrive at the phase space of the reduced system, $`_0=M_c/\{๐ฅ,\}`$, whose dimension is given by the expression
$$\mathrm{dimM}_0=2NKK_0.$$
(2.67)
In contrast to the commuting case, the reduced system could depend on the parameters $`\stackrel{~}{c}_{\stackrel{~}{\alpha }}`$ only.
Notice that the Hamiltonian system could also possess a discrete symmetry. In this case the reduced system has the same dimension as the previous one. To be more precise, the reduction by the discrete symmetry group could be described by a local canonical transformation. However, the quantum mechanical counterpart of this canonical transformation could yield a system with non-trivial physical properties.
Below, we shall illustrate the procedure of (Hamiltonian) reduction by discrete, commutative, and noncommutative symmetry generators on examples related to Hopf maps.
### Zero Hopf map. Magnetic flux tube
The transformation of the Hamiltonian system associated with the zero Hopf map corresponds to the reduction of the system by the discrete group $`Z_2`$. It is a (local) canonical transformation. As a consequence, the resulting system has the same dimension as the initial one.
Let us consider the Hamiltonian system with four-dimensional phase space, parameterized by the pair of canonically conjugated complex coordinates, $`(\omega =d\pi dz+d\overline{\pi }d\overline{z},)`$, which is invariant under the following action of $`Z_2`$ group:
$$(z,\overline{z},\pi ,\overline{\pi })=(z,\overline{z},\pi ,\overline{\pi }),\omega (\pi ,\overline{\pi },z,\overline{z})=\omega (\pi ,\overline{\pi },z,\overline{z}).$$
We can pass now to the coordinates, which are invariant under this transformation (clearly, it is associated with the zero Hopf map)
$`w=z^2,p=\pi /2z`$ (2.68)
$`\omega =d\pi dz+d\overline{\pi }d\overline{z}=dpdw+d\overline{p}d\overline{w}.`$ (2.69)
However, one can see that the angular momentum of the initial systems looks as a doubled angular momentum of the transformed one
$$J=i(z\pi \overline{z}\overline{\pi })=2i(wp\overline{w}\overline{p}).$$
(2.70)
This indicates that the global properties of these two systems could be essentially different. This difference has to be reflected in the respective quantum-mechanical systems.
Let us consider the Schrรถdinger equation
$$(\pi ,\overline{\pi },z,\overline{z})\mathrm{\Psi }(z,\overline{z})=E\mathrm{\Psi }(z,\overline{z}),\pi =i_z,\overline{\pi }=i_{\overline{z}},$$
(2.71)
with the wavefunction which obeys the condition
$$\mathrm{\Psi }(|z|,\mathrm{arg}z+2\pi )=\mathrm{\Psi }(|z|,\mathrm{arg}z).$$
(2.72)
Let us reduce it by the action of $`Z_2`$ group, restricting ourselves to even ($`\sigma =0`$) or odd ($`\sigma =\frac{1}{2}`$) solutions of Eq. (2.71)
$$\mathrm{\Psi }_\sigma (z,\overline{z})=\psi _\sigma (z^2,\overline{z}^2)\mathrm{e}^{2i\sigma \mathrm{arg}z},\sigma =0,1/2,$$
(2.73)
and then perform the Bohlin transformation (2.68). According to Eq.(2.73), the wave functions $`\psi _\sigma `$ satisfy the condition
$$\psi _\sigma (|w|,\mathrm{arg}w+2\pi )=\psi _\sigma (|w|,\mathrm{arg}w),$$
(2.74)
which implies that the range of definition $`\mathrm{arg}w[0,4\pi )`$ can be restricted, without loss of generality, to $`\mathrm{arg}w[0,2\pi )`$. In terms of $`\psi _\sigma `$ the Schrรถdinger equation (2.71) reads
$$(\widehat{p}_\sigma ,\widehat{p}_\sigma ^+,w,\overline{w})\psi _\sigma (w,\overline{w})=E\psi _\sigma ,\widehat{p}_\sigma =i_w\frac{i\sigma }{w}.$$
(2.75)
Equation (2.75) can be interpreted as the Schrรถdinger equation of a particle with electric charge $`e`$ in the static magnetic field given by the potential $`A_w=\frac{i\sigma }{ew}`$, $`\sigma =0,1/2`$. It is a potential of an infinitely thin solenoid- โmagnetic flux tubeโ (or magnetic vortex, in the two-dimensional interpretation): it has zero strength of the magnetic field $`B=rotA_w=0`$ ($`w\dot{\text{I}\mathrm{C}}`$) and nonzero magnetic flux $`2\pi \sigma /e`$.
In accordance with (2.70), the angular momentum transforms as follows:
$`J`$ $`2J_\sigma ,J_\sigma ={\displaystyle \frac{i}{\mathrm{}}}\left(w\widehat{p}_\sigma \overline{w}\widehat{p}_\sigma ^+\right),`$ (2.76)
where $`J_\sigma `$ is the angular momentum operator of the reduced system. Hence, the eigenvalues of the angular momenta of the reduced and initial systems, $`m_\sigma `$ and $`M`$, are related by the expression $`M=2m_\sigma `$, from which it follows that
$$m_\sigma =\pm \sigma ,\pm (1+\sigma ),\pm (2+\sigma ),\mathrm{}.$$
(2.77)
Hence, the $`Z_2`$-reduction related to zero Hopf map transforms the even states of the initial system to the complete basis of the resulting one. The odd states of the initial system yield the wave functions of the resulting system in the presence of magnetic flux generating spin $`1/2`$. Similarly to the above consideration, one can show that the reduction of the two-dimensional system by the $`Z_N`$ group yields the $`N`$ systems with the fractional spin $`\sigma =0,1/N,2/N,\mathrm{},(N1)/N`$ (see ).
### $`1`$st Hopf map. Dirac monopole
Now we consider the Hamiltonian reduction by the action of the $`U(1)`$ group, which is associated with the first Hopf map. It is known under the name of Kustaanheimo-Stiefel transformation.
Let us consider the Hamiltonian system on the four-dimensional Hermitean space $`(M_0,g_{a\overline{b}}dz^ad\overline{z}^b)`$, $`\mathrm{dim}_{\text{I}\mathrm{C}}M_0=2`$,
$$T^{}M_0,\omega =dz^ad\pi _a+d\overline{z}^ad\overline{\pi }_a,=g^{a\overline{b}}\pi _a\overline{\pi }_b+V(z,\overline{z}).$$
(2.78)
We define, on the $`T^{}M_0`$ space, the Hamiltonian action of the $`U(2)`$ group given by the generators
$`๐=iz๐\pi i\overline{\pi }๐\overline{z},J_0=iz\pi i\overline{z}\overline{\pi }:`$ (2.79)
$`\{J_0,J_k\}=0,\{J_k,J_l\}=2ฯต_{klm}J_m,`$ (2.80)
where $`๐`$ are Pauli matrices.
Let us consider the Hamiltonian reduction of the phase space $`(T^{}M_0,\omega )`$ by the (Hamiltonian) action of the $`U(1)=S^1`$ group given by the generator $`J_0`$. Since $`J_0`$ commutes with $`J_i`$, the latter will generate the Hamiltonian action of the $`su(2)=so(3)`$ algebra on the reduced space as well.
In order to perform the Hamiltonian reduction, we have to fix the level surface
$$J_0=2s,$$
(2.81)
and then factorize it by the action of the vector field $`\{J_0,\}`$.
The resulting six-dimensional phase space $`T^{}M^{\mathrm{red}}`$ could be parameterized by the following $`U(1)`$-invariant functions:
$$๐ฒ=z๐\overline{z},๐
=\frac{z๐\pi +\overline{\pi }๐\overline{z}}{2z\overline{z}}:\{๐ฒ,J_0\}=\{๐
,J_0\}=0.$$
(2.82)
In these coordinates the reduced symplectic structure and the generators of the angular momentum are given by the expressions (compare with (1.17),(1.20))
$$\mathrm{\Omega }_{\mathrm{red}}=d๐
d๐ฒ+s\frac{๐ฒ(d๐ฒ\times d๐ฒ)}{2|๐ฒ|^3},๐_{red}=๐/2=๐
\times ๐ฒ+s\frac{๐ฒ}{|๐ฒ|}.$$
Hence, we get the phase space of the Hamiltonian system describing the motion of a nonrelativistic scalar particle in the magnetic field of the Dirac monopole.
Let $`M_0`$ be a $`U(2)`$-invariant Kรคhler space with a metric generated by the Kรคhler potential $`K(z\overline{z})`$
$$g_{a\overline{b}}=\frac{^2K(z\overline{z})}{z^a\overline{z}^b}=a(z\overline{z})\delta _{a\overline{b}}+a^{}(z\overline{z})\overline{z}^az^b,$$
(2.83)
where
$$a(y)=\frac{dK(y)}{dy},a^{}(y)=\frac{d^2K(y)}{dy^2}.$$
Let the potential be also $`U(2)`$-invariant, $`V=V(z\overline{z})`$, so that $`U(2)`$ is a symmetry of the Hamiltonian: $`\{J_0,\}=\{J_i,\}=0`$.
Hence, the Hamiltonian could also be restricted to the reduced six-dimensional phase space. The reduced Hamiltonian looks as follows:
$$_{red}=\frac{1}{a}\left[y๐
^2b(๐ฒ๐
)^2\right]+s^2\frac{1by}{ay}+V(y),$$
where
$$y|๐ฒ|,b=\frac{a^{}(y)}{a+ya^{}(y)}.$$
Let us perform the canonical transformation $`(๐ฒ,๐
)(๐ฑ,๐ฉ)`$ to the conformal-flat metric
$$๐ฑ=f(y)๐ฒ,๐
=f๐ฉ+\frac{df}{dy}\frac{(\mathrm{๐ฒ๐ฉ})}{y}๐ฒ,$$
where
$$\left(1+\frac{yf^{}(y)}{f}\right)^2=1+\frac{ya^{}(y)}{a}\left(\frac{d\mathrm{log}x}{dy}\right)^2=\frac{d\mathrm{log}ya(y)}{ydy},x<1.$$
In the new coordinates the Hamiltonian takes the form
$$_{red}=\frac{x^2(y)}{ya(y)}๐ฉ^\mathrm{๐}+\frac{s^2}{y(a+ya^{}(y))}+V\left(y(x)\right).$$
In order to express the $`y`$, $`a(y)`$, $`a^{}(y)`$ via $`x`$, it is convenient to introduce the function
$$\stackrel{~}{A}(y)(a+ya^{}(y))yf(y)๐y$$
and consider its Legendre transform $`A(x)`$,
$$A(x)=A(x,y)|_{A(x,y)/y},A(x,y)=xa(y)y\stackrel{~}{A}(y).$$
Then, we immediately get
$$\frac{dA(x)}{dx}=a(y)y,x\frac{d^2A}{dx^2}=y\sqrt{a(a+ya^{}(y))}.$$
By the use of these expressions, we can represent the reduced Hamiltonian as follows:
$$_{red}=\frac{x^2}{N^2}๐ฉ^2+\frac{s^2}{\left(2xN^{}(x)\right)^2}+V\left(y(x)\right),N^2(x)\frac{dA}{dx}.$$
(2.84)
The Kรคhler potential of the initial system is connected with $`N`$ via the equations
$$\frac{dK}{dx}=\frac{N^3(x)}{2x^2N^{}(x)},\frac{d\mathrm{log}y}{dx}=\frac{N}{2x^2N^{}(x)}.$$
(2.85)
Hence, for $`s=0`$ we shall get the system (1.27). However, when $`s0`$, by comparing the reduced system with (1.28), we conclude that the only Kรคhler space which yields a โwell-defined system with monopoleโ is flat space.
### $`\text{I}\mathrm{C}^{\mathrm{N}+1}\text{I}\mathrm{CP}^\mathrm{N}`$ and $`T^{}\text{I}\mathrm{C}^{\mathrm{N}+1}\mathrm{T}^{}\text{I}\mathrm{CP}^\mathrm{N}`$
Now, we consider the Hamiltonian reduction of the the space $`(\text{I}\mathrm{C}^{\mathrm{N}+1},\omega =\mathrm{du}^0\mathrm{d}\overline{\mathrm{u}}^0+\mathrm{du}^\mathrm{a}\mathrm{d}\overline{\mathrm{u}}^\mathrm{a}`$), to the complex projective space $`\text{I}\mathrm{CP}^\mathrm{N}`$.
The $`U(N+1)=U(1)\times SU(N)`$ isometries of this space are defined by the following Killing potentials:
$$J_0=u\overline{u},J_{su(N+1)}=u\widehat{T}\overline{u},\{J_0,J_{su(N+1)}\}=0,$$
where $`T=T^{}`$, $`\mathrm{Tr}T=0`$ are $`(N+1)\times (N+1)`$ dimensional traceless matrices defining the $`su(N+1)`$ algebra.
The Poisson brackets, corresponding to the Kรคhler structure, are defined by the relations $`\{u^0,\overline{u}^0\}=i`$, $`\{u,\overline{u}^b\}=i\delta ^{ab}`$.
Let us perform the Hamiltonian reduction by the action of $`J_0`$. The reduced phase space is a $`2N`$ dimensional one. Let us choose for this space the following local complex coordinates:
$$z^a=\frac{u^a}{u^0}:\{z^a,J_0\}=0,a=1,\mathrm{},N$$
(2.86)
and fix the level surface
$$J_0=r_0^2|u^0|^2=\frac{r_0^2}{1+z\overline{z}}.$$
(2.87)
Then, we immediately get the Poisson brackets for the reduced space
$$\{z^a,\overline{z}^b\}=\frac{i}{r_0^2}(1+z\overline{z})(\delta ^{ab}+z^a\overline{z}^b),\{z^a,z^b\}=\{\overline{z}^a,\overline{z}^b\}=0.$$
(2.88)
Hence, the reduced Poisson bracket are associated with the Kรคhler structure. It could be easily seen, that this Kรคhler structure is given by the Fubini-Study metric (1.40) multiplied on $`r_0^2`$. The restriction of the generators $`J_{su(N+1)}`$ on the level surface (2.87) yields the expressions (1.43).
In the above example $`\text{I}\mathrm{C}^{\mathrm{N}+1}`$ and $`\text{I}\mathrm{CP}^\mathrm{N}`$ appeared as the phase spaces. Now, let us show, how to reduce the $`T^{}\text{I}\mathrm{C}^{\mathrm{N}+1}`$ to $`T^{}\text{I}\mathrm{CP}^\mathrm{N}`$, i.e. let us consider the case when $`\text{I}\mathrm{C}^{\mathrm{N}+1}`$ and $`\text{I}\mathrm{CP}^\mathrm{N}`$ play the role of the configuration spaces of the mechanical systems. Since the dimension of $`T^{}\text{I}\mathrm{C}^{\mathrm{N}+1}`$ is $`4(N+1)`$, and the dimension of $`T^{}\text{I}\mathrm{C}^\mathrm{N}`$ is $`4N`$, the reduction has to be performed by two commuting generators.
Let us equip the initial space with the canonical symplectic structure (2.78), and perform the reduction of this phase space by the action of the generators
$$J_0=i\pi u\overline{\pi }\overline{u},h_0=u\overline{u}:\{J_0,h_0\}=0.$$
(2.89)
We choose the following local coordinates of the reduced space:
$$z^a=\frac{u^a}{u^0},p_a=g_{a\overline{b}}(z,\overline{z})\left(\frac{\overline{\pi }^a}{\overline{u}^0}\overline{z}^a\frac{\overline{\pi }^0}{\overline{z}^0}\right):$$
$$\{z^a,J_0\}=\{z^a,h_0\}=\{p_a,J_0\}=\{p_a,h_0\}=0,$$
where $`g_{a\overline{b}}`$ is defined by the expression (1.40). Then, calculating the Poisson brackets between these functions, and fixing the value of the generators $`J_0,h_0`$,
$$h_0=r_0^2,J_0=2s,$$
(2.90)
we get
$$\{p_a,z^b\}=\delta _a^b,\{p_a,\overline{p}_b\}=i\frac{s}{r_0^2}g_{a\overline{b}}(z,\overline{z}).$$
(2.91)
Hence, we arrive at the phase space structure of the particle moving on $`\text{I}\mathrm{CP}^\mathrm{N}`$ in the presence of a constant magnetic field with $`B_0=s/r_0^2`$ strength.
### $`2`$nd Hopf map. $`SU(2)`$ instanton
In the above examples we have shown that the zero Hopf map is related to the canonical transformation corresponding to the reduction of the two-dimensional system by the discrete group $`Z_2=S^0`$, and transforms the system with two-dimensional configuration space to the system of the same dimension, which has a spin $`\sigma =0,1/2`$. The first Hopf map corresponds to the reduction of the system with four-dimensional configuration space by the Hamiltonian action of $`U(1)=S^1`$ group, and yields the system moving on the three-dimensional space in the presence of the magnetic field of the Dirac monopole. Similarly, with the second Hopf map one can relate the Hamiltonian reduction of the cotangent bundle of eight-dimensional space (say, $`T^{}\text{I}\mathrm{C}^4=\mathrm{T}^{}\mathrm{IHI}^2`$) by the action of $`SU(2)=S^3`$ group. When the $`SU(2)`$ generators $`I_i`$ have non-zero values, $`I_i=c_i`$,$`_i|c_i|0`$, the reduced space is a $`(2831=)12`$\- dimensional one, $`T^{}\mathrm{IR}^5\times \mathrm{S}^2`$. It is the phase space of a coloured particle moving on $`\mathrm{IR}^5`$ in the presence of the $`SU(2)`$ Yang monopole (here $`S^2`$ appears as a isospin space).
When $`c_1=c_2=c_3=0`$, the $`J_i`$ generators commute with each other, and the reduced space is a $`(2823=)10`$-dimensional one, $`T^{}\mathrm{IR}^5`$. Such a reduction is also known under the name of Hurwitz transformation relating the eight-dimensional oscillator with the five-dimensional Coulomb problem.
We shall describe a little bit different reduction, associated with the fibration $`\text{I}\mathrm{CP}^3/\text{I}\mathrm{CP}^1=\mathrm{S}^4`$ . This fibration could be immediately obtained by factorization of the second Hopf map $`S^7/S^3=S^4`$ by $`U(1)`$. Indeed, the second Hopf map is described by the formulae (1.48),(1.49), where $`S^7`$ is embedded in the two-dimensional quaternionic space $`\mathrm{IHI}^2=\text{I}\mathrm{C}^4`$, parameterized by four complex (two quaternionic) Euclidean coordinates
$$u_i=v_i+jv_{i+1},i=1,2,u_1,๐ฎ_2\mathrm{IHI},\mathrm{v}_1,\mathrm{v}_2,\mathrm{v}_3,\mathrm{v}_4\text{I}\mathrm{C}.$$
(2.92)
Here $`S^4`$ is embedded in $`\mathrm{IR}^5`$ parameterized by the Eucludean coordinates ($`๐ฑ,x_5`$) given by (1.48). This embedding is invariant under the right action of a $`SU(2)`$ group given by (1.49), so that $`๐ `$ defines a three-sphere (1.52). The complex projective space $`\text{I}\mathrm{CP}^3`$ is defined as $`S^7/U(1)`$, while the inhomogeneous coordinates $`z_a`$ appearing in the Fubini-Study metric of $`\text{I}\mathrm{CP}^3`$, are related to the coordinates of $`\text{I}\mathrm{C}^4`$ as follows: $`z_a=v_a/v_4`$, $`a=1,2,3`$. The expressions (1.48) defining $`S^4`$ are invariant under $`U(1)`$-factorization, while $`S^3/U(1)=S^2`$. Thus, we arrive to the conclusion that $`\text{I}\mathrm{CP}^3`$ is the $`S^2`$-fibration over $`S^4=\mathrm{IHIP}^1`$. The expressions for $`z_a`$ yield the following definition of the coordinates of $`S^4`$:
$$w_1=\frac{\overline{z}_2+z_1\overline{z}_3}{1+z_3\overline{z}_3},w_2=\frac{z_2\overline{z}_3\overline{z}_1}{1+z_3\overline{z}_3}.$$
(2.93)
Choosing $`z_3`$ as a local coordinate of $`S^2=\text{I}\mathrm{CP}^1`$,
$$u=z_3,$$
(2.94)
we get the expressions
$$z_1=w_1u\overline{w}_2,z_2=w_2u+\overline{w}_1,z_3=u.$$
(2.95)
In these coordinates the Fubini-Study metric on $`\text{I}\mathrm{CP}^3`$ looks as follows:
$$g_{a\overline{b}}dz_ad\overline{z}_b=\frac{dzd\overline{z}}{1+z\overline{z}}\frac{(\overline{z}dz)(zd\overline{z})}{(1+z\overline{z})^2}=\frac{dw_id\overline{w}_i}{(1+w\overline{w})^2}+\frac{(du+๐)(d\overline{u}+\overline{๐})}{(1+u\overline{u})^2},$$
(2.96)
where
$$๐=\frac{(\overline{w}_1+w_2u)(udw_1d\overline{w}_2)+(\overline{w}_2w_1u)(udw_2+d\overline{w}_1)}{1+w\overline{w}}.$$
(2.97)
Hence, $`w_1,w_2`$ and $`u`$ are the conformal-flat complex coordinates of $`S^4=\mathrm{IHIP}^1`$ and $`S^2=\text{I}\mathrm{CP}^1`$, while the connection $`๐`$ defines the $`SU(2)`$ gauge field.
Now, let us consider the Hamiltonian system describing the motion of a free particle on $`\text{I}\mathrm{CP}^3`$
$$_{\text{I}\mathrm{CP}^3}=g^{a\overline{b}}\pi _a\overline{\pi }_b,\{z_a,\pi _b\}=i\delta _{ab}$$
(2.98)
Let us extend the coordinate transformation (2.95) to the $`T^{}\text{I}\mathrm{CP}^3`$, by the following transformation of momenta:
$$\pi _1=\frac{\overline{u}p_1\overline{p}_2}{1+u\overline{u}},\pi _2=\frac{\overline{u}p_2+p_1}{1+u\overline{u}},$$
$$\pi _3=p_u+\frac{\overline{p}_2w_1\overline{p}_1w_2\overline{u}(w_1p_1+w_2p_2)}{1+u\overline{u}}.$$
(2.99)
This extended transformation is a canonical transformation,
$$\{w_i,p_j\}=\delta _{ij},\{u,p_u\}=1.$$
(2.100)
In the new coordinates the Hamiltonian reads
$$_{\text{I}\mathrm{CP}^3}=(1+w\overline{w})^2P_i\overline{P}_i+(1+u\overline{u})^2p_u\overline{p}_u.$$
(2.101)
Here we introduced the covariant momenta
$$P_1=p_1i\frac{\overline{w}_1}{1+w\overline{w}}I_1\frac{w_2}{1+w\overline{w}}I_+,P_2=p_2i\frac{\overline{w}_2}{1+w\overline{w}}I_1+\frac{w_1}{1+w\overline{w}}I_+,$$
(2.102)
and the $`su(2)`$ generators $`I_\pm ,I_1`$ defining the isometries of $`S^2`$
$$\begin{array}{c}I_1=i(p_uu\overline{p}_u\overline{u}),I_{}=p_u+\overline{u}^2\overline{p}_{\overline{u}},I_+=\overline{p}_{\overline{u}}+u^2p_u\\ \{I_\pm ,I_1\}=iI_\pm ,\{I_+,I_{}\}=2iI_1.\end{array}$$
(2.103)
The nonvanishing Poisson brackets between $`P_i`$, $`w_i`$ are given by the following relations (and their complex conjugates):
$$\{w_i,P_j\}=\delta _{ij},\{P_1,P_2\}=\frac{2I_+}{(1+w\overline{w})^2},\{P_i,\overline{P}_j\}=i\frac{2I_1\delta _{ij}}{(1+w\overline{w})^2}.$$
(2.104)
The expressions in the r.h.s. define the strength of a homogeneous $`SU(2)`$ instanton (the โangular partโ of the $`SU(2)`$ Yang monopole), written in terms of conformal-flat coordinates of $`S^4=\mathrm{IHIP}^1`$. Hence, the first part of the Hamiltonian, i.e. $`๐_4=(1+w\overline{w})^2P_i\overline{P}_i`$, describes a particle on the four-dimensional sphere in the field of a $`SU(2)`$ instanton.
The Poisson brackets between $`P_i`$ and $`u,\overline{u},p_u,\overline{p}_u`$ are defined by the following nonzero relations and their complex conjugates:
$$\{P_i,p_u\}=\frac{\overline{w_i}+2ฯต_{ij}w_ju}{1+\overline{w}w}p_u,\{P_i,\overline{p}_u\}=\frac{\overline{w}_i\overline{p}_u}{1+\overline{w}w},$$
$$\{P_i,u\}=\frac{\left(\overline{w_i}+ฯต_{ij}w_ju\right)u}{1+\overline{w}w},\{\overline{P_i},u\}=\frac{ฯต_{ij}\overline{w_j}w_iu}{1+\overline{w}w}.$$
The second part of the Hamiltonian defines the motion of a free particle on the two-sphere. It could be represented as a Casimir of $`SU(2)`$
$$๐_{S^2}=(1+u\overline{u})^2p_u\overline{p}_u=I_+I_{}+I_1^2I^2.$$
(2.105)
It commutes with the Hamiltonian $`๐_0`$, as well as with $`I_1,I_\pm `$ and $`P_i,w_i`$
$$\{๐_{\text{I}\mathrm{CP}^3},I^2\}=\{P_i,I^2\}_B=\{w_i,I^2\}_B=\{_1,I^2\}_B=\{_\pm ,I^2\}_B=0.$$
(2.106)
Hence, we can perform a Hamiltonian reduction by the action of the generator $`๐_2`$, which reduces the initial twelve-dimensional phase space $`T_{}\text{I}\mathrm{CP}^3=\mathrm{T}^{}(\mathrm{S}^4\times \mathrm{S}^2)`$ to a ten-dimensional one. The relations (2.106) allow us to parameterize the reduced ten-dimensional phase space in terms of the coordinates $`P_i,w_i,I_\pm ,I_1`$, where the latter obey the relation
$$I_+I_{}+I_1^2I^2=const.$$
(2.107)
Thus, the reduced phase space is nothing but $`T^{}S^4\times S^2`$, where $`S^2`$ is the internal space of the instanton.
Let us collect the whole set of non-zero expressions defining the Poisson brackets on $`T_{}S^4\times S^2`$
$$\{w_i,P_j\}=\delta _{ij},$$
$$\{P_1,P_2\}=\frac{2I_+}{(1+w\overline{w})^2},$$
$$\{P_i,\overline{P}_j\}=i\frac{2I_1\delta _{ij}}{(1+w\overline{w})^2},$$
$$\{P_i,I_1\}=i\frac{ฯต_{ij}w_jI_+}{1+w\overline{w}}$$
(2.108)
$$\{P_i,I_+\}=\frac{\overline{w_i}I_+}{1+\overline{w}w},$$
$$\{P_i,I_{}\}=\frac{\overline{w_i}I_{}+2iฯต_{ij}w_jI_1}{1+\overline{w}w}$$
$$\{I_+,I_{}\}=2iI_1,\{I_\pm ,I_1\}=iI_\pm .$$
The reduced Hamiltonian is $`_{\text{I}\mathrm{CP}^3}^{red}=(1+w\overline{w})^2P\overline{P}+I^2.`$ So, the Hamiltonian of the coloured particle on $`S^4`$ interacting with the $`SU(2)`$ instanton is connected with the Hamiltonian of a particle on $`\text{I}\mathrm{CP}^3`$ as follows:
$$๐_{S^4}=๐_{\text{I}\mathrm{CP}^3}^{\mathrm{red}}I^2(>0).$$
(2.109)
This yields an intuitive explanation of the degeneracy in the ground state in the corresponding quantum system on $`S^4`$. Indeed, since the l.h.s. is positive, the ground state of the quantum system on $`S^4`$ corresponds to the excited state of a particle on $`\text{I}\mathrm{CP}^3`$, which is a degenerate one. On the other hand, the ground state of a particle on $`\text{I}\mathrm{CP}^3`$ can be reduced to the free particle on $`S^4`$, when $`I=0`$.
Now, let us consider a similar reduction for the particle on $`\text{I}\mathrm{CP}^3`$, in the presence of constant magnetic field (1.36).
Passing to the coordinates (2.95) and momenta (2.102) we get the Poisson brackets defined by the nonzero relations given by (2) and
$$\{p_u,\overline{p}_u\}_B=\frac{iB}{(1+u\overline{u})^2},$$
(2.110)
$$\{w_i,P_j\}_B=\delta _{ij},\{P_1,P_2\}_B=\frac{2_+}{(1+w\overline{w})^2},$$
(2.111)
$$\{P_i,\overline{P}_j\}_B=i\frac{2_1\delta _{ij}}{(1+w\overline{w})^2}.$$
(2.112)
where $`_\pm ,_1`$ are defined by the expressions
$$_1=I_1+\frac{B}{2}\frac{1u\overline{u}}{1+u\overline{u}},_{}=I_{}B\frac{i\overline{u}}{1+u\overline{u}},_+=I_++B\frac{iu}{1+u\overline{u}}$$
(2.113)
Notice that the expressions (2.112) are similar to (2.104) and the generators (2.113) form, with respect to the new Poisson brackets, the $`su(2)`$ algebra
$$\{_\pm ,_1\}_B=i_\pm ,\{_+,_{}\}=2i_1.$$
(2.114)
It is clear that these generators define the isometries of the โinternalโ two-dimensional sphere with a magnetic monopole located at the center.
Once again, as in the absence of a magnetic field, we can reduce the initial system by the Casimir of the $`SU(2)`$ group
$$^2_1^2+_+_{}=๐_{S^2}+B^2/4,B/2.$$
(2.115)
In order to perform the Hamiltonian reduction, we have to fix the value of $`^2`$, and then factorize by the action of the vector field $`\{^2,\}_B`$.
The coordinates (2.93), (2.102) commute with the Casimir (2.115),
$$\{P_i,^2\}_B=\{w_i,^2\}_B=\{_1,^2\}_B=\{_\pm ,^2\}_B=0.$$
(2.116)
Hence, as we did above, we can choose $`P_i`$, $`w_i`$, and $`_\pm `$ as the coordinates of the reduced, ten-dimensional phase space.
The coordinates $`_\pm `$, $`_1`$ obey the condition
$$_1^2+_{}_+=^2=const.$$
(2.117)
The resulting Poisson brackets are defined by the expressions (2.108), with $`I_1,I_\pm `$ replaced by $`_\pm ,_1`$.
Hence, the particle on $`\text{I}\mathrm{CP}^3`$ moving in the presence of a constant magnetic field reduces to a coloured particle on $`S^4`$ interacting with the instanton field. The Hamiltonians of these two systems are related as follows:
$$๐_{S^4}=๐_{\text{I}\mathrm{CP}^3}^{red}^2+B^2/4,B/2$$
(2.118)
Notice that, upon quantization, we must replace $`^2`$ by $`(+1)`$ and require that both $``$ and $`B`$ take (half)integer values (since we assume unit radii for the spheres, this means that the โmonopole numberโ obeys a Dirac quantization rule). The extension of this reduction to quantum mechanics relates the theories of the quantum Hall effect on $`S^4`$ and on $`\text{I}\mathrm{CP}^3`$ .
Notice that the third Hopf map could also be related with the generalized quantum Hall effect theory .
## 3 Generalized oscillators
Among the integrable systems with hidden symmetries the oscillator is the simplest one. In contrast to other systems with hidden symmetries (e.g. Coulomb systems), its symmetries form a Lie algebra. The $`N`$-dimensional oscillator on $`T^{}\mathrm{IR}^\mathrm{N}`$,
$$=\frac{1}{2}\left(p_ap_a+\alpha ^2q^aq^a\right),\omega _{can}=dp_adq^a,a=1,\mathrm{},N$$
(3.119)
besides the rotational symmetry $`so(N)`$, has also hidden ones, so that the whole symmetry algebra is $`su(N)`$. The symmetries of the oscillator are given by the generators
$$J_{ab}=p_aq^bp_bq^a,I_{ab}=p_ap_b+\alpha ^2q^aq^b.$$
(3.120)
The huge number of hidden symmetries allows us to construct generalizations of the oscillator on curved spaces, which inherit many properties of the initial system.
The generalization of the oscillator on the sphere was suggested by Higgs . It is given by the following Hamiltonian system:
$$=\frac{1}{2}g^{ab}p_ap_a+\frac{\alpha ^2}{2}q^aq^a,\omega =dp_adq^a,q^a=\frac{x_a}{x_0},$$
(3.121)
where $`x^a,x_0`$ are the Euclidean coordinates of the ambient space $`\mathrm{IR}^{\mathrm{N}+1}`$: $`x_0^2+x^ax^a=1`$, and $`g_{ab}dq^adq^b`$ is the metric on $`S^N`$. This system inherits the rotational symmetries of the flat oscillator given by (3.120), and possesses the hidden symmetries given by the following constants of motion (compare with (3.120)):
$$I_{ab}=J_aJ_b+\alpha ^2q^aq^b,$$
(3.122)
where $`J_a`$ are the translation generators on $`S^N`$.
In contrast to the flat oscillator, whose symmetry algebra is $`su(N)`$, the spherical (Higgs) oscillator has a nonlinear symmetry algebra.
This construction has been extended to the complex projective spaces in Ref. , where the oscillator on $`\text{I}\mathrm{CP}^\mathrm{N}`$ was defined by the Hamiltonian
$$=g^{\overline{a}b}\overline{\pi }_a\pi _b+\alpha ^2z\overline{z},$$
(3.123)
with $`z^a=u^a/u^0`$ denoting inhomogeneous coordinates of $`\text{I}\mathrm{CP}^\mathrm{N}`$ and $`g_{a\overline{b}}dz^ad\overline{z}^b`$ being Fubini-Study metric (1.40).
It is easy to see that this system has constants of motion given by the expressions
$$J_{a\overline{b}}=i(z^b\pi _a\overline{\pi }_b\overline{z}^a),I_{a\overline{b}}=J_a^+J_b^{}+\omega ^2\overline{z}^az^b,$$
(3.124)
where $`J_a^+=\pi _a+(\overline{z}\overline{\pi })\overline{z}^a`$, $`J_a^{}=\overline{J}_a^+`$ are the translation generators on $`\text{I}\mathrm{CP}^\mathrm{N}`$. The generators $`J_{a\overline{b}}`$ define the kinematical symmetries of the system and form a $`su(N)`$ algebra. When $`N>1`$, the generators $`I_{a\overline{b}}`$ are functionally independent of $``$, $`J_{a\overline{b}}`$ and define hidden symmetries. As in the spherical case, their algebra is a nonlinear one
$$\begin{array}{c}\{J_{\overline{a}b},J_{\overline{c}d}\}=i\delta _{\overline{a}d}J_{\overline{b}c}i\delta _{\overline{c}b}J_{\overline{a}d},\\ \{I_{a\overline{b}},J_{c\overline{d}}\}=i\delta _{c\overline{b}}I_{a\overline{d}}i\delta _{a\overline{d}}I_{c\overline{b}}\\ \{I_{a\overline{b}},I_{c\overline{d}}\}=i\alpha ^2\delta _{c\overline{b}}J_{a\overline{d}}i\alpha ^2\delta _{a\overline{d}}J_{c\overline{b}}+\\ +iI_{c\overline{b}}(J_{a\overline{d}}+J_0\delta _{a\overline{d}})iI_{a\overline{d}}(J_{c\overline{b}}+J_0\delta _{c\overline{b}}).\end{array}$$
(3.125)
Hence, it is seen that for $`N=1`$, i.e. in the case of the two-dimensional sphere $`S^2=\text{I}\mathrm{CP}^1`$, the suggested system has no hidden symmetries, as opposed to the Higgs oscillator on $`S^2`$. Nevertheless, this model is exactly solvable both for $`N=1`$ and $`N>1`$ . Moreover, it remains exactly solvable, even after inclusion of a constant magnetic field, for any $`N`$ (including $`N=1`$, when it has no hidden symmetries). The magnetic field does not break the symmetry algebra of the system! As opposed to the described model, the constant magnetic field breaks the hidden symmetries, as well as the exact solvability, of the Higgs oscillator on $`S^2=\text{I}\mathrm{CP}^1`$.
Remark.The Hamiltonian (3.123) could be represented as follows:
$$=g^{a\overline{b}}(\pi _a\overline{\pi }_b+\alpha ^2_aK\overline{}_bK),$$
(3.126)
where $`K(z,\overline{z})=\mathrm{log}(1+z\overline{z})`$ is the Kรคhler potential of the Fubini-Study metric.
Although this potential is not uniquely defined, it provides the system with some properties, which are general for the few oscillator models on Kรคhler spaces. By this reason we postulate it as an oscillator potential on arbitrary Kรคhler manifolds.
Now, let us compare these systems with the sequence which we like: real, complex, quaternionic numbers (and zeroth, first, second Hopf map). Let us observe, that the $`S^N`$-oscillator potential is defined, in terms of the ambient space $`\mathrm{IR}^{\mathrm{N}+1}`$, in complete similarity to the $`\text{I}\mathrm{CP}^\mathrm{N}`$-oscillator potential in terms of the โambientโ space $`\text{I}\mathrm{C}^{\mathrm{N}+1}`$. The latter system preserves its exact solvability in the presence of a constant magnetic ($`U(1)`$ gauge) field.
Hence, continuing this sequence, one can define on the quaternionic projective spaces $`\mathrm{IHIP}^\mathrm{N}`$ the oscillator-like system given by the potential
$$V_{\mathrm{IHIP}^\mathrm{N}}=\alpha ^2w^a\overline{w}^a=\alpha ^2\frac{u_1^a\overline{u}_1^a+u_2^a\overline{u}_2^a}{u_1^0\overline{u}_1^0+u_1^0\overline{u}_1^0},$$
(3.127)
where
$$w^a=\frac{u_1^a+ju_2^a}{u_1^0+ju_2^0},u_1^a\overline{u}_1^a+u_2^a\overline{u}_2^a+u_1^0\overline{u}_1^0+u_2^0\overline{u}_2^0=1.$$
Here $`w^a`$ are inhomogeneous (quaternionic) coordinates of the quaternionic projective space $`\mathrm{IHIP}^\mathrm{N}`$, and $`u_0^a+ju_1^a,u_1^0+ju_2^0`$ are the Euclidean coordinates of the โambientโ quaternionic space $`\mathrm{IHI}^{\mathrm{N}+1}=\text{I}\mathrm{C}^{2\mathrm{N}+2}`$.
One can expect that this system will be a superintegrable one and will be exactly solvable also in the presence of a $`SU(2)`$ instanton field.
In the simplest case of $`\mathrm{IHIP}^1=\mathrm{S}^4`$ we shall get the alternative (with respect to the Higgs) model of the oscillator on the four-dimensional sphere. In terms of the ambient space $`\mathrm{IR}^5`$, its potential will be given by the expression
$$V_{S^4}=\alpha ^2\frac{1x^0/x}{1+x^0/x}=\alpha ^2\frac{1\mathrm{cos}\theta }{1+\mathrm{cos}\theta }.$$
(3.128)
Checking this system for this simplest case, we found, that it is indeed exactly solvable in the presence of the instanton field .
Let us mention that the Higgs (spherical) oscillator could be straightforwardly extended to (one- and two-sheet) hyperboloids, and the $`\text{I}\mathrm{CP}^\mathrm{N}`$-oscillator - to the Lobachevsky spaces $`_N=SU(N+1)/U(N)`$. In both cases these systems have hidden symmetries.
Notice also that, on the spheres $`S^N`$, there exists the analog of the Coulomb system suggested by Schrรถdinger . It is given by the potential
$$V_{Coulomb}=\frac{\gamma }{r_0}\frac{y_{N+1}}{|๐ฒ|},y_{N+1}^2+|๐ฒ|^2=r_0^2.$$
(3.129)
This system inherits the hidden symmetry of the conventional Coulomb system on $`\mathrm{IR}^\mathrm{N}`$.
Probably, as in the case of the oscillator, one can define superintegrable analogs of the Coulomb system on the complex projective spaces $`\text{I}\mathrm{CP}^\mathrm{N}`$ and on the quaternionic projective spaces $`\mathrm{IHIP}^\mathrm{N}`$. However, up to now, this question has not been analyzed.
### Relation of the (pseudo)spherical oscillator and Coulomb systems
The oscillator and Coulomb systems, being the best known among the superintegrable mechanical systems, possess many similarities both at the classical and quantum mechanical levels. Writing down these systems in spherical coordinates, one can observe that the radial Schrรถdinger equation of the $`(p+1)`$-dimensional Coulomb system could be transformed in the Shrรถdinger equation of the $`2p`$-dimensional oscillator by the transformation (see, e.g. )
$$r=R^2,$$
where $`r`$ and $`R`$ are the radial coordinates of the Coulomb and oscillator systems, respectively.
Due to the existence of the Hopf maps, in the cases of $`p=1,2,4`$ one can establish a complete correspondence between these systems. Indeed, their angular parts are, respectively, $`p`$\- and $`(2p1)`$-dimensional spheres, while the above relation follows immediately from (1.54). Considering the Hamiltonian reductions related to the Hopf maps (as it was done in the previous section), one can deduce, that the $`(p+1)`$-dimensional Coulomb systems could be obtained from the $`2p`$\- dimensional oscillator, by a reduction under the $`G=S^{(p1)}`$ group. Moreover, for non-zero values of those generators we shall get generalizations of the Coulomb systems, specified by the presence of a magnetic flux ($`p=1`$), a Dirac monopole ($`p=2`$), a Yang monopole ($`p=4`$) . However, this procedure assumes a change in the roles of the coupling constants and the energy. To be more precise, these reductions convert the energy surface of the oscillator in the energy surface of the Coulomb-like system, while there is no one-to-one correspondence between their Hamiltonians.
As we have seen above, there exists well-defined generalizations of the oscillator systems on the spheres, hyperboloids, complex projective spaces and Lobachevsky spaces. The Coulomb system could also be generalized on the spheres and hyperboloids. Hence, the following natural question arises. Is it possible to relate the oscillator and Coulomb systems on the spheres and hyperboloids, similarly to those in the flat cases? The answer is positive, but it is rather strange. The oscillators on the $`2p`$-dimensional sphere and two-sheet hyperboloid (pseudosphere) result in the Coulomb-like systems on the $`(p+1)`$-dimensional pseudosphere, for $`p=1,2,4`$ .
Below, following , we shall show how to relate the oscillator and Coulomb systems on the spheres and two-sheet hyperboloids. In the planar limit this relation results in the standard correspondence between the conventional (flat) oscillator and the Coulomb-like system. We shall discuss mainly the $`p=1`$ case, since the treatment could be straightforwardly extended to the $`p=2,4`$ cases.
Let us introduce the complex coordinate $`z`$ parameterizing the sphere by the complex projective plane $`\text{I}\mathrm{CP}^1`$ and the two-sheeted hyperboloid by the Poincarรฉ disk (Lobachevsky plane, pseudosphere) $`)`$
$$๐ฑx_1+ix_2=R_0\frac{2z}{1+ฯตz\overline{z}},x_3=R_0\frac{1ฯตz\overline{z}}{1+ฯตz\overline{z}}.$$
(3.130)
In these coordinates the metric becomes conformally-flat
$$ds^2=R_0^2\frac{4dzd\overline{z}}{(1+ฯตz\overline{z})^2}.$$
(3.131)
Here $`ฯต=1`$ corresponds to the system on the sphere, and $`ฯต=1`$ to that on the pseudosphere. The lower hemisphere and the lower sheet of the hyperboloid are parameterized by the unit disk $`|z|<1`$, while the upper hemisphere and the upper sheet of the hyperboloid are specified by $`|z|>1`$, and transform one into another by the inversion $`z1/z`$. In the limit $`R_0\mathrm{}`$ the lower hemisphere (the lower sheet of the hyperboloid) turns into the whole two-dimensional plane. In these terms the oscillator and Coulomb potentials read
$$V_{osc}=\frac{2\alpha ^2R_0^2z\overline{z}}{(1ฯตz\overline{z})^2},V_C=\frac{\gamma }{R_0}\frac{1ฯตz\overline{z}}{2|z|},$$
(3.132)
Let us equip the oscillator phase space $`T^{}\text{I}\mathrm{CP}^1`$ ($`T^{})`$ with the symplectic structure
$$\omega =d\pi dz+d\overline{\pi }d\overline{z}$$
(3.133)
and introduce the rotation generators defining the $`su(2)`$ algebra for $`ฯต=1`$ and the $`su(1.1)`$ algebra for $`ฯต=1`$
$$๐\frac{iJ_1J_2}{2}=\pi +ฯต\overline{z}^2\overline{\pi },J\frac{ฯตJ_3}{2}=i(z\pi \overline{z}\overline{\pi }).$$
(3.134)
These generators, together with $`๐ฑ/R_0,x_3/R_0`$, define the algebra of motion of the (pseudo)sphere via the following non-vanishing Poisson brackets:
$$\begin{array}{c}\{๐,๐ฑ\}=2x_3,\{๐,x_3\}=ฯต\overline{๐ฑ},\{J,๐ฑ\}=i๐ฑ,\\ \{๐,\overline{๐}\}=2iฯตJ,\{๐,J\}=i๐.\end{array}$$
(3.135)
In these terms, the Hamiltonian of a free particle on the (pseudo)sphere reads
$$H_0^ฯต=\frac{๐\overline{๐}+ฯตJ^2}{2R_0^2}=\frac{(1+ฯตz\overline{z})^2\pi \overline{\pi }}{2R_0^2},$$
(3.136)
whereas the oscillator Hamiltonian is given by the expression
$$H_{osc}^ฯต(\alpha ,R_0|\pi ,\overline{\pi },z,\overline{z})=\frac{(1+ฯตz\overline{z})^2\pi \overline{\pi }}{2R_0^2}+\frac{2\alpha ^2R_0^2z\overline{z}}{(1ฯตz\overline{z})^2}.$$
(3.137)
It can be easily verified that the latter system possesses the hidden symmetry given by the complex (or vectorial) constant of motion
$$๐=I_1+iI_2=\frac{๐^2}{2R_0^2}+\frac{\alpha ^2R_0^2}{2}\frac{\overline{๐ฑ}^2}{x_3^2},$$
(3.138)
which defines, together with $`J`$ and $`H_{osc}`$ , the cubic algebra
$$\{๐,J\}=2i๐,\{\overline{๐},๐\}=4i\left(\alpha ^2J+\frac{ฯตJH_{osc}}{R_0^2}\frac{J^3}{2R_0^4}\right).$$
(3.139)
The energy surface of the oscillator on the (pseudo)sphere $`H_{osc}^ฯต=E`$ reads
$$\frac{\left(1(z\overline{z})^2\right)^2\pi \overline{\pi }}{2R_0^4}+2\left(\alpha ^2+ฯต\frac{E}{R_0^2}\right)z\overline{z}=\frac{E}{R_0^2}\left(1+(z\overline{z})^2\right).$$
(3.140)
Now, performing the canonical Bohlin transformation (2.68) one can rewrite the expression (3.140) as follows:
$$\frac{(1w\overline{w})^2p\overline{p}}{2r_0^2}\frac{\gamma }{r_0}\frac{1+w\overline{w}}{2|w|}=_C,$$
(3.141)
where we introduced the notation
$$r_0=R_0^2,\gamma =\frac{E}{2},2_C=\alpha ^2+ฯต\frac{E}{r_0}.$$
(3.142)
Comparing the l.h.s. of (3.141) with the expressions (3.132), (3.136) we conclude that (3.141) defines the energy surface of the Coulomb system on the pseudosphere with โradiusโ $`r_0`$, where $`w,p`$ denote the complex stereographic coordinate and its conjugated momentum, respectively. In the above, $`r_0`$ is the โradiusโ of the pseudosphere, while $`_C`$ is the energy of the system. Hence, we related classical isotropic oscillators on the sphere and pseudosphere with the classical Coulomb problem on the pseudosphere.
The constants of motion of the oscillators, $`J`$ and $`๐`$ (which coincide on the energy surfaces (3.140)) are converted, respectively, into the doubled angular momentum and the doubled Runge-Lenz vector of the Coulomb system
$$J2J_C,๐2๐,๐=\frac{iJ_C๐_C}{r_0}+\gamma \frac{\overline{๐ฑ}_C}{|๐ฑ_C|},$$
(3.143)
where $`๐_C`$, $`J_C`$, $`๐ฑ_C`$ denote the rotation generators and the pseudo-Euclidean coordinates of the Coulomb system.
We have shown above that, for establishing the quantum-mechanical correspondence, we have to supplement the quantum-mechanical Bohlin transformation with the reduction by the $`Z_2`$ group action, choosing either even ($`\sigma =0`$) or odd ($`\sigma =1/2`$) wave functions (2.73). The resulting Coulomb system is spinless for $`\sigma =0`$, and it possesses spin $`1/2`$ for $`\sigma =1/2`$.
The presented construction could be straightforwardly extended to higher dimensions, concerning the $`2p`$dimensional oscillator on the (pseudo)sphere and the $`(p+1)`$dimensional Coulomb-like systems, $`p=2,4`$. It is clear, that the $`p=2`$ case corresponds to the Hamiltonian reduction, associated with the first Hopf map, and the $`p=4`$ case is related to the second Hopf map. Indeed, the oscillator on the $`2p`$-dimensional (pseudo)sphere is also described by the Hamiltonian (3.137), where the following replacement is performed: $`(z,\pi )(z^a,\pi _a)`$, $`a=1,\mathrm{},p`$, with the summation over these indices understood. Consequently, the oscillator energy surfaces are again given by Eq. (3.140). Then, performing the Hamiltonian reduction, associated with the $`p`$-th Hopf maps (see the previous Section) we shall get the Coulomb-like system on the $`(p+1)`$-dimensional pseudosphere.
For example, if $`p=2`$, we reduce the system under consideration by the Hamiltonian action of the $`U(1)`$ group given by the generator $`J=i(z\pi \overline{z}\overline{\pi })`$. This reduction was described in detail in Section 2. For this purpose, we have to fix the level surface $`J=2s`$ and choose the $`U(1)`$-invariant stereographic coordinates in the form of the conventional Kustaanheimo-Stiefel transformation (2.82). The resulting symplectic structure takes the form (1.17). The oscillator energy surface reads
$$\frac{(1๐ช^2)^2}{8r_0^2}(๐ฉ^2+\frac{s^2}{๐ช^2})\frac{\gamma }{r_0}\frac{1+๐ช^2}{2|๐ช|}=_C,$$
(3.144)
where $`r_0`$, $`\gamma `$, $`_C`$ are defined by the expressions (3.142).
Interpreting $`๐ช`$ as the (real) stereographic coordinates of the three-dimensional pseudosphere
$$๐ฑ=r_0\frac{2๐ช}{1๐ช^2},x_4=r_0\frac{1+๐ช^2}{1๐ช^2},$$
(3.145)
we conclude that (3.144) defines the energy surface of the pseudospherical analog of a Coulomb-like system proposed in Ref. , which is also known under the name of โMIC-Keplerโ system.
In the $`p=4`$ case, we have to reduce the system by the action of the $`SU(2)`$ group and choose the $`SU(2)`$-invariant stereographic coordinates and momenta in the form corresponding to the standard Hurwitz transformation, which yields a pseudospherical analog of the so-called $`SU(2)`$-Kepler (or Yang-Coulomb) system . The potential term of the resulting system will be given by the expression
$$V_{SU(2)Kepler}=\frac{I^2}{r_0^2}\left(\frac{x_5^2}{2๐ฑ^2}2\right)\frac{\gamma }{2r_0}\frac{x_5}{|๐ฑ|}$$
(3.146)
where $`(๐ฑ,x_5)`$ are the (pseudo)Euclidean coordinates of the ambient space $`\mathrm{IR}^{1.5}`$ of the five-dimensional hyperboloid, $`|๐ฑ|^2x_5^2`$; $`I^2`$ is the value of the generator $`๐ฅ_i^2`$, under which the $`SU(2)`$ reduction has been performed. The constants $`r_0`$, $`\gamma `$ are defined by the expressions (3.142).
It is interesting to clarify, which systems will the $`\text{I}\mathrm{CP}^\mathrm{N}`$-oscillators, after similar reductions, result in. We have checked it only for the first Hopf map, corresponding to the case $`p=2`$ .To our surprise, we found that the oscillators on $`\text{I}\mathrm{CP}^2`$ and $`_2`$ also resulted, after reduction, in the pseudospherical MIC-Kepler system!
## 4 Supersymplectic structures
In the previous Sections we presented some elements of Hamiltonian formalism which, in our belief, could be useful in the study of supersymmetric mechanics.
In the present Section we shall briefly discuss the Hamiltonian formalism on superspaces (super-Hamiltonian formalism). The super-Hamiltonian formalism, in its main lines, is a straightforward extension of the ordinary Hamiltonian formalism to superspace, with a more or less obvious placement of sign factors. Probably, from the supergeometrical viewpoint, the only qualitative difference appears in the existence of the odd Poisson brackets (antibrackets), which have no analogs in ordinary spaces, and in the respect of the differential forms to integration. Fortunately, these aspects are inessential for our purposes.
The Poisson brackets of the functions $`f(x)`$ and $`g(x)`$ on superspaces are defined by the expression
$$\{f,g\}_\kappa =\frac{_rf}{x^A}\mathrm{\Omega }_\kappa ^{AB}(x)\frac{_lg}{x^B},\kappa =0,1.$$
(4.147)
They obey the conditions
$`p(\{f,g\}_\kappa )=p(f)+p(g)+\kappa (\mathrm{grading}),`$
$`\{f,g\}_\kappa =(1)^{(p(f)+\kappa )(p(g)+\kappa )}\{g,f\}_\kappa (\mathrm{"}\mathrm{antisymmetricity}\mathrm{"}),`$ (4.148)
$`(1)^{(p(f)+1)(p(h)+\kappa )}\{f,\{g,h\}_\kappa \}_\kappa +\mathrm{cycl}.\mathrm{perm}.(\mathrm{f},\mathrm{g},\mathrm{h})=0(\mathrm{Jacobi}\mathrm{id}.).`$ (4.149)
Here $`x^A`$ are local coordinates of superspace, while $`\frac{_r}{x^A}`$ and $`\frac{_l}{x^A}`$ denote right and left derivatives, respectively.
It is seen that the nondegenerate odd Poisson brackets can be defined on the $`(N.N)`$-dimensional superspaces, and the nondegenerate even Poisson brackets could be defined on the $`(2N.M)`$-dimensional ones. In this case the Poisson brackets are associated with the supersymplectic structure
$$\mathrm{\Omega }_\kappa =dz^A\mathrm{\Omega }_{(\kappa )AB}dz^B,d\mathrm{\Omega }_\kappa =0$$
(4.150)
where $`\mathrm{\Omega }_{(\kappa )AB}\mathrm{\Omega }_\kappa ^{BC}=\delta _A^C`$.
The generalization of the Darboux theorem states that locally, the nondegenerate Poisson brackets could be transformed to the canonical form. The canonical odd Poisson brackets look as follows:
$$\{f,g\}_1^{\mathrm{can}}=\underset{i=1}{\overset{N}{}}\left(\frac{_rf}{x^i}\frac{_lg}{\theta _i}\frac{_rf}{\theta _i}\frac{_lg}{x^i}\right),$$
(4.151)
where $`p(\theta _i)=p(x^i)+1=1`$. The canonical even Poisson brackets read
$$\{f,g\}_0=\underset{i=1}{\overset{N}{}}\left(\frac{f}{x^i}\frac{g}{x^{i+N}}\frac{f}{x^{i+N}}\frac{g}{x^i}\right)+\underset{\alpha =1}{\overset{M}{}}ฯต_\alpha \frac{_rf}{\theta ^\alpha }\frac{_lLg}{\theta ^\alpha },ฯต_\alpha =\pm 1.$$
(4.152)
Here $`x^i,x^{i+N}`$ denote even coordinates, $`p(x)=0`$, and $`\theta ^\alpha `$ are the odd ones $`p(\theta )=1`$.
In a completely similar way to the ordinary (non-โsuperโ) space, one can show that the vector field preserving the supersymplectic structure is a locally Hamiltonian one. Hence, both types of supersymplectic structures can be related with the Hamiltonian systems, which have the following equations of motion:
$$\frac{dx^A}{dt}=\{_\kappa ,x^A\}_\kappa ,p(_\kappa )=\kappa .$$
(4.153)
Any supermanifold $``$ underlied by the bosonic manifold $`M_0`$ can be associated with some vector bundle $`VM_0`$ of $`M_0`$ , in the following sense. One can choose on $``$ local coordinates $`(x^i,\theta ^\mu )`$, such that the transition functions from one chart (parameterized by $`(x^i,\theta ^\mu )`$) to the other chart (parameterized by $`(\stackrel{~}{x}^i,\stackrel{~}{\theta }^\mu )`$) look as follows:
$$\stackrel{~}{x}^i=\stackrel{~}{x}^i(x),\stackrel{~}{\theta }^\mu =A_\nu ^\mu (x)\theta ^\nu .$$
(4.154)
Changing the parity of $`\theta `$: $`p(\theta ^\mu )=1p(\theta ^\mu )=0`$, we shall get the vector bundle $`VM_0`$ of $`M_0`$.
Any supermanifold equipped with the odd symplectic structure, is associated with the cotangent bundle of $`M_0`$ , so that the odd symplectic structure could be globally transformed to the canonical form, with the odd Poisson bracket given by the expression (4.151). Hence, the functions on the odd symplectic manifold could be interpreted as contravariant antisymmetric tensors on $`M_0`$.
The structure of the even symplectic manifold is not so rigid: there is a variety of ways to extend the given symplectic manifold $`(M_0,\omega )`$ to the supersymplectic ones, associated with the vector bundle $`VM_0`$. On these supermanifolds one can (globally) define the even symplectic structure
$`\mathrm{\Omega }`$ $`=`$ $`\omega +d\left(\theta ^\mu g_{\mu \nu }(x)๐\theta ^\nu \right)`$ (4.155)
$`=`$ $`\omega +{\displaystyle \frac{1}{2}}R_{\nu \mu ki}\theta ^\nu \theta ^\mu dx^idx^k+g_{\mu \nu }๐\theta ^\nu ๐\theta ^\mu ,`$
Here $`x^i`$ are local coordinates of $`_0`$ and $`\theta ^\mu `$ are the (odd) coordinates in the bundle; $`g_{\mu \nu }=g_{\mu \nu }(x)`$ are the components of the metrics in the bundle, while $`๐\theta ^\mu =d\theta ^\mu +\mathrm{\Gamma }_{\nu i}^\mu \theta ^\nu dx^i`$, where $`\mathrm{\Gamma }_{i\nu }^\mu `$ are the connection components respecting the metric in the bundle
$$g_{\mu \nu ;k}=g_{\mu \nu ,k}g_{\mu \alpha }\mathrm{\Gamma }_{k\nu }^\alpha g_{\alpha \nu }\mathrm{\Gamma }_{k\mu }^\alpha =0.$$
(4.156)
We used the following notation as well: $`R_{\mu \nu ki}=g_{\mu \alpha }R_{\nu ki}^\alpha `$, where $`R_{\nu ki}^\mu `$ are the components of connectionโs curvature
$$R_{\alpha ki}^\nu =\mathrm{\Gamma }_{k\alpha ,i}^\nu +\mathrm{\Gamma }_{i\alpha ,k}^\nu +\mathrm{\Gamma }_{k\beta }^\nu \mathrm{\Gamma }_{i\alpha }^\beta \mathrm{\Gamma }_{i\beta }^\nu \mathrm{\Gamma }_{k\alpha }^\beta ;R_{\alpha ik}^\nu =R_{\alpha ki}^\nu .$$
Let us consider the coordinate transformation (4.154). With respect to this transformation, the connection components transform as follows:
$`\overline{\mathrm{\Gamma }}_{i\nu }^\mu =A_\lambda ^\mu \mathrm{\Gamma }_{k\alpha }^\lambda {\displaystyle \frac{_rx^k}{\overline{x}^i}}B_\nu ^\alpha A_{\alpha ,k}^\mu B_\nu ^\alpha {\displaystyle \frac{_rx^k}{\overline{x}^i}},A_\mu ^\nu B_\nu ^\lambda =\delta _\mu ^\lambda .`$ (4.157)
Since $`๐\theta ^\nu `$ transforms homogeneously under (4.154), $`๐\overline{\theta }^\nu =๐\theta ^\mu A_\mu ^\nu (x)`$, we conclude that the supersymplectic structure (4.155) is covariant under (4.154) as well.
The corresponding Poisson brackets look as follows:
$$\{f,g\}=(_if)\stackrel{~}{\omega }^{ij}(_jg)+\alpha \frac{_rf}{\theta ^\mu }g^{\mu \nu }\frac{_lg}{\theta ^\nu };$$
(4.158)
where
$$\stackrel{~}{\omega }^{im}(\omega _{mj}+\frac{1}{2}R_{\nu \mu mj}\theta ^\nu \theta ^\mu )=\delta _j^i,_i=\frac{}{x^i}\mathrm{\Gamma }_{ij}^k(x)\theta ^{ja}\frac{}{\theta ^{ka}}.$$
On the supermanifolds one can define also the analog of the Kรคhler structures. We shall call the complex symplectic supermanifold an even (odd) Kรคhler one, when the even (odd) symplectic structure is defined by the expression
$$\mathrm{\Omega }_\kappa =i(1)^{p_A(p_B+\kappa +1)}g_{(\kappa )A\overline{B}}dz^Ad\overline{z}^B,$$
(4.159)
where
$$g_{(\kappa )A\overline{B}}=(1)^{(p_A+\kappa +1)(p_B+\kappa +1)+\kappa +1}\overline{g_{(\kappa )B\overline{A}}},p(g_{(\kappa )A\overline{B}})=p_A+p_B+\kappa .$$
Here and in the following, the index $`\kappa =0(1)`$ denotes the even(odd) case.
The Kรคhler potential on the supermanifold is a local real even (odd) function $`K_\kappa (z,\overline{z})`$ defining the Kรคhler structure
$$g_{(\kappa )A\overline{B}}=\frac{_l}{z^A}\frac{_r}{\overline{z}^B}K_\kappa (z,\overline{z}).$$
(4.160)
As in the usual case, $`K_\kappa `$ is defined up to arbitrary holomorphic and antiholomorphic functions.
With the even (odd) form $`\mathrm{\Omega }_\kappa `$ one can associate the even (odd) Poisson bracket
$$\{f,g\}_\kappa =i\left(\frac{_rf}{\overline{z}^A}g^{(\kappa )\overline{A}B}\frac{_lg}{z^B}(1)^{(p_A+\kappa )(p_B+\kappa )}\frac{_rf}{z^A}g^{(\kappa )\overline{A}B}\frac{_lg}{\overline{z}^B}\right),$$
(4.161)
where
$$g^{(\kappa )\overline{A}B}g_{(\kappa )B\overline{C}}=\delta _{\overline{C}}^{\overline{A}},\overline{g^{(\kappa )\overline{A}B}}=(1)^{(p_A+\kappa )(p_B+\kappa )}g^{(\kappa )\overline{B}A}.$$
Example. Let us consider the supermanifold $`\mathrm{\Lambda }M`$ associated with the tangent bundle of the Kรคhler manifold $`M_0`$. On this supermanifold one can define the even and odd Kรคhler potentials
$$K_0=K(z,\overline{z})+F(ig_{a\overline{b}}\sigma ^a\overline{\sigma }^b),K_1=\frac{K(z,\overline{z})}{z^a}\sigma ^a+\frac{K(z,\overline{z})}{\overline{z}^a}\overline{\sigma }^a,$$
(4.162)
where $`K(z,\overline{z})`$ is a Kรคhler potential on $`M_0`$, $`g_{a\overline{b}}=^2K/z^a\overline{z}^b`$, and $`F(x)`$ is a real function which obeys the condition $`F^{}(0)0`$. It is clear that these functions define even and odd Kรคhler structures on $`\mathrm{\Lambda }M_0`$, respectively.
Finally, let us notice that the analog of the Liouville measure for the even supersymplectic symplectic structure $`\mathrm{\Omega }_0`$ reads
$$\rho =\sqrt{\mathrm{Ber}\mathrm{\Omega }_{(0)AB}},$$
(4.163)
while the odd symplectic structure has no similar invariant . Indeed, one can verify that the even super-Hamiltonian vector field is always divergenceless, $`\mathrm{str}\{H,\}_0=0`$ (similarly to the non-superHamiltonian vector field), while in the case of the odd super-Hamiltonian vector field this property of the Hamiltonian vector field fails. As a consequence, in the latter case the so-called $`\mathrm{\Delta }`$-operator can be defined , which plays a crucial role in the Batalin-Vilkovisky formalism (Lagrangian BRST quantization formalism) .
### Odd super-Hamiltonian mechanics
Let us consider the supermanifold $`\mathrm{\Lambda }M`$, associated with the tangent bundle of the symplectic manifold $`(M,\omega )`$, i.e. the external algebra of $`(M,\omega )`$. In other words, the odd coordinates $`\theta ^i`$ transform from one chart to another like $`dx^i`$, and they can be interpreted as the basis of the 1-forms on $`M`$. By the use of the $`\omega `$ we can equip $`\mathrm{\Lambda }M`$ with the odd symplectic structure
$$\mathrm{\Omega }_1=d\left(\omega _{ij}\theta ^jdx^i\right)=\omega _{ij}dx^id\theta ^j+\frac{1}{2}\omega _{ki,j}\theta ^jdx^kdx^i.$$
(4.164)
The corresponding odd Poisson brackets are defined by the following relations:
$$\{x^i,x^j\}_1=0,\{x^i,\theta ^j\}_1=\omega ^{ij},\{\theta ^i,\theta ^j\}_1=\frac{\omega ^{ij}}{x^k}\theta ^k,$$
(4.165)
where $`\omega ^{ij}\omega _{jk}=\delta _k^i`$.
Let us define, on $`\mathrm{\Lambda }M`$, the even function
$$F=\frac{1}{2}\theta ^i\omega _{ij}\theta ^j,:\{F,F\}_1=0,$$
(4.166)
where the latter equation holds due to the closeness of $`\omega `$. By making use of this function, one can define the map of any function on $`M`$ in the odd function on $`\mathrm{\Lambda }M`$
$$f(x)Q_f(x,\theta )=\{f(x),F(x,\theta )\}_1,$$
(4.167)
which possesses the following important property:
$$\{f(x),g(x)\}=\{f(x),Q_g(x,\theta )\}_1\mathrm{for}\mathrm{any}f(x),g(x).$$
(4.168)
In particular, (4.167) maps the Hamiltonian mechanics $`(M,\omega ,H(x))`$ in the following super-Hamiltonian one: $`\left(\mathrm{\Lambda }M,\mathrm{\Omega }_1,Q_H=\{H,F\}_1\right)`$, where $`Q_H`$ plays the role of the odd Hamiltonian on $`\mathrm{\Lambda }M`$.
The functions $`H,F,Q_H`$ form the superalgebra
$$\begin{array}{c}\{H\pm F,H\pm F\}_1=\pm 2Q_H,\\ \{H+F,HF\}_1=\{H\pm F,Q_H\}_1=\{Q_H,Q_H\}_1=0,\end{array}$$
(4.169)
i.e. the resulting mechanics possesses the supersymmetry transformation defined by the โsuperchargeโ $`H+F`$. This superalgebra has a transparent interpretation in terms of base manifold $`(M,\omega )`$
$$\{H,\}_1=\xi _H^i\frac{}{\theta ^i}\widehat{ฤฑ}_H\mathrm{contraction}\mathrm{with}\xi _H,$$
$$\{F,\}_1=\theta ^i\frac{}{x^i}\widehat{d}\mathrm{exterior}\mathrm{differential},$$
$$\{Q,\}_1=\xi _H^i\frac{}{x^i}+\xi _{H,k}^i\theta ^k\frac{}{\theta ^i}\widehat{}_H\mathrm{Lie}\mathrm{derivative}\mathrm{along}\xi _H,$$
while, using the Jacobi identity (4.149), we get
$$\{H,F\}_1=Q_H\widehat{d}\widehat{ฤฑ}_H+\widehat{ฤฑ}_H\widehat{d}=\widehat{}_H\mathrm{homotopy}\mathrm{formula}.$$
Hence, the above dynamics could be useful for the description of the differential calculus on the symplectic (and Poisson) manifolds. Particulary, it has a nice application in equivariant cohomology and related localization formulae (see and refs therein).
However, the presented supersymmetric model has no deep dynamical meaning, since the odd Poisson brackets do not admit any consistent quantization scheme. Naively, this is reflected in the fact that conjugated operators should have opposite Grassmann grading, so that the Planck constant must be a Grassmann-odd number.
Moreover, the presented supersymmetric mechanics is not interesting even from the classical viewpoint. Its equations of motion read
$$\frac{dx^i}{dt}=\{x^i,Q_H\}_1=\xi _H^i,\frac{d\theta ^i}{dt}=\{\theta ^i,Q_H\}_1=\frac{\xi _H^i}{x^j}\theta ^j,$$
i.e. the โfermionicโ degrees of freedom have no impact in the dynamics of the โbosonicโ degrees of freedom.
Nevertheless, the odd Poisson brackets are widely known, since 1981, in the theoretical physics community under the name of โantibracketsโ. That was the year, when Batalin and Vilkovisky suggested their Covariant Lagrangian BRST quantization formalism (which is known presently as the Batalin-Vilkovisky formalism) , where the antibrackets (odd Poisson brackets) play the key role. However, only decades after, this elegant formalism was understood in terms of conventional supergeometrical constructions . It seems that the Batalin-Vilkovisky formalism could also be useful for the geometrical (covariant) formulation of the superfield approach to the construction of supersymmetric Lagrangian field-theoretical and mechanical models .
We shall not touch upon these aspects of super-Hamiltonian systems, and will restrict ourselves to the consideration of supersymmetric Hamiltonian systems with even symplectic structure.
### Hamiltonian reduction: $`\text{I}\mathrm{C}^{\mathrm{N}+1.\mathrm{M}}\text{I}\mathrm{CP}^{\mathrm{N}.\mathrm{M}}`$, $`\mathrm{\Lambda }\text{I}\mathrm{C}^{\mathrm{N}+1}\mathrm{\Lambda }\text{I}\mathrm{CP}^\mathrm{N}`$
The procedure of super-Hamiltonian reduction is very similar to the Hamiltonian one. The main difference is in the counting of the dimensionality of the phase superspace. Namely, we should separately count the number of โfermionicโ and โbosonicโ degrees of freedom, which were eliminated during the reduction.
Instead of describing the extension of the Hamiltonian reduction to the supercase, we shall illustrate it by considering superextensions of the reduction $`\text{I}\mathrm{C}^{\mathrm{N}+1}\text{I}\mathrm{CP}^\mathrm{N}`$ presented in Third Section. These examples were considered in details in Ref. .
Let us consider the complex superspace $`\text{I}\mathrm{C}^{\mathrm{N}+1,\mathrm{M}}`$ parameterized by the complex coordinates $`(u^{\stackrel{~}{a}},\eta ^n)`$, $`\stackrel{~}{a}=0,1,\mathrm{},N`$, $`n=1,\mathrm{},M`$. Let us equip it with the canonical symplectic structure
$$\mathrm{\Omega }^0=i(du^{\stackrel{~}{a}}\overline{d}u^{\stackrel{~}{a}}id\eta ^nd\overline{\eta }^n)$$
and with the corresponding even Poisson bracket
$$\{f,g\}_0=i\left(\frac{f}{u^{\stackrel{~}{a}}}\frac{g}{\overline{u}^{\stackrel{~}{a}}}\frac{f}{\overline{u}^{\stackrel{~}{a}}}\frac{g}{u^{\stackrel{~}{a}}}\right)+\frac{_rf}{\eta ^n}\frac{_lg}{\overline{\eta }^n}+\frac{_rf}{\overline{\eta }^n}\frac{_lg}{\eta ^n}.$$
(4.170)
The (super-)Hamiltonian action of the $`U(1)`$ group is given, on this space, by the generator
$$๐ฅ_0=u^{\stackrel{~}{a}}\overline{u}^{\stackrel{~}{a}}i\eta ^n\overline{\eta }^n.$$
(4.171)
For the reduction of $`\text{I}\mathrm{C}^{\mathrm{N}+1.\mathrm{M}}`$ by this generator, we have to factorize the $`(2N+1.2M)_{\mathrm{IR}}`$-dimensional level supersurface
$$๐ฅ_0=r_0^2$$
(4.172)
by the even super-Hamiltonian vector field $`\{๐ฅ_0,\}`$ (which is tangent to that surface). Hence, the resulting phase superspace is a $`(2N.2M)_{\mathrm{IR}}`$-dimensional one.
Hence, for the role of local coordinates of the reduced phase space, we have to choose the $`N`$ even and $`M`$ odd complex functions commuting with $`๐ฅ_0`$. On the chart $`u^{\stackrel{~}{a}}0`$, appropriate functions are the following ones:
$$z_{(\stackrel{~}{a})}^A=\left(z_{(\stackrel{~}{a})}^a=\frac{u^a}{u^{\stackrel{~}{a}}},\theta _{(\stackrel{~}{a})}^k=\frac{\eta ^k}{u^{\stackrel{~}{a}}},a\stackrel{~}{a}\right):\{z_{(\stackrel{~}{a})}^A,๐ฅ_0\}_0=0.$$
(4.173)
The reduced Poisson brackets could be defined by the expression $`\{f,g\}_0^{\mathrm{red}}=\{f,g\}_0_{๐ฅ_0=r_0^2}`$, where $`f,g`$ are functions depending on the coordinates $`z_{(\stackrel{~}{a})}^A,\overline{z}_{(\stackrel{~}{b})}^A`$. Straightforward calculations yield the result
$$\{z^A,z^B\}_0^{\mathrm{red}}=\{\overline{z}^A,\overline{w}^B\}_0^{\mathrm{red}}=0,$$
$$\{z^A,\overline{z}^B\}_0^{\mathrm{red}}=(i)^{p_Ap_B+1}\frac{1+(i)^{p_C}z^C\overline{z}^C}{r_0^2}\left(\delta ^{AB}+(i)^{p_Ap_B}z^A\overline{w}^B\right).$$
It is seen that these Poisson brackets are associated with a Kรคhler structure. This Kรคhler structure is defined by the potential
$$K=r_0^2\mathrm{log}(1+(i)^{p_C}z^C\overline{z}^{\overline{C}}).$$
(4.174)
The transition functions from the $`\stackrel{~}{a}`$-th chart to the $`\stackrel{~}{b}`$-th one look as follows:
$$z_{(\stackrel{~}{a})}^{\stackrel{~}{c}}=\frac{z_{(\stackrel{~}{b})}^{\stackrel{~}{c}}}{z_{(\stackrel{~}{b})}^{\stackrel{~}{a}}},\theta _{(\stackrel{~}{a})}^k=\frac{\theta _{(\stackrel{~}{b})}^k}{z_{(\stackrel{~}{b})}^{\stackrel{~}{a}}},\mathrm{where}z_{(\stackrel{~}{b})}^{\stackrel{~}{a}}=\left(w_{(\stackrel{~}{b})}^a,w_{(\stackrel{~}{a})}^{\stackrel{~}{a}}=1\right).$$
(4.175)
Upon these transformations the Kรคhler potential changes on the holomorphic and anti-holomorphic functions, i.e. the reduced phase space is indeed a Kรคhler supermanifold. We shall refer to it as $`\text{I}\mathrm{CP}^{\mathrm{N}.\mathrm{M}}`$. The quantization of this supermanifold is considered in .
Now, let us consider the Hamiltonian reduction of the superspace $`\text{I}\mathrm{C}^{\mathrm{N}+1,\mathrm{N}+1}`$ by the action of the $`๐ฉ=2`$ superalgebra, given by the generators
$$\begin{array}{c}๐ฅ_0=u^{\stackrel{~}{a}}\overline{u}^{\stackrel{~}{a}}i\eta ^{\stackrel{~}{a}}\overline{\eta }^{\stackrel{~}{a}},\mathrm{\Theta }^+=u^{\stackrel{~}{a}}\overline{\eta }^{\stackrel{~}{a}},\mathrm{\Theta }^{}=\overline{u}^{\stackrel{~}{a}}\eta ^{\stackrel{~}{a}}:\\ \{\mathrm{\Theta }^+,\mathrm{\Theta }^{}\}=๐ฅ_0,\{\mathrm{\Theta }^\pm ,\mathrm{\Theta }^\pm \}=\{\mathrm{\Theta }^\pm ,๐ฅ_0\}=0.\end{array}$$
(4.176)
The equations
$$J_0=r_0^2,\mathrm{\Theta }^\pm =0$$
(4.177)
define the $`(2N+1.2N)`$-dimensional level surface $`M_{r_0^2,0,0}`$. The reduced phase superspace can be defined by the factorization of $`M_{r_0^2,0,0}`$ by the action of the tangent vector field $`\{๐ฅ,\}_0`$. Hence, the reduced phase superspace is a $`(2N.2N)_{\mathrm{IR}}`$-dimensional one. The conventional local coordinates of the reduced phase superspace could be chosen as follows (on the chart $`u^00`$):
$$\sigma ^a=i\{z^a,\mathrm{\Theta }^+\}=\theta ^a\theta ^0z^a,w^a=z^a+i\frac{\mathrm{\Theta }^{}}{๐ฅ_0}\sigma ^a,$$
(4.178)
where $`z^a,\theta ^0,\theta ^a`$ are defined by (4.173). The reduced Poisson brackets are defined as follows:
$$\{f,g\}_0^{\mathrm{red}}=\{f,g\}_0_{๐ฅ=r_0^2,\mathrm{\Theta }^\pm =0},$$
where $`f,g`$ are the functions on $`(w^a,\sigma ^a)`$. Straightforward calculations result in the following relations:
$`\{w^A,w^B\}_0^{\mathrm{red}}`$ $`=`$ $`\{\overline{w}^A,\overline{w}^B\}_0^{\mathrm{red}}=0,\mathrm{where}w^A=(w^a,\sigma ^a)`$
$`\{w^a,\overline{w}^b\}_0^{\mathrm{red}}`$ $`=`$ $`i{\displaystyle \frac{A}{r_0^2}}(\delta ^{ab}+w^a\overline{w}^b){\displaystyle \frac{\sigma ^a\overline{\sigma }^b}{r_0^2}},`$
$`\{w^a,\overline{\sigma }^b\}_0^{\mathrm{red}}`$ $`=`$ $`i{\displaystyle \frac{A}{r_0^2}}\left(w^a\overline{\sigma }^b+\mu (\delta ^{ab}+w^a\overline{w}^b)\right)`$ (4.179)
$`\{\sigma ^a,\overline{\sigma }^b\}_0^{\mathrm{red}}`$ $`=`$ $`{\displaystyle \frac{A}{r_0^2}}((1+i\mu \overline{\mu })\delta ^{ab}+w^a\overline{w}^b+i(\sigma ^a+\mu w^a)(\overline{\sigma }^b+\overline{\mu }\overline{w}^b),`$
and
$$A=1+w^a\overline{w}^ai\sigma ^a\overline{\sigma }^a+\frac{i\sigma ^a\overline{w}^a\overline{\sigma }^bw^b}{1+w^c\overline{w}^c},\mu =\frac{\overline{w}^a\sigma ^a}{1+w^b\overline{w}^b}.$$
These Poisson brackets are associated with the Kรคhler structure defined by the potential
$$\begin{array}{c}K=r_0^2\mathrm{log}A(w,\overline{w},\sigma ,\overline{\sigma })==r_0^2\mathrm{log}(1+w^a\overline{w}^a)+r_0^2\mathrm{log}(1ig_{a\overline{b}}\sigma ^a\overline{\sigma }^b).\end{array}$$
(4.180)
where $`g_{a\overline{b}}(w,\overline{w})`$ is the Fubini-Study metric on $`\text{I}\mathrm{CP}^\mathrm{N}`$.
The transition functions from the $`\stackrel{~}{a}`$-th chart to the $`\stackrel{~}{b}`$-th one reads
$$w_{(\stackrel{~}{b})}^{\stackrel{~}{c}}=\frac{w_{(\stackrel{~}{a})}^{\stackrel{~}{c}}}{w_{(\stackrel{~}{a})}^{\stackrel{~}{b}}},\sigma _{(\stackrel{~}{b})}^{\stackrel{~}{c}}=\frac{\sigma _{(\stackrel{~}{a})}^{\stackrel{~}{c}}x_{(\stackrel{~}{a})}^{\stackrel{~}{b}}w_{(\stackrel{~}{a})}^{\stackrel{~}{c}}\sigma _{(\stackrel{~}{a})}^{\stackrel{~}{b}}}{(w_{(\stackrel{~}{a})}^{\stackrel{~}{b}})^2},$$
where $`(w_{(\stackrel{~}{a})}^{\stackrel{~}{a}}=1,\sigma _{(\stackrel{~}{b})}^{\stackrel{~}{b}}=0)`$. Hence, $`\sigma ^a`$ transforms like $`dw^a`$, i.e. the reduced phase superspace is $`\mathrm{\Lambda }\text{I}\mathrm{CP}^\mathrm{N}`$, the external algebra of the the complex projective space $`\text{I}\mathrm{CP}^\mathrm{N}`$.
Remark 1. On $`\text{I}\mathrm{C}^{\mathrm{N}+1,\mathrm{N}+1}`$ one can define the odd Kรคhler structure as well, $`\mathrm{\Omega }^1=du^nd\overline{\eta }^n+d\overline{u}^nd\eta ^n`$. It could be reduced to the odd Kรคhler structure on $`\mathrm{\Lambda }\text{I}\mathrm{CP}^\mathrm{N}`$ by the action of the generators
$$J_0=z\overline{z},Q=z\overline{\eta }+\overline{z}\eta .$$
Remark 2. The generalization of the reduction $`T^{}\text{I}\mathrm{C}^2\mathrm{T}^{}\mathrm{IR}^3`$, where the latter is specified by the presence of a Dirac monopole, is also straightforward. One should consider the $`(4.M)_{\text{I}\mathrm{C}}`$\- dimensional superspace equipped with the canonical even symplectic structure $`\mathrm{\Omega }_0=d\pi dz+d\overline{\pi }d\overline{z}+d\eta d\overline{\eta }`$, and reduce it by the Hamiltonian action of the $`U(1)`$ group given by the generator $`๐ฅ=i\pi zi\overline{\pi }\overline{z}i\eta \overline{\eta }`$. The resulting space is a $`(6.2M)_{\mathrm{IR}}`$-dimensional one. Its even local coordinates could be defined by the same expressions, as in the bosonic case, Eq.(2.82), while the odd coordinates could be chosen as follows: $`\theta ^m=f(z\overline{z})\overline{z}_0\eta ^m`$.
## 5 Supersymmetric mechanics
In the previous Sections we presented some basic elements of the Hamiltonian and super-Hamiltonian formalism. We paid special attention to the examples, related with Kรคhler geometry, keeping in mind that the latter is of a special importance in supersymmetric mechanics. Indeed, the incorporation of the Kรคhler structure(s) is one of the standard ways to increase the number of supersymmetries of the system.
Our goal is to construct the supersymmetric mechanics with $`๐ฉ2`$ supersymmetries. This means that, on the given phase superspace equipped with even symplectic structure, we should construct the Hamiltonian $``$ which has $`๐ฉ=N`$ odd constants of motion $`Q_i`$ forming the superalgebra
$$\{Q_i,Q_j\}=2\delta _{ij},\{Q_i,\}=0.$$
(5.181)
This kind of mechanics is referred to as โ$`๐ฉ=N`$ supersymmetric mechanicsโ.
It is very easy to construct the $`๐ฉ=1`$ supersymmetric mechanics with single supercharges: we should simply take the square (under a given nondegenerate even Poisson bracket) of the arbitrary odd function $`Q_1`$, and consider the resulting even function as the Hamiltonian
$$\{Q_1,Q_1\}2_{SUSY}:\{Q_1,_{SUSY}\}=0.$$
(5.182)
However, the case of $`๐ฉ=1`$ supersymmetric mechanics is not an interesting system, both from the dynamical and field-theoretical viewpoints.
If we want to construct the $`๐ฉ>1`$ supersymmetric mechanics, we must specify both the underlying system and the structure of phase superspace.
Let us illustrate it on the simplest examples of $`๐ฉ=2`$ supersymmetric mechanics. For this purpose, it is convenient to present the $`๐ฉ=2`$ superalgebra as follows:
$$\{Q^+,Q^{}\}=,\{Q^\pm ,Q^\pm \}=0,$$
(5.183)
where $`Q^\pm =(Q_1\pm iQ_2)/\sqrt{2}`$. Hence, we have to find the odd complex function, which is nilpotent with respect to the given nondegenerate Poisson bracket, in order to construct the appropriate system.
Let us consider a particular example, when the underlying system is defined on the cotangent bundle $`T^{}M_0`$, and it is given by (1.12).
In order to supersymmetrize this system, we extend the canonical symplectic structure as follows:
$$\mathrm{\Omega }=dp_adx^a+\frac{1}{2}R_{abcd}\theta _+^a\theta _{}^bdx^cdx^d+g_{ab}D\theta _+^aD\theta _{}^b,$$
(5.184)
where $`D\theta _\pm ^ad\theta _\pm ^a+\mathrm{\Gamma }_{bc}^a\theta _\pm ^bdx^c`$, and $`\mathrm{\Gamma }_{bc}^a`$, $`R_{abcd}`$ are the components of the connection and curvature of the metrics $`g_{ab}dx^adx^b`$ on $`M_0`$.
We choose the following candidate for a complex supercharge:
$$Q_\pm =(p_a\pm iW_{,a})\theta ^a{}_{\pm }{}^{}:\{Q_\pm ,Q_\pm \}=0.$$
(5.185)
Hence, the supersymmetric Hamiltonian could be constructed by the calculation of the Poisson brackets of these supercharges.
$$\{Q_+,Q_{}\}=\frac{1}{2}g^{ab}(p_ap_b+W_{,a}W_{,b})+W_{a;b}\theta _+^a\theta _{}^b+R_{abcd}\theta _{}^a\theta _+^b\theta _{}^c\theta _+^d.$$
(5.186)
The โminimalโ coupling of the magnetic field, $`\mathrm{\Omega }\mathrm{\Omega }+F_{ab}dx^adx^b`$, breaks the $`๐ฉ=2`$ supersymmetry of the system
$$\{Q_\pm ,Q_\pm \}=F_{ab}\theta _\pm ^a\theta _\pm ^b,\{Q_+,Q_{}\}=+iF_{ab}\theta _+^a\theta _{}^b.$$
Notice that the Higgs oscillator on the sphere $`S^N`$, considered in Section 3, could be supersymmetrised in this way, choosing $`W=\frac{\alpha }{2}\mathrm{log}\frac{2+๐ช^\mathrm{๐}}{2๐ช^2},`$ with $`๐ช`$ being the conformal coordinates of the sphere.
One of the ways to extend this construction to $`๐ฉ=4`$ supersymmetric mechanics is the doubling of the number of odd degrees of freedom. It was considered, within the (Lagrangian) superfield approach in Ref.. In this paper the authors considered the $`(2N.2N)_{\mathrm{IR}}`$-dimensional superspace and the supercharges containing term cubic on odd variables. Calculating the Poisson brackets, the authors found that the admissible metrics of the configuration space of that system should have the following local form:
$$g_{ab}=\frac{^2A(x)}{x^ax^b}.$$
(5.187)
The admissible set of potentials looks, in this local coordinates, as follows: $`V=g_{ab}c^{ab}+g^{ab}d_{af}`$, where $`c^{ab}`$ and $`d_{ab}`$ are constant matrices.
So, considering the Hamiltonian system with generic phase spaces, we found that without any efforts it could be extended to $`๐ฉ=1`$ supersymmetric mechanics. For the construction of $`๐ฉ=2`$ supersymmetric mechanics we were forced to restrict ourselves to systems on the cotangent bundle of Riemann manifolds. Even after this strong restriction, we found that the inclusion of a magnetic field breaks the supersymmetry of the system. On the other hand, in trying to construct $`๐ฉ=4`$ supersymmetric mechanics, we found that in this case even the metric of the configuration space and the admissible set of potentials are strongly restricted.
In further examples we shall show that the transition to Kรคhler geometry makes these restrictions much weaker.
### $`๐ฉ=2`$ supersymmetric mechanics with Kรคhler phase space
Let us consider a supersymmetric mechanics whose phase superspace is the external algebra of the Kรคhler manifold $`\mathrm{\Lambda }M`$, where $`(M,g_{a\overline{b}}(z,\overline{z})dz^ad\overline{z}^{\overline{b}})`$ is the phase space of the underlying Hamiltonian mechanics . The phase superspace is $`(D|D)_{\text{I}\mathrm{C}}`$ dimensional supermanifold equipped with the Kรคhler structure
$$\mathrm{\Omega }=i\overline{}\left(Kig_{a\overline{b}}\theta ^a\overline{\theta }^{\overline{b}}\right)=i(g_{a\overline{b}}+iR_{a\overline{b}c\overline{d}}\theta ^c\overline{\theta }^{\overline{d}})dz^ad\overline{z}^{\overline{b}}+g_{a\overline{b}}D\theta ^aD\overline{\theta }^{\overline{b}},$$
where $`D\theta ^a=d\theta ^a+\mathrm{\Gamma }_{bc}^a\theta ^cdz^c`$, and $`\mathrm{\Gamma }_{bc}^a`$, $`R_{a\overline{b}c\overline{d}}`$ are the Cristoffel symbols and curvature tensor of the underlying Kรคhler metrics $`g_{a\overline{b}}=_a_{\overline{b}}K(z,\overline{z})`$, respectively.
The corresponding Poisson bracket can be presented in the form
$$\{,\}=i\stackrel{~}{g}^{a\overline{b}}_a\overline{}_{\overline{b}}+g^{a\overline{b}}\frac{}{\theta ^a}\frac{}{\overline{\theta }^{\overline{b}}}$$
(5.188)
where
$$_a=\frac{}{z^a}\mathrm{\Gamma }_{ab}^c\theta ^b\frac{}{\theta ^c},\stackrel{~}{g}_{a\overline{b}}^1=(g_{a\overline{b}}+iR_{a\overline{b}c\overline{d}}\theta ^c\overline{\theta }^{\overline{d}}).$$
On this phase superspace one can immediately construct $`๐ฉ=2`$ supersymmetric mechanics, defined by the supercharges
$$Q_+^0=_aK(z,\overline{z})\theta ^a,Q_{}^0=_{\overline{a}}K(z,\overline{z})\overline{\theta }^{\overline{a}}$$
(5.189)
where $`K(z,\overline{z})`$ is the Kรคhler potential of $`M`$, defined up to holomorphic and anti-holomorphic functions, $`K(z,\overline{z})K(z,\overline{z})+U(z)+\overline{U}(\overline{z})`$.
The Hamiltonian of the system reads
$$_0=g^{a\overline{b}}_aK_{\overline{b}}Kig_{a\overline{b}}\theta ^a\overline{\theta }^{\overline{b}}+i\theta ^cK_{c;a}\stackrel{~}{g}^{a\overline{b}}K_{\overline{b};\overline{d}}\overline{\theta }^{\overline{d}}$$
(5.190)
where $`K_{a;b}=_a_bK\mathrm{\Gamma }_{ab}^c_cK`$.
Another example of $`๐ฉ=2`$ supersymmetric mechanics is defined by the supercharges
$$Q_+^c=_aG(z,\overline{z})\theta ^a,Q_{}^c=_{\overline{a}}G(z,\overline{z})\overline{\theta }^{\overline{a}},$$
(5.191)
where the real function $`G(z,\overline{z})`$ is the Killing potential of the underlying Kรคhler structure
$$_a_bG\mathrm{\Gamma }_{ab}^c_cG=0,G^a(z)=g^{a\overline{b}}_{\overline{b}}G(z,\overline{z}).$$
(5.192)
In this case the Hamiltonian of system reads
$$^c=g_{a\overline{b}}G^aG^{\overline{b}}+i\overline{\theta }^{\overline{d}}G_{a\overline{d}}\stackrel{~}{g}^{a\overline{b}}G_{c\overline{b}}\theta ^c,$$
(5.193)
where $`G_{a\overline{b}}=_a_{\overline{b}}G(z,\overline{z})`$.
The commutators of the supercharges in these particular examples read
$$\{Q_\pm ^c,Q_\pm ^0\}=_\pm ,\{Q_\pm ^c,Q_{}^0\}=๐ต,$$
(5.194)
where
$$\stackrel{~}{๐ต}G(z,\overline{z})+iG_{a\overline{b}}(z,\overline{z})\theta ^a\overline{\theta }^{\overline{b}},_+=i\theta ^cK_{c;a}\stackrel{~}{g}^{a\overline{b}}G_{\overline{b};d}\theta ^d,_{}=\overline{}_+.$$
(5.195)
Hence, introducing the supercharges
$$\mathrm{\Theta }_\pm =Q_\pm ^0\pm iQ_{}^c,$$
(5.196)
we can define $`N=2`$ SUSY mechanics specified by the presence of the central charge $`๐ต`$
$$\begin{array}{c}\{\mathrm{\Theta }_+,\mathrm{\Theta }_{}\}=\stackrel{~}{},\{\mathrm{\Theta }_\pm ,\mathrm{\Theta }_\pm \}=\pm i๐ต\\ \{๐ต,\mathrm{\Theta }_\pm \}=0,\{\stackrel{~}{},\mathrm{\Theta }_{}\}=0,\{๐ต,\stackrel{~}{}\}=0.\end{array}$$
(5.197)
The Hamiltonian of this generalized mechanics is defined by the expression
$$\stackrel{~}{}=_0+_c+i_+i_{}.$$
(5.198)
A โfermionic numberโ is of the form
$$\stackrel{~}{}=ig_{a\overline{b}}\theta ^a\overline{\theta }^{\overline{b}}:\{\stackrel{~}{},\mathrm{\Theta }_\pm ,\}=\pm i\mathrm{\Theta }_\pm .$$
(5.199)
It seems that, on the external algebra of the hyper-Kahler manifold, in the same manner one could construct $`๐ฉ=4`$ supersymmetric mechanics. On the other hand, the hyper-Kรคhler manifolds are the cotangent bundle of the Kรคhler manifolds equipped with Ricci-flat metrics.
We shall demonstrate, in the next examples, that these restrictions can be too strong. Namely, choosing the underlying phase space to be the cotangent bundle of the Kรคhler manifold, we will double the number of supercharges and get the $`๐ฉ=4`$ supersymmetric mechanics on the cotangent bundles of generic Kรคhler manifolds and the $`๐ฉ=8`$ ones on the cotangent bundles of the special Kรคhler manifolds.
### $`๐ฉ=4`$ supersymmetric mechanics
Let us show that the Hamiltonian mechanics (1.12) could be easily extended to the $`๐ฉ=4`$ supersymmetric mechanics, when the configuration space $`M_0`$ is the Kรคhler manifold $`(M_0,g_{a\overline{b}}dz^ad\overline{z}^{\overline{b}})`$, $`g_{a\overline{b}}=^2K(z,\overline{z})/z^a\overline{z}^b`$, and the potential term has the form
$$V(z,\overline{z})=\frac{\overline{U}(\overline{z})}{\overline{z}^a}g^{\overline{a}b}\frac{U(z)}{z^b}.$$
For this purpose, let us define the supersymplectic structure
$$\begin{array}{c}\mathrm{\Omega }=\omega _0i\overline{}๐ =\\ =d\pi _adz^a+d\overline{\pi }_ad\overline{z}^a+R_{a\overline{b}c\overline{d}}\eta _i^a\overline{\eta }_i^bdz^ad\overline{z}^b+g_{a\overline{b}}D\eta _i^aD\overline{\eta }_i^b\end{array}$$
(5.200)
where
$$๐ =ig_{a\overline{b}}\eta ^a\sigma _0\overline{\eta }^b,D\eta _i^a=d\eta _i^a+\mathrm{\Gamma }_{bc}^a\eta _i^adz^a,i=1,2$$
$`\mathrm{\Gamma }_{bc}^a,R_{a\overline{b}c\overline{d}}`$ are the connection and curvature of the Kรคhler structure, respectively, and the odd coordinates $`\eta _i^a`$ belong to the external algebra $`\mathrm{\Lambda }M_0`$, i. e. they transform as $`dz^a`$. This symplectic structure becomes canonical in the coordinates $`(p_a,\chi ^k)`$
$$\begin{array}{c}p_a=\pi _a\frac{i}{2}_a๐ ,\chi _i^m=\mathrm{e}_b^m\eta _i^b:\\ \mathrm{\Omega }=dp_adz^a+d\overline{p}_{\overline{a}}d\overline{z}^{\overline{a}}+d\chi _i^md\overline{\chi }_i^{\overline{m}},\end{array}$$
(5.201)
where $`\mathrm{e}_a^m`$ are the einbeins of the Kรคhler structure: $`\mathrm{e}_a^m\delta _{m\overline{m}}\overline{\mathrm{e}}_{\overline{b}}^{\overline{m}}=g_{a\overline{b}}.`$ The corresponding Poisson brackets are defined by the following non-zero relations (and their complex-conjugates):
$$\begin{array}{c}\{\pi _a,z^b\}=\delta _a^b,\{\pi _a,\eta _i^b\}=\mathrm{\Gamma }_{ac}^b\eta _i^c,\\ \{\pi _a,\overline{\pi }_b\}=R_{a\overline{b}c\overline{d}}\eta _k^c\overline{\eta }_k^d,\{\eta _i^a,\overline{\eta }_j^b\}=g^{a\overline{b}}\delta _{ij}.\end{array}$$
Let us represent the $`๐ฉ=4`$ supersymmetry algebra as follows:
$$\{Q_i^+,Q_j^{}\}=\delta _{ij},\{Q_i^\pm ,Q_j^\pm \}=\{Q_i^\pm ,\}=0,i=1,2,$$
(5.202)
and choose the supercharges given by the functions
$$Q_1^+=\pi _a\eta _1^a+iU_{\overline{a}}\overline{\eta }_2^{\overline{a}},Q_2^+=\pi _a\eta _2^aiU_{\overline{a}}\overline{\eta }_1^{\overline{a}}.$$
(5.203)
Then, calculating the commutators (Poisson brackets) of these functions, we get that the supercharges (5.203) belong to the superalgebra (5.202), when the functions $`U_a,\overline{U}_{\overline{a}}`$ are of the form
$$U_a(z)=\frac{U(z)}{z^a},\overline{U}_{\overline{a}}(\overline{z})=\frac{\overline{U}(\overline{z})}{\overline{z}^a},$$
(5.204)
while the Hamiltonian reads
$$=g^{a\overline{b}}(\pi _a\overline{\pi }_b+U_a\overline{U}_{\overline{b}})iU_{a;b}\eta _1^a\eta _2^b+i\overline{U}_{\overline{a};\overline{b}}\overline{\eta }_1^{\overline{a}}\overline{\eta }_2^{\overline{b}}R_{a\overline{b}c\overline{d}}\eta _1^a\overline{\eta }_1^b\eta _2^a\overline{\eta }_2^d,$$
(5.205)
where $`U_{a;b}_a_bU\mathrm{\Gamma }_{ab}^c_cU`$.
Now, following , let us extend this system to $`๐ฉ=4`$ supersymmetric mechanics with central charge
$$\{\mathrm{\Theta }_i^+,\mathrm{\Theta }_j^{}\}=\delta _{ij}+๐ต\sigma _{ij}^3,\{\mathrm{\Theta }_i^\pm ,\mathrm{\Theta }_j^\pm \}=0,\{๐ต,\}=\{๐ต,\mathrm{\Theta }_k^\pm \}=0.$$
(5.206)
For this purpose one introduces the supercharges
$$\begin{array}{c}\mathrm{\Theta }_1^+=\left(\pi _a+iG_{,a}(z,\overline{z})\right)\eta _1^a+i\overline{U}_{,\overline{a}}(\overline{z})\overline{\eta }_2^{\overline{a}},\\ \mathrm{\Theta }_2^+=\left(\pi _aiG_{,a}(z,\overline{z})\right)\eta _2^ai\overline{U}_{,\overline{a}}(\overline{z})\overline{\eta }_1^{\overline{a}},\end{array}$$
(5.207)
where the real function $`G(z,\overline{z})`$ obeys the conditions (5.192) and $`_{\overline{a}}Gg^{\overline{a}b}U_b=0`$. So, $`G`$ is a Killing potential defining the isometry of the underlying Kรคhler manifold (given by the vector $`๐=G^a(z)_a+\overline{G}^a(\overline{z})\overline{}_a,G^a=ig^{a\overline{b}}\overline{}_bG`$) which leaves the holomorphic function $`U(z)`$ invariant
$$_๐U=0G^a(z)U_a(z)=0.$$
Calculating the Poisson brackets of these supercharges, we get explicit expressions for the Hamiltonian
$$\begin{array}{c}g^{a\overline{b}}\left(\pi _a\overline{\pi }_{\overline{b}}+G_{,a}G_{\overline{b}}+U_{,a}\overline{U}_{,\overline{b}}\right)\\ iU_{a;b}\eta _1^a\eta _2^b+i\overline{U}_{\overline{a};\overline{b}}\overline{\eta }_1^{\overline{a}}\overline{\eta }_2^{\overline{b}}+\frac{1}{2}G_{a\overline{b}}(\eta _k^a\overline{\eta }_k^{\overline{b}})R_{a\overline{b}c\overline{d}}\eta _1^a\overline{\eta }_1^b\eta _2^c\overline{\eta }_2^d\end{array}$$
(5.208)
and for the central charge
$$๐ต=i(G^a\pi _a+G^{\overline{a}}\overline{\pi }_{\overline{a}})+\frac{i}{2}_a\overline{}_{\overline{b}}G(\eta ^a\sigma _3\overline{\eta }^{\overline{b}}).$$
(5.209)
It can be checked by a straightforward calculation that the function $`๐ต`$ indeed belongs to the center of the superalgebra (5.206). The scalar part of each phase with standard $`๐ฉ=2`$ supersymmetry can be interpreted as a particle moving on the Kรคhler manifold in the presence of an external magnetic field, with strength $`F=iG_{a\overline{b}}dz^ad\overline{z}^{\overline{b}}`$, and in the potential field $`U_{,a}(z)g^{a\overline{b}}\overline{U}_{,\overline{b}}(\overline{z})`$.
Assuming that $`(M_0,g_{a\overline{b}}dz^ad\overline{z}^b)`$ is the hyper-Kรคhler metric, $`U(z)+\overline{U}(\overline{z})`$ is a tri-holomorphic function and $`G(z,\overline{z})`$ defines a tri-holomorphic Killing vector, one should get $`๐ฉ=8`$ supersymmetric mechanics. In this case, instead of the phase with standard $`๐ฉ=2`$ supersymmetry arising in the Kรคhler case, we shall get the phase with standard $`๐ฉ=4`$ supersymmetry. This system could be straightforwardly constructed by the dimensional reduction of the $`๐ฉ=2`$ supersymmetric $`(1+1)`$ dimensional sigma-model by Alvarez-Gaumรฉ and Freedman .
#### $`๐ฉ=8`$ mechanics
We have seen that the transition from the generic Riemann space to the generic Kรคhler space allows one to double the number of supersymmetries from $`๐ฉ=2`$ to $`๐ฉ=4`$, with the appropriate restriction of the admissible set of potentials.
On the other hand, we mentioned that the doubling of the number of odd variables and the restriction the Riemann metric allow one to construct the $`๐ฉ=4`$ supersymmetric mechanics . Now, following the paper , we shall show that a similar procedure, applied to the systems on Kรคhler manifolds, permits to construct the $`๐ฉ=8`$ supersymmetric mechanics, with the supersymmetry algebra <sup>1</sup><sup>1</sup>1We use the following convention: $`ฯต_{ij}ฯต^{jk}=\delta _i^k,ฯต_{12}=ฯต^{21}=1`$ .
$$\{Q_{i\alpha },Q_{j\beta }\}=\{\overline{Q}_{i\alpha },\overline{Q}_{j\beta }\}=0,\{Q_{i\alpha },\overline{Q}_{j\beta }\}=ฯต_{\alpha \beta }ฯต_{ij}_{SUSY},$$
(5.210)
where $`i,j=\mathrm{\hspace{0.33em}1},2`$, $`\alpha ,\beta =1,2`$.
We present the results for the mechanics without (bosonic) potential term. The respective systems with potential terms are constructed in .
In order to construct the $`๐ฉ=8`$ supersymmetric mechanics, let us define the $`(2d.4d)_{\text{I}\mathrm{C}}`$-dimensional symplectic structure
$$\mathrm{\Omega }=d๐=d\pi _adz^a+d\overline{\pi }_ad\overline{z}^aR_{a\overline{b}c\overline{d}}\eta _{i\alpha }^c\overline{\eta }^{d|i\alpha }dz^ad\overline{z}^b+g_{a\overline{b}}D\eta _{i\alpha }^aD\overline{\eta }^{b|i\alpha },$$
where
$$๐=\pi _adz^a+\overline{\pi }_ad\overline{z}^a+\frac{1}{2}\eta _{i\alpha }^ag_{a\overline{b}}D\overline{\eta }^{b|i\alpha }+\frac{1}{2}\overline{\eta }_{i\alpha }^bg_{a\overline{b}}D\eta ^{a|i\alpha },$$
(5.211)
and $`D\eta _{i\alpha }^a=d\eta _{i\alpha }^a+\mathrm{\Gamma }_{bc}^a\eta _{i\alpha }^bdz^c`$. The corresponding Poisson brackets are given by the following non-zero relations (and their complex-conjugates):
$$\begin{array}{c}\{\pi _a,z^b\}=\delta _a^b,\{\pi _a,\eta _{i\alpha }^b\}=\mathrm{\Gamma }_{ac}^b\eta _{i\alpha }^c,\\ \{\pi _a,\overline{\pi }_b\}=R_{a\overline{b}c\overline{d}}\eta _{i\alpha }^c\overline{\eta }^{d|i\alpha },\{\eta _{i\alpha }^a,\overline{\eta }^{b|j\beta }\}=g^{ab}\delta _i^j\delta _\alpha ^\beta .\end{array}$$
(5.212)
Let us search the supercharges among the functions
$$Q_{i\alpha }=\pi _a\eta _{i\alpha }^a+\frac{1}{3}\overline{f}_{abc}\overline{T}_{i\alpha }^{abc},\overline{Q}_{i\alpha }=\overline{\pi }_a\overline{\eta }_{i\alpha }^a+\frac{1}{3}f_{abc}T_{i\alpha }^{abc}$$
(5.213)
where $`T_{i\alpha }^{abc}\eta _{i\beta }^a\eta ^{bj\beta }\eta _{j\alpha }^c`$.
Calculating the mutual Poisson brackets of $`Q_{i\alpha },\overline{Q}_{i\alpha }`$ one can get, that they obey the $`๐ฉ=8`$ supersymmetry algebra, provided the following relations hold:
$$\frac{}{\overline{z}^d}f_{abc}=0,R_{a\overline{b}c\overline{d}}=f_{ace}g^{ee^{}}\overline{f}_{e^{}bd}.$$
(5.214)
The above equations guarantee, respectively, that the first and second equations in (5.210) are fulfilled. Then we could immediately get the $`๐ฉ=8`$ supersymmetric Hamiltonian
$$_{SUSY}=\pi _ag^{ab}\overline{\pi }_b+\frac{1}{3}f_{abc;d}\mathrm{\Lambda }^{abcd}+\frac{1}{3}\overline{f}_{abc;d}\overline{\mathrm{\Lambda }}^{abcd}+f_{abc}g^{c\overline{c}^{}}\overline{f}_{c^{}de}\mathrm{\Lambda }_0^{ab\overline{d}\overline{e}},$$
(5.215)
where
$$\mathrm{\Lambda }^{abcd}\frac{1}{4}\eta _{i\alpha }^a\eta ^{bi\beta }\eta _{k\beta }^c\eta ^{dk\alpha },\mathrm{\Lambda }_0^{ab\overline{c}\overline{d}}\frac{1}{2}(\eta _i^{a\alpha }\eta _{j\alpha }^b\overline{\eta }^{c\beta i}\overline{\eta }_\beta ^{dj}+\eta _\alpha ^{ai}\eta _{i\beta }^b\overline{\eta }^{cj\alpha }\overline{\eta }_j^{d\beta }),$$
and $`f_{abc;d}=f_{abc,d}\mathrm{\Gamma }_{da}^ef_{ebc}\mathrm{\Gamma }_{db}^ef_{aec}\mathrm{\Gamma }_{dc}^ef_{abe}`$ is the covariant derivative of the third rank covariant symmetric tensor.
The equations (5.214) precisely mean that the configuration space $`M_0`$ is a special Kรคhler manifold of the rigid type . Taking into account the symmetrizing of $`f_{abc}`$ over spatial indices and the explicit expression of $`R_{a\overline{b}c\overline{d}}`$ in terms of the metric $`g_{ab}`$, we can immediately find the local solution for equations (5.214)
$$f_{abc}=\frac{^3f(z)}{z^az^bz^c},g_{a\overline{b}}=\mathrm{e}^{i\nu }\frac{^2f(z)}{z^az^b}+\mathrm{e}^{i\nu }\frac{^2\overline{f}(\overline{z})}{\overline{z}^a\overline{z}^b},$$
(5.216)
where $`\nu =const\mathrm{IR}`$.
Redefining the local function, $`fi\mathrm{e}^{i\nu }f`$, we shall get the $`\nu `$-parametric family of supersymmetric mechanics, whose metric is defined by the Kรคhler potential of a special Kรคhler manifold of the rigid type. Surely, this local solution is not covariant under arbitrary holomorphic transformation, and it assumes the choice of a distinguished coordinate frame.
The special Kรคhler manifolds of the rigid type became widely known during last decade due to the so-called โT-duality symmetryโ: in the context of $`๐ฉ=2,d=4`$ super Yang-Mills theory, it connects the UV and IR limit of the theory. The โT-duality symmetryโ is expressed in the line below
$$(z^a,f(z))\left(u_a=\frac{f(z)}{z^a},\stackrel{~}{f}(u)\right),\mathrm{where}\frac{^2\stackrel{~}{f}(u)}{u_au_c}\frac{f}{z^cz^b}=\delta _b^a.$$
(5.217)
It is clear that the symplectic structure is covariant under the following holomorphic transformations:
$$\stackrel{~}{z}^a=\stackrel{~}{z}^a(z),\stackrel{~}{\eta }_{i\alpha }^a=\frac{\stackrel{~}{z}^a(z)}{z^b}\eta _{i\alpha }^b,\stackrel{~}{\pi }_a=\frac{z^b}{\stackrel{~}{z}^a}\pi _b,$$
(5.218)
By the use of (5.218), we can extend the duality transformation (5.217) to the whole phase superspace $`(\pi _a,z^a,\eta _{i\alpha }^a)(p^a,u_a,\psi _{a|i\alpha })`$
$$u_a=_af(z),p^a\frac{^2f}{z^az^b}=\pi _b,\psi _a^{i\alpha }=\frac{^2f}{z^az^b}\eta ^{b|i\alpha }.$$
(5.219)
Taking into account the expression of the symplectic structure in terms of the presymplectic one-form (5.211), we can easily perform the Legendre transformation of the Hamiltonian to the (second-order) Lagrangian
$`=๐(d/dt)_{SUSY}|_{\pi _a=g_{a\overline{b}}\dot{\overline{z}}^b}=`$
$`=g_{a\overline{b}}\dot{z}^a\dot{\overline{z}}^b+{\displaystyle \frac{1}{2}}\eta _{i\alpha }^ag_{a\overline{b}}{\displaystyle \frac{D\overline{\eta }^{b|i\alpha }}{dt}}+{\displaystyle \frac{1}{2}}\overline{\eta }_{i\alpha }^bg_{a\overline{b}}{\displaystyle \frac{D\eta ^{a|i\alpha }}{dt}}`$ (5.220)
$$\frac{1}{3}f_{abc;d}\mathrm{\Lambda }^{abcd}\frac{1}{3}\overline{f}_{abc;d}\overline{\mathrm{\Lambda }}^{abcd}f_{abc}g^{c\overline{c}^{}}\overline{f}_{c^{}de}\mathrm{\Lambda }_0^{ab\overline{d}\overline{e}}.$$
Here we denoted $`d/dt=\dot{z}^a/z^a+\dot{\eta }_{i\alpha }^a/\eta _{i\alpha }^a+c.c.`$. Clearly, the Lagrangian (5.220) is covariant under holomorphic transformations (5.218), and the duality transformation as well. The prepotential $`\stackrel{~}{f}(u)`$ is connected with $`f(z)`$ by the Legendre transformation
$$\stackrel{~}{f}(u)=\stackrel{~}{f}(u,z)|_{u_a=_af(z)},\stackrel{~}{f}(u,z)=u_az^af(z).$$
### Supersymmetric Kรคhler Oscillator
So far, the Kรคhler structure allowed us to double the number of supersymmetries in the system. One can hope that in some cases this could be preserved after inclusion of constant magnetic field, since this field usually respects the Kรคhler structure. We shall show, on the example of the Kรคhler oscillator (3.126), that it is indeed a case.
Let us consider, following , the supersymmetrization of a specific model of Hamiltonian mechanics on the Kรคhler manifold $`(M_0,g_{a\overline{b}}dz^ad\overline{z}^{\overline{b}})`$ interacting with the constant magnetic field $`B`$, viz
$$=g^{a\overline{b}}(\pi _a\overline{\pi }_b+\alpha ^2_aK\overline{}_bK),\mathrm{\Omega }_0=d\pi _adz^a+d\overline{\pi }_ad\overline{z}^a+iBg_{a\overline{b}}dz^ad\overline{z}^b,$$
(5.221)
where $`K(z,\overline{z})`$ is a Kรคhler potential of configuration space.
Remind, that the Kรคhler potential is defined up to holomorphic and antiholomorphic terms, $`KK+U(z)+\overline{U}(\overline{z})`$. Hence, in the limit $`\omega 0`$ the above Hamiltonian takes the form
$$=g^{a\overline{b}}(\pi _a\overline{\pi }_b+_aU(z)\overline{}_b\overline{U}(\overline{z})),$$
(5.222)
i.e.it admits, in the absence of magnetic field, a $`๐ฉ=4`$ superextension.
Notice, also, that in the โlarge mass limitโ, $`\pi _a0`$, this system results in the following one:
$$_0=\omega ^2g^{a\overline{b}}_aK\overline{}_bK,\mathrm{\Omega }_0=iBg_{a\overline{b}}dz^ad\overline{z}^b,$$
which could be easily extended to $`๐ฉ=2`$ supersymmetric mechanics.
We shall show that, although the system under consideration does not possess a standard $`๐ฉ=4`$ superextension, it admits a superextension in terms of a nonstandard superalgebra with four fermionic generators, including, as subalgebras, two copies of the $`๐ฉ=2`$ superalgebra. This nonstandard superextension respects the inclusion of a constant magnetic field.
We use the following strategy. At first, we extend the initial phase space to a $`(2N.2N)_{\text{I}\mathrm{C}}`$-dimensional superspace equipped with the symplectic structure
$$\mathrm{\Omega }=\mathrm{\Omega }_BiR_{a\overline{b}c\overline{d}}\eta _i^c\overline{\eta }_i^d)dz^ad\overline{z}^b+g_{a\overline{b}}D\eta ^a_iD\overline{\eta }^b_i,$$
(5.223)
where $`\mathrm{\Omega }_B`$ is given by (1.36). The corresponding Poisson brackets are defined by the following non-zero relations (and their complex-conjugates):
$$\begin{array}{c}\{\pi _a,z^b\}=\delta _a^b,\{\pi _a,\eta _i^b\}=\mathrm{\Gamma }_{ac}^b\eta _i^c,\\ \{\pi _a,\overline{\pi }_b\}=i(Bg_{a\overline{b}}+iR_{a\overline{b}c\overline{d}}\eta _i^c\overline{\eta }_i^d),\{\eta _i^a,\overline{\eta }_j^b\}=g^{a\overline{b}}\delta _{ij}.\end{array}$$
(5.224)
Then, in order to construct the system with the exact $`๐ฉ=2`$ supersymmetry (5.183), we shall search for the odd functions $`Q^\pm `$, which obey the equations $`\{Q^\pm ,Q^\pm \}=0`$ (we restrict ourselves to the supersymmetric mechanics whose supercharges are linear in the Grassmann variables $`\eta _i^a`$, $`\overline{\eta }_i^{\overline{a}}`$).
Let us search for the realization of supercharges among the functions
$$Q^\pm =\mathrm{cos}\lambda \mathrm{\Theta }_1^\pm +\mathrm{sin}\lambda \mathrm{\Theta }_2^\pm ,$$
(5.225)
where
$$\mathrm{\Theta }_1^+=\pi _a\eta _1^a+i\overline{}_aW\overline{\eta }_2^a,\mathrm{\Theta }_2^+=\overline{\pi }_a\overline{\eta }_2^a+i_aW\eta _1^a,\mathrm{\Theta }_{1,2}^{}=\overline{\mathrm{\Theta }}_{1,2}^+,$$
(5.226)
and $`\lambda `$ is some parameter.
Calculating the Poisson brackets of the functions, we get
$`\{Q^\pm ,Q^\pm \}=`$ $`i(B\mathrm{sin}2\lambda +2\alpha \mathrm{cos}2\lambda )_\pm ,`$ (5.227)
$`\{Q^+,Q^{}\}=`$ $`_{SUSY}^0+\left(B\mathrm{cos}2\lambda 2\alpha \mathrm{sin}2\lambda \right)_3/2.`$ (5.228)
Here and further, we use the notation
$$_{SUSY}^0=R_{a\overline{b}c\overline{d}}\eta _1^a\overline{\eta }_1^b\eta _2^c\overline{\eta }_2^diW_{a;b}\eta _1^a\eta _2^b+iW_{\overline{a};\overline{b}}\overline{\eta }_1^a\overline{\eta }_2^b+B\frac{ig_{a\overline{b}}\eta _i^a\overline{\eta }_i^b}{2},$$
(5.229)
where $``$ denotes the oscillator Hamiltonian (3.126), and
$$๐
=\frac{i}{2}g_{a\overline{b}}\eta _i^a\overline{\eta }_j^b๐_{ij},_\pm =F_1\pm F_2.$$
(5.230)
One has, then
$$\{Q^\pm ,Q^\pm \}=0B\mathrm{sin}\mathrm{\hspace{0.33em}2}\lambda +2\alpha \mathrm{cos}\mathrm{\hspace{0.33em}2}\lambda =0,$$
(5.231)
so that $`\lambda =\lambda _0+(i1)\pi /2`$, $`i=1,2.`$
Here the parameter $`\lambda _0`$ is defined by the expressions
$$\mathrm{cos}2\lambda _0=\frac{B/2}{\sqrt{\alpha ^2+(B/2)^2}},\mathrm{sin}2\lambda _0=\frac{\alpha }{\sqrt{\alpha ^2+(B/2)^2}}.$$
(5.232)
Hence, we get the following supercharges:
$$Q_\nu ^\pm =\mathrm{cos}\lambda _0\mathrm{\Theta }_1^\pm +(1)^\nu \mathrm{sin}\lambda _0\mathrm{\Theta }_2^\pm ,$$
(5.233)
and the pair of $`๐ฉ=2`$ supersymmetric Hamiltonians
$$_{SUSY}^i=\{Q_\nu ^+,Q_\nu ^{}\}=_{SUSY}^0(1)^i\sqrt{\alpha ^2+(B/2)^2}_3.$$
(5.234)
Notice that the supersymmetry invariance is preserved in the presence of a constant magnetic field.
Calculating the commutators of $`Q_1^\pm `$ and $`Q_2^\pm `$, we get
$$\{Q_1^\pm ,Q_2^\pm \}=2\sqrt{\alpha ^2+(B/2)^2}_\pm ,\{Q_1^+,Q_2^{}\}=0.$$
(5.235)
The Poisson brackets between $`_\pm `$ and $`Q_\nu ^\pm `$ look as follows:
$$\begin{array}{c}\{Q_i^\pm ,_\pm \}=0,\{Q_i^\pm ,_{}\}=\pm ฯต_{ij}Q_j^\pm ,\{Q_i^\pm ,_3\}=\pm iQ_i^\pm .\end{array}$$
(5.236)
In the notation $`S_1^\pm Q_1^\pm ,S_2^\pm =Q_2^{}`$ the whole superalgebra reads
$$\begin{array}{c}\{S_i^\pm ,S_j^{}\}=\delta _{ij}_{SUSY}^0+\mathrm{\Lambda }\sigma _{ij}^\mu _\mu ,\\ \{S_i^\pm ,_\mu \}=\pm ฤฑ\sigma _{ij}^\mu S_j^\pm ,\{_\mu ,_\nu \}=ฯต_{\mu \nu \rho }_\rho ,\end{array}$$
(5.237)
where
$$\mathrm{\Lambda }=\sqrt{\omega ^2+(B/2)^2}.$$
(5.238)
This is precisely the weak supersymmetry algebra considered by A. Smilga . In the particular case $`\omega =0`$ it yields the $`๐ฉ=4`$ supersymmetric mechanics broken by the presence of a constant magnetic field.
Let us notice the $`\alpha `$ and $`B`$ appear in this superalgebra in a symmetric way, via the factor $`\sqrt{\alpha ^2+(B/2)^2}`$.
Remark In the case of the oscillator on $`\text{I}\mathrm{C}^\mathrm{N}`$ we can smoothly relate the above supersymmetric oscillator with a $`N=\mathit{4}`$ oscillator, provided we choose
$$K=\mathrm{cos}\gamma z\overline{z}+\mathrm{sin}\gamma (z^2+\overline{z}^2)/2,\gamma [0,\pi /2].$$
Hence,
$$=\pi \overline{\pi }+\alpha _0^2z\overline{z}+\mathrm{sin}2\gamma \alpha _0^2(z^2+\overline{z}^2)/2,$$
i.e. for $`\gamma =0,\pi /2`$ we have a standard harmonic oscillator, while for $`\gamma 0,\pi /2`$ we get the anisotropic one, which is equivalent to two sets of $`N`$ one-dimensional oscillators with frequencies $`\alpha _0\sqrt{1\pm \mathrm{sin}2\gamma }`$. The frequency $`\alpha `$ appearing in the superalgebra, is of the form: $`\alpha =\alpha _0\mathrm{cos}\gamma `$.
## Conclusion
We presented some constructions of the Hamiltonian formalism related with Hopf maps and Kรคhler geometry, and a few models of supersymmetric mechanics on Kรคhler manifolds. One can hope that the former constructions could be useful in supersymmetric mechanics along the following lines. Firstly, one could try to extend the number of supersymmetries, passing from the Kรคhler manifolds to quaternionic ones. The model suggested in indicates that this could indeed work. One could also expect that the latter system will respect the inclusion of an instanton field. Secondly, one can try to construct the supersymmetric mechanics, performing the Hamiltonian reduction of the existing systems, related with the Hopf maps. In this way one could get new supersymmetric models, specified by the presence of Dirac and Yang monopoles, as well as with constant magnetic and instanton fields.
Acknowledgements I would like to thank Stefano Bellucci for organizing the Winter School on Modern trends in supersymmetric mechanics and inviting me to deliver (and to write down) these lectures, which are the result of extensive discussions with him, and with S.Krivonos and V.P.Nair. Most of the examples included in the paper were obtained in collaboration. I am indebted to all of my co-authors, especially to A.Yeranian, O.Khudaverdian, V.Ter-Antonyan, P.-Y.Casteill.
Special thanks to Stefano Bellucci for careful reading of manuscript and substantial improving of the English.
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# Tests of Lorentz violation in ๐ฬ_๐โ๐ฬ_๐ oscillations
## I Introduction
Lorentz symmetry is one of the most fundamental ideas of both relativistic local quantum field theory and general relativity. Early tests, such as the Michelson-Morley and Kennedy-Thorndike experiments have established that Lorentz symmetry is an exact symmetry of nature. So it is natural to assume that Lorentz symmetry is an exact symmetry in the standard model (SM) of particle physics. However, since the SM does not address gravity, a fundamental theory of Planck-scale physics ($`M_P10^{19}`$ GeV), including string theory spontaneous and quantum gravity Hawking , may violate Lorentz and CPT symmetry CPT04 .
If limited to conventional relativistic quantum mechanics, it is possible to establish a self-consistent low-energy effective theory with Lorentz and CPT violation; this is called the standard-model Extension (SME) Colladay . The minimal-SME formalism has all the conventional properties of the standard model including observer Lorentz covariance, power counting renormalizability, energy momentum conservation, quantized field, micro causality, and spin-statistics with particle Lorentz and CPT violation due to background Lorentz tensor fields of the universe. The minimal SME also has $`SU(3)_C\times SU(2)_L\times U(1)_Y`$ gauge invariance. Since the background Lorentz tensor fields are fixed in spacetime, by definition, they do not transform under an active transformation law. That implies rotation and boost dependence of physics in a specific frame. This formalism focuses on the inverse Planck-scale effect which is believed to be suppressed by at least one order of the inverse Planck mass ($`\frac{E}{M_P}`$, where $`E`$ is the energy scale of the system under consideration). Therefore, the physics quantities involved in the formalism are perturbative.
Surprisingly, atomic physics has achieved this sensitivity level, and extensive experimental studies have been done (see, for example, Ref. CPT04 ). A recent experiment Walsworth of this type reaches a sensitivity to a specific combination of SME coefficients to order $`10^{32}`$ GeV, well beyond a basic estimate of the scale of new physics. In addition, spectral polarimetry of distant cosmological sources yields a similar sensitivity for another combination of SME coefficients spectro . However, many of the SME coefficients still have no experimental bounds.
Similarly, quantum interference experiments, such as meson oscillations, are also sensitive to the small effect of Lorentz and CPT violation meson . Tests have been made using data from many experiments, including KTeV KTeV , FOCUS FOCUS , BaBar BaBar , BELLE BELLE , and OPAL OPAL . Recently, the SME formalism for neutrino oscillations, another type of quantum interference experiment, has become available oscillation .
## II The LSND Evidence for Neutrino Oscillations
The Liquid Scintillator Neutrino Detector (LSND) experiment LSND2 , completed at the Los Alamos National Laboratory (LANL), observed an excess of $`\overline{\nu }_e`$ in a beam of $`\overline{\nu }_\mu `$ created from $`\mu ^+`$ decay at rest. The data analysis used the sample of detected $`\overline{\nu }_epe^+n`$ events with positron energy $`20<E_{e^+}<60\mathrm{MeV}`$. If interpreted as $`\overline{\nu }_\mu `$ to $`\overline{\nu }_e`$ oscillations, this $`\overline{\nu }_e`$ excess implies a two-neutrino oscillation probability of $`(0.264\pm 0.067\pm 0.045)\%`$. Here the first error is statistical and the second error is systematic (neutrino flux, particle detection efficiency, cross sections, etc.). Despite the evidence for neutrino oscillations from solar neutrinos Homestake ; kamiokande ; GALLEX ; SAGE ; GNO ; SNO , atmospheric neutrinos Super-K ; MACRO , accelerator neutrinos K2K , and reactor neutrinos KamLAND , the oscillation signal observed at LSND remains a puzzle. Since the neutrino sector is thought as likely to reveal new physics, the LSND anomaly is often explained with new ideas such as a mass-difference CPT-violating model (see, for example Ref. Barenboim ) or sterile neutrino models (see Ref. Janet for a recent example). The MiniBooNE experiment MiniBooNE at Fermilab is currently taking data to test the LSND signal.
## III Lorentz Violating Neutrino Oscillations
Perhaps LSND is seeing the first signal of Planck-scale physics TKRTCPT04 . To describe neutrino oscillations, including Lorentz and CPT violation, a recently developed formalism for neutrino oscillations oscillation ; SBL using the SME framework Colladay is employed. This framework allows for a sidereal time variation of the neutrino oscillation probability.
Within the SME framework, the neutrino free field Lagrangian becomes,
$``$ $`=`$ $`{\displaystyle \frac{1}{2}}i\overline{\psi }_A\mathrm{\Gamma }_{AB}^\mu \stackrel{}{D_\mu }\psi _B\overline{\psi }_AM_{AB}\psi _B,`$ (1)
$`\mathrm{\Gamma }_{AB}^\mu `$ $`=`$ $`\gamma ^\mu \delta _{AB}+c_{AB}^{\mu \nu }\gamma _\nu +d_{AB}^{\mu \nu }\gamma _5\gamma _\nu `$ (2)
$`+e_{AB}^\mu +if_{AB}^\mu \gamma _5+{\displaystyle \frac{1}{2}}g_{AB}^{\mu \nu \lambda }\sigma _{\nu \lambda },`$
$`M_{AB}`$ $`=`$ $`m_{AB}+im_{5AB}\gamma _5`$ (3)
$`+a_{AB}^\mu \gamma _\mu +b_{AB}^\mu \gamma _5\gamma _\mu +{\displaystyle \frac{1}{2}}H_{AB}^{\mu \nu }\sigma _{\mu \nu }.`$
The first term of $`\mathrm{\Gamma }_{AB}^\mu `$ and the first and second terms of $`M_{AB}`$ are the only nonzero terms in the case of conventional neutrino oscillations. The remaining terms in this Lagrangian represent the physics of the background fields. In general, the background Lorentz tensor fields are an infinite series, but if the focus is on a low-energy effective theory, these eight additional fields are complete. Here, vacuum expectation values that contain $`c_{AB}^{\mu \nu }`$, $`d_{AB}^{\mu \nu }`$, and $`H_{AB}^{\mu \nu }`$ are CPT-even terms while $`e_{AB}^\mu `$, $`f_{AB}^\mu `$, $`g_{AB}^{\mu \nu \lambda }`$, $`a_{AB}^\mu `$, and $`b_{AB}^\mu `$ are CPT-odd by definition of the background fields. Notice that each background field has flavor indices (A and B) that, unlike other systems, bring additional complication for the neutrino sector.
This Lagrangian leads to the modified Dirac equation,
$$(i\mathrm{\Gamma }_{AB}^\mu _\mu M_{AB})\psi _B=0.$$
(4)
After some manipulation, this yields the effective Hamiltonian for active neutrino oscillations oscillation . In particular, the effective Hamiltonian for active antineutrino to antineutrino oscillations is,
$`(h_{\mathrm{eff}})_{ab}`$ $`=`$ $`|\stackrel{}{p}|\delta _{ab}+\frac{(\stackrel{~}{m}^2)_{ab}^{}}{2|\stackrel{}{p}|}`$ (5)
$`+\frac{1}{|\stackrel{}{p}|}[(a_L)^\mu p_\mu (c_L)^{\mu \nu }p_\mu p_\nu ]_{ab}^{}.`$
Here, the effective Hamiltonian is a $`3\times 3`$ flavor Majorana basis matrix of three active, right-handed, antineutrinos. The original effective Hamiltonian oscillation can describe $`\nu \nu `$, $`\overline{\nu }\overline{\nu }`$, and $`\nu \overline{\nu }`$ oscillations, but, in this work, lepton-number violating $`\nu \overline{\nu }`$ oscillations are not considered. Therefore, the neutrino and antineutrino sectors can be diagonalized separately. There is some coupling of SME coefficients, $`(a_L)_{ab}^\mu =(a)_{ab}^\mu +(b)_{ab}^\mu `$ and $`(c_L)_{ab}^{\mu \nu }=(c)_{ab}^{\mu \nu }+(d)_{ab}^{\mu \nu }`$. Also, other types of SME coefficients do not show up in this analysis. For the usual conventional neutrino oscillation case, the effective Hamiltonian (Eq. 5) contains only the first two terms. Then, the neutrino oscillation probability depends on $`\mathrm{\Delta }m^2`$ and the mixing matrix. But, in this general form, including possible Lorentz and CPT violation, the diagonalization of the effective Hamiltonian is more complicated and, in general, it can not be represented by $`\mathrm{\Delta }m^2`$ and the mixing matrix alone.
## IV The Short-Baseline Approximation
If the baseline of the neutrino beam is short compared with the neutrino oscillation length, $`L`$, the neutrino oscillation probability can be expanded with an effective Hamiltonian. Expressed to leading order in $`h_{\mathrm{eff}}`$ SBL ,
$$P_{\overline{\nu }_\mu \overline{\nu }_e}\frac{|(h_{\mathrm{eff}})_{\overline{e}\overline{\mu }}|^2L^2}{(\mathrm{}c)^2}.$$
(6)
Note that in this equation, unlike the equations above, $`\mathrm{}`$ and $`c`$ have been explicitly included. Since, in the effective Hamiltonian, $`p^\mu `$ contains information about the propagation direction of the neutrino, this oscillation probability depends on the neutrino propagation direction. In order to form a phenomenological expression for the neutrino oscillation probability, it is most convenient to use a coordinate system fixed to the experiment spectro ; Bluhm . The standard choice is a Sun-centered system (Fig.1a) that is, to a good approximation, an inertial frame for the experiment.
With this choice of coordinates, the neutrino oscillation probability becomes,
$`P_{\overline{\nu }_\mu \overline{\nu }_e}`$ $``$ $`{\displaystyle \frac{L^2}{(\mathrm{}c)^2}}|(๐)_{\overline{e}\overline{\mu }}+(๐_s)_{\overline{e}\overline{\mu }}\mathrm{sin}\omega _{}T_{}`$ (7)
$`+(๐_c)_{\overline{e}\overline{\mu }}\mathrm{cos}\omega _{}T_{}+(_s)_{\overline{e}\overline{\mu }}\mathrm{sin}2\omega _{}T_{}`$
$`+(_c)_{\overline{e}\overline{\mu }}\mathrm{cos}2\omega _{}T_{}|^2,`$
where $`\omega _{}`$ is the sidereal frequency (=2$`\pi `$/23h 56min 4.1s) and $`T_{}`$ is the sidereal time as measured from a standard origin. Note that $`P_{\overline{\nu }_\mu \overline{\nu }_e}`$ may depend on the sidereal time.
These parameters, $`(๐)_{\overline{e}\overline{\mu }}`$, $`(๐_s)_{\overline{e}\overline{\mu }}`$, $`(๐_c)_{\overline{e}\overline{\mu }}`$, $`(_s)_{\overline{e}\overline{\mu }}`$, and $`(_c)_{\overline{e}\overline{\mu }}`$, depend on the SME coefficients $`(a_L)^\mu `$ and $`(c_L)^{\mu \nu }`$ and the neutrino propagation direction unit vectors $`\widehat{N}^X`$, $`\widehat{N}^Y`$, and $`\widehat{N}^Z`$ in the Sun-centered system. The direction unit vectors depend on the colatitude $`\chi `$ of the experiment in the Earth-centered system (Fig.1b) and the zenith and azimuthal angles $`\theta `$ and $`\varphi `$ of the $`\overline{\nu }_\mu `$ beam in the experiment local coordinate system (Fig.1c).
$$\left(\begin{array}{c}\widehat{N}^X\\ \widehat{N}^Y\\ \widehat{N}^Z\end{array}\right)=\left(\begin{array}{c}\mathrm{cos}\chi \mathrm{sin}\theta \mathrm{cos}\varphi +\mathrm{sin}\chi \mathrm{cos}\theta \\ \mathrm{sin}\theta \mathrm{sin}\varphi \\ \mathrm{sin}\chi \mathrm{sin}\theta \mathrm{cos}\varphi +\mathrm{cos}\chi \mathrm{cos}\theta \end{array}\right)$$
(8)
For neutrinos from the Los Alamos Neutron Science Center (LANSCE) beam to the LSND detector, $`\chi =54.1^{}`$, $`\theta =99.0^{}`$ and $`\varphi =82.6^{}`$ GPS . The sidereal time is defined to be zero ($`T_{}=0`$) at LANL local midnight on the autumnal equinox (Fig.1d). At that time, the $`y`$ axis of the Earth-centered system coincides with the $`Y`$ axis of the Sun-centered system. The estimated error using this definition is three minutes, which is small compared to the time scale sensitivity of this analysis.
Combining these values for neutrino propagation unit vectors with the detailed expression of the parameters, $`(๐)_{\overline{e}\overline{\mu }}`$, $`(๐_s)_{\overline{e}\overline{\mu }}`$, $`(๐_c)_{\overline{e}\overline{\mu }}`$, $`(_s)_{\overline{e}\overline{\mu }}`$, and $`(_c)_{\overline{e}\overline{\mu }}`$ SBL , yields, for the particular case of the LSND experiment:
$`(๐)_{\overline{e}\overline{\mu }}`$ $`=`$ $`\frac{(\stackrel{~}{m}^2)_{\overline{e}\overline{\mu }}}{2E}+[(a_L)_{\overline{e}\overline{\mu }}^T+0.19(a_L)_{\overline{e}\overline{\mu }}^Z]`$ (9)
$`+E[1.48(c_L)_{\overline{e}\overline{\mu }}^{TT}0.39(c_L)_{\overline{e}\overline{\mu }}^{TZ}`$
$`+0.44(c_L)_{\overline{e}\overline{\mu }}^{ZZ}],`$
$`(๐_s)_{\overline{e}\overline{\mu }}`$ $`=`$ $`[0.98(a_L)_{\overline{e}\overline{\mu }}^X+0.053(a_L)_{\overline{e}\overline{\mu }}^Y]`$ (10)
$`+E[1.96(c_L)_{\overline{e}\overline{\mu }}^{TX}0.11(c_L)_{\overline{e}\overline{\mu }}^{TY}`$
$`0.38(c_L)_{\overline{e}\overline{\mu }}^{XZ}0.021(c_L)_{\overline{e}\overline{\mu }}^{YZ}],`$
$`(๐_c)_{\overline{e}\overline{\mu }}`$ $`=`$ $`[0.053(a_L)_{\overline{e}\overline{\mu }}^X0.98(a_L)_{\overline{e}\overline{\mu }}^Y],`$ (11)
$`+E[0.11(c_L)_{\overline{e}\overline{\mu }}^{TX}+1.96(c_L)_{\overline{e}\overline{\mu }}^{TY}`$
$`0.021(c_L)_{\overline{e}\overline{\mu }}^{XZ}+0.38(c_L)_{\overline{e}\overline{\mu }}^{YZ}],`$
$`(_s)_{\overline{e}\overline{\mu }}`$ $`=`$ $`E[0.052((c_L)_{\overline{e}\overline{\mu }}^{XX}(c_L)_{\overline{e}\overline{\mu }}^{YY})`$ (12)
$`+0.96(c_L)_{\overline{e}\overline{\mu }}^{XY}],`$
$`(_c)_{\overline{e}\overline{\mu }}`$ $`=`$ $`E[0.48((c_L)_{\overline{e}\overline{\mu }}^{XX}(c_L)_{\overline{e}\overline{\mu }}^{YY})`$ (13)
$`+0.10(c_L)_{\overline{e}\overline{\mu }}^{XY}].`$
In Eq. 9, the mass-squared term, $`\stackrel{~}{m}_{\overline{e}\overline{\mu }}^2`$, has been included. This allows for conventional massive-neutrino oscillations in addition to the Lorentz-violation oscillations. It is assumed that the size of this term does not invalidate the short-baseline approximation SBL .
## V Sidereal Time Distribution of the LSND Data
In conventional explanations of neutrino oscillations, the oscillation probability is independent of sidereal time and, therefore, the sidereal time distribution of oscillation events is expected to be constant. In the Lorentz and CPT violating model of neutrino oscillations considered here, nonzero values of the model parameters could exhibit themselves as modulations to the sidereal time distribution (as in Eq. 7). The sidereal time dependence of candidate oscillation events from the LSND data sample has been examined and subjected to statistical tests to quantify any evidence for a sidereal variation.
In the analysis of the final LSND data set LSND2 , 205 neutrino oscillation candidate events were reported with positron energy in the range $`20<E_{e^+}<60\mathrm{MeV}`$ and with an identified neutron-capture photon. There are two classes of background in the oscillation sample: beam-unrelated and beam-related ($`\nu `$-induced). The beam-unrelated backgrounds arise from cosmic ray processes. It is measured in beam-off data and then subtracted from the beam-on data. The beam-related backgrounds are calculated from known neutrino (nonoscillation) interactions.
The neutrino beam used for the LSND experiment was produced using protons from the LANSCE accelerator LSND-NIM . The proton beam was delivered at approximately 100Hz in pulses of 600 $`\mu `$s duration. The detector was triggered independently of the state of the beam and the beam status was recorded. In this manner, beam-off data was taken continuously in the time between beam pulses. The resulting beam-off data set was approximately 16 times larger than the beam-on data set. This allowed for an accurate measurement of the beam-unrelated background by weighting the beam-off data by the beam duty-factor (calculated for each run).
The estimated number of beam-unrelated and $`\nu `$-induced background events in this sample are $`106.8\pm 2.5`$ and $`39.2\pm 3.1`$, respectively. These events were collected during experimental running in 1993 through 1998. There were six sets of runs, one in each of these years. The GPS (Global Positioning System) time stamp, necessary for this analysis, was not included into the LSND data stream until midway through the 1994 run period. Because of this, only 186 of the 205 oscillation candidate events could be used in this analysis. The expected numbers of beam-off and $`\nu `$-induced backgrounds in this smaller sample are $`94.0\pm 2.3`$ and $`35.6\pm 2.8`$, respectively.
Ideally, an experiment to search for sidereal variations in a signal would run continuously throughout the calendar year so that one particular sidereal time bin would be drawn from the entire range of local time. This was not the case with LSND, but runs did cover the space of local time vs sidereal time with reasonable completeness, as can be seen in Figure 2. The Los Alamos (clock) time can be determined from Greenwich Mean (GM) time by subtracting 6 (7) hours in the summer (winter).
To quantify the statistical significance of any sidereal time variation in the data, we employed two different statistical tests: a Pearsonโs-$`\chi ^2`$ test Frodesen ; PDG and a Kolmogorov-Smirnov (KS) test Frodesen . In both of these tests, the data were compared to the (null) hypothesis that the event rate is constant in sidereal time. Note that this null hypothesis is not that no oscillation signal exists, but only that the signal is constant in time. We also examined the GM time distributions and applied these tests with the null hypothesis of an underlying distribution that is constant in GM time.
The Pearsonโs-$`\chi ^2`$ (P-$`\chi ^2`$), as implemented in this analysis, is
$$\mathrm{P}\chi ^2=\underset{i=1}{\overset{N}{}}\frac{(n_i\nu _i)^2}{\nu _i};\nu _i=n/N,$$
(14)
where $`n`$ is the total number of events in the sample, $`N`$ is the number of time bins, and $`n_i`$ is the measured number of events in time bin $`i`$. The predicted number of events in each time bin, $`\nu _i`$, is constant for each time bin. Note that this quantity is constructed with the variance of the expected number of events in the denominator.
The P-$`\chi ^2`$ statistic will follow, in the absence of sidereal time variations and with sufficient events per time bin, a $`\chi ^2`$ distribution with number of degrees of freedom equal to the number of time bins minus one Frodesen ; PDG . The standard criterion for sufficient events is that $`\nu _i5`$ Frodesen ; PDG . The binning for the beam-on data has been chosen to satisfy this. The $`p`$-value, ($`P(\chi ^2`$), one minus the $`\chi ^2`$ cumulative distribution) can be extracted and interpreted as a confidence level that the null hypothesis explains the data.
The KS test has the advantage that it works with unbinned data, thus eliminating the need to choose a binning. It involves a comparison between the data and the null hypothesis via cumulative distributions. Unlike the P-$`\chi ^2`$ test, it is sensitive to โrunsโ in the data, thus making the P-$`\chi ^2`$ and KS tests complementary. The KS statistics reported here are the maximum cumulative deviation, $`D_n`$, and the KS probability, $`P(\mathrm{KS})`$. The quantity $`P(\mathrm{KS})`$, which is obtained from the known distribution of $`D_n`$ Frodesen , can be interpreted as a confidence level that the data is explained by the null hypothesis.
In the LSND data set considered for this analysis, there were 1656 beam-off events passing the neutrino oscillation cuts. These events, after weighting for the beam-on duty-factor, determine the number and distributions of beam-unrelated background events in the beam-on sample. The distribution of sidereal and GM times in 37 time bins for these beam-off events is shown in Fig. 3. The number of time bins chosen for this distribution was obtained by applying the $`N>5`$ criterion for the beam-on data. The errors shown in Fig. 3 (and subsequent figures) are the square root of the number of counts in the bin. Note that these errors are not used in the calculation of the P-$`\chi ^2`$ (Eq. 14). The P-$`\chi ^2`$ is 29.6 for 37 sidereal time bins corresponding to $`P(\chi ^2)=0.77`$. The KS test on this same data yields $`D_n=0.019`$ and $`P(\mathrm{KS})=0.60`$. These results indicate that the beam-off data are in reasonable agreement with the null hypothesis (no sidereal time dependence). The GM time distribution yields a slightly low $`P(\mathrm{KS})=0.01`$, however, for this same distribution $`P(\chi ^2)=0.29`$. In addition, any GM time variations are distributed throughout a range in sidereal time. For these reasons, we conclude that there is no evidence for substantial environmental or โday-nightโ sidereal variations in the beam-unrelated backgrounds.
The sidereal and GM time distributions of the 186 oscillation candidate events are shown in Figure 4. The P-$`\chi ^2`$ for the sidereal time distributions is 44.8 for 37 time bins. The corresponding $`p`$-value for the sidereal time distribution is $`P(\chi ^2)=0.15`$. A KS test applied to these distributions yields $`P(\mathrm{KS})=0.234`$. The sidereal time distribution is slightly less compatible with no time variation as is evident in both of these statistical tests. However, the variation is not statistically significant. A KS test between beam-on and beam-off data was also applied and shows compatibility between the two data sets. The complete results from the statistical tests on the sidereal and GM time distributions for beam-on and beam-off data are summarized in Table 1.
To check the underlying assumption that the beam was delivered with equal efficiency throughout the sidereal day, a sample of $`{}_{}{}^{12}C(\nu _e,e^{})^{12}N_{\mathrm{g}.\mathrm{s}.}`$ events was obtained by applying cuts to select for subsequent $`\beta `$-decays of $`{}_{}{}^{12}N_{\mathrm{g}.\mathrm{s}.}^{}`$ (as described in Ref. LSND-nueC ). This procedure yields 722 beam-on events with a beam-unrelated background of 17.5 events. The sidereal time distribution of these beam-on events are shown in Fig. 5. The P-$`\chi ^2`$ for this sidereal time distributions is 29.3 for 37 time bins which corresponds to $`P(\chi ^2)=0.78`$. The Kolmogorov-Smirnov test yields $`D_n=0.020`$ and $`P(\mathrm{KS})=0.94`$. The values indicate that the assumption of constant beam delivery, averaged over the sidereal day and over all LSND runs, is consistent with the data.
## VI The Extraction of the SME Parameters
While the LSND oscillation data examined in the previous section shows no statistically significant sidereal time variation, it is interesting to examine the data in context of the SME model explained in Section IV. First, this model does not require a sidereal time variation and, second, the LSND data set does allow for some sidereal time variation.
A maximum-likelihood method, with Eq. 7 as a description of the oscillation signal, was performed to extract allowed values of the SME parameters, $`(๐)_{\overline{e}\overline{\mu }}`$, $`(๐_s)_{\overline{e}\overline{\mu }}`$, $`(๐_c)_{\overline{e}\overline{\mu }}`$, $`(_s)_{\overline{e}\overline{\mu }}`$, and $`(_c)_{\overline{e}\overline{\mu }}`$. In general, these parameters are complex โ the special case is considered here where these parameters are real. Also, the values extracted are an effective average over the energy range of the LSND data set, $`20<E_{e^+}<60\mathrm{MeV}`$.
The parameters were extracted using an unbinned likelihood function,
$`\mathrm{\Lambda }`$ $`=`$ $`{\displaystyle \frac{e^\mu }{N!}}{\displaystyle \underset{i=1}{\overset{N}{}}}(\mu _s_s+\mu _b_b)`$ (15)
$`\times {\displaystyle \underset{i=\mathrm{s},\mathrm{b}}{}}{\displaystyle \frac{1}{\sqrt{2\pi \sigma _{i}^{}{}_{}{}^{2}}}}\mathrm{exp}({\displaystyle \frac{(\mu _i\overline{\mu _i})^2}{2\sigma _{i}^{}{}_{}{}^{2}}})`$
where $`N`$ is the total number of events in the sample, $`\mu _s`$ is the total predicted oscillation signal events, $`\mu _b`$ is the estimated number of background events, and $`\mu =\mu _s+\mu _b`$. The shape of the data in sidereal time is described with the functions $`_s`$ and $`_b`$. $`_s`$ depends on the SME parameters as in Eq. 7 and $`_b`$ is assumed to be constant in sidereal time. The latter half of the likelihood function describes systematic errors on the signal and background events. In implementation, the natural log of the likelihood function, $`\mathrm{}`$ $`(=\mathrm{ln}\mathrm{\Lambda })`$, was used. Note that this function describes both the shape and the overall number of events.
Three different parameter combinations were considered.
* 1-parameter:
$`(๐)_{\overline{e}\overline{\mu }}0;(๐_s)_{\overline{e}\overline{\mu }},(๐_c)_{\overline{e}\overline{\mu }},(_s)_{\overline{e}\overline{\mu }},(_c)_{\overline{e}\overline{\mu }}=0`$
The โrotationally invariantโ case Coleman ; Pakvasa ; Bahcall ; Pena-Garay .
* 3-parameter:
$`(๐)_{\overline{e}\overline{\mu }},(๐_s)_{\overline{e}\overline{\mu }},(๐_c)_{\overline{e}\overline{\mu }}0;(_s)_{\overline{e}\overline{\mu }},(_c)_{\overline{e}\overline{\mu }}=0`$
Includes all of the CPT-odd terms of the minimal-SME model.
* 5-parameter:
$`(๐)_{\overline{e}\overline{\mu }},(๐_s)_{\overline{e}\overline{\mu }},(๐_c)_{\overline{e}\overline{\mu }},(_s)_{\overline{e}\overline{\mu }},(_c)_{\overline{e}\overline{\mu }}0`$
Full minimal-SME model including both CPT-odd and CPT-even terms.
Using each of these three parameter sets, the log likelihood, $`\mathrm{}`$, was calculated for the 186 candidate oscillation events as each of the parameters in the set was varied in a range around zero. The sidereal time for the parameter values that maximized $`\mathrm{}`$ is plotted together with the data in Figure 6. Note that the data in Fig. 6 is grouped into 24 time bins instead of 37 as was used in Fig. 4. This is to allow for the quality of the fit to be more easily seen.
The maximum-$`\mathrm{}`$ solutions are summarized below. The likelihood contours for the 3-parameter combination are shown in Fig. 7. The (1$`\sigma `$) errors were calculated by determining the parameter ranges where $`\mathrm{}`$ decreased by 0.5 (1-parameter), 1.77 (3-parameters), or 3.0 (5-parameters) from the maximum value.
* 1-parameter:
$`(๐)_{\overline{e}\overline{\mu }}`$ $`=`$ $`3.3\pm 0.4\pm 0.2`$ (16)
* 3-parameter:
There are two solutions within the $`1\sigma `$ likelihood region (see Fig. 7).
Solution 1 (maximum-$`\mathrm{}`$):
$`(๐)_{\overline{e}\overline{\mu }}`$ $`=`$ $`0.2\pm 1.0\pm 0.3,`$
$`(๐_s)_{\overline{e}\overline{\mu }}`$ $`=`$ $`4.0\pm 1.3\pm 0.4,`$
$`(๐_c)_{\overline{e}\overline{\mu }}`$ $`=`$ $`1.9\pm 1.8\pm 0.4.`$ (17)
Solution 2:
$`(๐)_{\overline{e}\overline{\mu }}`$ $`=`$ $`3.3\pm 0.5\pm 0.3,`$
$`(๐_s)_{\overline{e}\overline{\mu }}`$ $`=`$ $`0.1\pm 0.6\pm 0.2,`$
$`(๐_c)_{\overline{e}\overline{\mu }}`$ $`=`$ $`0.5\pm 0.6\pm 0.2.`$ (18)
* 5-parameter: Multiple (connected) solutions exist in the 5-parameter case making a numerical extraction of errors impossible. The maximum-$`\mathrm{}`$ solution is:
$`(๐)_{\overline{e}\overline{\mu }}`$ $`=`$ $`0.7,`$
$`(๐_s)_{\overline{e}\overline{\mu }}`$ $`=`$ $`3.7,`$
$`(๐_c)_{\overline{e}\overline{\mu }}`$ $`=`$ $`2.3,`$
$`(_s)_{\overline{e}\overline{\mu }}`$ $`=`$ $`0.9,`$
$`(_c)_{\overline{e}\overline{\mu }}`$ $`=`$ $`0.6.`$ (19)
All parameters have units of $`10^{19}`$ GeV and the errors quoted above are in the form $`\pm `$(statistical)$`\pm `$(systematic).
In all of these results, duplicate solutions exist with opposite signs for all of the parameters. Note that in the 3-parameter case, the two solutions correspond to a large value for $`(๐)_{\overline{e}\overline{\mu }}`$ with a small value for $`(๐_s)_{\overline{e}\overline{\mu }},(๐_c)_{\overline{e}\overline{\mu }}`$ and vice versa. The small-$`(๐)_{\overline{e}\overline{\mu }}`$, large-$`(๐_s)_{\overline{e}\overline{\mu }},(๐_c)_{\overline{e}\overline{\mu }}`$ solution is only slightly favored over the large-$`(๐)_{\overline{e}\overline{\mu }}`$, small-$`(๐_s)_{\overline{e}\overline{\mu }},(๐_c)_{\overline{e}\overline{\mu }}`$ solution. This is because the sinusoidal terms in Eq. 7 improve the description of the data in sidereal time, although, an oscillation probability that is constant in sidereal time is consistent with the data (as was reported in Sec. V). Note also that the solution where $`(๐)_{\overline{e}\overline{\mu }}`$ is the only nonzero term can be identified with the conventional neutrino oscillation description via the first term of Eq. 9. A solution with all parameters $`0`$ is highly disfavored. This is equivalent to the statement that the LSND oscillation excess is statistically significant.
Since the oscillation probability depends on the SME parameters squared, the results for the SME parameters obtained above are more easily compared to the measured oscillation probability from LSND via combinations of the squares of the parameters. The value resulting from the 1-parameter solution is
$$|(๐)_{\overline{e}\overline{\mu }}|^2=10.7\pm 2.6\pm 1.3(10^{19}\text{GeV})^2.$$
(20)
The values for the parameter square sum resulting from the multiparameter combinations are more highly constrained than for individual parameters. The value extracted from the 3-parameter solution is
$`|(๐)_{\overline{e}\overline{\mu }}|^2+\frac{1}{2}|(๐_s)_{\overline{e}\overline{\mu }}|^2+\frac{1}{2}|(๐_c)_{\overline{e}\overline{\mu }}|^2`$
$`=9.9\pm 2.3\pm 1.4(10^{19}\text{GeV})^2,`$ (21)
and from the 5-parameter solution,
$`|(๐)_{\overline{e}\overline{\mu }}|^2+\frac{1}{2}|(๐_s)_{\overline{e}\overline{\mu }}|^2+\frac{1}{2}|(๐_c)_{\overline{e}\overline{\mu }}|^2`$
$`+\frac{1}{2}|(_s)_{\overline{e}\overline{\mu }}|^2+\frac{1}{2}|(_c)_{\overline{e}\overline{\mu }}|^2`$
$`=10.5\pm 2.4\pm 1.4(10^{19}\text{GeV})^2.`$ (22)
These results for the combination of SME parameters are consistent with the previously reported oscillation probability from LSND LSND2 and with the estimate presented in Ref. SBL .
## VII A High-Energy Subset of the LSND Data
A high-energy subset of the LSND data with a positron energy cut, $`36<E_{e^+}<60\mathrm{MeV}`$, is interesting to examine separately. The $`\nu `$-induced background is reduced in this sample LSND2 . Furthermore, the $`1/E`$ prefactor to the mass term in Eq. 9 would suppress the conventional oscillation terms relative to any Lorentz-violation terms present.
This reduced data set consists of 73 beam-on events with expected beam-unrelated and $`\nu `$-induced background events of $`31.3\pm 0.8`$ and $`10.0\pm 0.8`$, respectively. The sidereal and GM time distributions of the 571 beam-off events passing these high-energy cuts are shown in Figure 8. The P-$`\chi ^2`$ is 12.2 for 14 sidereal time bins, corresponding to $`P(\chi ^2)=0.51`$. The resulting $`p`$-value from the KS test to this distribution is $`P(\mathrm{KS})=0.080`$. Again, these values show no reason to reject the null hypothesis for beam-off data.
The sidereal time distribution of the high-energy beam-on data is shown in Figure 9. The P-$`\chi ^2`$ is 20.4 for 14 sidereal time bins corresponding to $`P(\chi ^2)=0.09`$. The resulting $`p`$-value for the KS test is $`P(\mathrm{KS})=0.178`$. Although these values indicate a slightly reduced agreement with the null hypothesis, they do not indicate a statistically significant sidereal variation. The complete results from the statistical tests on the sidereal and GM time distributions for the high-energy data are summarized in Table 2.
The maximum-likelihood procedure was applied to this high-energy data set using the 1- and 3-parameter combinations described in Section VI. The limited data sample did not allow for the 5-parameter combination. Figure 10 (with a reduced bin size) shows 1- and 3-parameter maximum-$`\mathrm{}`$ solutions superimposed on the high-energy data. Both parameter combinations produce acceptable descriptions of the data.
The values for the parameter square sums extracted with the maximum-likelihood method are summarized below. The likelihood contours for the 3-parameter combination are shown in Fig. 11.
* 1-parameter:
$`|(๐)_{\overline{e}\overline{\mu }}|^2`$ $`=`$ $`10.7\pm 2.9\pm 1.5(10^{19}\text{GeV})^2`$ (23)
* 3-parameter:
$`|(๐)_{\overline{e}\overline{\mu }}|^2+\frac{1}{2}|(๐_s)_{\overline{e}\overline{\mu }}|^2+\frac{1}{2}|(๐_c)_{\overline{e}\overline{\mu }}|^2`$
$`=10.2\pm 2.7\pm 1.3(10^{19}\text{GeV})^2`$ (24)
As can be seen by comparing the results from the high-energy subset with the entire data set, there are no significant differences. The time distributions from the high-energy subset are consistent with no sidereal variation and the results from the SME-parameter extraction are consistent with those obtained from the entire data set.
## VIII A Global Solution of Neutrino Oscillations?
To determine the implications of the allowed SME-parameter values extracted from the LSND data, consider the situation where the only nonzero term is $`(๐_s)_{\overline{e}\overline{\mu }}`$. This term was the largest in the maximum-$`\mathrm{}`$ solutions that allowed for sidereal variation. In this case, one or more of the SME coefficients, $`(a_L)_{\overline{e}\overline{\mu }}^X`$, $`(a_L)_{\overline{e}\overline{\mu }}^Y`$, $`(c_L)_{\overline{e}\overline{\mu }}^{TX}`$, $`(c_L)_{\overline{e}\overline{\mu }}^{TY}`$, $`(c_L)_{\overline{e}\overline{\mu }}^{XZ}`$ and $`(c_L)_{\overline{e}\overline{\mu }}^{YZ}`$ would be nonzero (as can be seen from Eq. 10).
A simple interpretation is that one of the $`a_L`$-type SME coefficients is of order $`10^{19}`$ GeV or one of the $`c_L`$-type is of order $`10^{17}`$ (or $`E\times c_L10^{19}`$ GeV, where $`E`$ is the neutrino energy). These values would have significant implications in other neutrino oscillation experiments and produce effects that have not been observed. In the simplest class of models, the acceptable maximum scale of Lorentz and CPT violation for reactor neutrino oscillations is $`a_L10^{21}`$ GeV, $`c_L10^{22}`$, and, for long-baseline neutrino oscillations, $`a_L10^{22}`$ GeV, $`c_L10^{19}`$ oscillation ; Coleman ; Pakvasa ; Bahcall ; Pena-Garay . However, there is no theoretical motivation that nature has chosen a simple solution for neutrino oscillations bicycle . A global solution of neutrino oscillations with Lorentz and CPT violation that accommodates all the data may yet be obtainable within this SME framework.
For this reason, it is important to search for sidereal variations in other short-baseline neutrino oscillation experiments. The data can be analyzed with the same method as presented here.
The currently running MiniBooNE experiment MiniBooNE , with a different beam energy and with a $`\nu _\mu `$ beam, will be able to test these LSND solutions for Lorentz and CPT violating neutrino oscillations with a high-statistics appearance measurement. In particular, a measurement from MiniBooNE would provide an additional five constraints (Eqs. 9-13) on the SME coefficients. The neutrino propagation vectors are different for MiniBooNE as the neutrino beamline is oriented toward compass north (as opposed to east for LSND). Also, the SME coefficients would be transformed for the neutrino case SBL . If MiniBooNE collects a significant set of data with a $`\overline{\nu }_\mu `$ beam, an additional set of constraints with antineutrino coefficients would also be obtained.
Of course, results from other neutrino oscillation experiments would add further valuable information. This has been investigated for Super-Kamiokande MMCPT04 and MINOS BRCPT04 .
## IX Conclusions
The neutrino oscillation candidate events from the LSND experiment have been examined for evidence of sidereal time variation โ a possible signal for Lorentz violation in the neutrino sector. The oscillation excess is consistent with no sidereal time variation. An examination of a high-energy subset of the data yields the same conclusion.
A โsmoking-gunโ for Lorentz violation has not been found in the LSND signal. However, the data are adequately described within the SME neutrino oscillation formalism that includes both Lorentz and CPT violation oscillation ; SBL . A maximum-likelihood method was used to determine allowed parameter regions for SME parameter combinations. They indicate values on the order of $`10^{19}`$ GeV for $`a_L`$ and $`E\times c_L`$. These values are in the range expected for Planck-scale effects in the neutrino sector. Future results from high-statistics oscillation experiments will allow more stringent tests of the SME framework.
## X Acknowledgments
This work was conducted under the auspices of the US Department of Energy, supported in part by funds provided by the University of California for the conduct of discretionary research by Los Alamos National Laboratory. This work was also supported by the National Science Foundation.
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# Continuum-discretized coupled-channels method for four-body breakup reactions 1footnote 11footnote 1 A talk given at the Workshop on Reaction Mechanisms for Rare Isotope Beams, Michigan State University, March 9-12, 2005 (to appear in an AIP Conference Proceedings).
## I INTRODUCTION
In the study of reactions induced by unstable nuclei, analysis of the case where the projectile is considered to be composed of three-clusters such as <sup>6</sup>He and <sup>11</sup>Li becomes quite important. For this purpose, along the diagram in Fig. 1, we have developed the three-body CDCC (Continuum-Discretized Coupled-Channels) for nuclear breakup of two-body projectiles cdcc into the four-body CDCC for Coulomb and nuclear breakup of three-body projectiles.
The momentum-bin method to discretize the continuum states of the two-body projectiles (such as $`{}_{}{}^{6}\mathrm{Li}=\alpha +d,^8\mathrm{B}+p`$, etc.) is not practically available to the case of three-body projectiles. On the basis of the Gaussian expansion method (GEM)GEM , we proposed, in Ref.ps-cdcc , the pseudo-state (PS) method to discretize the continuum states and examined it in the case of two-body projectiles (three-body CDCC); this is Step A in Fig.1. In the PS method we diagonalized the two-body Hamiltonian of the internal motion of the projectile using the Gaussian basis functions GEM and obtained dense distribution of the pseudo-states, namely discretized continuum states. An advantage of this method is that it can easily be extended to the case of three-body projectiles by using the GEM. Another advantage of the PS method is that we can derive continuous $`S`$-matrix elements as a smooth function of the momentum of the projectile breakup states. We found ps-cdcc that the $`S`$-matrix elements obtained by the PS method agrees well with the $`S`$-matrix elements by the momentum-bin method with very precise bins.
As Step B in Fig.1, we extended the three-body CDCC (for nuclear breakup) to the four-body CDCC (for nuclear breakup) using the three-body Gaussian basis functions of GEM to obtain bound and pseudo-states of the three-body projectiles fb-cdcc . The GEM is very suitable for describing bound and pseudo-states of three- and four-body systems; it is extensively reviewed in Ref.GEM . The four-body CDCC was applied to the <sup>6</sup>He+<sup>12</sup>C at 18 and 229.8 MeV. The differential cross sections of the elastic scattering were well reproduced by using the double-folding CC potentials.
In Step C, we improved the three-body CDCC for nuclear breakup to that for Coulomb and nuclear breakup psc-cdcc by using the PS method with the complex-range Gaussian basis functions GEM instead of the (real-range) Gaussian basis functions adopted in the previous Steps. Due to the long-ranged Coulomb coupling-potentials, the modelspace required for CDCC is very large. Particularly, one must prepare the internal wave functions of the projectile, both in bound and continuum states, for a wide range of the internal coordinate, say 0โ100 fm, which is in general difficult for PS methods. This can easily be achieved by using the complex-range Gaussian basis in the case of two-body projectile.
In order to treat both Coulomb and nuclear breakup processes at intermediate energies with high accuracy and computational speed, a new method was proposed in Ref. ecdcc ; namely, a hybrid calculation with the three-body CDCC method and the eikonal-CDCC (E-CDCC) method. E-CDCC describes the center-of-mass motion of the projectile relative to the target by straight-line approximation (or by using Coulomb wave functions instead of plane waves) and treats the excitation of the projectile explicitly by CDCC with the momentum-bin method or the PS method. E-CDCC drastically reduces computation time and eliminates many problems concerned with huge angular momentum in solving coupled-channel equations. Thus, the hybrid calculation is expected to be opening the door to the systematic analysis of Coulomb (plus nuclear) dissociation of projectiles in the wide range of beam energies.
Finally, by Step Bโ (or by Step Cโ) we can reach the four-body CDCC for Coulomb and nuclear breakup. This step was not reported in the time of the RIA workshop but was recently accomplished and successfully applied to the <sup>6</sup>He+<sup>209</sup>Bi scattering at 19.0 and 22.5 MeV fbc-cdcc .
## II METHOD OF PSEUDO-STATE CDCC FOR TWO-BODY PROJECTILES
In the method of CDCC, the total wave function of the scattering state $`\mathrm{\Psi }_{JM}`$ is expanded in terms of a finite number of internal wave functions $`\mathrm{\Psi }_{nIm}(\xi )`$ of the projectile:
$$\mathrm{\Psi }^{JM}(\xi ,๐)=\underset{nI,L}{}[\mathrm{\Phi }_{nI}(\xi )\chi _{nI,L}^J(๐)]_{JM},$$
(1)
where $`๐`$ is the coordinate of the center-of-mass of the projectile relative to the target, and $`\xi `$ is the internal coordinates of the projectile. $`I`$ is the total spin of the projectile and $`n`$ stands for the $`n`$th eigenstate. $`\chi _{nI,L}^J`$ represents the relative motion between the projectile and the target; $`L`$ is the orbital angular momentum regarding $`๐`$. The unknown function $`\chi _{nI,L}^J(๐)`$ are solved using the usual framework of the coupled-channel method for discrete excited states.
The projectile internal wave functions $`\mathrm{\Phi }_{nI}(\xi )`$ include both bound states and discretized continuum states. To calculate the wave functions of the latter states the momentum-bin method has widely been utilized in the usual three-body CDCC calculations. In the method the exact scattering wave functions are averaged within each narrow intervals of momentum between the two constituents in the projectile. But, this method is not practically suitable for discretizing the breakup states of the three-body projectile.
In the pseudo-state (PS) method cdcc ; PS1 ; PS2 , on the other hand, wave functions of the discretized breakup states are obtained by diagonalizing the internal Hamiltonian of the projectile, which describes the relative motion of the two constituents, using $`L^2`$-type basis functions. Since the wave functions of such pseudo breakup states have wrong asymptotic forms, the PS method was mainly used in the past to describe virtual breakup processes in the intermediate stage of elastic scattering PS2 and ($`d,p`$) reactions cdcc .
In the work of Ref.ps-cdcc , however, we proposed the new method of pseudo-state (PS) discretization for two-body projectiles. It can be used not only for virtual breakup processes in elastic scattering but also for breakup reactions. In order to diagonalize the Hamiltonian of the two-body projectile, we employed two types of basis functions. One is the conventional real-range Gaussian functions
$$\varphi _j\mathrm{}(r)=r^{\mathrm{}}\mathrm{exp}\left[(r/a_j)^2\right],(j=1\text{}n)$$
(2)
where $`\{a_j\}`$ are assumed to increase in a geometric progression Kamimura88 ; GEM :
$$a_j=a_1(a_n/a_1)^{(j1)/(n1)}.$$
(3)
The other is an extension of (2) introduced in Ref. GEM , i.e., the following pairs of functions:
$`\varphi _j\mathrm{}^\mathrm{C}(r)`$ $`=`$ $`r^{\mathrm{}}\mathrm{exp}\left[(r/a_j)^2\right]\mathrm{cos}\left[b(r/a_j)^2\right],`$
$`\varphi _j\mathrm{}^\mathrm{S}(r)`$ $`=`$ $`r^{\mathrm{}}\mathrm{exp}\left[(r/a_j)^2\right]\mathrm{sin}\left[b(r/a_j)^2\right],(j=1\text{}n).`$ (4)
Here, $`b`$ is a free parameter, in principle, but numerical test showed that $`b=\pi /2`$ is recommendable. Both $`\varphi _j\mathrm{}^\mathrm{C}`$ and $`\varphi _j\mathrm{}^\mathrm{S}`$ are to be used simultaneously; the total number of basis is thus $`2n`$. The basis functions (4) can also be expressed as
$`\varphi _j\mathrm{}^\mathrm{C}(r)`$ $`=`$ $`\{\psi _j\mathrm{}^{}(r)+\psi _j\mathrm{}(r)\}/2,`$
$`\varphi _j\mathrm{}^\mathrm{S}(r)`$ $`=`$ $`\{\psi _j\mathrm{}^{}(r)\psi _j\mathrm{}(r)\}/(2i),`$ (5)
with
$$\psi _j\mathrm{}(r)=r^{\mathrm{}}\mathrm{exp}[\eta _jr^2],\eta _j=(1+ib)/a_j^2,$$
(6)
i.e., Gaussian functions with a complex-range parameter. We thus refer to the basis $`\varphi _j\mathrm{}^\mathrm{C}`$ and $`\varphi _j\mathrm{}^\mathrm{S}`$ as the complex-range Gaussian basis.
The complex-range Gaussian basis functions are oscillating with $`r`$. They are therefore expected to simulate the oscillating pattern of the continuous breakup state wave functions better than the real-range Gaussian basis functions do. Moreover, numerical calculation with the complex-range Gaussians can be done using essentially the same computer programs as for the real-range Gaussians, just replacing real variables for $`a_j`$ of Eq. (3) by complex ones. Usefulness of the real- and complex-range Gaussian basis functions in few-body calculations are extensively presented in the review work GEM .
Here, we explore a typical example in which the complex-range Gaussian basis functions reproduce highly oscillatory functions with high accuracy. A good test is to calculate the wave functions of highly excited states in a harmonic oscillator potential; note that this potential is not specially advantageous for the Gaussian bases. We take the case of a nucleon with angular momentum $`l=0`$ in a potential having $`\mathrm{}\omega =15.0`$ MeV. Parameters of the complex-range Gaussian basis functions are $`\{\mathrm{\hspace{0.17em}2}n=28,a_1=1.4\mathrm{fm},a_n=5.8\mathrm{fm},b=\frac{\pi }{2}\frac{1}{1.2^2}=1.09\}`$. For the sake of comparison, we also tested the Gaussian basis functions with the parameters $`\{n=28,a_1=0.5\mathrm{fm},a_n=11.3`$fm }. Optimized $`a_1`$ and $`a_n`$ are quite different between the two types of bases though the total numbers of basis functions are the same. In Table 1, we compare the calculated energy eigenvalues with the exact ones. It is evident that the complex-range Gaussians can reproduce the energy up to much more highly excited states than the Gausssians do. For the Gaussian basis, even if the number of basis functions is increased, the result is not significantly improved, because the number of oscillation does not increase. On the other hands, for the complex-range Gaussian functions, as the number is increased, the result becomes better so long as the number of oscillation is not too larger than $``$ 20. Figure 2 demonstrates good accuracy of the wave function of the 19-th excited state having 38 quanta. Error is within a few %, much smaller than the thickness of the line. The figure suggests that the basis functions is also suitable for describing pseudo-states used for Coulomb breakup reactions.
We here emphasize that even in the case where the projectile is assumed to be three-body system, the Gaussian basis functions with real and complex ranges are easily utilized in the CDCC calculation with the PS method. We discuss this point in the next section.
Another advantage of the PS method, in the case of two-body projectiles, is that the discrete breakup $`S`$-matrix elements, say $`S_{nIL,0I_0L_0}`$, for the transition from $`\mathrm{\Phi }_{0I_0}(\xi )`$ to $`\mathrm{\Phi }_{nI}(\xi )`$ can be accurately transformed to smooth $`S`$-matrix elements, say $`\stackrel{~}{S}_{IL,I_0L_0}(k)`$, as following ps-cdcc , since the two-body PS basis functions can form in the good approximation a complete set in the finite region which is important for the breakup processes:
$$\stackrel{~}{S}_{IL,0I_0L_0}(k)=\underset{n}{}\stackrel{~}{\mathrm{\Phi }}_I(k,\xi )|\mathrm{\Phi }_{nI}(\xi )_\xi S_{nIL,0I_0L_0},$$
(7)
where $`\stackrel{~}{\mathrm{\Phi }}_I(k,\xi )`$ is the exact wave function of the internal motion of the two-body projectile.
Example 1 $`:^6`$Li+<sup>40</sup>C scattering at 156 MeV.
Here, we briefly show results of test calculations done in ps-cdcc <sup>6</sup>Li+<sup>40</sup>C scattering at 156 MeV. The $`\alpha d`$ continuum of the <sup>6</sup>Li projectile is discretized as in Fig. 4 using the real-range Gaussian bases and the complex-range Gaussian bases. The modelspace sufficient for describing breakup processes in this scattering is $`k_{\mathrm{max}}=2.0`$ fm<sup>-1</sup> and $`\mathrm{}_{\mathrm{max}}=2`$; the modelspace is composed of two $`k`$-continua for $`s`$-state and $`d`$-state. There exists a $`d`$-state resonance. The resonance is automatically taken care by the PS method by the lowest-lying several pseudo-states. On the other hand, in the momentum-bin method, the d-state $`k`$-continuum is further divided in the momentum-bin method into the resonant part $`[0<k<0.55\mathrm{fm}^1]`$ and the non-resonant part $`[0.55<k<2.0\mathrm{fm}^1]`$. In the former region the $`k`$ continuum $`d`$-state wave function varies rapidly with $`k`$. The momentum-bin method can simulate this rapid change with bins of an extremely small width. In fact clear convergence is found for both the elastic and the breakup $`S`$-matrix elements, when the resonance part is described by 30 bins and the non-resonance part of the $`d`$-state and the $`s`$-state $`k`$-continua by 20 bins.
Figure 4 represents breakup $`S`$-matrix elements at grazing total angular momentum $`J=43`$; (a) $`s`$-state breakup and (b) $`d`$-state breakup in the case of $`L=J2`$. The real- and complex-range Gaussian PS discretization well reproduce the โexactโ solution calculated by the momentum-bin method with dense bins. The results of the two PS methods turn out to coincide within the thickness of the line. The resonance peak can be expressed by only 8 (12) breakup channels in the complex-range (real-range) Gaussian PS method, while the corresponding number of breakup channels is 30 in the momentum-bin method, as mentioned above. Thus, one can conclude that the real- and complex-range Gaussian PS methods are very useful for describing not only non-resonant states but also resonant ones.
The PS method has at least two advantages over the widely used momentum bin average method. One is that it does not need the exact wave function of the projectile over the entire region of $`r`$. This is important from a theoretical point of view. The other is that with the real- and complex-range Gaussian bases one can calculate all the coupling potentials semi-analytically GEM , which is very useful in actual calculations; note that the Gaussian bases are very suitable for transforming wave functions and interactions from a Jacobian coordinate system to other ones. Furthermore, if the projectile has resonances in its excitation spectrum, the PS method discretizes the complicated spectrum with a reasonable number of the pseudo-states, without distinguishing the resonance states from non-resonant continuous states. These advantages of the PS method are extremely helpful, sometimes even essential, in applying CDCC to four-body breakup effects of unstable nuclei such as <sup>6</sup>He and <sup>11</sup>Li.
Example 2 $`:^8`$B+<sup>58</sup>Ni scattering at 25.8 MeV.
Here, we briefly show results of test calculation in psc-cdcc for Coulomb breakup process of <sup>8</sup>B+<sup>58</sup>Ni scattering at 25.8 MeV. The <sup>7</sup>Be$`p`$ continuum in the <sup>8</sup>B projectile is discretized as in Fig. 5 by the PS method with the real-range Gaussian bases and the complex-range Gaussian bases. In the PS method, the number of channels included in the CDCC calculation, was 18 for both the $`s`$\- and $`p`$-states at $`k<k_{\mathrm{max}}=0.66`$ fm<sup>-1</sup>, which give a satisfactory convergence of the result. The resulting wave functions with positive eigenenergies turned out to oscillate up to about 100 fm. In the momentum-bin method, the modelspace with $`k_{\mathrm{max}}=0.66`$ fm<sup>-1</sup> and $`\mathrm{\Delta }k=0.66/16`$ (0.66/32) fm<sup>-1</sup> for $`p`$-state ($`s`$-state) gives convergence of the resulting total breakup cross section. The maximum internal coordinate $`r_{\mathrm{max}}`$ was taken to be 100 fm.
Figure 6 shows the result of the comparison of $`|S_{\mathrm{}}(k)|^2`$ at $`J=150`$, which corresponds to the scattering angle of 10 assuming the classical path. It was found that CDCC calculation with only Coulomb coupling potentials gives a peak at 10 in the total breakup cross section. Thus, it can be assumed that Fig. 6 corresponds to the most-Coulomb-like breakup process; in any case, the feature of the result was found to be almost independent of $`J`$. In each panel of Fig. 6, one sees that the result of PS-CDCC (solid line) very well reproduces the โexactโ solution (step line by the momentum-bin method) for all $`k`$ being significant for the <sup>8</sup>B Coulomb breakup.
## III GAUSSIAN EXPANSION METHOD FOR FEW-BODY SYSTEMS
In this section we briefly explain the Gaussian expansion method (GEM) for few-body systems. The method was proposed by Kamimura in 1988 Kamimura88 for three-body systems and was much developed by Hiyama using the infinitesimally-shifted Gaussian basis functions even for four-body systems (reviewed in GEM ).
A good example to show the accuracy and usefulness of the method is the determination of upper limit of the difference between the masses of proton and antiproton, $`m_p`$ and $`m_{\overline{p}}`$, respectively. The first recommended upper limit of $`|m_{\overline{p}}m_p|/m_p`$ by the Particle Data Group listed in Particle Listings 2000 Listing2000 was $`5\times 10^7`$, which could be used for a test of $`CPT`$ invariance. This number was extracted from a high-resolution laser experiment involving metastable states of antiprotonic helium atom (He$`{}_{}{}^{2+}+e^{}+\overline{p}`$) Torii99 by Kino et al. Kino99 through a theoretical analysis of the highly excited states of the Coulomb three-body system using GEM. The ratio was improved to $`|m_{\overline{p}}m_p|/m_p<1\times 10^8`$, as listed in the Particle Listings 2004, by later, more extensive experiments and additional calculations (cf. Ref.GEM )
In the Gaussian expansion method GEM , wave functions of the projectile, $`\mathrm{\Phi }_{nIm}`$ in (1), is written as a sum of component functions in the Jacobian coordinates for rearrangement channels $`c=13`$ in Fig. 7 as
$$\mathrm{\Phi }_{nIm}(\xi )=\underset{c=1}{\overset{3}{}}\psi _{nIm}^{(c)}(\xi ),$$
(8)
Each $`\psi _{nIm}^{(c)}`$ is expanded in terms of the Gaussian basis functions:
$`\psi _{nIm}^{(c)}(\xi )`$ $`=`$ $`\phi ^{(\alpha )}{\displaystyle \underset{\lambda \mathrm{}\mathrm{\Lambda }S}{}}{\displaystyle \underset{i=1}{\overset{i_{\mathrm{max}}}{}}}{\displaystyle \underset{j=1}{\overset{j_{\mathrm{max}}}{}}}A_{i\lambda j\mathrm{}\mathrm{\Lambda }S}^{(c)nI}y_c^\lambda r_c^{\mathrm{}}e^{(y_c/\overline{y}_i)^2}e^{(r_c/\overline{r}_j)^2}`$ (9)
$`\times \left[\left[Y_\lambda (\widehat{๐ฒ}_c)Y_{\mathrm{}}(\widehat{๐ซ}_c)\right]_\mathrm{\Lambda }\left[\eta _{\frac{1}{2}}^{(n_1)}\eta _{\frac{1}{2}}^{(n_2)}\right]_S\right]_{Im},`$
where $`\lambda `$ ($`\mathrm{}`$) is the angular momentum regarding the Jacobian coordinates $`๐ฒ_c`$ ($`๐ซ_c`$), and $`\eta _{1/2}`$ is the spin wave function of each valence neutron ($`n_1`$ or $`n_2`$). <sup>4</sup>He has been treated as an inert core with the $`(0s)^4`$ internal configuration, $`\phi ^{(\alpha )}`$. The Gaussian range parameters are taken to lie in geometric progression:
$`\overline{y}_i=\overline{y}_1(\overline{y}_{\mathrm{max}}/\overline{y}_1)^{(i1)/(i_{\mathrm{max}}1)},`$ (10)
$`\overline{r}_j=\overline{r}_1(\overline{r}_{\mathrm{max}}/\overline{r}_1)^{(j1)/(j_{\mathrm{max}}1)}.`$ (11)
$`\mathrm{\Phi }_{nIm}`$ is antisymmetrized for the exchange between $`n_1`$ and $`n_2`$. Meanwhile, the exchange between each valence neutron and each nucleon in <sup>4</sup>He is treated approximately by the orthogonality condition. The eigenenergies $`ฯต_{nI}`$ of <sup>6</sup>He and the corresponding expansion-coefficients $`A_{i\lambda j\mathrm{}\mathrm{\Lambda }S}^{(c)nI}`$ are determined by diagonalizing the Hamiltonian of the interrenal motion of <sup>6</sup>He GEM6He1 ; GEM6He using a large number of three-body Gaussian basis functions. Detailed information on the basis is listed in Ref.fb-cdcc . The calculated $`ฯต_{nI}`$ are $`0.98`$ MeV for the $`0^+`$ ground state and 0.72 MeV for the $`2^+`$ resonance state; here, we took the Bonn A potential between the valence nucleons and increased the depth of the $`n\alpha `$ potential by a few percent so that the ground-state energy is reproduced.
In the four-body CDCC calculation of <sup>6</sup>He+<sup>12</sup>C shown in a later section, we take $`I^\pi =0^+`$ and $`2^+`$ states for <sup>6</sup>He. Here we omit the $`1^{}`$ state that does not contribute to the nuclear breakup processes (but they are included in the calculation of Coulomb and nuclear breakup in Ref.fbc-cdcc ). In order to demonstrate the convergence of the four-body CDCC solution with increasing the number of the Gaussian basis functions, we prepare three sets of the basis functions, i.e., sets I, II and III listed in Table II of fb-cdcc . Resultant energy levels of the ground and pseudo-states are shown in (a), (b) and (c) in Fig. 8, respectively. For <sup>6</sup>He+<sup>12</sup>C scattering at 18 MeV (229.8 MeV) which will be discussed in the next section, high-lying states with $`ฯต_{nI}>12`$ MeV ($`ฯต_{nI}>25`$ MeV) are found to give no effect on the elastic and breakup $`S`$-matrix elements. Thus, the effective number of the eigenstates of <sup>6</sup>He, is reduced much for each of cases (a), (b), (c) as shown in Fig. 8. The case (b) was found to be sufficient to obtain a good convergence. In the GEM, computation time to obtain the wave functions of the bound and pseudo states is very short; for example, all the wave functions of the states in Fig. 8(c) is obtained in 10 minutes on FUJITSU VPP5000, a supercomputer.
It is to be noted that the bound and pseudo-states obtained with the GEM calculations construct an approximate complete sets for each $`J(=0,1,2)`$ in a finite region which is responsible for the reaction; this was examined by checking that those states (below 100 MeV) satisfies 99.9 % of the energy-weighted cluster sum-rule limit for monopole, dipole and quadrupole transitions.
## IV Four-body CDCC analysis of <sup>6</sup>He+<sup>12</sup>C scattering at 18 and 229.8 MeV
In this section, we briefly introduce the results obtained in the work of Ref.fb-cdcc . We performed the four-body CDCC calculation for <sup>6</sup>He+<sup>12</sup>C scattering at 18 and 229.8 MeV using the wave functions of the bound state and the pseudo-states of <sup>6</sup>He obtained above.
The real part of the CC potentials, say $`V_{nIL,n^{}I^{}L^{}}^J(R)`$, was constructed by using the double-folding model dfm ; the potentials were calculated by folding the DDM3Y $`NN`$ interaction into the transition densities between the states $`\mathrm{\Phi }_{nI}(\xi )`$ and $`\mathrm{\Phi }_{n^{}I^{}}(\xi )`$ (cf. Ref.fb-cdcc for details) and the ground-state density of <sup>12</sup>Kamimura12C . The imaginary part was assumed, as usually done cdcc , to be given as (together with the real part)
$$(N_R+iN_I)V_{nIL,n^{}I^{}L^{}}^J(R),$$
(12)
where $`N_R=1.0`$ with no renormarization of the real part. The only parameter $`N_I`$ is searched for to reproduce the observed elastic cross section as well as possible. In the analysis of the <sup>6</sup>He+<sup>12</sup>C scattering, Coulomb breakup effect is ignored since it is negligible for this light target; the Coulomb potential is assumed to work between the center-of-mass of the target and that of the projectile.
Calculated and observed elastic cross sections for <sup>6</sup>He+<sup>12</sup>C scattering at 18 MeV are shown in Fig. 9. The optimum value of $`N_I`$ is 0.5, which is the same as that for <sup>6</sup>Li scattering at various incident energies cdcc . The dotted lines represent the elastic cross sections due to the single-channel calculation. Then, the difference between the solid and dotted lines shows the effect of the four-body breakup on the elastic cross section. The effect is sizable and indispensable to explain the behavior of the angular distribution. The case at 229.8 MeV is shown in Fig.10 and the optimum value of $`N_I`$ is 0.3. The breakup effect in this case is also important in reproducing the data. The origin of the small $`N_I`$ value for the <sup>6</sup>He scattering at 229.8 MeV is not clear at this moment, so more systematic experimental data are highly desirable for <sup>6</sup>He scattering.
We calculated the dynamical polarization (DP) potential induced by the four-body breakup processes, in order to understand effects of the processes on the elastic scattering. The DP potential is given by the deviation of the so-called wave-function-equivalent local potential derived using the elastic channel amplitude in the solution of the CDCC equation from the double-folding potential of the elastic channel. From the analysis fb-cdcc of the DP potential, one sees that inclusion of the four-body breakup processes makes the real part of the <sup>6</sup>Heโ<sup>12</sup>C potential shallower and the imaginary one deeper compared with the double-folding potential of the elastic channel. In particular, the latter effect is important and can be assumed to come from the Borromean structure of <sup>6</sup>He. This is consistent with the fact that the total reaction cross section is enhanced by the Borromean structurefb-cdcc .
## V conclusion and near-future problems
In conclusion, a fully quantum-mechanical method of treating four-body breakup is presented by extending CDCC. The validity of the method called four-body CDCC is confirmed by clear convergence of the calculated elastic and energy-integrated breakup cross sections with respect to extending the modelspace. The four-body CDCC is found to explain well the <sup>6</sup>He+<sup>12</sup>C scattering at 18 and 229.8 MeV in which <sup>6</sup>He easily breaks up into two neutrons and <sup>4</sup>He. For the elastic scattering, the four-body breakup processes make, in particular, the imaginary part of the <sup>6</sup>Heโ<sup>12</sup>C potential deeper, which is originated in the Borromean structure of <sup>6</sup>He.
In the analysis of fb-cdcc , four-body Coulomb breakup is neglected. However, it is possible to treat it within the four-body CDCC framework (cf. Fig. 1). Actually, after this RIA workshop, we reported in Ref.fbc-cdcc our four-body CDCC calculation of the <sup>6</sup>He+<sup>209</sup>Bi scattering at 19.0 and 22.5 MeV taking both the Coulomb and nuclear breakup effects into account. The elastic cross sections were well reproduced by the calculation. So, the same framework will be applicable to other cases of three-body projectiles with Coulomb and nuclear breakup.
In order to treat both Coulomb and nuclear breakup processes at intermediate energies, Ref. ecdcc proposed a new method, namely a hybrid calculation with the three-body CDCC method and the eikonal-CDCC (E-CDCC) method. This hybrid calculation is expected to be opening the door to the systematic analysis of Coulomb (plus nuclear) dissociation of projectiles in the wide range of beam energies. For example, the method was recently applied to the analysis of <sup>8</sup>B dissociation measurements to determine the astrophysical factor $`S_{17}(0)`$ accurately astro .
There are some important unstable nuclei that are considered to be composed of four-body constituents. For reactions in which such a four-body nucleus is a projectile, a five-body CDCC calculation is required. The GEM was already severely and successfully tested for the bound states and pseudo-states of four-body systems. A good example is seen in a calculation of four-nucleon system (<sup>4</sup>He) in Ref. Hiyama04 . The four-body GEM calculation with a realistic $`NN`$ force (AV8โ) and a phenomenological $`NNN`$ force (which is adjusted to reproduce the ground-state energy) reproduced the energy of the second $`0^+`$ state and the $`{}_{}{}^{4}\mathrm{He}(e,e^{})^4\mathrm{He}(0_2^+)`$ form factor. Furthermore, some 3000 $`\mathrm{\hspace{0.17em}0}^+`$ pseudo-states below 300-MeV excitation satisfied the energy-weighted monopole sum rule by 99.9% (with saturation) and made clear, for the first time, that the major part of the monopole sum rule limit, which had been long unknown, was distributed into low-lying four-body non-resonant continuum states. So, it may be said that it is ready to perform five-body CDCC calculations for reactions induced by four-body projectiles.
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# Stability of quasi-two-dimensional Bose-Einstein condensates with dominant dipole-dipole interactions
## Abstract
We consider quasi-two-dimensional atomic/molecular Bose-Einstein condensates with both contact and dipole-dipole interactions. It is shown that, as a consequence of the dimensional reduction, and within mean-field theory, the condensates do not develop unstable excitation spectra, even when the dipole-dipole interaction completely dominates the contact interaction.
Dilute Bose-Einstein condensates with long-range interactions offer promising opportunities to explore the potentially strong correlations induced by the interaction, which go beyond the comparatively weak correlations induced by the local contact interaction pseudopotential conventionally sufficient to describe most atomic condensates. The recent realization of a concrete physical system in which atomic magnetic dipoles play a significant role Griesmaier , is a first experimental step towards the exploration of dipolar condensate physics Goral ; YiYouII ; SantosZoller ; ODell ; Baranov . To investigate dipole-dominated physics, greater potential than by magnetic dipoles is offered by heteronuclear molecules with an electric dipole moment YiYou ; Kotochigova , due to the fact that the electric dipole interaction strength is larger than the magnetic one by a factor of about $`10^4`$ for typical atomic/molecular dipole moments of 1 Bohr magneton or 1 Debye, respectively.
An interesting property of three-dimensional (3D) condensates is the existence of a โrotonโ minimum, which stems from the fact that the Fourier transform of the dipole-dipole interaction, occurring in the Bogoliubov mean-field spectrum, assumes negative values in certain directions in momentum space RotonDipolar ; RotonsDipole ; EPJD . This property however also causes a problem of 3D dipolar condensates, as they become dynamically unstable for dipole-dipole interactions dominating the contact interaction, at densities corresponding to the Thomas-Fermi limit of large condensates. The roton minimum quickly deepens with increasing density and/or dipole coupling to hit the zero of energy, and a further increase is impossible because the system develops imaginary excitation energies.
In the present study an analytical proof is given that, in contrast to 3D condensates, quasi-two-dimensional (quasi-2D) dipolar condensates are stable, even when the dipole coupling completely dominates the contact coupling. The experimental realization of purely dipolar condensates is thus possible only in the quasi-2D regime. The physical reason for the stabilization is that dimensional reduction entails that the mutual dipole-dipole interaction of the atoms/molecules, with dipoles oriented perpendicular to the plane in the strongly confined direction, can effectively sample an exponentially smaller region in configuration space where the interaction assumes negative values. It is, furthermore, shown that a strongly correlated regime of purely dipolar condensates, where large quantum depletion of the condensate occurs and the mean-field description in terms of a single condensate wave function breaks down, takes place in a crossover regime from quasi-2D to 3D. For an electric-dipole-dominated dilute gas, the densities corresponding to this crossover regime to 3D turn out to be quite small, even within the strongly enlarged stability window of strongly anisotropic pancake-shaped condensates.
The mean-field description of the system we start with is based on a nonlocal 3D Gross-Pitaevskiว equation for the order parameter $`\mathrm{\Psi }`$ (we put $`\mathrm{}=m=1`$, where $`m`$ is the mass of the atoms/molecules) Goral ; YiYou :
$`i{\displaystyle \frac{}{t}}\mathrm{\Psi }(๐,t)`$ $`=`$ $`\{{\displaystyle \frac{^2}{2}}+{\displaystyle \frac{1}{2}}\omega _z^2z^2+g_{3\mathrm{D}}|\mathrm{\Psi }(๐,t)|^2`$
$`+{\displaystyle }d^3r^{}V_{dd}(๐๐^{})|\mathrm{\Psi }(๐^{},t)|^2\}\mathrm{\Psi }(๐,t).`$
The trapping in the plane is assumed to be negligible as compared with the strong confinement, of harmonic trapping frequency $`\omega _z`$, in the $`z`$ direction. The term in the second line contains the dipole-dipole interaction of the atoms/molecules, with all dipoles oriented by an applied field along the strongly confining $`z`$ direction:
$$V_{dd}(๐)=\frac{3g_d}{4\pi }\frac{13z^2/|๐|^2}{|๐|^3}.$$
(2)
The dipole coupling reads $`g_d=\mu _0d_m^2/3`$ for magnetic and $`g_d=d_e^2/3ฯต_0`$ for electric dipoles. In the units which we employ, the $`g_d`$ are having dimensions of length like the 3D contact interaction $`g_{3\mathrm{D}}=4\pi a_s`$, for which the length scale is set by the $`s`$-wave scattering length $`a_s`$. We note that the value of the contact interaction coupling strength generally depends on the value of $`g_d`$ (in an effective-dimension dependent manner), because the long-range dipole-dipole interaction affects the short-range scattering processes YiYouII . For the statements made in the present paper, the actual ratio $`g_{3\mathrm{D}}/g_d`$ and the absolute value of the (unscreened) $`g_d`$ will be relevant.
For a quasi-2D Bose-Einstein condensate, the motion of the atoms/molecules is by definition restricted to zero point oscillations in a harmonic oscillator potential. We then take as a general ansatz for the density
$$\rho (๐)=|\mathrm{\Psi }(๐)|^2=\frac{1}{\sqrt{\pi d_z^2}}\mathrm{exp}\left[\frac{z^2}{d_z^2}\right]n(x,y),$$
(3)
where the density in the plane, $`n(x,y)`$, is normalized to the total number of particles $`N`$. We treat $`d_z`$ as a parameter minimizing the Gross-Pitaevskiว ground state energy YiYou ; EPJD . Assuming homogeneous density in the plane, the equation determining $`d_z`$ reads $`\omega _z^2=d_z^4+(g_{3\mathrm{D}}+2g_d)nNd_z^3/\sqrt{2\pi }`$. If the right-hand side is dominated by the first (kinetic energy) term, $`d_z`$ equals the harmonic oscillator length $`1/\sqrt{\omega _z}`$, the system is quasi-2D, and the above Gaussian gives the density profile exactly. Defining a parameter $`\alpha \omega _zd_z^2`$, we have $`\alpha =1`$ if the system is quasi-2D and $`\alpha 1`$ deep into the 3D regime, where $`\alpha `$ becomes interaction dependent.
We now calculate the total dipole-dipole energy given that the density profile in $`z`$-direction is prescribed by the above Gaussian. In accordance with Eq. (Stability of quasi-two-dimensional Bose-Einstein condensates with dominant dipole-dipole interactions), the Gross-Pitaevskiว energy functional of the dipole-dipole interaction generally reads
$`H_{dd}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle d^3rd^3r^{}\rho (๐)V_{dd}(๐๐^{})\rho (๐^{})}`$ (4)
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{1}{(2\pi )^3}}{\displaystyle d^3k\stackrel{~}{\rho }(๐)\stackrel{~}{V}_{dd}(๐)\stackrel{~}{\rho }(๐)},`$
where the second line employs a convolution to Fourier space. The Fourier transform of the dipole-dipole interaction (2) takes the form $`\stackrel{~}{V}_{dd}(๐)=g_d[3k_z^2/(k_x^2+k_y^2+k_z^2)1]`$; using (3), we obtain the Fourier transform of the density $`\stackrel{~}{\rho }(๐)=\mathrm{exp}\left[\frac{1}{4}k_z^2d_z^2\right]\stackrel{~}{n}(k_x,k_y)`$. Integrating over the $`k_z`$ direction, the effective dipole-dipole energy is then given by a two-dimensional integral in $`(k_x,k_y)`$ space
$`H_{dd}`$ $`=`$ $`{\displaystyle \frac{g_d}{2}}{\displaystyle \frac{1}{(2\pi )^2}}{\displaystyle d^2k\stackrel{~}{n}(k_x,k_y)\stackrel{~}{n}(k_x,k_y)}`$ (5)
$`\times \left\{{\displaystyle \frac{2}{\sqrt{2\pi }d_z}}{\displaystyle \frac{3}{2}}\mathrm{exp}\left[{\displaystyle \frac{k^2d_z^2}{2}}\right]k\mathrm{Erfc}({\displaystyle \frac{kd_z}{\sqrt{2}}})\right\}`$
where the complementary error function Erfc $`(z)=1\mathrm{erf}(z)=1(2/\sqrt{\pi })_0^z\mathrm{exp}(t^2)๐t`$ and $`k=(k_x^2+k_y^2)^{1/2}`$; $`\stackrel{~}{n}(k_x,k_y)`$ is the Fourier transform of the 2D density. Employing the same procedure of integrating out the Gaussian in the $`z`$ direction, for the contact interaction part in the Gross-Pitaevskiว energy functional, yields the well-known result for the 2D effective coupling $`g_{2\mathrm{D}}=2\sqrt{2\pi }a_s/d_z`$ (valid in the limit that the 3D $`s`$-wave scattering length $`a_sd_z`$ Holzmann ).
From the relation (5) for the dipole-dipole contribution in the interaction energy, it follows that the Fourier transform of the total, contact plus dipole, interaction potential assumes the form
$$\stackrel{~}{V}_{\mathrm{tot}}^{2\mathrm{D}}(\zeta )=\frac{A}{\rho (0)\sqrt{\pi }d_z^3}\left\{1\frac{3R}{2}\zeta w\left[\frac{\zeta }{\sqrt{2}}\right]\right\},$$
(6)
where we made use of the $`w`$-function, related to Erfc$`(z)`$ by $`w(z)=\mathrm{exp}[z^2]\mathrm{Erfc}(z)`$ Abramowitz ; note . The dimensionless wavenumber $`\zeta =kd_z`$, and the two dimensionless parameters occurring in $`\stackrel{~}{V}_{\mathrm{tot}}^{2\mathrm{D}}(\zeta )`$, using the central 3D density $`\rho (0)`$, are defined to be
$$A=\frac{\rho (0)\sqrt{\pi }d_z^2g_d}{R},R=\frac{\sqrt{\pi /2}}{1+g_{3\mathrm{D}}/2g_d}.$$
(7)
The value of the parameter $`R`$ ranges from $`R=0`$ if $`g_d/g_{3\mathrm{D}}0`$, to $`R=\sqrt{\pi /2}`$ for $`g_d/g_{3\mathrm{D}}\mathrm{}`$. The second term in the curly brackets in Eq. (6) rapidly decreases as a function of $`\zeta `$ and approaches a constant, which is due to the fact that $`w(\zeta /\sqrt{2})\sqrt{2/\pi }\zeta ^1`$ for $`\zeta \mathrm{}`$, cf. the plot of $`\stackrel{~}{V}_{\mathrm{tot}}^{2\mathrm{D}}(\zeta )`$ in Fig. 1. For small $`\zeta `$, the quasi-2D Fourier transform behaves like $`\stackrel{~}{V}_{\mathrm{tot}}^{2\mathrm{D}}(\zeta )=\frac{A}{\rho (0)\sqrt{\pi }d_z^3}\left[1\frac{3R}{2}\zeta +\frac{3R}{\sqrt{2\pi }}\zeta ^2+๐ช(\zeta ^3)\right]`$. Observe that $`\stackrel{~}{V}_{\mathrm{tot}}^{2\mathrm{D}}(\zeta )`$ possesses a well-defined value at the origin $`\zeta =k=0`$, as opposed to the 3D Fourier transform of the dipole-dipole interaction potential.
From the Fourier transform of the interaction (6), we conclude that the squared Bogoliubov spectrum Bogoliubov ; Foldy , for excitations confined to the plane, $`\omega ^2=\rho (0)\sqrt{\pi }d_z\stackrel{~}{V}_{\mathrm{tot}}^{2\mathrm{D}}(k)k^2+k^4/4`$ is, in units of $`1/d_z^4`$, given by
$$ฯต^2(\zeta )=A\zeta ^2\left(1\frac{3R}{2}\zeta w\left[\frac{\zeta }{\sqrt{2}}\right]\right)+\frac{\zeta ^4}{4}.$$
(8)
The main observation of the present analysis is that, as opposed to the 3D case (or, in an exacerbated manner, the quasi-1D case EPJD ), the Bogoliubov spectrum (8) does not necessarily become unstable if the dipole interaction coupling exceeds the contact interaction coupling. Assuming $`g_{3\mathrm{D}}/g_d0`$, i.e. $`R\sqrt{\pi /2}`$, the above squared spectrum can assume negative values, and hence the excitation energies become imaginary, when $`A`$ exceeds the critical value $`A_c=3.446`$, cf. Fig. 2, where we show a sequence of four spectra for different $`A`$ at constant $`R=\sqrt{\pi /2}`$.
The critical value of $`A_c=3.446`$ corresponds to a critical dipole coupling given by $`(g_d)_c=2.436/\rho (0)d_z^2`$. Using numbers appropriate for atomic magnetic moments $`\mu _m=N_m\mu _B`$, we find that $`g_d=5.4\times 10^6\mu `$m $`MN_m^2`$ ($`M`$ is the mass of the atoms or molecules in units of the atomic mass unit $`=1.66\times 10^{27}`$ kg). In the case of electric dipoles, with moment $`d_e=N_e`$ Debye, we find $`g_d=6.2\times 10^2\mu `$m $`MN_e^2`$. The critical coupling $`(g_d)_c`$ for dominant dipole-dipole interactions then translates into a critical central 3D density $`\rho _c(0)=4.5\times 10^{16}\mathrm{cm}^3N_m^2\alpha ^1\omega _z[2\pi \times \mathrm{kHz}]`$ and $`\rho _c(0)=3.8\times 10^{12}\mathrm{cm}^3N_e^2\alpha ^1\omega _z[2\pi \times \mathrm{kHz}]`$ in the case of magnetic and electric dipoles, respectively.
The quantity $`\sqrt{\pi }\rho (0)g_d/R`$, a measure of the energy per particle (the chemical potential), is equal to $`A/d_z^2`$. For the system to be quasi-2D, we therefore need $`A\omega _zd_z^2=\alpha 1`$. Furthermore, the Bogoliubov spectrum in Fig. 2 does not possess any points where $`dฯต^2/d\zeta ^2`$=0 for values $`A<A_{\mathrm{min}}=1.249`$ (when $`R=\sqrt{\pi /2}`$). That is, a โrotonโ minimum cannot develop within the regime of quasi-2D. Thus we conclude that a quasi-2D purely dipolar system of bosons is always stable, and that the instability of the condensate takes place in the crossover region to 3D. By contrast, deep inside the 3D regime, $`g_{3\mathrm{D}}g_d`$ implies collapse of the condensate for any reasonable value of the density, and in the Thomas-Fermi limit the condensate will be unstable for any $`g_d`$ slightly exceeding $`g_{3\mathrm{D}}`$ RotonDipolar .
The critical value of $`A_c(R)`$, at which the gas becomes unstable, obtained by numerically finding the large momentum zeros of the Bogoliubov spectrum (8), exponentially increases for smaller $`R`$ (increasing $`g_{3\mathrm{D}}/g_d`$), and diverges at $`R=\frac{2}{3}\sqrt{\pi /2}`$ (i.e., for $`g_{3\mathrm{D}}=g_d`$), where the Fourier transform $`\stackrel{~}{V}_{\mathrm{tot}}^{2\mathrm{D}}`$ becomes positive everywhere. Comparing the critical values of $`A`$ thus obtained to those from a full 3D solution of the Bogoliubov-de Gennes equations RotonDipolar , the present approach reproduces the critical $`A`$ sufficiently accurate in the dipole-dominated case, for which the instability takes place in the crossover regime from quasi-2D to 3D. Deep into the 3D regime, the critical value $`A_c`$ calculated from (8) is strongly overestimated, mainly because it is exponentially sensitive on the exact form of the spectrum.
At nonzero temperature $`T`$, 2D Bose-Einstein condensates do not exist in the homogeneous case Hohenberg , while trapping in the plane enables the existence of a condensate also at finite $`T`$ Fischer . On the other hand, in the presently discussed $`T=0`$ case, it is a well-established fact that even without trapping, Bose-Einstein condensation occurs in two spatial dimensions, see, e.g., PitaLongRange . We next turn to a discussion of the zero temperature value of the number density of excitations above the condensate, the so-called quantum depletion. To this end, we use a mode expansion for the annihilation operators $`\widehat{\chi }_๐`$ of the original bosons in terms of the Bogoliubov quasiparticle operators $`\widehat{a}_๐,\widehat{a}_๐^{}`$ Schutzhold :
$$\widehat{\chi }_๐=\sqrt{\frac{๐^2}{2ฯต_๐}}\left[\left(\frac{1}{2}+\frac{ฯต_๐}{๐^2}\right)\widehat{a}_๐+\left(\frac{1}{2}\frac{ฯต_๐}{๐^2}\right)\widehat{a}_๐^{}\right].$$
(9)
The above form of the Bogoliubov transformation results, after inversion, in the usual phonon quasiparticle operators at low momenta, and gives $`\widehat{\chi }_๐=\widehat{a}_๐`$ at $`๐\mathrm{}`$, i.e., the quasiparticles and the bare bosons become, as required, identical at large momenta.
The quantum depletion density at zero temperature is calculated by evaluating the expectation of $`\widehat{\chi }^{}\widehat{\chi }`$ in the quasiparticle vacuum defined by $`\widehat{a}_๐|\mathrm{vac}=0`$:
$`\widehat{\chi }^{}\widehat{\chi }`$ $`=`$ $`{\displaystyle \frac{1}{2\pi d_z^2}}{\displaystyle _0^{\mathrm{}}}๐\zeta {\displaystyle \frac{\zeta ^3}{2ฯต(\zeta )}}\left({\displaystyle \frac{1}{2}}{\displaystyle \frac{ฯต(\zeta )}{\zeta ^2}}\right)^2.`$ (10)
The instability generally happens because in-plane excitations become of imaginary frequency RotonDipolar ; above the calculated critical value of $`A=A_c`$, the in-plane excitations have vanishing energy, and the quantum depletion diverges. Using the in-plane momentum integral above to calculate the depletion, with the effective dispersion relation (8), assumes that the corrections due to the neglected out-of-plane Bogoliubov excitations, characterized by transverse quantum numbers, are small. This is justified because the dominant contribution to quantum depletion is, for a dilute Bose gas, from large momentum excitations with (approximately) vanishing energy. Out-of-plane excitations have large energies at the relevant in-plane momenta of order $`1/d_z`$; they thus do not contribute significantly to the depletion. In Fig. 3, we show the result for the quantum depletion in the case of dipole-interaction dominated condensates, for which we can expect the in-plane spectrum (8) to be sufficiently accurate up to $`AA_c`$. At the critical value $`A=A_c(\sqrt{\pi /2})3.4`$ (dotted vertical line in Fig. 3), the condensate depletion diverges, and the mean-field condensate will yield to a new quantum phase.
In conclusion, we have shown that dipolar quasi-2D Bose-Einstein condensates are extremely stable systems as compared to their 3D counterparts. This offers the potential of approaching, starting from the mean-field physics of condensates, a strongly correlated regime of dilute atomic/molecular gases with long-range interactions. A conceivable experimental procedure is to start from a quasi-2D dipolar condensate, and, by decreasing $`\omega _z`$, to enter the crossover regime to 3D, where the quantum depletion of the condensate becomes large. If the dipole-dipole interaction dominates, i.e., in the limit $`R\sqrt{\pi /2}`$, the quantity $`A/\sqrt{2}=\rho (0)d_z^2g_d=2.4\rho (0)/\rho _c(0)`$, where $`\rho _c(0)`$ is the critical 3D density for Bogoliubov excitations above the mean-field condensate to be stable, discussed in the above. Hence $`\rho (0)d_z^2g_d`$ measures the diluteness of the system, that is the proximity to the critical region. Numerical values are
$`\rho (0)d_z^2g_d`$ $`=`$ $`5.5\times 10^3\alpha N_m^2{\displaystyle \frac{\rho (0)[10^{14}\mathrm{cm}^3]}{\omega _z[2\pi \times \mathrm{kHz}]}},`$
$`\rho (0)d_z^2g_d`$ $`=`$ $`0.63\alpha N_e^2{\displaystyle \frac{\rho (0)[10^{12}\mathrm{cm}^3]}{\omega _z[2\pi \times \mathrm{kHz}]}},`$ (11)
in the magnetic and electric cases, respectively. According to the first equation, magnetically dipolar gases will allow for the realization of a purely dipolar condensate, such that the depletion still remains small. On the other hand, it is evident from the second line of Eq. (11), that for electric dipoles the densities at which strong depletion of the condensate sets in are rather small, even for strongly increased axial trapping.
Among further lines of research offered by the stability potential of quasi-2D dipolar condensates are the phenomena expected when they are set in rotation. When both contact and dipole interaction are present, various of these phenomena have been explored in Cooper ; quantum Hall states of purely dipolar Fermi gases were studied in BaranovQH . The interplay of quasi-two-dimensionality and rotation may generally yield interesting new physics Sinha . It will, furthermore, be of interest to determine the change in the systemโs stability region when the dipoles are no longer locked in one direction, and thus to investigate the influence of spin waves on the stability of dipolar quantum gases of bosons.
I thank C. Zimmermann, N. Schopohl, L. Santos, C. Iniotakis, and N. R. Cooper for helpful discussions.
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# DCPT-05/27 EMPG-05-12 hep-th/0506092 On T-folds, G-structures and Supersymmetry.
## I Introduction
The usual procedure for constructing background field configurations of string and M theory involves two steps. Firstly, the low energy fields are defined on coordinate patches in an atlas covering some manifold. A consistent global picture is then created in the second step by insisting that where coordinate patches overlap the different field configurations describe the same physical setup. This is usually achieved by relating the relevant local fields on different coordinate patches by diffeomorphisms and gauge transformations. We shall refer to such backgrounds as โgeometricโ as the resulting configurations include a globally well defined metric tensor.
One can carry out the second step of requiring a consistent global picture in other ways (see for the first general consideration and discussion of this). Any symmetry of the theory can be used to join together coordinate patches while still fulfilling the criterion that the various field configurations on overlapping patches describe the same physics. For example, in a toroidally compactified type II string theory one could use elements of the $`O(d,d)`$ T-duality group as transition functions . This results in what Hull has called a โT-foldโ . In general, there is no globally defined metric tensor on a T-fold. This is a consequence of the fact that T-duality transformations mix up the internal components of metric with those of the NS-NS two form. We will therefore follow the literature in referring to such spaces where the metric โjumpsโ between coordinate patches in this manner as โnon-geometricโ. The idea of patching a manifold together with duality transformations, to obtain these so-called โduality-foldsโ , has been pursued by a number of authors over the last few years (see for this and related work). In passing we note that it appears that the use of supersymmetry in performing the โsecond stepโ above has not yet been considered in the literature.
Given their unusual nature it is helpful to briefly review some of the evidence that string theory can make sense on such backgrounds. Firstly some T-folds can be shown to be T-dual to various, more standard, geometric spaces. For example, the T-dual of certain torus compactifications with NS-NS flux (with an appropriate domain wall solution in the external space) are T-folds (see for example ). Thus one is left with two possibilities. Either the T-fold obtained in such a manner is as good a string background as the torus compactification with flux, or T-duality is not a feature of the full theory in that, for some reason, the torus compactificationโs T-dual is not a good string background.
A second piece of evidence which indicates that we should take T-fold backgrounds seriously within the context of string theory is provided by their relation to certain well defined conformal field theories (CFTs) . It is well known that certain smooth geometric backgrounds of string theory have limits in which they become orbifolds of flat space. The associated world sheet CFT can then be described in detail. In a similar manner there are limits of certain T-fold compactifications which are associated with well understood consistent CFTs . In this case the relevant orbifolds are asymmetric examples in which the left and right moving degrees of freedom on the string world sheet are acted on differently by the orbifolding .
Once it is accepted that it makes sense to consider string theory on these T-fold backgrounds the next obvious question is why are these configurations interesting? There are a multitude of interesting consequences which follow from considering string and M-theory on various symmetry folds. One strong motivation for studying such backgrounds is that they have been found to typically possess fewer unstabilised moduli than more conventional compactifications . It is reasonably straight forward to see why this is so. Unstabilised moduli in lower dimensional theories correspond to unspecified integration constants in the associated higher dimensional vacuum solutions. The use of a larger set of symmetries (such as the full T-duality group instead of just some โgeometricโ subgroup) in patching together a compactification manifold effectively constitutes a more restrictive set of boundary conditions on the possible vacuum solutions. As such the solutions tend to have fewer unspecified integration constants, leading to fewer unstabilised moduli in the lower dimensional theories. However, within the context of phenomenology, where considerations of moduli stabilisation would be most important, the use of these spaces raises various problems. This is because, for all values of the moduli, the known examples of non-trivial T-folds contain cycles of string scale size. This leads to concerns about the lack of a mass gap above the standard supergravity modes as well as concern that non-perturbative effects could well be important. These difficulties may well be resolvable in some cases. For example, it is easy to build examples of T-folds where one can show that, despite these comments, the size of non-perturbative effects in the superpotential is controlled by the volume of the base manifold and not that of the fibre . In addition, there are such examples where only a finite number of extra light states would have to be included before a mass gap with the rest of the spectrum is obtained. However, in the interests of being conservative, and to show that our work is of interest whether or not these issues are resolvable, we shall state our motivation in this paper in terms of a different use for T-fold compactification.
An uncontroversial use of T-fold compactification is to regard it as a formal tool for finding new massive supergravity theories. Compactification of massless higher dimensional type II supergravities on T-folds results in examples of lower dimensional massive supergravities, which in some cases were previously unknown . The supergravities obtained in this manner can be thought of as the completion of the class of such theories obtained by compactification on geometric spaces (for examples based upon Scherk Schwarz reduction see ). In particular, these new sugras fit nicely into certain classifications of such theories, one of which is based on an $`O(d,d)`$ group of transformations .
One of the most basic questions one would like to answer about these lower dimensional supergravity theories is how much supersymmetry they possess. This reduces to a question about the nature of the truncation of the higher dimensional fields which is used in the reduction. More precisely, the question is which of the higher dimensional supersymmetries are compatible with the different forms of truncation that one might use? This is the question we shall answer in part in this paper - how many supersymmetries are enjoyed by a theory obtained by T-fold compactification with a certain natural choice of truncation?
Our discussion will make it clear how to calculate the number of such supersymmetries associated with a large class of theories resulting from T-fold compactification. However, in the interests of clarity we will illuminate our considerations with the example of compactifications of type II string theory on T-folds which take the form of a $`T^d`$ bundle over $`S^1`$, where the torus experiences a monodromy in $`SO(d,d,Z)`$ around the base $`S^1`$. Similar compactifications have been considered in Refs. . We will extend this work by incorporating extra terms originating from the NS B-field in the dimensional reduction and will then proceed to calculate the amount of supersymmetry that is manifest in the compactified theories obtained. To do this we will use and extend work in the literature where the rules for the transformations of supersymmetry parameters under T-duality have been given. We will then use a modified version of the criterion for preservation of supersymmetry in dimensional reduction, phrased in terms of G-structures , to determine the number of supersymmetries manifest in the theory. The approach we will follow, based on the work of , will lead us to consider the supersymmetry of theories which result from compactification upon T-folds constructed using a certain minimal version of compactified supergravity. In particular, we will not consider cases arising from constructions where a coset reformulation of the theory obtained by reducing on the fibre of the T-fold is utilised in forming the overall compact space. See Appendix A and later sections of the paper for more details.
The plan of this paper is as follows. In section II we shall introduce the T-fold backgrounds that we will be using as examples in the rest of the paper. We will then review dimensional reduction on such a background to illustrate the kind of massive supergravity that can be obtained in this manner. In section III we shall describe how to calculate the amount of supersymmetry associated with a class of theories obtained by T-fold reduction using the examples of the previous section to illustrate our method. We will describe first how the usual arguments proceed in the case of geometric dimensional reduction before going on to describe how this analysis changes in the T-fold case. In section IV we briefly conclude.
## II Examples of Massive Supergravities from T-fold Reduction.
### A Examples of T-fold vacua for dimensional reduction.
The examples of duality-folds that we will consider are the class of T-folds which have been studied in . We start by considering a type II superstring theory on a $`T^d`$ bundle over $`M_{9d}\times S^1`$ (the discussion and notation of this subsection follows that of ). However, we are not going to take the trivial vacuum on this space. Instead we consider a situation where the various fields in the theory have a dependence on the $`S^1`$, which has coordinate $`yy+2\pi `$, of the following form,
$`\psi (y)=g(y)[\psi ]`$ (1)
Here $`\psi `$ is a general field which is taken to be independent of the $`T^d`$ directions. We then make a Scherk-Schwarz ansatz for $`g(y)`$,
$`g(y)=\mathrm{exp}\left({\displaystyle \frac{yT}{2\pi }}\right)`$ (2)
where $`T`$ is in the Lie algebra of $`SO(d,d)`$, which is a subgroup of the T-duality group associated with the $`T^d`$ fibre. This ansatz ensures that a solution to the lower dimensional theory, obtained by compactification on this vacuum, is also a solution to the higher dimensional equations of motion up to the relevant approximations. As such it guarantees that we do not need to worry about such subtleties as compensators in the dimensional reduction. The resulting field configurations are not periodic on traversing the $`S^1`$. Instead there is a monodromy,
$`(g)=e^T`$ (3)
This is simply an element of $`SO(d,d)`$. In other words, when we traverse the circle once the fields are not identified in the usual way, but instead come back to themselves up to an $`SO(d,d)`$ element of the T-duality group. In fact we should choose the monodromy to be within $`SO(d,d,Z)`$ as this is the relevant symmetry subgroup of the full theory when massive states, which we have truncated, are included. Note that since we are making this Scherk-Schwarz ansatz we will only be able to consider elements of the full T-duality group, $`O(d,d,Z)`$, which are continuously connected to the identity in $`O(d,d,R)`$.
The resulting vacuum is an example of a T-fold . The field configurations at $`y=0`$ and $`y=2\pi `$ are not simply related by diffeomorphisms and gauge transformations. A less trivial element of the T-duality group is required to transform them into one another. This vacuum is then non-geometric in the sense that there is not a globally well defined metric in ten dimensions. This is because the transition functions mix up the metric and NS-NS two form on the toroidal fibre when we go once around the circle.
### B Massive supergravities from dimensional reduction on T-folds.
We will now describe how the low energy effective action associated with a dimensional reduction on one of the T-folds of the previous subsection is obtained. Many, but not all, of the terms in the reduced theory that we will present have been obtained elsewhere . However, it is useful to provide a discussion of the dimensional reduction here both to keep this paper reasonably self contained and also to complete the dimensional reduction of the ten dimensional NS-NS sector. To our knowledge some terms of this reduction, descending from the three form field strength, have yet to appear in the literature. Since this paper first appeared the relation between reductions of the type we are going to present here and so called โtwisted torusโ reductions (see for example ) has been demonstrated in .
For the sake of brevity we will consider a $`T^2`$ fibre, i.e. we will compactify on a $`T^2`$ bundle over $`S^1`$. The generalisation of this to the $`T^d`$ case is straight-forward. The starting point for our dimensional reduction is the ten dimensional low energy effective action of type II superstring theory. Again, for the sake of brevity, and because it is common to both IIA and IIB, we only consider the NS-NS sector.
$`S_{10}={\displaystyle _{^{10}}}๐x\sqrt{g}e^\varphi \left[R+(\varphi )^2{\displaystyle \frac{1}{12}}H^2\right]`$ (4)
The simplest procedure to follow for the reduction is to first reduce on the $`T^2`$ fibre and then perform a further reduction on the base space incorporating our non-trivial twist. As such we require an effective eight dimensional action for the reduction of the above theory on the $`T^2`$ fibre in which the $`O(2,2)`$ duality group is manifest. Such an eight dimensional action is given in the seminal paper by Maharana and Schwarz . We will not repeat their calculation here but shall simply quote the results. They obtain,
$`S_{8D}={\displaystyle _^8}d^8x\sqrt{g}e^{\varphi ^{(8)}}\left[R^{(8)}+(\varphi ^{(8)})^2+{\displaystyle \frac{1}{8}}\text{tr}(M^1M){\displaystyle \frac{1}{4}}_{\mu \nu }^iM_{ij}^1^{\mu \nu j}{\displaystyle \frac{1}{12}}H^2\right],`$ (5)
where the three form field strength and eight dimensional dilaton are defined by,
$`H_{\mu \nu \rho }`$ $`=`$ $`_\mu B_{\nu \rho }{\displaystyle \frac{1}{2}}๐_\mu ^i\eta _{ij}_{\nu \rho }^i+\text{cyclic permutations},`$ (6)
$`\varphi ^{(8)}`$ $`=`$ $`\varphi {\displaystyle \frac{1}{2}}\mathrm{log}detG.`$ (7)
In these expression we have also defined,
$`M=\left(\begin{array}{cc}G^1& G^1B\\ BG^1& GBG^1B\end{array}\right)`$ (10)
where $`B`$ and $`G`$ are the NS two form and metric on the fibre respectively. Finally, if we denote the two coordinate indices on the fibre as $`\theta ^i`$, the gauge potentials are given by
$`๐_\mu ^i`$ $`=`$ $`G_{\mu \theta ^i}i=1,2,`$ (12)
$`๐_\mu ^i`$ $`=`$ $`B_{\mu \theta ^{i2}}+B_{\theta ^{i2}\theta ^j}๐_\mu ^ji=3,4,`$ (13)
and the field strengths are simply $`^i=d๐^i`$.
In this formulation, the various fields transform under an $`O(2,2)`$ transformation as follows,
$`M`$ $``$ $`\mathrm{\Omega }M\mathrm{\Omega }^T`$ (14)
$`๐_\mu `$ $``$ $`\mathrm{\Omega }๐_\mu `$ (15)
Here $`\mathrm{\Omega }`$ is the $`O(2,2)`$ matrix, satisfying $`\mathrm{\Omega }^T\eta \mathrm{\Omega }=\eta `$, where in our conventions the invariant metric is given by,
$`\eta =\left(\begin{array}{cc}\hfill 0& \hfill 1\\ \hfill 1& \hfill 0\end{array}\right).`$ (18)
Given the $`O(2,2)`$ covariant form of the action in (5), we now make an ansatz for the dimensional reduction on the base of our T-fold. From the discussion in the previous subsection our ansatz takes the following form,
$`ds_8^2`$ $`=`$ $`g_{\alpha \beta }^7(x^\gamma )dx^\alpha dx^\beta +e^{2\alpha }\left(dy+A_\alpha ^B(x^\gamma )dx^\alpha \right)^2`$ (19)
$`\varphi ^{(8)}`$ $`=`$ $`\varphi ^{(8)}(x^\alpha )`$ (20)
$`M`$ $`=`$ $`\mathrm{\Omega }(y)M_0(x^\alpha )\mathrm{\Omega }^T(y)`$ (21)
$`๐^i`$ $`=`$ $`\mathrm{\Omega }_j^i(y)๐_\alpha ^{0j}(x^\beta )dx^\alpha +\mathrm{\Omega }_j^i(y)๐_y^{0j}(x^\beta )dy`$ (22)
$`B`$ $`=`$ $`{\displaystyle \frac{1}{2}}B_{\alpha \beta }(x^\gamma )dx^\alpha dx^\beta +B_\alpha ^B(x^\gamma )dx^\alpha dy`$ (23)
where $`x^\alpha `$ are coordinates on the 7-dimensional non-compact space and
$`\mathrm{\Omega }(y)=e^{\frac{y}{2\pi }T}`$ (24)
defines the monodromy. We can now perform the reduction on the base space.
The reduction of the Ricci scalar and dilaton terms to seven dimensions is unaffected by our duality twists. As such the expressions for these terms are the standard ones.
$`S_1={\displaystyle _^7}๐x\sqrt{g^7}e^{\varphi ^{(7)}}\left[R^{(7)}(\alpha )^2+(\varphi ^{(7)})^2{\displaystyle \frac{1}{4}}e^{2\alpha }(F^B)^2\right]`$ (25)
Here we have defined the usual seven dimensional dilaton $`\varphi ^{(7)}=\varphi ^{(8)}\alpha `$.
The reduction of the remaining scalar kinetic terms does get modified by the duality twist. The resulting kinetic terms are the usual ones but the derivative involved is replaced by a covariant derivative with a non-trivial connection. In addition a potential for these fields is obtained. This potential and its properties have been discussed in some detail in . The relevant terms are,
$`S_2={\displaystyle _^7}๐x\sqrt{g^7}e^{\varphi ^{(7)}}\left[{\displaystyle \frac{1}{8}}\text{tr}\left(DM_0^1DM_0\right){\displaystyle \frac{1}{4(2\pi )^2}}\text{tr}(T^TM_0^1TM_0+T^2)\right].`$ (26)
Here we have defined a covariant derivative of $`M_0`$ as follows.
$`D_\alpha M_0=_\alpha M_0{\displaystyle \frac{TM_0}{2\pi }}A_\alpha ^B{\displaystyle \frac{M_0T^T}{2\pi }}A_\alpha ^B`$ (27)
where $`M_0(x^\alpha )`$ is defined by $`M(x^\alpha ,y)=\mathrm{\Omega }(y)M_0(x^\alpha )\mathrm{\Omega }(y)^T`$.
The reduction of the vector field strength terms is a little more subtle. If we wish the result of the dimensional reduction to take an elegant form then we must ensure that we have made reasonable choices for our definitions of the seven dimensional fields .
For example, consider the behaviour of the vector fields under an $`x^\alpha `$ dependent shift of the $`y`$ coordinate, $`dydy^{}=dy+_\alpha \omega ^ydx^\alpha `$.
$`๐๐^{}=๐_\alpha ^{0i}dx^\alpha +๐_y^{0i}\left(dy+_\alpha \omega ^ydx^\alpha \right)`$ (28)
Thus we find that, while $`๐_y^{0i}`$ is unchanged by such a transformation, $`๐_\alpha ^{0i}`$ is not a gauge invariant definition of a seven dimensional field.
$`๐_\alpha ^{0i}๐_\alpha ^{0i}=๐_\alpha ^{0i}_\alpha \omega ^y๐_y^{0i}`$ (29)
The only seven dimensional gauge field we wish to transform non-trivially under a shift of the base coordinate is the associated Kaluza Klein gauge field $`A^B`$. As such, instead of using $`๐_\alpha ^{0i}`$ as a seven dimensional field we make a field redefinition to obtain a quantity that has the properties we desire.
$`๐_\alpha ^{7i}=๐_\alpha ^{0i}A_\alpha ^B๐_y^{0i}`$ (30)
Having decided on a set of definitions for the seven dimensional fields we may proceed with the dimensional reduction of the vector field strength terms. These terms are affected by the duality twists and so we obtain modifications to the result obtained in an untwisted reduction which depend on the generator $`T`$. We find for the components of the field strength,
$`_{\alpha \beta }^j`$ $`=`$ $`\mathrm{\Omega }_i^j(_{\alpha \beta }^{7i}+2A_{[\beta }^B_{\alpha ]}๐_y^{0i}+๐_y^{0i}F_{\alpha \beta }^B)`$ (31)
$`_{\alpha y}^j`$ $`=`$ $`\mathrm{\Omega }_i^j(_\alpha ๐_y^{0i}{\displaystyle \frac{T}{2\pi }}A_\alpha ^B๐_y^{0i}{\displaystyle \frac{T}{2\pi }}๐_\alpha ^{7i}).`$ (32)
Here we have defined a field strength for $`๐_\alpha ^{7i}`$ in the usual way. As always in such reductions the easiest way to proceed is to switch to using an orthonormal basis during the reduction and then return to a coordinate basis once the result has been obtained. We define the following vielbeins (barred indices denote those in an orthonormal basis).
$`e_{\overline{\nu }}^\mu =\left(\begin{array}{cc}\hfill e_{\overline{\alpha }}^{7\alpha }& \hfill e_{\overline{\alpha }}^{7\beta }A_\beta ^B\\ \hfill 0& \hfill e^\alpha \end{array}\right)e_\nu ^{\overline{\mu }}=\left(\begin{array}{cc}\hfill e_\alpha ^{7\overline{\alpha }}& \hfill e^\alpha A_\alpha ^B\\ \hfill 0& \hfill e^\alpha \end{array}\right)`$ (37)
In these expressions, $`e_\alpha ^{7\overline{\gamma }}e_\beta ^{7\overline{\gamma }}=g_{\alpha \beta }^7`$, etc. in the usual manner. Using these definitions we then find the following.
$`{\displaystyle \frac{1}{4}}^iM_{ij}^1^j={\displaystyle \frac{1}{4}}_{\overline{\alpha }\overline{\beta }}^iM_{ij}^1^{j\overline{\alpha }\overline{\beta }}{\displaystyle \frac{1}{2}}_{\overline{\alpha }\overline{y}}^iM_{ij}^1^{j\overline{\alpha }\overline{y}}`$ (38)
where,
$`_{\overline{\alpha }\overline{\beta }}^i`$ $`=`$ $`e_{\overline{\alpha }}^{7\alpha }e_{\overline{\beta }}^{7\beta }_{\alpha \beta }^i2e_{[\overline{\alpha }}^{7\alpha }e_{\overline{\beta }]}^{7\gamma }A_\gamma ^B_{\alpha y}^i`$ (39)
$`_{\overline{\alpha }\overline{y}}^i`$ $`=`$ $`e_{\overline{\alpha }}^{7\alpha }e^\alpha _{\alpha y}^i.`$ (40)
After a little algebra the following result is obtained for the reduction of the vector field strength terms.
$`S_3={\displaystyle _^7}๐x\sqrt{g^7}e^{\varphi ^{(7)}}\left[{\displaystyle \frac{1}{4}}\stackrel{~}{}_{\alpha \beta }^iM_{0ij}^1\stackrel{~}{}^{j\alpha \beta }{\displaystyle \frac{1}{2}}e^{2\alpha }D_\alpha ๐_y^{0i}M_0^{1ij}D^\alpha ๐_y^{0j}\right]`$ (41)
Here we have used the following definitions of field strengths and covariant derivatives.
$`\stackrel{~}{}_{\alpha \beta }^i`$ $`=`$ $`\stackrel{~}{}_{\alpha \beta }^{7i}+๐_y^{0i}F_{\alpha \beta }^B`$ (42)
$`\stackrel{~}{}_{\alpha \beta }^{7i}`$ $`=`$ $`_{\alpha \beta }^{7i}+2A_{[\beta }^B{\displaystyle \frac{T_j^i}{2\pi }}๐_{\alpha ]}^{7j}`$ (43)
$`D_\alpha ๐_y^{0i}`$ $`=`$ $`_\alpha ๐_y^{0i}{\displaystyle \frac{T_j^i}{2\pi }}A_\alpha ^B๐_y^{0j}{\displaystyle \frac{T_j^i}{2\pi }}๐_\alpha ^{7j}`$ (44)
This just leaves us with the dimensional reduction of the anti-symmetric tensor terms to perform. Despite the fact that the field strength $`H`$ is itself invariant under the duality transformations we are using, the matrix $`T`$ still enters this part of the reduction. The reason for this is that the field strength, while invariant overall, is made up of components which transform non trivially as can be seen from its definition (6).
As was the case for the reduction of the vector field part of the action we need to be careful to choose โgauge invariantโ definitions of our seven dimensional degrees of freedom. We find that $`B_\alpha ^B`$ has the properties we desire but that $`B_{\alpha \beta }`$ does not. Therefore we define a new two form potential as follows.
$`B_{\alpha \beta }^7=B_{\alpha \beta }+A_{[\alpha }^BB_{\beta ]}^B`$ (45)
Then we can proceed with the dimensional reduction of these terms making use of our orthonormal basis as we did for the vector field strength piece. We find,
$`{\displaystyle \frac{1}{12}}H^2={\displaystyle \frac{1}{12}}H_{\overline{\alpha }\overline{\beta }\overline{\gamma }}H_{\overline{\alpha }\overline{\beta }\overline{\gamma }}{\displaystyle \frac{1}{4}}H_{\overline{\alpha }\overline{\beta }\overline{y}}H_{\overline{\alpha }\overline{\beta }\overline{y}}`$ (46)
where,
$`H_{\overline{\alpha }\overline{\beta }\overline{\gamma }}`$ $`=`$ $`e_{\overline{\alpha }}^{7\alpha }e_{\overline{\beta }}^{7\beta }e_{\overline{\gamma }}^{7\gamma }H_{\alpha \beta \gamma }(e_{\overline{\alpha }}^{7\delta }e_{\overline{\beta }}^{7\beta }e_{\overline{\gamma }}^{7\gamma }A_\delta ^BH_{y\beta \gamma }+\text{2 perms.})`$ (47)
$`H_{\overline{\alpha }\overline{\beta }\overline{y}}`$ $`=`$ $`e_{\overline{\alpha }}^{7\alpha }e_{\overline{\beta }}^{7\beta }e^\alpha H_{\alpha \beta \gamma }`$ (48)
and
$`H_{\alpha \beta \gamma }`$ $`=`$ $`_\alpha B_{\beta \gamma }^7{\displaystyle \frac{1}{2}}F_{\alpha \beta }^BB_\gamma ^B+{\displaystyle \frac{1}{2}}H_{\alpha \beta }^BA_\gamma ^B{\displaystyle \frac{1}{2}}(๐_\alpha ^{7i}+A_\alpha ^B๐_y^{0i})\eta _{ij}(_{\beta \gamma }^{7i}+๐_y^{0i}F_{\beta \gamma }^B)(๐_\alpha ^{7i}+A_\alpha ^B๐_y^{0i})\eta _{ij}A_{[\gamma }^B_{\beta ]}๐_y^{0i}`$ (49)
$`H_{\alpha \beta y}`$ $`=`$ $`H_{\alpha \beta }^B{\displaystyle \frac{1}{2}}๐_y^{0i}\eta _{ij}(_{\alpha \beta }^j+2A_{[\beta }^B_{\alpha ]}๐_y^{0i}+๐_y^{0i}F_{\alpha \beta }^B)(๐_{[\alpha }^{7i}+A_{[\alpha }^B๐_{|y|}^{0i})\eta _{ij}D_{\beta ]}๐_y^{0i}.`$ (50)
In the above $`H^B`$ is the field strength associated with $`B^B`$ defined in the usual manner. After a little algebra we then find the following for the final piece of our seven dimensional action.
$`S_4={\displaystyle _^7}๐x\sqrt{g^7}e^{\varphi ^{(7)}}\left[{\displaystyle \frac{1}{12}}(H^7)^2{\displaystyle \frac{1}{4}}e^{2\alpha }(H_{\alpha \beta }^B{\displaystyle \frac{1}{2}}๐_y^{0i}\eta _{ij}\stackrel{~}{}_{\alpha \beta }^j๐_{[\alpha }^{7i}\eta _{ij}D_{\beta ]}๐_y^{0i})^2\right]`$ (51)
In this expression we have used the following definition of the seven dimensional three-form field strength.
$`H_{\alpha \beta \gamma }^7=_\alpha B_{\beta \gamma }^7{\displaystyle \frac{1}{2}}F_{\alpha \beta }^BB_\gamma ^B{\displaystyle \frac{1}{2}}H_{\alpha \beta }^BA_\gamma ^B{\displaystyle \frac{1}{2}}๐_\alpha ^{7i}\eta _{ij}\stackrel{~}{}_{\beta \gamma }^j+\text{ cyclic perms.}`$ (52)
Combining all the pieces we have obtained in this subsection we then obtain the following for the reduction of the NS-NS sector of type II string theory on this class of T-folds.
$`S_7={\displaystyle _^7}\sqrt{g^7}e^{\varphi ^{(7)}}[R^{(7)}(\alpha )^2+(\varphi ^{(7)})^2{\displaystyle \frac{1}{4}}e^{2\alpha }(F^B)^2`$ (53)
$`+{\displaystyle \frac{1}{8}}\text{tr}\left(DM_0^1DM_0\right){\displaystyle \frac{1}{4(2\pi )^2}}\text{tr}(T^TM_0^1TM_0+T^2)`$ (54)
$`{\displaystyle \frac{1}{4}}\stackrel{~}{}_{\alpha \beta }^iM_{0ij}^1\stackrel{~}{}^{j\alpha \beta }{\displaystyle \frac{1}{2}}e^{2\alpha }D_\alpha ๐_y^{0i}M_0^{1ij}D^\alpha ๐_y^{0j}`$ (55)
$`{\displaystyle \frac{1}{12}}(H^7)^2{\displaystyle \frac{1}{4}}e^{2\alpha }(H_{\alpha \beta }^B{\displaystyle \frac{1}{2}}๐_y^{0i}\eta _{ij}\stackrel{~}{}_{\alpha \beta }^j๐_{[\alpha }^{7i}\eta _{ij}D_{\beta ]}๐_y^{0i})^2]`$ (56)
In this action we have the following field strength and covariant derivative definitions.
$`F_{\alpha \beta }^B`$ $`=`$ $`2_{[\alpha }A_{\beta ]}^B`$ (57)
$`H_{\alpha \beta }^B`$ $`=`$ $`2_{[\alpha }B_{\beta ]}^B`$ (58)
$`\stackrel{~}{}_{\alpha \beta }^i`$ $`=`$ $`2_{[\alpha }๐_{\beta ]}^{7i}+๐_y^{0i}F_{\alpha \beta }^B+2A_{[\beta }^B{\displaystyle \frac{T_j^i}{2\pi }}๐_{\alpha ]}^{7j}`$ (59)
$`H_{\alpha \beta \gamma }^7`$ $`=`$ $`_\alpha B_{\beta \gamma }^7{\displaystyle \frac{1}{2}}F_{\alpha \beta }^BB_\gamma ^B{\displaystyle \frac{1}{2}}H_{\alpha \beta }^BA_\gamma ^B{\displaystyle \frac{1}{2}}๐_\alpha ^{7i}\eta _{ij}\stackrel{~}{}_{\beta \gamma }^j+\text{ cyclic perms.}`$ (60)
$`D_\alpha M_0`$ $`=`$ $`_\alpha M_0{\displaystyle \frac{TM_0}{2\pi }}A_\alpha ^B{\displaystyle \frac{M_0T^T}{2\pi }}A_\alpha ^B`$ (61)
$`D_\alpha ๐_y^{0i}`$ $`=`$ $`_\alpha ๐_y^{0i}{\displaystyle \frac{T_j^i}{2\pi }}A_\alpha ^B๐_y^{0j}{\displaystyle \frac{T_j^i}{2\pi }}๐_\alpha ^{7j}`$ (62)
It is easy to show that this action is that obtained by reducing the NS-NS part of the ten dimensional action on a three torus, supplemented by various mass and charge terms which are determined by the duality twist generator $`T`$. This is an example of the kind of massive supergravity we will consider in the next section. In particular, we will calculate the amount of supersymmetry this type of theory possesses.
In the low energy supergravity limit we have discussed here, the compactifications which have monodromies within the same $`O(d,d,R)`$ conjugacy class are equivalent and so give rise to the same reduced theory . This point will tie in nicely with the discussion of supersymmetry that will appear in the next section. It should also be pointed out that the above reduction implicitly uses the fact that the usual rules for flux quantisation in string theory are modified when the compact cycles being considered are nongeometric. There are scalar fields in the lower dimensional theory which originate from the components of the higher dimensional NS two form with indices on the compact space. For some of these fields, the fact that these degrees of freedom appear as lower dimensional fields which are able to vary in value depends crucially on the fact that the integral of the three-form field strength over certain cycles is not quantised in the usual manner.
## III Supersymmetry and T-fold reduction
This section is split into three main parts. Firstly we shall briefly review the standard discussion of how one determines the amount of supersymmetry associated with a geometric dimensional reduction. In the second subsection we will describe the transformation of various quantities under T-duality. Finally, we will use this information to describe how the standard discussion for the degree of supersymmetry of a dimensionally reduced theory is modified in the T-fold case.
A geometrical dimensional reduction proceeds in a series of steps. First one chooses a Riemannian manifold $`X`$, with metric and ansatze for the other fields, on which to compactify the extra dimensions. One then proceeds to rewrite the fields in the higher dimensional theory in a manner that respects the symmetries of this background. The fields are split up into representations of the subgroup of the full symmetry group of the theory that is preserved by the background. Up until this point nothing in the theory has actually changed. One has simply relabeled various fields in such a way as to make the rest of the dimensional reduction easier. The change to the content of the theory comes from the rest of the process. The extra dimensional dependence of the fields is expanded in a series of modes about the background and various parts of the expanded and rewritten higher dimensional fields, such as some of these modes, are discarded. The fields are then substituted into the higher dimensional action, compensators if needed are calculated, and the extra dimensions are integrated out to leave the lower dimensional theory. All of this has to be performed in a consistent manner so that every solution to the resulting low energy theory can be associated to a solution of the higher dimensional one, up to any approximations that are made in the truncation - which defines the lower dimensional theory. Furthermore, one might wish the truncation to be โphysicalโ in that the modes one discards in this process should describe fluctuations which are more massive than those which are kept. The classic example of this is Kaluza Klein reduction on a torus. In general, however, figuring out how physically to define the truncation is a highly non-trivial task.
Fortunately, for the Scherk-Schwarz reductions considered in the previous section a consistent truncation is easily defined. In fact, the ansatze for the fields already includes this truncation of the full possible field content of the theory.
The question we are interested in is how much supersymmetry is preserved by this sort of process (note we are interested in the amount of supersymmetry preserved by the lower dimensional theory as opposed to any specific vacuum of it). Supersymmetry, like any symmetry, describes the fact that if one mixes up the various fields in the theory in a certain manner then one gets back to the same action as one began with. Given the above discussion, the crucial question becomes whether this mixing up of the fields is consistent with the expansion and truncation that has been made of the fields - no other procedure in dimensional reduction changes the theory and so no other procedure can break supersymmetry. Therefore, a supersymmetry will be preserved in a dimensional reduction if this process is such that the supersymmetry transformations only mix up parts of the higher dimensional fields which were kept in the truncated action with other parts that were kept. If this is not the case then the supersymmetry can not be preserved in the lower dimensional theory.
For commonly used geometrical compactifications and truncations there is a single mathematical rule which tells us which supersymmetries are preserved by dimensional reduction on direct product spacetimes. A supersymmetry is preserved if its parameter, $`ฯต`$, can be written as a product of an external and an internal spinor, the latter being a singlet under the structure group of the spin bundle associated to the internal manifold <sup>3</sup><sup>3</sup>3Since in a geometric compactification the internal manifold is Riemannian it has a metric tensor which constitutes an $`O(d)`$ structure for a $`d`$ dimensional compact geometric space. The structure group of the frame bundle of the manifold is reduced from $`Gl(d)`$ to (a subgroup of) $`O(d)`$ in this case. Thus the frame bundle admits a principle sub-bundle which is the orthonormal frame bundle. If the manifold is spin this principle bundle has an associated spin bundle - sections of which form the spinor fields on the manifold.. In other words, the internal part of the supersymmetry parameter should be globally defined and nowhere vanishing on the internal space in order for the associated supersymmetry to be preserved.
This rule is certainly true for the standard Kaluza Klein toroidal reduction where one truncates to zero modes - in this case all supersymmetry is then preserved. We shall show shortly that this rule is also true for the kind of Scherk Schwarz reductions we considered in the previous section when the monodromy is taken to lie in a geometric subgroup of $`O(d,d)`$. It is true for truncation to the massless modes on a Calabi-Yau reduction. In more complicated cases the same rule also holds; indeed in a truncation was proposed for reductions on manifolds of $`SU(3)`$ structure with non-vanishing intrinsic torsion which was constructed to obey this rule.
In the general case where the structure group of the frame bundle fills out the whole of $`O(d)`$, there are no singlets in the decomposition of the internal spinors under the structure group of the spin bundle. However, if the structure group only fills out part of $`O(d)`$ then singlets may exist.
To illustrate this consider dimensional reduction on a Calabi-Yau threefold, which is a six dimensional manifold with $`SU(3)`$ structure. In fact Calabi-Yau manifolds also have $`SU(3)`$ holonomy but this additional constraint is not important here as it relates to the amount of supersymmetry preserved by the four dimensional Minkowski space vacuum rather than the amount associated with the theory. Since the structure group of the frame bundle is $`SU(3)`$ the structure group of the associated spin bundle is also $`SU(3)Spin(6)`$ where $`Spin(6)`$ is the double cover of $`SO(6)`$. A spinor of $`Spin(6)`$ is in the fundamental of $`SU(4)`$ and so it is easy to see, if one imagines the internal spinor as a four component column vector, that the $`SU(3)`$ structure will only leave one of the four independent components of this spinor representation invariant. More formally the $`\mathrm{๐}`$ of $`SU(4)`$ decomposes under $`SU(4)SU(3)\times U(1)`$ as $`\mathrm{๐}=(\mathrm{๐},\mathrm{๐})+(\mathrm{๐},\mathrm{๐})`$ where the first numbers in the brackets are the dimension of the relevant $`SU(3)`$ representations and the second numbers are the $`U(1)`$ charges. Thus we obtain one $`SU(3)`$ singlet under this decomposition. Therefore, standard compactification on a Calabi-Yau threefold (or indeed any six dimensional manifold of SU(3) structure where the truncation used has been appropriately defined) preserves one quarter of the supersymmetry of the higher dimensional theory.
More generally, this condition for preservation of a certain number of supersymmetries in type II reductions would involve products of G-structures (see for example ). This is because the two higher dimensional supersymmetry parameters, $`ฯต_\pm `$, can be decomposed using two different internal spinors $`\chi _\pm `$ which are associated to different G-structures. For example, if one follows this rule, continuing to examine the conditions for $`N=2`$ supersymmetry in a reduction to four dimensions, the relevant condition on the internal manifold is that it should have $`SU(3)\times SU(3)`$ structure. The two internal spinors, $`\chi _\pm `$, which we require to be present by demanding $`N=2`$ supersymmetry, each define an $`SU(3)`$ structure. Locally, these two spinors need not be parallel and so the two $`SU(3)`$ structures are different and locally define an $`SU(2)`$ structure. However, if the $`SU(2)`$ structure is globally defined then the dimensionally reduced action could instead be written as an $`N=4`$ theory in four dimensions. Therefore, to avoid obtaining a lower dimensional theory that can be written in an $`N=4`$ manner, the two spinors must be somewhere parallel. The $`SU(3)`$ structure example mentioned above is a special case of $`SU(3)\times SU(3)`$ structure where the two spinors are everywhere parallel (see for example ).
An important point to note here is that although in geometric reductions the higher dimensional supersymmetry parameters can be associated to a product of G-structures, this product is defined on a single spin bundle (and hence we can consider the intersection of the two structures). So even though the spinors $`ฯต_\pm `$ are associated to different G-structures, they are sections of the same spin bundle. This is because there is only one orthonormal frame bundle with one associated spin bundle. This is one of the features of the discussion which will be modified in the forthcoming sections when we consider T-folds.
### A The behaviour of various quantities under T-duality.
We wish to propose a modification to the rule for the amount of supersymmetry preserved by geometrical reductions, given above, to the case of non-geometric T-fold compactifications. However, before we can describe these modifications we need to know how various quantities, in particular the supersymmetry parameters, behave under T-duality transformations. <sup>4</sup><sup>4</sup>4See appendix A for some comments about local reformulations of the theories under consideration which are relevant to this and following sections..
For simplicity, in the rest of this paper we will consider the case where, in the lower dimensional theories we are interested in, supermultiplets containing fields coming from metric and $`B`$ field components with a single fibre index are truncated. In terms of the massive supergravity presented in the last section for example this would correspond to the (consistent) truncation where the supermultiplets containing the calligraphic gauge fields (i.e. $`๐_y^{0i}`$ and $`๐_\alpha ^{7i}`$) are set to zero. This truncation of the general case provides a vast simplification in the discussion that will be pursued in this section while still including many of the novel mathematical (non-) geometric structures that arise in these contexts. It would be of considerable interest to consider the more general case where these terms are not zero, however, and this will be pursued by the authors in future work.
Much of what we need to know about the behaviour of the supersymmetry parameters of type II theories under T-duality has been given in the literature . In this work the relevant transformation rules were obtained by examining how the supersymmetry variations change under T-duality. Using the well known results for how the NS-NS sector transforms, such considerations are enough to fix the transformation properties of the spinor parameters.
In this paper we are interested in transformations whose generator, $`T`$, is not simply that of some geometric transformation such as a rotation. Of the $`d(2d1)`$ generators of the group $`SO(d,d)`$, $`d^2`$ correspond to general linear coordinate transformations, and $`d(d1)/2`$ correspond to B-field shifts. As such we shall follow in only considering a subgroup of the full $`O(d,d)`$ T-duality group which contains the interesting new cases and which corresponds to symmetries of a given theory, rather than dualities which swap IIA and IIB. In fact we shall consider an $`SO(d)\times SO(d)`$ subgroup which includes an $`SO(d)`$ group which is simply the ordinary rotations. We consider this slightly larger than necessary subgroup as this choice makes the following discussion somewhat clearer. The generalisation to the $`SO(d,d)`$ case is straightforward, and will be discussed in appendix B. Our T-duality group elements will take the following form,
$`O={\displaystyle \frac{1}{2}}\left(\begin{array}{cc}\hfill S+R& \hfill SR\\ \hfill SR& \hfill S+R\end{array}\right)`$ (65)
where $`S`$ and $`R`$ are $`(d+b)\times (d+b)`$ dimensional matrices which are related to $`SO(d)`$ matrices $`๐ฎ`$ and $``$ as follows,
$`S=\left(\begin{array}{cc}\hfill ๐ฎ& \\ & \hfill 1_b\end{array}\right)R=\left(\begin{array}{cc}\hfill & \\ & \hfill 1_b\end{array}\right)`$ (70)
Here $`d`$ is the dimension of the fibre and $`b`$ is the dimension of the base. In this basis our $`O(d,d)`$ invariant metric is given by the following $`2(d+b)\times 2(d+b)`$ matrix,
$`\eta =\left(\begin{array}{cc}\hfill 0& \hfill 1\\ \hfill 1& \hfill 0\end{array}\right)`$ (73)
and so $`O^T\eta O=\eta `$ as required. The elements corresponding to ordinary rotations of the fibre are those where $`S=R`$. Note that $`O`$ is analogous to $`\mathrm{\Omega }`$ as introduced in section II B. However we have encoded the same information in higher dimensional matrices so that we can describe the relevant transformations as actions on points living in the entire compact space and not just the fibre. As such the dimension of $`O`$ corresponds to the whole fibre bundle dimension, while $`\mathrm{\Omega }`$ had dimension $`2d`$, and thus was associated only to the fibre.
We now consider the transformation rule for the metric and NS two form under $`O`$. As in the 2-dimensional case, the metric and two form on the fibre bundle (which has overall dimension $`d+b`$) can be combined into a matrix $`M`$ which takes the same form as before, i.e.
$`M=\left(\begin{array}{cc}G^1& G^1B\\ BG^1& GBG^1B\end{array}\right)`$ (76)
where now the blocks are $`(d+b)\times (d+b)`$ dimensional. Under T-duality $`MOMO^T`$. Using the explicit form for $`O`$ we find the following transformation property for the inverse metric :
$`G^1Q_{}G^1Q_{}^T=Q_+G^1Q_+^T`$ (77)
In the above we have used the following definitions of the $`(d+b)`$ dimensional matrices $`Q`$:
$`Q_{}={\displaystyle \frac{1}{2}}\left[(S+R)+(SR)(G+B)\right]`$ (78)
$`Q_+={\displaystyle \frac{1}{2}}\left[(S+R)(SR)(GB)\right]`$ (79)
These equations imply two possible different transformations for the vielbein of the compact space under the transformation (65), either
$`e_{\overline{a}}^M`$ $``$ $`\widehat{e}_{()\overline{a}}^M=Q_N^Me_{\overline{a}}^N`$ (80)
$`\text{or}e_{\overline{a}}^M`$ $``$ $`\widehat{e}_{(+)\overline{a}}^M=Q_{+N}^Me_{\overline{a}}^N.`$ (81)
Here $`M`$ is a $`(d+b)`$-dimensional spacetime index and $`\overline{a}`$ an orthonormal basis index. These two vielbeins are then related by a local Lorentz transformation as follows,
$`\widehat{e}`$ $`{}_{(+)\overline{b}}{}^{M}=\widehat{e}_{()\overline{a}}^M\mathrm{\Lambda }_{\overline{b}}^{\overline{a}}\text{where}`$ (82)
$`\mathrm{\Lambda }`$ $`=e^1Q_{}^1Q_+e`$ (83)
Physically the situation is as follows. The vielbein $`e_{\overline{a}}^M`$ can be considered as arising from either the left $`(+)`$ or right $`()`$ moving sector of the string worldsheet. Under a T-duality transformation it transforms to either $`\widehat{e}_{(+)\overline{a}}^M`$ or $`\widehat{e}_{()\overline{a}}^M`$ depending on its worldsheet origin. These two local Lorentz frames are twisted relative to each other by an amount $`\mathrm{\Lambda }`$. This generalises the statement that a single T-duality on a flat background is just parity on one of the world sheet sectors to the case of non-trivial backgrounds and more complicated transformations.
We see then from equation (80) that the two vielbeins transform differently, and both transformations are nonlinear. Thus from the point of view of discussing vielbeins on a T-fold we are left with the following picture. We define the left and right moving vielbeins separately. The transition functions seen by these vielbeins are a nonlinear realisation of $`SO(d)`$ transformations. Thus there are objects which resemble an orthonormal frame bundle. The differences to the geometric case lie in the fact that there are two โorthonormal frame bundlesโ with (in general) different transition functions - one for the left and one for the right moving sectors. In addition the transition functions on these objects act non-linearly. In the limit where we take $`S=R`$ (which corresponds to ordinary rotations) we find that $`Q_{}=Q_+=S`$ and the frame bundles reduce to the usual single orthonormal frame bundle with linearly acting transition functions.
Given these transformations for the orthonormal frames, how do the supersymmetry parameters transform under the T-duality element (65)? We denote the ten dimensional supersymmetry transformation parameters of our type II theories by $`ฯต_\pm `$. We have followed here in that the subscripts โ$`\pm `$โ denote the worldsheet sector to which the spinor is associated. We choose $`ฯต_{}`$ to have positive chirality in both type II theories which then fixes the chirality of $`ฯต_+`$ in both cases.
We define the spinor representation associated with $`\mathrm{\Lambda }`$ in the usual manner:
$`\mathrm{\Sigma }_{\text{relative}}^1\mathrm{\Gamma }^a\mathrm{\Sigma }_{\text{relative}}=\mathrm{\Lambda }_b^a\mathrm{\Gamma }^b`$ (84)
The transformation rules for the supersymmetry parameters have then been found in the paper by Hassan, , to be,
$`ฯต_{}`$ $``$ $`\mathrm{\Sigma }_{\text{total}}ฯต_{}`$ (85)
$`ฯต_+`$ $``$ $`\mathrm{\Sigma }_{\text{total}}\mathrm{\Sigma }_{\text{relative}}ฯต_+`$ (86)
where the explicit expression for $`\mathrm{\Sigma }_{\text{relative}}`$ is given by,
$`\mathrm{\Sigma }_{\text{relative}}=2^{\frac{d}{2}}\sqrt{{\displaystyle \frac{det(๐ฌ_{}+๐ฌ_+)}{det๐ฌ_{}}}}\left\{1+{\displaystyle \underset{p=1}{\overset{[d/2]}{}}}{\displaystyle \frac{(1)^p}{p!2^p}}๐^{i_1i_2}\mathrm{}๐^{i_{2p1}i_{2p}}\mathrm{\Gamma }_{i_1i_2\mathrm{}i_{2p1}i_{2p}}\right\}.`$ (87)
where $`๐ฌ_\pm `$ are the $`d\times d`$ blocks of the matrices $`Q_\pm `$, which are associated to the fibre, and the quantities $`๐^{ij}`$ are components of the following matrix:
$`๐^{ij}=\left[(1_d๐ฎ^1)^1(1_d+๐ฎ^1)+\right]_{ij}^1`$ (88)
where the $`\{ij\}`$ indices are raised by matrix inversion, and $``$ is the $`B`$-field on the $`d`$ dimensional fibre.
Now the work by Hassan is concerned solely with local considerations, and therefore the overall rotation on the spinors, $`\mathrm{\Sigma }_{\text{total}}`$, can be set to the identity using the local Lorentz symmetry. This makes the spinor $`ฯต_{}`$ invariant under the T-duality transformations. However, in our case we wish to use the spinor transformations as transition functions. Therefore, we are no longer free to make the choice $`\mathrm{\Sigma }_{\text{total}}=1`$. We must specify the total action of the transformations on the spinors up to a conjugacy class, as this is the degree to which the transition functions are defined given the presence of local Lorentz symmetry. This should be compared with the similar situation which exists in the geometric case. There $`\mathrm{\Sigma }_{\text{relative}}=1`$, and $`\mathrm{\Sigma }_{\text{total}}`$ implements the rotation on the spinors associated to the rotation of the underlying frame.
There are a number of ways in which we can isolate the overall rotation on the spinors. The structure of the โorthonormal frame bundlesโ discussed above makes the correct choice relatively clear - one should pick the overall rotation so that there is a correspondence between the action on the orthonormal frames and the action on the spinors which preserves group structure. In other words the left and right moving spin bundles should be โassociated toโ the left and right moving frame bundles. One consequence of this link is that, due to the nonlinear and $`B`$-dependent nature of the transition functions on the frame bundles, the transition functions on the spin bundles will in general depend on $`G`$ and $`B`$ as well as $`S`$ and $`R`$. There are a number of symmetries we expect the spinor transition functions to obey, which follow from the association with the orthonormal frame bundles. Firstly, in the limit of flat space with no $`B`$ field (i.e. $`G+B=11`$), $`Q_+`$ and $`Q_{}`$ are simply the linearly acting rotations $`R`$ and $`S`$ respectively. Therefore, we expect the transformations on the spinors to be simply given by the spinor representations of $`R`$ and $`S`$ in that case. Secondly, we notice that for a general background there is a symmetry $`Q_+Q_{}`$ if we take $`BB`$ and exchange $`SR`$. This symmetry exchanges the left and right moving world sheet sectors and so the two frame bundles. Therefore, we expect this same symmetry to be present in the spinor transition functions; in particular it should exchange the the transition functions for $`ฯต_+`$ and $`ฯต_{}`$.
In practical terms, how do we actually construct the overall rotation? The simplest method is the following. Consider a background that is a $`T^d`$ fibration over an $`S^1`$ base with some monodromy associated with the non-contractible loop of the base. We take the monodromy to lie in the non-trivial subgroup $`SO(d)\times SO(d)`$ that we have been discussing. As we traverse the $`S^1`$, starting at $`y=0`$ and ending at $`y=2\pi `$, the orthonormal frames $`e_\pm `$ undergo field-dependent transformations implemented by $`Q_\pm `$. Due to the monodromy in the background these frames do not in general come back to themselves at $`y=2\pi `$, but are related to the original frames by some T-duality transformation. Now consider parallel transporting the spinors $`ฯต_\pm `$ around paths where all of the coordinates are constant except for $`y`$. The left moving spinor is transported using the spin connection constructed from the left moving vielbein, $`e_+`$, and similarly for the right moving sector. We find that when we compare the two spinors with themselves at $`y=0=2\pi `$ they have undergone a rotation. This rotation is defined up to a conjugacy class due to the local Lorentz invariance on the coordinate patch. We get different rotations on the left and right moving spinors and these have precisely the properties that we require for those associated to the monodromy element we have picked - they obey all of the symmetries mentioned above (this is guaranteed because our parallel transport operators involve spin connections derived from the two vielbeins $`e_\pm `$, and the vielbeins satisfy these properties) and, in particular, give the correct relative rotation, $`\mathrm{\Sigma }_{\text{relative}}`$, given in (87), which was derived by Hassan using another method.
By carrying out this construction for an arbitrary monodromy within our $`SO(d)\times SO(d)`$ subgroup we can find all the previously unknown transition functions on the supersymmetry parameters. These can then be used to construct the spin bundles associated to a given T-fold. We emphasise that the construction we have outlined here need not be based upon the T-fold we are finally interested in. It constitutes instead merely a trick for finding transition functions which, once we have them, can be used in any appropriate T-fold context we desire.
#### Example.
To give a concrete illustration of the above discussion we will now present the details for a simple example. We shall examine the case where we have a $`T^2`$ fibre on an $`S^1`$ base with a $`U(1)\times U(1)`$ monodromy in a โnontrivialโ subgroup of $`SO(2,2,Z)`$.
First, we label the coordinates on the $`T^2`$ fibre by $`z^i`$, $`i=1,2`$, and the coordinate on the base $`S^1`$ by $`y[0,2\pi ]`$. We will also use the combined coordinate $`z^A`$, $`A=1,2,3`$, where $`z^3y`$. In these coordinates, the metric on the fibre at $`y=0`$ is given by the usual torus metric:
$`๐ข_0={\displaystyle \frac{\rho _2}{\tau _2}}\left(\begin{array}{cc}1& \tau _1\\ \tau _1& |\tau |^2\end{array}\right)`$ (91)
where $`\tau =\tau _1+i\tau _2`$ is the complex structure modulus, and $`\rho _2`$ is the volume modulus. It is convenient to combine the volume modulus into a complex field $`\rho =\rho _1+i\rho _2`$, where $`\rho _1`$ is related to the B-field at $`y=0`$ as follows,
$`_0=\left(\begin{array}{cc}0& \rho _1\\ \rho _1& 0\end{array}\right)`$ (94)
We recall that we are considering the truncation where there are no off-diagonal terms in the metric or $`B`$-field between the fibre and the base, i.e. the components which have one base and one fibre index are set to zero. Therefore, the metric on the fibre bundle at $`y=0`$ is given by
$`ds^2={\displaystyle \frac{\rho _2}{\tau _2}}\left|dz^1+\tau dz^2\right|^2+e^{2\alpha }dy^2`$ (95)
where $`\alpha `$ is a function of the lower-dimensional space. From the form of the metric above, a natural vielbein and inverse vielbein to take at $`y=0`$ are
$`e_A^{\overline{a}}=\tau _2^{1/2}\left(\begin{array}{ccc}\rho _2^{1/2}& 0& 0\\ \rho _2^{1/2}\tau _1& \rho _2^{1/2}\tau _2& 0\\ 0& 0& e^\alpha \tau _2^{1/2}\end{array}\right)e_{\overline{a}}^A=\tau _2^{1/2}\left(\begin{array}{ccc}\rho _2^{1/2}\tau _2& 0& 0\\ \rho _2^{1/2}\tau _1& \rho _2^{1/2}& 0\\ 0& 0& e^\alpha \tau _2^{1/2}\end{array}\right)`$ (102)
Now the torus fibre experiences a monodromy around the base $`S^1`$. To illustrate our method we choose this mondromy (and so the transition function action we are calculating) to take the form
$`O=\left(\begin{array}{cccccc}\mathrm{cos}(a)& 0& 0& 0& \mathrm{sin}(a)& 0\\ 0& \mathrm{cos}(a)& 0& \mathrm{sin}(a)& 0& 0\\ 0& 0& 1& 0& 0& 0\\ 0& \mathrm{sin}(a)& 0& \mathrm{cos}(a)& 0& 0\\ \mathrm{sin}(a)& 0& 0& 0& \mathrm{cos}(a)& 0\\ 0& 0& 0& 0& 0& 1\end{array}\right)`$ (109)
which corresponds to the following rotation matrices $`๐ฎ`$ and $``$ (see (65)),
$`๐ฎ=^T=\left(\begin{array}{cc}\mathrm{cos}(a)& \mathrm{sin}(a)\\ \mathrm{sin}(a)& \mathrm{cos}(a)\end{array}\right)`$ (112)
If we compare with Ref. this monodromy lies in the elliptic conjugacy class of $`SL(2)_\rho `$. In a stringy application we should take $`a=\pi `$ in order to obtain a non-trivial monodromy within $`SL(2,Z)_\rho `$. We now require a $`y`$-dependent $`O`$ which implements this monodromy as we go from $`y=0`$ to $`y=2\pi `$. The obvious choice is to replace $`aay/2\pi `$ in the above expressions. However, it is more convenient to choose the following simpler expressions for the base space dependent matrices,
$`๐ฎ(x(y))๐ฎ(x)={\displaystyle \frac{1}{\sqrt{1+x^2}}}\left(\begin{array}{cc}1& x\\ x& 1\end{array}\right),(x)=๐ฎ(x)^T`$ (115)
where $`x(y)=\mathrm{tan}(ay/2\pi )`$ and so $`x=0`$ and $`x=\mathrm{tan}(a)`$ are the endpoints of the loop in the base space (we assume $`a\pi /2`$, so that the parameterisation makes sense). Note that using $`x`$ is simply a coordinate choice for the base, and the final answers will not depend on this choice (we will see this explicitly from our expressions later on). These $`x`$-dependent matrices allow us to calculate the fibre metric and $`B`$-field at arbitrary points on the base space using (14), i.e.
$`M(x)=O(x)M_0O(x)^T`$ (116)
where $`M_0`$ is the matrix given in (76) constructed from the fibre bundle metric (95) and $`B`$-field (94). Then using all of these $`x`$-dependent quantities we can calculate the matrices $`Q_+(x)`$ and $`Q_{}(x)`$ which describe the transformations of the left and right moving frames as a function of the $`S^1`$ baseโs coordinate:
$`e_{(\pm )}^{}{}_{\overline{a}}{}^{A}(x)=Q_{\pm }^{}{}_{}{}^{A}{}_{B}{}^{}(x)e_{\overline{a}}^B`$ (117)
where $`e_{\overline{a}}^B`$ is the inverse vielbein at $`y=0`$ given in (102), and $`Q_\pm (x)`$ are constructed from $`๐ฎ(x)`$, $`(x)`$, and $`M(x)`$.
We now use the expressions for $`e_{(\pm )}(x)`$ to calculate the components of the two associated three dimensional spin connections. In fact, we do not need to determine all of the components of the connections but simply those which appear in the expressions governing parallel transport of the spinors around the relevant closed loops. Now the spinors $`ฯต_\pm `$, which arise from the left and right moving sectors on the world-sheet, split up into external 7-dimensional spinors, $`\eta _\pm `$, and internal 3-dimensional spinors, $`\chi _\pm `$. Since we are interested in the internal space, and because we have made the simplification mentioned at the beginning of section IIA, we will only need to deal with the spinors $`\chi _\pm `$ from now on. The equations for parallel transport around a path where the fibre coordinates are constant are
$$_x^{(\pm )}\chi _\pm =_x\chi _\pm +\frac{1}{4}\omega _{(\pm )}^{}{}_{x}{}^{\overline{a}\overline{b}}\mathrm{\Gamma }_{\overline{a}\overline{b}}\chi _\pm =0$$
(118)
where summation is assumed over the orthonormal indices $`\overline{a},\overline{b}`$. It can easily be shown that the expression for the spin connection components is
$`\omega _{(\pm )}^{}{}_{x}{}^{\overline{a}\overline{b}}={\displaystyle \underset{A}{}}e_{(\pm )}^{}{}_{}{}^{A[\overline{a}}_xe_{(\pm )}^{}{}_{A}{}^{\overline{b}]}`$ (119)
Using the expressions (117) for the left and right moving vielbeins as functions of $`x`$ this becomes
$`\omega _{(\pm )}^{}{}_{x}{}^{\overline{a}\overline{b}}=e_A^{[\overline{a}}\left(_x(Q_{(\pm )}^1(x))_{}^{A}{}_{B}{}^{}(Q_{(\pm )}(x))_{}^{B}{}_{C}{}^{}\right)e^{\overline{b}]C}`$ (120)
with summation understood over $`A,B,C`$. Notice that the spin connection is now written in terms of the original vielbein at $`x=0`$. Using this equation we find that the only non-vanishing spin connection components are
$$\omega _{()}^{}{}_{x}{}^{\overline{1}\overline{2}}=\omega _{(+)}^{}{}_{x}{}^{\overline{1}\overline{2}}=\frac{2x^2(x^21)|\rho |^2+8x^3\rho _1x^4+4x^2+1}{|1+x\rho |^2(4x^2|\rho |^24x(x^21)\rho _1+x^42x^2+1)}$$
(121)
The simple relation between $`\omega _+`$ and $`\omega _{}`$ is due to the relationship $`๐ฎ=^T`$ for this monodromy. Note also that this spin connection component does not depend on the complex structure modulus of the torus. This is because the monodromy is in the $`SL(2)_\rho `$ subgroup of $`SO(2,2)`$. We now have the required information to explicitly parallel transport the spinors around our paths with $`z^1=z^2=\text{ constant}`$. The resulting spinors at $`x=\mathrm{tan}(a)`$ are given by
$`\chi _\pm =\mathrm{exp}\left({\displaystyle \frac{1}{2}}{\displaystyle _0^{\mathrm{tan}a}}\omega _{(\pm )}^{}{}_{x}{}^{\overline{1}\overline{2}}\mathrm{\Gamma }_{\overline{1}\overline{2}}๐x\right)\chi _\pm ^0`$ (122)
where $`\chi _\pm ^0`$ are the spinors at $`x=0`$. Note that from the above expression it is clear, given that the spin connection contains one $`x`$ derivative, that we could have chosen any parameterisation for the closed loop and obtained the same answer. From our explicit expression for $`\omega _{(\pm )}^{}{}_{x}{}^{\overline{1}\overline{2}}`$ we obtain
$`{\displaystyle \frac{1}{2}}{\displaystyle _0^{\mathrm{tan}(a)}}\omega _{(\pm )}^{}{}_{x}{}^{\overline{1}\overline{2}}\mathrm{\Gamma }_{\overline{1}\overline{2}}๐x=`$ (123)
$`\pm {\displaystyle \frac{1}{2}}\left\{\mathrm{arctan}\left({\displaystyle \frac{\rho _1}{\rho _2}}\right)+\mathrm{arctan}\left({\displaystyle \frac{|\rho |^2\mathrm{tan}(a)+\rho _1}{\rho _2}}\right)+\mathrm{arctan}\left({\displaystyle \frac{2\rho _2\mathrm{tan}(a)}{\mathrm{tan}^2(a)12\mathrm{tan}(a)\rho _1}}\right)\right\}\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)`$ (126)
where we have chosen the following realisation for the $`\mathrm{\Gamma }`$ matrices: $`\mathrm{\Gamma }_{\overline{1}}=\sigma _1`$, $`\mathrm{\Gamma }_{\overline{2}}=\sigma _3`$, where $`\sigma _i`$ are the usual Pauli matrices. The exponential of this matrix then gives the explicit matrix which takes $`\chi _\pm ^0`$ to $`\chi _\pm ^{2\pi }`$ for this monodromy. Since we therefore now have the action of the monodromy on both kinds of spinor we have also determined the overall rotation we have been looking for. Clearly the matrix multiplying $`\chi _+^0`$ is the inverse of that multiplying $`\chi _{}^0`$. This is simply a consequence of the choice of monodromy we have made in this example. Despite this simple relation, it should be noted that formally the two transition functions associated to $`\chi _\pm `$ live in the structure groups of two different spinor bundles.
Note that if we take the flat space limit at $`y=0`$, i.e. take $`๐ข_0+_0=11`$ (which corresponds to taking $`\rho =\tau =i`$, although note that the $`\tau `$ modulus does not actually appear, so its value doesnโt make a difference here), then the factor multiplying the constant matrix in (126) becomes $`a/2`$. This is $`1/2`$ times the angle which appears in the monodromy matrix. This makes sense because in the flat space limit $`Q_{}=S`$, $`Q_+=R`$ and so the spinors transform via the spin representation of these rotation matrices. From (126) we see that another expected feature of these transformations also holds, namely that taking $`aa`$ and $`\rho _1\rho _1`$ exchanges the transformation matrices for $`\chi _{}`$ and $`\chi _+`$.
As a brief aside before we move on to discuss the preservation of supersymmetry, we note that in the 2-torus case, the matrix multiplying $`\chi _\pm ^0`$ in (122) can be written in a nicer form as the spin representation of the following matrix
$`A_\pm =\left(\begin{array}{cc}\frac{1}{\sqrt{det(๐ฌ_\pm )}}\stackrel{~}{e}^1๐ฌ_\pm \stackrel{~}{e}& 0\\ 0& 1\end{array}\right)`$ (129)
where $`๐ฌ_\pm `$ is evaluated at $`x=\mathrm{tan}(a)`$ and $`\stackrel{~}{e}`$ is the $`2\times 2`$ part of the vielbein $`e`$. It turns out that in the $`O(2,2)`$ case, these matrices are orthogonal. In other words,
$`\mathrm{\Sigma }(A_\pm )=\mathrm{exp}\left({\displaystyle \frac{1}{2}}{\displaystyle _0^{\mathrm{tan}a}}\omega _{(\pm )}^{}{}_{x}{}^{\overline{1}\overline{2}}\mathrm{\Gamma }_{\overline{1}\overline{2}}๐x\right)`$ (130)
where $`\mathrm{\Sigma }`$ denotes the spin representation, obtained in the usual way. Similar comments do not hold in the case of higher dimensional fibres however.
In any case we now set aside our example and return to the general discussion to show how to determine the amount of supersymmetry preserved by appropriately defined T-fold reductions.
### B โG-structuresโ and supersymmetry in the non-geometric case.
We are now going to describe a proposal for a simple rule for determining the amount of supersymmetry preserved by T-fold reductions of the form of a $`T^d`$ bundle over $`S^1`$ when a sensible truncation is defined in the dimensional reduction process. Our discussion generalises in an obvious way to more complicated T-fold reductions. As discussed earlier, we are interested in describing the supersymmetry preserved by the compactified theory rather than by any particular lower dimensional vacuum. Therefore, we will follow the analogue of the G-structures story that we outlined at the start of this section. In addition we repeat that we are considering the case where supermultiplets with fields originating from higher dimensional metric and $`B`$ field components with a single fibre index have been truncated and that we are not considering the โnon-minimalโ cases discussed in the appendix A.
The rule we propose for the amount of supersymmetry preserved by T-fold reduction in fact follows along fairly similar lines to the geometric case if the situation is phrased in terms of structure groups of spin bundles. As in the geometric case we decompose the spinor parameters into sums of internal and external pieces. We then consider how the internal pieces of these spinors transform under the structure groups of the relevant spin bundles (of which there are now two). We propose that the number of singlets in the decomposition of each internal spinor into representations of the relevant structure group then determines how many lower dimensional supersymmetries are associated to it in the case of an appropriately defined truncation.
However, there are crucial differences between the non-geometric and geometric cases. One obvious difference is that on T-folds the transition functions on the spinor supersymmetry parameters are generally $`B`$ and $`G`$ dependent, as we have seen explicitly in the previous section. This point deserves a little further discussion. At a first glance one might naively assume that this fact, coupled with the above comments, means that the amount of supersymmetry of the lower dimensional theory is a moduli dependent quantity! This clearly can not be the case. In fact since we are talking about the theory rather than about any particular lower dimensional vacuum (and so moduli values) the correct amount of supersymmetry is the minimum that is obtained when arbitrary values of the moduli are allowed.
We could however create a different lower dimensional theory if a consistent truncation were available where some of the lower dimensional fields were set to certain constant values (indeed we have already done something very similar with the fields we have already truncated). This could conceivably result in a truncated theory with more supersymmetry than its parent theory. Such phenomena are well known within supersymmetric compactifications and it should not come as a surprise that examples could also appear in this context.
The fact that we have two different spin bundles with different transition functions on them is major difference between the case at hand and that of geometric compactification. For example, the existence of two globally well defined, nowhere parallel internal spinors in the geometric case would imply the existence of 4 supersymmetries in the dimensionally reduced theory if the usual rule where to hold (i.e. two from each of $`ฯต_+`$ and $`ฯต_{}`$). However, in the T-fold case, the same need not necessarily be true. The two internal spinors could be sections of different spin bundles. This would lead to us obtaining one lower dimensional supersymmetry from each higher dimensional spinor parameter - giving us a total of two rather than four.
Let us examine now the examples of section II where we performed a Scherk-Schwarz reduction and truncation. Do these examples obey the rule we have proposed for the number of supersymmetries in the lower dimensional theory? In fact, this is a relatively easy question to answer.
We have, by the nature of our reduction ansatz (1), truncated the higher dimensional theory to consider only field configurations of the form
$`\psi (y)=g(y)[\psi _0]`$ (131)
Which supersymmetries are compatible with this truncation? Clearly from our earlier discussion only those supersymmetry transformations which, when applied to field configurations of this form give back another field configuration also in this form are compatible. Applying a supersymmetry transformation to a field configuration of this form we find the following,
$`(1+\delta _ฯต)\psi (y)`$ $`=`$ $`(1+\delta _ฯต)g(y)[\psi _0]`$ (132)
$`=`$ $`g(y)[(1+\delta _{\mathrm{\Sigma }^1ฯต})\psi _0]`$ (133)
Thus only if $`(1+\delta _{\mathrm{\Sigma }^1ฯต})\psi _0=\psi _0^{}`$, where $`\psi _0^{}`$ is independent of $`y`$, do we get back to a field configuration of the same form. Given $`\mathrm{\Sigma }`$โs dependence on $`y`$ in our examples this means that only supersymmetries of the form $`ฯต=\mathrm{\Sigma }ฯต_0`$, where $`ฯต_0`$ is not dependent on the internal coordinates, are compatible with the truncation. However, such spinors are not globally well defined unless $`\mathrm{\Sigma }ฯต_0=ฯต_0`$. In other cases the spinor $`ฯต`$ will be double valued at $`y=0,2\pi `$. Thus as our rule states, the supersymmetries which are preserved in Scherk-Schwarz reduction arise from internal spinors which are singlets under the monodromy (and so the structure group).
Clearly there is a lot of interesting (non)-geometrical structure in the above. In addition we have now answered the question we set out to address: we know how to calculate the amount of supersymmetry preserved by certain T-fold reductions.
#### Some Examples.
Let us illustrate our discussion by working out the amount of supersymmetry associated with the appropriate truncations (dropping the supermultiplets containing the calligraphic fields of section II) of the massive supergravity theories discussed earlier in the paper. We therefore consider the case where we have a $`T^d`$ fibration over an $`S^1`$ base. The structure group of the spin bundles is simply given by the monodromy $`e^T`$ which is introduced into the background. Therefore, in all of our cases the monodromy simply constitutes a single operator which then lies in a $`U(1)`$ subgroup of $`O(d,d)`$. One can then use the discussion of this section to determine, for any given monodromy, how this $`U(1)`$ acts on the left and right moving spinors. For the case of a $`T^2`$ fibre the $`\mathrm{๐๐}`$ of $`SO(9,1)`$ breaks up into $`(\mathrm{๐},\mathrm{๐})`$ representations of $`SO(6,1)\times SO(3)`$. Now clearly because the spinor of $`SO(3)`$ is in the fundamental of $`SU(2)`$, a $`U(1)`$ structure group will leave none of the degrees of freedom of the $`\mathrm{๐}`$ invariant. For general values of $`B`$ and $`G`$ the monodromy will result in a non-trivial structure group for both spin bundles, as can be seen from our example of the previous subsection. Therefore in general we will be left with none of the higher dimensional supersymmetries being present in the lower dimensional theory. Thus in the general case the resulting 7-dimensional theory, as presented in section II, has no supersymmetry.
However, as discussed earlier, if one takes a truncation of the lower dimensional theory then it is possible that this truncated theory could contain more supersymmetry than its parent. Let us consider the particular example presented in section III A where we chose a certain monodromy for the $`T^2`$ fibre (109) and see if this occurs in this case. By considering the scalar potential (26) associated to this monodromy, one finds that a consistent truncation can be made by setting $`\rho =i`$, as this minimises the potential for the relevant moduli. However in this case, even after such a truncation, we see from our expressions at the end of section III that the monodromy still acts on both types of spinor in a non-trivial way. Thus for this particular choice of monodromy the truncated theory also has no supersymmetry and the phenomenon of enhancement of symmetry on truncation does not occur.
For the case of a $`T^3`$ fibre, the story for the general case is slightly different. The $`\mathrm{๐๐}`$ of $`SO(9,1)`$ breaks up into $`(\mathrm{๐},\mathrm{๐})+(\mathrm{๐}^{},\mathrm{๐}^{})`$ representations of $`SO(5,1)\times SO(4)`$. Now the two Weyl spinors of $`SO(4)`$ that arise here transform under different $`SU(2)`$ subgroups and so one of these will remain invariant if we arrange matters so that the $`U(1)`$ is a subgroup of the other $`SU(2)`$. The analysis of the other $`SU(2)`$ will follow as in the case of the $`T^2`$ fibre. Therefore, in general half the supersymmetries associated to $`ฯต_+`$ and $`ฯต_{}`$ are preserved, and so half of the possible supersymmetries in the six dimensional theory are preserved, so we obtain either $`(1,1)`$ or $`(2,0)`$ supersymmetry in 6 dimensions.
As before, a consistent truncation of the lower dimensional theory could result in more supersymmetry with different numbers of supersymmetries coming from the two higher dimensional supersymmetry parameters for certain choices of monodromy.
Similar analyses can be pursued for all higher dimensional fibres. Clearly the arguments described in this section can also easily be extended to deal with more general T-fold constructions, where the base is not simply $`S^1`$.
## IV Conclusions.
In this paper we have considered the dimensional reduction of type II supergravity theories on T-folds. We began by briefly discussing the dimensional reduction of higher dimensional actions to obtain lower dimensional massive supergravity theories. We then went on to discuss some underlying (non-) geometric structures associated with these compactifications in order to calculate how much supersymmetry the lower dimensional theories possess.
We showed that associated with these spaces there are two โorthonormal frame bundlesโ, one associated with each of the left and right moving sectors of the string worldsheet. These have non-linearly acting, field dependent transition functions on them. Associated to these orthonormal frame bundles, in some sense, are two spin bundles. Due to the nature of the action of T-duality on the relevant spinors these also have field dependent transition functions. By examining the subset of the full field dependent representation of $`SO(d)`$ that these transition functions take in any one case, we were able to identify any spinors that are invariant under the structure group of the bundle of which they are sections. As in the geometric case this information then told us how much supersymmetry to expect in an appropriately defined dimensionally reduced theory. All of the structure we described reduces to the usual structure of a single orthonormal frame and spin bundle etc. when the structure group is taken to lie within a geometric subgroup. We finally illustrated our discussion by calculating the amount of supersymmetry associated with truncations of various massive supergravity theories that have been presented in the literature and in the early part of this paper. It should be noted that we have been concerned in this paper with the amount of supersymmetry associated to the compactified theory rather than any one particular lower dimensional vacuum.
There is clearly much work still to be done on the investigation of various properties of T-folds. One of the more interesting directions for future work would be to look for connections between the structure presented here and mathematical frameworks such as generalised complex geometry.
## V Acknowledgements
We would like to thank the following people for extremely useful discussions and emails, J. Figueroa-OโFarrill, S.F. Hassan, V. Jejjala and S. Ross. We would also like to especially thank S. Morris for collaboration during the early stages of this project, D. Smith and C. Hull for very many helpful discussions and comments and D. Waldram for very helpful comments on the manuscript. J.G. is funded by PPARC and E. H.-J. by EPSRC.
## Appendix A: Coset reformulations.
In this paper when we have considered type II supergravity compactified on a torus we have used the formulation of this theory with the minimum number of degrees of freedom in it. There are various other formulations of these theories. Indeed two such reformulations are given in the paper of Maharana and Schwarz . These descriptions of the theory essentially constitute introducing extra degrees of freedom into the theory as well as extra auxiliary gauge symmetries in order to remove them again. This kind of procedure might be carried out, for example, in order to linearise the action of the T-duality transformations acting on the fields. Locally these reformulations of the theory are all physically equivalent. One can recover the description of the theory used in this paper from one of these coset formulations merely by stipulating a specific gauge choice for the auxiliaries. Globally, however, the formulations can differ. The presence of the auxiliary gauge symmetries can lead to a richer global structure than is present in the minimal version of the theory. One could, for example, have Wilson lines in these gauge fields. It would then be impossible to choose a gauge to recover the minimal formulation on every coordinate patch at the same time.
In this paper we only consider the standard formulation and the supersymmetry of compactifications on T-folds which are constructed within this framework. It would clearly be of interest to study the coset reformulations of this theory in this context as well. However, this is beyond the scope of this paper. It should be pointed out that our results do apply to a subset of the configurations that can result from these other formulations. The relevant subset is that where the global structure of the auxiliary gauge fields is such that we can make the gauge choice to restore the minimal formulation simultaneously on every coordinate patch <sup>5</sup><sup>5</sup>5We would like to thank Chris Hull for pointing out the importance of these reformulations of the theory in this context..
## Appendix B: Extension to more general $`SO(d,d)`$ monodromies
In the main body of this paper we have considered monodromies which lie in a $`d(d1)/2`$-dimensional subgroup of $`SO(d,d)`$ and take the โnon-trivialโ form given in (65). However, we can consider more general elements of $`SO(d,d)`$ by combining these non-trivial twists with general coordinate transformations and B-shifts. In the case $`d=2`$ we can generate a general element of $`SO(2,2)`$ by the following product of group elements,
$`O_{\text{total}}=\mathrm{\Lambda }BO`$ (134)
where $`\mathrm{\Lambda }`$ is a coordinate transformation matrix, $`B`$ defines a B-shift and $`O`$ is a non-trivial $`SO(2)\times SO(2)`$ matrix of the form (65). Explicitly, the B-shifts take the following form,
$`B_i=\left(\begin{array}{cc}1& 0\\ b_i& 1\end{array}\right)`$
and the matrix $`\mathrm{\Lambda }`$ is given by
$`\mathrm{\Lambda }=\left(\begin{array}{cc}\lambda & 0\\ 0& (\lambda ^T)^1\end{array}\right)`$
where $`\lambda `$ is the $`2\times 2`$ matrix associated to the coordinate transformation. We now consider the action of the element $`O_{\text{total}}`$ on the spinors $`ฯต_+`$ and $`ฯต_{}`$. Firstly, we note that B-shifts have no effect on either spinor parameters, but the value of $`B`$ does enter the matrices $`\mathrm{\Sigma }^{(\pm )}`$, which are the spinor transformation matrices associated to the non-trivial element $`O`$. So the effective action on the spinors $`ฯต_\pm `$ is
$`ฯต_+`$ $``$ $`\mathrm{\Sigma }^{(+)}ฯต_+`$ (135)
$`ฯต_{}`$ $``$ $`\mathrm{\Sigma }^{()}ฯต_{}`$ (136)
where $``$ is the spin representation of the coordinate transformation $`\mathrm{\Lambda }`$ (note that the spinors will only โseeโ the rotation part of the matrix $`\mathrm{\Lambda }`$). For $`d3`$ the situation is slightly different as we can generate all elements of $`SO(d,d)`$ from B-shifts, coordinate rotations (as opposed to general coordinate transformations) and non-trivial twists. This can be seen by considering the Lie algebra elements associated to each of the transformations, and by calculating their commutators. One finds that a general element of $`SO(d,d)`$, $`d3`$, can be written as
$`O_{\text{total}}=LO_3B_3O_2B_2O_1`$ (137)
where $`L`$ is a coordinate rotation. The corresponding spinor transformations are then
$`ฯต_+`$ $``$ $`\mathrm{\Sigma }_3^{(+)}\mathrm{\Sigma }_2^{(+)}\mathrm{\Sigma }_1^{(+)}ฯต_+`$ (138)
$`ฯต_{}`$ $``$ $`\mathrm{\Sigma }_3^{()}\mathrm{\Sigma }_2^{()}\mathrm{\Sigma }_1^{()}ฯต_{}`$ (139)
Note that the fact that $`B`$ changes is important in these transformations, as a different value of $`B`$ will enter each $`\mathrm{\Sigma }_i`$. With this information we could now tackle a T-fold with monodromy in the full connected subgroup of $`SO(d,d)`$.
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# A test of tau neutrino interactions with atmospheric neutrinos and K2K
## I Introduction
Despite the remarkable successes of the Standard Model (SM), it is widely believed that there is new physics at the TeV scale, which stabilizes the Higgs mass against large radiative corrections. The search for this physics has been the goal of many key particle physics experiments of the last two decades. This search has been, and is being, carried out in two, often complementary, directions: (i) efforts involving colliding particles at highest achievable energies, and (ii) efforts involving precision measurements at low energies.
Here, we concentrate on the second possibility. There are many well-known examples of the low-energy techniques being very effective. For instance, precision measurements of atomic parity violation have resulted in an accurate determination of the interactions between the electron and $`u`$ and $`d`$ quarks mediated by the $`Z`$-boson AtomicPhysRep . Searches for exotic decays of the muon Brooks et al. (1999) and precision measurements of its anomalous dipole moment gminus2 are placing important constraints on new physics at the TeV scale. Finally, searches for proton decay SKproton are sensitive to certain types of new physics all the way up to the scale of Grand Unification.
This paper deals with another probe of this type, namely, the process of neutrino refraction in matter. More specifically, we will consider the refraction of the atmospheric neutrinos in the matter of the Earth and investigate the sensitivity of this process to new physics above the electroweak scale <sup>1</sup><sup>1</sup>1The refraction process is indeed a low-energy one, not because of the energies of the particles themselves โ which reach hundreds of GeVs for atmospheric neutrinos โ but because refraction involves the neutrino-matter forward-scattering amplitude.. The focus on neutrinos is motivated by the fact that they are the least tested particles of the Standard Model. While the charged lepton sector has been extensively probed for many exotic modes (such as $`\mu e+\gamma `$, which has an upper limit of $`1.2\times 10^{11}`$ Brooks et al. (1999) on its branching ratio), our knowledge about the neutrino sector is not nearly at the same level. While in certain classes of models it is possible to relate the properties of the two sectors, such relations do not have general character Berezhiani and Rossi (2002), making it necessary to probe the interactions of neutrinos directly.
The most direct limits on the couplings of neutrinos with matter are provided by experiments of scattering of neutrino beams on a target. Scattering tests the neutrino-matter cross section; results of this type are given by CHARM Vilain et al. (1994) and NuTeV Zeller et al. (2002) and constrain mainly the coupling of the muon neutrino to matter, as described in Sec. II.3. The interactions of the tau neutrino, however, and of the electron neutrino in some channels are still poorly restricted, being allowed at the same order as the Standard Model ones Berezhiani and Rossi (2002); Davidson et al. (2003). To this end, oscillation experiments may possess an important advantage: while it is hard to produce a beam of tau neutrinos in a laboratory, a very abundant flux of tau neutrinos and/or antineutrinos is produced by oscillation from muon or electron-flavored fluxes. It is also important that scattering results are subject to obvious degeneracies, for example in the complex phases of flavor-changing couplings, since they measure probabilities rather than amplitudes. Neutrino oscillations probe different combinations of NSI parameters with respect to accelerators and thus can help to resolve the degeneracies.
The realization that oscillation experiments are sensitive to non-standard interactions (NSI) of the neutrino is, or course, very old. In fact, already in his seminal work Wolfenstein (1978), which laid the foundations for the MSW effect Wolfenstein (1978); Mikheev and Smirnov (1986, 1985), Lincoln Wolfenstein focused on the possibility that non-standard neutrino interactions may change the flavor composition of the solar neutrino flux. The idea was further developed in Valle (1987); Guzzo et al. (1991) and many other subsequent works. Recently, the idea to use both solar and atmospheric neutrino oscillation as a way to *measure* the neutrino-matter interactions has been receiving progressively more attention Fornengo et al. (2002); Guzzo et al. (2002); Friedland:2004pp ; Guzzo et al. (2004); Gonzalez-Garcia and Maltoni (2004); Miranda et al. (2004); de Gouvea and Pena-Garay (2004); Friedland:2004ah . The underlying physical argument is the following: both solar and atmospheric neutrino fluxes have been measured over a range of energies and are fit very well by neutrino oscillations. One may hope, then, that the system is already overconstrained and the introduction of new physics would break the fit, thus opening the possibility to constrain physics beyond the Standard Model strongly.
The case of atmospheric neutrinos looks especially promising in this respect. The Super-Kamiokande (SK) experiment has collected data on the neutrino survival probability as a function of the zenith angle over five decades in neutrino energy, $`E_\nu [0.1,10^4]`$ GeV. All these data are very well fit with neutrino oscillations, with only two parameters Fukuda et al. (1998, 2000); Ashie et al. (2004). Furthermore, this result is confirmed by an independent measurement of the neutrinos from an accelerator beam at the K2K experiment Aliu et al. (2004). Thus, it is very natural to use atmospheric neutrinos as probes of NSI.
The first investigations of this type Fornengo et al. (2002); Guzzo et al. (2002); Gonzalez-Garcia and Maltoni (2004) indeed gave very strong bounds. It was shown that if the analysis is restricted to two flavors Fornengo et al. (2002); Gonzalez-Garcia and Maltoni (2004), $`\nu _\mu `$ and $`\nu _\tau `$, one can constrain NSI down to the level of a few percent of the standard weak interaction couplings. It was shown later, however, that it is essential that the analysis be done with all three flavors Friedland:2004ah : the physical arguments for reducing the problem down to two flavors used in the case of the standard oscillation analysis no longer apply, in general, when one introduces NSI on top of vacuum oscillations.
Within the framework of a three-flavor analysis it was shown Friedland:2004ah that the bounds derived in the two-flavor regime are relaxed when the $`\nu _e`$ generation is included. The details of the argument are as follows. The high-energy sample is the most sensitive to matter effects. This happens since the vacuum oscillation Hamiltonian is inversely proportional to the neutrino energy, while the matter contribution to the Hamiltonian is energy independent. It turns out, however, that in a certain region of the parameter space the oscillations of high-energy atmospheric neutrinos regain the character of vacuum oscillations. The low-energy neutrinos remain in the vacuum oscillation regime and an overall satisfactory fit to the data can be achieved. For very large NSI, the fit is eventually broken because the values of the oscillation parameters preferred by the high- and low-energy parts of the data become incompatible with each other.
Our analysis in Friedland:2004ah was essentially limited to demonstrating the above point. In this work, we present the first comprehensive study of the effects of NSI in the $`e\tau `$ sector on the oscillations of atmospheric neutrinos. Our analysis here has both a numerical and an analytical part. We scan the full three dimensional space of the (effective) NSI couplings $`ฯต_{ee},ฯต_{e\tau },ฯต_{\tau \tau }`$, and present the allowed region (marginalized over the vacuum oscillations parameters) in the space of these quantities. We also discuss several generalizations, including subdominant effects like those of the non-zero $`\theta _{13}`$ mixing angle and of the smaller, โsolarโ mass splitting. A third important aspect is that our analysis updates the previous ones by including the most recent results from the K2K experiment Aliu et al. (2004).
The text contains generalities in Sect. II, where we give a general review of atmospheric neutrinos and neutrino oscillations with NSI. In Sect. III we treat the problem of atmospheric neutrino oscillations in the presence of NSI. We describe various reductions to two-flavor oscillations (Sect. III.1) and show how these help to understand the physics of the sensitivity of atmospheric neutrinos to NSI (Sect. III.2). In Sect. IV we present a detailed numerical analysis of the problem, including the discussions of dominant (Sect. IV.2) and subdominant effects (Sect. IV.3). Our summary and conclusions follow in Sect. V.
## II Generalities
### II.1 Neutrino oscillations, masses and mixings
The results of nearly all <sup>2</sup><sup>2</sup>2A possible exception is the LSND result Athanassopoulos et al. (1995, 1996), currently being tested at MiniBOONE Ray (2004). available neutrino experiments can be explained assuming that the three known neutrinos, $`\nu _e,\nu _\mu ,\nu _\tau `$, undergo flavor oscillations. As is well known, this means that neutrinos have flavor mixing and non-zero masses.
Let us denote by $`\nu _i`$ ($`i=1,3`$) the neutrino mass eigenstates, and by $`m_i`$ their masses. In vacuum, the time evolution of these states is described by the kinetic Hamiltonian, which is given by the relativistic dispersion relation: $`E=\sqrt{p^2+m^2}`$. The phenomenon of oscillations depends on the difference of the quantum phases of the mass eigenstates; therefore, for the purpose of describing oscillations one can neglect the overall constant in the Hamiltonian. The Hamiltonian in the mass basis can be then written as follows:
$$H_{\mathrm{vac}}^{\mathrm{diag}}=\mathrm{Diag}(\mathrm{\Delta }_{}\mathrm{\Delta },\mathrm{\Delta }_{}\mathrm{\Delta },\mathrm{\Delta }_{}+\mathrm{\Delta }).$$
(1)
We use the notation $`\mathrm{\Delta }\mathrm{\Delta }m_{32}^2/(4E)`$, $`\mathrm{\Delta }_{}\mathrm{\Delta }m_{21}^2/(4E)`$, with $`E`$ being the neutrino energy and $`\mathrm{\Delta }m_{ij}^2m_i^2m_j^2`$. Here $`m_iE`$ has been assumed.
The connection to the flavor basis $`\nu _\alpha `$ is provided by the matrix $`U`$, defined as $`\nu _\alpha =U_{\alpha i}\nu _i`$. We adopt the standard parameterization for $`U`$ (see e.g. Krastev and Petcov (1988)), which contains a Dirac phase, $`\delta `$, and the three mixing angles $`\theta _{ij}`$:
$$U=\left(\begin{array}{ccc}1& 0& 0\\ 0& \mathrm{cos}\theta _{23}& \mathrm{sin}\theta _{23}\\ 0& \mathrm{sin}\theta _{23}& \mathrm{cos}\theta _{23}\end{array}\right)\left(\begin{array}{ccc}\mathrm{cos}\theta _{13}& 0& \mathrm{sin}\theta _{13}e^{i\delta }\\ 0& 1& 0\\ \mathrm{sin}\theta _{13}e^{i\delta }& 0& \mathrm{cos}\theta _{13}\end{array}\right)\left(\begin{array}{ccc}\mathrm{cos}\theta _{12}& \mathrm{sin}\theta _{12}& 0\\ \mathrm{sin}\theta _{12}& \mathrm{cos}\theta _{12}& 0\\ 0& 0& 1\end{array}\right).$$
(2)
While the mass squared splittings, $`\mathrm{\Delta }m_{ij}^2`$, determine the frequency of the oscillations, the mixing angles control their amplitudes. In vacuum, no oscillations happen if the matrix $`U`$ equals the identity (zero mixing angles) or if the three masses $`m_i`$ are equal (zero mass splittings).
In many physical situations, and in the assumption of purely standard interactions, observations happen to depend mainly on one mixing and one mass square splitting, while the other parameters give small corrections. Conventionally, $`\theta _{12}`$ and $`\mathrm{\Delta }m_{21}^2`$ are assigned to describe the oscillations of solar neutrinos, while $`\theta _{23}`$ and $`\mathrm{\Delta }m_{23}^2`$ are used for atmospheric neutrinos. The third angle, $`\theta _{13}`$, gives small effects on both solar and atmospheric neutrinos; it is hoped to be tested with future precision experiments involving neutrino beams. The sign of $`\mathrm{\Delta }m_{23}^2`$ distinguishes between the two physically different configurations of the mass spectrum: the โnormalโ neutrino mass hierarchy ($`\mathrm{\Delta }m_{23}^2>0`$) and the โinvertedโ one ($`\mathrm{\Delta }m_{23}^2<0`$).
It should be stressed that the possibility of reducing both solar and atmospheric neutrino evolution to two-state oscillations is an absolutely non-trivial result and should not be taken for granted. It relies on the measured smallness of the $`\theta _{13}`$ mixing angle, and of $`\mathrm{\Delta }m_{21}^2`$ relative to $`\mathrm{\Delta }m_{23}^2`$. Moreover, it crucially depends on the matter interactions being standard. In the presence of nonstandard interactions, while the solar neutrino analysis can still be done with two states Friedland:2004pp , the atmospheric neutrino case *requires* a full three-neutrino treatment Friedland:2004ah .
### II.2 Fluxes of atmospheric neutrinos and experimental results
Let us summarize the essential features of the fluxes of atmospheric neutrinos, and of the available data on them. We also mention tests of oscillations with other neutrino sources that are relevant to our discussion. The reader is referred to the literature for a more complete review Gaisser:1990vg ; Gonzalez-Garcia and Nir (2003); Fukugita and Yanagida ; Kajita (2004).
Atmospheric neutrinos are products of the absorption of cosmic rays in the atmosphere of the Earth. They proceed from pion and kaon decays and, thus, are produced in the muon and electron species. Complex numerical models (see e.g. Gaisser et al. (1996); Bugaev et al. (1998); Fiorentini et al. (2001); Battistoni et al. (2003); Barr et al. (2004); Honda et al. (2004)) have been developed to predict how the flux of these neutrinos develops in the atmosphere; here we summarize the main features of these models.
Both neutrinos and antineutrinos are produced in similar abundances. The energy spectrum of neutrinos detected at Super-Kamiokande spans several orders of magnitudes, from $`0.1`$ GeV to over a TeV. In the absence of neutrino oscillations, for energies higher than $`1`$ GeV, the muon (electron) neutrino flux is predicted to decrease as $`E^3`$ ($`E^{3.5}`$) (see e.g. Honda et al. (2004)). At lower energy the spectrum is made more complicated by several factors, such as geomagnetic effects and solar modulations.
At energies E
<
1
<
๐ธ1E\mathrel{\vbox{\hbox{$<$}\nointerlineskip\hbox{$\sim$}}}1 GeV most muons decay in flight before reaching the ground. Correspondingly, for these energies, one expects at ground level a ratio of fluxes in the muon and electron flavors of about 2. This ratio increases with neutrino energy, as more and more muons reach the ground without decaying. Numerical models give the muon-to-electron neutrino ratio with a $`5\%`$ accuracy. By comparison, the uncertainty on the individual fluxes is significantly larger, $`20\%`$.
After a long and extensive effort of many different collaborations, the Super-Kamiokande experiment has conclusively demonstrated that atmospheric muon neutrinos undergo oscillations Fukuda et al. (1998, 2000); Ashie et al. (2004). The indication of oscillations comes from the observation of an energy- and zenith-dependent muon neutrino and antineutrino fluxes. The different muon data samples at SK, which refer to different energy windows, show that at low energy (sub-GeV events) and intermediate energy (multi-GeV events) the muon (anti)neutrino flux is suppressed at large zenith angles, suggesting a distance-dependent effect. The distance-dependent suppression is best seen in the multi-GeV events, due to the better alignment of the direction of the detected lepton with that of the incoming neutrino. In the highest energy sample (through going muons, $`E10^110^4`$ GeV) the suppression is reduced in size for all zenith angles. It is important to notice that the electron neutrino flux is not tested in the same interval of energy: the e-like event sample reaches at most $`E10`$ GeV, beyond which absorption in the rock prevents detection.
Detailed analyses of the SK data strongly favor oscillations of $`\nu _\mu `$ into $`\nu _\tau `$, and, for purely standard interactions, give the parameters
$`|\mathrm{\Delta }m_{32}^2|`$ $``$ $`(1.73.6)10^3\mathrm{eV}^2,`$
$`\mathrm{sin}^22\theta _{23}`$ $``$ $`0.851.`$ (3)
These numbers are taken from the 99% confidence level (C.L.) contours of the recent Super-Kamiokande paper, ref. Ashie et al. (2005); they agree with the results of our analysis (see Fig. 3 later in Sect. IV.2).
Several alternative neutrino conversion scenarios have been ruled out. For example, conversion into a purely sterile neutrino ($`\nu _s`$) has been excluded, mainly due to the non-observations of the $`\nu _\mu \nu _s`$ matter effects inside the Earth Fukuda et al. (2000). A scenario in which neutrinos have no masses but oscillate in the Earth due to NSI is strongly disfavored too Lipari and Lusignoli (1999); Fogli et al. (1999); Fornengo et al. (2002). Non-oscillations mechanisms, like decoherence or neutrino decay are in strong tension with the data, especially with the presence of a minimum in the $`L/E`$ event distribution at SK Ashie et al. (2004) (here $`L`$ is the distance traveled by the neutrinos from production to detection).
Given the picture of $`\nu _\mu \nu _\tau `$ oscillations, one may expect $`\nu _\tau `$-induced events in the SK detector and, in fact, there is a hint of the presence of such events Saji .
Additional support for atmospheric neutrino oscillations comes from the MACRO Ambrosio et al. (1998, 2000, 2001) and Soudan2 Allison et al. (1997); Sanchez et al. (2003) atmospheric neutrino experiments, and โ very importantly โ from the K2K neutrino beam experiment Ahn et al. (2003); Aliu et al. (2004). K2K measures the flux and energy distribution (centered at $`E1`$ GeV ) of muon neutrinos produced by an accelerator at a distance of 250 km from the detector. The results evidence an oscillatory disappearance of $`\nu _\mu `$, with a region of oscillation parameters compatible with atmospheric results:
$`|\mathrm{\Delta }m_{32}^2|`$ $``$ $`(1.25)10^3\mathrm{eV}^2,`$
$`\mathrm{sin}^22\theta _{23}`$ $``$ $`0.251,`$ (4)
(99% C.L. interval). Small regions with larger mass squared splitting ($`few10^2\mathrm{eV}^2`$) are also allowed by K2K at 99% C.L.
In addition to $`\theta _{23}`$ and $`\mathrm{\Delta }m_{32}^2`$, the other vacuum parameters are relevant to atmospheric neutrinos, as they contribute to subdominant effects. For this reason, we briefly summarize the status of the tests of these quantities.
The parameters $`\theta _{12}`$ and $`\mathrm{\Delta }m_{21}^2`$ are measured with solar neutrinos and the KamLAND experiment Eguchi et al. (2003); Araki:2004mb . A combined analysis of their data, with standard interactions only, gives Araki:2004mb ; SNO2005 :
$`\mathrm{\Delta }m_{21}^2`$ $``$ $`(710)10^5\mathrm{eV}^2,`$
$`\mathrm{tan}^2\theta _{12}`$ $``$ $`0.40.55,`$ (5)
which corresponds to the Large Mixing Angle (LMA) solution of the solar neutrino problem. Interestingly, in presence of NSI other regions of the $`\theta _{12}`$-$`\mathrm{\Delta }m_{21}^2`$ plane become allowed Friedland:2004pp ; Guzzo et al. (2004); Miranda et al. (2004).
The mixing $`\theta _{13}`$ is strongly constrained by the non-observation results from short base-line reactor experiments. The most conservative limit compatible with the allowed interval of $`|\mathrm{\Delta }m_{32}^2|`$ is:
sin2ฮธ13
<
0.02,
<
superscript2subscript๐130.02\sin^{2}\theta_{13}\mathrel{\vbox{\hbox{$<$}\nointerlineskip\hbox{$\sim$}}}0.02~{}, (6)
as given by the CHOOZ experiment at 90% C.L. Apollonio et al. (1999, 2003), and supported by the results of Palo Verde Boehm et al. (2001).
### II.3 NSI and the three neutrino oscillation Hamiltonian in matter
When neutrinos propagate in a medium, an interaction term has to be added to the Hamiltonian to account for neutrino refraction in matter. For this, we consider the Standard Model (SM) weak interactions, and possible NSI, both flavor changing (FC) and flavor preserving (FP). We can write the NSI Lagrangian in the form of effective four-fermion terms:
$`L^{NSI}`$ $`=`$ $`2\sqrt{2}G_F(\overline{\nu }_\alpha \gamma _\rho \nu _\beta )(ฯต_{\alpha \beta }^{f\stackrel{~}{f}L}\overline{f}_L\gamma ^\rho \stackrel{~}{f}_L+ฯต_{\alpha \beta }^{f\stackrel{~}{f}R}\overline{f}_R\gamma ^\rho \stackrel{~}{f}_R)+h.c.,`$ (7)
where $`ฯต_{\alpha \beta }^{f\stackrel{~}{f}L}`$ ($`ฯต_{\alpha \beta }^{f\stackrel{~}{f}R}`$) denotes the strength of the NSI between the neutrinos $`\nu `$ of flavors $`\alpha `$ and $`\beta `$ and the left-handed (right-handed) components of the fermions $`f`$ and $`\stackrel{~}{f}`$.
The operators in Eq. (7) may arise from physics at a high energy scale, with new, heavy scalars and gauge bosons. While a theorist may prefer to see a concrete model leading to neutrino NSI, from the experimental point of view it is impractical to test any given model in isolation. The advantage of using effective low-energy operators is that they encompass the experimental effects of a variety of models, including, very importantly, those that have not been thought of yet. The four-fermion form is appropriate since propagator corrections are negligible even at the highest energies of the atmospheric neutrino spectrum, just like for the SM interactions.
The effect of standard and non-standard interactions on neutrino propagation is given by the sum over the contributions of the individual scatterers. This results in the Hamiltonian:
$$H_{\mathrm{mat}}=\sqrt{2}G_Fn_e\left(\begin{array}{ccc}1+ฯต_{ee}& ฯต_{e\mu }^{}& ฯต_{e\tau }^{}\\ ฯต_{e\mu }& ฯต_{\mu \mu }& ฯต_{\mu \tau }^{}\\ ฯต_{e\tau }& ฯต_{\mu \tau }& ฯต_{\tau \tau }\end{array}\right),$$
(8)
in the flavor basis, up to an irrelevant identity term. Here $`n_e`$ is the electron number density of the medium and the definition $`ฯต_{\alpha \beta }_{f=u,d,e}ฯต_{\alpha \beta }^fn_f/n_e`$ accounts for above mentioned sum. We use $`ฯต_{\alpha \beta }^fฯต_{\alpha \beta }^{fL}+ฯต_{\alpha \beta }^{fR}`$ and $`ฯต_{\alpha \beta }^{fP}ฯต_{\alpha \beta }^{ffP}`$, because matter effects are sensitive only to the interactions that preserve the flavor of the background fermion $`f`$ (required by coherence Friedland:2003dv ) and, furthermore, only to the vector part of that interaction.
As mentioned in the introduction, the $`ee`$, $`e\tau `$ and $`\tau \tau `$ NSI are the least constrained by direct measurements on neutrinos. Of these, we quote some results (valid at $`90\%`$ C.L.) from Davidson et al. (2003), where each NSI coupling was analyzed isolated from the others:
|ฯตฯฯuL|
<
1.4
<
subscriptsuperscriptitalic-ฯต๐ข๐ฟ๐๐1.4\displaystyle|\epsilon^{uL}_{\tau\tau}|\mathrel{\vbox{\hbox{$<$}\nointerlineskip\hbox{$\sim$}}}1.4
0.6
<
ฯตeedR
<
0.5
<
0.6subscriptsuperscriptitalic-ฯต๐๐
๐๐
<
0.5\displaystyle-0.6\mathrel{\vbox{\hbox{$<$}\nointerlineskip\hbox{$\sim$}}}\epsilon^{dR}_{ee}\mathrel{\vbox{\hbox{$<$}\nointerlineskip\hbox{$\sim$}}}0.5
|ฯตeฯdL|
<
0.5.
<
subscriptsuperscriptitalic-ฯต๐๐ฟ๐๐0.5\displaystyle|\epsilon^{dL}_{e\tau}|\mathrel{\vbox{\hbox{$<$}\nointerlineskip\hbox{$\sim$}}}0.5~{}. (9)
We also have |ฯตฮผฯeL|
<
0.1
<
subscriptsuperscriptitalic-ฯต๐๐ฟ๐๐0.1|\epsilon^{eL}_{\mu\tau}|\mathrel{\vbox{\hbox{$<$}\nointerlineskip\hbox{$\sim$}}}0.1 and |ฯตฮผฯeR|
<
0.1
<
subscriptsuperscriptitalic-ฯต๐๐
๐๐0.1|\epsilon^{eR}_{\mu\tau}|\mathrel{\vbox{\hbox{$<$}\nointerlineskip\hbox{$\sim$}}}0.1, while all the other epsilons have stronger bounds: |ฯตฮผedL|
<
8104
<
subscriptsuperscriptitalic-ฯต๐๐ฟ๐๐8superscript104|\epsilon^{dL}_{\mu e}|\mathrel{\vbox{\hbox{$<$}\nointerlineskip\hbox{$\sim$}}}8\cdot 10^{-4}, |ฯตฮผฮผeR|
<
0.03
<
subscriptsuperscriptitalic-ฯต๐๐
๐๐0.03|\epsilon^{eR}_{\mu\mu}|\mathrel{\vbox{\hbox{$<$}\nointerlineskip\hbox{$\sim$}}}0.03 Davidson et al. (2003).
We note that these bounds are much weaker than the corresponding ones in the charged lepton sector. In many models, the latter can be carried over to the neutrinos by the $`SU(2)`$ symmetry. Since the $`SU(2)`$ symmetry is violated, however, it is also possible that the NSI couplings receive $`SU(2)`$-violating contributions. For example, it was shown that certain dimension eight operators involving the Higgs field Berezhiani and Rossi (2002); Davidson et al. (2003) affect the neutrinos but not the charged leptons. Thus, it is important to seek direct, model-independent limits on the neutrino interactions, and our work, it is hoped, contributes in this direction.
Motivated by the loose limits in Eq. (9), we focus on NSI in the $`\nu _e\nu _\tau `$ sector, take $`ฯต_{e\mu }=ฯต_{\mu \tau }=ฯต_{\mu \mu }=0`$, and study in detail the oscillations of atmospheric neutrinos with the NSI described by $`ฯต_{ee},ฯต_{e\tau },ฯต_{\tau \tau }`$. Setting $`ฯต_{e\mu }`$ and $`ฯต_{\mu \mu }`$ to zero is certainly justified in view of the strong direct limits on these couplings Davidson et al. (2003). The case for $`ฯต_{\mu \tau }`$ is a bit more subtle. Arguments can be made that even for $`ฯต_{e\tau }0`$ the sensitivity of the data to $`ฯต_{\mu \tau }`$ is essentially the same as in the two neutrino analysis Fornengo et al. (2002), and therefore the bound |ฯตฮผฯ|
<
๐ช(102)|\mathrel{{\epsilon_{\mu\tau}}}|\mathrel{\vbox{\hbox{$<$}\nointerlineskip\hbox{$\sim$}}}{\cal O}(10^{-2}), found in Fornengo et al. (2002), should apply. We hope to return to this point in a future work.
## III Atmospheric neutrinos and NSI: analytical treatment
### III.1 Reduction to two neutrinos
As follows from Sect. II.1 and II.3, the oscillations of neutrinos in matter are described, in the flavor basis, by the Hamiltonian
$$H=UH_{\mathrm{vac}}^{\mathrm{diag}}U^{}+H_{mat}.$$
(10)
In general, the density $`n_e`$ varies along the neutrino trajectory and the resulting time dependence of the Hamiltonian is too complicated to allow an exact solution of the Schroedinger equation. Moreover, in presence of NSI on quarks there is a further dependence on the nucleon-to-electron ratio, and therefore on the chemical composition of the medium. In spite of this, for purely standard interactions the oscillations of atmospheric neutrinos in the Earth is well approximated by a two neutrino oscillation, $`\nu _\mu \nu _\tau `$ , on most of the energy spectrum. This is thanks to the smallness of $`\theta _{13}`$ and to the hierarchy of the mass splittings ( $`|\mathrm{\Delta }m_{21}^2/\mathrm{\Delta }m_{32}^2|1`$). With NSI the problem is intrinsically different, and requires a different approach, as illustrated below. The analysis is carried out with $`\theta _{13}=0`$ for simplicity; corrections due to $`\theta _{13}`$ will be described in sec. IV.3 and Appendix A.
Let us start by considering the low energy part of the atmospheric neutrino spectrum: $`E0.11`$ GeV. Here we have |ฮ|2GFne
>
ฮmuch-greater-thanฮ2subscript๐บ๐นsubscript๐๐
>
subscriptฮdirect-product|\Delta|\gg\sqrt{2}G_{F}n_{e}\mathrel{\vbox{\hbox{$>$}\nointerlineskip\hbox{$\sim$}}}\Delta_{\odot}, and a two neutrino reduction is rather simple: the observations can be described in terms of dominant $`\nu _\mu \nu _\tau `$ (vacuum) oscillations driven by $`\mathrm{\Delta }`$, with small corrections due to $`\nu _e\nu _\mu /\nu _\tau `$ oscillations driven by the solar scale $`\mathrm{\Delta }_{}`$ and matter effects. We will give more details on these in Sect. IV.3.
At higher energy, $`E15`$ GeV, we have $`\sqrt{2}G_Fn_e|\mathrm{\Delta }|\mathrm{\Delta }_{}`$. We can neglect the smaller mass splitting, and put $`\mathrm{\Delta }_{}=0`$. In general, however, this approximations is not enough to reduce to a two-neutrino problem, and so oscillations cannot be studied analytically. This follows from the fact that the mixing $`\theta _{23}`$ couples the $`\nu _\mu `$ and $`\nu _\tau `$ flavors, and the flavor-changing NSI term $`ฯต_{e\tau }`$ couples $`\nu _\tau `$ with $`\nu _e`$. This is an important difference with respect to the case of SM interactions (or, more generally, flavor-preserving interactions), where having $`ฯต_{e\tau }=0`$ allows to decouple the $`\nu _e`$ state and reduce to a $`\nu _\mu `$-$`\nu _\tau `$ system.
An important, nontrivial two-neutrino reduction is possible in the *highest energy limit*: $`\sqrt{2}G_Fn_e|\mathrm{\Delta }|\mathrm{\Delta }_{}`$, which is well realized for E
>
10
>
๐ธ10E\mathrel{\vbox{\hbox{$>$}\nointerlineskip\hbox{$\sim$}}}10 GeV. As mentioned in Sect. II.2, at these energies the observed signal is just $`\nu _\mu `$ and $`\overline{\nu }_\mu `$ disappearance, due to the absorption of electrons. Thus, the focus here is primarily on the conversion of $`\nu _\mu `$ and $`\overline{\nu }_\mu `$.
To see the two neutrino reduction, it convenient to introduce the eigenvalues of $`H_{mat}`$:
$`\lambda _e^{}={\displaystyle \frac{\sqrt{2}G_Fn_e}{2}}[1+ฯต_{ee}+ฯต_{\tau \tau }+\sqrt{(1+ฯต_{ee}ฯต_{\tau \tau })^2+4|ฯต_{e\tau }|^2}],`$
$`\lambda _\mu ^{}=0,`$
$`\lambda _\tau ^{}={\displaystyle \frac{\sqrt{2}G_Fn_e}{2}}[1+ฯต_{ee}+ฯต_{\tau \tau }\sqrt{(1+ฯต_{ee}ฯต_{\tau \tau })^2+4|ฯต_{e\tau }|^2}].`$ (11)
and the matter angles $`\beta `$ and $`\psi `$:
$`\mathrm{tan}2\beta `$ $`=`$ $`2|ฯต_{e\tau }|/(1+ฯต_{ee}ฯต_{\tau \tau }),`$
$`2\psi `$ $`=`$ $`Arg(ฯต_{e\tau }).`$ (12)
First, consider a situation in which both of the matter eigenvalues dominate over the vacuum terms: $`|\lambda _\tau ^{}|,|\lambda _e^{}|\mathrm{\Delta }\mathrm{\Delta }_{}`$. In this case, the mixing $`\nu _\mu `$ in the eigenstates of the Hamiltonian is suppressed by $`\mathrm{\Delta }/|\lambda _e^{}|1`$. This means the muon neutrino will not oscillate in the matter of the Earth, in conflict with the data. This case then can be excluded with confidence.
Now, let us consider a very important case when the epsilons compensate to a certain degree to give a *hierarchical scheme* of the type $`|\lambda _e^{}|\mathrm{\Delta }|\lambda _\tau ^{}|`$ or $`|\lambda _\tau ^{}|\mathrm{\Delta }|\lambda _e^{}|`$. This is the case when the mention two-state reduction is realized. Indeed, the effects of the matter interactions decouple one of the neutrino states, while the vacuum term $`\mathrm{\Delta }`$ drives the oscillations between the remaining two.
Let us illustrate this for the situation $`|\lambda _e^{}|\mathrm{\Delta }|\lambda _\tau ^{}|`$, since it is smoothly connected to the standard (no NSI) case and therefore appears more natural. Our results can be easily generalized to the second scenario.
It is convenient to take the matter eigenstates:
$`\nu _{e}^{}{}_{}{}^{}=\mathrm{cos}\beta \nu _e+\mathrm{sin}\beta e^{2i\psi }\nu _\tau ,`$
$`\nu _{\mu }^{}{}_{}{}^{}=\nu _\mu ,`$
$`\nu _{\tau }^{}{}_{}{}^{}=\mathrm{sin}\beta e^{2i\psi }\nu _e+\mathrm{cos}\beta \nu _\tau .`$ (13)
In the basis of the primed states, the Hamiltonian has the form:
$$H=\mathrm{\Delta }\left(\begin{array}{ccc}c_\beta ^2+s_\beta ^2c_{2\theta }+\lambda _e^{}/\mathrm{\Delta }& s_\beta s_{2\theta }e^{2i\psi }& c_\beta s_\beta (1+c_{2\theta })e^{2i\psi }\\ s_\beta s_{2\theta }e^{2i\psi }& c_{2\theta }& s_{2\theta }c_\beta \\ c_\beta s_\beta (1+c_{2\theta })e^{2i\psi }& s_{2\theta }c_\beta & s_\beta ^2+c_\beta ^2c_{2\theta }+\lambda _\tau ^{}/\mathrm{\Delta }\end{array}\right).$$
(14)
We see that the mixing of the state $`\nu _{e}^{}{}_{}{}^{}`$ with the other two is suppressed, being of the order $`\mathrm{\Delta }/|\lambda _e^{}|1`$. Thus, this state decouples, and the problem reduces to $`\nu _\mu `$-$`\nu _{\tau }^{}{}_{}{}^{}`$ oscillations described by the 2-3 block of the Hamiltonian (14). The unsuppressed oscillations of $`\nu _\mu `$ into a combination of $`\nu _e`$ and $`\nu _\tau `$ can account for the observed $`\nu _\mu `$ disappearance at high energy.
The two-state $`\nu _\mu `$-$`\nu _{\tau }^{}{}_{}{}^{}`$ system is quite simple to deal with and indeed we may apply to it the known formalism for two-neutrino oscillations in matter, (see e.g. Mikheyev and Smirnov (1989); Kuo and Pantaleone (1989)), with $`\lambda _\tau ^{}`$ playing the role of matter potential <sup>3</sup><sup>3</sup>3The analogy with the standard description of neutrino propagation in matter is accurate if $`\beta `$ is constant along the neutrino trajectory. This is realized for NSI on electrons, or for NSI on quarks if the neutrons to protons ratio is constant. In other cases the dependence on time (distance) of the reduced 2$`\times `$2 Hamiltonian will be more complicated due to the t-dependence of $`\beta `$.. We find the effective mixing and mass splitting in matter:
$`\mathrm{tan}2\theta _m={\displaystyle \frac{2s_{2\theta }c_\beta }{c_{2\theta }(1+c_\beta ^2)s_\beta ^2+\lambda _\tau ^{}/\mathrm{\Delta }}},`$
$`\mathrm{\Delta }_m={\displaystyle \frac{\mathrm{\Delta }}{2}}\left[\left(c_{2\theta }(1+c_\beta ^2)s_\beta ^2+{\displaystyle \frac{\lambda _\tau ^{}}{\mathrm{\Delta }}}\right)^2+4s_{2\theta }^2c_\beta ^2\right]^{\frac{1}{2}},`$ (15)
and the equation for the oscillation probability in medium with constant density:
$$P(\nu _\mu \nu _{\tau }^{}{}_{}{}^{})=\mathrm{sin}^22\theta _m\mathrm{sin}^2\left(\mathrm{\Delta }_mL\right).$$
(16)
The expressions (15) give $`\theta _m=\theta `$ and $`\mathrm{\Delta }_m=\mathrm{\Delta }`$ if $`ฯต_{e\tau }=ฯต_{\tau \tau }=0`$. Notice that along the direction $`ฯต_{e\tau }=0`$ the matter eigenstates (13) coincide with the flavor ones, and we recover the 2$`\times `$2 problem of Refs. Fornengo et al. (2002); Gonzalez-Garcia and Maltoni (2004), in which $`\nu _e`$ is decoupled, resulting in no sensitivity to $`ฯต_{ee}`$.
The properties of the probability $`P(\nu _\mu \nu _{\tau }^{}{}_{}{}^{})`$ in Eq. (16) follow those of the MSW effect, with the possibility of resonant amplification of the oscillation amplitude ($`\mathrm{sin}2\theta _m=1`$) when the condition $`c_{2\theta }(1+c_\beta ^2)s_\beta ^2+\lambda _\tau ^{}/\mathrm{\Delta }=0`$ is realized. It is worth noticing that $`\theta _m`$ and $`\mathrm{\Delta }_m`$, and thus the probability $`P(\nu _\mu \nu _{\tau }^{}{}_{}{}^{})`$, do not depend on the phase of the $`e\tau `$ NSI term, $`\psi `$.
Given its particular relevance for the analysis of atmospheric neutrinos, we now comment on the limit of small $`\lambda _\tau ^{}`$, i.e. $`|\lambda _e^{}|\mathrm{\Delta }|\lambda _\tau ^{}|`$. This corresponds to a parabola in the space of the NSI couplings for fixed $`ฯต_{ee}`$:
$$ฯต_{\tau \tau }|ฯต_{e\tau }|^2/(1+ฯต_{ee}).$$
(17)
If we take $`\lambda _\tau ^{}=0`$ in Eqs. (15) and (16), we see that here the phase of the $`\nu _\mu `$-$`\nu _{\tau }^{}{}_{}{}^{}`$ oscillations has same dependence on the product $`\mathrm{\Delta }L`$ of vacuum oscillations, while at the same time the interaction with matter is far from negligible, as it enters the probability through the angle $`\beta `$. We also notice that the conversion does not depend on the sign of $`\mathrm{\Delta }`$ (mass hierarchy). It is also independent of the overall sign of the matter Hamiltonian, therefore neutrinos and antineutrinos have the same oscillation probability.
### III.2 Expected sensitivity
What values of $`ฯต_{ee},ฯต_{e\tau },ฯต_{\tau \tau }`$ are compatible with the atmospheric and K2K neutrino data? And, how does the presence of NSI change the allowed region in the space of $`\theta _{23}`$ and $`\mathrm{\Delta }m_{32}^2`$? While to find the allowed region in the five-dimensional parameter space requires a numerical scan, a good part of the relevant features of this region can be understood from analytical arguments. Here we illustrate those. For simplicity, we first consider how the atmospheric neutrino data put constraints on NSI, and successively analyze how the K2K signal contributes to tighten those limits.
Let us consider several different regimes.
FP couplings only: $`ฯต_{e\tau }=0`$. Here we expect that practically any value of $`ฯต_{ee}`$ will be compatible with the data, due to the decoupling of $`\nu _e`$, as explained in Sect. III.1. In contrast, the data strongly constrain $`ฯต_{\tau \tau }`$. Indeed, $`ฯต_{\tau \tau }`$ influences the oscillations in the $`\nu _\mu \nu _\tau `$ sector, with possible conflict with observations. More specifically, an upper bound on $`ฯต_{\tau \tau }`$ comes from requiring that the $`\nu _\mu \nu _\tau `$ oscillation amplitude remains maximal over a large interval of energy, $`E10^110^2`$ GeV, as indicated by the data. This can only be realized if $`\theta _{23}`$ is maximal and if the matter term remains subdominant to the vacuum ones even at the highest energies<sup>4</sup><sup>4</sup>4 One could devise a resonant MSW-like solution, where $`\theta _{23}`$ is away from maximal mixing but the mixing in matter is maximal at $`EE_0`$ in one channel (neutrinos or antineutrinos, but not both due to the different sign of the matter potential) due to the cancellation of vacuum and matter terms: $`\sqrt{2}ฯต_{\tau \tau }G_Fn_e\mathrm{\Delta }m^2/(2E_0)0`$. This scenario is not viable due to the suppression of mixing in the other channel and also because it would poorly fit the sub-GeV data, which require maximal mixing.. For neutrinos going through the center of the Earth, the highest energy at which an oscillation minimum occurs in the standard case is around
$$E_02030\text{ GeV}.$$
(18)
If we require 2ฯตฯฯGFne
<
ฮm322/(2E0)subscriptitalic-ฯต๐๐2subscript๐บ๐นsubscript๐๐
<
ฮsubscriptsuperscript๐2322subscript๐ธ0\sqrt{2}\mathrel{{\epsilon_{\tau\tau}}}G_{F}n_{e}\mathrel{\vbox{\hbox{$<$}\nointerlineskip\hbox{$\sim$}}}\Delta m^{2}_{32}/(2E_{0}), and use $`\mathrm{\Delta }m_{32}^22.510^3\mathrm{eV}^2`$, we find the bound $`ฯต_{\tau \tau }0.2`$. This agrees well with numerical results Friedland:2004ah (see sec. IV). We also expect no modifications of the allowed region of $`\theta _{23}`$ and $`\mathrm{\Delta }m_{32}^2`$ with respect to the standard case.
Both FP and FC couplings: no cancellations. In the presence of flavor changing interactions, $`ฯต_{e\tau }0`$, a more general bound on the NSI follows from the analysis of Sect. III.1. Let us begin with the โgenericโ scenario in which $`|\lambda _\tau ^{}|,|\lambda _e^{}|\mathrm{\Delta }`$ for $`EE_0`$. This case is clearly excluded. Indeed, in this case the $`\nu _\mu `$ mixing is suppressed (see Sect. III.1), and thus the muon neutrino flux remains unoscillated, in clear conflict with the data. This scenario yields what could be called a โgenericโ bound on the epsilons: for example, we get |ฯตeฯ|
<
0.5|\mathrel{{\epsilon_{e\tau}}}|\mathrel{\vbox{\hbox{$<$}\nointerlineskip\hbox{$\sim$}}}0.5 when $`ฯต_{ee}0`$. The bound depends on $`ฯต_{ee}`$ in a way that will be described later.
Both FP and FC couplings: hierarchical scenario. We now analyze the โhierarchicalโ scenario, where $`|\lambda _e^{}|\mathrm{\Delta }|\lambda _\tau ^{}|`$ for $`EE_0`$. This configuration turns out to fit the data even for rather large values of $`ฯต_{e\tau }`$ and $`ฯต_{\tau \tau }`$. To understand the reason, we can first focus on the limiting case $`\lambda _\tau ^{}=0`$, corresponding to the parabolic direction in Eq. (17). As explained in Sect. III.1, in this circumstance the muon neutrino evolution allows a two-state $`\nu _\mu \nu _{\tau }^{}{}_{}{}^{}`$ reduction and the resulting disappearance of $`\nu _\mu `$ has the all features of vacuum oscillations that are known to fit the data, namely the same $`L/E`$ dependence and equal survival probabilities for neutrinos and antineutrinos. The fact that here $`\nu _\mu `$ oscillates into a combination of $`\nu _\tau `$ and $`\nu _e`$, and not into pure $`\nu _\tau `$ as in the standard scenario, is inconsequential, since there are no e-like data available at E
>
10
>
๐ธ10E\mathrel{\vbox{\hbox{$>$}\nointerlineskip\hbox{$\sim$}}}10 GeV, as mentioned in sec. II.2. Moreover, at $`EE_0`$ vacuum terms are dominant and so vacuum oscillation features are recovered in the lower energy data samples. The result is an allowed region in the space of the NSI couplings centered along the parabola (17).
Let us study the allowed region of the hierarchical case, and in particular its width and extent along the parabolic direction. The width of the region is given by the condition $`|\lambda _\tau ^{}|<\mathrm{\Delta }m_{32}^2/(2E_0)`$, or, numerically:
$$|1+ฯต_{ee}+ฯต_{\tau \tau }\sqrt{(1+ฯต_{ee}ฯต_{\tau \tau })^2+4|ฯต_{e\tau }|^2}|0.4.$$
(19)
To determine the extent of the region, more complicated considerations are necessary. As a first step, let us address the question of how the NSI change the vacuum oscillations parameters reconstructed from the data. We expect that, if the multi-GeV and through-going muon events have a significant weight in the global fit to the data, for a given set of NSI in the region (19) the allowed region in the space of $`\theta _{23}`$ and $`\mathrm{\Delta }m_{32}^2`$ changes with respect to the standard case. Indeed, if NSI are present, but not included in the data analysis, a fit of the highest energy atmospheric data, i.e. the through-going muon ones, would give $`\mathrm{\Delta }m_m^2`$ and $`\theta _m`$ instead of the corresponding vacuum quantities. If we require that the measured $`\theta _m`$ is maximal, as favored by the data, Eq. (15) gives:
$`\mathrm{cos}2\theta _{23}s_\beta ^2/(1+c_\beta ^2),`$ (20)
$`\mathrm{\Delta }m_{32}^2\mathrm{\Delta }m_m^2(1+\mathrm{cos}^2\beta )/2,`$ (21)
where $`\beta `$ is the rotation angle between the NSI eigenbasis and the flavor basis, Eq. (12). Interestingly, Eqs. (20) and (21) tell us that the vacuum mixing $`\theta _{23}`$ would not be maximal and that $`\mathrm{\Delta }m_{32}^2`$ would be larger than the measured value. The fit to the high-energy dataset can be achieved, but only *at the expense of modifying the vacuum oscillation parameters*. At the same time, the low-energy (sub-GeV) dataset still has negligible matter effect and hence is best fit by maximal $`\theta _{23}`$. Thus, for sufficiently large NSI, there will be *a tension* between the values of the oscillation parameters preferred by the low- and high-energy datasets. Avoiding this tension leads to a constraint on $`\beta `$, which in turn translates into a constraint on the epsilons.
More specifically, one should impose that the angle given by Eq. (20) is larger than the minimum value of $`\theta _{23}`$ โ let us call it $`\theta _{min}`$ โ allowed by the low energy sample. This gives:
cos2ฮฒ
>
tan2ฮธmin,
>
superscript2๐ฝsuperscript2subscript๐๐๐๐\cos^{2}\beta\mathrel{\vbox{\hbox{$>$}\nointerlineskip\hbox{$\sim$}}}\tan^{2}\theta_{min}~{}, (22)
which will cut off the potentially infinite parabola (19) down to the shape of a smile (see sec. IV). Fits to the sub-GeV data sample give $`\theta _{min}0.52`$ Gonzalez-Garcia et al. (2004). Using this value we get cos2ฮฒ
>
0.3
>
superscript2๐ฝ0.3\cos^{2}\beta\mathrel{\vbox{\hbox{$>$}\nointerlineskip\hbox{$\sim$}}}0.3; this result is relaxed if we allow $`\theta _m`$ to be slightly away from $`\pi /4`$ in the derivation of Eq. (20).
Similarly, we have to require that the right hand side of Eq. (21) does not exceed the maximal mass splitting $`\mathrm{\Delta }m_{max}^2`$ allowed by the sub-GeV data, and obtain:
$$\mathrm{cos}^2\beta \left[\frac{2\mathrm{\Delta }m_{max}^2}{\mathrm{\Delta }m_m^2}1\right]^1.$$
(23)
For $`\mathrm{\Delta }m_{max}^2=5.010^3\mathrm{eV}^2`$ Gonzalez-Garcia et al. (2004) and $`\mathrm{\Delta }m_m^22.510^3\mathrm{eV}^2`$ Eq. (23) gives cos2ฮฒ
>
0.3
>
superscript2๐ฝ0.3\cos^{2}\beta\mathrel{\vbox{\hbox{$>$}\nointerlineskip\hbox{$\sim$}}}0.3, comparable to what given by Eq. (22).
While Eq. (19) was given in our earlier work Friedland:2004ah , the results in Eqs. (22) and (23) are presented here for the first time.
Let us check if the constraints in Eqs. (22) and (23) become stronger if we combine the atmospheric neutrino data with those from K2K. For the neutrino energies used at K2K and ฯตฮฑฮฒ
<
1
<
subscriptitalic-ฯต๐ผ๐ฝ1\epsilon_{\alpha\beta}\mathrel{\vbox{\hbox{$<$}\nointerlineskip\hbox{$\sim$}}}1, matter effects are negligible, and therefore K2K measures the vacuum parameters $`\theta _{23}`$ and $`|\mathrm{\Delta }m_{32}^2|`$. From the K2K limit on $`\theta _{23}`$ we get the condition to avoid the tension between K2K and atmospheric data. The analogue of Eq. (22) for K2K turns out to be looser than that from sub-GeV atmospheric events, while the condition on the mass splitting, Eq. (23) is important: using the K2K limit, $`\mathrm{\Delta }m^24.010^3\mathrm{eV}^2`$ Aliu et al. (2004), we find cos2ฮฒ
>
0.45
>
superscript2๐ฝ0.45\cos^{2}\beta\mathrel{\vbox{\hbox{$>$}\nointerlineskip\hbox{$\sim$}}}0.45. Thus, the present K2K data constrains NSI by limiting the allowed range of $`\mathrm{\Delta }m^2`$.
If the opposite hierarchy of eigenvalues is realized, $`|\lambda _\tau ^{}|\mathrm{\Delta }|\lambda _e^{}|`$, similar considerations to those above apply. The allowed region in the space of the NSI couplings is now given by the requirement |ฮปe|
<
ฮm322/(2E0)
<
subscript๐superscript๐ฮsubscriptsuperscript๐2322subscript๐ธ0|\lambda_{e^{\prime}}|\mathrel{\vbox{\hbox{$<$}\nointerlineskip\hbox{$\sim$}}}\Delta m^{2}_{32}/(2E_{0}), or, numerically:
$$|1+ฯต_{ee}+ฯต_{\tau \tau }+\sqrt{(1+ฯต_{ee}ฯต_{\tau \tau })^2+4|ฯต_{e\tau }|^2}|0.4.$$
(24)
This condition can be satisfied for $`ฯต_{ee}1`$. It describes a region centered around the parabola (17) oriented in the negative $`ฯต_{\tau \tau }`$ direction. Physically, this means that the contribution of $`ฯต_{ee}`$ overcompensates the standard matter term and changes the sign of the $`e`$-$`e`$ entry of the matter Hamiltonian. This is somewhat extreme but still compatible with accelerator limits (Eq. (9)) and with the combination of solar and KamLAND data, provided that a certain combination of the other parameters (NSI and vacuum ones) is realized Guzzo et al. (2004); Miranda et al. (2004).
As we decrease $`ฯต_{ee}`$ from $`ฯต_{ee}0`$ down to $`ฯต_{ee}2`$ a transition from one limiting case ($`|\lambda _e^{}|\mathrm{\Delta }|\lambda _\tau ^{}|`$) to the other ($`|\lambda _\tau ^{}|\mathrm{\Delta }|\lambda _e^{}|`$) occurs. We expect this transition to be continuous, meaning that for each fixed value of $`ฯต_{ee}`$ in this interval a region in the $`ฯต_{e\tau }`$-$`ฯต_{\tau \tau }`$ plane that is compatible with the data exists and varies smoothly with $`ฯต_{ee}`$. This conclusion is justified by our earlier comment that atmospheric neutrino oscillations are insensitive to NSI in the direction $`ฯต_{e\tau }=ฯต_{\tau \tau }=0`$. Numerical results indeed confirm the existence of such a continuous transition (see Sect. IV).
## IV Atmospheric neutrinos with NSI: numerical analysis
### IV.1 Combined analysis of atmospheric and K2K data
We performed a quantitative analysis of the atmospheric neutrino data with five parameters: two โvacuumโ ones, $`(\mathrm{\Delta }m_{32}^2,\theta _{23})`$, and three NSI quantities $`(ฯต_{ee},ฯต_{e\tau },ฯต_{\tau \tau })`$, where $`ฯต_{e\tau }`$ has been treated as real for simplicity. The parameter space was scanned and a goodness-of-fit analysis was performed for each grid point.
In the analysis, we have used two types of codes. For many of the preliminary investigations we used our own code, which made several simplifying assumptions, but was designed to capture the relevant physical features of the atmospheric neutrinos in different energy ranges. For the final fits, we used a binary of the atmospheric neutrino program kindly provided to us by Michele Maltoni (SUNY, Stony Brook). This binary is essentially the same program used in our earlier paper Friedland:2004ah . It uses the complete 1489-day charged current Super-Kamiokande phase I data set Hayato (2004), including the $`e`$-like and $`\mu `$-like data samples of sub- and multi-GeV contained events (each grouped into 10 bins in zenith angle) as well as the stopping (5 angular bins) and through-going (10 angular bins) upgoing muon data events. This amounts to a total of 55 data points. For the calculation of the expected rates the code adopts the three-dimensional atmospheric neutrino fluxes given in Ref. Honda et al. (2004). The statistical analysis of the data follows the appendix of Ref. Gonzalez-Garcia and Maltoni (2004). The binary underwent extensive testing in the course of this project and the feedback was reported back to the author.
We have included the results of the K2K data analysis in our study, by adding our atmospheric $`\chi ^2`$ to the K2K $`\chi ^2`$ in the space of $`\theta _{23}`$ and $`\mathrm{\Delta }m_{32}^2`$. The latter was provided by the K2K collaboration in tabular form, and refers to the published analysis of Ref. Aliu et al. (2004) (Fig. 4 there).
It is important to point out that in our analysis we use only NSI *propagation* effects. The *detection* effects are purposefully left out. The reason for this is that the changes in the detection cross sections due to NSI do not uniquely follow from the propagation effects, but depend on additional parameters and assumptions. As a result, the detection effects can vary significantly, from large to unobservable, depending on the underlying model of the NSI.
#### IV.1.1 NSI and detection effects
Let us give a detailed argument for why possible NSI effects on the detection cross sections are not directly related to the propagation effects. A reader primarily interested in the results of our numerical analysis may wish to skip to Sect IV.2.
Consider how the Super-Kamiokande collaboration extracts NC information from their data. One method is to analyze a multi-ring dataset specifically enriched with NC events through a careful sequence of cuts Fukuda et al. (2000); phd1 . Another method is to compare the rate of single $`\pi ^0`$ events, most of which are produced in NC interactions, to the muon rate phd2 . In principle, only the collaboration can reliably model these data. Nevertheless, in what follows we set this consideration aside and ask: Can one in principle use these data to improve the bounds on NSI derived from propagation effects?
The multi-ring dataset. The main contribution to this dataset is from multi-pion production in the energy range $`0.510`$ GeV phd1 . While the exact expression for the cross-section is quite complicated (SK uses a semi-empirical formula), a qualitative estimate may be obtained by assuming deep inelastic neutrino-parton scattering. In this case, the cross-section depends on the squares of the left- and right-handed couplings, $`g_L^2`$ and $`g_R^2`$ Fukugita and Yanagida . It is immediately obvious that the result depends on how the NSI effects are distributed between the $`u`$ and $`d`$ quarks. The value of the axial coupling, $`g_A=g_Lg_R`$, is also crucial (the refraction effects fix only the vector part, $`g_V=g_L+g_R`$). For example, for the point used in Friedland:2004pp , splitting NSI evenly between $`g_L`$ and $`g_R`$ and assuming the standard value for $`g_A`$, we find an increase in the cross-section of only 20%. The effect becomes even smaller once the axial coupling is adjusted to compensate the increase given by the vector part. Hence, no rigorous constraint can be obtained from this sample.
The single $`\pi ^0`$ sample. The main contribution to the single $`\pi ^0`$ sample is from incoherent single pion production, followed by coherent single pion production phd2 . The first process is dominated by the $`\mathrm{\Delta }`$ resonance and is largely controlled by the size of the axial coupling Fukugita and Yanagida . The second one is also axial: the neutral current creates a pion that scatters on the nucleus Rein and Shegal (1999). Hence, this sample would not constrain the vector interaction that is responsible for the modified matter effect. Also, it should be mentioned that the statistics in this sample is not sufficient to even separate the active-active from active-sterile scenarios. (The latter is disfavored by only $`1.3\sigma `$, see phd2 .)
Lastly, a simple observation is that if the NSI are assigned to electrons, they have no effect on the NC event rate in the detector, since the latter is dominated by scattering on nuclei. This makes especially clear that the propagation and detection effects need not be correlated. Only if strong model-dependent assumptions are made, like in the case of a sterile neutrino, can the two be used together to exclude an oscillation scenario.
We emphasize that it is important for Super-Kamiokande and other neutrino oscillation experiments to be looking for any anomalous NC signal: an observation of such signal would imply the discovery of NSI. At the same time, as the preceding examples show, a lack of such anomalous signal would not guarantee there are no NSI propagation effects.
### IV.2 Main results: NSI, mixing and mass splitting
Our main result involves a scan of the $`(ฯต_{ee},ฯต_{e\tau },ฯต_{\tau \tau },\mathrm{\Delta }m_{32}^2,\theta _{23})`$ space for real $`ฯต_{e\tau }`$, inverted mass hierarchy ($`\mathrm{\Delta }m_{32}^2<0`$), $`\theta _{13}=0`$ and $`\mathrm{\Delta }m_{21}^2=0`$. Variations of the three latter parameters represent subdominant effects; the cases of different mass hierarchy and nonzero $`\theta _{13}`$ are treated later, in Sect. IV.3.
The scan yields a five-dimensional allowed region. Various projections of this region are described next, in Figs. 1-5.
The first result is the region allowed by the data in the space of the NSI. It was found by marginalizing the $`\chi ^2`$ function over the vacuum oscillation parameters. The resulting function, $`\chi ^2(ฯต_{ee},ฯต_{e\tau },ฯต_{\tau \tau })`$, gives a 3-dimensional allowed region, two-dimensional $`(ฯต_{e\tau },ฯต_{\tau \tau })`$ sections of which are shown in Fig. 1 (the values of $`ฯต_{ee}`$ for each section are shown in the Figure). As the region is invariant under $`ฯต_{e\tau }ฯต_{e\tau }`$, only the positive $`ฯต_{e\tau }`$ halves are shown for each section. The symmetry in the sign (phase) of $`ฯต_{e\tau }`$ is understood in terms of the oscillations probabilities, that were shown to depend on the absolute value of $`ฯต_{e\tau }`$ only (see sec. III.2).
The curves in the Figure are isocontours of $`\mathrm{\Delta }\chi ^2\chi ^2\chi _{min}^2=7.81,11.35,14.16`$ (from the inner to the outer), corresponding to $`95\%,99\%`$ and $`3\sigma `$ C.L. The parabolic flat direction in the function $`\chi ^2(ฯต_{ee},ฯต_{e\tau },ฯต_{\tau \tau })`$ predicted by Eq. (17) is clearly seen. This direction is essentially determined by the atmospheric data, and is not altered significantly by the contribution of K2K (which, however, changes the extent of the region, as will be explained). The width of the parabolic region matches well the condition on the matter eigenvalues, Eq. (19), as was shown explicitly in Fig. 1 of ref. Friedland:2004ah . The points where the region (at a given C.L.) ends along the parabola well follow isocontours of the angle $`\beta `$ (specifically, $`\mathrm{cos}^2\beta 0.30,0.35,0.42`$ for the three confidence level contours in the figure), as expected from the โcutoffโ conditions given by the low energy atmospheric sample and by the K2K results, Eqs. (22) and (23).
Along the parabola, the function $`\chi ^2(ฯต_{ee},ฯต_{e\tau },ฯต_{\tau \tau })`$ varies very slowly near the origin, and starts to increase appreciably only at $`ฯต_{e\tau }`$ of about 0.5 or so. For example, we have: $`\chi ^2(0,0,0)=148.11`$ (no NSI) and $`\chi ^2(0.73,0.35,0.07)=148.07`$. The latter happens to be the absolute minimum, $`\chi _{min}^2`$, but clearly has practically the same goodness of fit as the origin. The curves Fig. 2 show how $`\mathrm{\Delta }\chi ^2`$ varies with $`ฯต_{e\tau }`$ along the parabola (17) for three fixed values of $`ฯต_{ee}`$.
As a side comment, we notice that the agreement with Eqs. (17) and (19) becomes worse with the decrease of $`ฯต_{ee}`$. This makes sense because as $`ฯต_{ee}`$ approaches $`1`$, the two matter eigenvalues $`\lambda _e^{}`$ and $`\lambda _\tau ^{}`$ have comparable size, thus breaking the approximation of hierarchy of the eigenvalues used in our derivations (sec. III.1). Notice that the panel with $`ฯต_{ee}=1`$ shows a hint of transition from an upward to a downward parabolic region. This transition is expected to happen with the change of sign of the term $`1+ฯต_{ee}`$ in the matter Hamiltonian.
The second result is the allowed region in the space of $`\theta _{23}`$ and $`|\mathrm{\Delta }m_{31}^2|`$, obtained by marginalizing over the NSI parameters. This region is shown in Fig. 3. In the marginalization procedure we have restricted $`ฯต_{ee}`$ in the interval $`1.6ฯต_{ee}1.6`$. This serves as an example, and corresponds to the CHARM bound if the NSI are present exclusively as flavor-preserving interactions of $`\nu _e`$ on the right-handed down quark Davidson et al. (2003) <sup>5</sup><sup>5</sup>5A completely rigorous way to incorporate the CHARM bound would be to reanalyze the CHARM results with all the relevant NSI couplings in the $`e\tau `$ sector simultaneously, and make a global fit of atmospheric, K2K and accelerator data. This is beyond the scope of this work. . We have marginalized also over the sign of $`\mathrm{\Delta }m_{31}^2`$, thus including both mass hierarchies. Expectedly, with respect to the standard case, the allowed region is bigger, and extends to smaller mixing ($`\mathrm{sin}^2\theta _{23}0.2`$ instead then $`\mathrm{sin}^2\theta _{23}0.31`$) and slightly larger $`|\mathrm{\Delta }m_{31}^2|`$, in agreement with the analytic considerations (sec. III). The absolute minimum lies at $`|\mathrm{\Delta }m_{31}^2|=2.6\times 10^3`$ eV<sup>2</sup>, $`\mathrm{sin}^2\theta =0.5`$.
We find that the contours in the figure practically coincide with those obtained for normal hierarchy with $`ฯต_{ee}`$ fixed at the upper limit, $`ฯต_{ee}=1.6`$. Since this upper limit is determined by accelerator experiments, improvements of the accelerator NSI bounds would lead to a better knowledge of the oscillation parameters.
In the Table 1 we list the intervals allowed, at different confidence levels, for each vacuum parameter after marginalizing over all the other quantities. They exhibit features analogous to those of fig. 3. The results for the standard case are given, too, for comparison.
The K2K results play an important role in restricting both the region of the oscillations parameters and that of the NSI couplings. This is shown in Figs. 4 and 5. Figure 4 shows a representative section of $`\chi ^2(ฯต_{ee},ฯต_{\tau \tau },ฯต_{e\tau })`$ (the same function as in Fig. 1) along the plane $`ฯต_{ee}=0.3`$. The lower panel refers to the full atmospheric+K2K fit (like all the results in Figs. 1-3), while the upper one is obtained with atmospheric data only. The reduction due to K2K is evident: for example, the edge of the $`3\sigma `$ contour changes from $`(ฯต_{e\tau },ฯต_{\tau \tau })(2.8,6.0)`$ to $`(ฯต_{e\tau },ฯต_{\tau \tau })(2.2,3.8)`$. The effect is smaller for smaller NSI, as one can see that the $`95\%`$ C.L. contour is practically unchanged between the two panels.
Fig. 5 is a series of variations of Fig. 3: the results of the fit in the plane $`|\mathrm{\Delta }m_{31}^2|`$-$`\mathrm{sin}^2\theta _{23}`$ are shown with and without NSI, each with and without the inclusion of the K2K results. We see that with NSI the K2K data contribute to restrict the parameters, especially in the region of large $`|\mathrm{\Delta }m_{31}^2|`$ and small mixing. The $`3\sigma `$ contour of this region is moved from $`(\mathrm{sin}^2\theta _{23},|\mathrm{\Delta }m_{31}^2|)(0.15,610^3\mathrm{eV}^2)`$ to $`(\mathrm{sin}^2\theta _{23},|\mathrm{\Delta }m_{31}^2|)(0.21,410^3\mathrm{eV}^2)`$. We also observe a restriction of $`|\mathrm{\Delta }m_{31}^2|`$ from below, which is essentially the same in the cases with and without NSI.
### IV.3 Subdominant effects: mass hierarchy, $`\theta _{13}`$, three neutrino corrections
Here we generalize our results to include subdominant effects, namely the effect of the sign of the mass hierarchy, of a non-zero $`\theta _{13}`$ and of the 1-2, โsolarโ, oscillation parameters.
Fig. 6 shows the allowed region in the space of the NSI parameters for normal mass hierarchy. Analogously to the case of inverted hierarchy shown in Fig. 1, the region follows the parabola (17) and its endpoints follow isocontours of $`\beta `$. The main difference is that for normal hierarchy these isocontours are more restricted, i.e., the allowed range of NSI is smaller. Along the parabola, the $`\chi ^2`$ grows faster with the epsilons for normal hierarchy, as shown in Fig. 7. For example, if $`ฯต_{ee}=0`$, the $`3\sigma `$ C.L. contour ends at $`(ฯต_{e\tau },ฯต_{\tau \tau })(1.0,1.0)`$ for normal hierarchy, while we have $`(ฯต_{e\tau },ฯต_{\tau \tau })(1.8,2.5)`$ for inverted hierarchy. In terms of $`\mathrm{cos}^2\beta `$, the contours in Fig. 6 correspond to $`\mathrm{cos}^2\beta 0.47,0.53,0.65`$ (compared with $`\mathrm{cos}^2\beta 0.30,0.35,0.42`$ for Fig. 1).
From our analytics, we can see two sources of difference between the results for the two hierarchies. One traces back to the term $`\lambda _\tau ^{}/\mathrm{\Delta }`$ in the 3-3 entry of the Hamiltonian (14) (see also Eq. (15)). This can be as large as $`\lambda _\tau ^{}/\mathrm{\Delta }0.2`$ in the allowed region of parameters (Eq. (19)) and so is expected to contribute at the subdominant level. Secondly, corrections that depend on the sign of $`\mathrm{\Delta }`$ arise also from the small, but not zero, coupling of the state $`\nu _{e}^{}{}_{}{}^{}`$ with the other two. This depends on the 1-1 entry of the Hamiltonian (14) and therefore on the relative sign of $`\lambda _e^{}`$ and $`\mathrm{\Delta }`$. Considering that the data are dominated by neutrinos over antineutrinos (due to the larger detection cross section, see e.g. Fukugita and Yanagida ), it makes sense that a larger allowed region of NSI is obtained for inverted hierarchy, where, for neutrinos, matter and vacuum terms have the same sign and thus suppress the mixing of $`\nu _{e}^{}{}_{}{}^{}`$ more than in the normal hierarchy case.
Fig. 8 shows a generalization of Fig. 1 to non-zero $`\theta _{13}`$. Only the plane $`ฯต_{ee}=0`$ is shown for illustration. The plot was obtained for $`\theta _{13}`$ at the reactor limit, Eq. (6), with $`\theta _{13}>0`$. From the comparison of Figs. 8 and 1 it is clear that $`\theta _{13}`$ breaks the symmetry in the sign of $`ฯต_{e\tau }`$. While the parabolic direction of the region is unchanged, the position where the region ends along this direction is affected. These features can be understood analytically, as shown in Appendix A. The effect of non-zero $`\theta _{13}`$ on the extent of the allowed region in the space of $`|\mathrm{\Delta }m_{31}^2|`$-$`\mathrm{sin}^2\theta _{23}`$ is small and is not shown here.
Let us now comment on conversion effects in the lowest energy part of the atmospheric neutrino spectrum, corresponding to the sub-GeV events. Here oscillations driven by the $`\mathrm{\Delta }m_{21}^2`$ vacuum terms and matter terms occur on top of faster vacuum oscillations due to $`\mathrm{\Delta }m_{32}^2`$. The result is a deviation of the sub-GeV e-like events from the unoscillated prediction. This effect has been discussed in detail for the standard, no-NSI case Peres and Smirnov (1999, 2004); Gonzalez-Garcia et al. (2004). It was found that the ratio of the fluxes of neutrinos of muon and electron flavors depends on $`\mathrm{\Delta }m_{21}^2`$ through the probability $`P(\nu _2\nu _e)`$ of conversion of the mass eigenstate $`\nu _2`$ into $`\nu _e`$ inside the Earth. This probability is multiplied by the flux factor $`=r\mathrm{cos}^2\theta _{23}1`$, with $`r`$ being the unoscillated ratio on muon and electron neutrino fluxes. The numerical coincidence $`\mathrm{cos}^2\theta _{23}1/2`$ and $`r2`$ produces a strong suppression of $``$ and thus of the conversion effect Peres and Smirnov (1999).
The generalization of this to NSI is immediate. Given the smallness of the effect on absolute scale, with an impact on the fit to the data not larger than few per cents, we choose not to discuss it in detail. We mention two sources of enhancement of the conversion effect. The first is a possibly larger flux factor, due to a smaller $`\theta _{23}`$, see sec. IV.2. The second is a larger probability $`P(\nu _2\nu _e)`$ due to NSI. Indeed, if flavor-changing NSI are present, $`P(\nu _2\nu _e)`$ converges to a non-zero value when matter terms dominate over the vacuum ones, in contrast with the standard scenario (see e.g. Friedland:2004pp ).
Finally, it should be pointed out that there exists an important identification in the parameter space. Indeed, physical results depend only on the *relative sign* of the vacuum and matter terms of the Hamiltonian. Because of this, for vanishing $`\theta _{13}`$, the case of one hierachy maps onto the case of the other hierarchy with $`ฯต_{ee}2ฯต_{ee}`$, $`ฯต_{e\tau }ฯต_{e\tau }`$, $`ฯต_{\tau \tau }ฯต_{\tau \tau }`$. This, even though we have presented only the cases $`ฯต_{ee}1`$, our results extend to $`ฯต_{ee}<1`$. Explicitly, normal hierachy with $`ฯต_{ee}<1`$ maps to inverted hierarchy with $`ฯต_{ee}>1`$, and, likewise, inverted hierachy with $`ฯต_{ee}<1`$ maps to normal hierarchy with $`ฯต_{ee}>1`$. The cases $`ฯต_{ee}=1`$ for the two hierarchies, shown in Figs. 1 and 6, are clearly related by the transformation $`ฯต_{\tau \tau }ฯต_{\tau \tau }`$ (and also $`ฯต_{e\tau }ฯต_{e\tau }`$, but the region is symmetric with respect to this tranformation).
## V Summary and conclusions
We have explored the sensitivity of the atmospheric and K2K neutrino data to neutrino NSI in the $`e\tau `$ sector, and investigated how the presence of NSI can change the allowed region of the vacuum oscillation parameters. The results can be summarized as follows:
1. The region of the NSI parameters $`ฯต_{ee},ฯต_{\tau \tau },ฯต_{e\tau }`$ allowed by the data is essentially determined by two conditions on the neutrino-matter interaction Hamiltonian. The first condition is that one of the eigenvalues not exceed the vacuum term at high energy ($`E2030`$ GeV): |ฮปฯ|/ฮ
<
1
<
subscript๐superscript๐ฮ1|\lambda_{\tau^{\prime}}|/\Delta\mathrel{\vbox{\hbox{$<$}\nointerlineskip\hbox{$\sim$}}}1 or, numerically, |ฮปฯ|
<
51014
<
subscript๐superscript๐5superscript1014|\lambda_{\tau^{\prime}}|\mathrel{\vbox{\hbox{$<$}\nointerlineskip\hbox{$\sim$}}}5\cdot 10^{-14} eV. It explains the presence of a โflat directionโ in the $`\chi ^2`$, which extends the allowed region to large NSI couplings. This direction is a parabola on planes of constant $`ฯต_{ee}`$. Transversely to it, the $`\chi ^2`$ function grows rapidly, as can be seen from our figures. The second condition is on the mixing angle that describes the flavor composition of the eigenstates of the matter Hamiltonian: ฮฒ
<
0.3ฯ57
<
๐ฝ0.3๐similar-to-or-equalssuperscript57\beta\mathrel{\vbox{\hbox{$<$}\nointerlineskip\hbox{$\sim$}}}0.3\pi\simeq 57^{\circ}. It follows from requiring consistency between atmospheric data samples of different energy and/or between the atmospheric data and the (practically) matter-free K2K results. This condition explains the worsening of the fit along the parabolic flat direction.
In terms of the epsilons, we see that the allowed range of $`ฯต_{e\tau }`$ and $`ฯต_{\tau \tau }`$ strongly depends on the value of $`ฯต_{ee}`$, which in turn is unconstrained by atmospheric neutrinos. Both $`ฯต_{e\tau }`$ and $`ฯต_{\tau \tau }`$ are most constrained for $`ฯต_{ee}=1`$, and the constraint rapidly weakens as $`ฯต_{ee}`$ is increased. For $`ฯต_{ee}=1.5`$ and inverted mass hierarchy, values as large as $`ฯต_{\tau \tau }5`$ and $`|ฯต_{e\tau }|3.5`$ are allowed along the parabola. These would translate into very loose limits on the NSI parameters of the effective four-fermion Lagrangian. Still, such limits generally improve on existing accelerator bounds.
2. The inclusion of NSI in the analysis modifies the allowed region of the vacuum oscillation parameters, $`\mathrm{sin}^2\theta _{23}`$ and $`\mathrm{\Delta }m_{31}^2`$. As our analytical treatment shows, the fit to large NSI is achieved at the expense of changing the values of the vacuum oscillation parameters. After marginalizing over the NSI couplings and the sign of the mass squared splitting, we find that the region in the space of $`\mathrm{sin}^2\theta _{23}`$ and $`|\mathrm{\Delta }m_{31}^2|`$ is larger than that of the standard, no-NSI case. Smaller mixing and larger mass splitting are allowed. If we fit the data to one parameter at a time, and marginalize over all the others (see Table 1) we find the intervals $`\mathrm{sin}^2\theta _{23}=0.240.68`$ and $`|\mathrm{\Delta }m_{31}^2|=(1.73.9)10^3\mathrm{eV}^2`$ at $`3\sigma `$ C.L. ( to be compared to the results of the standard case: $`\mathrm{sin}^2\theta _{23}=0.320.66`$ and $`|\mathrm{\Delta }m_{31}^2|=(1.73.6)10^3\mathrm{eV}^2`$).
3. The recent K2K results play an important role in limiting NSI. This stems from the fact that for the K2K setup matter effects are negligible, and therefore K2K measures the true vacuum oscillations parameters. The K2K constraint on the oscillation parameters, particularly on $`|\mathrm{\Delta }m_{31}^2|`$, translates in a constraint on NSI, by limiting the range over which the oscillation parameters could be varied to compensate for the effects of large NSI. As an example, for $`ฯต_{ee}๐ช(10^1)`$, the addition of the K2K results restricts the region of the NSI couplings by about $`25\%`$ in the direction of $`ฯต_{\tau \tau }`$ (Fig. 4), with respect to the analysis of atmospheric neutrinos only.
4. We have studied the dependence of the results on the mass hierarchy (sign of $`\mathrm{\Delta }m_{31}^2`$). The allowed region of the NSI couplings has the same shape for both hierarchies, but, for ฯตee
>
1\mathrel{{\epsilon_{ee}}}\mathrel{\vbox{\hbox{$>$}\nointerlineskip\hbox{$\sim$}}}-1, it is more extended for the inverted hierarchy by up to a factor of $`2`$ in the direction of $`ฯต_{\tau \tau }`$. Subdominant effects due to $`\theta _{13}`$ have been analyzed. They mainly break the degeneracy in the sign (phase) of $`ฯต_{e\tau }`$. Corrections due to the smaller mass squared splitting, $`\mathrm{\Delta }m_{21}^2`$, turn out to be quite small.
Our results represent a step toward the reconstruction of the region in the space of NSI couplings that is compatible with all existing data. One of the next steps is the extension of the analysis to include the solar neutrino and KamLAND data. While it is known that the latter do not exclude large NSI couplings along the parabolic region (17) Friedland:2004pp , a detailed study has not been done before; it is presented in a companion paper of this work, soon to be completed Friedland and Lunardini .
The fact that still large NSI in the $`e\tau `$ sector are not experimentally excluded has important implications for other aspects of neutrino physics. First, it is an important motivation for neutrino experiments with man-made sources. Experiments with neutrino beams of short or intermediate base-line, such as MINOS Thomson (2005) or OPERA Autiero (2005) will have negligible matter effects, and will increase the precision of the measurement of $`\theta _{23}`$ and $`|\mathrm{\Delta }m_{31}^2|`$. This increased precision will leave even less room for NSI, or give indication of their existence, depending on whether the measured vacuum parameters are in agreement or in tension (especially if the tension is in the direction of smaller mixing and/or larger mass splitting) with the analysis of atmospheric neutrino with standard interactions only.
Neutrino beams with energy $`E110`$ GeV and long base-line (of the order of thousands of Km) like the proposed Fermilab-to-Soudan design, for example, would exhibit dramatic effects of NSI in the disappearance of muon neutrinos (or antineutrinos) as well as in the appearance channel $`\nu _\mu \nu _e`$ (or $`\overline{\nu }_\mu \overline{\nu }_e`$). In the disappearance channel NSI would produce an irregular pattern of oscillations minima and maxima, due to all three neutrinos being involved in the oscillations, in contrast with the simpler two-neutrino $`\nu _\mu \nu _\tau `$ oscillations expected in absence of NSI. The appearance channel $`\nu _\mu \nu _e`$ would be particularly characteristic in the fact that the $`\nu _e`$ component could be much larger than what allowed by standard interactions and subdominant effects (those of $`\theta _{13}`$ and of $`\mathrm{\Delta }m_{21}^2`$) and would not be suppressed at high energy โ in contrast with the case of the MSW effect with standard interactions โ as a consequence of flavor-changing NSI. In the high-energy limit, i.e. where the matter potential in the Earth dominates over vacuum terms, the amplitude of the $`\nu _\mu \nu _e`$ oscillation would be controlled by the matter mixing $`\beta `$:
$$P(\nu _\mu \nu _e)\mathrm{sin}^2\beta P(\nu _\mu \nu _\tau ^{}),$$
(25)
where $`P(\nu _\mu \nu _{\tau }^{}{}_{}{}^{})`$ is given in Eq. (16) and is unsuppressed at high energy for NSI along the parabola (17).
The presence of NSI can also alter significantly the physics of supernova neutrinos, in a way that may be tested with data from a future galactic supernova. Firstly, in the outer regions of the collapsing star, the NSI couplings may produce a richer structure of level-crossings with respect to the two MSW resonances of the standard case. New resonances may, in principle, appear depending on the chemical composition of the medium in the star and on what scatterers are responsible for the NSI. A less trivial effect will be on the evolution of trapped neutrinos inside the core of the protoneutron star. Here, NSI may affect the explosion itself, by changing the dynamics of the core and the energetics of the shock wave. This fundamental effect has been overlooked until recently Amanik et al. (2004). More work is required to understand the full implications of NSI, including possible changes in the r-processes and in the energy deposition by neutrinos in the matter of the star.
## Acknowledgments
We thank the K2K collaboration (J. Wilkes and R. Gran in particular) for useful information and for sharing with us the results of the K2K data analysis. We acknowledge the effort of M. Maltoni in the initial stage of this work, and thank him for providing the numerical executable used for our calculations. A special thank goes to the Oak Ridge National Laboratory for allowing the use of their numerical resources and to A. Mezzacappa and W. Haxton for facilitating the contact with ORNL. A.F. acknowledges support from the Department of Energy, under contract number W-7405-ENG-36. C.L. thanks the IAS of Princeton and LANL for hospitality during the preparation of this work. She also acknowledges the INT-SCiDAC grant number DE-FC02-01ER41187 for financial support.
## Appendix A Estimating corrections due to $`\theta _{13}`$
As seen in Fig. 8, the effect of nonzero $`\theta _{13}`$ is to break the $`ฯต_{e\tau }ฯต_{e\tau }`$ symmetry, while preserving the general parabolic shape of the region. These features can be understood by generalizing our analytical description of Sec. III.2. The corrections due to $`\theta _{13}`$ enter the vacuum part of the Hamiltonian, and therefore influence those predictions that depend on vacuum terms, like the mixing and mass splitting in matter, Eq. (15), and the โcutoffโ conditions (22) and (23). One expects an interplay between $`ฯต_{e\tau }`$ terms and $`\theta _{13}`$ terms, since both couple $`\nu _e`$ to $`\nu _\tau `$. This interplay has the form of an additive interference, analogous to the one involving the vacuum terms and the standard interaction in the MSW effect. This is the origin of the breaking of the symmetry in the sign (phase) of $`ฯต_{e\tau }`$.
Expanding the Hamiltonian in Eq. (14) along the parabolic direction (17) to first order in $`\mathrm{sin}2\theta _{13}`$, we find the generalizations of Eqs. (20) and (21):
$`\mathrm{cos}\theta {\displaystyle \frac{1}{\sqrt{1+c_\beta ^2}}}+{\displaystyle \frac{\mathrm{sin}2\theta _{13}s_{2\beta }\mathrm{cos}(2\psi \delta )}{4(1+c_\beta ^2)}},`$
$`\mathrm{\Delta }m^2\mathrm{\Delta }m_m^2{\displaystyle \frac{1+\mathrm{cos}^2\beta }{2}}\left[1+{\displaystyle \frac{\mathrm{sin}2\theta _{13}\mathrm{tan}\beta }{(1+\mathrm{cos}^2\beta )^{3/2}}}\right].`$ (26)
As in Eqs. (20) and (21), these expressions give the values of the vacuum parameters that correspond to $`\theta _m\pi /4`$ and a given value of the mass splitting in matter, $`\mathrm{\Delta }m_m^2`$.
The generalized form of the condition (22) is
1tan2ฮธ1sin2ฮธ13cฮธcos(2ฯฮด)
<
tanฮฒ
<
+1tan2ฮธ1sin2ฮธ13cฮธcos(2ฯฮด).
<
1superscript2๐12subscript๐13subscript๐๐2๐๐ฟ๐ฝ
<
1superscript2๐12subscript๐13subscript๐๐2๐๐ฟ\displaystyle-\sqrt{\frac{1}{\tan^{2}\theta}-1}-\sin 2\theta_{13}c_{\theta}\cos(2\psi-\delta)\mathrel{\vbox{\hbox{$<$}\nointerlineskip\hbox{$\sim$}}}\tan\beta\mathrel{\vbox{\hbox{$<$}\nointerlineskip\hbox{$\sim$}}}+\sqrt{\frac{1}{\tan^{2}\theta}-1}-\sin 2\theta_{13}c_{\theta}\cos(2\psi-\delta)~{}. (27)
Eqs. (26)-(27) are accurate to the order $`\mathrm{sin}2\theta _{13}`$, while $`๐ช(\mathrm{sin}^2\theta _{13})`$ or higher terms are not under control. We have neglected terms proportional to $`\lambda _\tau ^{}/\mathrm{\Delta }1`$ (which is accurate along the parabola). The $`\theta _{13}`$ correction to Eq. (23) is more complicated and will not be given here. We notice that the $`\theta _{13}`$ term in Eqs. (26) comes as a product with $`\mathrm{sin}\beta `$ from an expansion of terms of the type $`|A\mathrm{sin}\beta /\mathrm{\Delta }+\mathrm{sin}\theta _{13}e^{i\delta }|^2`$; this confirms the MSW-like interference mentioned above. One can also see how results depend on the *relative* phase $`2\psi \delta `$ of the matter and vacuum terms. For $`\theta _{23}`$ in the first octant, $`\psi =\delta =0`$, and $`\theta _{13}>0`$ Eq. (27) gives a stronger restriction for positive rather than negative $`ฯต_{e\tau }`$. Similarly, from (26) we see that a positive $`ฯต_{e\tau }`$ would increase the ratio $`\mathrm{\Delta }m^2/\mathrm{\Delta }m_m^2`$, thus making the tension with the low-energy sample and/or with K2K stronger. Both of these features correspond to the trend observed in Fig. 8. We have checked that there is also quantitative agreement between numerical and analytical results (Eqs. (26)-(27)) at the order of magnitude level.
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# The Vlasov limit and its fluctuations for a system of particles which interact by means of a wave field
## 1 Introduction
In recent years, significant progress has been made on the Cauchy problem of relativistic kinetic theory.<sup>1</sup><sup>1</sup>1Beside these physical Vlasov models, also the โrelativistic VlasovโPoisson equationsโ \[GlSch85\] and more recently the VlasovโNordstrรถm equations \[CaRe03, CaRe04\] have been studied. The special-relativistic VlasovโMaxwell equations \[LuVl50, Vla61\], with applications in theories of astrophysical plasma waves \[SchJ73, Jan77\], are treated with rigor in \[Hor86, GlSt86, GlSt87a, GlSt87b, GlSch88, dPLi89b, Rei90, BGP00, BGP03, KlSt03\]; the general-relativistic VlasovโEinstein equations \[Ehl71, Ehl73\], which play a rรดle in models of cosmological evolutions \[Ber88\], have been treated rigorously in \[ReRe92, Ren94, RRSch95, Rei95\]; see also \[AnTo99, Ang00\]. The relativistic VlasovโMaxwell and VlasovโEinstein equations both reduce in their strictly non-relativistic limits to corresponding VlasovโPoisson equations \[Vla38\], for which the classical Cauchy problem has been settled \[Pfa89, Sch91, Pfa92\]. Much of the special-relativistic material is reviewed in \[Gla96\], the non-relativistic results in \[Rei97\].
Progress on the microscopic foundations of all these Vlasov models has been lagging behind in comparison. Regularized VlasovโPoisson equations have been derived through a continuum limit for a family of classical $`N`$-body problems with regularized Coulomb and Newton interactions, see \[NeWi74, Neu85\] and \[BrHe77\]. In \[BrHe77\] also a law of large numbers (LLN) and a central limit theorem (CLT) for the fluctuations around this Vlasov limit are proven; see also \[Spo91, CIP91\] for further discussions. The global regularity results of \[Pfa89\] should definitely allow one to remove the regularization after the Vlasov limit of the regularized $`N`$-body dynamics has been taken and to obtain the proper (i.e. non-regularized) VlasovโPoisson equations \[Vla38\], but we are not aware of work where this has been done explicitly. In any event, while mathematically clean, physically such a derivation of the proper VlasovโPoisson equations is still far from satisfactory, for it uses the wrong order of limits, physically speaking. The regularization should actually be removed while taking the Vlasov limit for the regularized dynamical system, which likewise seems feasible with current techniques, but as far as we know has not yet been done either; however, see \[KuRe01a, KuRe01b\] for relevant work on the expected radiation-reaction corrections to VlasovโPoisson and other Vlasov models. Another open question is whether one can obtain the proper VlasovโPoisson equations directly from the classical Newtonian $`N`$-body problem for Newton or Coulomb interactions without any regularization, essentially because the classical $`N`$-body problem is still not well-controlled. For a derivation of the classical VlasovโPoisson equations from a regularized quantum mechanical $`N`$-particles model, see \[NaSe81\]; we also mention a derivation of a SchrรถdingerโPoisson model from an $`N`$-particles quantum model without regularization, see \[BEGMY02\]. While there is thus plenty of mathematical work left to be done on the microscopic foundations of the non-relativistic VlasovโPoisson equations, their status is moderately well established. The microscopic foundations of the relativistic VlasovโMaxwell and VlasovโEinstein equations by contrast seem not to have been established with rigor in any form.
To bring about a modest change in the state of affairs of the microscopic foundations of relativistic Vlasov theory, in this paper we prove a LLN and a CLT for a regularization of the following (almost) special-relativistic generalization of the VlasovโPoisson equations for a self-gravitating system,<sup>2</sup><sup>2</sup>2We use natural dimensionless units to avoid burdening the equations with irrelevant dimensional constants. Conversion of equations (1)โ(4) to the more conventional Gaussian units for a โgravitationalโ system is effected by replacing $`tNct`$, $`xNx`$, $`vv/c`$, $`pp/(mc)`$, $`\varphi \varphi /c^2`$, $`\rho 4\pi Gm\rho /(Nc^2)`$, $`f4\pi Gcm^4f/N`$; here, $`c`$ is the speed of light, $`G`$ is Newtonโs constant of universal gravitation, $`N`$ is the total number of particles in the system, and $`m`$ is the empirical mass of a particle. Note that $`\rho `$ and $`f`$ retain their normalization as probability densities on $`^3`$ and $`^6`$. comprising the continuity equation
$$\left(_t^{}+v_x^{}_x\varphi (x,t)_p\right)f(x,p,t)=0$$
(1)
on $`x,p`$ phase space, where
$$v=\frac{p}{\sqrt{1+|p|^2}}$$
(2)
is the velocity of a generic particle with momentum $`p`$ and empirical mass of unity, and the inhomogeneous wave equation
$$\mathrm{}\varphi (x,t)=\rho (x,t)$$
(3)
on $`x`$ space, where $`\mathrm{}=_t^2+_x^2`$ is the dโAlembertian,<sup>3</sup><sup>3</sup>3We write $`_x^2`$ div grad rather than $`\mathrm{\Delta }`$, for $`\mathrm{\Delta }`$ is used with a different meaning later on. and where
$$\rho (x,t)=_^3f(x,p,t)dp$$
(4)
is the normalized density of particles attributed to the space point $`x^3`$ at time $`t`$. Clearly, $`\varphi (x,t)`$ is a wave-like generalization of the Newtonian gravity potential generated by $`\rho (x,t)`$, and $`f(x,p,t)`$ in turn is the normalized density of particles attributed to the phase-space point $`(x,p)^3\times ^3`$ at time $`t`$. We remark that although a normalized density $`f(.,.,t)`$ formally satisfies the definition of a probability density function, $`f(.,.,t)`$ is more properly thought of as (an approximation to) the actual empirical phase space density of particles for an individual system.
It is to be noted that our set of equations (1), (2), (3), (4) is not meant to be taken physically seriously in itself; in particular, the model is not manifestly Lorentz-covariant (more on that in a moment). Its derivation from a microscopic model mainly serves as a simpler primer for the derivation of the special-relativistic VlasovโMaxwell equations, which we undertake in a sequel to this paper. Indeed, the model (1), (2), (3), (4) is a simple truncation of the usual set of special-relativistic VlasovโMaxwell equations for a single species of (say, positive unit charge) particles, obtained as follows:<sup>4</sup><sup>4</sup>4To obtain this โtruncated VlasovโMaxwell systemโ in the conventional Gaussian units, replace $`tNct`$, $`xNx`$, $`vv/c`$, $`pp/(mc)`$, $`\varphi e\varphi /(mc^2)`$, $`\rho 4\pi e^2\rho /(Nmc^2)`$, $`f4\pi e^2cm^2f/N`$ in our dimensionless equations (1), (2), (3), (4); here, $`c`$, $`m`$, and $`N`$ have the same meaning as for the โgravitationalโ system, while $`e`$ is the empirical unit of charge of a particle. in the VlasovโMaxwell equations, the electromagnetic fields $`E`$ and $`B`$ are expressed in terms of the electromagnetic potentials $`\varphi `$ and $`A`$ as $`B=\times A`$ and $`E=_tA\varphi `$, gauged by the LorentzโLorenz condition $`_t\varphi +A=0`$; one then purges the inhomogeneous vector wave equation for $`A`$ and all terms involving $`A`$ (or rather its derivatives) in the Lorentz force. Curiously, and contrary to what one might have naively thought, this mutilation of the VlasovโMaxwell equations does not result in a model which approximates quasi-electrodynamical behavior without magnetic fields, but in one which rather mimics some quasi-gravitational system, for in the strictly non-relativistic limit the model formally reduces to the VlasovโPoisson equations for a Newtonian gravitational system.
We remark that the replacement $`_x\varphi (x,t)_x\varphi (x,t)/\sqrt{1+|p|^2}`$ in (1) results in an essentially Lorentz-covariant model with scalar interaction field $`\varphi `$. We say โessentiallyโ because this modification of equations (1)โ(4) is still not manifestly Lorentz-covariant when $`\varphi `$ is interpreted as a Lorentz scalar field, for the right-hand side of (3) when taken face value is the time component of a Minkowski vector. However, the model becomes manifestly Lorentz-covariant when this set of equations is supplemented by the constraint $`V_^3j(x,t)dx=0`$, where
$$j(x,t)=_^3vf(x,p,t)dp$$
(5)
is the mass current vector density, and the right-hand side of (3) is interpreted as the Minkowski scalar $`M\rho Vj`$ in the center-of-mass frame, in which $`V=0`$ and $`M_^3\rho (x,t)dx=1`$. As does our truncated VlasovโMaxwell model, the Vlasov model with a factor $`1/\sqrt{1+|p|^2}`$ multiplying $`_x\varphi (x,t)`$ in (1) formally reduces to the VlasovโPoisson equations for a Newtonian gravitational system in the strictly non-relativistic limit. While this model has a number of interesting features, we will not pursue it further here because it is less relevant to the VlasovโMaxwell equations.
Ideally, we would like to prove that the kinetic equations (1)โ(4) constitute a LLN for the dynamics of an atomistic system of $`N`$ classical point particles that interact by means of a wave gravity field. The natural candidate for this atomistic system is suggested by โatomizingโ the characteristic system for (1), which reads
$`{\displaystyle \frac{\mathrm{d}q}{\mathrm{d}t}}={\displaystyle \frac{p}{\sqrt{1+|p|^2}}},`$ (6)
$`{\displaystyle \frac{\mathrm{d}p}{\mathrm{d}t}}=\varphi (q,t),`$ (7)
with $`\varphi (x,t)`$ the wave field for (1)โ(4). Thus, interpreting $`f`$ as an empirical atomic measure of $`N`$ classical point particles, having positions $`q_i^{\left(N\right)}(t)`$ and momenta $`p_i^{\left(N\right)}(t)`$ at time $`t`$, these particle motions satisfy the characteristic equations of motion, viz.
$`\dot{q}_i^{\left(N\right)}(t)={\displaystyle \frac{p_i^{\left(N\right)}(t)}{\sqrt{1+\left|p_i^{\left(N\right)}(t)\right|^2}}},`$ (8)
$`\dot{p}_i^{\left(N\right)}(t)=\varphi ^{\left(N\right)}(q_i^{\left(N\right)}(t),t),`$ (9)
for a $`\varphi ^{\left(N\right)}`$ which satisfies the inhomogeneous wave equation
$$\mathrm{}\varphi ^{\left(N\right)}(x,t)=\frac{1}{N}_{i=1}^N\delta _{q_i^{(N)}(t)}(x).$$
(10)
Unfortunately, this system of equations has only a symbolic character, at best. Since $`\varphi ^{\left(N\right)}`$ is a distributional solution of (10) which is not in $`H^1(B)`$ for any open ball $`B`$ containing the location of a point particle, each particle is surrounded by an infinite field energy which equips the particles with an infinite inertia (via Einsteinโs $`E=mc^2`$); furthermore, the right-hand side of (9) is not well-defined. Infinite self-interaction terms are encountered also if one applies the above line of reasoning in the context of the microscopic foundations of the strictly non-relativistic VlasovโPoisson equations, but in that case the self-interactions are not dynamical, and simply discarding them formally yields a locally well-defined and consistent dynamical system. In a local relativistic theory such a formal omission of the self-interaction terms is not justified because of the dynamical radiation-reaction. Hence, before any classical microscopic derivation based on point particles can be attempted for the relativistic VlasovโMaxwell, VlasovโEinstein, and for that matter also for the simpler Vlasov equations considered here, one first has to overcome the even harder conceptual problem of setting up a well-defined microscopic relativistic model with point particles. While this is being sorted out,<sup>5</sup><sup>5</sup>5For recent progress on relativistic microscopic classical electromagnetic theory, see \[Kie04\]. it is still of interest to study the obvious amelioration of the infinite self-interactions dilemma by regularizing the ill-defined point particle models \[Spo04\].
In this vein, we follow \[KSpK97, KKSp99\], who discussed a regularized version of the symbolic equations (8), (9), (10) with $`N=1`$. They โsmearโ the instantaneous location $`q(t)^3`$ of a particle with a probability density function $`\varrho (.):^3_+`$. Consistency requires that in Newtonโs equation, the gradient of $`\varphi `$ for point sources is replaced by the $`\varrho `$-average of the gradient of a $`\varphi _\varrho `$ for $`\varrho `$-averaged point sources. The global existence and uniqueness of the dynamics for the regularized microscopic model with $`N=1`$ \[KSpK97, KKSp99\] is readily extended to arbitrary $`N`$, with uniform control in $`t`$. It should be noted that the regularization just described is non-relativistic.
Interestingly, one of the caveats of the similarly regularized electromagnetic models discussed in \[KoSp00, KuSp00a, KuSp00b, KuSp00c, BaDu01\] that was pointed out in \[Kie99\] does not occur in the regularized scalar model of \[KSpK97, KKSp99\]. Namely, in contrast to what is the case in the electromagnetic models, the a-priori density function $`\varrho `$ does *not* act as a โsource or sinkโ for the conventional scalar-field angular momentum. Thus, conservation of angular momentum holds in its conventional form and does not need to be rescued through the cosmetical surgery of associating to each particle a spin variable (cf. Appendix A.3 of \[ApKi01\] for the electromagnetic models).
Our main objective in this paper then is to show that the corresponding regularization of the Vlasov model (1)โ(4) governs a LLN for the regularization of the microscopic $`N`$ particles model with wave gravity interactions (8), (9), (10). To achieve this goal we adapt the strategies of \[NeWi74, Neu85\], and \[BrHe77\] from the VlasovโPoisson to our system of equations; see \[Spo91\] for an account of Neunzertโs proof, and \[FiEl98\] for an application to a wave modes truncation of the Vlasov equations of plasma physics. The limit $`N\mathrm{}`$ not only yields a LLN for the regularized Vlasov equations, but also their well-posedness globally in time. By adapting the strategy of \[BrHe77\] we also establish a CLT for the fluctuations around the Vlasov limit. It goes without emphasis that these โadaptationsโ involve plenty of technical and conceptional innovations.
The removal of the regularization has to be addressed at a later time. We expect violations of Lorentz symmetry caused by the finite support of $`\varrho `$ to vanish when the regularization is removed, either after the Vlasov limit has been taken or along with it. Should this expectation turn out to be unfounded, it would become pointless to try to derive the relativistic VlasovโMaxwell equations along the lines developed here.
## 2 The regularized field & $`N`$-body problem
Let $`C_\mathrm{c}^{\mathrm{}}(^3)`$ denote the infinitely many times continuously differentiable functions with compact support. In the following, it is assumed that $`\varrho (.)C_\mathrm{c}^{\mathrm{}}(^3)`$. For convenience we will also demand that $`\varrho `$ is radially symmetric and decreasing. For technical reasons \[KSpK97\] a Wiener condition (positive Fourier transform) needs to be imposed on $`\varrho `$.
We introduce the abbreviation $`=_^3`$ and the convolution notations
$`(\varrho g)(x)`$ $`={\displaystyle \varrho (yx)g(y)dy},`$ (11)
$`(\varrho g)(x)`$ $`={\displaystyle \varrho (yx)_y^{}g(y)\mathrm{d}y},`$ (12)
$`(\varrho \mathrm{Id}\times g)(x)`$ $`={\displaystyle \varrho (yx)(yx)\times _y^{}g(y)\mathrm{d}y},`$ (13)
where $`g:^3`$ is any scalar function the derivative of which is in $`L^2(^3)`$.
### 2.1 The dynamical system
We begin by listing the first-order evolution equations which define the regularized microscopic dynamical model. Incidentally, the model can be viewed as a Hamiltonian system, on which we briefly comment at the end of the next subsection.
Regularizing the inhomogeneous wave equation for the microscopic wave gravity potential with point particle sources gives an inhomogeneous wave equation for the regularized wave gravity potential. Recast as a first-order system for the canonically conjugate scalar field variables<sup>6</sup><sup>6</sup>6We recall that the homogeneous Sobolev spaces $`\dot{H}^k(^d)`$ are defined as the closure of $`C_c^k(^d)`$ w.r.t. $`u_{\dot{H}^k}^2=_{|\alpha |=k}D^\alpha u_{L^2}^2`$, where $`C_c^k(^d)`$ in turn denotes the $`k`$ times classically differentiable functions with compact support, and $`\alpha `$ is a multi-index \[GiTr01\]. The reason for why we do not work with $`H^1(^3)`$ is (14): functions in $`\dot{H}^1(^3)`$ satisfying (14) are not in $`L^2(^3)`$. However, alternatively we could work with the affine Sobolev space $`\{\psi :\psi +\frac{1}{4\pi }\varrho |.|^1L^2(^3)\&\psi L^2(^3)\}`$ with seminorm $`|||\psi |||^2=\psi _{L^2}^2+\psi +\frac{1}{4\pi }\varrho |.|^1_{L^2}^2`$. $`\psi ^{\left(N\right)}(.,t)\dot{H}^1(^3)`$ and $`\varpi ^{\left(N\right)}(.,t)L^2(^3)`$ satisfying
$`\psi ^{\left(N\right)}(x,0)`$ $`=`$ $`1/4\pi |x|`$ (14)
$`\varpi ^{\left(N\right)}(x,0)`$ $`=`$ $`0`$ (15)
outside a closed ball $`B_R^3`$ which contains the initial locations of the $`N`$ particles and the supports of their regularizations, the inhomogeneous wave equation becomes
$`_t^{}\psi ^{\left(N\right)}(x,t)=\varpi ^{\left(N\right)}(x,t),`$ (16)
$`_t^{}\varpi ^{\left(N\right)}(x,t)=_x^2\psi ^{\left(N\right)}(x,t)\left(\varrho \rho _t^{\left(N\right)}\right)(x),`$ (17)
with
$$\rho _t^{\left(N\right)}(.)=\frac{1}{N}_{i=1}^N\delta _{q_i^{(N)}(t)}(.);$$
(18)
derivatives are understood in the sense of distributions. Note that for given trajectories $`tq_i^{\left(N\right)}(t),i=1,\mathrm{},N`$, we have just $`\psi ^{\left(N\right)}(x,t)=(\varrho \varphi ^{\left(N\right)}(.,t))(x)`$ with $`\varphi ^{\left(N\right)}`$ solving (10). For $`i=1,\mathrm{},N`$, the evolution equations for the $`i`$-th particleโs canonically conjugate positions $`q_i^{\left(N\right)}(t)^3`$ and momenta $`p_i^{\left(N\right)}(t)^3`$ at time $`t`$, are Einsteinโs law relating relativistic momentum to velocity,
$$\dot{q}_i^{\left(N\right)}(t)=\frac{p_i^{\left(N\right)}(t)}{\sqrt{1+|p_i^{\left(N\right)}(t)|^2}},$$
(19)
and Newtonโs law of motion,
$$\dot{p}_i^{\left(N\right)}(t)=(\varrho \psi ^{\left(N\right)}(.,t))\left(q_i^{\left(N\right)}(t)\right).$$
(20)
A complete specification at time $`t`$ of all the first-order evolutionary variables $`(q_1^{\left(N\right)}(t),p_1^{\left(N\right)}(t);\mathrm{};q_N^{\left(N\right)}(t),p_N^{\left(N\right)}(t);\psi ^{\left(N\right)}(.,t),\varpi ^{\left(N\right)}(.,t))`$ constitutes a physical state in this model. To shorten the notation, we frequently write $`z_k^{\left(N\right)}(t)`$ for the particle variables $`(q_k^{\left(N\right)}(t),p_k^{\left(N\right)}(t))`$ and $`๐ณ_t^{\left(N\right)}`$ for $`(z_1^{\left(N\right)}(t),\mathrm{},z_N^{\left(N\right)}(t))`$; furthermore $`\zeta _t^{\left(N\right)}`$ for the wave field variables $`(\psi ^{\left(N\right)}(.,t),\varpi ^{\left(N\right)}(.,t))`$, yet sometimes $`\zeta [๐ณ_0^{\left(N\right)}]`$ rather than $`\zeta _0^{\left(N\right)}`$ for the initial fields when we want to emphasize their dependence on the initial data $`๐ณ_0^{\left(N\right)}`$ rather than merely on $`N`$; finally, we frequently write $`๐ท_t^{\left(N\right)}`$ for the physical state at time $`t`$, viz.
$$๐ท_t^{\left(N\right)}:=(๐ณ_t^{\left(N\right)},\zeta _t^{\left(N\right)}).$$
(21)
The space of all possible physical states is known as the system phase space. To conveniently adapt some results of \[KSpK97\], $`\mathrm{\Gamma }^{\left(N\right)}`$ is given Hilbert space topology by taking the Hilbert space direct sum of the particle and the field Hilbert spaces,
$$\mathrm{\Gamma }^{\left(N\right)}=\underset{2N\mathrm{terms}}{\underset{}{^3\mathrm{}^3}}\dot{H}^1(^3)L^2(^3),$$
(22)
equipped with the conventional Hilbert space inner product $`.,.`$ implied by (22). The subset of $`\mathrm{\Gamma }^{\left(N\right)}`$ on which (14), (15) is satisfied is denoted $`\mathrm{\Gamma }_\mathrm{B}^{\left(N\right)}`$.
The Hilbert space topology of $`\mathrm{\Gamma }^{\left(N\right)}`$ is of course equivalent to the Banach space topology for (22) interpreted as a Banach space direct sum, but the Hilbert space topology is indeed more natural for the $`N`$-body plus field dynamics. In contrast, a Banach space topology is the natural one for the Vlasov model which we discuss in section 3.
We remark that, while $`\dot{H}^1(^3)`$ and $`L^2(^3)`$ allow quite rough fields $`\psi ^{\left(N\right)}(.,t)`$ and $`\varpi ^{\left(N\right)}(.,t)`$, any roughness would be inherited from the initial data. To have strong solutions of the wave equation in our case, we demand $`\psi (.,0)(\dot{H}^1\dot{H}^2)(^3)`$, rather than the usual $`\psi (.,0)H^2(^3)`$; cf. \[Ika00\]. Higher regularity, e.g. as for classical solutions, can also be obtained by the usual bootstrapping, if desired.
### 2.2 The conservation laws
The conventional conservation laws for mass, momentum, angular momentum, and energy are satisfied for sufficiently regular solutions of the dynamical system. To state the conservation laws, we introduce several functionals on the system phase space of generic states $`(z_1,\mathrm{},z_N,\zeta )=(๐ณ^{\left(N\right)},\zeta )=:๐ท^{\left(N\right)}\mathrm{\Gamma }^{\left(N\right)}`$.
The mass functional, for $`\rho ^{\left(N\right)}`$ given in (18) with generic $`q_i`$, is given by
$$\left(๐ท^{\left(N\right)}\right)=\varrho \rho ^{\left(N\right)}dx,$$
(23)
the momentum functional by
$$๐ซ\left(๐ท^{\left(N\right)}\right)=\frac{1}{N}\underset{i=1}{\overset{N}{}}p_i\varpi _x^{}\psi \mathrm{d}x,$$
(24)
the angular momentum functional by
$$๐ฅ\left(๐ท^{\left(N\right)}\right)=\frac{1}{N}\underset{i=1}{\overset{N}{}}q_i\times p_i(x\times _x^{}\psi )\varpi dx,$$
(25)
and the energy functional by
$$\left(๐ท^{\left(N\right)}\right)=\frac{1}{N}\underset{i=1}{\overset{N}{}}(\sqrt{1+|p_i|^2}+(\varrho \psi )\left(q_i\right))+\frac{1}{2}\left(|_x^{}\psi |^2+|\varpi |^2\right)dx.$$
(26)
We note that $`,๐ซ,`$ are well-defined on all of $`\mathrm{\Gamma }^{\left(N\right)}`$, while $`๐ฅ`$ is well-defined only on a subset of $`\mathrm{\Gamma }^{\left(N\right)}`$; in particular, $`๐ฅ`$ is well-defined on $`\mathrm{\Gamma }_\mathrm{B}^{\left(N\right)}`$.
###### Remark 2.1
The energy functional (26) furnishes the Hamiltonian for the regularized dynamical system. It is readily verified that the Hamiltonian system, $`\frac{\mathrm{d}}{\mathrm{d}t}q_i^{\left(N\right)}=/p_i^{\left(N\right)}`$ and $`\frac{\mathrm{d}}{\mathrm{d}t}p_i^{\left(N\right)}=/q_i^{\left(N\right)}`$, together with $`_t^{}\psi ^{\left(N\right)}=\delta /\delta \varpi ^{\left(N\right)}`$ and $`_t^{}\varpi ^{\left(N\right)}=\delta /\delta \psi ^{\left(N\right)}`$, coincides with the evolution equations for the wave gravity potential and the particles.
A map $`t๐ท_t^{\left(N\right)}C^1(,\mathrm{\Gamma }^{\left(N\right)})`$ satisfying our microscopic scalar wave gravity equations will be called a $`\mathrm{\Gamma }^{\left(N\right)}`$-strong solution.
###### Proposition 2.2
For any sufficiently regular (in particular, a $`\mathrm{\Gamma }^{\left(N\right)}`$-strong) solution $`t๐ท_t^{\left(N\right)}`$ of the microscopic scalar wave gravity system, we have
$`\left(๐ท_t^{\left(N\right)}\right)`$ $`=`$ $`M,`$ (27)
$`๐ซ\left(๐ท_t^{\left(N\right)}\right)`$ $`=`$ $`P,`$ (28)
$`\left(๐ท_t^{\left(N\right)}\right)`$ $`=`$ $`E,`$ (29)
with $`M`$, $`P`$, $`E`$ independent of time; in particular, $`M=1`$. If $`๐ท_0^{\left(N\right)}\mathrm{\Gamma }_\mathrm{B}^{\left(N\right)}`$, then also
$$๐ฅ\left(๐ท_t^{\left(N\right)}\right)=J,$$
(30)
with $`J`$ independent of time.
Proof of Proposition 2.2. Proposition 2.2 is proved in the appendix as a special case of the conservation laws in Theorems 3.2 and 3.3 of subsubsection 3.2.1. Q.E.D.
###### Remark 2.3
One may contemplate attaching also an Euler spin variable $`s_i^{\left(N\right)}(t)^3`$ at time $`t`$ to the $`i`$-th particle, the (non-relativistic) evolution equations for $`s_i^{\left(N\right)}(t)`$ being just Eulerโs equations for a degenerate gyroscope, viz.
$$\dot{s}_i^{\left(N\right)}(t)=(\varrho \mathrm{Id}\times \psi ^{\left(N\right)}(.,t))\left(q_i^{\left(N\right)}(t)\right).$$
(31)
However, standard identities of vector analysis and the radial symmetry of $`\varrho `$ yield for the (negative of the) field torque on the $`i`$-th particle
$$(\varrho \mathrm{Id}\times \psi ^{\left(N\right)}(.,t))(x)=_y^{}\times [\varrho (yx)(yx)\psi ^{\left(N\right)}(y,t)]\mathrm{d}y0,$$
(32)
the vanishing as a result of one of Greenโs theorems and the compact support of $`\varrho `$. Hence, each $`s_i^{\left(N\right)}(t)`$ is itself a constant of the motion. Moreover, since there is no feedback loop from $`s_i^{\left(N\right)}(t)`$ to the particle-field dynamics, the introduction of spin into this model is uncalled for.
### 2.3 Global existence and uniqueness
In this subsection we extend the single particle global existence and regularity results of \[KSpK97\], \[KKSp99\] to the many-body problem. To get started, one needs decent a-priori bounds on the norms of the various dynamical quantities.
#### 2.3.1 A-priori bounds without invoking conservation laws
We begin with a-priori bounds that can be obtained without invoking the conservation laws. It is trivially clear by the upper bound 1 on their speeds that the positions of the particles are bounded above linearly in $`t`$. In the following we recall the familiar linear in $`t`$ a-priori estimate for the field norms, and a bound on the momenta quadratic in $`t`$.
###### Lemma 2.4
Let $`(q_1^{\left(N\right)}(.),\mathrm{},q_N^{\left(N\right)}(.))C^{0,1}(,^{3N})`$ be a given Lipschitz-continuous curve, its components having Lipschitz constant $`<1`$, and let $`\zeta _.C^1(,(\dot{H}^1L^2)(^3))`$ be a strong solution of (16), (17), satisfying conditions (14) (15). Then at any $`t`$,
$$\mathrm{max}\{\psi ^{\left(N\right)}(.,t)_{\dot{H}^1},\varpi ^{\left(N\right)}(.,t)_{L^2}\}(2_\mathrm{W}(\zeta _0^{\left(N\right)}))^{1/2}+\varrho _{L^2}|t|,$$
(33)
where $`_\mathrm{W}(\zeta _t)=\frac{1}{2}(|_x^{}\psi (.,t)|^2+|\varpi (.,t)|^2)\mathrm{d}x`$ is the wave field energy at time $`t`$.
###### Remark 2.5
The a-priori bound (33) extends to the strong solution of the wave equation for any subluminal source $`\varrho \rho C^0(,C_0^{\mathrm{}}(^3))`$.
Proof of Lemma 2.4: By hypothesis, $`t\zeta _t^{\left(N\right)}C^1(,(\dot{H}^1L^2)(^3))`$ is a strong solution of the wave equation with source $`t\varrho \rho _t^{\left(N\right)}C^0(,C_c^{\mathrm{}}(^3))`$ moving at speeds less than light; hence, $`t_\mathrm{W}(\zeta _t^{\left(N\right)})`$ is differentiable. We have
$`{\displaystyle \frac{\mathrm{d}}{\mathrm{d}t}}_\mathrm{W}(\zeta _t^{\left(N\right)})={\displaystyle \varpi ^{\left(N\right)}(t,x)(\varrho \rho _t^{\left(N\right)})(x)dx}`$ $``$ $`\varpi ^{\left(N\right)}(.,t)_{L^2}\varrho \rho _t^{\left(N\right)}_{L^2}`$ (34)
$``$ $`\varrho _{L^2}\left(2_\mathrm{W}(\zeta _t^{\left(N\right)})\right)^{1/2},`$ (35)
the first inequality by CauchyโSchwarz, while for the second one we used the estimate $`\varrho \rho _t^{\left(N\right)}_{L^2}^2sup_{x^3}\{(\varrho \varrho )(x)\}=\varrho _{L^2}^2`$, as well as the obvious estimate $`\varpi ^{\left(N\right)}(.,t)_{L^2}\left(2_\mathrm{W}(\zeta _t^{\left(N\right)})\right)^{1/2}`$ implied by the definition of $`_\mathrm{W}`$. Dividing (35) by $`\left(2_\mathrm{W}(\zeta _t^{\left(N\right)})\right)^{1/2}`$ and integrating over $`t`$ gives
$$(2_\mathrm{W}(\zeta _t^{\left(N\right)}))^{1/2}(2_\mathrm{W}(\zeta _0^{\left(N\right)}))^{1/2}+\varrho _{L^2}t.$$
(36)
The definition of $`_\mathrm{W}`$ given in Lemma 2.4 now shows that (36) implies (33). Q.E.D.
###### Lemma 2.6
Let $`(q_1^{\left(N\right)}(.),\mathrm{},q_N^{\left(N\right)}(.))C^{0,1}(,^{3N})`$ be a given Lipschitz-continuous curve, its components having Lipschitz constant $`<1`$, and let $`\zeta _.C^1(,(\dot{H}^1L^2)(^3))`$ be a strong solution of the wave equation with source $`t\varrho \rho _t^{\left(N\right)}C^0(,C_c^{\mathrm{}}(^3))`$. Suppose $`t(p_1^{\left(N\right)}(t),\mathrm{},p_N^{\left(N\right)}(t))`$ is a classical solution of (20). Then the momenta at $`t`$, $`p_k^{\left(N\right)}(t)`$, $`k=1,\mathrm{},N`$, are bounded by
$$\underset{1kN}{\mathrm{max}}\{p_k^{\left(N\right)}(t)\}\underset{1kN}{\mathrm{max}}\{p_k^{\left(N\right)}(0)\}+\varrho _{L^2}(2_\mathrm{W}(\zeta _0^{\left(N\right)}))^{1/2}|t|+\frac{1}{2}\varrho _{L^2}^2t^2.$$
(37)
Proof of Lemma 2.6: Use $`p(t)=p(0)+_0^t\dot{p}(\stackrel{~}{t})d\stackrel{~}{t}`$, take absolute values, use the triangle inequality, then invoke (20) and estimate
$$|(\varrho \psi ^{\left(N\right)}(.,t))\left(q_i^{\left(N\right)}(t)\right)|\varrho _{L^2}\psi ^{\left(N\right)}(.,t)_{\dot{H}^1},$$
(38)
then recall Lemma 2.4. Q.E.D.
###### Remark 2.7
We note that (36) is far from optimal, which is evident from the fact that no details of the time dependence of $`\rho _t^{\left(N\right)}`$ enter (36); in any event, squaring gives an upper bound on the wave field energy quadratic in $`t`$. Similarly, (37) is far from optimal; in any event, its right-hand side provides a quadratic-in-$`t`$ upper bound on the kinetic energy, โl.h.s.(48)$`1`$.โ These bounds together with the bound (38) and the asymptotics (14) now also imply that $`\left|\frac{1}{N}_{i=1}^N(\varrho \psi ^{\left(N\right)}(.,t))\left(q_i^{\left(N\right)}(t)\right)\right|`$, and therefore finally also the total energy, are both bounded above by $`a+b|t|+ct^2`$.
This does not yet exhaust our list of bounds that obtain without invoking conservation laws. The next such bound is nevertheless given its own subsection, for the special role it plays subsequently.
#### 2.3.2 A lower bound on the Hamiltonian functional
To state our lower bound on the Hamiltonian functional, we first define
$$E_{}:=1\frac{1}{8\pi }\frac{\varrho (x)\varrho (y)}{|xy|}dxdy.$$
(39)
Note that the energy value $`E_{}`$ depends only on the regularization but not on $`N`$.
###### Proposition 2.8
The Hamiltonian functional is bounded below by
$$\left(๐ท^{\left(N\right)}\right)E_{},$$
(40)
independently of $`N`$. The bound is attained when $`๐ท^{\left(N\right)}`$ is any translation in space of $`๐ท_{}^{\left(N\right)}`$, the state in which for all $`k=1,\mathrm{},N`$ we have $`q_k=0`$, $`p_k=0`$, and furthermore $`\varpi 0`$ and $`\psi \psi _\varrho `$, with
$$\psi _\varrho (x)=\frac{1}{4\pi }(|.|^1\varrho )(x).$$
(41)
(However, note that only the standard ground state satisfies (14).)
The state $`๐ท_{}^{\left(N\right)}`$ will be called the standard ground state of the regularized dynamical system, and (39) will be called the ground state energy.
Proof of Proposition 2.8. For later purposes, we will prove the bound (40) as an upper limit of a one-parameter family of bounds to $`\left(๐ท^{\left(N\right)}\right)`$. Thus, let $`\kappa (0,1]`$. Then
$`\left(๐ท^{\left(N\right)}\right)\frac{1\kappa }{2}\psi _{\dot{H}^1}^2=\frac{1}{N}\underset{i=1}{\overset{N}{}}(\sqrt{1+|p_i|^2}+\left(\varrho \psi \right)\left(q_i\right))+\frac{1}{2}\varpi _{L^2}^2+\frac{\kappa }{2}\psi _{\dot{H}^1}^2.`$ (42)
Discarding the manifestly positive momentum contributions we obtain
$$\left(๐ท^{\left(N\right)}\right)(1\kappa )\frac{1}{2}\psi _{\dot{H}^1}^21+\kappa \frac{1}{2}\psi _{\dot{H}^1}^2+\frac{1}{N}\underset{i=1}{\overset{N}{}}\left(\varrho \psi \right)\left(q_i\right).$$
(43)
Minimizing the right-hand side of (43) with respect to $`\psi `$ now gives
$$\left(๐ท^{\left(N\right)}\right)(1\kappa )\frac{1}{2}\psi _{\dot{H}^1}^21\frac{1}{N^2}\underset{i=1}{\overset{N}{}}\underset{k=1}{\overset{N}{}}\frac{1}{8\pi \kappa }\frac{\varrho \left(q_ix\right)\varrho \left(q_ky\right)}{|xy|}dxdy.$$
(44)
The right-hand side of (44) can be minimized w.r.t. the $`\left\{q_i\right\}_{i=1}^N`$ by equi-measurable, radially symmetric rearrangement of $`_{n=1}^N\varrho (.q_n)`$ centered at the origin. Since $`\varrho `$ is itself radially symmetric and decreasing, this is achieved by simply translating all $`q_n`$ to the same position, in particular to the origin. This gives, for all $`\kappa (0,1]`$,
$$\left(๐ท^{\left(N\right)}\right)(1\kappa )\frac{1}{2}\psi _{\dot{H}^1}^21\frac{1}{8\pi \kappa }\frac{\varrho (x)\varrho (y)}{|xy|}dxdy.$$
(45)
The bound (40) now obtains by taking $`\kappa =1`$ in (45) and recalling (39). Straightforward computation of $`\left(๐ท_{}^{\left(N\right)}\right)`$ proves that (40) is attained at $`๐ท_{}^{\left(N\right)}`$ and, by the translation invariance in position space of $`\left(๐ท^{\left(N\right)}\right)`$, also at any translate of $`๐ท_{}^{\left(N\right)}`$. Q.E.D.
#### 2.3.3 Bounds invoking conservation laws
Using energy conservation of sufficiently regular solutions, we next bootstrap from the proof of Proposition 2.8 to uniform bounds in $`t`$ and $`N`$ on the four major additive contributions to $`\left(๐ท_t^{\left(N\right)}\right)`$.
###### Lemma 2.9
Let $`t๐ท_t^{\left(N\right)}`$ be a sufficiently regular (e.g. $`\mathrm{\Gamma }^{\left(N\right)}`$-strong) solution of the dynamical system (16)โ(20) conserving energy. Then, uniformly in $`t`$ and $`N`$, we have
$`\psi ^{\left(N\right)}(.,t)_{\dot{H}^1}^2`$ $`4+4E8E_{},`$ (46)
$`\varpi ^{\left(N\right)}(.,t)_{L^2}^2`$ $`2E2E_{},`$ (47)
$`{\displaystyle \frac{1}{N}}\underset{i=1}{\overset{N}{}}\sqrt{1+|p_i^{\left(N\right)}(t)|^2}`$ $`1+EE_{},`$ (48)
$$6E_{}3E3\frac{1}{N}\underset{i=1}{\overset{N}{}}(\varrho \psi ^{\left(N\right)}(.,t))\left(q_i^{\left(N\right)}(t)\right)E1$$
(49)
Proof of Lemma 2.9. Since $`\left(๐ท_t^{\left(N\right)}\right)=E`$ is fixed by the Cauchy data $`๐ท_0^{\left(N\right)}`$, a simple rewriting of (45) with $`\kappa =1/2`$, using the definition (39), gives us (46). As to (47) and (48), $`\left(๐ท_t^{\left(N\right)}\right)=E`$ and the definition (26) of $`\left(๐ท^{\left(N\right)}\right)`$ give us the identity
$`{\displaystyle \frac{1}{N}}\underset{i=1}{\overset{N}{}}\sqrt{1+|p_i^{\left(N\right)}(t)|^2}+`$ $`\frac{1}{2}\varpi ^{\left(N\right)}(.,t)_{L^2}^2=`$ (51)
$`E\frac{1}{2}\psi ^{\left(N\right)}(.,t)_{\dot{H}^1}^2{\displaystyle \frac{1}{N}}\underset{i=1}{\overset{N}{}}(\varrho \psi ^{\left(N\right)}(.,t))\left(q_i^{\left(N\right)}(t)\right).`$
Recalling the minimization steps that lead from (43) to (45) (here with $`\kappa =1`$), and the definition (39), we see that the right-hand side of (51) is bounded above, giving
$$\frac{1}{N}\underset{i=1}{\overset{N}{}}\sqrt{1+|p_i^{\left(N\right)}(t)|^2}+\frac{1}{2}\varpi ^{\left(N\right)}(.,t)_{L^2}^21+EE_{}.$$
(52)
Now (47) follows at once from (52) by estimating $`|p_i^{\left(N\right)}(t)|0`$; to get (48), we instead use $`\varpi ^{\left(N\right)}(.,t)_{L^2}^20`$ in (52). Finally, to obtain (49), rewrite the definition (26) of $`\left(๐ท^{\left(N\right)}\right)`$ into an identity for $`\frac{1}{N}_{i=1}^N(\varrho \psi ^{\left(N\right)}(.,t))\left(q_i^{\left(N\right)}(t)\right)`$, then use $`\left(๐ท_t^{\left(N\right)}\right)=E`$; now the bounds (46), (47), (48) give the first, the positivity of $`\psi ^{\left(N\right)}(.,t)_{\dot{H}^1}^2`$, $`\varpi ^{\left(N\right)}(.,t)_{L^2}^2`$, and $`|p_i^{\left(N\right)}(t)|`$ the second inequality in (49). Q.E.D.
###### Remark 2.10
The bounds (47) and (48) happen to be asymptotically sharp when $`EE_{}`$, in which case they correctly imply that $`\varpi ^{\left(N\right)}(.,t)_{L^2}0`$ and $`|p_i^{\left(N\right)}(t)|0`$ for all $`i=1,\mathrm{},N`$. It is to be doubted though that (47) and (48) are sharp for $`E>E_{}`$; in any event, certainly (46) and (49) are not sharp (for instance, (46) misses the correct ground state value by a factor 2). Of course, it is a straightforward matter to improve on (46) and (49) by optimizing w.r.t. $`\kappa `$ (N.B.: $`\kappa =1/2`$ is the optimizer for $`E=1`$), and while this does lead to asymptotically sharp upper bounds as $`EE_{}`$ (in which case $`\kappa 1`$), for $`E>E_{}`$ these bounds are still not sharp, but now more cumbersome than (46) and (49). Fortunately, for our purposes any a-priori bounds uniform in $`t`$ and $`N`$ will do; hence, we gain by sticking to the simple ones given in Lemma 2.9.
As a corollary to (48) the particle momenta are bounded above in magnitude. This has an easy but important corollary for the particle speeds, which we state explicitly.
###### Corollary 2.11
The particle speeds are bounded away from the speed of light, viz.
$$\underset{i\{1,\mathrm{},N\}}{\mathrm{max}}|\dot{q}_i^{\left(N\right)}(t)|\sqrt{1\left(1+N(EE_{})\right)^2},$$
(53)
uniformly in $`t`$. In particular, when $`E=E_{}`$, then $`|\dot{q}_i^{\left(N\right)}(t)|=0`$ for all $`i`$ and $`N`$.
Proof of Corollary 2.11. We rewrite (48) as
$$\frac{1}{N}\underset{i=1}{\overset{N}{}}\left(\sqrt{1+|p_i^{\left(N\right)}(t)|^2}1\right)EE_{}.$$
(54)
Since $`\sqrt{1+|p|^2}10`$, the bound (54) now implies that for all $`i`$,
$$\sqrt{1+|p_i^{\left(N\right)}(t)|^2}1N(EE_{}),$$
(55)
and solving for $`|p_i^{\left(N\right)}(t)|`$ gives, uniformly in $`t`$,
$$\underset{i\{1,\mathrm{},N\}}{\mathrm{max}}|p_i^{\left(N\right)}(t)|\sqrt{\left(1+N(EE_{})\right)^21},$$
(56)
which now yields (53) by inverting the monotone map $`|p||v|`$ given in (2). Q.E.D.
Note that for any $`E>E_{}`$, (53) does not imply boundedness away from the speed of light of the $`|\dot{q}_i^{\left(N\right)}(t)|`$ uniformly in $`N`$; only $`\mathrm{max}_i|\dot{q}_i^{\left(N\right)}(t)|1`$ holds uniformly in $`N`$.
#### 2.3.4 Global existence and uniqueness of solutions
Lemma 2.9 and Corollary 2.11 imply that any energy-conserving solution is represented by a point moving in a weakly compact subset of $`\mathrm{\Gamma }^{\left(N\right)}`$, and such solutions do exist.
###### Theorem 2.12
For every $`๐ท_0^{\left(N\right)}\mathrm{\Gamma }_\mathrm{B}^{\left(N\right)}`$ there exists a unique, global strong solution $`t๐ท_t^{\left(N\right)}C^1(,\mathrm{\Gamma }^{\left(N\right)})`$ of the Hamiltonian field & $`N`$-body problem (16)โ(20), satisfying $`lim_{t0}๐ท_t^{\left(N\right)}=๐ท_0^{\left(N\right)}`$, and conserving mass, energy, momentum and angular momentum as stated in Proposition 2.2. For more regular initial data one can bootstrap to correspondingly higher regularity of $`t๐ท_t^{\left(N\right)}`$.
Proof of Theorem 2.12. The proof is a largely straightforward adaption to our many-body problem of the proof for a single particle system in \[KSpK97\]. We remark that our Wiener condition for $`\varrho `$ is only needed to adapt their proof. The strategy is to first construct local weak solutions conserving energy, then to use the uniform bounds on the norms of the various dynamical quantities that follow from energy conservation (see Lemma 2.9 and its Corollary) to continue to all times. Strong solutions obtain by restricting $`\psi (.,0)(\dot{H}^1\dot{H}^2)(^3)`$. Proofs of the conservation laws, which are stated without proof in \[KSpK97\], are provided in our appendix, for the convenience of the reader. Q.E.D.
###### Remark 2.13
For the proof of Theorem 2.12 the a-priori bounds in Lemma 2.9 based on energy conservation suffice. The other conservation laws provide additional bounds that may be useful in different contexts. For instance, momentum conservation and the CauchyโSchwarz inequality give us the uniform bound in $`t`$,
$$\frac{1}{N}|\underset{i=1}{\overset{N}{}}p_i^{\left(N\right)}(t)||P|+\psi ^{\left(N\right)}(.,t)_{\dot{H}^1}\varpi ^{\left(N\right)}(.,t)_{L^2},$$
(57)
while angular momentum conservation, the CauchyโSchwarz inequality, and the finite wave speed give us the linear bound in $`t`$,
$$\frac{1}{N}|\underset{i=1}{\overset{N}{}}p_i^{\left(N\right)}(t)\times q_i^{\left(N\right)}(t)||J|+(R+|t|)\psi ^{\left(N\right)}(.,t)_{\dot{H}^1}\varpi ^{\left(N\right)}(.,t)_{L^2},$$
(58)
where $`R`$ is the radius of the ball $`B_R`$ containing the initial positions of all particles and the supports of their regularizations, and outside of which (14) and (15) hold.
###### Remark 2.14
As is the case for the regularized VlasovโPoisson equations \[Spo91\], the solutions described by Theorem 2.12 map one-to-one into generalized solutions of the regularized wave gravity Vlasov model in which derivatives of $`f`$ are meant in the sense of distributions.
## 3 The regularized Vlasov model
In this section we discuss the regularized wave gravity Vlasov model. First, we present the Vlasov equations formally as a continuum model. Next we recall the concept of generalized (distributional) solutions, for which we introduce two suitable topologies, one based on the vague and one on a strong Banach space topology for distributions. The solutions to the field & $`N`$-body model of the previous section furnish particular generalized solutions of our Vlasov model in either of the just mentioned topologies. We then prove global existence and uniqueness in the strong Banach space topology of generalized solutions to our regularized wave gravity Vlasov model.
### 3.1 The dynamical continuum system
As first order system, the inhomogeneous wave equation for the regularized wave gravity potential $`\psi (.,t)\dot{H}^1(^3)`$ and its conjugate variable $`\varpi (.,t)L^2(^3)`$ now reads
$`_t^{}\psi (x,t)=\varpi (x,t)`$ (59)
$`_t^{}\varpi (x,t)=_x^2\psi (x,t)(\varrho \rho (.,t))(x).`$ (60)
The initial data $`\psi (.,0)\psi _0(.)(\dot{H}^1\dot{H}^2)(^3)`$ and $`\varpi (.,0)\varpi _0(.)L^2(^3)`$ satisfy
$`\psi _0(x)=1/4\pi |x|,`$ (61)
$`\varpi _0(x)=0`$ (62)
outside some closed ball $`B_R^3`$. The density $`\rho (x,t)`$ on the r.h.s. in (60) is given by
$$\rho (x,t)=f(x,p,t)dp,$$
(63)
where $`f(.,.,t)`$ is the normalized particle density function at time $`t`$, satisfying the following (continuity) equation on time-position-momentum space $`\times ^3\times ^3`$,
$$_t^{}f(x,p,t)=(_p^{}\sqrt{1+|p|^2}_x^{}_x^{}\left(\varrho \stackrel{}{}\psi (.,t)\right)(x)_p^{})f(x,p,t),$$
(64)
with $`x^3`$ being the space and $`p^3`$ the momentum variable. Initial data $`f(.,.,0)f_0(.,.)`$ for (64) are restricted by the requirement that $`\varrho \rho (.,0)`$ is supported in $`B_R`$. As to the appropriate function space, we re-emphasize that in the form stated above, one should think of Vlasovโs $`f(.,.,t)`$ as a continuum approximation to the empirical $`x,p`$ phase space density of particles for an actual individual $`N`$-body system in the large $`N`$ regime, when fine details of the particlesโ behaviors become irrelevant on the โmacroscopicโ scales so that the empirical atomic measure can be well approximated by a function $`f(.,.,t)L_{+,1}^1(^6)`$ โ the subset of $`L^1(^6)`$ consisting of the RadonโNikodym derivatives $`f`$ of Borel probability measures $`\mu ^f(\mathrm{d}x\mathrm{d}p)`$ which are absolutely continuous w.r.t. Lebesgue measure. In fact, such functions $`f(.,.,t)`$, the fields $`\psi (.,t)`$ and their formal time derivatives $`\varpi (.,t)`$, would even be expected to have time and space derivatives in the classical sense for all time whenever their initial data are chosen sufficiently regular.
###### Remark 3.1
It is known, but perhaps not well-known, that for sufficiently regular solutions $`f`$ (say, classical with rapid decay at infinity), the continuity equation (64) can readily be associated with a Hamiltonian $`_\mathrm{C}`$, given $`\psi (.,t)`$ for all $`t`$. To obtain the Hamiltonian $`_\mathrm{C}`$ for $`f`$ given $`\psi (.,t)`$, multiply r.h.s.(64) with a test function $`g(.,.,t)C^1(^6)`$ of at most polynomial growth in $`x,p`$ whose $`t`$-dependence is yet to be determined, and integrate over $`^6`$; $`f`$ and $`g`$ can now be viewed as conjugate variables, with $`_t^{}f=\delta _\mathrm{C}(f,g)/\delta g`$, $`_t^{}g=\delta _\mathrm{C}(f,g)/\delta f`$. Interestingly, the equation for $`g`$ is just (64) with $`g`$ in place of $`f`$, and in this sense (64) already is the Hamiltonian system, given the fields. The inhomogeneous wave equation (59), (60) for the fields $`\psi (.,t)`$ and $`\varpi (.,t)`$ is a Hamiltonian dynamical system, given $`f`$. The full set of equations (59), (60), (64) becomes a Hamiltonian system with the help of non-canonical Lie brackets, cf. \[Mor80, WeMo81, MMW84\].
Our goal is to validate the continuum approximation to the microscopic atomistic dynamics by means of a continuum limit in $`x,p`$ space (the โVlasov limitโ), supplemented by a law of large numbers and a central limit theorem. To pave the way for the continuum validation, we next recall the concept of generalized solutions.
### 3.2 Distributional form of the regularized Vlasov model
In order to think of $`f(.,.,t)`$ as the actual atomic measure of an individual $`N`$-body system, one has to interpret the derivatives in the sense of distributions. Thus, for given $`\psi (.,t)\dot{H}^1(^3)`$ and $`\varpi (.,t)L^2(^3)`$, we implement the idea of distributional derivatives of $`f`$ in the usual way by multiplying (all of) (64) with any real test function $`g(.,.)C_0^1(^6)`$ and integrate over $`^6`$ by parts to transfer the partial derivatives w.r.t. $`x,p`$ onto the smooth $`g`$; also, the partial derivative w.r.t. $`t`$ is pulled out of the integral. So far, $`f(.,.,t)`$ had to be a sufficiently regular function, but nothing now prevents us from allowing $`fL_{+,1}^1(^6)`$, the RadonโNikodym derivative of an absolutely continuous measure $`\mu ^f`$. The so integrated and manipulated form of (64) remains well-defined even if we replace $`\mu ^f(\mathrm{d}x\mathrm{d}p)`$ with any Borel probability measure $`\mu _t(\mathrm{d}x\mathrm{d}p)`$. Indeed, let $`Dg`$ denote any of the partial derivatives of $`g`$. Then $`DgC_0^0(^6)`$, where $`C_0^0(^6)`$ is equipped with the uniform norm (a.k.a. sup-norm) $`Dg_\mathrm{u}=sup_{z^6}|Dg(z)|`$. On the other hand, the Borel probability measures $`P(^6)`$ are a subset of $`M(^6)`$, the Banach space of finite signed Radon measures $`\sigma `$ (which on $`^6`$ coincide with the finite regular signed Borel measures $`\sigma `$) equipped with the total variation (TV) norm $`\sigma _{\mathrm{TV}}=(|\sigma _+|+|\sigma _{}|)(^6)`$, and $`M(^6)`$ is isometrically isomorphic to $`C_0^0(^6)^{}`$, the dual space for (real) $`C_0^0(^6)`$.
In the above we used $`z`$ to denote a generic point $`(x,p)^6`$. In the same vein, we sometimes write $`\zeta `$ for the generic wave variables $`(\psi ,\varpi )`$.
A physical generalized state of the regularized Vlasov model constitutes a complete specification of all its first-order evolutionary variables. We accordingly define the set $`\mathrm{\Gamma }`$ of all possible physical generalized states at time $`t`$ to be the subset of points
$$_t:=(\mu _t,\zeta _t)M(^6)\dot{H}^1(^3)L^2(^3)$$
(65)
for which $`\mu _tP(^6)`$. The subset $`\mathrm{\Gamma }_\mathrm{B}`$ of $`\mathrm{\Gamma }`$ denotes those physical generalized states which satisfy (61) and (62), and for which $`\mathrm{supp}(\mu _0(\mathrm{d}x\times ^3))B_R`$.
In (65), the first direct sum is clearly in the sense of Banach spaces while the second may be either in Banach or Hilbert space sense (with the understanding then that Hilbert binds stronger than Banach on its left); since it is a little awkward to have two direct sum symbols with different meanings in a single expression, *for the Vlasov model* we use the Banach space meaning throughout.
In this paper we are only interested in systems with finite energy, momentum, and angular momentum (the mass of a system is finite by default, namely unity). Thus, for a suitable subset of generic physical generalized states $``$ in $`\mathrm{\Gamma }_\mathrm{B}`$ we formally define
the mass functional
$$\left(\right)=\mu (\mathrm{d}z),$$
(66)
the momentum functional
$$๐ซ\left(\right)=p\mu (\mathrm{d}z)\varpi _x^{}\psi \mathrm{d}x,$$
(67)
the angular momentum functional
$$๐ฅ\left(\right)=x\times p\mu (\mathrm{d}z)\varpi x\times _x^{}\psi \mathrm{d}x,$$
(68)
and the energy functional
$$\left(\right)=\left(\sqrt{1+|p|^2}+(\varrho \psi )\right)\mu (\mathrm{d}z)+\frac{1}{2}\left(|_x^{}\psi |^2+|\varpi |^2\right)dx.$$
(69)
Here, in keeping with our already stipulated abbreviations, $`\mathrm{d}z`$ denotes the Lebesgue measure $`\mathrm{d}x\mathrm{d}p`$ on $`^6`$ and $`\mathrm{}\mu (\mathrm{d}z)`$ stands for $`\mathrm{}\mu (\mathrm{d}x\mathrm{d}p)`$. We restrict the set of physical generalized states to measures with finite expected values of $`|x|`$ and $`|p|`$. Now $`|x|`$ and $`|p|`$, understood as functions on $`^6`$, are not in $`C_0^0(^6)`$, but they are lower semi-continuous and therefore Radon measurable; hence, our condition of finite expected values of $`|x|`$ and $`|p|`$ (equivalently, of $`|z|`$) defines a proper subset $`P_1(^6)`$ of the Borel probability measures. The corresponding subset of the physical states $`\mathrm{\Gamma }_\mathrm{B}`$ is denoted by $`\mathrm{\Gamma }_{\mathrm{B},1}`$; the energy, momentum, and angular momentum are well-defined on $`\mathrm{\Gamma }_{\mathrm{B},1}`$.
It remains to stipulate a suitable topology on $`\mathrm{\Gamma }_{\mathrm{B},1}`$ in which the maps $`t_t`$ for $`t`$ are continuous curves in $`\mathrm{\Gamma }_{\mathrm{B},1}`$ that qualify as generalized solutions of (64) (given the fields). Unfortunately, the Banach space topology which $`\mathrm{\Gamma }_{\mathrm{B},1}`$ naturally inherits as a subset of $`C_0^0(^6)^{}\dot{H}^1(^3)L^2(^3)`$ is too strong to study families of empirical atomic measures, for any two (atomic) empirical measures with disjoint supports are always at a distance 2 from each other in the metric induced by the TV topology.
A more suitable topology that immediately comes to mind is the vague (a.k.a. weak) topology on $`M(^6)C_0^0(^6)^{}`$ induced by $`C_0^0(^6)`$. The set $`\mathrm{\Gamma }_{\mathrm{B},1}`$ with the vague topology on $`M(^6)`$ in place of the TV topology is denoted by $`\mathrm{\Gamma }_{\mathrm{B},1}^v`$. However, since we are interested only in the subset $`P_1(^6)M(^6)`$, we can do somewhat better and equip $`P_1(^6)`$ with the the standard KantorovichโRubinstein topology<sup>7</sup><sup>7</sup>7The relationship between the various topologies is summarized in Appendix A.1. induced by the dual Lipschitz distance in $`P_1(^6)`$ (a map on $`P_1(^6)\times P_1(^6)`$),
$$\mathrm{dist}_\mathrm{L}^{}(\mu _1,\mu _2):=\underset{gC^{0,1}(^6)}{sup}\left\{\right|g\mathrm{d}(\mu _1\mu _2)|:\mathrm{Lip}\left(g\right)1\}.$$
(70)
We write $`\mu _n\mu `$ if $`\mathrm{dist}_\mathrm{L}^{}(\mu _n,\mu )0`$.
Since it is convenient for the presentation to have a Banach space, we note that the metric $`\mathrm{dist}_\mathrm{L}^{}(.,.)`$ defines a norm on $`P_1P_1`$ by $`\sigma _\mathrm{L}^{}:=\mathrm{dist}_\mathrm{L}^{}(\sigma _+,\sigma _{})`$ for $`\sigma (P_1P_1)(^6)`$. As described in Appendix A.1, $`._\mathrm{L}^{}`$ can be extended<sup>8</sup><sup>8</sup>8This extension will not be needed for any of our technical estimates. to a norm $`._{\stackrel{~}{\mathrm{L}^{}}}`$ on the linear span of $`P_1(^6)`$, such that $`\sigma _{\stackrel{~}{\mathrm{L}^{}}}=\sigma _\mathrm{L}^{}`$ whenever $`\sigma (^6)=0`$. The completion of the linear span of $`P_1(^6)`$ w.r.t. $`._{\stackrel{~}{\mathrm{L}^{}}}`$, denoted $`\stackrel{~}{M}_1(^6)`$, is a Banach space. We also write $`\stackrel{~}{P}_1(^6)`$ for $`P_1(^6)\stackrel{~}{M}_1(^6)`$. By $`\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1}`$ we denote the closed subset of generic phase space points
$$=(\mu ;\zeta )\stackrel{~}{M}_1(^6)(\dot{H}^1L^2)(^3)$$
(71)
for which $`\mu \stackrel{~}{P}_1(^6)`$, and for which $`\zeta `$ satisfies (61), (62); once again, the Banach space direct sum is meant on the right-hand side, so that in particular the norm of $`_t`$ reads
$$:=\mu _{\stackrel{~}{\mathrm{L}^{}}}+\zeta _{HL}$$
(72)
with $`\zeta _t_{HL}=\psi _{\dot{H}^1}+\varpi _{L^2}`$.
#### 3.2.1 Generalized solutions w.r.t. to the vague topology for $`\mu `$
Considering first the field variables $`\zeta _.C^1(,(\dot{H}^1L^2)(^3))`$ as given, we will call $`M`$-vague solution of (64) a $`C^1`$ map $`t\mu _t`$ satisfying (64) with $`\mu _t`$ in place of $`f`$, for all $`t`$ integrated against any test function $`gC_0^1(^6)`$, and with the $`_t^{}`$ pulled in front of the corresponding integral. Accordingly, a map $`t_t=(\mu _t;\zeta _t)`$ $`C^1(,\mathrm{\Gamma }_{\mathrm{B},1}^v)`$ will be called $`M`$-vague $`HL`$-strong solution of (59), (60), (64).
With the help of the concept of the $`M`$-vague $`HL`$-strong solution of (59), (60), (64), we can now immediately reformulate Theorem 2.12 into an existence result for what we call $`M`$-vague $`N`$-body solutions of the regularized Vlasov model (59), (60), (64).
###### Theorem 3.2
Let $`t(๐ณ_t^{\left(N\right)},\zeta _t^{\left(N\right)})=๐ท_t^{\left(N\right)}C^1(,\mathrm{\Gamma }_\mathrm{B}^{\left(N\right)})`$, with $`\underset{t0}{lim}๐ท_t^{\left(N\right)}=(๐ณ_0^{\left(N\right)},\zeta [๐ณ_0^{\left(N\right)}])`$, be the unique strong solution of the Hamiltonian field & $`N`$-body problem (16)โ(20), and denote the empirical measure associated to $`๐ณ_t^{\left(N\right)}`$ by
$$\epsilon [๐ณ_t^{\left(N\right)}](\mathrm{d}x\mathrm{d}p)=\frac{1}{N}_{k=1}^N\delta _{q_k^{(N)}(t)}(\mathrm{d}x)\times \delta _{p_k^{(N)}(t)}(\mathrm{d}p)$$
(73)
Then $`(t(\epsilon [๐ณ_t^{\left(N\right)}];\zeta _t^{\left(N\right)})=_t^{\left(N\right)})`$ $`C^1(,\mathrm{\Gamma }_{\mathrm{B},1}^v)`$ is an $`M`$-vague $`N`$-body solution of the regularized wave gravity Vlasov equations (59), (60), (64), satisfying the Cauchy data $`\underset{t0}{lim}_t^{\left(N\right)}=_0^{\left(N\right)}`$, and conserving mass, momentum, angular momentum, and energy:
$$\left(_t^{\left(N\right)}\right)=M,$$
(74)
$$๐ซ\left(_t^{\left(N\right)}\right)=P,$$
(75)
$$๐ฅ\left(_t^{\left(N\right)}\right)=J,$$
(76)
$$\left(_t^{\left(N\right)}\right)=E,$$
(77)
with $`M,P,J,E`$ independent of time; in particular, $`M=1`$.
Note that by Thm. 3.2 the set $`\mathrm{\Gamma }_\mathrm{B}^{\left(N\right)}`$ becomes identified with a subset of $`\mathrm{\Gamma }_{\mathrm{B},1}^v`$.
#### 3.2.2 Generalized solutions w.r.t. the KantorovichโRubinstein topology for $`\mu `$
Since $`\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1}`$ is equipped with a Banach space topology, a map $`t_t=(\mu _t;\zeta _t)`$ $`C^1(,\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1})`$ satisfying (59), (60), (64) is properly called a $`\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1}`$-strong generalized solution of our regularized Vlasov model. Such solution satisfy the conventional conservation laws. Particular $`\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1}`$-strong generalized solutions, called $`\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1}`$-strong $`N`$-body solutions, are generated by the solutions of the field & $`N`$-body model of section 2. We summarize this in
###### Theorem 3.3
Let $`t_t=(\mu _t;\zeta _t)`$ $`C^1(,\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1})`$ be a $`\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1}`$-strong generalized solution of (59), (60), (64) with Cauchy data $`lim_{t0}_t=_0`$. Then mass, momentum, angular momentum, and energy are conserved; i.e. (74), (75), (76), (77) hold. In particular, let $`t๐ท_t^{\left(N\right)}C^1(,\mathrm{\Gamma }_\mathrm{B}^{\left(N\right)})`$ and $`\epsilon [๐ณ_t^{\left(N\right)}](\mathrm{d}x\mathrm{d}p)`$ be given as in Theorem 3.2. Then $`t_t^{\left(N\right)}=(\epsilon [๐ณ_t^{\left(N\right)}];\zeta _t^{\left(N\right)})`$ $`C^1(,\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1})`$ is a $`\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1}`$-strong $`N`$-body solution of the regularized wave gravity Vlasov equations (59),(60),(64), with Cauchy data $`\underset{t0}{sup}_t^{\left(N\right)}=_0^{\left(N\right)}`$.
Note that by Theorem 3.3 the set $`\mathrm{\Gamma }_\mathrm{B}^{\left(N\right)}`$ becomes identified with a subset of $`\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1}`$.
We next show that arbitrary initial data $`_0\mathrm{\Gamma }_{\mathrm{B},1}`$ launch a unique $`\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1}`$-strong generalized solution $`t_tC^1(,\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1})`$ of our Vlasov model. Since the vague topology on $`P_1`$ is controlled by the standard KantorovichโRubinstein topology, solutions of the type $`C^1(,\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1})`$ are automatically solutions of the type $`C^1(,\mathrm{\Gamma }_{\mathrm{B},1}^v)`$.
### 3.3 The Cauchy problem for $`\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1}`$-strong solutions
To study the general Cauchy problem for (59), (60), (64) in the $`\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1}`$-strong topology, we rewrite (59), (60), (64) together with their Cauchy data as a fixed point problem,
$$_.=F_{.,0}(_.|_0),$$
(78)
where $`F_{.,0}`$ is a continuous map from $`C^0(,\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1})`$ into $`C^0(,\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1})`$, conditioned on $`_0\mathrm{\Gamma }_{\mathrm{B},1}`$. We will show that, w.r.t. a suitably weighted sup-norm, a truncated version of $`F`$ is a Lipschitz map, with Lipschitz constant $`<1`$, from a closed subset of weighted $`C^0(,\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1})`$ into itself. Existence of a unique fixed point of the truncated $`F`$ then follows from the standard contraction mapping theorem. By bootstrapping regularity, fixed points of the full $`F`$ will then be shown to exist and to be in $`C^1(,\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1})`$, thus furnishing unique $`\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1}`$-strong solutions of (59), (60), (64) that conserve mass (66), momentum (67), angular momentum (68), and energy (69).
#### 3.3.1 Definition of the fixed point map
Given any $`_0=(\mu _0;\zeta _0)\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1}`$, for each $`t`$ the map $`F_{t,0}(.;.|\mu _0;\zeta _0)`$ is given by
$$F_{t,0}(\mu _.;\zeta _.|\mu _0;\zeta _0)(\mathrm{\Pi }_{t,0}^{}[\zeta _.](\mu _0);\mathrm{\Phi }_{t,0}^{}[\mu _.](\zeta _0)),$$
(79)
where $`\mathrm{\Pi }_{t,0}^{}[\zeta _.](\mu _0)\mu _0\mathrm{\Pi }_{0,t}^{}[\zeta _.]`$, and where $`\mathrm{\Pi }_{.,.}`$ and $`\mathrm{\Phi }_{.,.}`$ are two-parameter groups, flows on the phase subspaces of the particles and the fields, respectively. Given a trajectory $`\zeta _.`$ in field space, $`\mathrm{\Pi }_{.,.}[\zeta _.]`$ is the particle phase space flow, and given a trajectory $`\mu _.`$ in probability measure space, $`\mathrm{\Phi }_{.,.}[\mu _.]`$ is the field phase space flow.
As to the flow on particle phase space, let $`t\zeta _tC_b^0(,(\dot{H}^1L^2)(^3))`$ be a generic, bounded continuous curve in $`(\dot{H}^1L^2)(^3)`$. Given $`t\zeta _t`$, the characteristic equations for (64) are the Hamiltonian equations for test particle motion $`\mathrm{d}z/\mathrm{d}t=๐_z^{}(z,\zeta _t)`$, with $`z=(q,p)`$, where $`๐`$ is the symplectic matrix, and $`(z,\zeta _t)`$ is the Hamiltonian (26) for $`N=1`$ and with $`(z,\zeta _t)`$ substituted for $`๐ท^{(1)}`$. Explicitly,
$$๐_z^{}(z,\zeta _t)=(_p^{}\sqrt{1+|p|^2},_x^{}(\varrho \stackrel{}{}\psi (.,t))(q));$$
(80)
note that only the $`\psi `$ part of $`\zeta `$ enters in (80). The particle phase space flow $`\mathrm{\Pi }_{.,.}[\zeta _.]`$ is now defined implicitly as follows: given $`t\psi (.,t)`$, for each solution $`z_.C^1(,^6)`$ of the characteristic equations the integrated characteristic equations give the identity
$$z_t=z_t^{}+_t^{}^t๐_z^{}(z_\tau ,\zeta _\tau )\mathrm{d}\tau =:\mathrm{\Pi }_{t,t^{}}[\zeta _.]\left(z_t^{}\right),$$
(81)
the r.h.s. of which being the transition function from some $`z`$ at time $`t^{}`$ to another $`z`$ at time $`t`$ for all $`t`$, $`t^{}`$; considering the totality of all $`t`$, $`t^{}`$ gives the particle flow.
Similarly, to define the flow on field phase space, suppose $`t\mu _tC^0(,\stackrel{~}{P}_1(^6))`$ is given, and let $`\rho _t(\mathrm{d}x)=\mu _t(\mathrm{d}x\mathrm{d}p)`$, and $`(\varrho \rho _t)(x)=\varrho (yx)\rho _t(\mathrm{d}y)`$. Then $`\frac{1}{2T}_T^T\varrho \rho _t_{L^2(^3)}^2dt\varrho \varrho _\mathrm{u}\left(\varrho _{L^2}^2\right)`$. Given such $`\varrho \rho _.L^2([T,T],L^2(^3))`$ for any $`T>0`$, the solution $`\zeta _.=(\psi ,\varpi )(.,.)`$ to the wave equation with field source $`\varrho \rho _.`$ defines the flow $`\mathrm{\Phi }_{.,.}[\mu _.]`$ on field space through the transition function
$$\zeta _t=:\mathrm{\Phi }_{t,t^{}}[\mu _.](\zeta _t^{}).$$
(82)
An explicit representation of (82) in terms of Fourier & Laplace transforms is available. However, by the higher regularity of $`\varrho \rho _.`$, the $`\dot{H}^1`$ and $`L^2`$ estimates of $`\mathrm{\Phi }_{.,.}[\mu _.]`$ are conveniently obtained from Kirchhoffโs explicit pointwise expressions for classical solutions. In components, $`\mathrm{\Phi }_{.,.}[\mu _.](\mathrm{\Phi }_{.,.}^\psi [\mu _.],\mathrm{\Phi }_{.,.}^\varpi [\mu _.])`$ reads (\[Bre93\], \[Ika00\], \[ShSt00\])
$`\psi (x,t)`$ $`=`$ $`{\displaystyle _{\mathrm{SS}^2}}([1+(tt^{})\mathrm{\Omega }]\psi (x^{},t^{})+(tt^{})\varpi (x^{},t^{}).`$ (84)
$`.{\displaystyle _t^{}^t}(tt^{\prime \prime })(\varrho \rho _{t^{\prime \prime }})(x^{\prime \prime })\mathrm{d}t^{\prime \prime })\mathrm{d}\mathrm{\Omega }`$
$`=:`$ $`\mathrm{\Phi }_{t,t^{}}^\psi [\mu _.](\zeta _t^{})(x)`$ (85)
$`\varpi (x,t)`$ $`=`$ $`{\displaystyle _{\mathrm{SS}^2}}([1+2(tt^{})\mathrm{\Omega }]\mathrm{\Omega }\psi (x^{},t^{})+[1+(tt^{})\mathrm{\Omega }]\varpi (x^{},t^{}).`$ (87)
$`.{\displaystyle _t^{}^t}[1+(tt^{\prime \prime })\mathrm{\Omega }](\varrho \rho _{t^{\prime \prime }})(x^{\prime \prime })\mathrm{d}t^{\prime \prime })\mathrm{d}\mathrm{\Omega }`$
$`=:`$ $`\mathrm{\Phi }_{t,t^{}}^\varpi [\mu _.](\zeta _t^{})(x),`$ (88)
where $`x^{\mathrm{or}\prime \prime }=x+(tt^{\mathrm{or}\prime \prime })\mathrm{\Omega }`$, where $`\mathrm{\Omega }\mathrm{SS}^2`$, and where $`_{\mathrm{SS}^2}`$ is short for $`\frac{1}{4\pi }_{\mathrm{SS}^2}`$.
Having defined $`\mathrm{\Pi }_{.,.}[\zeta _.]`$ and $`\mathrm{\Phi }_{.,.}[\mu _.]`$, we are now ready to analyse equation (78).
#### 3.3.2 Statement of the main fixed point results
So far, (78) has been defined purely formally as a rewriting of (59), (60), (64), with Cauchy data imposed. Our first duty should be to show that (78) in fact makes sense, viz. that $`F`$ maps a relevant, closed subset of $`C^0(,\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1})`$ into itself, indeed. We prove this as a byproduct of the auxiliary result that a truncated version of $`F_{.,0}(.|_0)`$ is a Lipschitz map, with Lipschitz constant $`<1`$, from a closed subset of $`C^0(,\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1})`$ into itself, where โclosedโ is meant w.r.t. (a suitably weighted) sup-norm.
We note that by a density argument for the curves of empirical measures $`t\rho _t^{\left(N\right)}`$, the a-priori estimates of section 2.3.1 extend to our regularized Vlasov model. Hence, any Vlasov solution $`t_t`$ must be in some subset of $`C^0(,\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1})`$ satisfying $`_tc_0+c_1|t|+c_2t^2`$ for some positive constants $`c_0`$, $`c_1`$ and $`c_2`$. This suggests to work with the closure of the bounded continuous functions from $``$ to $`\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1}`$, denoted $`C_b^0(,\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1})`$, w.r.t. a weighted sup-norm of $`_.`$ given by $`sup_t\left([c_0+c_1|t|+c_2t^2]^1_t\right)`$; however, for technical reasons it is more convenient to close $`C_b^0(,\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1})`$ w.r.t. the weighted sup-norm
$$_._\mathrm{w}=\underset{t}{sup}\left(e^{w|t|}_t\right)$$
(89)
for some $`w>0`$; eventually we will restrict $`w`$ to $`w>\underset{ยฏ}{w}>0`$. The closure of $`C_b^0(,\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1})`$ w.r.t. norm (89) is a Banach space, denoted $`C_\mathrm{w}^0(,\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1})`$; the subscript w can be read as meaning both โweightedโ and reference to the parameter $`w`$ in the definition (89). We also introduce the Banach space $`C_\mathrm{w}^0(_+,\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1}|_0)`$ with norm $`_._\mathrm{w}=sup_{t0}\left(e^{wt}_t\right)`$. By the time reversal symmetry of (78) it suffices to limit the discussion to $`t0`$.
Moreover, as regards the argument $`_.(=(\mu _.;\zeta _.))`$ of $`F_{.,0}(.|_0)(=F_{.,0}(.;.|\mu _0;\zeta _0))`$ given in (79), it is not a-priori required that $`lim_{t0}_t`$ of $`_.`$ coincides with the given $`_0`$; in fact, it is not even a-priori necessary that the measure component of $`_.`$ is in $`P_1`$ but could as well be in $`\stackrel{~}{M}_1`$. However, for the solution of (78) this must be so, for $`F_{.,0}(.;.|\mu _0;\zeta _0)`$ has been constructed such that $`F_{0,0}(\mu _.;\zeta _.|\mu _0;\zeta _0)=(\mu _0;\zeta _0)=_0`$, as is readily verified by inspection of $`\mathrm{\Pi }_{0,0}`$ and $`\mathrm{\Phi }_{0,0}`$. Therefore, with the exception of some technical estimates that we will highlight explicitly, we only need to apply $`F_{.,0}(.|_0)`$ to those $`_.C^0(,\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1})`$ which satisfy $`lim_{t0}_t=_0`$. We denote the corresponding subsets of $`C_\mathrm{w}^0(,\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1})`$ and $`C_\mathrm{w}^0(_+,\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1})`$ by $`C_\mathrm{w}^0(,\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1}|_0)`$ and $`C_\mathrm{w}^0(_+,\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1}|_0)`$, respectively. Furthermore we denote the free evolution of the initial data $`_0`$ by $`_.^0:=F_{.,0}(0_.|_0)C_\mathrm{w}^0(,\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1}|_0)`$; here, $`0_.`$ is the trivial constant map $`(t0)C_\mathrm{w}^0(,\stackrel{~}{M}_1(^6)\dot{H}^1(^3)L^2(^3))`$. We will work with certain closed subsets of $`C_\mathrm{w}^0(_{(+)},\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1}|_0)`$. By $`B_{\stackrel{~}{R}}(_.^0)C_\mathrm{w}^0(_{(+)},\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1}|_0)`$ we denote a closed ball of radius $`\stackrel{~}{R}`$ centered at $`_.^0`$ (with $`t`$ or $`_+`$). Furthermore, we shall need the closed subsets defined by the condition $`sup_{t0}\psi (.,t)_{\dot{H}^1}C_\psi `$ with $`C_\psi \psi _0_{\dot{H}^1}`$. In this vein, we also introduce a truncation of $`F_{.,0}(.|_0)`$, denoted $`\overline{F}_{.,0}(.|_0)`$, which for each $`t>0`$ is obtained from $`F_{t,0}(.|_0)`$ by replacing $`\mathrm{\Phi }_{t,0}^\psi [\mu _.](\zeta _0)`$ by
$$\overline{\mathrm{\Phi }}_{t,0}^\psi [\mu _.](\zeta _0):=\mathrm{min}\{1,C_\psi \mathrm{\Phi }_{t,0}^\psi [\mu _.](\zeta _0)_{\dot{H}^1}^1\}\mathrm{\Phi }_{t,0}^\psi [\mu _.](\zeta _0)$$
(90)
###### Proposition 3.4
For every $`_0\mathrm{\Gamma }_{\mathrm{B},1}`$, there exist $`\underset{ยฏ}{C}_\psi \psi _0_{\dot{H}^1}`$ and $`\underset{ยฏ}{w}>0`$, such that $`\overline{F}_{.,0}(.|_0)`$ is a Lipschitz map with Lipschitz constant $`<1`$ which maps the closed subsets of balls $`B_{\stackrel{~}{R}}(_.^0)C_\mathrm{w}^0(_+,\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1}|_0)`$ for which $`sup_{t0}\psi (.,t)_{\dot{H}^1}C_\psi `$ into themselves whenever $`\stackrel{~}{R}_.^0_\mathrm{w}`$, $`w>\underset{ยฏ}{w}`$, and $`C_\psi \underset{ยฏ}{C}_\psi \psi _0_{\dot{H}^1}`$.
By the standard contraction mapping theorem, an immediate corollary to Proposition 3.4 is the existence of a unique fixed point $`(t_t)C_\mathrm{w}^0(_+,\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1})`$, with $`(t\psi (.,t))C_b^0(_+,\dot{H}^1(^3))`$, of the fixed point equation with the truncated $`F`$,
$$_.=\overline{F}_{.,0}(_.|_0).$$
(91)
By bootstrapping regularity, fixed points of the untruncated $`F`$ will then be shown to exist and to actually be in $`C^1(,\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1})`$, furnishing unique $`\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1}`$-strong Vlasov solutions. Thus we may state our main existence and uniqueness theorem of this section.
###### Theorem 3.5
For every $`_0\mathrm{\Gamma }_{\mathrm{B},1}`$ there exists $`\underset{ยฏ}{w}>0`$ such that whenever $`w>\underset{ยฏ}{w}`$, the Vlasov fixed point equation (78) with Cauchy data $`lim_{t0}_t=_0`$ is solved by a unique curve $`t_tC_\mathrm{w}^0(,\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1})`$; since also $`\psi _0(\dot{H}^1\dot{H}^2)(^3)`$, the map $`t_t(C_\mathrm{w}^0C^1)(,\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1})`$, and thus it is the unique $`\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1}`$-strong solution to (59), (60), (64) conserving mass (66), momentum (67), angular momentum (68), and energy (69).
#### 3.3.3 Proof of Proposition 3.4
We begin with auxiliary results concerning the flow on the particle sub-phase space.
###### Lemma 3.6
Given any curve $`\zeta _.C^k(,(\dot{H}^1L^2)(^3))`$, $`k=0,1,\mathrm{}`$, we have
(i) $`๐_z^{}(.,\zeta _.)C^k(\times ^6,^6)`$,
(ii) $`๐_z^{}(.,\zeta _t)C^{\mathrm{}}(^6,^6)`$;
(iii) $`_z^{}๐_z^{}(.,\zeta _t)0`$;
(iv) $`|๐_z^{}(.,\zeta _t)|1+\varrho _{L^2}\psi (.,t)_{\dot{H}^1}`$.
Proof of Lemma 3.6: Regularity (i), (ii), and incompressibility (iii), are obvious. The bound (iv) obtains by using the triangle inequality, then $`|p|\sqrt{1+|p|^2}`$ for the momentum part, respectively for the space part the CauchyโSchwarz inequality to get $`|\varrho \psi |(x)\varrho _{L^2}\psi _{\dot{H}^1}`$ for all $`x`$; cf. (38). Q.E.D.
As a straightforward spin-off of Lemma 3.6, we have
###### Corollary 3.7
If $`\zeta _.C^k(,(\dot{H}^1L^2)(^3))`$, then $`\mathrm{\Pi }_{.,.}[\zeta _{}]C^k(\times \times ^6,^6)`$, and $`\mathrm{\Pi }_{t,t^{}}[\zeta _{}]`$ is a symplectomorphism $`t,t^{}`$; in particular, $`det_z\mathrm{\Pi }_{t,t^{}}[\zeta _{}](z)=1`$.
Proof of Corollary 3.7: This is a standard corollary. See, e.g. \[HiSm74\]. Q.E.D.
Controlling the field space component of $`\overline{F}_{.,0}`$ requires only the following Lemma:
###### Lemma 3.8
If $`\mu _.C^0(,\stackrel{~}{P}_1)`$, then $`\mathrm{\Phi }_{.,.}[\mu _.]C^0(\times \times (\dot{H}^1L^2)(^3),(\dot{H}^1L^2)(^3))`$.
Proof of Lemma 3.8: For $`\zeta _0`$ classical: straightforward calculation, for (85), (88) are quite explicit. Then apply the Hahn-Banach theorem.Q.E.D.
Proof of Proposition 3.4: We first show that, given any $`_0\mathrm{\Gamma }_{\mathrm{B},1}`$, the map $`F_{.,0}(.|_0)`$ is Lipschitz-continuous from a closed subset of $`C_\mathrm{w}^0(_+,\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1}|_0)`$, defined by the condition $`sup_{t0}\psi (.,t)_{\dot{H}^1}C_\psi `$ with $`C_\psi \underset{ยฏ}{C}_\psi \psi _0_{\dot{H}^1}`$, to $`C_\mathrm{w}^0(_+,\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1}|_0)`$ whenever $`w>\underset{ยฏ}{w}`$, with $`\underset{ยฏ}{w}`$ depending at most on $`\varrho ,C_\psi `$, and the Lipschitz constant at most on $`\varrho ,C_\psi ,w`$. We emphasize that the conditioning $`lim_{t0}_t=_0`$ and $`lim_{t0}\stackrel{~}{}_t=_0`$ implied by the definition of $`C_\mathrm{w}^0(_+,\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1}|_0)`$ do not enter our estimates.
To break up the proof into two parts, we use the triangle inequality in the form
$`F_{.,0}(\mu _.;\zeta _.|_0)F_{.,0}(\stackrel{~}{\mu }_.;\stackrel{~}{\zeta }_.|_0)_\mathrm{w}`$ $``$ $`F_{.,0}(\mu _.;\zeta _.|_0)F_{.,0}(\stackrel{~}{\mu }_.;\zeta _.|_0)_\mathrm{w}`$ (93)
$`+F_{.,0}(\stackrel{~}{\mu }_.;\zeta _.|_0)F_{.,0}(\stackrel{~}{\mu }_.;\stackrel{~}{\zeta }_.|_0)_\mathrm{w}.`$
Given $`_0=(\mu _0,\zeta _0)`$ and $`C_\psi \psi _0_{\dot{H}^1}`$, we show that, (a) given $`\zeta _.C^0(_+,(\dot{H}^1L^2)(^3))`$ satisfying $`\psi _t_{\dot{H}^1}C_\psi `$ for all $`t>0`$, for any two $`\mu _.`$ and $`\stackrel{~}{\mu }_.`$ in $`C^0(_+,\stackrel{~}{P}_1)`$ and all $`w>0`$ we have
$$F_{.,0}(\mu _.;\zeta _.|_0)F_{.,0}(\stackrel{~}{\mu }_.;\zeta _.|_0)_\mathrm{w}L_1[\varrho ;w]\underset{t0}{sup}(e^{wt}\mu _t\stackrel{~}{\mu }_t_\mathrm{L}^{}),$$
(94)
and (b), given $`\mu _.C^0(_+,\stackrel{~}{P}_1)`$, for any two $`\zeta _.`$ and $`\stackrel{~}{\zeta }_.`$ in $`C^0(_+,(\dot{H}^1L^2)(^3))`$, satisfying $`\mathrm{max}\{\psi _t_{\dot{H}^1},\stackrel{~}{\psi }_t_{\dot{H}^1}\}C_\psi `$ for all $`t>0`$, and for all $`w>\underset{ยฏ}{w}[\varrho ;C_\psi ]`$ we have
$$F_{.,0}(\mu _.;\zeta _.|_0)F_{.,0}(\mu _.;\stackrel{~}{\zeta }_.|_0)_\mathrm{w}L_2[\varrho ;w,\underset{ยฏ}{w}]\underset{t0}{sup}(e^{wt}\zeta _t\stackrel{~}{\zeta }_t_{HL});$$
(95)
for then it follows from (93), (94), (95) that, given any $`_0`$ and $`C_\psi \underset{ยฏ}{C}_\psi \psi _0_{\dot{H}^1}`$,
$$F_{.,0}(_.|_0)F_{.,0}(\stackrel{~}{}_.|_0)_\mathrm{w}L[\varrho ;w,\underset{ยฏ}{w}]_.\stackrel{~}{}_._\mathrm{w}$$
(96)
whenever $`w>\underset{ยฏ}{w}[\varrho ,C_\psi ]`$, with $`L[\varrho ;w,\underset{ยฏ}{w}]:=\mathrm{max}\{L_1[\varrho ;w],L_2[\varrho ;w,\underset{ยฏ}{w}]\}`$.
*Part a)* To prove (94), we fix $`_0`$ and $`\zeta _.`$ and note that in this case
$$F_{t,0}(\mu _.;\zeta _.|_0)F_{t,0}(\stackrel{~}{\mu }_.;\zeta _.|_0)=\mathrm{\Phi }_{t,0}[\mu _.](\zeta _0)\mathrm{\Phi }_{t,0}[\stackrel{~}{\mu }_.](\zeta _0)_{HL},$$
(97)
where, in components,
$`\mathrm{\Phi }_{t,0}[\mu _.](\zeta _0)\mathrm{\Phi }_{t,0}[\stackrel{~}{\mu }_.](\zeta _0)_{HL}`$ $`=`$ $`\mathrm{\Phi }_{t,0}^\psi [\mu _.](\zeta _0)\mathrm{\Phi }_{t,0}^\psi [\stackrel{~}{\mu }_.](\zeta _0)_{\dot{H}^1}`$ (99)
$`+\mathrm{\Phi }_{t,0}^\varpi [\mu _.](\zeta _0)\mathrm{\Phi }_{t,0}^\varpi [\stackrel{~}{\mu }_.](\zeta _0)_{L^2}.`$
Furthermore, using (85) and then the definition of $`._{\dot{H}^1}`$, we have
$$\mathrm{\Phi }_{t,0}^\psi [\mu _.](\zeta _0)\mathrm{\Phi }_{t,0}^\psi [\stackrel{~}{\mu }_.](\zeta _0)_{\dot{H}^1}^2=\left|_0^t_{\mathrm{SS}^2}(tt^{})[\varrho (\rho _t^{}\stackrel{~}{\rho }_t^{})](x^{})d\mathrm{\Omega }dt^{}\right|^2dx,$$
(100)
while with (88) and the definition of $`._{L^2}`$, we find
$`\mathrm{\Phi }_{t,0}^\varpi [\mu _.](\zeta _0)\mathrm{\Phi }_{t,0}^\varpi [\stackrel{~}{\mu }_.](\zeta _0)_{L^2}^2`$ (102)
$`={\displaystyle \left|_0^t_{\mathrm{SS}^2}\left(1+(tt^{})\mathrm{\Omega }\right)[\varrho (\rho _t^{}\stackrel{~}{\rho }_t^{})](x^{})d\mathrm{\Omega }dt^{}\right|^2dx}.`$
As to (100), triangle and Jensenโs inequalities, and Fubini, yield the estimate
$$\mathrm{\Phi }_{t,0}^\psi [\mu _.](\zeta _0)\mathrm{\Phi }_{t,0}^\psi [\stackrel{~}{\mu }_.](\zeta _0)_{\dot{H}^1}^2_{\mathrm{SS}^2}\left[_0^t(tt^{})\left|[\varrho \stackrel{}{}(\rho _t^{}^{}\stackrel{~}{\rho }_t^{})](x^{})\right|dt^{}\right]^2dxd\mathrm{\Omega }.$$
(103)
Now multiply (103) by $`e^{2wt}`$, pull $`e^{2wt}`$ under the square in rhs(103) and note that
$`{\displaystyle \left[e^{wt}_0^t(tt^{})\left|[\varrho \stackrel{}{}(\rho _t^{}^{}\stackrel{~}{\rho }_t^{})](x^{})\right|dt^{}\right]^2dx}`$ (104)
$`=`$ $`{\displaystyle [_0^t\left(e^{\frac{1}{2}w(tt^{})}(tt^{})\right)\left(e^{\frac{1}{2}w(tt^{})}e^{wt^{}}\left|[\varrho \stackrel{}{}(\rho _t^{}^{}\stackrel{~}{\rho }_t^{})](x^{})\right|\right)dt^{}]^2dx}`$ (105)
$``$ $`{\displaystyle _0^t}e^{w(tt^{\prime \prime })}(tt^{\prime \prime })^2dt^{\prime \prime }{\displaystyle _0^t}e^{w(tt^{})}e^{2wt^{}}{\displaystyle \left|[\varrho \stackrel{}{}(\rho _t^{}^{}\stackrel{~}{\rho }_t^{})](x^{})\right|^2dxdt^{}}`$ (106)
$``$ $`\frac{2}{w^4}\underset{t^{}0}{sup}(e^{2wt^{}}{\displaystyle \left|[\varrho \stackrel{}{}(\rho _t^{}^{}\stackrel{~}{\rho }_t^{})](x^{})\right|^2dx}),`$ (107)
the first inequality by CauchyโSchwarz, followed by Fubini, the second inequality by Hรถlder followed by $`_0^te^{w(tt^{\prime \prime })}(tt^{\prime \prime })^2dt^{\prime \prime }_0^te^{w(tt^{})}dt^{}_0^{\mathrm{}}e^{w\tau }\tau ^2d\tau _0^{\mathrm{}}e^{w\tau }d\tau =2/w^4`$. We next estimate the remaining $`\mathrm{d}x`$ integral by itself. For this we first rewrite it with the help of one of Greenโs identities, a change of integration variables $`xx^{}`$, and Fubiniโs theorem, exchanging the $`\mathrm{d}x^{}`$ integration with one of the convolution integrations ($`\mathrm{d}y`$, say); we then apply the KantorovichโRubinstein duality twice to obtain generalized Hรถlder estimates, then use the estimate $`\rho \stackrel{~}{\rho }_\mathrm{L}^{}\mu \stackrel{~}{\mu }_\mathrm{L}^{}`$ for $`\rho (\mathrm{d}x)=\mu (\mathrm{d}x\times ^3)`$ (similarly for $`\stackrel{~}{\rho }`$). Thus, independently of $`\mathrm{\Omega }`$, we have
$`{\displaystyle \left|\left[\varrho \stackrel{}{}(\rho _t^{}^{}\stackrel{~}{\rho }_t^{})\right](x^{})\right|^2dx}`$ $`=`$ $`{\displaystyle (\varrho \stackrel{}{}^2\varrho \stackrel{}{}(\rho _t^{}\stackrel{~}{\rho }_t^{}))(y)(\rho _t^{}\stackrel{~}{\rho }_t^{})(\mathrm{d}y)}`$ (108)
$``$ $`\mathrm{Lip}\left(\varrho \stackrel{}{}^2\varrho \stackrel{}{}(\rho _t^{}\stackrel{~}{\rho }_t^{})\right)\rho _t^{}^{}\stackrel{~}{\rho }_t^{}_\mathrm{L}^{}`$ (109)
$``$ $`\mathrm{Lip}^2\left(\varrho ^2\varrho \right)\rho _t^{}^{}\stackrel{~}{\rho }_t^{}_\mathrm{L}^{}^2`$ (110)
$``$ $`\mathrm{Lip}^2\left(\varrho ^2\varrho \right)\mu _t^{}^{}\stackrel{~}{\mu }_t^{}_\mathrm{L}^{}^2,`$ (111)
where $`\mathrm{Lip}^2\left(\varrho ^2\varrho \right)`$ is the *iterated Lipschitz constant*<sup>9</sup><sup>9</sup>9The iterated Lipschitz constant of $`f`$ is given by $`\mathrm{Lip}^2\left(f\right)=\underset{xy}{sup}\underset{\stackrel{~}{x}\stackrel{~}{y}}{sup}{\displaystyle \frac{|f(x\stackrel{~}{x})+f(y\stackrel{~}{y})f(x\stackrel{~}{y})f(y\stackrel{~}{x})|}{|xy||\stackrel{~}{x}\stackrel{~}{y}|}}.`$ If $`fC^2(^d)`$, then $`\mathrm{Lip}^2\left(f\right)=sup_{x^d}^2f(x)_{\mathrm{}}`$, where $`^2f(x)`$ is the Hessian of $`f`$ at $`x`$ and $`M_{\mathrm{}}`$ the sup norm (i.e. spectral radius) of a real symmetric matrix $`M`$. of $`\varrho ^2\varrho `$. We estimate rhs(107) with the help of (111), which in turn estimates ($`e^{2wt}`$rhs(103)) in such a way that the integration over $`\mathrm{d}\mathrm{\Omega }`$ now factors out, yielding the factor unity. Thus, taking $`sup_{t0}(e^{2wt}\mathrm{l}.\mathrm{h}.\mathrm{s}.(\text{103}))`$ and then square roots yields
$$\underset{t0}{sup}\left(e^{wt}\mathrm{\Phi }_{t,0}^\psi [\mu _.](\zeta _0)\mathrm{\Phi }_{t,0}^\psi [\stackrel{~}{\mu }_.](\zeta _0)_{\dot{H}^1}\right)\sqrt{\mathrm{Lip}^2\left(\varrho ^2\varrho \right)\frac{2}{w^4}}\underset{t0}{sup}(e^{wt}\mu _t^{}\stackrel{~}{\mu }_t_\mathrm{L}^{}).$$
(112)
As for (102), we proceed similarly, except that after the Cauchy-Schwarz and Fubini steps we now use also that
$$\left|\left(1+(tt^{})\mathrm{\Omega }\right)g(x)\right|^2dx=\left(g(x)\left(1(tt^{})^2(\mathrm{\Omega })^2\right)g(x)\right)dx$$
(113)
where $`g(x)=[\varrho \stackrel{}{}(\rho _t^{}^{}\stackrel{~}{\rho }_t^{})](x^{})`$, and obtain
$`\underset{t0}{sup}`$ $`\left(e^{wt}\mathrm{\Phi }_{t,0}^\varpi [\mu _.](\zeta _0)\mathrm{\Phi }_{t,0}^\varpi [\stackrel{~}{\mu }_.](\zeta _0)_{L^2}\right)`$ (115)
$`\sqrt{\mathrm{Lip}^2(\varrho \varrho )\frac{1}{w^2}+\mathrm{Lip}^2(\varrho \left(\mathrm{\Omega }_0\right)^2\varrho )\frac{2}{w^4}}\underset{t0}{sup}\left(e^{wt}\mu _t^{}\stackrel{~}{\mu }_t_\mathrm{L}^{}\right),`$
where $`\mathrm{\Omega }_0\mathrm{SS}^2`$ is arbitrary.
We now recall (99). Noting that, by triangle inequality, $`sup_{t0}(e^{wt}\mathrm{l}.\mathrm{h}.\mathrm{s}(\text{99}))`$ is not bigger than the sum of (112) and (115), we arrive at
$$\underset{t0}{sup}\left(e^{wt}\mathrm{\Phi }_{t,0}[\mu _.](\zeta _0)\mathrm{\Phi }_{t,0}[\stackrel{~}{\mu }_.](\zeta _0)_{HL}\right)L_1[\varrho ;w]\underset{t0}{sup}\left(e^{wt}\mu _t^{}\stackrel{~}{\mu }_t_\mathrm{L}^{}\right),$$
(116)
with
$$L_1[\varrho ;w]=\sqrt{\mathrm{Lip}^2\left(\varrho ^2\varrho \right)\frac{2}{w^4}}+\sqrt{\mathrm{Lip}^2(\varrho \varrho )\frac{1}{w^2}+\mathrm{Lip}^2(\varrho \left(\mathrm{\Omega }_0\right)^2\varrho )\frac{2}{w^4}}.$$
(117)
Finally recalling (97), we see that our proof of (94) is concluded.
*Part b)* To prove (95), we fix $`_0`$, and $`\mu _.`$ and note that in this case
$$F_{t,0}(\mu _.;\zeta _.|_0)F_{t,0}(\mu _.;\stackrel{~}{\zeta }_.|_0)=\mu _0\mathrm{\Pi }_{0,t}[\zeta _.]\mu _0\mathrm{\Pi }_{0,t}^{}[\stackrel{~}{\zeta }_.]_\mathrm{L}^{}.$$
(118)
Recalling (70) and Corollary 3.7, we note next that
$`\mu _0\mathrm{\Pi }_{0,t}[\zeta _.]\mu _0\mathrm{\Pi }_{0,t}^{}[\stackrel{~}{\zeta }_.]_\mathrm{L}^{}`$ $`=`$ $`\underset{\genfrac{}{}{0pt}{}{gC^{0,1}(^6)}{\mathrm{Lip}\left(g\right)1}}{sup}\left|{\displaystyle g\mathrm{d}(\mu _0\mathrm{\Pi }_{0,t}[\zeta _.]\mu _0\mathrm{\Pi }_{0,t}^{}[\stackrel{~}{\zeta }_.])}\right|`$ (119)
$`=`$ $`\underset{\genfrac{}{}{0pt}{}{gC^{0,1}(^6)}{\mathrm{Lip}\left(g\right)1}}{sup}\left|{\displaystyle (g\mathrm{\Pi }_{t,0}[\zeta _.]g\mathrm{\Pi }_{t,0}^{}[\stackrel{~}{\zeta }_.])d\mu _0}\right|.`$ (120)
Pulling $`|.|`$ under the last integral in (120) and using $`\mathrm{Lip}\left(g\right)1`$ gives the estimate
$$\mu _0\mathrm{\Pi }_{0,t}[\zeta _.]\mu _0\mathrm{\Pi }_{0,t}^{}[\stackrel{~}{\zeta }_.]_\mathrm{L}^{}|\mathrm{\Pi }_{t,0}[\zeta _.](z)\mathrm{\Pi }_{t,0}^{}[\stackrel{~}{\zeta }_.](z)|\mu _0(\mathrm{d}z).$$
(121)
By (81), for $`z_.,\stackrel{~}{z}_.C^1(,^6)`$ solving the characteristic equations for given $`t\psi (.,t)`$ and $`t\stackrel{~}{\psi }(.,t)`$, respectively, with<sup>10</sup><sup>10</sup>10The initial data condition $`z_0=\stackrel{~}{z}_0`$ derives from the $`_0`$ in $`F_{.,0}(.;.|_0)`$. $`z_0=z=\stackrel{~}{z}_0`$, we have
$$\mathrm{\Pi }_{t,0}[\zeta _.](z)\mathrm{\Pi }_{t,0}^{}[\stackrel{~}{\zeta }_.](z)=_0^t\left(๐_z^{}(z_\tau ,\zeta _\tau )๐_z^{}(\stackrel{~}{z}_\tau ,\stackrel{~}{\zeta }_\tau )\right)d\tau $$
(122)
We now insert (122) in the right-hand side of (121), estimate the resulting expression by pulling the absolute bars under the $`t`$-integral and applying the triangle inequality, next simplify by noting that $`๐`$ is an isometry on $`^6`$, and use Fubiniโs theorem to exchange $`\mathrm{d}\tau `$ and $`\mathrm{d}\mu _0`$ integrations. Thus we obtain the estimate
$`\mu _0\mathrm{\Pi }_{0,t}[\zeta _.]\mu _0\mathrm{\Pi }_{0,t}^{}[\stackrel{~}{\zeta }_.]_\mathrm{L}^{}`$ $`{\displaystyle _0^t}{\displaystyle \left|_z^{}(z_\tau ,\zeta _\tau )_z^{}(z_\tau ,\stackrel{~}{\zeta }_\tau )\right|\mu _0(\mathrm{d}z)d\tau }`$ (123)
$`+`$ $`{\displaystyle _0^t}{\displaystyle \left|_z^{}(z_\tau ,\stackrel{~}{\zeta }_\tau )_z^{}(\stackrel{~}{z}_\tau ,\stackrel{~}{\zeta }_\tau )\right|\mu _0(\mathrm{d}z)d\tau }.`$ (124)
Since we want an estimate for $`sup_{t0}(e^{wt}\mathrm{l}.\mathrm{h}.\mathrm{s}.(\text{124}))`$, we next consider the exponentially weighted suprema of the two integrals on r.h.s.(124) separately.
The first integral on r.h.s.(124) is estimated as follows. By (80) with $`z_\tau =(q_\tau ,p_\tau )`$ and by the CauchyโSchwarz inequality, we have
$`|_z^{}(z_\tau ,\zeta _\tau )_z^{}(z_\tau ,\stackrel{~}{\zeta }_\tau )|`$ $`=`$ $`\left|\left(\varrho \stackrel{}{}[\psi (.,\tau )\stackrel{~}{\psi }(.,\tau )]\right)(q_\tau )\right|`$ (125)
$``$ $`\varrho _{L^2}\psi (.,\tau )\stackrel{~}{\psi }(.,\tau )_{\dot{H}^1}.`$ (126)
Since r.h.s.(126) is independent of $`z`$, integration w.r.t. $`\mathrm{d}\mu _0`$ factors out and yields 1. Proceeding now similarly as in estimate (107), we obtain
$`\underset{t0}{sup}(e^{wt}{\displaystyle _0^t}\psi (.,\tau )\stackrel{~}{\psi }(.,\tau )_{\dot{H}^1}\mathrm{d}\tau )\frac{1}{w}\underset{t0}{sup}(e^{wt}\psi (.,t)\stackrel{~}{\psi }(.,t)_{\dot{H}^1}).`$ (127)
Multiplying (127) by $`\varrho _{L^2}`$ gives an upper bound for $`sup_{t0}(e^{wt}_0^t\mathrm{r}.\mathrm{h}.\mathrm{s}.(\text{126})\mathrm{d}\mu _0\mathrm{d}\tau )`$, which in turn is an upper bound for $`sup_{t0}(e^{wt}`$ first integral on r.h.s.(124)).
As to the second integral on r.h.s.(124), its integrand is rewritten using (80), with $`z_\tau =(q_\tau ,p_\tau )`$ and $`\stackrel{~}{z}_\tau =(\stackrel{~}{q}_\tau ,\stackrel{~}{p}_\tau )`$, then estimated by the triangle inequality, giving
$`|_z^{}(z_\tau ,\stackrel{~}{\zeta }_\tau )_z^{}(\stackrel{~}{z}_\tau ,\stackrel{~}{\zeta }_\tau )|`$ $``$ $`|(\varrho \stackrel{}{}\stackrel{~}{\psi }(.,\tau ))(q_\tau )(\varrho \stackrel{}{}\stackrel{~}{\psi }(.,\tau ))(\stackrel{~}{q}_\tau )|`$ (129)
$`+\left|{\displaystyle \frac{p_\tau }{\sqrt{1+|p_\tau |^2}}}{\displaystyle \frac{\stackrel{~}{p}_\tau }{\sqrt{1+|\stackrel{~}{p}_\tau |^2}}}\right|.`$
The two expressions on the right-hand side of (129) are now estimated separately.
The first term on r.h.s.(129) is estimated as follows. Spelling out the convolutions and factoring out $`\stackrel{~}{\psi }`$ in the integrand, then pulling $`|.|`$ into the convolution integral, then using the CauchyโSchwarz inequality, we get
$$|(\varrho \stackrel{~}{\psi }(.,\tau ))(q)(\varrho \stackrel{~}{\psi }(.,\tau ))(\stackrel{~}{q})|\varrho (.q)\varrho (.\stackrel{~}{q})_{L^2}\stackrel{~}{\psi }(.,\tau )_{\dot{H}^1}.$$
(130)
Now recall that for any two equi-measurable translates $`\mathrm{\Delta }_1`$ and $`\mathrm{\Delta }_2`$ of a bounded domain $`\mathrm{\Delta }`$ one has $`\chi _{_{\mathrm{\Delta }_1}}+\chi _{_{\mathrm{\Delta }_2}}\chi _{_{\mathrm{\Delta }_1\mathrm{\Delta }_2}}_{L^2}\sqrt{2|\mathrm{\Delta }|}`$, where $`\chi __\mathrm{\Delta }`$ is the characteristic function of $`\mathrm{\Delta }`$. This, the compact support of $`\varrho `$, and its Lipschitz continuity, then yield
$$|(\varrho \stackrel{}{}\stackrel{~}{\psi }(.,\tau ))(q_\tau )(\varrho \stackrel{}{}\stackrel{~}{\psi }(.,\tau ))(\stackrel{~}{q}_\tau )|C_\varrho C_\psi |q_\tau \stackrel{~}{q}_\tau |,$$
(131)
where $`C_\varrho =\sqrt{2|\mathrm{supp}(\varrho )|}\mathrm{Lip}\left(\varrho \right)`$, and where we also used $`sup_{\tau 0}\stackrel{~}{\psi }(.,\tau )_{\dot{H}^1}C_\psi `$. We next estimate $`|q_\tau \stackrel{~}{q}_\tau |`$, for which purpose we (i) use the integrated characteristic equations for $`q_\tau `$ and $`\stackrel{~}{q}_\tau `$, with $`q(0)=\stackrel{~}{q}(0)`$, (ii) pull $`|.|`$ under the time integral, (iii) use that $`_p^{}\sqrt{1+|p|^2}C_b^{0,1}(^3)`$ with $`\mathrm{Lip}\left(_p^{}\sqrt{1+|p|^2}\right)=1`$, and obtain
$`|q_\tau \stackrel{~}{q}_\tau |`$ $``$ $`{\displaystyle _0^\tau }|{\displaystyle \frac{p_\tau ^{}}{\sqrt{1+|p_\tau ^{}|^2}}}{\displaystyle \frac{\stackrel{~}{p}_\tau ^{}}{\sqrt{1+|\stackrel{~}{p}_\tau ^{}|^2}}}|\mathrm{d}\tau ^{}`$ (132)
$``$ $`{\displaystyle _0^\tau }|p_\tau ^{}\stackrel{~}{p}_\tau ^{}|d\tau ^{}.`$ (133)
Repeating steps (i) and (ii) now for $`p_\tau ^{}`$ and $`\stackrel{~}{p}_\tau ^{}`$ with $`p_0=\stackrel{~}{p}_0`$, then applying the triangle inequality to the resulting integrand, followed by applications of (131) and $`sup_{\tau ^{\prime \prime }0}\psi (.,\tau ^{\prime \prime })_{\dot{H}^1}C_\psi `$, respectively using (126), gives the string of estimates
$`|p_\tau ^{}\stackrel{~}{p}_\tau ^{}|`$ $``$ $`{\displaystyle _0^\tau ^{}}|(\varrho \stackrel{}{}\psi (.,\tau ^{\prime \prime }))(q_{\tau ^{\prime \prime }})(\varrho \stackrel{}{}\stackrel{~}{\psi }(.,\tau ^{\prime \prime }))(\stackrel{~}{q}_{\tau ^{\prime \prime }})|\mathrm{d}\tau ^{\prime \prime }`$ (134)
$``$ $`{\displaystyle _0^\tau ^{}}|(\varrho \stackrel{}{}\psi (.,\tau ^{\prime \prime }))(q_{\tau ^{\prime \prime }})(\varrho \stackrel{}{}\psi (.,\tau ^{\prime \prime }))(\stackrel{~}{q}_{\tau ^{\prime \prime }})|\mathrm{d}\tau ^{\prime \prime }`$ (136)
$`+{\displaystyle _0^\tau ^{}}|(\varrho \stackrel{}{}\psi (.,\tau ^{\prime \prime })\varrho \stackrel{}{}\stackrel{~}{\psi }(.,\tau ^{\prime \prime }))(\stackrel{~}{q}_{\tau ^{\prime \prime }})|\mathrm{d}\tau ^{\prime \prime }`$
$``$ $`C_\varrho C_\psi {\displaystyle _0^\tau ^{}}|q_{\tau ^{\prime \prime }}\stackrel{~}{q}_{\tau ^{\prime \prime }}|\mathrm{d}\tau ^{\prime \prime }+\varrho _{L^2}{\displaystyle _0^\tau ^{}}\psi (.,\tau ^{\prime \prime })\stackrel{~}{\psi }(.,\tau ^{\prime \prime })_{\dot{H}^1}\mathrm{d}\tau ^{\prime \prime },`$ (137)
with $`C_\varrho `$ given below (131), and $`C_\psi \psi _0_{\dot{H}^1}`$ chosen later. Inserting (137) into (133), recalling from (124) that $`\tau t`$, then employing a second order variant of the Gronwall lemma (see Appendix A.2) with $`z_0=\stackrel{~}{z}_0`$, we find for all $`\tau t`$ that
$$|q_\tau \stackrel{~}{q}_\tau |\varrho _{L^2}_0^\tau \mathrm{cosh}\left[\underset{ยฏ}{w}(\tau \tau ^{})\right]_0^\tau ^{}\psi (.,\tau ^{\prime \prime })\stackrel{~}{\psi }(.,\tau ^{\prime \prime })_{\dot{H}^1}\mathrm{d}\tau ^{\prime \prime }\mathrm{d}\tau ^{},$$
(138)
with $`\underset{ยฏ}{w}=\sqrt{C_\varrho C_\psi }`$. With (138) and (131) we have the relevant estimates for the first term on the right-hand side of (129).
To estimate the second term on r.h.s.(129), we recall that $`_p^{}\sqrt{1+|p|^2}C_b^{0,1}(^3)`$ with $`\mathrm{Lip}\left(_p^{}\sqrt{1+|p|^2}\right)=1`$, then recall (137). Estimating $`|q_\tau \stackrel{~}{q}_\tau |`$ by r.h.s.(138) and inserting this estimate into (137), we now find that for all $`\tau t`$ we have
$`\left|{\displaystyle \frac{p_\tau }{\sqrt{1+|p_\tau |^2}}}{\displaystyle \frac{\stackrel{~}{p}_\tau }{\sqrt{1+|\stackrel{~}{p}_\tau |^2}}}\right|`$ $``$ $`\varrho _{L^2}{\displaystyle _0^\tau }\psi (.,\tau ^{})\stackrel{~}{\psi }(.,\tau ^{})_{\dot{H}^1}\mathrm{d}\tau ^{}`$ (140)
$`+\varrho _{L^2}\underset{ยฏ}{w}^2{\displaystyle _0^\tau }{\displaystyle _0^\tau ^{}}\mathrm{cosh}\left[\underset{ยฏ}{w}(\tau ^{}\tau ^{\prime \prime })\right]{\displaystyle _0^{\tau ^{\prime \prime }}}\psi (.,\stackrel{ห}{\tau })\stackrel{~}{\psi }(.,\stackrel{ห}{\tau })_{\dot{H}^1}\mathrm{d}\stackrel{ห}{\tau }\mathrm{d}\tau ^{\prime \prime }\mathrm{d}\tau ^{},`$
which provides an upper bound to the second term on r.h.s.(129).
The bounds on the two terms of r.h.s.(129), i.e. (138) with (131), and (140), combine into an estimate of l.h.s.(129) which is independent of $`z_\tau `$ and $`\stackrel{~}{z}_\tau `$, the solutions to the characteristic equations for given fields $`\psi `$ and $`\stackrel{~}{\psi }`$ with initial data $`z_0=z=\stackrel{~}{z}_0`$, respectively. We have, independently of $`z`$,
$`|_z^{}(z_\tau ,\stackrel{~}{\zeta }_\tau )_z^{}(\stackrel{~}{z}_\tau ,\stackrel{~}{\zeta }_\tau )|`$ $``$ $`\varrho _{L^2}{\displaystyle _0^\tau }\psi (.,\tau ^{})\stackrel{~}{\psi }(.,\tau ^{})_{\dot{H}^1}\mathrm{d}\tau ^{}`$ (143)
$`+\varrho _{L^2}\underset{ยฏ}{w}^2{\displaystyle _0^\tau }\mathrm{cosh}\left[\underset{ยฏ}{w}(\tau \tau ^{})\right]{\displaystyle _0^\tau ^{}}\psi (.,\tau ^{\prime \prime })\stackrel{~}{\psi }(.,\tau ^{\prime \prime })_{\dot{H}^1}\mathrm{d}\tau ^{\prime \prime }\mathrm{d}\tau ^{}`$
$`+\varrho _{L^2}\underset{ยฏ}{w}^2{\displaystyle _0^\tau }{\displaystyle _0^\tau ^{}}\mathrm{cosh}\left[\underset{ยฏ}{w}(\tau ^{}\tau ^{\prime \prime })\right]{\displaystyle _0^{\tau ^{\prime \prime }}}\psi (.,\stackrel{ห}{\tau })\stackrel{~}{\psi }(.,\stackrel{ห}{\tau })_{\dot{H}^1}\mathrm{d}\stackrel{ห}{\tau }\mathrm{d}\tau ^{\prime \prime }\mathrm{d}\tau ^{}.`$
We integrate (143) w.r.t. $`\mu _0(\mathrm{d}z)`$; due to the $`z`$-independence of the integrand on the right-hand side that integral factors out there and equals 1. It thus remains to integrate (143) w.r.t. $`\mathrm{d}\tau `$ from $`0`$ to $`t`$, to multiply by $`e^{wt}`$ and to take the supremum over $`t0`$. The three terms on r.h.s.(143) are estimated by repeating the strategy used in (107) a total of 9 times (however, one of the estimates is just (127) again). For all $`w>\underset{ยฏ}{w}`$ we thereby arrive at the desired estimate for $`sup_{t0}(e^{wt}`$ second integral on r.h.s.(124)),
$`\underset{t0}{sup}(e^{wt}{\displaystyle _0^t}{\displaystyle \left|_z^{}(z_\tau ,\stackrel{~}{\zeta }_\tau )_z^{}(\stackrel{~}{z}_\tau ,\stackrel{~}{\zeta }_\tau )\right|d\mu _0d\tau })`$ (144)
$`\varrho _{L^2}({\displaystyle \frac{1}{w^2}}+{\displaystyle \frac{\underset{ยฏ}{w}^2}{2w^2(w\underset{ยฏ}{w})}}+{\displaystyle \frac{\underset{ยฏ}{w}^2}{2w^3(w\underset{ยฏ}{w})}})\underset{t0}{sup}(e^{wt}\psi (.,t)\stackrel{~}{\psi }(.,t)_{\dot{H}^1}),`$ (145)
with $`\underset{ยฏ}{w}=\sqrt{C_\varrho C_\psi }`$ from (138). The estimates given by (145) and by (127) (and ensuing text), together with (124), now give the estimate
$$\underset{t0}{sup}(e^{wt}\mu _0\mathrm{\Pi }_{0,t}[\zeta _.]\mu _0\mathrm{\Pi }_{0,t}^{}[\stackrel{~}{\zeta }_.]_\mathrm{L}^{})L_2\underset{t0}{sup}(e^{wt}\psi (.,t)\stackrel{~}{\psi }(.,t)_{\dot{H}^1})$$
(146)
whenever $`w>\underset{ยฏ}{w}`$, with
$$L_2[\varrho ;w,\underset{ยฏ}{w}]=(\frac{1}{w}+\frac{1}{w^2}+\frac{(1+w)\underset{ยฏ}{w}^2}{2w^3(w\underset{ยฏ}{w})})\varrho _{L^2}.$$
(147)
Since $`sup_{t0}(e^{wt}\psi (.,\tau )\stackrel{~}{\psi }(.,\tau )_{\dot{H}^1})sup_{t0}(e^{wt}\zeta _t\stackrel{~}{\zeta }_t_{HL})`$, and because of (118), we see that (146) proves (95). Part b) is completed.
We have thus proved that $`F_{.,0}(.|_0)`$ is a Lipschitz map from a closed subset of $`C_\mathrm{w}^0(_+,\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1}|_0)`$, defined by the condition $`sup_{t0}\psi (.,t)_{\dot{H}^1}C_\psi `$, to $`C_\mathrm{w}^0(_+,\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1}|_0)`$ whenever $`C_\psi \underset{ยฏ}{C}_\psi \psi _0_{\dot{H}^1}`$ and $`w>\underset{ยฏ}{w}=\sqrt{C_\varrho C_\psi }`$. The Lipschitz constant is $`L[\varrho ;w,\underset{ยฏ}{w}]=\mathrm{max}\{L_1[\varrho ;w],L_2[\varrho ;w,\underset{ยฏ}{w}]\}`$, with $`L_1`$ given in (117) and $`L_2`$ in (147). Finally, we note that everything proven so far for $`F_{.,0}(.|_0)`$ remains valid for its truncation $`\overline{F}_{.,0}(.|_0)`$, obtained for all $`t>0`$ by replacing $`\mathrm{\Phi }_{t,0}^\psi [\mu _.]`$ by its upper truncation $`\overline{\mathrm{\Phi }}_{t,0}^\psi [\mu _.]`$ given in (90). By time reversal symmetry, the same conclusions hold with $`C_\mathrm{w}^0(_+,\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1}|_0)`$ replaced by $`C_\mathrm{w}^0(,\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1}|_0)`$. This concludes the Lipschitz continuity part of the proof.
We show next that for sufficiently large $`\underset{ยฏ}{C}_\psi \psi _0_{\dot{H}^1}`$ and $`\stackrel{~}{R}_.^0_\mathrm{w}`$, the map $`F_{.,0}(.|_0)`$ sends a closed subset of a closed ball $`B_{\stackrel{~}{R}}(_.^0)C_\mathrm{w}^0(_+,\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1}|_0)`$, satisfying $`sup_{t0}\psi (.,t)_{\dot{H}^1}C_\psi `$, into $`B_{\stackrel{~}{R}}(_.^0)`$ whenever $`C_\psi \underset{ยฏ}{C}_\psi `$ and $`w>\underset{ยฏ}{w}=\sqrt{C_\varrho C_\psi }`$; recall that $`_.^0=F_{.,0}(0_.|_0)C_\mathrm{w}^0(,\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1}|_0)`$ denotes the free evolution of $`_0(\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1})`$, where $`0_.`$ is the trivial constant map $`(t0)C_\mathrm{w}^0(,\stackrel{~}{M}_1(^6)\dot{H}^1(^3)L^2(^3))`$. Note that l.h.s.(96) is well-defined if we substitute $`0_.`$ for $`\stackrel{~}{}_.`$; clearly, $`0_.C_\mathrm{w}^0(,\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1}|_0)`$. However, since $`\mu _.0_.P_1P_1`$, r.h.s.(96) cannot be directly used to estimate this particular version of l.h.s.(96). One possible way out is to derive an analog of (96) for more general measures involving our extension $`._{\stackrel{~}{\mathrm{L}^{}}}`$ of the dual Lipschitz distance, see Appendix A.1. A more direct way out is as follows. We invoke the triangle inequality to estimate $`\mathrm{\Phi }_{t,0}[\mu _.](\zeta _0)\mathrm{\Phi }_{t,0}[0_.](\zeta _0)_{HL}\mathrm{\Phi }_{t,0}[\mu _.](\zeta _0)\mathrm{\Phi }_{t,0}[\mu _.^0](\zeta _0)_{HL}+\mathrm{\Phi }_{t,0}[\mu _.^0](0)_{HL}`$, where we used that $`\mathrm{\Phi }_{t,0}[\mu _.^0](\zeta _0)\mathrm{\Phi }_{t,0}[0_.](\zeta _0)=\mathrm{\Phi }_{t,0}[\mu _.^0](0)`$, and also to estimate
$`\mu _0\mathrm{\Pi }_{0,t}[\zeta _.]\mu _0\mathrm{\Pi }_{0,t}^{}[0_.]_\mathrm{L}^{}`$ $``$ $`\mu _0\mathrm{\Pi }_{0,t}[\zeta _.]\mu _0\mathrm{\Pi }_{0,t}^{}[\zeta _.^0]_\mathrm{L}^{}`$ (148)
$`+\mu _0\mathrm{\Pi }_{0,t}[\zeta _.^0]\mu _0\mathrm{\Pi }_{0,t}^{}[0_.]_\mathrm{L}^{}`$
Recalling that $`F_{.,0}(0_.|_0)=_.^0`$, it is straightforward to verify that (96) modifies to
$$F_{.,0}(_.|_0)_.^0_\mathrm{w}L_._.^0_\mathrm{w}+K.$$
(149)
whenever $`w>\underset{ยฏ}{w}=\sqrt{C_\varrho C_\psi }`$, with $`L[\varrho ;w,\underset{ยฏ}{w}]=\mathrm{max}\{L_1[\varrho ;w],L_2[\varrho ;w,\underset{ยฏ}{w}]\}`$ as before, and
$$K=\underset{t0}{sup}\left(e^{wt}\left(\mathrm{\Phi }_{t,0}[\mu _.^0](0)_{HL}+\mu _0\mathrm{\Pi }_{0,t}[\zeta _.^0]\mu _0\mathrm{\Pi }_{0,t}^{}[0_.]_\mathrm{L}^{}\right)\right).$$
(150)
Since $`\mathrm{\Phi }_{t,0}^\psi [\mu _.^0](0)_{\dot{H}^1}\varrho _{L^2}t`$ and $`\mathrm{\Phi }_{t,0}^\varpi [\mu _.^0](0)_{L^2}\varrho _{L^2}t`$ by (33), and since
$$\underset{t0}{sup}\left(e^{wt}\mu _0\mathrm{\Pi }_{0,t}[\zeta _.^0]\mu _0\mathrm{\Pi }_{0,t}^{}[0_.]_\mathrm{L}^{}\right)L_2\underset{t0}{sup}\left(e^{wt}\mathrm{\Phi }_{t,0}^\psi [\mu _.^0](\zeta _0)_{\dot{H}^1}\right)$$
(151)
by (146), with $`\mathrm{\Phi }_{t,0}^\psi [\mu _.^0](\zeta _0)_{\dot{H}^1}(2_\mathrm{W}(\zeta _0))^{1/2}+\varrho _{L^2}t`$ (by (33) again), we find that
$$K\frac{1}{w}\frac{2}{e}\varrho _{L^2}+L_2\mathrm{max}\{\frac{1}{w}\varrho _{L^2},(2_\mathrm{W}(\zeta _0))^{1/2}\}.$$
(152)
Recalling that $`L[\varrho ;w,\underset{ยฏ}{w}]=\mathrm{max}\{L_1[\varrho ;w],L_2[\varrho ;w,\underset{ยฏ}{w}]\}`$ with $`L_1`$ given in (117) and $`L_2`$ in (147), we see that both $`L`$ and $`K`$ are monotonically decreasing functions of $`w(>\underset{ยฏ}{w})`$, with asymptotic decay to zero $`1/w`$ for large $`w`$. Now let $`w`$ be large enough such that $`L1/2`$. Pick an $`R_{}`$, independent of $`w`$, such that $`KLR_{}`$ (clearly such an $`R_{}`$ exists). Now, either $`_._.^0_{\stackrel{~}{\mathrm{w}}}R_{}`$ or $`R_{}_._.^0_{\stackrel{~}{\mathrm{w}}}`$. In the former case, $`F_{.,0}(_.|_0)_.^0_\mathrm{w}2LR_{}`$, i.e. $`F_{.,0}(.|_0)`$ maps any closed subset of the closed ball $`B_R_{}(_.^0)`$ satisfying $`sup_{t0}\psi (.,t)_{\dot{H}^1}C_\psi `$ into $`B_{2LR_{}}(_.^0)`$; clearly, since by assumption $`w`$ is large enough so that $`L1/2`$, the ball $`B_{2LR_{}}(_.^0)B_R_{}(_.^0)`$. In the latter case on the other hand we have $`F_{.,0}(_.|_0)_.^0_\mathrm{w}2L_._.^0_{\stackrel{~}{\mathrm{w}}}`$, with $`L1/2`$. Thus we conclude that, as long as $`\stackrel{~}{R}R_{}`$, the fixed point map $`F_{.,0}(.|_0)`$ sends any closed subset of $`B_{\stackrel{~}{R}}(_.^0)`$ which satisfies $`sup_{t0}\psi (.,t)_{\dot{H}^1}C_\psi `$ into $`B_{\stackrel{~}{R}}(_0)`$, given that $`L1/2`$. Again, we note that everything proven in this paragraph for $`F_{.,0}(.|_0)`$ remains valid for its truncation $`\overline{F}_{.,0}(.|_0)`$.
It remains to notice that the truncated map $`\overline{F}_{.,0}(.|_0)`$ sends any closed subset of $`C_\mathrm{w}^0(_+,\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1}|_0)`$ satisfying $`sup_{t0}\psi (.,t)_{\dot{H}^1}C_\psi `$, with $`C_\psi \psi _0_{\dot{H}^1}`$, to itself. Hence, for $`\stackrel{~}{R}R_{}`$ and $`\underset{ยฏ}{C}_\psi \psi _0_{\dot{H}^1}`$, the truncated map $`\overline{F}_{.,0}(.|_0)`$ sends the intersection of any closed ball $`B_{\stackrel{~}{R}}(_.^0)C_\mathrm{w}^0(_+,\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1}|_0)`$ with any closed subsets of $`C_\mathrm{w}^0(_+,\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1}|_0)`$ satisfying $`sup_{t0}\psi (.,t)_{\dot{H}^1}C_\psi `$ to itself whenever $`C_\psi \underset{ยฏ}{C}_\psi `$ and $`w>\underset{ยฏ}{w}=\sqrt{C_\varrho C_\psi }`$ are large enough so that $`L1/2`$. For instance, this can be achieved as follows. Setting $`w=2\underset{ยฏ}{w}`$, the Lipschitz constant becomes
$$L=\mathrm{max}\{\sqrt{\frac{\mathrm{Lip}^2(\varrho ^2\varrho )}{16\underset{ยฏ}{w}^4}}+\sqrt{\frac{\mathrm{Lip}^2(\varrho \varrho )}{4\underset{ยฏ}{w}^2}+\frac{\mathrm{Lip}^2(\varrho (\mathrm{\Omega }_0){}_{}{}^{2}\varrho )}{16\underset{ยฏ}{w}^4}},\left(1+\frac{1}{2\underset{ยฏ}{w}}\right)\frac{5\varrho _{L^2}}{8\underset{ยฏ}{w}}\},$$
(153)
and also setting now $`C_\psi =\underset{ยฏ}{C}_\psi `$ so that $`\underset{ยฏ}{w}=\sqrt{C_\varrho \underset{ยฏ}{C}_\psi }`$, we see that there is a unique $`\underset{ยฏ}{C}_\psi ^{}[\varrho ]`$ such that r.h.s.(153)$`1/2`$ for $`\underset{ยฏ}{C}_\psi \underset{ยฏ}{C}_\psi ^{}[\varrho ]`$; hence, choosing $`\underset{ยฏ}{C}_\psi \mathrm{max}\{\underset{ยฏ}{C}_\psi ^{}[\varrho ],\psi _0_{\dot{H}^1}\}`$ will do. This completes the proof of Proposition 3.4. Q.E.D.
#### 3.3.4 Proof of Theorem 3.5
It suffices in the following to continue to work with $`C_\psi =\underset{ยฏ}{C}_\psi `$ and $`w=2\underset{ยฏ}{w}`$, as done at the end of the proof of Proposition 3.4. Thus, we need to show that Theorem 3.5 is true for sufficiently large $`\underset{ยฏ}{C}_\psi \mathrm{max}\{\underset{ยฏ}{C}_\psi ^{}[\varrho ],\psi _0_{\dot{H}^1}\}`$.
Proof of Theorem 3.5: Let $`C_\psi =\underset{ยฏ}{C}_\psi \mathrm{max}\{\underset{ยฏ}{C}_\psi ^{}[\varrho ],\psi _0_{\dot{H}^1}\}`$, and $`w=2\underset{ยฏ}{w}`$. Then, as shown in the proof of Proposition 3.4, $`\overline{F}_{.,0}(_.|_0)`$ is a contraction map, with Lipschitz constant $`L1/2`$, from the intersection of any closed ball $`B_{\stackrel{~}{R}}(_.^0)C_\mathrm{w}^0(,\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1}|_0)`$ of radius $`\stackrel{~}{R}R_{}`$ with any closed subset of $`C_\mathrm{w}^0(,\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1}|_0)`$ defined by the condition $`sup_{t0}\psi (.,t)_{\dot{H}^1}C_\psi (=\underset{ยฏ}{C}_\psi )`$, to itself. The standard contraction mapping theorem now guarantees the existence of a unique fixed point $`t_tC_\mathrm{w}^0(,\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1})`$, with $`t\psi (.,t)C_b^0(,\dot{H}^1(^3))`$, of the fixed point equation with the truncated $`F`$,
$$_.=\overline{F}_{.,0}(_.|_0).$$
(154)
We will now show that for sufficiently large $`\underset{ยฏ}{C}_\psi `$ the solutions of (154) are fixed points for $`F`$, moreover of type $`C^1(,\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1})`$, thus furnishing unique $`\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1}`$-strong Vlasov solutions.
To this effect, choose $`\underset{ยฏ}{C}_\psi >\mathrm{max}\{\underset{ยฏ}{C}_\psi ^{}[\varrho ],\sqrt{2_\mathrm{W}(\zeta _0)},\sqrt{4+4E_08E_{}}\}`$, where $`_\mathrm{W}(\zeta _0)`$ is the initial field energy, $`E_0=(_0)`$ is the total energy of the initial state, and where $`E_{}`$ is the ground state energy of the $`N`$-body Hamiltonian given in (39); note that the ground state energy is $`N`$-independent and therefore identical to the ground state of the continuum (Vlasov) limit energy functional (69). Note that automatically we have $`\underset{ยฏ}{C}_\psi >\psi _0_{\dot{H}^1}^{}`$, for it is trivially obvious that $`\psi _0_{\dot{H}^1}^{}\sqrt{2_\mathrm{W}(\zeta _0)}`$ and easily shown that $`\psi _0_{\dot{H}^1}^{}\sqrt{4+4E_08E_{}}`$. With the so chosen $`\underset{ยฏ}{C}_\psi `$, there then exists at least a small neighborhood of $`t=0`$ such that for all $`t`$ in this neighborhood, the fixed points of (154), which are of type $`C^0(,\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1})`$, by continuity satisfy
$$_t=F_{t,0}(_.|_0).$$
(155)
Now recall that $`\psi _0(\dot{H}^1\dot{H}^2)(^3)`$, ensuring a strong solution of the wave equation, by the HilleโYosida theorem. Next recall the remark after Lemma 2.4; viz.
$$\psi (.,t)_{\dot{H}^1}(2_\mathrm{W}(\zeta _0))^{1/2}+\varrho _{L^2}|t|,$$
(156)
for the strong solution of the wave equation given any subluminal source $`\varrho \rho C^0(,C_0^{\mathrm{}}(^3))`$. Clearly, there is a unique $`\underset{ยฏ}{T}>0`$ for which
$$\underset{ยฏ}{C}_\psi =(2_\mathrm{W}(\zeta _0))^{1/2}+\varrho _{L^2}\underset{ยฏ}{T},$$
(157)
such that $`\psi (.,t)_{\dot{H}^1}<\underset{ยฏ}{C}_\psi `$ strictly for all $`|t|<\underset{ยฏ}{T}`$, by (156). This now implies that there exist $`T\underset{ยฏ}{T}`$ such that the fixed point $`_.`$ of (154) satisfies (155) for all $`t[T,T]`$. We now show that $`sup\{T:(\text{155})`$ holds for all $`|t|T\}=\mathrm{}`$.
Thus, suppose that $`sup\{T:(\text{155})`$ holds for all $`|t|T\}=T_{}<\mathrm{}`$. Then for either $`t=T_{}^+`$ or $`t=T_{}^{}`$ (or both), $`_t`$ is given by (154) but not by (155). For the sake of concreteness, assume that this is so for some $`t`$ in a right neighborhood of $`T_{}`$. This then means that for all $`t(T_{},T_{}+ฯต)`$ we have $`\mathrm{\Phi }_{t,0}^\psi _{\dot{H}^1}\underset{ยฏ}{C}_\psi >\sqrt{4+4E8E_{}}`$, which in particular implies that $`lim_{tT_{}}\mathrm{\Phi }_{t,0}^\psi _{\dot{H}^1}\underset{ยฏ}{C}_\psi >\sqrt{4+4E8E_{}}`$. On the other hand, for all $`t[T_{},T_{}]`$, the solution $`_.`$ of (154) satisfies (155), and $`\zeta _.`$ is a strong solution of the wave equation. But then, by Corollary 3.7, $`t_tC^1([T_{},T_{}],\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1})`$ is a $`\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1}`$-strong solution, which by Theorem 3.3 conserves energy. As a consequence of energy conservation, for all $`t[T_{},T_{}]`$, and in particular for $`t=T_{}`$, we have
$$\mathrm{\Phi }_{t,0}^\psi _{\dot{H}^1}\sqrt{4+4E8E_{}}.$$
(158)
Since $`t\mathrm{\Phi }_{t,0}^\psi C^0(,\dot{H}^1(^3))`$, we thus have a contradiction to the previously concluded $`\mathrm{\Phi }_{t,0}^\psi _{\dot{H}^1}\underset{ยฏ}{C}_\psi >\sqrt{4+4E8E_{}}`$ for $`t>T_{}`$. Hence, $`T_{}=\mathrm{}`$. Q.E.D.
###### Remark 3.9
By the proof of Theorem 3.5 it suffices to work with $`w2\underset{ยฏ}{w}^{}`$, where
$$\underset{ยฏ}{w}^{}=\sqrt{C_\varrho \mathrm{max}\{\underset{ยฏ}{C}_\psi ^{}[\varrho ],\sqrt{4+4E8E_{}}\}}.$$
(159)
###### Remark 3.10
In the proof of Theorem 3.5 we only made use of the a priori bound (158) following from Theorem 3.3 and the analog of the proof of our Proposition 2.8 for the regularized Vlasov model. The other bounds expressed in Proposition 2.8 are carried over as follows.
Let $`t(t)C^0(,\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1})`$ be a generalized solution of the wave gravity Vlasov equations which conserves energy $`E`$, momentum $`P`$, and angular momentum $`J`$, and of course mass $`M=1`$. Then, beside (158), uniformly in $`t`$ we have
$`\varpi (.,t)_{L^2}^2`$ $``$ $`2E2E_{},`$ (160)
$`{\displaystyle \sqrt{1+|p|^2}\mu _t(\mathrm{d}z)}`$ $``$ $`1+EE_{},`$ (161)
$`6E_{}3E3{\displaystyle }(\varrho \psi (.,t))(x)\mu _t(\mathrm{d}z)`$ $``$ $`E1`$ (162)
Moreover, (57) and (58) extend to
$$|p\mu _t(\mathrm{d}z)||P|+\psi ^{\left(N\right)}(.,t)_{\dot{H}^1}\varpi ^{\left(N\right)}(.,t)_{L^2},$$
(163)
$$|p\times x\mu _t(\mathrm{d}z)||J|+(R+|t|)\psi ^{\left(N\right)}(.,t)_{\dot{H}^1}\varpi ^{\left(N\right)}(.,t)_{L^2}.$$
(164)
## 4 The limit $`N\mathrm{}`$
We prove first that the $`\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1}`$-strong $`N`$-body generalized solutions of the Vlasov model converge $`._\mathrm{w}`$-strongly to solutions when $`N\mathrm{}`$. We then specify when these limit solutions are continuum solutions. Finally we discuss the probabilistic import in terms of a law of large numbers and a central limit theorem.
### 4.1 The $`\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1}`$-strong limit of the $`N`$-body generalized solutions
Suppose the family of initial empirical measures converges $`\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1}`$-strongly when $`N\mathrm{}`$, written $`\epsilon [๐ณ_0^{\left(N\right)}](\mathrm{d}z)\mu _0(\mathrm{d}z)`$. Then the microscopic โdensityโ $`\rho ^{\left(N\right)}(.,0)`$, as given in (18) with $`t=0`$, converges strongly (in the marginal measuresโ subspace) to the โdensityโ $`\rho (.,0)`$. Finally, assume that $`\zeta [๐ณ_0^{\left(N\right)}]\zeta _0\dot{H}^1(^3)L^2(^3)`$ satisfying (61), (62) with $`\psi _0(\dot{H}^1\dot{H}^2)(^3)`$. Our goal is to show that, when $`N\mathrm{}`$, the generalized solution $`t(\epsilon [๐ณ_t^{\left(N\right)}];\zeta _t^{\left(N\right)})(C_\mathrm{w}^0C^1)(,\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1})`$ associated with this converging family of initial data in turn converges in $`._\mathrm{w}`$ norm to a solution $`t(\mu _t;\zeta _t)(C_\mathrm{w}^0C^1)(,\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1})`$ of the regularized wave gravity Vlasov fixed point equation. In the following, $`_t^{\left(N\right)}\stackrel{\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1}}{}_t`$ means $`\psi ^{\left(N\right)}(.,t)\stackrel{\dot{H}^1}{}\psi (.,t)`$ satisfying (61), $`\varpi ^{\left(N\right)}(.,t)\stackrel{L^2}{}\varpi (.,t)`$ satisfying (62), and $`\epsilon [๐ณ_t^{\left(N\right)}]\mu _t`$ in $`\stackrel{~}{P}_1`$.
The main result is an immediate consequence of the following theorem, which states that the $`._\mathrm{w}`$ induced distance between any two $`C_\mathrm{w}^0(,\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1})`$ solutions of our Vlasov fixed point equation is controlled by the $`\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1}`$ distance of their initial states in $`\mathrm{\Gamma }_{\mathrm{B},1}`$.
###### Proposition 4.1
Let $`๐`$ or $`๐`$ be an index set, and let $`\{_.^{(\alpha )}C_\mathrm{w}^0(,\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1})\}_{\alpha ๐}`$ be a family of solutions of the Vlasov fixed point equation (78), having initial data $`_0^{(\alpha )}\mathrm{\Gamma }_{\mathrm{B},1}`$ with $`\psi _0^{(\alpha )}(\dot{H}^1\dot{H}^2)(^3)`$, for which $`E^{}:=sup_{\alpha ๐}\{(_0^{(\alpha )})\}`$ exists. Define
$$\overline{w}=\sqrt{C_\varrho \mathrm{max}\{\underset{ยฏ}{C}_\psi ^{}[\varrho ],\sqrt{4+4E^{}8E_{}}\}},$$
(165)
with $`C_\varrho =\sqrt{2|\mathrm{supp}(\varrho )|}\mathrm{Lip}\left(\varrho \right)`$ (cf. text below (131)). Then for all $`w2\overline{w}`$ there exists a constant $`L_0[\overline{w}]`$ such that for any $`(\alpha ,\stackrel{~}{\alpha })๐^2`$,
$$_.^{(\alpha )}_.^{(\stackrel{~}{\alpha })}_\mathrm{w}L_0_0^{(\alpha )}_0^{(\stackrel{~}{\alpha })}.$$
(166)
Before we prove this proposition, we state and prove its main corollary.
###### Theorem 4.2
Let $`t_t^{\left(N\right)}(C_\mathrm{w}^0C^1)(,\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1})`$ be the $`\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1}`$-strong $`N`$-body solution of the regularized wave gravity Vlasov equations (59), (60), (64) with Cauchy data $`_0^{\left(N\right)}=lim_{t0}_t^{\left(N\right)}`$ described in Theorem 3.3. Suppose $`_0^{\left(N\right)}\stackrel{\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1}}{}_0`$, with $`_0`$ having mass $`M(=1)`$, energy $`E`$, momentum $`P`$, and angular momentum $`J`$, and with $`\psi _0(\dot{H}^1\dot{H}^2)(^3)`$. Then $`_.^{\left(N\right)}_._\mathrm{w}0`$, where $`t_tC_\mathrm{w}^0(,\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1})`$ is the unique solution of (78) described in Theorem 3.5. Beside mass, $`t_t`$ also conserves energy, momentum, and angular momentum. Furthermore, since $`\psi _0(\dot{H}^1\dot{H}^2)(^3)`$, we also have $`_.C^1(,\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1})`$.
Proof of Theorem 4.2: By Theorem 3.5, there exist unique type $`C_\mathrm{w}^0(,\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1})`$ solutions $`_.^{\left(N\right)},_.`$ of the fixed point equation (78) for each Cauchy data $`_0^{\left(N\right)},_0\mathrm{\Gamma }_{\mathrm{B},1}`$, with $`\psi _0^{\left(N\right)},\psi _0`$ in $`(\dot{H}^1\dot{H}^2)(^3)`$, respectively. The latter restriction upgrades the solutions to the wave equation to be strong, which by Lemma 3.6 implies solutions of type $`(C_\mathrm{w}^0C^1)(,\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1})`$ which conserve mass, momentum, angular momentum, and energy.
Let $`_.^{(\mathrm{})}_.`$, and set $`๐=\{\mathrm{}\}`$. Since energy is conserved by each solution, $`E^{}=sup_{\alpha ๐}\{(_0^{(\alpha )})\}`$ exists. Thus, $`\overline{w}`$ exists. Pick any $`w>2\overline{w}`$. Now Proposition 4.1 applies to our family $`\{_.^{(\alpha )}\}_{\alpha ๐}`$, and since $`_0^{\left(N\right)}_00`$ by hypothesis, Proposition 4.1 now implies that $`_.^{\left(N\right)}_.^{(\mathrm{})}_\mathrm{w}0`$. Q.E.D.
To prepare the proof of Proposition 4.1, we will need the following lemmata.
###### Lemma 4.3
Let $`\zeta _.C_b^0(,(\dot{H}^1L^2)(^3))`$, with $`sup_{t0}\psi (.,t)_{\dot{H}^1}C_\psi `$, and let $`\underset{ยฏ}{w}=\sqrt{C_\varrho C_\psi }`$. Then $`\mathrm{\Pi }_{t,t^{}}[\zeta _{}]C^{0,1}(^6,^6)`$, with Lipschitz constant<sup>11</sup><sup>11</sup>11Incidentally, by Lemma 4.3, the largest Liapunov exponent for $`\mathrm{\Pi }_{t,t^{}}[\zeta _{}]`$ is bounded above by $`\underset{ยฏ}{w}`$.
$$\mathrm{Lip}\left(\mathrm{\Pi }_{t,t^{}}[\zeta _{}]\right)=\frac{1}{\sqrt{2}}(2+\mathrm{max}\{\underset{ยฏ}{w},\mathrm{\hspace{0.17em}1}/\underset{ยฏ}{w}\})e^{\underset{ยฏ}{w}|tt^{}|}.$$
(167)
Proof of Lemma 4.3: Let $`\psi (.,t)C_b^0(,\dot{H}^1(^3))`$ be given, with $`sup_{t0}\psi (.,t)_{\dot{H}^1}C_\psi `$. To unburden notation, let $`\psi _t(.)`$ stand for $`\psi (.,t)`$. Let $`tz_t^6`$ and $`t\stackrel{~}{z}_t^6`$ be two distinct solutions of (8), (9) for this $`\psi _.`$. Proceeding analogously to the steps taken in (133) and (137), this time for $`\psi =\stackrel{~}{\psi }`$, but now allowing $`z_0\stackrel{~}{z}_0`$, we find
$`|\stackrel{~}{q}_tq_t|`$ $``$ $`|\stackrel{~}{q}_t^{}q_t^{}|+{\displaystyle _t^{}^t}|{\displaystyle \frac{\stackrel{~}{p}_\tau }{\sqrt{1+|\stackrel{~}{p}_\tau |^2}}}{\displaystyle \frac{p_\tau }{\sqrt{1+|p_\tau |^2}}}|\mathrm{d}\tau `$ (168)
$``$ $`|\stackrel{~}{q}_t^{}q_t^{}|+{\displaystyle _t^{}^t}|\stackrel{~}{p}_\tau p_\tau |d\tau ,`$ (169)
respectively
$`|\stackrel{~}{p}_tp_t|`$ $``$ $`|\stackrel{~}{p}_t^{}p_t^{}|+{\displaystyle _t^{}^t}|\left(\varrho \stackrel{}{}\psi _\tau \right)(\stackrel{~}{q}_\tau )\left(\varrho \stackrel{}{}\psi _\tau \right)(q_\tau )|d\tau `$ (170)
$``$ $`|\stackrel{~}{p}_t^{}p_t^{}|+C_\varrho C_\psi {\displaystyle _t^{}^t}|\stackrel{~}{q}_\tau q_\tau |d\tau ,`$ (171)
with $`C_\varrho `$ and $`C_\psi `$ as stated in the lemma. Inserting (171) into (169) and using the second order variant of Gronwallโs lemma gives
$$|\stackrel{~}{q}_tq_t||\stackrel{~}{q}_t^{}q_t^{}|\mathrm{cosh}\left[\underset{ยฏ}{w}(tt^{})\right]+|\stackrel{~}{p}_t^{}p_t^{}|\frac{1}{\underset{ยฏ}{w}}\mathrm{sinh}\left[\underset{ยฏ}{w}|tt^{}|\right],$$
(172)
with $`\underset{ยฏ}{w}=\sqrt{C_\varrho C_\psi }`$. Back-substituting (172) into (171) and integrating then gives
$$|\stackrel{~}{p}_tp_t||\stackrel{~}{p}_t^{}p_t^{}|\mathrm{cosh}\left[\underset{ยฏ}{w}(tt^{})\right]+|\stackrel{~}{q}_t^{}q_t^{}|\underset{ยฏ}{w}\mathrm{sinh}\left[\underset{ยฏ}{w}|tt^{}|\right].$$
(173)
To get from (172) and (173) to the conclusion of Lemma 4.3, use $`\mathrm{cosh}(x)e^{|x|}`$ and $`\mathrm{sinh}(|x|)e^{|x|}/2`$, as well as the familiar $`\stackrel{}{v}_2\stackrel{}{v}_1\sqrt{2}\stackrel{}{v}_2`$ for $`\stackrel{}{v}^n`$. Q.E.D.
The next lemma transfers control about the flow $`\mathrm{\Pi }_{t,t^{}}`$ on $`^6`$ to the flow $`\mathrm{\Pi }_{t,t^{}}^{}`$ on $`\stackrel{~}{P}_1`$.
###### Lemma 4.4
For any symplectomorphism $`\mathrm{\Pi }`$ on $`^6`$ which in addition is a Lipschitz map with Lipschitz constant $`\mathrm{\Lambda }`$, the adjoint map $`\mathrm{\Pi }^{}:M(^6)M(^6)`$, defined by $`\mathrm{\Pi }^{}(\sigma ):=\sigma \mathrm{\Pi }^1`$, is a positivity- and $`._{\mathrm{TV}}`$-preserving smooth automorphism of $`M(^6)`$, and it is a Lipschitz homeomorphism on $`\stackrel{~}{M}_1(^6)`$ for $`._{\stackrel{~}{\mathrm{L}^{}}}`$ with Lipschitz constant $`\mathrm{\Lambda }`$.
Proof of Lemma 4.4: First, since $`\mathrm{\Pi }`$ is a symplectomorphism, by way of the definition of its adjoint, $`\mathrm{\Pi }^{}`$ maps $`M(^6)`$ smoothly onto $`M(^6)`$, and it preserves (a) $`\sigma _{\mathrm{TV}}`$ for $`\sigma M`$ and (b) the positivity of $`\mu M_+`$. Furthermore, since $`\mathrm{\Pi }`$ is invertible, so is $`\mathrm{\Pi }^{}`$.
To see that $`\mathrm{\Pi }^{}`$ is a homeomorphism of $`\stackrel{~}{M}_1(^6)`$, we only need to show that $`\mathrm{\Pi }^{}`$ maps $`\stackrel{~}{M}_1`$ into $`\stackrel{~}{M}_1`$. Thus, let $`z_{}^6`$ be the unique element of $`\mathrm{ker}\mathrm{\Pi }`$. Then note that by the definition of $`\mathrm{\Pi }^{}`$ and the Lipschitz property of $`\mathrm{\Pi }`$ we have
$$||z|\sigma \mathrm{\Pi }^1(\mathrm{d}z)|=||\mathrm{\Pi }(z)\mathrm{\Pi }(z_{})|\sigma (\mathrm{d}z)|\mathrm{\Lambda }|zz_{}||\sigma |(\mathrm{d}z),$$
(174)
where $`|\sigma |`$ is the total variation of $`\sigma `$; the last integral exists for $`\sigma \stackrel{~}{M}_1`$.
As for the Lipschitz continuity of the adjoint flow, let $`\widehat{\sigma },\stackrel{ห}{\sigma }\stackrel{~}{M}_1(^6)`$. We have
$`\mathrm{\Pi }^{}(\widehat{\sigma })\mathrm{\Pi }^{}(\stackrel{ห}{\sigma })_\mathrm{L}^{}`$ $`=`$ $`\underset{gC^{0,1}(^6)}{sup}\left\{\right|{\displaystyle }g\mathrm{d}(\widehat{\sigma }\mathrm{\Pi }^1\stackrel{ห}{\sigma }\mathrm{\Pi }^1)|:\mathrm{Lip}\left(g\right)1\}`$ (175)
$`=`$ $`\underset{gC^{0,1}(^6)}{sup}\left\{\right|{\displaystyle }g\mathrm{\Pi }\mathrm{d}(\widehat{\sigma }\stackrel{ห}{\sigma })|:\mathrm{Lip}\left(g\right)1\}`$ (176)
$`=`$ $`\mathrm{\Lambda }\underset{gC^{0,1}(^6)}{sup}\left\{\right|{\displaystyle }{\displaystyle \frac{1}{\mathrm{\Lambda }}}g\mathrm{\Pi }\mathrm{d}(\widehat{\sigma }\stackrel{ห}{\sigma })|:\mathrm{Lip}\left(g\right)1\}`$ (177)
$``$ $`\mathrm{\Lambda }\widehat{\sigma }\stackrel{ห}{\sigma }_\mathrm{L}^{}.`$ (178)
In the last step, we used that $`\mathrm{\Lambda }^1g\mathrm{\Pi }C^{0,1}(^6)`$ with $`\mathrm{Lip}\left(\mathrm{\Lambda }^1g\mathrm{\Pi }\right)1`$. Q.E.D.
Proof of Proposition 4.1: Pick $`w>2\overline{w}`$, with $`\overline{w}`$ defined in (165), and pick $`_.,\stackrel{~}{}_.C_\mathrm{w}^0(,\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1})`$ from the family of solutions $`_.^{(\alpha )}`$ of the Vlasov fixed point equation (78) specified in Proposition 4.1. Then
$$_.\stackrel{~}{}_._\mathrm{w}=F_{.,0}(_.|_0)F_{.,0}(\stackrel{~}{}_.|\stackrel{~}{}_0)_\mathrm{w}.$$
(179)
By the triangle inequality,
$`F_{.,0}(_.|_0)F_{.,0}(\stackrel{~}{}_.|\stackrel{~}{}_0)_\mathrm{w}`$ $``$ $`F_{.,0}(_.|_0)F_{.,0}(_.|\stackrel{~}{}_0)_\mathrm{w}`$ (181)
$`+F_{.,0}(_.|\stackrel{~}{}_0)F_{.,0}(\stackrel{~}{}_.|\stackrel{~}{}_0)_\mathrm{w}.`$
Now, $`F_{.,0}(_.|\stackrel{~}{}_0)F_{.,0}(\stackrel{~}{}_.|\stackrel{~}{}_0)_\mathrm{w}`$ was estimated already in the proof of Proposition 3.4, see (96) (recall that the conditioning $`lim_{t0}_t=_0=lim_{t0}\stackrel{~}{}_t`$ that entered the statement of Proposition 3.4 did not enter the estimates for (96) themselves). Furthermore, with $`w>2\overline{w}`$ it follows that the parameter conditions in the proof of Theorem 3.5 are fulfilled for each $`\stackrel{~}{}_0`$; hence, in (96) we have $`L[\varrho ;w,\underset{ยฏ}{w}]1/2`$ for each $`\stackrel{~}{}_0`$. Thus, by (179), (181), and (96), and with $`(1L[\varrho ;w,\underset{ยฏ}{w}])^12`$, we arrive at the estimate
$$_.\stackrel{~}{}_._\mathrm{w}2F_{.,0}(_.|_0)F_{.,0}(_.|\stackrel{~}{}_0)_\mathrm{w},$$
(182)
uniformly for all $`\stackrel{~}{}_0\{_0^{(\alpha )}\}_{\alpha ๐}`$.
The proof of Proposition 4.1 has thus been reduced to proving Lipschitz continuity of $`F_{.,0}`$ in its second argument, given the first. Since, by the triangle inequality,
$`F_{.,0}(_.|\mu _0;\zeta _0)F_{.,0}(_.|\stackrel{~}{\mu }_0;\stackrel{~}{\zeta }_0)_\mathrm{w}`$ $``$ $`F_{.,0}(_.|\mu _0;\zeta _0)F_{.,0}(_.|\stackrel{~}{\mu }_0;\zeta _0)_\mathrm{w}`$ (184)
$`+F_{.,0}(_.|\stackrel{~}{\mu }_0;\zeta _0)F_{.,0}(_.|\stackrel{~}{\mu }_0;\stackrel{~}{\zeta }_0)_\mathrm{w},`$
it suffices to show that for given $`_.`$ and $`\zeta _0`$, we have
$$F_{.,0}(_.|\mu _0;\zeta _0)F_{.,0}(_.|\stackrel{~}{\mu }_0;\zeta _0)_\mathrm{w}L_1^{}\mu _0\stackrel{~}{\mu }_0_\mathrm{L}^{},$$
(185)
and for given $`_.`$ and $`\stackrel{~}{\mu }_0`$,
$$F_{.,0}(_.|\stackrel{~}{\mu }_0;\zeta _0)F_{.,0}(_.|\stackrel{~}{\mu }_0;\stackrel{~}{\zeta }_0)_\mathrm{w}L_2^{}\zeta _0\stackrel{~}{\zeta }_0_{HL},$$
(186)
with $`L_1^{},L_2^{}`$ depending at most on $`\overline{w}`$. For then it follows from (184), (185), (186) that
$$F_{.,0}(_.|_0)F_{.,0}(_.|\stackrel{~}{}_0)_\mathrm{w}L^{}_0\stackrel{~}{}_0,$$
(187)
with $`L^{}[\overline{w}]:=\mathrm{max}\{L_1^{},L_2^{}\}`$, completing the proof of Proposition 4.1, with $`L_0=2L^{}`$.
As to (185), for all $`t`$ we have
$`F_{t,0}(_.|\mu _0;\zeta _0)F_{t,0}(_.|\stackrel{~}{\mu }_0;\zeta _0)`$ $`=`$ $`\mathrm{\Pi }_{t,0}^{}[\zeta _.](\mu _0)\mathrm{\Pi }_{t,0}^{}[\zeta _.](\stackrel{~}{\mu }_0)_\mathrm{L}^{}`$ (188)
$``$ $`\frac{2+\mathrm{max}\{\overline{w},1/\overline{w}\}}{\sqrt{2}}e^{\overline{w}|t|}\mu _0\stackrel{~}{\mu }_0_\mathrm{L}^{},`$ (189)
the inequality by Lemma 4.3 and Lemma 4.4. Since $`w2\overline{w}`$, the $`sup_t\left(e^{w|t|}(\text{189})\right)`$ exists. Estimating it further with $`w\overline{w}\overline{w}`$ for $`w2\overline{w}`$ now gives (185), with
$$L_1^{}[\overline{w}]=\frac{2+\mathrm{max}\{\overline{w},1/\overline{w}\}}{\sqrt{2}}.$$
(190)
As to (186), for all $`t`$ we have
$$F_{t,0}(_.|\stackrel{~}{\mu }_0;\zeta _0)F_{t,0}(_.|\stackrel{~}{\mu }_0;\stackrel{~}{\zeta }_0)=\mathrm{\Phi }_{t,0}^\psi [0_.](\zeta _0\stackrel{~}{\zeta }_0)_{\dot{H}^1}+\mathrm{\Phi }_{t,0}^\varpi [0_.](\zeta _0\stackrel{~}{\zeta }_0)_{L^2},$$
(191)
where $`\mathrm{\Phi }_{t,0}^\psi [0_.](.)`$ and $`\mathrm{\Phi }_{t,0}^\varpi [0_.](.)`$ are the free propagators obtained from Kirchhoffโs formulas (85) and (88) by replacing $`\mu _.0_.`$; note that $`\mathrm{\Phi }_{t,0}^\psi [0_.](.)`$ and $`\mathrm{\Phi }_{t,0}^\varpi [0_.](.)`$ are linear operators. For initial data $`\psi _0(\dot{H}^1\dot{H}^2)(^3)`$, the freely propagating wave is a $`HL`$-strong solution of the homogeneous wave equation and its field energy $`_\mathrm{W}(\zeta _.^{\mathrm{free}})`$ is conserved. This implies the bounds
$$\mathrm{\Phi }_{t,0}^\psi [0_.](\zeta _0\stackrel{~}{\zeta }_0)_{\dot{H}^1}+\mathrm{\Phi }_{t,0}^\varpi [0_.](\zeta _0\stackrel{~}{\zeta }_0)_{L^2}2\sqrt{_\mathrm{W}(\zeta _0\stackrel{~}{\zeta }_0)}\sqrt{2}\zeta _0\stackrel{~}{\zeta }_0_{HL}.$$
(192)
Hence,
$$L_2^{}=\sqrt{2}.$$
(193)
Estimates (185) with (189), and (186) with (193) now combine to
$$F_{.,0}(_.|\mu _0;\zeta _0)F_{.,0}(_.|\stackrel{~}{\mu }_0;\stackrel{~}{\zeta }_0)_\mathrm{w}L^{}_0\stackrel{~}{}_0$$
(194)
with
$$L^{}[\overline{w}]=\mathrm{max}\{\sqrt{2},\frac{2+\mathrm{max}\{\overline{w},1/\overline{w}\}}{\sqrt{2}}\}$$
(195)
for all $`w>2\overline{w}`$.
The proof of Proposition 4.1 is complete, with $`L_0[\overline{w}]=2L^{}[\overline{w}]`$. Q.E.D.
### 4.2 The continuum limit
Note that so far nothing prevents the measure $`\mu _0(\mathrm{d}z)`$, which obtains as limit of the $`\epsilon [๐ณ_0^{\left(N\right)}](\mathrm{d}z)`$ when $`N\mathrm{}`$, from being as singular as the $`\epsilon [๐ณ_0^{\left(N\right)}](\mathrm{d}z)`$ are. In particular, we may even allow $`\epsilon [๐ณ_0^{\left(N\right)}](\mathrm{d}z)\delta _{z_0}(\mathrm{d}z)`$. Since in physical applications of Vlasov theory one is typically interested in continuum solutions, we now suppose that when $`N\mathrm{}`$, the familiy of initial empirical measures $`\epsilon [๐ณ_0^{\left(N\right)}](\mathrm{d}z)`$ converges $`\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1}`$-strongly to a measure $`\mu _0(\mathrm{d}z)`$ which is absolutely continuous w.r.t. Lebesgue measure. We write $`\mu (\mathrm{d}z)=\mu ^f(\mathrm{d}z)=f(z)\mathrm{d}z`$ for the absolutely continuous measures in $`P_1(^6)`$. The set of their RadonโNikodym derivatives is denoted $`L_{+,1}^{1,1}(^6)`$; thus $`fL_{+,1}^{1,1}(^6)`$. We now show that when $`\mu _0(\mathrm{d}z)=\mu ^{f_0}(\mathrm{d}z)`$, then $`\mu _t=\mu _t^f`$, with $`f(.,.,t)L_{+,1}^{1,1}(^6)`$ for all $`t`$.
###### Proposition 4.5
If $`(\mu _.,\zeta _.)(C_\mathrm{w}^0C^1)(,\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1})`$ solves the Vlasov fixed point equation (78) with $`\mu _0=\mu ^{f_0}`$, $`f_0(L_{+,1}^{1,1}L^p)(^6)`$ for some $`p1`$, then $`\mu _.=\mu ^{f(.,.,t)}`$ with $`f(.,.,t)(L_{+,1}^{1,1}L^p)(^6)`$ for all $`t`$; note that $`p1`$ includes the case that $`f_0L_{+,1}^{1,1}(^6)`$ while $`f_0L^p(^6)`$ for any $`p>1`$.
Proof of Proposition 4.5: Suppose $`\mu _.C^1(,\stackrel{~}{P}_1)`$ is a strong generalized solution of the Vlasov continuity equation (64) for given $`\zeta _.(C_b^0C^1)(,(\dot{H}^1L^2)(^3))`$, with Cauchy data $`\mu _0=\mu ^{f_0}`$, $`f_0(L_{+,1}^{1,1}L^p)(^6)`$ for some $`p1`$. Then $`\mu _.=\mu ^{f(.,.,t)}`$ with $`f(.,.,t)(L_{+,1}^{1,1}L^p)(^6)`$ for all $`t`$. But this follows from the definition of a generalized solution, a straightforward change of variables from $`z`$ to $`\mathrm{\Pi }_{t,0}[\zeta _.](z)`$ under the integral, noting the properties of the flow $`\mathrm{\Pi }_{.,.}`$ summarized in Corollary 3.7. Q.E.D.
#### 4.2.1 Additional conservation laws for continuum solutions
The argument used to prove Proposition 4.5 has the useful corollary that continuum solutions $`_.^f`$ with $`f_0(L_{+,1}^{1,1}L^p)(^6)`$ for some $`p>1`$ enjoy additional conservation laws. Here we wrote $`_.^f`$ for $`_.=(\mu _.,\zeta _.)`$ with $`\mu _t=\mu ^{f(.,.,t)}`$.
For any $`g:_+`$, we define the $`g`$-Casimir functional of $`^f`$ by
$$๐^{(g)}\left(^f\right)=gfdz,\mathrm{whenever}gfL^1(^6).$$
(196)
For $`g=\mathrm{id}`$, we obtain the mass functional (66) for absolutely continuous $`\mu _t=\mu _t^f`$; for $`g(.)=(\mathrm{id}(.))^p`$, $`p>1`$, we get the $`p`$-th power of the $`L^p`$ norm of $`f`$; the case $`g(.)=\mathrm{id}(.)\mathrm{log}(\mathrm{id}(.)/f_{})`$, gives the entropy of $`f`$ relative to $`f_{}`$,
$$๐^{(\mathrm{id}\mathrm{log}(\mathrm{id}/f_{}))}\left(^f\right)=f\mathrm{ln}(f/f_{})dz๐ฎ(f|f_{});$$
(197)
here, $`f_{}L_{+,1}^{1,1}(^6)`$ is an otherwise arbitrary probability density function. In particular, $`๐ฎ(f|f_{})`$ is well-defined if $`f(L_{+,1}^{1,1}L^{1+ฯต})(^6)`$ for some $`ฯต>0`$.
###### Proposition 4.6
Let $`t_t(C_\mathrm{w}^0C^1)(,\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1})`$ be a generalized solution of the regularized wave gravity Vlasov model for which $`\mu _t=\mu _t^f`$ is absolutely continuous. Then, beside the conservation laws (75), (76), (77), whenever $`๐^{(g)}\left(_0^f\right)`$ exists, also
$$๐^{(g)}\left(_.^f\right)=๐^{(g)}\left(_0^f\right).$$
(198)
In particular, if $`f_0(L_{+,1}^{1,1}L^{1+ฯต})(^6)`$ for some $`ฯต>0`$, then the relative entropy of $`f(.,.,t)L_{+,1}^{1,1}(^6)`$ is conserved, i.e.
$$๐ฎ(f|f_{})=๐ฎ(f_0|f_{}).$$
(199)
### 4.3 Law of large numbers and central limit theorem
#### 4.3.1 The law of large numbers
Convergence in KR topology of probability measures $`\mu ^{\left(N\right)}`$ as $`N\mathrm{}`$ implies the convergence in probability of a family of random variables with laws $`\mu ^{\left(N\right)}`$; see \[Dud02\]. Since at time $`t`$ the empirical measures $`\epsilon [๐ณ_t^{\left(N\right)}]`$ do converge in KR topology to $`\mu _t`$ if they do so at $`t=0`$, our theorem about the $`N\mathrm{}`$ limit of the $`N`$-body generalized solutions to the regularized wave gravity Vlasov model is equivalent to the following law of large numbers.
###### Theorem 4.7
For $`N`$, let $`\mu _0^{\left(N\right)}P_1(^6)`$ be given, with $`\mathrm{supp}(\mu _0(\mathrm{d}x\times ^3))B_R`$, and suppose $`\mu _0^{\left(N\right)}\mu _0P_1(^6)`$. Moreover, let $`๐ณ_0^{\left(N\right)}^{6N}`$ be a random vector whose components in $`^6`$ are (not necessarily independent) random variables $`z_1^{\left(N\right)}(0),\mathrm{},z_N^{\left(N\right)}(0)`$ with common law $`\mu _0^{\left(N\right)}`$. To each $`๐ณ_0^{\left(N\right)}`$ assign a unique $`\zeta [๐ณ_0^{\left(N\right)}]((\dot{H}^1\dot{H}^2)L^2)(^3)`$ satisfying (61),(62), such that $`\zeta [๐ณ_0^{\left(N\right)}]\zeta _0=(\psi (.,0),\varpi (.,0))((\dot{H}^1\dot{H}^2)L^2)(^3)`$ $`HL`$-strongly when $`N\mathrm{}`$, with $`\psi (.,0)`$, $`\varpi (.,0)`$ satisfying (61), (62).
Let $`(\mu _.,\zeta _.)(C_\mathrm{w}^0C^1)(,\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1})`$ be the unique generalized strong solution of the regularized wave gravity Vlasov equations for initial data $`(\mu _0,\zeta _0)`$. Let $`t๐ณ_t^{\left(N\right)}^{6N}`$ be the path in particle phase space whose $`N`$ components in $`^6`$ jointly solve the EinsteinโNewton equations of motion (19), (20) with the initial data $`๐ณ_0^{\left(N\right)}`$, and with $`t\psi ^{\left(N\right)}(.,t)`$ and $`t\varpi ^{\left(N\right)}(.,t)`$ solving the regularized wave gravity equations (16) and (17) with initial data $`\zeta [๐ณ_0^{\left(N\right)}]`$. Then, for any $`gC^{0,1}(^6)`$ and for all $`t`$, in probability we have
$$\frac{1}{N}\underset{i=1}{\overset{N}{}}g\left(z_i^{\left(N\right)}(t)\right)\stackrel{N\mathrm{}}{}g(z)\mu _t(\mathrm{d}z).$$
(200)
###### Remark 4.8
By invoking the extremal decomposition of permutation invariant probability measures on $`(^6)^{}`$ \[HeSa55\], our LLN (200) can be generalized to
$$\left(\genfrac{}{}{0pt}{}{N}{n}\right)^1\underset{1i_1<\mathrm{}<i_nN}{}g(z_{i_1}^{\left(N\right)}(t),\mathrm{},z_{i_n}^{\left(N\right)}(t))\stackrel{N\mathrm{}}{}g(z_1,\mathrm{},z_n)\mu _t^{\times n}(\mathrm{d}z_1\mathrm{}\mathrm{d}z_n)$$
(201)
for any permutation invariant $`gC^{0,1}(^{6n})`$ and for all $`t`$.
###### Remark 4.9
One actually should also allow the field initial data $`\zeta _0^{\left(N\right)}`$ to be random variables independently of the particle random variables for each $`N`$, but this would require a whole new setup involving probability measures on field space, the choice of an adequate topology on that space, beyond what has been developed in this paper.
Even though our LLN does not demand that $`\mu _.`$ be a continuum solution, in applications this is typically so. While our LLN implies that the Vlasov continuum approximation to the sampling of $`N`$ body systems becomes exact for all $`t`$ in the limit of infinite $`N`$, for any particular physical system $`N`$ is fixed and may vary only from system to system. Thus, take $`\mu _0^{\left(N\right)}=\mu _0`$ for all $`N`$, with $`\mu _0=\mu ^{f_0}`$. By assumption $`\epsilon [๐ณ_0^{\left(N\right)}]\mu _0`$ when $`N\mathrm{}`$. Yet for any finite $`N`$ we have $`\epsilon [๐ณ_0^{\left(N\right)}]\mu _0_\mathrm{L}^{}>0`$, and then our estimates of section 4.1 show that at any other time $`t`$ we only have $`\epsilon [๐ณ_t^{\left(N\right)}]\mu _t_\mathrm{L}^{}e^{C|t|}\epsilon [๐ณ_0^{\left(N\right)}]\mu _0_\mathrm{L}^{}`$. Hence, we can only conclude that the physical mean values l.h.s.(200) at time $`t`$ can be computed in acceptable approximation by their Vlasov continuum analog, i.e. r.h.s.(200) with $`\mu _t=\mu ^{f(.,.,t)}`$, if $`|t|`$ is โnot too large,โ a notion which depends on $`N`$ and on how good the approximation is initially.
#### 4.3.2 The central limit theorem
Having obtained the law of large numbers, we next inquire into the fluctuations around the deterministic mean. Our goal is to derive a central limit theorem for the dynamical variables $`(๐ณ_t^{\left(N\right)},\zeta _t^{\left(N\right)})`$, which are random variables through their dependence on the random initial data for the $`N`$ particles, viz. $`๐ณ_0^{\left(N\right)}`$.
We adapt the technique of \[BrHe77\], who studied fluctuations of particle motions for non-relativistic Vlasov equations. This is done in two steps. First we study the differences of (primarily) test particle motions generated by the finite $`N`$ versus the continuum flows, and of similar type differences of field evolutions. The attribute โprimarilyโ in parentheses refers to the fact that almost all (w.r.t. Lebesgue measure) initial conditions launch a test particle evolution, with finitely many exceptions which are upgraded and follow the proper evolution. In our case all these particle evolutions are generated by the adjoint flows on particle phase space of the Vlasov flows on $`\stackrel{~}{P}_1`$ whose fixed points are the proper Vlasov evolutions (both of course coupled to the same wave gravity equations). We prove the convergence of the characteristic function of the fluctuation process to a Gaussian characteristic function in a suitable norm. In the second step we extract from this analysis the fluctuations for the proper evolutions.
We recall that initial data $`๐ณ_0^{\left(N\right)}=(z_1^{\left(N\right)}(0),\mathrm{},z_N^{\left(N\right)}(0))^{6N}`$ with $`q_k^{\left(N\right)}(0)B_R`$ for $`k=1,2,\mathrm{},N`$ uniquely define initial data $`(\epsilon [๐ณ_0^{\left(N\right)}],\zeta [๐ณ_0^{\left(N\right)}])\mathrm{\Gamma }_{\mathrm{B},1}`$ which launch the unique strong solution $`(\epsilon [๐ณ_.^{\left(N\right)}],\zeta _.^{\left(N\right)})(C_\mathrm{w}^0C^1)(,\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1})`$ of our regularized wave gravity Vlasov equations. By (81), (82), the solution $`(\epsilon [๐ณ_.^{\left(N\right)}],\zeta _.^{\left(N\right)})`$ generates a single particle flow $`\mathrm{\Pi }_{.,.}[\zeta _.^{\left(N\right)}](.)`$ on $`^6`$ giving single particle evolutions
$$z_t(z_0;๐ณ_0^{\left(N\right)})=\mathrm{\Pi }_{t,0}[\zeta _.^{\left(N\right)}](z_0),$$
(202)
and a flow $`\mathrm{\Phi }_{.,.}[\epsilon [๐ณ_.^{\left(N\right)}]](.)`$ on field space giving field evolutions
$$\zeta _t(\zeta _0;๐ณ_0^{\left(N\right)})=\mathrm{\Phi }_{t,0}[\epsilon [๐ณ_.^{\left(N\right)}]](\zeta _0),$$
(203)
with data $`z_0^6`$ and $`\zeta _0(\dot{H}^1L^2)(^3)`$; note, however, that we are *exclusively* considering data $`(z_0,\zeta _0)\mathrm{\Gamma }_\mathrm{B}^{(1)}`$. As for the notation, by the r.h.s. of (202) the dynamics $`z_.`$ depends on $`z_0`$ and on $`\zeta _.^{\left(N\right)}`$, but the latter in turn is implicitly fixed by $`๐ณ_0^{\left(N\right)}`$ (and our Vlasov equations); similarly the notation in (203) is explained. As for their dynamical significance, the dynamics $`z_.`$ solves the *characteristic equations* of the Vlasov continuity equation given the fields $`\zeta _.^{\left(N\right)}`$. For almost all data $`z_0^6`$ this is a *test particle* dynamics, the exception being when $`z_0\{z_1^{\left(N\right)}(0),\mathrm{},z_N^{\left(N\right)}(0)\}`$, in which case $`z_.`$ coincides with one of the components $`z_k^{\left(N\right)}(.)`$ of the unique solution $`(๐ณ_.^{\left(N\right)},\zeta _.^{\left(N\right)})`$ of equations (16), (17), (18), (19), (20) with Cauchy data $`(๐ณ_0^{\left(N\right)},\zeta [๐ณ_0^{\left(N\right)}])`$. The wave dynamics $`\zeta _.`$ in turn solves the linear inhomogeneous wave equation given the source term $`\varrho \rho _.^{\left(N\right)}`$ obtained from $`\epsilon [๐ณ_.^{\left(N\right)}]`$ Note that the dynamics of $`z_.`$ and $`\zeta _.`$ are in general independent of each other; the exception occurs when $`\zeta _0=\zeta [๐ณ_0^{\left(N\right)}]`$, in which case $`\zeta _.=\zeta _.^{\left(N\right)}`$.
When $`N\mathrm{}`$ such that $`\epsilon [๐ณ_t^{\left(N\right)}]\epsilon [๐ณ_t^{(\mathrm{})}]`$ and $`\zeta _t^{\left(N\right)}\zeta _t^{(\mathrm{})}`$ ($`HL`$-strongly) for all $`t`$, the flows for (202), (203) converge to the corresponding flows generated by the Vlasov solution $`t(\epsilon [๐ณ_t^{(\mathrm{})}],\zeta _t^{(\mathrm{})})`$; note that flows analogous to those for (202), (203) are defined for any admissible solution $`(C_\mathrm{w}^0C^1)(,\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{B},1})`$ of our Vlasov equations. The discussion of the finite-$`N`$ fluctuations amounts to analyzing the difference of the evolutions (202), (203) for finite $`N`$ versus $`N=\mathrm{}`$. Note that in either case the Vlasov dynamics that generates (202) and (203) is determined by $`๐ณ_0^{\left(N\right)}`$, respectively $`๐ณ_0^{(\mathrm{})}`$.
As to our notation, above we use the symbol $`\zeta _.^{(\mathrm{})}`$ to distinguish the limit fields of $`\zeta _.^{\left(N\right)}`$ from the fields $`\zeta _.`$ solving (203) in general. However, eventually we will choose $`\zeta _0\zeta _0^{(\mathrm{})}`$ for the sake of concreteness in the CLT. Note furthermore that $`๐ณ_0^{(\mathrm{})}`$ could be anything from a single point $`z_{}`$ to a continuous function, i.e. we may have $`๐ณ_0^{(\mathrm{})}=f_0(z)`$ (in distribution). Therefore, $`\epsilon [๐ณ_0^{(\mathrm{})}]`$ could be any probability measure $`\mu _0P_1`$ with $`\mathrm{supp}(\mu _0(\mathrm{d}x\times ^3))B_R`$; in particular, if $`๐ณ_0^{(\mathrm{})}=f_0(z)`$ is an empirical continuum density, the empirical continuum measure $`\mu ^{f_0}(\mathrm{d}z)=\epsilon [๐ณ_0^{(\mathrm{})}](\mathrm{d}z)`$.
We stipulate further notation. Recall that $`(z,\zeta )=๐ท\mathrm{\Gamma }_\mathrm{B}^{(1)}`$ denotes generic points in $`\mathrm{\Gamma }_\mathrm{B}^{(1)}`$. In this vein, for solutions of (202), (203) we use whenever possible the shorthand
$$(z_t,\zeta _t)=๐ท_t$$
(204)
but when necessary to discuss the components in more detail, we write
$$(z_t,\zeta _t)=(q_t,p_t,\psi _t,\varpi _t).$$
(205)
Note that $`๐ท_t`$ is a function of $`๐ท_0`$ and $`๐ณ_0^{\left(N\right)}`$; in components, $`q_t`$ and $`p_t`$ are points in $`^3`$ which are functions of $`z_0`$ and $`๐ณ_0^{\left(N\right)}`$ while $`\psi _t`$ and $`\varpi _t`$ are points in $`\dot{H}^1\dot{H}^2`$, respectively $`L^2`$, which are functions of $`\zeta _0`$ and $`๐ณ_0^{\left(N\right)}`$. We also recycle some of our previously stipulated abbreviations. Thus, $`..\mu (\mathrm{d}z)`$ stands for $`_^6..\mu (\mathrm{d}x\mathrm{d}p)`$ and $`..\mu (\mathrm{d}z)\nu (\mathrm{d}z^{})`$ for $`_{^6\times ^6}..\mu (\mathrm{d}x\mathrm{d}p)\nu (\mathrm{d}x^{}\mathrm{d}p^{})`$. So much for notation.
As a technical primer which will allow us to perform estimates needed for the main theorem, we will first show that $`๐ท_t(๐ท_0,๐ณ_0^{\left(N\right)})`$ is regular as a function of $`๐ณ_0^{\left(N\right)}`$, with bounded derivatives. Since $`๐ณ_0^{\left(N\right)}`$ is uniquely identified with the empirical measure $`\epsilon [๐ณ_0^{\left(N\right)}]`$ and similarly $`๐ณ_0^{(\mathrm{})}`$ is uniquely identified with the empirical measure $`\epsilon [๐ณ_0^{(\mathrm{})}]`$ (with the possibility $`\epsilon [๐ณ_0^{(\mathrm{})}](\mathrm{d}z)=\mu ^{f_0}(\mathrm{d}z)`$, the continuum case), we really mean regularity of $`๐ท_t`$ as a function of $`\epsilon [๐ณ_0^{\left(N\right)}]`$, respectively $`\epsilon [๐ณ_0^{(\mathrm{})}]`$, with bounded derivatives w.r.t. $`\epsilon [๐ณ_0^{\left(N\right)}]`$ (respectively $`\epsilon [๐ณ_0^{(\mathrm{})}]`$). Actually, it will be necessary to take derivatives not just restricted to the subset $`\stackrel{~}{P}_1`$ of $`\stackrel{~}{M}_1`$. For this purpose, we note that by a simple scaling of our wave gravity Vlasov equations (59), (60), (64), we can first generalize the initial data $`\stackrel{~}{P}_1(^6)`$ to those $`\stackrel{~}{M}_{1,+}(^6)`$, and by the linearity of (64) in $`f`$ given $`\psi `$ together with the linearity of the map $`f\rho `$, we can even allow the initial data for (64) to be a signed measure $`\sigma _0\stackrel{~}{M}_1(^6)`$; however, we always demand that $`\mathrm{supp}(\sigma _0(\mathrm{d}x\times ^3))B_R`$. As to the initial data for (60), (64), we suitably extend the unique map $`๐ณ_0^{\left(N\right)}\zeta [๐ณ_0^{\left(N\right)}]`$ to $`\sigma _0\zeta [\sigma _0](\dot{H}^1\dot{H}^2)(^3)L^2(^3)`$ (slightly abusing notation), obeying (61), (62), and we write $`๐ท_t(๐ท_0,๐ณ_0^{\left(N\right)})`$ by $`๐ท_t(๐ท_0,\sigma _0)`$. Note that $`๐ท_0`$ itself does *not* depend on $`\sigma _0`$.
The (Gateaux) derivative $`D^1g(\sigma ,.):^6๐
`$ with respect to a finite measure $`\sigma \stackrel{~}{M}_1(^6)`$ of a function $`g(\sigma )๐
`$ (any Banach space) is defined by the identity
$$D^1g(\sigma ,z)\nu (\mathrm{d}z)=\underset{s0}{lim}\frac{g(\sigma +s\nu )g(\sigma )}{s},$$
(206)
valid for all $`\nu \stackrel{~}{M}_1(^6)`$; we here will restrict $`\nu `$ to satisfy $`\mathrm{supp}(\nu _0(\mathrm{d}x\times ^3))B_R`$. Analogously one defines the higher derivatives $`D^jg(\sigma ,.)`$ on $`^{6j}`$. To have a shorthand we write $`๐ณ^{(j)}`$ for generic points in $`^{6j}`$; to achieve a more uniform notation in the presentation we will also write $`๐ณ^{(1)}`$ for $`z`$ in (206) and other first derivatives.
Next we define several auxiliary norms. Below, whenever we take the sup over $`z_0,\zeta _0,\sigma _0,๐ณ^{(j)}`$, it is understood that $`z_0B_R\times ^3`$, that $`\zeta _0(\dot{H}^1\dot{H}^2)(^3)L^2(^3)`$ obeying (61), (62), that $`\sigma _0\stackrel{~}{M}_1`$ has $`\mathrm{supp}(\sigma _0(\mathrm{d}x\times ^3))B_R`$ and $`|\sigma _0|1`$, and that $`๐ณ^{(j)}(B_R\times ^3)^j`$. Thus, we define (noticing that $`D^1\psi _.=D^1\psi _.`$)
$`D^jz_t_{(\mathrm{u})}`$ $`:=`$ $`\underset{z_0,\sigma _0,๐ณ^{(j)}}{sup}|D^jz_t(z_0,\sigma _0,๐ณ^{(j)})|`$ (207)
$`D^j\psi _t_{(L)}`$ $`:=`$ $`\underset{\zeta _0,\sigma _0,๐ณ^{(j)}}{sup}D^j\psi _t(\zeta _0,\sigma _0,๐ณ^{(j)})_{L^2}`$ (208)
$`D^j\varpi _t_{(L)}`$ $`:=`$ $`\underset{\zeta _0,\sigma _0,๐ณ^{(j)}}{sup}D^j\varpi _t(\zeta _0,\sigma _0,๐ณ^{(j)})_{L^2}`$ (209)
$`D^j\psi _t_{(H)}`$ $`:=`$ $`\underset{\zeta _0,\sigma _0,๐ณ^{(j)}}{sup}D^j\psi _t(\zeta _0,\sigma _0,๐ณ^{(j)})_{\dot{H}^1},`$ (210)
and we define $`D^j\zeta _t_{(HL)}`$ by setting
$$D^j\zeta _t_{(HL)}^2:=D^j\psi _t_{(H)}^2+D^j\varpi _t_{(L)}^2$$
(211)
###### Proposition 4.10
Given Cauchy data $`\sigma _0\stackrel{~}{M}_1(^6)`$ with $`\mathrm{supp}(\sigma _0(\mathrm{d}x\times ^3))B_R`$ for (64), and associated with it Cauchy data $`\zeta [\sigma _0]((\dot{H}^1\dot{H}^2)L^2)(^3)`$ for (59), (60), satisfying (61), (62), this Vlasovian Cauchy problem has a unique strong solution $`(\sigma _.,\zeta _.^{(\sigma )})(C_\mathrm{w}^0C^1)(,\stackrel{~}{M}_1((\dot{H}^1\dot{H}^2)L^2)(^3))`$. This solution generates evolutions $`(z_.,\zeta _.)=๐ท_.`$ defined by (202)-(203) which $`k`$ are $`k`$ times continuously differentiable with respect to $`\sigma _0`$. Moreover, for all $`\sigma _0\zeta [\sigma _0]`$ such that there exist functions $`B^j(.)C^0(^+)`$, $`j=1,\mathrm{},k`$, depending on $`\sigma _0`$ only through $`(\sigma _0,\zeta [\sigma _0])`$ and $`|\sigma _0|`$ such that $`D^j\psi ^{(\sigma )}(.,t^{},.)_{(H)}B^j(t)`$, there exist functions $`B_{\mathrm{}}^j(.)C^0(^+)`$, $`j=1,\mathrm{},k`$, and $`\mathrm{}=1,2,3`$, depending on $`\sigma _0`$ only through $`(\sigma _0,\zeta [\sigma _0])`$ and $`|\sigma _0|`$, such that for all $`t`$,
$`D^jz_t_{(\mathrm{u})}`$ $``$ $`B_1^j(t),`$ (212)
$`D^j\zeta _t_{(HL)}`$ $``$ $`B_2^j(t),`$ (213)
$`D^j\psi _t_{(L)}`$ $``$ $`B_3^j(t).`$ (214)
Proof. By inspection one verifies that our existence and uniqueness proof for the Vlasov model carries over to these more general data with only miniscule changes.
As to the existence and uniqueness of $`D^j๐ท_.`$, this follows in the standard way \[Die60\]. The first variation of the evolution equations yields $`D^1๐ท_.`$ thus
$`D^1q_t(z_0,๐ณ^{(1)})`$ $`=`$ $`{\displaystyle _0^t}\left[{\displaystyle \frac{D^1p_t^{}(z_0,๐ณ^{(1)})}{(1+|p_t^{}(z_0)|^2)^{1/2}}}{\displaystyle \frac{p_t^{}(z_0)p_t^{}(z_0)D^1p_t^{}(z_0,๐ณ^{(1)})}{(1+|p_t^{}(z_0)|^2)^{3/2}}}\right]dt^{}`$ (215)
$`D^1p_t(z_0,๐ณ^{(1)})`$ $`=`$ $`{\displaystyle _0^t}{\displaystyle }([\varrho (yq_t^{}(z_0))D^1q_t^{}(z_0,๐ณ^{(1)})]\psi ^{(\sigma )}(y,t^{})`$ (217)
$`\varrho (yq_t^{}(z_0))D^1\psi ^{(\sigma )}(y,t^{},๐ณ^{(1)}))\mathrm{d}y\mathrm{d}t^{}`$
$`D^1\psi _t(\zeta _0,๐ณ^{(1)})(x)`$ $`=`$ $`{\displaystyle _0^t}(tt^{}){\displaystyle _{\mathrm{SS}^2}}[{\displaystyle }[\varrho (x^{}q_t^{}(z_0))D^1q_t^{}(z_0,๐ณ^{(1)})]\sigma _0(\mathrm{d}z_0).`$ (219)
$`..\varrho (x^{}q_t^{}(๐ณ^{(1)}))]\mathrm{d}\mathrm{\Omega }\mathrm{d}t^{}`$
$`D^1\varpi _t(\zeta _0,๐ณ^{(1)})(x)`$ $`=`$ $`{\displaystyle _0^t}{\displaystyle _{\mathrm{SS}^2}}[{\displaystyle }(D^1q_t^{}(z_0,๐ณ^{(1)}))(1+(tt^{})\mathrm{\Omega })\varrho (x^{}q_t^{}(z_0))\sigma _0(\mathrm{d}z_0)`$ (221)
$`.[1+(tt^{})\mathrm{\Omega }]\varrho (x^{}q_t^{}(๐ณ^{(1)}))]\mathrm{d}\mathrm{\Omega }\mathrm{d}t^{}`$
$`D^1\psi _t(\zeta _0,๐ณ^{(1)})(x)`$ $`=`$ $`{\displaystyle _0^t}(tt^{}){\displaystyle _{\mathrm{SS}^2}}[{\displaystyle }(D^1q_t^{}(z_0,๐ณ^{(1)}))\varrho (x^{}q_t^{}(z_0))\sigma _0(\mathrm{d}z_0)`$ (223)
$`.\varrho (x^{}q_t^{}(๐ณ^{(1)}))]\mathrm{d}\mathrm{\Omega }\mathrm{d}t^{},`$
where once again we used the notation $`x^{}=x+(tt^{})\mathrm{\Omega }`$, with $`\mathrm{\Omega }\mathrm{SS}^2`$, and $`_{\mathrm{SS}^2}`$ to denote $`\frac{1}{4\pi }_{\mathrm{SS}^2}`$; moreover, to shorten the presentation, we have not displayed the dependence on $`\sigma _0`$ of $`(z_t,\zeta _t)`$. Note also that in (219), (221), (223) we have transferred the time-dependence from $`\sigma _t^{}`$ to the adjoint time-dependence of $`q_t^{}`$ by change of variables.
Next, as to (219), one easily obtains the bound
$`{\displaystyle _0^t}(tt^{}){\displaystyle _{\mathrm{SS}^2}}\varrho (.^{}q_t^{}(๐ณ^{(1)}))\mathrm{d}\mathrm{\Omega }\mathrm{d}t^{}_{(L)}`$ $``$ $`{\displaystyle _0^t}(tt^{})\underset{\genfrac{}{}{0pt}{}{|\sigma _0|1}{๐ณ^{(1)}B_R\times ^3}}{sup}{\displaystyle _{\mathrm{SS}^2}}\varrho (.^{}q_t^{}(๐ณ^{(1)}))\mathrm{d}\mathrm{\Omega }_{L^2}\mathrm{d}t^{}`$ (224)
$``$ $`{\displaystyle _0^t}(tt^{})\varrho _{L^2}dt^{}=t^2\frac{1}{2}\varrho _{L^2},`$ (225)
where $`.^{}=.+(tt^{})\mathrm{\Omega }`$, and similarly easily, using (207), one obtains
$`{\displaystyle _0^t}(tt^{}){\displaystyle _{\mathrm{SS}^2}}{\displaystyle }\varrho (.^{}q_t^{}(๐ณ^{(1)}))D^1q_t^{}(z_0,๐ณ^{(1)})\sigma _0(\mathrm{d}z_0)\mathrm{d}\mathrm{\Omega }\mathrm{d}t^{}_{(L)}`$ (226)
$``$ $`|\sigma _0|{\displaystyle _0^t}(tt^{})\underset{\genfrac{}{}{0pt}{}{|\sigma _0|1}{๐ณ^{(1)}B_R\times ^3}}{sup}{\displaystyle _{\mathrm{SS}^2}}{\displaystyle }\varrho (.^{}q_t^{}(๐ณ^{(1)}))D^1q_t^{}(z_0,๐ณ^{(1)})\sigma _0(\mathrm{d}z_0)\mathrm{d}\mathrm{\Omega }_{L^2}\mathrm{d}t^{}`$ (227)
$``$ $`|\sigma _0|{\displaystyle _0^t}(tt^{})\varrho _{L^2}D^1q_t^{}_{(\mathrm{u})}dt^{}.`$ (228)
These bounds and similar ones for (221) and (223), and some obvious estimates for (215) and (217), yield
$`D^1q_t_{(\mathrm{u})}`$ $``$ $`{\displaystyle _0^t}D^1p_t^{}_{(\mathrm{u})}dt^{}`$ (229)
$`D^1p_t_{(\mathrm{u})}`$ $``$ $`{\displaystyle _0^t}[\varrho _{L^2}\psi ^{(\sigma )}(.,t^{})_{\dot{H}^1}D^1q_t^{}_{(\mathrm{u})}`$ (231)
$`+\varrho _{L^2}D^1\psi ^{(\sigma )}(.,t^{},.)_{(H)}]\mathrm{d}t^{}`$
$`D^1\psi _t_{(H)}`$ $``$ $`t^2\frac{1}{2}\varrho _{L^2}+^2\varrho _{L^2}|\sigma _0|{\displaystyle _0^t}(tt^{})D^1q_t^{}_{(\mathrm{u})}dt^{}`$ (232)
$`D^1\varpi _t_{(L)}`$ $``$ $`t\varrho _{L^2}+t^2\frac{1}{2}\varrho _{L^2}`$ (234)
$`+|\sigma _0|{\displaystyle _0^t}\left[\varrho _{L^2}+(tt^{})^2\varrho _{L^2}\right]D^1q_t^{}_{(\mathrm{u})}dt^{}`$
$`D^1\psi _t_{(L)}`$ $``$ $`t^2\frac{1}{2}\varrho _{L^2}+\varrho _{L^2}|\sigma _0|{\displaystyle _0^t}(tt^{})D^1q_t^{}_{(\mathrm{u})}dt^{},`$ (235)
where $`^2\varrho _{L^2}^2=|^2\varrho |^2(x)dx`$, with $`|M|`$ the familiar Euclidean norm of a real symmetric $`3\times 3`$ matrix as an element of $`^9`$.
We now recall that $`\psi ^{(\sigma )}(.,t^{})_{\dot{H}^1}`$ is bounded uniformly for $`t`$ by a constant depending on $`\sigma _0`$ only through $`(\sigma _0,\zeta [\sigma _0])`$; cf. Remark 3.10.
Next, as a special case, assume first that $`\zeta [\sigma _0]=\zeta _0`$ independent of $`\sigma _0`$. In this case $`\psi ^{(\sigma )}(.,t)=\psi _t`$, and the bounds (212), (213) for $`j=1`$ follow from (215)โ(221) by variants of the Gronwall lemma. The bound (214) for $`j=1`$ then follows immediately. The bounds (212), (213) for general values of $`j`$ follow by applying $`D^{j1}`$ to (215), (217), (219), (221), (223). In particular, the bound $`D^j\psi ^{(\sigma )}(.,t^{},.)_{(H)}B^j(t)`$, $`j=1,\mathrm{},k`$, holds, with $`B^j(.)C^0(^+)`$ depending on $`\sigma _0`$ only through $`(\sigma _0,\zeta [\sigma _0])`$ and $`|\sigma _0|`$, for $`\psi ^{(\sigma )}(.,t)=\psi _t`$ in this case.
Finally, the bounds clearly generalize to field data $`\sigma _0\zeta [\sigma _0]`$ for which there exist functions $`B^j(.)C^0(^+)`$, $`j=1,\mathrm{},k`$, depending on $`\sigma _0`$ only through $`(\sigma _0,\zeta [\sigma _0])`$ and $`|\sigma _0|`$ such that $`D^j\psi ^{(\sigma )}(.,t^{},.)_{(H)}B^j(t)`$. That this hypothesis is legitimate we just showed, for its supposition is valid in particular when $`\zeta [\sigma _0]=\zeta _0`$. Q.E.D.
###### Remark 4.11
Unfortunately we do not yet know how general the class of initial data $`\sigma _0\zeta [\sigma _0]`$ is which validate our hypothesis that $`D^j\psi ^{(\sigma )}(.,t^{},.)_{(H)}B^j(t)`$, $`j=1,\mathrm{},k`$, with $`B^j(.)C^0(^+)`$ depending on $`\sigma _0`$ only through $`(\sigma _0,\zeta [\sigma _0])`$ and $`|\sigma _0|`$.
Because of Remark 4.11, in the following we will restrict the initial conditions for the fields to the special case $`\zeta [๐ณ_0^{\left(N\right)}]=\zeta _0`$ independent of $`๐ณ_0^{\left(N\right)}`$; the case $`N=\mathrm{}`$ is included. In this vein, we may now analyze the limit $`N\mathrm{}`$ of the finite-$`N`$ fluctuations
$`\mathrm{\Delta }_z(t,z_0;๐ณ_0^{\left(N\right)},\mu _0)`$ $`:=`$ $`\sqrt{N}[z_t(z_0;๐ณ_0^{\left(N\right)})z_t(z_0,\mu _0)],`$ (236)
$`\mathrm{\Delta }_\zeta (t,\zeta _0;๐ณ_0^{\left(N\right)},\mu _0)`$ $`:=`$ $`\sqrt{N}[\zeta _t(\zeta _0;๐ณ_0^{\left(N\right)})\zeta _t(\zeta _0,\mu _0)].`$ (237)
We write $`\mathrm{\Delta }_๐ท=(\mathrm{\Delta }_z,\mathrm{\Delta }_\zeta )`$. Recall that by Theorem 4.2, if $`\epsilon [๐ณ_0^{\left(N\right)}]\mu _0`$ as $`N\mathrm{}`$ and $`\zeta [๐ณ_0^{\left(N\right)}]=\zeta _0^{(\mathrm{})}\zeta _0`$ for all $`N`$, then
$`\epsilon [๐ณ_t^{\left(N\right)}]`$ $``$ $`\mu _t`$ (238)
$`\psi _t(\zeta _0;๐ณ_0^{\left(N\right)})`$ $`\stackrel{\dot{H}^1}{}`$ $`\psi _t(\zeta _0;\mu _0)`$ (239)
$`\varpi _t(\zeta _0;๐ณ_0^{\left(N\right)})`$ $`\stackrel{L^2}{}`$ $`\varpi _t(\zeta _0;\mu _0)`$ (240)
and, due to Corollary 3.7,
$$|z_t(z_0,๐ณ_0^{\left(N\right)})z_t(z_0,\mu _0)|0.$$
(241)
We now are ready to state the central limit theorem.
###### Theorem 4.12
For $`\mu _0P_1(^6)`$, define the sequence of particle product measures $`\mu _0^{\times N}(\mathrm{d}๐ณ_0^{\left(N\right)})=_{i=1}^N\mu _0(\mathrm{d}z_i^{(N)}(0))`$, and consider field initial data $`\zeta _0(\dot{H}^1\dot{H}^2)(^3)L^2(^3)`$. For any $`K1`$ and $`1kK`$, let $`Z_k\mathrm{\Gamma }^{(1)}`$, $`z_k^6`$, $`t_k`$ and define $`\stackrel{}{Z}=(Z_1,\mathrm{},Z_K)(\mathrm{\Gamma }^{(1)})^K`$, $`\stackrel{}{z}=(z_1,\mathrm{},z_k)^{6K}`$, and $`\stackrel{}{t}=(t_1,\mathrm{},t_K)^K`$. Moreover, let $`\stackrel{}{\mathrm{\Delta }_๐ท}(\stackrel{}{t},\stackrel{}{z},\zeta _0;๐ณ_0^{\left(N\right)})(\mathrm{\Gamma }^{(1)})^K`$ have $`k`$-th component $`\mathrm{\Delta }_๐ท(t_k,z_k,\zeta _0;๐ณ_0^{\left(N\right)})`$. Then
$$\underset{N\mathrm{}}{lim}e^{\mathrm{i}\stackrel{}{Z}|\stackrel{}{\mathrm{\Delta }_๐ท}(\stackrel{}{t},\stackrel{}{z},\zeta _0;๐ณ_0^{(N)})}\mu _0^{\times N}(\mathrm{d}๐ณ_0^{\left(N\right)})=e^{\frac{1}{2}\stackrel{}{Z}|\stackrel{}{Q}|\stackrel{}{Z}}$$
(242)
where $`.|.`$ is the scalar product in the Hilbert space $`(\mathrm{\Gamma }^{(1)})^K`$ (i.e. the sum over $`K`$ scalar products $`.,.`$ in $`\mathrm{\Gamma }^{(1)}`$ indexed by <sub>k</sub>), and the operator $`\stackrel{}{Q}`$ has $`k,k^{}`$ component
$`Q_{k,k^{^{}}}`$ $`=`$ $`{\displaystyle \underset{\kappa \{k,k^{}\}}{}D^1๐ท_{t_\kappa }(z_\kappa ,\zeta _0,\mu _0,๐ณ^{(1)})\mu _0(\mathrm{d}๐ณ^{(1)})}`$ (244)
$`{\displaystyle \underset{\kappa \{k,k^{}\}}{}}{\displaystyle D^1๐ท_{t_\kappa }(z_\kappa ,\zeta _0,\mu _0,๐ณ_\kappa ^{(1)})\mu _0(\mathrm{d}๐ณ_\kappa ^{(1)})}.`$
For $`\stackrel{}{t}`$ and $`\stackrel{}{Z}`$ in bounded sets, the convergence is uniform in $`\stackrel{}{z}`$.
The stochastic process $`\eta =lim_N\mathrm{}\mathrm{\Delta }_๐ท`$ on $`\mathrm{\Gamma }^{(1)}`$, with vanishing expectation and covariance (244), can be represented as
$$\eta _t(z_0,\mu _0)=D^1๐ท_t(z_0,\mu _0,๐ณ^{(1)})\phi (\mu _0,\mathrm{d}๐ณ^{(1)})$$
(245)
and satisfies the equations obtained integrating (215), (217), (219), (221) with respect to $`\phi (\mu _0,\mathrm{d}๐ณ^{(1)})`$ with initial conditions $`\eta _0(z_0,\mu _0)=0`$. Here, the random measure $`\phi (\mu ,\mathrm{d}๐ณ^{(1)})M(^6)`$ with Gaussian law is defined by
$`๐ผ_{\mu _0^{\times N}}\phi (\mu _0,\mathrm{}_1)`$ $`=`$ $`0,`$ (246)
$`๐ผ_{\mu _0^{\times N}}[\phi (\mu _0,\mathrm{}_1)\phi (\mu _0,\mathrm{}_2)]`$ $`=`$ $`\mu _0(\mathrm{}_1\mathrm{}_2)\mu _0(\mathrm{}_1)\mu _0(\mathrm{}_2),`$ (247)
for measurable $`\mathrm{}_1,\mathrm{}_2^6`$.
Proof. Apart from the different kind of convergence for the fluctuations of the potential, the proof is analogous to the one in \[BrHe77\]. We carry out some calculations to clarify the procedure. Let $`\nu [๐ณ_0^{\left(N\right)}]:=\epsilon [๐ณ_0^{\left(N\right)}]\mu _0`$. The following calculations are valid whenever $`D^1๐ท_.`$ exists.
Obviously,
$$๐ท_t(z_0,\epsilon _{๐ณ_0^{(N)}})๐ท_t(z_0,\mu _0)=_0^1D^1๐ท_t(z_0,\zeta _0,\mu _0+r\nu [๐ณ_0^{\left(N\right)}],๐ณ^{(1)})\nu [๐ณ_0^{\left(N\right)}](\mathrm{d}๐ณ^{(1)})dr.$$
(248)
We now define
$`\mathrm{\Xi }_k^{(1)}(\sigma ,๐ณ^{(1)})`$ $`=`$ $`\sqrt{N}D^1๐ท_{t_k}(z_k,\zeta _0,\sigma ,๐ณ^{(1)}),`$ (249)
$`\mathrm{\Xi }_k^{(2)}(\sigma ,๐ณ^{(2)})`$ $`=`$ $`\sqrt{N}D^2๐ท_{t_k}(z_k,\zeta _0,\sigma ,๐ณ^{(2)}),`$ (250)
$`\xi _k^{(1)}(\sigma )`$ $`=`$ $`{\displaystyle \mathrm{\Xi }_k^{(1)}(\sigma ,๐ณ^{(1)})\nu [๐ณ_0^{\left(N\right)}](\mathrm{d}๐ณ^{(1)})},`$ (251)
$`\xi _k^{(2)}(\sigma )`$ $`=`$ $`{\displaystyle \mathrm{\Xi }_k^{(2)}(\sigma ,๐ณ^{(2)})\nu [๐ณ_0^{\left(N\right)}]^{\times 2}(\mathrm{d}๐ณ^{(2)})},`$ (252)
and group the corresponding $`K`$ components into the vectors $`\stackrel{}{\mathrm{\Xi }}^{(1)}`$, $`\stackrel{}{\mathrm{\Xi }}^{(2)}`$, $`\stackrel{}{\xi }^{(1)}`$ and $`\stackrel{}{\xi }^{(2)}`$. We also define
$$\stackrel{~}{\xi }_{kj}(\sigma )=\frac{1}{\sqrt{N}}\left(\mathrm{\Xi }_k^{(1)}(\sigma ,z_j^{\left(N\right)}(0))\mathrm{\Xi }_k^{(1)}(\sigma ,๐ณ^{(1)})\sigma (\mathrm{d}๐ณ^{(1)})\right)$$
(253)
where $`z_j^{\left(N\right)}(0)`$ is the $`j`$-th component of $`๐ณ_0^{\left(N\right)}`$, and note that
$$\xi _k^{(1)}=\frac{1}{\sqrt{N}}\underset{j=1}{\overset{N}{}}\stackrel{~}{\xi }_{kj}$$
(254)
and that
$`๐ผ_{\mu _0^{\times N}}(\stackrel{~}{\xi }_{kj})`$ $`=`$ $`0,`$ (255)
$`๐ผ_{\mu _0^{\times N}}(\stackrel{~}{\xi }_{kj}\stackrel{~}{\xi }_{k^{}j^{}})`$ $`=`$ $`Q_{k,k^{}}\delta _{j,j^{}},`$ (256)
where $`\delta _{j,j^{}}`$ is the Kronecker symbol. Hence the central limit theorem applies to $`\xi _k^{(1)}`$.
We may now write
$$e^{\mathrm{i}\stackrel{}{Z}|\stackrel{}{\mathrm{\Delta }_๐ท}}=e^{\mathrm{i}\stackrel{}{Z}|\stackrel{}{\xi }^{(1)}}+_0^1\frac{\mathrm{d}}{\mathrm{d}s}\mathrm{exp}\left(\mathrm{i}_0^1\stackrel{}{Z}|\stackrel{}{\xi }^{(1)}(\mu _0+rs\nu [๐ณ_0^{\left(N\right)}])dr\right)ds.$$
(257)
When $`N\mathrm{}`$, the expectation of the first term in the right-hand side converges to the right-hand side of (242).
The integrand in the second term becomes, after differentiation w.r.t. $`s`$,
$`\mathrm{i}{\displaystyle _0^1}\stackrel{}{Z}|\stackrel{}{\xi }^{(2)}(\mu _0+rs\nu [๐ณ_0^{\left(N\right)}])rdr\mathrm{exp}\left(\mathrm{i}{\displaystyle _0^1}\stackrel{}{Z}|\stackrel{}{\xi }^{(1)}(\mu _0+r^{}s\nu [๐ณ_0^{\left(N\right)}])dr^{}\right)`$ (258)
$`=`$ $`\frac{\mathrm{i}}{N\sqrt{N}}{\displaystyle \underset{i,j=1}{\overset{N}{}}}[๐ข(z_i^{(N)}(0),z_j^{(N)}(0),\epsilon [๐ณ_0^{\left(N\right)}]){\displaystyle }๐ข(z_i^{(N)}(0),๐ณ^{(1)},\epsilon [๐ณ_0^{\left(N\right)}])\mu _0(\mathrm{d}๐ณ^{(1)})`$ (261)
$`{\displaystyle }๐ข(๐ณ^{(1)},z_j^{(N)}(0),\epsilon [๐ณ_0^{\left(N\right)}])\mu _0(\mathrm{d}๐ณ^{(1)})+{\displaystyle }{\displaystyle }๐ข(๐ณ^{(2)},\epsilon [๐ณ_0^{\left(N\right)}])\mu _0^{\times 2}(\mathrm{d}๐ณ^{(2)})]`$
$`\mathrm{exp}\left({\displaystyle \frac{\mathrm{i}}{\sqrt{N}}}{\displaystyle \underset{i^{}=1}{\overset{N}{}}}[(z_i^{}^{\left(N\right)}(0),\epsilon [๐ณ_0^{\left(N\right)}]){\displaystyle }(\stackrel{~}{๐ณ}{}_{}{}^{(1)},\epsilon [๐ณ_0^{\left(N\right)}])\mu _0(\mathrm{d}\stackrel{~}{๐ณ}{}_{}{}^{(1)})]\right)`$
where (for given $`\stackrel{}{Z}`$, $`\stackrel{}{z}`$, $`\stackrel{}{t}`$, $`\zeta _0,\mu _0,s`$)
$`\sqrt{N}(๐ณ^{(1)},\sigma _0)`$ $`=`$ $`{\displaystyle _0^1}\stackrel{}{Z}|\stackrel{}{\mathrm{\Xi }}^{(1)}(\mu _0+rs(\sigma _0\mu _0),๐ณ^{(1)})dr,`$ (262)
$`\sqrt{N}๐ข(๐ณ^{(2)},\sigma _0)`$ $`=`$ $`{\displaystyle _0^1}\stackrel{}{Z}\stackrel{}{\mathrm{\Xi }}^{(2)}(\mu _0+rs(\sigma _0\mu _0),๐ณ^{(2)})rdr.`$ (263)
The regularization ensures that $`๐ข`$ and $``$ are differentiable to any order with respect to $`\sigma _0`$, with bounded derivatives in the sense of Proposition 4.10. We may then evaluate the size of the expectation of (261) in the following way.
Expression (261) can be split in two components : the โdiagonalโ part, which is obtained from the terms in the sum such that $`i=j`$, gives with trivial estimates a contribution of order $`N^{1/2}`$ to the expectation of (261) ; an estimate of the size of the โnon-diagonalโ component needs some more manipulation.
Consider the measure (positive with mass $`12/N`$)
$$\mu ^{ij}[๐ณ_0^{\left(N\right)}]=\frac{1}{N}\underset{\stackrel{k=1}{ki,j}}{\overset{N}{}}\delta _{z_k^{(N)}(0)}=\epsilon [๐ณ_0^{\left(N\right)}]\frac{1}{N}(\delta _{z_i^{(N)}(0)}+\delta _{z_j^{(N)}(0)}).$$
(264)
Given the result in Proposition 4.10, we may write
$`(๐ณ^{(1)},\epsilon [๐ณ_0^{\left(N\right)}])`$ $`=`$ $`(๐ณ^{(1)},\mu ^{ij}[๐ณ_0^{\left(N\right)}])+\frac{1}{N}\underset{k\{i,j\}}{}D^1(๐ณ^{(1)},\mu ^{ij}[๐ณ_0^{\left(N\right)}],z_k^{(N)}(0))+\mathrm{O}(\frac{1}{N^2}),`$
$`๐ข(๐ณ^{(2)},\epsilon [๐ณ_0^{\left(N\right)}])`$ $`=`$ $`๐ข(๐ณ^{(2)},\mu ^{ij}[๐ณ_0^{\left(N\right)}])+\frac{1}{N}\underset{k\{i,j\}}{}D^1๐ข(\stackrel{~}{\stackrel{~}{๐ณ}}{}_{}{}^{(1)},\mu ^{ij}[๐ณ_0^{\left(N\right)}],z_k^{(N)}(0))+\mathrm{O}(\frac{1}{N^2}).`$
For a given pair $`\{i,j\}`$, we then obtain
$`\mathrm{exp}\left({\displaystyle \frac{\mathrm{i}}{\sqrt{N}}}{\displaystyle \underset{n=1}{\overset{N}{}}}\left[(z_n^{(N)}(0),\epsilon [๐ณ_0^{\left(N\right)}]){\displaystyle (๐ณ^{(1)},\epsilon [๐ณ_0^{\left(N\right)}])\mu _0(\mathrm{d}๐ณ^{(1)})}\right]\right)`$ (265)
$`=\left[1+\frac{\mathrm{i}}{\sqrt{N}}\left[_1(\widehat{z}_i^{(N)}(0))+_1(\widehat{z}_j^{(N)}(0))\right]+\mathrm{O}\left(\frac{1}{N}\right)\right]_2(\widehat{z}_i^{(N)}(0),\widehat{z}_j^{(N)}(0))`$ (266)
where we denote by $`g(\widehat{x})`$ a function such that $`_xg=0`$, which is bounded according to Proposition 4.10 and need not be further specified. Moreover,
$`๐ข(z_i^{(N)}(0),z_j^{(N)}(0),\epsilon [๐ณ_0^{\left(N\right)}])`$ $``$ $`{\displaystyle }๐ข(z_i^{(N)}(0),\stackrel{~}{๐ณ}{}_{}{}^{(1)},\epsilon [๐ณ_0^{\left(N\right)}])\mu _0(\mathrm{d}\stackrel{~}{๐ณ}{}_{}{}^{(1)})`$ (267)
$``$ $`{\displaystyle ๐ข(๐ณ^{(1)},z_j^{(N)}(0),\epsilon [๐ณ_0^{\left(N\right)}])\mu _0(\mathrm{d}๐ณ^{(1)})}`$ (268)
$`+`$ $`{\displaystyle ๐ข(๐ณ^{(2)},\epsilon [๐ณ_0^{\left(N\right)}])\mu _0^{\times 2}(\mathrm{d}๐ณ^{(2)})}=๐_{ij}+\mathrm{O}\left(\frac{1}{N}\right)`$ (269)
where
$`๐_{ij}=๐ข(z_i^{(N)}(0),z_j^{(N)}(0),\mu ^{ij}[๐ณ_0^{\left(N\right)}])`$ $``$ $`{\displaystyle }๐ข(z_i^{(N)}(0),\stackrel{~}{๐ณ}{}_{}{}^{(1)},\mu ^{ij}[๐ณ_0^{\left(N\right)}])\mu _0(\mathrm{d}\stackrel{~}{๐ณ}{}_{}{}^{(1)})`$
$``$ $`{\displaystyle ๐ข(๐ณ^{(1)},z_j^{(N)}(0),\mu ^{ij}[๐ณ_0^{\left(N\right)}])\mu _0(\mathrm{d}๐ณ^{(1)})}`$
$`+`$ $`{\displaystyle }{\displaystyle }๐ข(๐ณ^{(2)},\mu ^{ij}[๐ณ_0^{\left(N\right)}])]\mu _0^{\times 2}(\mathrm{d}๐ณ^{(2)})`$
so that
$`๐ผ_{\mu _0^{\times N}}[๐_{ij}]`$ $`=`$ $`0,`$ (270)
$`๐ผ_{\mu _0^{\times N}}[๐_{ij}g(\widehat{z}_k^{(N)}(0))]`$ $`=`$ $`0,`$ (271)
for any bounded function $`g(\widehat{z}_k^{(N)}(0))`$ with $`k\{i,j\}`$. The expectation of the non-diagonal component of (261) is then given by
$`\frac{\mathrm{i}}{N\sqrt{N}}{\displaystyle \underset{\stackrel{i,j=1}{ij}}{\overset{N}{}}}{\displaystyle }[๐_{ij}+\mathrm{O}\left(\frac{1}{N}\right)][1+\frac{\mathrm{i}}{\sqrt{N}}(_1(\widehat{z}_i^{(N)}(0))+_1(\widehat{z}_j^{(N)}(0)))+\mathrm{O}\left(\frac{1}{N}\right)]\times `$ (272)
$`\times _2(\widehat{z}_i^{(N)}(0),\widehat{z}_j^{(N)}(0))\mu _0^{\times N}(\mathrm{d}๐ณ_0^{\left(N\right)})=\frac{N(N1)}{N\sqrt{N}}\mathrm{O}\left(\frac{1}{N}\right)`$ (273)
so that the second term in (257) is $`\mathrm{O}(N^{1/2})`$.
The identification of the limit stochastic process is obtained from (215)-(217)-(219)-(221). Q.E.D.
Acknowledgement The work of M.K. was supported in parts by the National Science Foundation under Grant No. DMS-0103808, and in parts by CNRS through a poste rouge to M.K. while visiting CNRS-Universitรฉ de Provence; the work of Y.E. by Universitรฉ de Provence through a congรฉ pour recherche, the work of V.R., when the collaboration started, by the Foundation BLANCEFLOR Boncompagni Ludovisi nรฉe Bildt. The participation by A. Nouri in the early stages of this work is gratefully acknowledged. M.K. thanks C. Lancellotti, M. Kunze, and H. Spohn for valuable discussions.
Appendix
### A.1 Nested modes of convergence of probability measures
A certain frustration about the absence of an authoritative survey of the relationships of various important notions of convergence that are used in the probability literature has already been expressed \[GiSu02\], where that gap has been filled to some extent. Unfortunately, \[GiSu02\] does not cover all our needs. Furthermore, when addressing a mixed readership of mathematical physicists, analysts and probabilists, the frustration can get compounded by the various โcompetingโ terminologies and notations that are in use in these areas of activity. In view of this, it seems advisable to be more explicit about how the notions of convergence that we use. The following general notions hold (and are formulated) for any dimension $`d1`$.
We recall that, if $`\{\mu _n\}_n`$ is a sequence of Borel probability measures on $`^d`$ and $`\mu P(^d)`$, too, and if $`fd\mu _nfd\mu `$ for every bounded continuous function $`fC_b^0(^d)`$, then one says that $`\mu _n`$ converges to $`\mu `$ in law,<sup>12</sup><sup>12</sup>12In the probability literature, convergence in law is usually called โweak convergenceโ of probability measures; however, this notion generally differs from the analystsโ notion of weak convergence on $`M`$. written $`\mu _n\stackrel{}{}\mu `$; see p. 292 of \[Dud02\]. Clearly, since $`C_0^0(^d)C_b^0(^d)`$, convergence in law $`\mu _n\stackrel{}{}\mu `$ implies vague convergence $`\mu _n\mu `$. Moreover, convergence in law $`\mu _n\stackrel{}{}\mu `$ is equivalent to convergence in probability of the underlying family of random variables having laws $`\mu _n`$ to a random variable with law $`\mu `$, a notion we need for our law of large numbers.
Convergence in law can be metrized as follows. Let $`C_b^{0,\alpha }(^d)`$ denote the subset of the bounded continuous functions on $`^d`$ which are also Hรถlder continuous with exponent $`\alpha (0,1]`$. Now $`C_b^{0,\alpha }(^d)`$ is not a closed subspace of $`C_b^0(^d)`$ w.r.t. $`._\mathrm{u}`$, but
$$g_{\mathrm{u},\alpha }\mathrm{max}\{g_\mathrm{u},\mathrm{H}\ddot{\mathrm{o}}\mathrm{l}_\alpha \left(g\right)\},$$
(274)
where
$$\mathrm{H}\ddot{\mathrm{o}}\mathrm{l}_\alpha \left(g\right)\underset{\xi \xi ^{}^d}{sup}\frac{|g(\xi )g(\xi ^{})|}{|\xi \xi ^{}|^\alpha }$$
(275)
is the $`\alpha `$-Hรถlder seminorm of $`g`$, turns $`C_b^{0,\alpha }(^d)`$ into a (non-separable) Banach space. The positive cone in $`C_b^{0,\alpha }(^d)`$ is denoted by $`C_{b,+}^{0,\alpha }(^d)`$. If the suffix <sub>b</sub> is replaced by the suffix <sub>0</sub>, we mean the corresponding subsets of these functions that vanish at infinity. In much of what follows, we will need $`C_b^{0,1}(^d)`$, the space of bounded Lipschitz functions on $`^d`$, and we write<sup>13</sup><sup>13</sup>13Since $`\mathrm{Lip}(.)`$ is a seminorm, we prefer this notation over $`._\mathrm{L}`$, which is also in use in the literature. $`\mathrm{Lip}\left(g\right)`$ for $`\mathrm{H}\ddot{\mathrm{o}}\mathrm{l}_\alpha \left(g\right)`$ when $`\alpha =1`$.<sup>14</sup><sup>14</sup>14We recall that if $`gC^1(^d)`$, then $`\mathrm{Lip}\left(g\right)=sup_{x^d}|g(x)|`$.
Now let $`\mu _1P(^d)`$ and $`\mu _2P(^d)`$ be two Borel probability measures on $`^d`$. We define the dual bounded-Lipschitz distance between $`\mu _1`$ and $`\mu _2`$ as<sup>15</sup><sup>15</sup>15The * at $`\mathrm{dist}_{\mathrm{bL}^{}}(,)`$ refers to the KantorovichโRubinstein duality theorems; see below.
$$\mathrm{dist}_{\mathrm{bL}^{}}(\mu _1,\mu _2):=\underset{gC_b^{0,1}(^d)}{sup}\left\{\right|g\mathrm{d}(\mu _1\mu _2)|:g_{\mathrm{u},1}1\}.$$
(276)
Our dual bounded-Lipschitz distance, though not identical to, is equivalent to the FortetโMourier $`\beta `$-distance (p.395 of \[Dud02\]), which instead of $`g_{\mathrm{u},1}1`$ works with the equivalent condition $`g_\mathrm{u}+\mathrm{Lip}\left(g\right)1`$. Therefore, by Proposition 11.3.2 of \[Dud02\], $`\mathrm{dist}_{\mathrm{bL}^{}}(,)`$ is a metric on the convex set $`P(^d)`$, and by Corollary 11.5.5 of \[Dud02\], $`P(^d)`$ is complete for $`\mathrm{dist}_{\mathrm{bL}^{}}(,)`$. Furthermore, by Theorem 11.3.3 of \[Dud02\], if $`\{\mu _n\}_n`$ is a sequence of Borel probability measures on $`^d`$, and $`\mu P(^d)`$, too, then $`\mathrm{dist}_{\mathrm{bL}^{}}(\mu _n,\mu )0`$ as $`n\mathrm{}`$ is equivalent to $`\mu _n\stackrel{}{}\mu `$ as $`n\mathrm{}`$. Hence, $`\mathrm{dist}_{\mathrm{bL}^{}}(,)`$ metrizes convergence in law of the Borel probability measure on $`^d`$.
Our dual bounded-Lipschitz distance $`\mathrm{dist}_{\mathrm{bL}^{}}(,)`$ is equivalent, but not identical, to the distance obtained by restricting $`g`$ to $`C_{b,+}^{0,\alpha }(^d)`$, here denoted $`\mathrm{d}_{\mathrm{bL}^{}}(,)`$ (following \[Spo91\], Def. 2.2; actually, Spohn writes $`\mathrm{d}_{\mathrm{bL}}(,)`$, but we here better keep the \*). Clearly, $`\mathrm{dist}_{\mathrm{bL}^{}}(\mu _n,\mu )0`$ implies $`\mathrm{d}_{\mathrm{bL}^{}}(\mu _n,\mu )0`$. The converse of this follows from three simple observations: first, the integral on the r.h.s. of (276) is invariant under $`gg+g_\mathrm{u}`$, so that in our definition of $`\mathrm{dist}_{\mathrm{bL}^{}}(,)`$ we can replace $`C_b^{0,1}(^d)`$ by $`C_{b,+}^{0,1}(^d)`$ and simultaneously replace the condition $`g_{\mathrm{u},1}1`$ with the condition $`\mathrm{max}\{\frac{1}{2}g_\mathrm{u},\mathrm{Lip}\left(g\right)\}1`$; second, $`\{gC_{b,+}^{0,1}(^d):g_\mathrm{u}2,\mathrm{Lip}\left(g\right)1\}`$ is a strict subset of $`\{gC_{b,+}^{0,1}(^d):g_\mathrm{u}2,\mathrm{Lip}\left(g\right)2\}`$; third, the simple scaling $`g2g`$ reveals that the sup of $`\left|g\mathrm{d}(\mu _1\mu _2)\right|`$ over $`\{gC_{b,+}^{0,1}(^d):g_\mathrm{u}2,\mathrm{Lip}\left(g\right)2\}`$ is twice the sup of $`\left|g\mathrm{d}(\mu _1\mu _2)\right|`$ over $`\{gC_{b,+}^{0,1}(^d):g_\mathrm{u}1,\mathrm{Lip}\left(g\right)1\}`$. These three facts together imply that $`\mathrm{dist}_{\mathrm{bL}^{}}(\mu _1,\mu _2)2\mathrm{d}_{\mathrm{bL}^{}}(\mu _1,\mu _2)`$, and this means that $`\mathrm{dist}_{\mathrm{bL}^{}}(\mu _n,\mu )0`$ whenever $`\mathrm{d}_{\mathrm{bL}^{}}(\mu _n,\mu )0`$.
Recall that the general KantorovichโRubinstein distance<sup>16</sup><sup>16</sup>16Also associated with the names of Monge and Wasserstein. is defined as
$$\mathrm{dist}_{\mathrm{KRc}}(\mu _1,\mu _2):=\underset{\mu P_c(^{2d}|\mu _1,\mu _2)}{inf}\left\{\mathrm{cost}(\xi _1,\xi _2)\mu (\mathrm{d}\xi _1\mathrm{d}\xi _2)\right\},$$
(277)
where $`\mathrm{cost}(\xi ,\xi ^{})=\mathrm{dist}_{\mathrm{KRc}}(\delta _\xi ,\delta _\xi ^{})`$ for $`\xi ,\xi ^{}^d`$ is the โcost (per transport unit) function,โ and where $`P_c(^{2d}|\mu _1,\mu _2)`$ is the set of Borel probability measures $`\mu `$ on $`^d\times ^d`$ satisfying $`\mu (\mathrm{d}\xi _1\times ^d)=\mu _1(\mathrm{d}\xi _1)`$ and $`\mu (^d\times \mathrm{d}\xi _2)=\mu _2(\mathrm{d}\xi _2)`$, with $`\mu _1`$ and $`\mu _2`$ satisfying $`\mathrm{cost}(\xi _1,\xi )\mu _1(\mathrm{d}\xi _1)<\mathrm{}`$ and $`\mathrm{cost}(\xi ,\xi _2)\mu _2(\mathrm{d}\xi _2)<\mathrm{}`$ for some $`\xi ^d`$.
By the KantorovichโRubinstein theorem (\[Dud02\], Theorem 11.8.2), $`\mathrm{dist}_{\mathrm{bL}^{}}(\mu _1,\mu _2)`$ is identical to the KantorovichโRubinstein distance for $`\mathrm{cost}(\xi _1,\xi _2)=\mathrm{min}\{2,|\xi _1\xi _2|\}`$. Incidentally, $`\mathrm{cost}(\xi _1,\xi _2)=\mathrm{min}\{1,|\xi _1\xi _2|\}`$ is the cost function for the particular KantorovichโRubinstein distance identical to $`\mathrm{d}_{\mathrm{bL}^{}}(,)`$. The dual bounded-Lipschitz distance ($`\mathrm{d}_{\mathrm{bL}^{}}`$) is used in \[NeWi74, BrHe77, Neu85, Spo91\] and \[FiEl98\].
However, if one is only interested, as we are, in the subset $`P_1(^d)P(^d)`$, it is rather prudent to work with the dual Lipschitz distance in $`P_1(^d)`$, given by
$$\mathrm{dist}_\mathrm{L}^{}(\mu _1,\mu _2):=\underset{gC^{0,1}(^d)}{sup}\left\{\right|g\mathrm{d}(\mu _1\mu _2)|:\mathrm{Lip}\left(g\right)1\},$$
(278)
which is identical with the standard<sup>17</sup><sup>17</sup>17The word โstandardโ refers to the custom in the probability community that, by default, the cost function is identified with the metric of the underlying complete metric space on which the Borel probability measures are defined; in standard Euclidean $`^d`$ this gives $`\mathrm{cost}(\xi _1,\xi _2)=|\xi _1\xi _2|`$. KantorovichโRubinstein distance, given by
$$\mathrm{dist}_{\mathrm{KR}}(\mu _1,\mu _2):=\underset{\mu P_1(^{2d}|\mu _1,\mu _2)}{inf}\{|\xi _1\xi _2|\mu (\mathrm{d}\xi _1\mathrm{d}\xi _2)\},$$
(279)
where $`P_1(^{2d}|\mu _1,\mu _2)`$ is the set of Borel probability measures $`\mu `$ on $`^d\times ^d`$ satisfying $`\mu (\mathrm{d}\xi _1\times ^d)=\mu _1(\mathrm{d}\xi _1)P_1(^d)`$ and $`\mu (^d\times \mathrm{d}\xi _2)=\mu _2(\mathrm{d}\xi _2)P_1(^d)`$. We write $`\mu _n\mu `$ if $`\mathrm{dist}_\mathrm{L}^{}(\mu _n,\mu )0`$. Clearly, $`\mathrm{dist}_\mathrm{L}^{}(\mu _n,\mu )0`$ implies<sup>18</sup><sup>18</sup>18The converse is not true. In particular, Dudley gives the following counterexample for $`d=1`$: $`\mu _n=(1n^1)\delta _0+n^1\delta _n`$ and $`\mu =\delta _0`$, for which $`\mathrm{dist}_\mathrm{L}^{}(\mu _n,\mu )=1`$ while $`\mathrm{dist}_{\mathrm{bL}^{}}(\mu _n,\mu )2n^10`$. $`\mathrm{dist}_{\mathrm{bL}^{}}(\mu _n,\mu )0`$.
We note that the metric $`\mathrm{dist}_\mathrm{L}^{}(.,.)`$ defines a norm $`._\mathrm{L}^{}`$ on $`(P_1P_1)M`$ by<sup>19</sup><sup>19</sup>19In particular, if $`\sigma =\mu _1\mu _2`$ with $`\mu _1,\mu _2P_1`$, then $`\mathrm{dist}_\mathrm{L}^{}(\sigma _+,\sigma _{})=\mathrm{dist}_\mathrm{L}^{}(\mu _1,\mu _2)`$; note, however, that generally $`\mu _1(\mu _1\mu _2)_+`$ and $`\mu _2(\mu _1\mu _2)_{}`$. $`\sigma _\mathrm{L}^{}:=\mathrm{dist}_\mathrm{L}^{}(\sigma _+,\sigma _{})`$. This definition extends identically to $`\lambda (P_1P_1)`$ for any $`\lambda `$. To extend $`._\mathrm{L}^{}`$ to the linear span of $`P_1`$ for $`\sigma `$ lsp $`P_1`$ we define
$$\sigma _{\stackrel{~}{\mathrm{L}^{}}}:=\mathrm{dist}_\mathrm{L}^{}((\sigma \sigma (^d)\stackrel{~}{\mu })_+,(\sigma \sigma (^d)\stackrel{~}{\mu })_{})+|\sigma (^d)|$$
(280)
where $`\stackrel{~}{\mu }P_1(^d)`$ is arbitrary but fixed; e.g. $`\stackrel{~}{\mu }=\delta _0`$. Clearly, for $`\sigma P_1P_1`$, such that $`\sigma (^d)=0`$, (280) reduces to $`\sigma _{\stackrel{~}{\mathrm{L}^{}}}=\mathrm{dist}_\mathrm{L}^{}(\sigma _+,\sigma _{})`$, i.e. $`\sigma _{\stackrel{~}{\mathrm{L}^{}}}=\sigma _\mathrm{L}^{}`$ whenever $`\sigma (^d)=0`$. It is straightforward to verify that $`._{\stackrel{~}{\mathrm{L}^{}}}`$ is a norm on lsp $`P_1`$. The completion of the linear span of $`P_1(^d)`$ w.r.t. (280), denoted $`\stackrel{~}{M}_1(^d)`$, is a Banach space with norm $`._{\stackrel{~}{\mathrm{L}^{}}}`$ given in (280). We write $`\stackrel{~}{P}_1(^d)`$ for $`P_1(^d)\stackrel{~}{M}_1(^d)`$.
### A.2 The second order variant of the Gronwall lemma
The standard Gronwall lemma provides a simple upper bound on a function $`tu(t)`$ satisfying the first order differential inequality
$$\frac{\mathrm{d}}{\mathrm{d}t}uf(t)u+g(t)$$
(281)
for all $`t_+`$, with $`u(0)=u_0>0`$, and with $`f(t)`$ and $`g(t)`$ given positive continuous functions; namely, with the help of an integrating factor one finds right away that $`u`$ is bounded by
$$u(t)u_0\mathrm{exp}\left(_0^tf(\tau )d\tau \right)+_0^t\mathrm{exp}\left(_\tau ^tf(\stackrel{~}{\tau })d\stackrel{~}{\tau }\right)g(\tau )d\tau .$$
(282)
In particular, if $`f(t)\gamma >0`$ is a constant, then
$$u(t)u_0\mathrm{exp}(\gamma t)+_0^t\mathrm{exp}[\gamma (t\tau )]g(\tau )d\tau .$$
(283)
However, (282) does not suit our purposes; instead, we need the following second order variant of (282):
Lemma A1: Let $`\gamma >0`$ be a given constant and $`g(t)`$ a given positive continuous function. Suppose $`tu(t)`$ satisfies the second order differential inequality
$$\frac{\mathrm{d}^2}{\mathrm{d}t^2}u\gamma ^2u+g(t)$$
(284)
for all $`t_+`$, with $`u(0)=u_00`$ and $`u^{}(0)=v_00`$. Then $`u`$ is bounded by
$$u(t)u_0\mathrm{cosh}(\gamma t)+v_0\frac{1}{\gamma }\mathrm{sinh}(\gamma t)+_0^t\mathrm{cosh}[\gamma (t\tau )]_0^\tau g(\stackrel{~}{\tau })d\stackrel{~}{\tau }d\tau $$
(285)
for all $`t_+`$.
Proof of Lemma A.1: Denote r.h.s.(285) $`=U(t)=U_{hom}(t)+U_{inh}(t)`$, where $`U_{inh}(t)`$ is the term linear in $`g`$. By direct computation one verifies that the function $`tU(t)`$ satisfies (284) with โ$`=`$โ instead of $``$, and $`U(0)=u_0`$ and $`U^{}(0)=v_0`$. Since the Cauchy problem for (284) with positive data has a unique positive solution, it follows that $`u(t)U(t)`$ by the usual subsolution argument. Q.E.D.
### A.3 Proof of the conservation laws
We prove the conservation laws for the regularized wave gravity Vlasov equations. The laws for the microscopic regularized field & $`N`$-body systems are included as a special case. For general background material on conservation laws, see \[SuMu74\].
Proof of Proposition 4.6. The conservation of $`๐^{(\alpha )}()`$ holds because (64) is a continuity equation in $`^6`$, and because a Hamiltonian vector field is divergence-free. Q.E.D.
Proof of the conservation laws of Theorems 3.2 and 3.3. As to the conservation of $`()`$, for the time derivative of the matter energy (i.e. kinetic plus rest) we have
$`{\displaystyle \frac{\mathrm{d}}{\mathrm{d}t}}{\displaystyle \sqrt{1+|p|^2}f(x,p,t)dxdp}`$ $`=`$ $`{\displaystyle \sqrt{1+|p|^2}(\varrho \psi )(x,t)_p^{}f(x,p,t)\mathrm{d}p\mathrm{d}x}`$ (286)
$`=`$ $`{\displaystyle (\varrho \psi )(x,t)vf(x,p,t)dpdx},`$
where we first pulled the time derivative into the integral, then used (64) to rewrite the integrand, noted that $`x`$ divergences integrate to zero, then integrated by parts w.r.t. $`p`$, and finally used that $`_p^{}\sqrt{1+|p|^2}=v`$ is the velocity of a particle with unit mass, having momentum $`p`$. On the other hand, for the wave field energy, we have
$`{\displaystyle \frac{\mathrm{d}}{\mathrm{d}t}}{\displaystyle \frac{1}{2}}{\displaystyle (|_x^{}\psi |^2+|\varpi |^2)(x,t)dx}`$ $`=`$ $`{\displaystyle (_x^2\psi +_t^{}\varpi )(x,t)\varpi (x,t)dx}`$ (287)
$`=`$ $`{\displaystyle }(\varrho {\displaystyle }f(.,p,t)\mathrm{d}p)(x)\varpi (x,t)\mathrm{d}x`$
$`=`$ $`{\displaystyle f(x,p,t)(\varrho \varpi )(x,t)dxdp},`$
where we pulled the time derivative into the integral, used (60) to rewrite the integrand, and invoked Fubini. Finally, for the regularized coupling energy, we have
$`{\displaystyle \frac{\mathrm{d}}{\mathrm{d}t}}{\displaystyle (\varrho \psi )(x,t)f(x,p,t)dxdp}`$ $`=`$ $`{\displaystyle (\varrho \varpi )(x,t)f(x,p,t)dxdp}`$ (288)
$`+{\displaystyle (\varrho \psi )(x,t)_t^{}f(x,p,t)\mathrm{d}x\mathrm{d}p}.`$
The last expression in (287) cancels against the first term on r.h.s.(288). It remains to show that the second term on r.h.s.(288) cancels against the final expression in (286). We use (64) to rewrite the integrand of the second term on r.h.s.(288), note that $`p`$ divergences integrate to zero, then invoke Fubini and integrations by parts. Thus
$`{\displaystyle (\varrho \psi )(x,t)_t^{}f(x,p,t)\mathrm{d}x\mathrm{d}p}`$ $`=`$ $`{\displaystyle }(\varrho {\displaystyle }vf(.,p,t)\mathrm{d}p)(x)\psi (x,t)\mathrm{d}x`$ (289)
$`=`$ $`{\displaystyle v_x^{}f(x,p,t)\mathrm{d}p\left(\varrho \psi \right)(x,t)\mathrm{d}x}`$
$`=`$ $`{\displaystyle \left(_x^{}vf(x,p,t)dp\right)(\varrho \psi )(x,t)dx}`$
$`=`$ $`{\displaystyle vf(x,p,t)dp_x^{}(\varrho \psi )(x,t)\mathrm{d}x}`$
$`=`$ $`{\displaystyle vf(x,p,t)dp(\varrho \psi )(x,t)dx}.`$
Thus conservation of the energy $``$ is proved.
The conservation of $`๐ซ()`$ is shown similarly. For the matter momentum, we have
$`{\displaystyle \frac{\mathrm{d}}{\mathrm{d}t}}{\displaystyle pf(x,p,t)dpdx}`$ $`=`$ $`{\displaystyle p(\varrho \psi )(x,t)_p^{}f(x,p,t)\mathrm{d}p\mathrm{d}x}`$ (290)
$`=`$ $`{\displaystyle (\varrho \psi )(x,t)f(x,p,t)dpdx},`$
the last step through integration by parts, using the identity $`(_x^{}u(x)_p^{})p=_x^{}u(x)`$. On the other hand, for the field momentum we have
$`{\displaystyle \frac{\mathrm{d}}{\mathrm{d}t}}\left({\displaystyle _x^{}\psi (x,t)\varpi (x,t)\mathrm{d}x}\right)`$ $`=`$ $`{\displaystyle _t^{}\varpi (x,t)_x^{}\psi (x,t)\mathrm{d}x}`$ (291)
$`=`$ $`{\displaystyle \left(_t^{}\varpi _x^2\psi \right)(x,t)_x^{}\psi (x,t)\mathrm{d}x}`$
$`=`$ $`{\displaystyle }(\varrho {\displaystyle }f(.,p,t)\mathrm{d}p)(x)_x^{}\psi (x,t)\mathrm{d}x`$
$`=`$ $`{\displaystyle f(x,p,t)(\varrho \psi )(x,t)dpdx},`$
where we used the identity $`2\varpi _x^{}\varpi =_x^{}\left(|\varpi |^2\right)`$ in the first step, and the identity
$$_x^{}\psi _x^2\psi =_x^{}\left(_x^{}\psi _x^{}\psi \right)\frac{1}{2}_x^{}|_x^{}\psi |^2$$
(292)
in the second step, and noting the vanishing of โsurface integrals at infinity.โ Adding (290) and (291) we see that total momentum $`๐ซ`$ is conserved.
As to the conservation of $`๐ฅ()`$, for the orbital matter angular momentum we have
$`{\displaystyle \frac{\mathrm{d}}{\mathrm{d}t}}{\displaystyle x\times pf(x,p,t)dpdx}`$ $`=`$ $`{\displaystyle (x\times p)v_x^{}f(x,p,t)\mathrm{d}p\mathrm{d}x}`$ (294)
$`+{\displaystyle (x\times p)(\varrho \psi )(x,t)_p^{}f(x,p,t)\mathrm{d}p\mathrm{d}x}`$
$`=`$ $`{\displaystyle (\varrho \psi )(x,t)\times xf(x,p,t)dpdx}`$ (295)
$`=`$ $`{\displaystyle }_x^{}\psi (x,t)\times x(\varrho {\displaystyle }f(.,p,t)\mathrm{d}p)(x)\mathrm{d}x`$ (297)
$`{\displaystyle }_x^{}\psi (x,t)\times (\varrho \mathrm{Id}{\displaystyle }f(.,p,t)\mathrm{d}p)(x)\mathrm{d}x.`$
In the first step we used the continuity equation (64), in the second step integrations by parts and the identities $`(v_x^{})(x\times p)=v\times p=0`$ and $`(_x^{}u(x)_p^{})(x\times p)=x\times _x^{}u(x)`$; the last step is Fubini and a trivial rewriting. The last integral in (297) gives
$`{\displaystyle }_x^{}\psi (x,t)\times (\varrho \mathrm{Id}{\displaystyle }f(.,p,t)\mathrm{d}p)(x)\mathrm{d}x={\displaystyle }{\displaystyle }(\varrho \mathrm{Id}\times \psi (.,t))(x)f(x,p,t)\mathrm{d}p\mathrm{d}x`$ (298)
where again we used Fubini. Finally, by some standard identities of vector analysis and the radial symmetry of $`\varrho `$, the (negative of the) field torque evaluates to
$$(\varrho \mathrm{Id}\times \psi (.,t))(x)=_y^{}\times \left((yx)\varrho (yx)\psi (y,t)\right)\mathrm{d}y=0.$$
(299)
The last integral vanishes by one of Greenโs theorems and the compact support of $`\varrho `$.
Lastly, for the field angular momentum we have
$`{\displaystyle \frac{\mathrm{d}}{\mathrm{d}t}}{\displaystyle (x\times _x^{}\psi (x,t))\varpi (x,t)dx}`$ $`=`$ $`{\displaystyle (x\times _x^{}\psi (x,t))_t^{}\varpi (x,t)\mathrm{d}x}`$ (300)
$`=`$ $`{\displaystyle (x\times _x^{}\psi (x,t))(_x^2\psi _t^{}\varpi )(x,t)dx}`$ (301)
$`=`$ $`{\displaystyle }_x^{}\psi (x,t)\times x(\varrho {\displaystyle }f(.,p,t)\mathrm{d}p)(x)\mathrm{d}x`$ (302)
where we used the identity $`2\varpi x\times _x^{}\varpi =_x^{}\times \left(x|\varpi |^2\right)`$ in the first step, and the identity
$$(x\times _x^{}\psi )_x^2\psi =_x^{}\left((x\times _x^{}\psi )_x^{}\psi \right)_x^{}\times \left(\frac{1}{2}x|_x^{}\psi |^2\right)$$
(303)
in the second, and (60) in the third, noting the vanishing of โsurface integrals.โ
Adding (297) and (302), noting (298), (303), proves conservation of $`๐ฅ`$. Q.E.D.
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# Three-Dimensional Chern-Simons and BF Theories
## 1 Introduction.
Chern-Simons theory gives an interesting example of topological field theory. Its Lagrangian 3-form lives on a principal $`G`$-bundle and after pulling back to space-time (base) manifold provides, in general, a family of local, non-covariant Lagrangian densities . <sup>1</sup><sup>1</sup>1However, the corresponding Euler-Lagrange equations of motion have well-defined global meaning. Because of this, it is also more difficult to analyze, in this case, Nรถether conserved quantities . A more standard approach to the problem of symmetries and conservation laws has been applied in the so called covariant formalism . It exploits the transgression 3-form as a global and covariant Chern-Simons Lagrangian with two dynamical gauge fields. This formalism has been used for the calculation of Nรถther currents and their identically vanishing parts - superpotentials. Augmented variational principle and relative conservation laws have been recently proposed in . Our aim in the present note, which can be viewed as an appendix to , is to explain a link between covariant Chern-Simons theory and the so called BF theories .
## 2 Change of variables.
Let us consider a principal bundle $`P(M,G)`$ over a three-dimensional manifold $`M`$ with a (semisimple) structure group $`G`$. Let $`\omega _i`$ ($`i=0,1`$) be two principal connection 1-forms with the corresponding curvature 2-forms
$`\mathrm{\Omega }_i=d\omega _i+\omega _i^2=d\omega _i+{\displaystyle \frac{1}{2}}[\omega _i,\omega _i].`$ (1)
Denote by $`\alpha =\omega _1\omega _0`$, a tensorial 1-form.
The transgression 3-form is given by the well known formula
$`Q(\omega _1,\omega _0)=Q(\omega _0,\omega _1)=tr\left(2\mathrm{\Omega }_0\alpha +D_0\alpha \alpha +{\displaystyle \frac{2}{3}}\alpha ^3\right)`$
$`=tr\left(2\mathrm{\Omega }_1\alpha D_1\alpha \alpha +{\displaystyle \frac{2}{3}}\alpha ^3\right)`$ (2)
where $`D_i\alpha =d\alpha +[\omega _i,\alpha ]`$ denotes the covariant derivative of $`\alpha `$ with respect to the connection $`\omega _i`$. Thus $`Q(\omega _1,\omega _0)`$ is a tensorial (covariant) object which well-defines the corresponding global 3-form on $`M`$. It undergoes a non-covariant splitting as a difference of two Chern-Simons Lagrangians
$`Q(\omega _1,\omega _0)=CS(\omega _1)CS(\omega _0)+dtr\left(\omega _0\omega _1\right)`$ (3)
where
$`CS(\omega )=tr\left(\mathrm{\Omega }\omega {\displaystyle \frac{1}{3}}\omega ^3\right)`$ (4)
stands for non-covariant Chern-Simons Lagrangian.
Notice that in the case of two connections one has
$`2\mathrm{\Omega }_0+D_0\alpha =2\mathrm{\Omega }_1D_1\alpha =`$
$`\mathrm{\Omega }_0+\mathrm{\Omega }_1+{\displaystyle \frac{1}{2}}\left(D_0\alpha D_1\alpha \right)=\mathrm{\Omega }_0+\mathrm{\Omega }_1\alpha ^2`$ (5)
The last equality entitles us to rewrite
$`Q(\omega _1,\omega _0)=2tr\left(\overline{\mathrm{\Omega }}\alpha +{\displaystyle \frac{1}{12}}\alpha ^3\right)`$ (6)
where $`\overline{\omega }=\frac{1}{2}\left(\omega _1+\omega _0\right)`$ is a new (average) connection and $`\overline{\mathrm{\Omega }}=d\overline{\omega }+\overline{\omega }^2`$. Of course, one has $`\omega _1=\overline{\omega }+\frac{1}{2}\alpha `$ , $`\omega _0=\overline{\omega }\frac{1}{2}\alpha `$.
Thus the Lagrangian $`Q(\omega _1,\omega _0)`$ can be treated in three different (but equivalent) ways:
* with two (flat) connections $`\omega _0,\omega _1`$ as dynamical variables; see (3). In this case $`\mathrm{\Omega }_0=0`$ and $`\mathrm{\Omega }_1=0`$ are equations of motion. This point of view was presented in .
* with a (flat) connection $`\omega _1`$ and tensorial 1-form $`\alpha `$ as dynamical variables; see (2). In this case equations of motion are $`\mathrm{\Omega }_1=0`$ and $`D_1\alpha =\alpha ^2`$.
* with an โaverageโ (non-flat) connection $`\overline{\omega }=\frac{1}{2}\left(\omega _1+\omega _2\right)`$ and tensorial 1-form $`\alpha `$ as independent dynamical variables; see (6). In this case $`\overline{\mathrm{\Omega }}=\frac{1}{4}\alpha ^2`$ and $`\overline{D}\alpha =0`$ are equations of motion. This is the so called BF theory with a cosmological constant $`\mathrm{\Lambda }=1`$ (see e.g references ).
More generally, one can define a new connection $`\omega _t=t\omega _1+\left(1t\right)\omega _0=\omega _0+t\alpha `$ as a convex combination of two connections with parameter $`0t1`$. The inverse transformation is $`\omega _0=\omega _tt\alpha `$ , $`\omega _1=\omega _t+(1t)\alpha `$. In this case
$$\mathrm{\Omega }_t=t\mathrm{\Omega }_1+\left(1t\right)\mathrm{\Omega }_0t\left(1t\right)\alpha ^2.$$
Now the equation (2) can be replaced by the more general one
$`2\mathrm{\Omega }_1D_1\alpha =2\mathrm{\Omega }_0+D_0\alpha =`$ (7)
$`=2\mathrm{\Omega }_t2t\left(1t\right)\alpha ^2\left(2t1\right)D_t\alpha `$
(notice that $`t\omega _1+\left(t1\right)\omega _0=\left(2t1\right)\omega _t+2t\left(1t\right)\alpha `$).
In this new variables $`(\omega _t,\alpha )`$ we obtain
$`Q(\omega _1,\omega _0)=2tr\left(\mathrm{\Omega }_t\alpha \left(t{\displaystyle \frac{1}{2}}\right)D_t\alpha \alpha +({\displaystyle \frac{1}{3}}t+t^2)\alpha ^3\right)`$ (8)
The corresponding equations of motion are:
$`\mathrm{\Omega }_t=t(1t)\alpha ^2,D_t\alpha =(2t1)\alpha ^2`$ (9)
Thus the choices $`t=\frac{1}{2},0,1`$ lead to the simplest formulae.
It is interesting to observe that the superpotential related to (infinitesimal) gauge transformation $`\chi `$ remains independent of the choice of variables $`(\omega _t,\alpha )`$ (compare formula (19) in ), i.e.:
$$U(\chi )=tr(\alpha \chi ).$$
Instead, an explicite expression for the superpotential related to (infinitesimal) diffeomorphism transformation driven by a vectorfield $`\xi `$ does depend on the variables $`(\omega _t,\alpha )`$ and equals to (compare formulae (23,25) in ):
$$U(\xi )=tr[\alpha (2\omega _t(\xi )+(12t)\alpha (\xi ))].$$
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# Adiabatic pumping through interacting quantum dots
## Abstract
We present a general formalism to study adiabatic pumping through interacting quantum dots. We derive a formula that relates the pumped charge to the local, instantaneous Greenโs function of the dot. This formula is then applied to the infinite-$`U`$ Anderson model both for weak and strong tunnel-coupling strengths.
Introduction. The idea of producing a DC current at zero bias voltage by changing some parameters of a conductor periodically in time dates back to the work of Thouless thouless . This method of exploiting the explicit time dependence of the Hamiltonian of the system is known as pumping. If the parameters change slowly as compared to all internal time scales of the system, pumping is adiabatic, and the average transmitted charge does not depend on the detailed time dependence of the parameters. For non-interacting mesoscopic systems, Brouwer brouwer , using the concept of emissivity proposed by Bรผttiker et al. emissivity , related the charge pumped in a period to the derivatives of the instantaneous scattering matrix of the conductor with respect to the time-varying parameters. In case of noninteracting electrons, a general framework for the computation of the pumped charge has been developed brouwer ; zhou . Pumping through open quantum dots has also been investigated experimentally marcus .
The situation is profoundly different for pumping through interacting systems. In fact, there are only few works that address this problem pothier ; aono ; cota with methods suited to tackle specific systems or regimes. As far as pumping through interacting quantum dots is concerned, the work by Aono aono exploits the zero-temperature mapping of the Kondo problem raikh to a noninteracting system and uses the noninteracting formalism. On the other hand, Cota et al. cota study adiabatic pumping in a double-dot system in the sequential tunneling limit.
The aim of this Letter is to derive a formula for the charge pumped through an interacting quantum dot, which is valid from the high-temperature limit where sequential tunneling dominates down to low temperatures where Kondo correlations are relevant.
Model and formalism. We consider a single-level quantum dot coupled to two noninteracting leads. The system is described by the Hamiltonian
$$H=H_{\text{leads}}+H_{\text{dot}}+H_{\text{tun}},$$
(1)
with $`H_{\text{leads}}=_{k,\sigma ,\alpha }ฯต_\alpha (k)c_{\sigma k\alpha }^{}c_{\sigma k\alpha },`$ where $`c_{\sigma k\alpha }`$ ($`c_{\sigma k\alpha }^{}`$) is the fermionic annihilation (creation) operator for an electron with spin $`\sigma =,`$ and momentum $`k`$ in lead $`\alpha =\text{L},\text{R}`$. The leads are assumed to be in thermal equilibrium with the same chemical potential and to have flat bands with constant density of states $`\rho _\alpha `$.
The quantum dot is described by $`H_{\text{dot}}=[ฯต+\mathrm{\Delta }ฯต(t)]_\sigma n_\sigma +Un_{}n_{}`$ with $`n_\sigma =d_\sigma ^{}d_\sigma `$, where $`d_\sigma `$ ($`d_\sigma ^{}`$) is the fermionic annihilation (creation) operator for a dot electron with spin $`\sigma `$. The level position of the dot contains a time-independent part $`ฯต`$ and a time-dependent part, $`\mathrm{\Delta }ฯต(t)`$. Coulomb interaction in the dot is described by the on-site energy $`U`$. Tunneling is modeled by $`H_{\text{tun}}=_{k,\sigma ,\alpha }\left[V_\alpha (t)c_{\sigma k\alpha }^{}d_\sigma +\text{H.c.}\right]`$ with time-dependent tunnel matrix elements $`V_\alpha (t)`$. We only allow for the modulus, but not the phase, of $`V_\alpha (t)`$ to vary in time, since a time-dependent phase would correspond to a bias voltage.
By periodically changing (at least two of) the three quantities $`V_\mathrm{L}(t)`$, $`V_\mathrm{R}(t)`$, and $`\mathrm{\Delta }ฯต(t)`$, a finite charge can be pumped through the quantum dot. The charge $`Q`$ that is pumped after one cycle $`๐ฏ`$ is connected to the time-dependent current $`J_\text{L}(t)`$ flowing through the left barrier via the relation $`Q=_0^๐ฏJ_\text{L}(\tau )๐\tau `$. The starting point for our analysis is the exact relation that expresses the current in terms of the dot Greenโs function jauho
$`J_\text{L}(t)`$ $`=`$ $`{\displaystyle \frac{2e}{\mathrm{}}}{\displaystyle \underset{\sigma }{}}\text{Im}[{\displaystyle \frac{\mathrm{\Gamma }_\text{L}(t,t)}{2}}G_{\sigma \sigma }^<(t,t)+{\displaystyle }{\displaystyle \frac{d\omega }{2\pi }}f(\omega )`$ (2)
$`{\displaystyle }dt^{}e^{i\omega (t^{}t)/\mathrm{}}\mathrm{\Gamma }_\text{L}(t^{},t)G_{\sigma \sigma }^\text{r}(t,t^{})],`$
with $`\mathrm{\Gamma }_\text{L}(t_1,t)=2\pi \rho _\text{L}V_\text{L}(t)V_\text{L}^{}(t_1)`$, and $`f(\omega )`$ is the Fermi function. The lesser, retarded, and advanced Greenโs function are defined as usual, $`G_{\sigma \sigma }^<(t,t^{})=id_\sigma ^{}(t^{})d_\sigma (t)`$, $`G_{\sigma \sigma }^\text{r}(t,t^{})=i\theta (tt^{})\{d_\sigma (t),d_\sigma ^{}(t^{})\}`$, and $`G_{\sigma \sigma }^\text{a}(t,t^{})=\left[G_{\sigma \sigma }^\text{r}(t^{},t)\right]^{}`$. The Greenโs functions are diagonal in spin space since tunneling is spin conserving. Furthermore, spin degeneracy yields $`G_{}(t,t^{})=G_{}(t,t^{})G(t,t^{})`$. We remark that the Greenโs functions $`G(t,t^{})`$ are defined with a Hamiltonian that explicitly depends on time. They are determined by the Dyson equation
$$\stackrel{ห}{G}(t,t^{})=\stackrel{ห}{g}(t,t^{})+๐t_1๐t_2\stackrel{ห}{G}(t,t_1)\stackrel{ห}{\mathrm{\Sigma }}(t_1,t_2)\stackrel{ห}{g}(t_2,t^{}),$$
(3)
in matrix notation $`\stackrel{ห}{A}=\left(\begin{array}{cc}A^\text{r}& A^<\\ 0& A^\text{a}\end{array}\right)`$ for the bare, $`\stackrel{ห}{g}`$, and full Greenโs function $`\stackrel{ห}{G}`$, and the self-energy $`\stackrel{ห}{\mathrm{\Sigma }}`$. The latter takes into account the tunnel coupling $`\mathrm{\Gamma }_\alpha (t)`$ to the leads, the Coulomb interaction $`U`$ in the dot,and the time-dependent part $`\mathrm{\Delta }ฯต(t)`$ of the level position. Note that $`\stackrel{ห}{\mathrm{\Sigma }}(t_1,t_2)=\stackrel{ห}{\mathrm{\Sigma }}(t_1,t_2,\left\{H(\tau )\right\}_{\tau [t_1,t_2]})`$ is a functional of the time-dependent Hamiltonian $`H(\tau )`$ on the interval $`[t_1,t_2]`$.
We are interested in the behavior of the self-energy and, thus, the Greenโs function for a slowly varying Hamiltonian $`H(\tau )`$. This means that the time scale over which the system parameters are varying is large compared to the lifetime of the system. To construct the adiabatic expansion of the self-energy we first linearize the time dependence of the Hamiltonian, $`H(\tau )H(t_0)+(\tau t_0)\dot{H}(t_0)`$, with respect to some fixed time $`t_0`$, and expand the self-energy up to linear order in the time derivative, where the time ordering in $`H(t_0)`$ is still done with respect to time $`\tau `$. The relation $`_{\tau _1}^{\tau _2}\tau \dot{H}(t_0)๐\tau =(\tau _1+\tau _2)/2_{\tau _1}^{\tau _2}\dot{H}(t_0)๐\tau `$, valid for each segment of time evolution between two vertices at times $`\tau _1`$ and $`\tau _2`$ in the self-energy, motivates a global replacement of the time variable $`\tau `$ with the average time $`(t_1+t_2)/2`$ in the self-energy. This replacement defines an approximation, which we refer to as the average-time approximationnote\_corrections . As a result, the dependence of the self-energy on the function $`H(\tau )`$ over the interval $`[t_1,t_2]`$ is replaced by the dependence on the three times $`t_0`$, $`t_1`$, and $`t_2`$ only, and we arrive at the adiabatic expansion $`\stackrel{ห}{\mathrm{\Sigma }}(t_1,t_2,\left\{H(\tau )\right\}_{\tau [t_1,t_2]})\stackrel{ห}{\mathrm{\Sigma }}_0(t_1,t_2,t_0)+\stackrel{ห}{\mathrm{\Sigma }}_1(t_1,t_2,t_0)`$ with
$`\stackrel{ห}{\mathrm{\Sigma }}_0(t_1,t_2,t_0)`$ $`=`$ $`\stackrel{ห}{\mathrm{\Sigma }}(t_1,t_2,\left\{H(t_0)\right\}),`$ (4)
$`\stackrel{ห}{\mathrm{\Sigma }}_1(t_1,t_2,t_0)`$ $`=`$ $`\left({\displaystyle \frac{t_1+t_2}{2}}t_0\right){\displaystyle \frac{\stackrel{ห}{\mathrm{\Sigma }}_0(t_1,t_2,t_0)}{t_0}}.`$ (5)
The lowest term in the adiabatic expansion corresponds to replacing the time-dependent Hamiltonian $`H(\tau )`$ with the constant value $`H(t_0)`$. Then, $`\stackrel{ห}{\mathrm{\Sigma }}_0(t_1,t_2,t_0)`$ depends on $`t_1`$ and $`t_2`$ only via the difference $`t_1t_2`$, and we can introduce the Fourier transform $`\stackrel{ห}{\mathrm{\Sigma }}_0(\omega ,t_0)=d(t_1t_2)\mathrm{exp}[i\omega (t_1t_2)/\mathrm{}]\stackrel{ห}{\mathrm{\Sigma }}_0(t_1,t_2,t_0)`$.
The adiabatic expansion $`\stackrel{ห}{G}(t,t^{})\stackrel{ห}{G}_0(t,t^{},t_0)+\stackrel{ห}{G}_1(t,t^{},t_0)`$ for the Greenโs function follows from that for the self-energy via the Dyson equation Eq. (3). Again, we can introduce Fourier transforms $`\stackrel{ห}{G}_{0/1}(\omega ,t_0)=d(tt^{})\mathrm{exp}[i\omega (tt^{})/\mathrm{}]\stackrel{ห}{G}_{0/1}(t,t^{},t_0)`$. Since our goal is an adiabatic expansion of the current at time $`t`$ as given in Eq. (2), we choose from now on $`t_0=t`$. This results in
$`\stackrel{ห}{G}_0(\omega ,t)`$ $`=`$ $`\left[\left(\stackrel{ห}{g}(\omega )\right)^1\stackrel{ห}{\mathrm{\Sigma }}_0(\omega ,t)\right]^1,`$ (6)
$`\stackrel{ห}{G}_1(\omega ,t)`$ $`=`$ $`i\mathrm{}{\displaystyle \frac{\stackrel{ห}{G}_0(\omega ,t)}{\omega }}{\displaystyle \frac{\stackrel{ห}{\mathrm{\Sigma }}_0(\omega ,t)}{t}}\stackrel{ห}{G}_0(\omega ,t)`$ (7)
$`+`$ $`{\displaystyle \frac{i\mathrm{}}{2}}\stackrel{ห}{G}_0(\omega ,t){\displaystyle \frac{^2\stackrel{ห}{\mathrm{\Sigma }}_0(\omega ,t)}{\omega t}}\stackrel{ห}{G}_0(\omega ,t).`$
We specify these matrix equations for the retarded and lesser part and make use of the equilibrium relations $`\mathrm{\Sigma }_0^<(\omega ,t)=2if(\omega )\text{Im}\mathrm{\Sigma }_0^\text{r}(\omega ,t)`$ and $`G_0^<(\omega ,t)=2if(\omega )\text{Im}G_0^\text{r}(\omega ,t)`$, where $`G_0^\text{r}(\omega ,t)=\left[\omega ฯต\mathrm{\Sigma }_0^\text{r}(\omega ,t)\right]^1`$. Furthermore, the adiabatic expansion for $`\mathrm{\Gamma }_\text{L}(t^{},t)`$ can be constructed as $`\mathrm{\Gamma }_\text{L}(t^{},t)\mathrm{\Gamma }_\text{L}(t)\frac{tt^{}}{2}\dot{\mathrm{\Gamma }}_\text{L}(t)`$ with $`\mathrm{\Gamma }_\text{L}(t)\mathrm{\Gamma }_\text{L}(t,t)`$. Plugging everything into Eq. (2) we find that the zeroth-order term of the adiabatic expansion for the current vanishes (as it should since it is equivalent to time-independent problem at equilibrium). The first-order correction is given by
$`J_\text{L}(t)`$ $`=`$ $`{\displaystyle \frac{e}{\pi }}{\displaystyle }d\omega ({\displaystyle \frac{f}{\omega }})\text{Re}[{\displaystyle \frac{d}{dt}}\left[\mathrm{\Gamma }_\text{L}(t)G_0^\text{r}(\omega ,t)\right]`$ (8)
$`\left(G_0^\text{r}(\omega ,t)\right)^1G_0^\text{a}(\omega ,t)].`$
A factor 2 accounts for the spin degeneracy. Equation (8) is the central result of this Letter note3 . It generalizes Brouwerโs formula brouwer to interacting quantum dots. We emphasize that this result relies on the average-time approximation for the self-energy. The latter is exact whenever the self-energy contains two vertices (either tunneling or interaction) only. This is the case for $`U=0`$ but also for $`U\mathrm{}`$ as long as the self-energy is calculated up to linear order in the tunneling coupling $`\mathrm{\Gamma }`$, as well as for arbitrary interaction at zero temperature, where the interacting problem can be mapped to a noninteracting one. We now specialize Eq. (8) to the case of weak pumping due to time-dependent tunneling barriers, $`\mathrm{\Gamma }_\alpha (t)=\overline{\mathrm{\Gamma }}_\alpha +\mathrm{\Delta }\mathrm{\Gamma }_\alpha (t)`$, where $`\left|\mathrm{\Delta }\mathrm{\Gamma }_\alpha (t)\right|\overline{\mathrm{\Gamma }}_\alpha `$ at any time $`t`$. To the lowest order in $`\mathrm{\Delta }\mathrm{\Gamma }_\alpha (t)`$ the charge $`Q=_0^๐ฏJ_\text{L}(\tau )๐\tau `$ in one period $`๐ฏ`$ is
$`Q={\displaystyle \frac{e\eta \overline{\mathrm{\Gamma }}}{\pi \overline{\mathrm{\Gamma }}_\text{L}\overline{\mathrm{\Gamma }}_\text{R}}}{\displaystyle ๐\omega \left(\frac{f}{\omega }\right)\frac{\overline{\delta }(\omega )}{\overline{\mathrm{\Gamma }}}\overline{T}(\omega )},`$ (9)
where $`\overline{\mathrm{\Gamma }}=\overline{\mathrm{\Gamma }}_\text{L}+\overline{\mathrm{\Gamma }}_\text{R}`$, and $`\eta =_0^๐ฏ\dot{\mathrm{\Delta }\mathrm{\Gamma }_\text{L}}(t)\mathrm{\Delta }\mathrm{\Gamma }_\text{R}(t)๐t`$. The symbol $`\overline{\delta }(\omega )`$ denotes the phase of the Greenโs function $`\overline{G}_0^\text{r}(\omega )=|\overline{G}_0^\text{r}(\omega )|\mathrm{exp}[i\overline{\delta }(\omega )]`$ computed with $`\mathrm{\Gamma }_\alpha (t)`$ replaced by $`\overline{\mathrm{\Gamma }}_\alpha `$, and $`\overline{T}(\omega )=2\overline{\mathrm{\Gamma }}_\text{L}\overline{\mathrm{\Gamma }}_\text{R}/\overline{\mathrm{\Gamma }}\text{Im}[\overline{G}_0^\text{r}(\omega )]`$ can be interpreted as the transmission probability through an interacting quantum dot transmission .
Examples. We now consider only weak pumping with the barriers, and we restrict ourselves to the case $`\overline{\mathrm{\Gamma }}_\text{L}=\overline{\mathrm{\Gamma }}_\text{R}=\overline{\mathrm{\Gamma }}/2`$. We start by studying the noninteracting single-level quantum dot, using Eq. (9), with $`\overline{G}_0^\text{r}(\omega )=(\omega ฯต+\frac{i}{2}\overline{\mathrm{\Gamma }})^1`$. From inspection of Eq. (9) it is clear that there is no pumping in the noninteracting case if the level is resonant ($`ฯต=0`$). In the high-temperature limit $`\beta \overline{\mathrm{\Gamma }}1`$ (with $`\beta =1/k_\text{B}T`$), the pumped charge reads $`Q=\frac{e\eta }{2}f^{\prime \prime }(ฯต)`$.
We now turn our attention to weak pumping with the barriers in the limit of large electron-electron interaction, $`U\mathrm{}`$. For temperatures larger than the Kondo temperature (defined below), we approximate the instantaneous Greenโs function of the dot within the equation-of-motion method eom . Replacing $`\mathrm{\Gamma }_\alpha (t)`$ by $`\overline{\mathrm{\Gamma }}_\alpha `$ one finds $`\overline{G}_0^\text{r}(\omega )=\left(1\overline{n}\right)\left(\omega ฯต\frac{\overline{\mathrm{\Gamma }}}{2}A(\omega )+i\frac{\overline{\mathrm{\Gamma }}}{2}[1+f(\omega )]\right)^1`$, where $`A(\omega )=\frac{1}{\pi }\left\{\psi \left(\frac{1}{2}+\frac{\beta E_\text{c}}{2\pi }\right)\text{Re}\left[\psi \left(\frac{1}{2}+i\frac{\beta \omega }{2\pi }\right)\right]\right\}`$, and $`\overline{n}=\frac{d\omega }{\pi }\text{Im}\left\{\overline{G}_0^\text{r}(\omega )\right\}f(\omega )`$ is the occupation of the level per spin, $`\psi `$ the Digamma function, and $`E_\text{c}`$ a high-energy cutoff. For the high-temperature limit, $`\beta \overline{\mathrm{\Gamma }}1`$, we obtain the analytical expression
$$Q=\frac{e\eta }{2}\left[f^{\prime \prime }(ฯต)+\frac{f^{}(ฯต)}{1+f(ฯต)}\left(f^{}(ฯต)+\frac{2A(ฯต)/\overline{\mathrm{\Gamma }}}{1+f(ฯต)}\right)\right].$$
(10)
Figure 1 shows the pumped charge as a function of the level position in the high-temperature limit. The enhancement of the pumped charge as compared to the noninteracting case, is mainly due to the fact that, in the presence of interactions, the bare level is renormalized by an amount which depends on $`\mathrm{\Gamma }(t)`$, and hence the level position becomes time dependent. The oscillation of the level increases the pump effect \[third term in Eq. (10)\] note0 . Also the fact that the level width is energy dependent, $`\frac{\overline{\mathrm{\Gamma }}}{2}[1+f(\omega )]`$, has some small effect on the pumped charge \[second term in Eq. (10)\]. We note that for $`\beta 0`$, the third term in Eq. (10) goes as $`\beta ^2\overline{\mathrm{\Gamma }}E_\text{c}`$, while the other two terms go as $`(\beta \overline{\mathrm{\Gamma }})^2`$ note1 . The shape of the curves in Fig. 1 are easily understood from the dependence of $`\overline{\delta }/\overline{\mathrm{\Gamma }}`$ around $`|\omega |k_\text{B}T`$ on the bare level position $`ฯต`$. In the noninteracting case, the scale on which the phase $`\overline{\delta }`$ varies around the level position $`ฯต`$ increases linearly with $`\overline{\mathrm{\Gamma }}`$; i.e., $`\overline{\delta }/\overline{\mathrm{\Gamma }}`$ changes sign when tuning the level position through the Fermi energy. In presence of interaction, though, the dominant mechanism is the variation of the level renormalization, which shifts $`\overline{\delta }(\omega )`$ along the $`\omega `$ axis, with no sign change in $`\overline{\delta }/\overline{\mathrm{\Gamma }}`$.
The equation-of-motion method gives qualitative, reliable results down to the Kondo temperature, given by $`k_\text{B}Tk_\text{B}T_\text{K}=\sqrt{E_\text{c}\overline{\mathrm{\Gamma }}}/2\mathrm{exp}\left(\pi |ฯต|/\overline{\mathrm{\Gamma }}\right)`$. In Fig. 2, we show the temperature dependence of the pumped charge for the interacting quantum dot, obtained by numerical integration of Eq. (9) and for comparison, the noninteracting result. At very high temperatures the pumped charge tends to zero. Decreasing the temperature the charge exhibits a maximum both in the interacting and noninteracting case. It occurs when the level position and the temperature are of the same order. Its position is determined by the spectral weight of the integrand function in Eq. (9) which falls in the energy window set by temperature through the derivative of the Fermi function. Approaching the Kondo temperature, the pumped charge in the interacting system increases rapidly, indicating Kondo correlations.
To address the limit $`TT_K`$ we resort to the slave-boson method slaveboson in the mean-field approximation in the boson field. The instantaneous dot Greenโs function can be written as $`G_0^\text{r}(\omega ,t)=\frac{\mathrm{\Gamma }_{\text{pf}}(t)}{\mathrm{\Gamma }(t)}G_{\text{pf}}^\text{r}(\omega ,t)`$, where the pseudofermion Greenโs function is $`G_{\text{pf}}^\text{r}(\omega ,t)=(\omega ฯต_{\text{pf}}(t)+i\mathrm{\Gamma }_{\text{pf}}(t)/2)^1`$. We are interested in $`\overline{G}_{\text{pf}}^\text{r}(\omega )`$, where $`\mathrm{\Gamma }_\alpha (t)`$ is replaced by $`\overline{\mathrm{\Gamma }}_\alpha `$. The renormalized level position $`\overline{ฯต}_{\text{pf}}`$ and the renormalized rate $`\overline{\mathrm{\Gamma }}_{\text{pf}}`$ have to be found as the solutions of the non-linear system of equations:
$`2\overline{\mathrm{\Gamma }}{\displaystyle \frac{d\omega }{2\pi }\text{Re}\left\{\overline{G}_{\text{pf}}^\text{r}(\omega )\right\}f(\omega )}+\overline{ฯต}_{\text{pf}}ฯต=0,`$ (11a)
$`4{\displaystyle \frac{d\omega }{2\pi }\text{Im}\left\{\overline{G}_{\text{pf}}^\text{r}(\omega )\right\}f(\omega )}=1{\displaystyle \frac{\overline{\mathrm{\Gamma }}_{\text{pf}}}{\overline{\mathrm{\Gamma }}}}.`$ (11b)
At zero temperature the pumped charge can be expressed by means of Friedelโs sum rule langreth and Eq. (9), as
$$Q=\frac{4e\eta }{\overline{\mathrm{\Gamma }}}\frac{\overline{n}}{\overline{\mathrm{\Gamma }}}\mathrm{sin}^2\left(\pi \overline{n}\right),$$
(12)
which relates the pumped charge to the average occupation per spin $`\overline{n}`$ only. The full knowledge of the latter, e.g. from numerical renormalization group, would establish an exact solution of the problem. Within the mean-field slave-boson approach we get $`\overline{n}=1/2(1\overline{\mathrm{\Gamma }}_{\text{pf}}/\overline{\mathrm{\Gamma }})`$ hewson , and $`\overline{\mathrm{\Gamma }}_{\text{pf}}`$ is computed from Eqs. (11). In the unitary limit ($`ฯต\overline{\mathrm{\Gamma }}_{\text{pf}}`$ and $`TT_\text{K}`$, such that $`\overline{n}1/2`$) the level is renormalized to resonance $`\overline{ฯต}_{\text{pf}}0`$, and the pumped charge is zero. This result is consistent with the fact that in the unitary limit the problem maps to the noninteracting dot with the level shifted to resonance, and that for the free-electron case there is no pumped charge when the level is at the Fermi energy. On the other hand, in experimentally relevant situations the renormalized level is not exactly at the Fermi energy and non-negligible charge pumping occurs.
In Fig. 3 we show the charge pumped at zero temperature obtained solving numerically Eqs. (11). As expected the charge tends to zero when the level is deep enough below the Fermi energy.
The behavior of the pumped charge around $`TT_\text{K}`$ can be obtained by performing a Sommerfeld expansion in Eq. (9), and in Eqs. (11). The pumped charge goes as $`T^2`$. Comparing the temperature behavior of the charge for $`TT_\text{K}`$ with the one for $`T`$ just above $`T_\text{K}`$ \[see inset of Fig. 2\], we expect a maximum at around $`T_\text{K}`$. Roughly speaking, this extremum is analogous to the one that occurs at higher temperatures.
Acknowledgments. We acknowledge useful discussions with E. Mucciolo, Y. Oreg, E. Sela, and F. Taddei, and support from DFG via SFB491 and GRK726 (J.K.) and from EC through grants EC-RTN Nano, EC-RTN Spintronics and EC-IST-SQUIBIT2 (M.G., J.S., and R.F.).
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# cover page title of paper
## Introduction
Let $`M`$ be a compact boundaryless Riemannian manifold of dimension $`n`$ and $`E`$ a smooth vector bundle based on $`M`$. For a classical pseudodifferential operator ($`\psi \mathrm{do}`$) $`A`$ with non-integer order acting on smooth sections of $`E`$ one can define following Kontsevich and Vishik \[KV\] and Lesch \[Le\] the canonical trace of $`A`$
$$\mathrm{TR}(A):=_M๐x\text{TR}_x(A),\text{TR}_x(A):=_{T_x^{}M}\mathrm{tr}_x(\sigma _A(x,\xi ))๐\overline{}\xi ,$$
in terms of a local classical symbol $`\sigma _A`$ and a finite-part integral $`_{T_x^{}M}`$ over the cotangent space $`T_x^{}M`$ at $`xM`$. Here, $`d\overline{}\xi =(2\pi )^nd\xi `$ with $`d\xi `$ Lebesgue measure on $`T_x^{}M^n`$, while $`\mathrm{tr}_x`$ denotes the fibrewise trace. Since the work of Seeley \[Se1\] and later of Guillemin \[Gu\], Wodzicki \[Wo\] and then Kontsevich and Vishik \[KV\], it has been known that given a holomorphic family $`zA(z)`$ of classical $`\psi \mathrm{do}`$s parametrized by a domain $`W\mathrm{I}\mathrm{C}`$, with holomorphic order $`\alpha :W`$ such that $`\alpha ^{}`$ does not vanish on
$$P:=\alpha ^1(\mathrm{Z}\mathrm{Z}[n,+\mathrm{}[),$$
then the map $`z\mathrm{TR}(A(z))`$ is a meromorphic function with no more than simple poles located in $`P`$. The complex residue at $`z_0P`$ is given by a local expression \[Wo\], \[Gu\], \[KV\]
$$\mathrm{Res}_{z=z_0}\mathrm{TR}(A(z))=\frac{1}{\alpha ^{}(z_0)}\mathrm{res}(A(z_0)),$$
(0.1)
where for a classical pseudodifferential operator $`B`$ with symbol $`\sigma _B`$
$$\mathrm{res}(B):=_M๐x\mathrm{res}_x(B),\mathrm{res}_x(B):=_{S_x^{}M}\mathrm{tr}_x\left((\sigma _B)_n(x,\xi )\right)๐\overline{}_S\xi $$
is the residue trace of $`B`$. Here, $`d\overline{}_S\xi =(2\pi )^nd_S(\xi )`$ with $`d_S(\xi )`$ the sphere measure on $`S_x^{}M=\{|\xi |=1|\xi T_x^{}M\}`$, while the subscript refers to the positively homogeneous component of the symbol of order $`n`$.
In this paper, extending the identification (0.1), we provide a complete solution to the problem of giving exact formulae for all coefficients in the Laurent expansion of $`\mathrm{TR}(A(z))`$ around each pole in terms of locally defined canonical trace and residue trace densities.
For a meromorphic function $`G`$, define its finite-part $`\mathrm{fp}_{z=z_0}G(z)`$ at $`z_0`$ to be the constant term in the Laurent expansion of $`G(z)`$ around $`z_0`$. Let $`A^{(r)}(z)=_z^rA(z)`$ be the derivative $`\psi \mathrm{do}`$ with symbol $`\sigma _{A^{(r)}(z)}:=_z^r\sigma _{A(z)}`$.
Theorem Let $`zA(z)`$ be a holomorphic family of classical $`\psi \mathrm{do}`$s of order $`\alpha (z)=qz+b`$. If $`z_0P`$ and $`q0`$, then $`\mathrm{TR}(A(z))`$ has Laurent expansion for $`z`$ near $`z_0`$
$$\mathrm{TR}(A(z))=\frac{\mathrm{res}(A(z_0))}{q}\frac{1}{(zz_0)}+\underset{k=0}{\overset{\mathrm{}}{}}\mathrm{fp}_{z=z_0}\mathrm{TR}(A^{(k)}(z))\frac{(zz_0)^k}{k!}.$$
(0.2)
Furthermore,
$$\left(\text{TR}_x\left(A^{(k)}(z_0)\right)\frac{1}{q(k+1)}\text{res}_{x,0}\left(A^{(k+1)}(z_0)\right)\right)dx$$
(0.3)
defines a global density on $`M`$ and
$$\mathrm{fp}_{z=z_0}\mathrm{TR}(A^{(k)}(z))=_M๐x\left(\text{TR}_x\left(A^{(k)}(z_0)\right)\frac{1}{q(k+1)}\text{res}_{x,0}\left(A^{(k+1)}(z_0)\right)\right).$$
(0.4)
At a point $`z_0P`$ the function $`\mathrm{TR}(A(z))`$ is holomorphic near $`z_0`$ and the Laurent expansion (0.2) reduces to the Taylor series
$$\mathrm{TR}(A(z))=\mathrm{TR}(A(z_0))+\underset{k=1}{\overset{\mathrm{}}{}}\mathrm{TR}(A^{(k)}(z_0))\frac{(zz_0)^k}{k!}.$$
It is to be emphasized here that $`A^{(r)}(z)`$ cannot be a classical $`\psi \mathrm{do}`$ for $`r>0`$, but in local coordinates is represented for $`|\xi |>0`$ by a log-polyhomogeneous symbol of the form
$$\sigma _{A^{(r)}(z)}(x,\xi )\underset{j0}{}\underset{l=0}{\overset{r}{}}\sigma (A^{(r)}(z))_{\alpha (z)j,l}(x,\xi )\mathrm{log}^l|\xi |$$
with $`\sigma (A^{(r)})_{\alpha (z)j,l}(x,\xi )`$ positively homogeneous in $`\xi `$ of degree $`\alpha (z)j`$. It follows that individually the terms in (0.3)
$$\text{TR}_x\left(A^{(k)}(z_0)\right)dx:=_{T_x^{}M}\mathrm{tr}_x(\sigma _{A^{(k)}(z_0)}(x,\xi ))๐\overline{}\xi ๐x$$
(0.5)
and
$$\text{res}_{x,0}\left(A^{(k+1)}(z_0)\right)dx:=_{S_x^{}M}\mathrm{tr}_x\left(\sigma (A^{(k+1)}(z_0))_{n,0}(x,\xi )\right)๐\overline{}_S\xi ๐x$$
(0.6)
do not in general determine globally defined densities on the manifold $`M`$ when $`r>0`$, rather it is then only the sum of terms (0.3) which integrates to a global invariant of $`M`$. (In particular, it is important to distinguish (0.6) from the higher residue trace density of \[Le\], see Remark(1.5) here.) When $`\alpha (z)=qz+b`$ is not integer valued it is known that $`\text{TR}_x(A^{(k)}(z))dx`$ does then define a global density on $`M`$; in this case, $`\text{res}_{x,0}(A^{(k+1)}(z))`$ is identically zero and (0.4) reduces to the canonical trace $`\mathrm{TR}(A^{(k)}(z))=_M๐x\text{TR}_x(A^{(k)}(z))`$ on non-integer order $`\psi \mathrm{do}`$s with log-polyhomogeneous symbol \[Le\].
These results hold more generally when $`\alpha (z)`$ is an arbitrary holomorphic function with $`\alpha ^{}(z_0)0`$ at $`z_0P`$. Then the local residue term in (0.4) is replaced by the local residue of an explicitly computable polynomial in the symbols of the operators $`A^{(k+1)}(z_0),\mathrm{},A(z_0)`$. A general formula is given in Theorem 1.20, here we state the formula just for the constant term in the Laurent expansion of $`\mathrm{TR}(A(z))`$: one has
$$\mathrm{fp}_{z=z_0}\mathrm{TR}(A(z))=_M๐x\left(\text{TR}_x\left(A(z_0)\right)\frac{1}{\alpha ^{}(z_0)}\text{res}_{x,0}\left(A^{^{}}(z_0)\right)\right)$$
(0.7)
$$+\frac{\alpha ^{\prime \prime }(z_0)}{2\alpha ^{}(z_0)^2}\mathrm{res}\left(A(z_0)\right).$$
Thus compared to (0.4), the constant term (0.7) in the expansion acquires an additional residue trace term. Moreover, the identification implies that
$$\left(\text{TR}_x\left(A(z_0)\right)\frac{1}{\alpha ^{}(z_0)}\text{res}_{x,0}\left(A^{}(z_0)\right)\right)dx$$
(0.8)
defines a global density on $`M`$ independently of the order $`\alpha (z_0)`$ of $`A(z_0)`$ (for $`z_0P`$ the locally defined residue term vanishes and (0.8) reduces to the usual canonical trace density). Though this follows from general properties of holomorphic families of canonical traces, we additionally give an elementary direct proof in Appendix A.
Applied to $`\psi \mathrm{do}`$ zeta-functions this yields formulae for a number of widely studied spectral geometric invariants. For $`Q`$ an elliptic classical pseudodifferential operator of order $`q>0`$ and with spectral cut $`\theta `$, its complex powers $`Q_\theta ^z`$ are well defined \[Se1\] and to a classical pseudodifferential operator $`A`$ of order $`\alpha \mathrm{I}\mathrm{R}`$ one can associate the holomorphic family $`A(z)=AQ_\theta ^z`$ with order function $`\alpha (z)=\alpha qz`$. The generalized zeta-function
$$z\zeta _\theta (A,Q,z):=\mathrm{TR}(AQ_\theta ^z)$$
is meromorphic on $``$ with at most simple poles in $`P:=\{\frac{\alpha j}{q}|j[n,\mathrm{})\mathrm{Z}\mathrm{Z}\}`$. It has been shown by Grubb and Seeley \[GS, Gr1\] that $`\mathrm{\Gamma }(s)\zeta _\theta (A,Q,s)`$ has pole structure
$$\mathrm{\Gamma }(s)\zeta _\theta (A,Q,s)\underset{jn}{}\frac{c_j}{s+\frac{j\alpha }{q}}\frac{\text{Tr }(A\mathrm{\Pi }_Q)}{s}+\underset{l0}{}\left(\frac{c_l^{^{}}}{(s+l)^2}+\frac{c_l^{^{\prime \prime }}}{(s+l)}\right),$$
(0.9)
where the coefficients $`c_j`$ and $`c_l^{^{}}`$ are locally determined, by finitely many homogeneous components of the local symbol, while the $`c_l^{^{\prime \prime }}`$ are globally determined. In particular, whenever
$$\frac{j\alpha }{q}:=l[0,\mathrm{})\mathrm{Z}\mathrm{Z}$$
it is shown that the sum of terms
$$c_l^{^{\prime \prime }}+c_{\alpha +lq}$$
(0.10)
is defined invariantly on the manifold $`M`$, while individually the coefficients $`c_l^{^{\prime \prime }}`$ and $`c_{\alpha +lq}`$ (which contain contributions from the terms (0.5) and (0.6) respectively) depend on the symbol structure in each local trivialization. Here, in Theorem 2.2, we compute the Laurent expansion around each of the poles of the meromorphically continued Schwartz kernel $`K_{AQ_\theta ^z}(x,x)|^{\mathrm{mer}}:=_{T_x^{}M}\sigma _{AQ_\theta ^z}(x,\xi )d\overline{}\xi `$ giving the following exact formula for (0.10). One has
$$c_l^{^{\prime \prime }}+c_{\alpha +lq}=\frac{(1)^l}{l!}_M๐x\left(\text{TR}_x(AQ^l)\frac{1}{q}\text{res}_{x,0}(AQ^l\mathrm{log}_\theta Q)\right).$$
(0.11)
The remaining coefficients in (0.9) occur as residue traces of the form (0.1). By a well known equivalence, see for example \[GS, Gr1\], when $`Q`$ is a Laplace-type operator these formulae acquire a geometric character as coefficients in the asymptotic heat trace expansion $`\text{Tr }(Ae^{tQ})_{jn}c_jt^{\frac{j\alpha }{q}}+_{l0}(c_l^{^{}}\mathrm{log}t+c_l^{^{\prime \prime }})t^l`$ as $`t0+`$.
From (0.1) (\[Wo, Gu, KV\]) $`\zeta _\theta (A,Q,z)`$ has a simple pole at $`z=0`$ with residue $`\frac{1}{q}\mathrm{res}(A)`$, which vanishes if $`\alpha \mathrm{Z}\mathrm{Z}`$. The coefficients of the full Laurent expansion of $`\zeta _\theta (A,Q,s)`$ around $`z=0`$ are given by the following formulae (Theorem 2.5).
Theorem For $`k`$, let $`\zeta _\theta ^{(k)}(A,Q,0)`$ denote the coefficient of $`z^k/k!`$ in the Laurent expansion of $`\zeta _\theta (A,Q,z)`$ around $`z=0`$. Then
$`\zeta _\theta ^{(k)}(A,Q,0)`$ $`=`$ $`(1)^k{\displaystyle _M}๐x\left(\text{TR}_x(A\mathrm{log}_\theta ^kQ){\displaystyle \frac{1}{q(k+1)}}\text{res}_{x,0}(A\mathrm{log}_\theta ^{k+1}Q)\right)`$ (0.12)
$`+(1)^{k+1}\mathrm{tr}\left(A\mathrm{log}_\theta ^kQ\mathrm{\Pi }_Q\right),`$
where $`\mathrm{\Pi }_Q`$ is a smoothing operator projector onto the generalized kernel of $`Q`$. Specifically, for a classical $`\psi \mathrm{do}`$ $`A`$ of arbitrary order
$$\left(\text{TR}_x(A)\frac{1}{q}\text{res}_{x,0}(A\mathrm{log}_\theta Q)\right)dx$$
(0.13)
is a globally defined density on $`M`$ and, setting $`\zeta _\theta (A,Q,0):=\zeta _\theta ^{(0)}(A,Q,0)`$, the constant term in the expansion around $`z=0`$ is
$$\zeta _\theta (A,Q,0)=_M๐x\left(\text{TR}_x(A)\frac{1}{q}\text{res}_{x,0}(A\mathrm{log}_\theta Q)\right)\mathrm{tr}\left(A\mathrm{\Pi }_Q\right).$$
(0.14)
When $`A`$ is a differential operator $`\zeta _\theta (A,Q,0)=lim_{z0}\zeta _\theta (A,Q,z)`$ and equation (0.14) becomes
$$\zeta _\theta (A,Q,0)=\frac{1}{q}\mathrm{res}(A\mathrm{log}_\theta Q)\mathrm{tr}\left(A\mathrm{\Pi }_Q\right).$$
(0.15)
When $`Q`$ is a differential operator and $`m`$ a non-negative integer, setting $`\zeta _\theta (Q,m):=\mathrm{fp}_{z=m}\zeta _\theta (I,Q,z)`$, one has
$$\zeta _\theta (Q,m)=\frac{1}{q}\mathrm{res}(Q^m\mathrm{log}_\theta Q)\mathrm{tr}\left(Q^m\mathrm{\Pi }_Q\right).$$
(0.16)
If $`A`$ is a $`\psi \mathrm{do}`$ of non-integer order $`\alpha \mathrm{Z}\mathrm{Z}`$ then $`0P`$ and from \[Le\] the canonical trace of $`A\mathrm{log}_\theta ^kQ`$ is defined. Then (0.12) reduces to
$$\zeta _\theta ^{(k)}(A,Q,0)=(1)^k\mathrm{TR}(A\mathrm{log}_\theta ^kQ)(1)^k\mathrm{tr}\left(A\mathrm{log}_\theta ^kQ\mathrm{\Pi }_Q\right)$$
(0.17)
and, in particular, in this case
$$\zeta _\theta (A,Q,0)=\mathrm{TR}(A)\mathrm{tr}\left(A\mathrm{\Pi }_Q\right).$$
(0.18)
Notice, that in (0.15) the term $`\mathrm{res}(A\mathrm{log}_\theta Q)=\zeta _\theta (A,Q,0)+\mathrm{tr}\left(A\mathrm{\Pi }_Q\right)`$ is locally determined, depending on only finitely many of the homogeneous terms in the local symbols of $`A`$ and $`Q`$ (\[GS\] Thm 2.7, see also \[Sc\] Prop 1.5). In the case $`A=I`$ the identity (0.15) was shown for pseudodifferential $`Q`$ in \[Sc\] and \[Gr2\], and in the particular case where $`Q`$ is an invertible positive differential operator (0.16) can be inferred from \[Lo\]. The identity (0.18) is known from \[Gr1\](Rem. (1.6)). A resolvent proof of (0.14) has been given recently in \[Gr3\].
If, on the other hand, one considers, for example, $`A(z)=AQ_\theta ^{\frac{z}{1+\mu z}}`$, then the corresponding โzeta functionโ $`\mathrm{TR}(AQ_\theta ^{\frac{z}{1+\mu z}})`$ has simple and real poles in $`\backslash \{1/\mu \}`$ and by (0.7) the constant term at $`z=0`$ has, compared to (0.14), an extra term
$$\mathrm{fp}_{z=0}\mathrm{TR}(AQ_\theta ^{\frac{z}{1+\mu z}})=_M๐x\left(\text{TR}_x(A)\frac{1}{q}\text{res}_{x,0}(A\mathrm{log}_\theta Q)\right)\mathrm{tr}\left(A\mathrm{\Pi }_Q\right)$$
$$+\frac{\mu }{q}\mathrm{res}\left(A\right).$$
The appearance here of $`\frac{\mu }{q}\mathrm{res}\left(A\right)`$ corresponds to additional terms that occur as a result of a rescaling of the cut-off parameter when expectation values are computed from Feynman diagrams using a momentum cut-off procedure, see \[Gro\]. See also Remark 1.23.
One view point to adopt on (0.7) is that it provides a defect formula for regularized traces and indeed most well known trace defect formulas \[MN, O1, CDMP, Gr2\] are an easy consequence of it. On the other hand, new more precise formulae also follow. In particular, though $`\mathrm{TR}`$ is not in general defined on the bracket $`[A,B]`$ when the bracket is of integer order, we find (Theorem 2.18) that in this case the following exact global formula holds
$$_M๐x\left(\text{TR}_x\left([A,B]\right)\frac{1}{q}\text{res}_{x,0}\left([A,B\mathrm{log}_\theta Q]\right)\right)=0$$
(0.19)
independently of the choice of $`Q`$; when $`[A,B]`$ is not of integer order (0.19) reduces to the usual trace property of the canonical trace $`\mathrm{TR}([A,B])=0`$, see Section 2.
Looking at the next term up in the Laurent expansion of $`\mathrm{TR}(Q^z)`$ at zero, equation (0.12) provides an explicit formula for the $`\zeta `$-determinant
$$\mathrm{det}_{\zeta ,\theta }Q=\mathrm{exp}(\zeta _\theta ^{}(Q,0)),$$
where $`\zeta _\theta ^{}(Q,0)=_z\zeta _\theta (Q,z))_{|_{z=0}}`$, of an invertible elliptic classical pseudodifferential operator $`Q`$ of positive order $`q`$ and with spectral cut $`\theta `$. The zeta determinant is a complicated non-local invariant which has been studied in diverse mathematical contexts. From (0.12) one finds (Theorem 2.11):
Theorem
$$\mathrm{log}\mathrm{det}_{\zeta ,\theta }(Q)=_M๐x\left(\text{TR}_x\left(\mathrm{log}_\theta Q\right)\frac{1}{2q}\text{res}_{x,0}\left(\mathrm{log}_\theta ^2Q\right)\right).$$
(0.20)
A slightly modified formula holds for non-invertible $`Q`$. Notice, here, that $`\mathrm{TR}`$ of $`\mathrm{log}_\theta Q`$ does not generally exist; if $`\text{res}_{x,0}(\mathrm{log}_\theta ^2Q)=0`$ pointwise it is defined, and then $`\mathrm{log}\mathrm{det}_{\zeta ,\theta }(Q)=\mathrm{TR}(\mathrm{log}_\theta Q)`$, which holds for example for odd-class operators of even order, such as differential operators of even order on odd-dimensional manifolds \[KV, O2\]. Equation (0.20) leads to explicit formulae for the multiplicative anomaly.
## 1. Finite-part integrals (and canonical traces) of holomorphic families of classical symbols (and pseudodifferential operators)
### 1.1. Classical and log-polyhomogeneous symbols
We briefly recall some notions concerning symbols and pseudodifferential operators and fix the corresponding notations. Classical references for the polyhomogeneous symbol calculus are e.g. \[Gi\], \[GS\], \[Ho\], \[Se2\], \[Sh\], and for the extension to log-polyhomogeneous symbols \[Le\]. $`E`$ denotes a smooth hermitian vector bundle based on some closed Riemannian manifold $`M`$. The space $`C^{\mathrm{}}(M,E)`$ of smooth sections of $`E`$ is endowed with the inner product $`\psi ,\varphi :=_M๐\mu (x)\psi (x),\varphi (x)_x`$ induced by the hermitian structure $`,_x`$ on the fibre over $`xM`$ and the Riemannian measure $`\mu `$ on $`M`$. $`H^s(M,E)`$ denotes the $`H^s`$-Sobolev closure of the space $`C^{\mathrm{}}(M,E)`$.
Given an open subset $`U`$ of $`\mathrm{I}\mathrm{R}^n`$ and an auxiliary (finite-dimensional) normed vector space $`V`$, the set of symbols $`\mathrm{S}^r(U,V)`$ on $`U`$ of order $`r`$ consists of those functions $`\sigma (x,\xi )`$ in $`C^{\mathrm{}}(T^{}U,\mathrm{End}(V))`$ such that $`_x^\mu _\xi ^\nu \sigma (x,\xi )`$ is $`O((1+|\xi |)^{r|\nu |})`$ for all multi-indices $`\mu ,\nu `$, uniformly in $`\xi `$, and, on compact subsets of $`U`$, uniformly in $`x`$. We set $`\mathrm{S}(U,V):=_{r\mathrm{I}\mathrm{R}}\mathrm{S}^r(U,V)`$ and $`\mathrm{S}^{\mathrm{}}(U,V):=_{r\mathrm{I}\mathrm{R}}\mathrm{S}^r(U,V)`$. A classical (1-step polyhomogeneous) symbol of order $`\alpha `$ means a function $`\sigma (x,\xi )`$ in $`C^{\mathrm{}}(T^{}U,\mathrm{End}(V))`$ such that for each $`N\mathrm{I}\mathrm{N}`$ and each integer $`0jN`$ there exists $`\sigma _{\alpha j}C^{\mathrm{}}(T^{}U,\mathrm{End}(V))`$ which is homogeneous in $`\xi `$ of degree $`\alpha j`$ for $`|\xi |1`$, so $`\sigma _{\alpha j}(x,t\xi )=t^{\alpha j}\sigma _{\alpha j}(x,\xi )`$ for $`t1,|\xi |1`$, and a symbol $`\sigma _{(N)}\mathrm{S}^{\mathrm{Re}(\alpha )N1}(U,V)`$ such that
$$\sigma (x,\xi )=\underset{j=0}{\overset{N}{}}\sigma _{\alpha j}(x,\xi )+\sigma _{(N)}(x,\xi )(x,\xi )T^{}U.$$
(1.1)
We then write $`\sigma (x,\xi )_{j=0}^{\mathrm{}}\sigma _{\alpha j}(x,\xi ).`$ Let $`\text{CS}(U,V)`$ denote the class of classical symbols on $`U`$ with values in $`V`$ and let $`\text{CS}^\alpha (U,V)`$ denote the subset of classical symbols of order $`\alpha `$. When $`V=\mathrm{I}\mathrm{C}`$, we write $`\mathrm{S}^r(U)`$, $`\text{CS}^\alpha (U)`$, and so forth; for brevity we may omit the $`V`$ in the statement of some results. A $`\psi \mathrm{do}`$ which for a given atlas on $`M`$ has a classical symbol in the local coordinates defined by each chart is called classical, this is independent of the choice of atlas. Let $`\text{Cl}(M,E)`$ denote the algebra of classical $`\psi \mathrm{do}`$s acting on $`C^{\mathrm{}}(M,E)`$ and let $`\text{Ell}(M,E)`$ be the subalgebra of elliptic operators. For any $`\alpha \mathrm{I}\mathrm{C}`$ let $`\text{Cl}^\alpha (M,E)`$, resp. $`\text{Ell}^\alpha (M,E)`$, denote the subset of operators in $`\text{Cl}(M,E)`$, resp. $`\text{Ell}(M,E)`$, of order $`\alpha `$. With $`\mathrm{I}\mathrm{R}_+=(0,\mathrm{})`$, set $`\text{Ell}_{\mathrm{ord}>0}(M,E):=_{r\mathrm{I}\mathrm{R}_+}\text{Ell}^r(M,E)`$.
To deal with derivatives of complex powers of classical $`\psi \mathrm{do}`$s one considers the larger class of $`\psi \mathrm{do}`$s with log-polyhomogeneous symbols. Given an open subset $`UM`$, a non-negative integer $`k`$ and a complex number $`\alpha `$, a symbol $`\sigma `$ lies in $`\text{CS}^{\alpha ,k}(U,V)`$ and is said to have order $`\alpha `$ and log degree $`k`$ if
$$\sigma (x,\xi )=\underset{j=0}{\overset{N}{}}\sigma _{\alpha j}(x,\xi )+\sigma _{(N)}(x,\xi )(x,\xi )T^{}U$$
(1.2)
where $`\sigma _{(N)}\mathrm{S}^{\mathrm{Re}(\alpha )N1+ฯต}(U,V)`$ for any $`ฯต>0`$, and
$$\sigma _{\alpha j}(\xi )=\underset{l=0}{\overset{k}{}}\sigma _{\alpha j,l}(x,\xi )\mathrm{log}^l[\xi ]\xi T_x^{}U$$
with $`\sigma _{\alpha j,l}`$ homogeneous in $`\xi `$ of degree $`\alpha j`$ for $`|\xi |1`$, and (in the notation of \[Gr1\]) $`[\xi ]`$ a strictly positive $`C^{\mathrm{}}`$ function in $`\xi `$ with $`[\xi ]=|\xi |`$ for $`|\xi |1`$. As before, in this case we write
$$\sigma (x,\xi )\underset{j=0}{\overset{\mathrm{}}{}}\sigma _{\alpha j}(x,\xi )=\underset{j=0}{\overset{\mathrm{}}{}}\underset{l=0}{\overset{k}{}}\sigma _{\alpha j,l}(x,\xi )\mathrm{log}^l[\xi ].$$
(1.3)
Then $`\text{CS}^,(U,V):=_{k=0}^{\mathrm{}}\text{CS}^{,k}(U,V)`$, where $`\text{CS}^{,k}(U,V)=_{\alpha \mathrm{I}\mathrm{C}}\text{CS}^{\alpha ,k}(U,V),`$ defines the class filtered by $`k`$ of log-polyhomogeneous symbols on $`U`$. In particular, $`\text{CS}(U,V)`$ coincides with $`\text{CS}^{,0}(U,V)`$.
Given a non-negative integer $`k`$, let $`\text{Cl}^{\alpha ,k}(M,E)`$ denote the space of pseudodifferential operators on $`C^{\mathrm{}}(M,E)`$ which in any local trivialization $`E_{|U}U\times V`$ have symbol in $`\text{CS}^{\alpha ,k}(U,V)`$. Set $`\text{Cl}^{,k}(M,E):=_{\alpha \mathrm{I}\mathrm{C}}\text{Cl}^{\alpha ,k}(M,E).`$
The following subclasses of symbols and $`\psi \mathrm{do}`$s will be of importance in what follows.
###### Definition 1.1.
A log-polyhomogeneous symbol (1.3) with integer order $`\alpha \mathrm{Z}\mathrm{Z}`$ is said to be even-even (or, more fully, to have even-even alternating parity) if for each $`j0`$
$$\sigma _{\alpha j,l}(x,\xi )=(1)^{\alpha j}\sigma _{\alpha j,l}(x,\xi )\mathrm{for}|\xi |1,$$
(1.4)
and the same holds for all derivatives in $`x`$ and $`\xi `$. It is said to be even-odd (or, more fully, to have even-odd alternating parity) if for each $`j0`$
$$\sigma _{\alpha j,l}(x,\xi )=(1)^{\alpha j1}\sigma _{\alpha j,l}(x,\xi )\mathrm{for}|\xi |1,$$
(1.5)
and the same holds for all derivatives in $`x`$ and $`\xi `$. A $`\psi \mathrm{do}`$ $`A\text{Cl}^{\alpha ,k}(M,E)`$ will be said to be even-even (resp. even-odd) if in each local trivialization any local symbol $`\sigma _A(x,\xi )\text{CS}^{\alpha ,k}(U,V)`$ representing $`A`$ (modulo smoothing operators) has even-even (resp. even-odd) parity.
Thus, an even-even symbol with even integer degree is even in $`\xi `$, while an even-odd symbol with even integer degree is odd in $`\xi `$; a similar statement holds if the symbol has odd integer degree.
###### Remark 1.2.
The terminology in Definition (1.1) follows \[Gr4\]. Kontsevich-Vishik \[KV\] studied even-even classical $`\psi \mathrm{do}`$s on odd-dimensional manifolds, calling them odd-class operators. Odd-class operators (or symbols) form an algebra and include differential operators and their parametrices. The class of operators with even-odd parity symbols on even-dimensional manifolds, which includes the modulus operator $`|A|=(A^2)^{1/2}`$ for $`A`$ a first-order elliptic self-adjoint differential operator, was introduced and studied by Grubb \[Gr1\]; this class admits similar properties with respect to traces on $`\psi \mathrm{do}`$s as the odd-class operators, though they do not form an algebra. In \[O2\] Okikiolu uses the terminology โregular parityโ and โsingular parityโ for (1.4) and (1.5).
### 1.2. Finite part integrals of symbols and the canonical trace
In order to make sense of $`_{T_x^{}M}\sigma (x,\xi )๐\overline{}\xi `$ when $`\sigma \text{CS}^{\alpha ,}(U,V)`$ is a log-polyhomogeneous symbol (the integral diverges a priori if $`\mathrm{Re}(\alpha )n`$) on a local subset $`U\mathrm{I}\mathrm{R}^n`$, one can extract a โfinite partโ when $`R\mathrm{}`$ from the integral $`_{B_x^{}(0,R)}\sigma (x,\xi )๐\overline{}\xi `$ where $`B_x^{}(0,R)T_x^{}U`$ denotes the ball centered at $`0`$ with radius $`R`$ for a given point $`xU`$.
First, though, we introduce the local residue density on log-polyhomogeneous symbols, which acts as an obstruction to the finite part integral of a classical symbol defining a global density on $`M`$ and measures the anomalous contribution to the Laurent coefficients at the poles of the finite part integral when evaluated on holomorphic families of symbols.
###### Definition 1.3.
Given an open subset $`U\mathrm{I}\mathrm{R}^n`$, the local Guillemin-Wodzicki residue is defined for $`\sigma \text{CS}^\alpha (U,V)`$ by
$$\text{res}_x(\sigma )=_{S_x^{}U}\mathrm{tr}_x\left(\sigma _n(x,\xi )\right)๐\overline{}_S\xi $$
and extends to a map $`\text{res}_{x,0}:\text{CS}^{\alpha ,k}(U,V)`$ by the same formula
$`\text{res}_{x,0}(\sigma )`$ $`=`$ $`{\displaystyle _{S_x^{}U}}\mathrm{tr}_x\left(\sigma _n(x,\xi )\right)๐\overline{}_S\xi `$
$`=`$ $`{\displaystyle \underset{l=0}{\overset{k}{}}}{\displaystyle _{S_x^{}U}}\mathrm{tr}_x\left(\sigma _{n,l}(x,\xi )\right)\mathrm{log}^l|\xi |d\overline{}_S\xi `$
$`=`$ $`{\displaystyle _{S_x^{}U}}\mathrm{tr}_x\left(\sigma _{n,0}(x,\xi )\right)๐\overline{}_S\xi .`$
When $`k>0`$ the extra subscript is included in the notation $`\text{res}_{x,0}(\sigma )`$ as a reminder that it is the residue of the log degree zero component of the symbol that is being computed. The distinction is made because when $`k>0`$ the local densities $`\text{res}_{x,0}(\sigma )dx`$ do not in general define a global density on $`M`$, due to cascading derivatives of powers of logs when changing local coordinates. When $`k=0`$, Guillemin \[Gu\] and Wodzicki \[Wo\] showed the following remarkable properties.
###### Proposition 1.4.
Let $`A\text{Cl}^\alpha (M,E)`$ be a classical $`\psi \mathrm{do}`$ represented in a local coordinate chart $`U`$ by $`\sigma \text{CS}^\alpha (U,V)`$. Then $`\mathrm{res}_x(\sigma )dx`$ determines a global density on $`M`$, that is, an element of $`C^{\mathrm{}}(M,|\mathrm{\Omega }|)`$, which defines the projectively unique trace on $`\text{Cl}^{,0}(M,E)`$.
Proofs may be found in loc.cit., and in Section 2 here. The first property means that $`\mathrm{res}_x(\sigma )dx`$ can be integrated over $`M`$. The resulting number
$$\mathrm{res}(A):=_M\mathrm{res}_x(\sigma )๐x=_M๐x_{S_x^{}M}\mathrm{tr}_x\left(\sigma _n(x,\xi )\right)๐\overline{}_S\xi $$
(1.6)
is known as the residue trace of $`A`$. The terminology refers to the trace property in Proposition 1.4 that if the manifold $`M`$ is connected and has dimension larger than $`1`$, then up to a scalar multiple (1.6) defines on $`\text{Cl}^{}(M,E)`$ the unique linear functional vanishing on commutators
$$\mathrm{res}([A,B])=0,A,B\text{Cl}^{}(M,E).$$
Notice, from its definition, that the residue trace also vanishes on operators of order $`<n`$ and on non-integer order operators.
###### Remark 1.5.
The residue trace was extended by Lesch \[Le\] to $`A\text{Cl}^{\alpha ,k}(M,E)`$ with $`k>0`$ by defining $`\mathrm{res}_k(A):=(k+1)!_M๐x_{S_x^{}M}\mathrm{tr}_x\left(\sigma _{n,k}(x,\xi )\right)๐\overline{}_S\xi .`$ For an operator with log-polyhomogeneous symbol of log degree $`k>0`$ the form $`\sigma _{n,k}(x,\xi )dx`$ defines a global density on $`M`$, a property which is not generally true for the lower log degree densities $`\sigma _{n,0}(x,\xi )dx,\mathrm{},\sigma _{n,k1}(x,\xi )dx`$ which depend on the symbol structure in each local coordinate chart. We emphasize that the higher residue is not being used in the Laurent expansions we compute here, rather the relevant object is the locally defined form $`\sigma _{n,0}(x,\xi )dx`$ which for suitable $`A`$ defines one component of a specific local density which does determine an element of $`C^{\mathrm{}}(M,\mathrm{End}(E)|\mathrm{\Omega }|).`$
It was, on the other hand, observed by Kontsevich and Vishik \[KV\] that the usual $`L^2`$-trace on $`\psi \mathrm{do}`$s of real order $`<n`$ extends to a functional on the space $`\text{Cl}^{\mathrm{I}\mathrm{C}\backslash \mathrm{Z}\mathrm{Z}}(M,E)`$ of $`\psi \mathrm{do}`$s of non-integer order and vanishes on commutators of non-integer order. Lesch \[Le\] subsequently showed that the resulting canonical trace can be further extended to
$$\text{Cl}^{\mathrm{I}\mathrm{C}\backslash \mathrm{Z}\mathrm{Z},}(M,E):=\underset{\alpha \mathrm{I}\mathrm{C}\backslash \mathrm{Z}\mathrm{Z}}{}\text{Cl}^{\alpha ,}(M,E)$$
in the following way.
###### Lemma 1.6.
Let $`U`$ be an open subset of $`\mathrm{I}\mathrm{R}^n`$ and let $`\sigma \text{CS}^{\alpha ,k}(U,V)`$ be a log-polyhomogeneous symbol of order $`\alpha `$ and log-degree $`k`$. Then for any $`xU`$ the integral $`_{B_x^{}(0,R)}\sigma (x,\xi )๐\overline{}\xi `$ has an asymptotic expansion as $`R\mathrm{}`$
$$_{B_x^{}(0,R)}\sigma (x,\xi )๐\overline{}\xi _R\mathrm{}$$
$$C_x(\sigma )+\underset{j=0,\alpha j+n0}{\overset{\mathrm{}}{}}\underset{l=0}{\overset{k}{}}P_l(\sigma _{\alpha j,l})(\mathrm{log}R)R^{\alpha j+n}+\underset{l=0}{\overset{k}{}}\frac{1}{l+1}_{S_x^{}U}\sigma _{n,l}(x,\xi )๐\overline{}_S\xi \mathrm{log}^{l+1}R$$
(1.7)
where $`P_l(\sigma _{\alpha j,l})(X)`$ is a polynomial of degree $`l`$ with coefficients depending on $`\sigma _{\alpha j,l}`$. Here $`B_x^{}(0,R)`$ stands for the ball of radius $`R`$ in the cotangent space $`T_x^{}M`$ and $`S_x^{}U`$ the unit sphere in the cotangent space $`T_x^{}U`$.
Discarding the divergences, we can therefore extract a finite part from the asymptotic expansion of $`_{B(0,R)}\sigma (x,\xi )๐\overline{}\xi `$:
###### Definition 1.7.
The finite-part integral <sup>1</sup><sup>1</sup>1This concept is closely related to partie finie of Hadamard \[Ha\], hence the terminology used here. However in the physics literature this is also known as โcut-off regularizationโ. of $`\sigma \text{CS}^{\alpha ,k}(U,V)`$ is defined to be the constant term in the asymptotic expansion (1.7)
$$_{T_x^{}U}\sigma (x,\xi )๐\overline{}\xi :=\mathrm{LIM}_R\mathrm{}_{B_x^{}(0,R)}\sigma (x,\xi )๐\overline{}\xi :=C_x(\sigma ).$$
(1.8)
The proof of the following formula \[Gr1\],\[Pa\] and of Lemma 1.6 is included in Appendix B.
###### Lemma 1.8.
For $`\sigma \text{CS}^{\alpha ,k}(U,V)`$
$`{\displaystyle _{T_x^{}U}}\sigma (x,\xi )๐\overline{}\xi `$ $`=`$ $`{\displaystyle \underset{j=0}{\overset{N}{}}}{\displaystyle _{B_x^{}(0,1)}}\sigma _{\alpha j}(x,\xi )๐\overline{}\xi +{\displaystyle _{T_x^{}U}}\sigma _{(N)}(x,\xi )๐\overline{}\xi `$ (1.9)
$`+{\displaystyle \underset{j=0,\alpha j+n0}{\overset{N}{}}}{\displaystyle \underset{l=0}{\overset{k}{}}}{\displaystyle \frac{(1)^{l+1}l!}{(\alpha j+n)^{l+1}}}{\displaystyle _{S_x^{}U}}\sigma _{\alpha j,l}(x,\xi )๐\overline{}_S\xi .`$
It is independent of $`N>\mathrm{Re}(\alpha )+n1`$.
The residue terms on the right-side of (1.9) measure anomalous behaviour in the finite-part integral. Specifically, (1.9) implies that for a rescaling $`R\mu R`$
$`\mathrm{LIM}_R\mathrm{}{\displaystyle _{B_x^{}(0,\mu R)}}\sigma (x,\xi )๐\overline{}\xi `$ $`=`$ $`\mathrm{LIM}_R\mathrm{}{\displaystyle _{B_x^{}(0,R)}}\sigma (x,\xi )๐\overline{}\xi `$ (1.10)
$`+{\displaystyle \underset{l=0}{\overset{k}{}}}{\displaystyle \frac{\mathrm{log}^{l+1}\mu }{l+1}}{\displaystyle _{S_x^{}U}}\sigma _{\alpha j,l}(x,\xi )๐\overline{}_S\xi `$
(cf. Appendix B) and hence that the finite-part integral is independent of a rescaling if $`_{S_x^{}U}\sigma _{n,l}(x,\xi )๐\overline{}_S\xi `$ vanishes for each integer $`0lk`$. More generally, just as ordinary integrals obey the transformation rule $`|\mathrm{det}C|_{\mathrm{I}\mathrm{R}^n}f(C\xi )๐\overline{}\xi =_{\mathrm{I}\mathrm{R}^n}f(\xi )๐\overline{}\xi `$, one hopes for a similar transformation rule for the regularized integral $`_{\mathrm{I}\mathrm{R}^n}\sigma (\xi )๐\overline{}\xi `$ when $`\sigma `$ is a log-polyhomogeneous symbol in order to obtain a globally defined density on $`M`$. That, however, is generally not the case in the presence of a residue as the following theorem shows.
###### Proposition 1.9.
\[Le\] The finite part integral of $`\sigma \text{CS}^{,k}(U)`$ is generally not invariant under a transformation $`CGl_n(T_x^{}U)`$. One has,
$`|\mathrm{det}C|{\displaystyle _{T_x^{}U}}\sigma (x,C\xi )d\overline{}\xi `$ $`=`$ $`{\displaystyle _{T_x^{}U}}\sigma (x,\xi )๐\overline{}\xi `$
$`+{\displaystyle \underset{l=0}{\overset{k}{}}}{\displaystyle \frac{(1)^{l+1}}{l+1}}{\displaystyle _{S_x^{}U}}\sigma _{n,l}(x,\xi )\mathrm{log}^{l+1}|C^1\xi |d\overline{}\xi .`$
###### Proof.
We refer the reader to the proof of Proposition 5.2 in \[Le\]. โ
As a consequence, whenever $`_{S_x^{}U}\sigma _{n,l}(x,\xi )\mathrm{log}^{l+1}|C^1\xi |d\overline{}\xi `$ vanishes for each integer $`0lk`$ and $`xU`$ one then recovers the usual transformation property
$$|\mathrm{det}C|_{T_x^{}U}\sigma (x,C\xi )d\overline{}\xi =_{T_x^{}U}\sigma (x,\xi )d\overline{}\xi .$$
(1.12)
With respect to a trivialization $`E_{|U}U\times V`$, a localization of $`A\text{Cl}^{\alpha ,k}(M,E)`$ in $`\text{Cl}^{\alpha ,k}(U,V)`$ can be written
$$Af(x)=_{\mathrm{I}\mathrm{R}^n}_Ue^{i(xy).\xi }\text{a}(x,y,\xi )f(y)๐y๐\overline{}\xi $$
with amplitude $`\text{a}\text{CS}^\alpha (U\times U,V)`$. Then with
$$\sigma _A(x,\xi ):=\text{a}(x,x,\xi )\text{CS}^\alpha (U,V)$$
we define
$$\text{TR}_x(A)dx:=_{T_x^{}M}\mathrm{tr}_x(\sigma _A(x,\xi ))๐\overline{}\xi ๐x.$$
If (1.12) holds for $`\sigma =\sigma _A`$ in each localization it follows that $`\text{TR}_x(A)dx`$ is independent of the choice of local coordinates. This is known in the following cases.
###### Proposition 1.10.
Let $`A\text{Cl}^{\alpha ,k}(M,E)`$. In each of the following cases $`\text{TR}_x(A)dx`$ defines an element of $`C^{\mathrm{}}(M,|\mathrm{\Omega }|)`$, that is, a global density on $`M`$:
1. $`\alpha [n,\mathrm{})\mathrm{Z}\mathrm{Z}`$,
2. $`A`$ (of integer order) is even-even and $`M`$ is odd-dimensional,
3. $`A`$ (of integer order) is even-odd and $`M`$ is even-dimensional.
Cases (1) and (2) were shown in \[KV\], where the canonical trace was first introduced, in terms of homogeneous distributions. Case (1) was reformulated in \[Le\] in terms of finite part integrals and extended to log-polyhomogeneous symbols $`k0`$. Case (3) was introduced in \[Gr1\] where it was shown that (2) and (3) may be included in the finite part integral formulation. We refer there for details. Notice though that it is easily seen that the integrals $`_{S_x^{}U}\sigma _{n,l}(x,\xi )\mathrm{log}^{l+1}|C^1\xi |d\overline{}\xi `$ vanish in each case; for (1), there is no homogeneous component of the symbol of degree $`n`$ and so the integrals vanish trivially, while setting $`g(\xi ):=\sigma _{n,l}(x,\xi )\mathrm{log}^{l+1}|C^1\xi |`$, for cases (2) and (3) one has $`g(\xi )=g(\xi )`$ and so the vanishing is immediate by symmetry.
###### Definition 1.11.
For a $`\psi \mathrm{do}`$ $`A\text{Cl}^{\alpha ,}(M,E)`$ satisfying one of the criteria (1),(2),(3) in Proposition 1.10 the canonical trace is defined by
$$\mathrm{TR}(A):=_M๐x\text{TR}_x(A).$$
The case of $`\psi \mathrm{do}`$s of non-integer order is all that is needed for the general formulae we prove here, cases (2) and (3) of Proposition 1.10 will be relevant only for applications and refinements. Case (2) in particular includes differential operators on odd-dimensional manifolds, though this holds by default in so far as $`\mathrm{TR}`$ vanishes on differential operators in any dimension (noted also in \[Gr1\]):
###### Proposition 1.12.
Let $`A\text{Cl}(M,E)`$ be a differential operator with local symbol $`\sigma _A`$, then for any $`xM`$
$$\text{TR}_x(A):=_{T_x^{}M}\mathrm{tr}_x(\sigma _A(x,\xi ))๐\overline{}\xi =0.$$
###### Proof.
Since $`A`$ is a differential operator, $`\sigma _A(x,\xi )=_{|k|=0}^{\mathrm{ord}A}\sigma _k(x,\xi )`$ with $`k=(k_1,\mathrm{},k_n)`$ a multi-index with $`k_i`$ and $`\sigma _k(x,\xi )=a_k(x)\xi ^k`$ positively homogeneous (with the previous notations we have $`\sigma _{(N)}=0`$ provided $`N\mathrm{ord}A`$). Its finite-part integral on the cotangent space at $`xM`$ therefore reads
$`{\displaystyle _{T_x^{}M}}\sigma _A(x,\xi )๐\overline{}\xi `$ $`=`$ $`{\displaystyle \underset{|k|=0}{\overset{\mathrm{ord}A}{}}}a_k(x)\mathrm{LIM}_R\mathrm{}{\displaystyle _{B_x^{}(0,R)}}\xi ^k๐\overline{}\xi `$
$`=`$ $`{\displaystyle \underset{|k|=0}{\overset{\mathrm{ord}A}{}}}a_k(x)\mathrm{LIM}_R\mathrm{}\left({\displaystyle _0^R}r^{|k|+n1}๐r\right){\displaystyle _{S_x^{}M}}\xi ^k๐\overline{}\xi `$
which vanishes since $`\mathrm{LIM}_R\mathrm{}\frac{R^{|k|+n}}{|k|+n}=0`$. โ
On commutators the canonical trace has the following more substantial vanishing properties \[KV\], \[MN\], \[Le\], \[Gr1\], providing some justification for its name.
###### Proposition 1.13.
Let $`A\text{Cl}^{a,k}(M,E)`$, $`B\text{Cl}^{b,l}(M,E)`$. In each of the cases
1. $`\alpha +\beta [n,\mathrm{})\mathrm{Z}\mathrm{Z}`$,
2. $`A`$ and $`B`$ are both even-even or are both even-odd and $`M`$ is odd-dimensional,
3. $`A`$ is even-even, $`B`$ is even-odd, and $`M`$ is even-dimensional,
the canonical trace is then defined on the commutator $`[A,B]`$ and is equal to zero,
$$\mathrm{TR}([A,B])=0.$$
The canonical trace extends the usual operator trace defined on the subalgebra $`\text{Cl}^{\mathrm{ord}<n}(M,E)`$ of $`\psi \mathrm{do}`$s of real order $`\mathrm{Re}(\alpha )<n`$, in so far as for $`\psi \mathrm{do}`$s with (real) order less than $`n`$ finite part integrals coincide with ordinary integrals. More precisely, if $`K_A(x,y)`$ denotes the Schwartz kernel of $`A\text{Cl}^{\mathrm{ord}<n}(M,E)`$ in a given localization, then $`\sigma _A(x,\xi )`$ is integrable in $`\xi `$ and $`K_A(x,x)dx=\left(_{T_x^{}M}\sigma _A(x,\xi )๐\overline{}\xi \right)dx`$ determines a global density on $`M`$, and one has
$$\mathrm{tr}(A)=_M๐x\mathrm{tr}_x\left(K_A(x,x)\right)=_M๐x_{T_x^{}M}\mathrm{tr}_x(\sigma _A(x,\xi ))๐\overline{}\xi =\mathrm{TR}(A).$$
### 1.3. Holomorphic families of symbols
We consider next families of symbols depending holomorphically on a complex parameter $`z`$. The definition is somewhat more delicate than that used in \[KV\] (or \[Le\]) since growth conditions are imposed on each $`z`$-derivative of the symbol. This is in order to maintain control of the full Laurent expansion.
First, the meaning here of holomorphic dependence on a parameter is as follows. Let $`W\mathrm{I}\mathrm{C}`$ be a complex domain, let $`Y`$ be an open subset of $`\mathrm{I}\mathrm{R}^m`$, and let $`V`$ be a vector space. A function $`p(z,\eta )C^{\mathrm{}}(W\times Y,\mathrm{End}(V))`$ is holomorphic at $`z_0W`$ if for fixed $`\eta `$ with
$$p^{(k)}(z_0,\eta )=_z^k(p(z,\eta ))|_{z=z_0}$$
there is a Taylor expansion in a neighbourhood $`N_{z_0}`$ of $`z_0`$
$$p(z,\eta )=\underset{k=0}{\overset{\mathrm{}}{}}p^{(k)}(z_0,\eta )\frac{(zz_0)^k}{k!}$$
(1.13)
which is convergent, uniformly on compact subsets of $`N_{z_0}`$, with respect to the (metrizable) topology on $`C^{\mathrm{}}(W\times Y,\mathrm{End}(V))`$ associated with the family of semi-norms
$$q_{m,K_1,K_2}=\underset{\stackrel{(z,\eta )K_1\times K_2}{r+|\mu |m}}{sup}|_z^r_\eta ^\mu q(z,\eta )|$$
(1.14)
defined for $`m\mathrm{I}\mathrm{N}`$ and compact subsets $`K_1W`$, $`K_2\mathrm{I}\mathrm{R}^m`$.
###### Definition 1.14.
Let $`m`$ be a non-negative integer, let $`U`$ be an open subset of $`\mathrm{I}\mathrm{R}^n`$, and let $`W`$ be a domain in $`\mathrm{I}\mathrm{C}`$. A holomorphic family of log-polyhomogeneous symbols parametrized by $`W`$ of order $`\alpha C^{\mathrm{}}(W,\mathrm{I}\mathrm{C})`$ and of log-degree $`m`$ means a function
$$\sigma (z)(x,\xi ):=\sigma (z,x,\xi )C^{\mathrm{}}(W\times U\times \mathrm{I}\mathrm{R}^n,\mathrm{End}V)$$
for which:
1. $`\sigma (z)(x,\xi )`$ is holomorphic at $`zW`$ as an element of $`C^{\mathrm{}}(W\times U\times \mathrm{I}\mathrm{R}^n,\mathrm{End}V)`$ and
$$\sigma (z)(x,\xi )\underset{j0}{}\sigma (z)_{\alpha (z)j}(x,\xi )\text{CS}^{\alpha (z),m}(U,V),$$
(1.15)
where the function $`\alpha :W\mathrm{I}\mathrm{C}`$ is holomorphic;
2. for any integer $`N1`$ the remainder
$$\sigma _{(N)}(z)(x,\xi ):=\sigma (z)(x,\xi )\underset{j=0}{\overset{N1}{}}\sigma _{\alpha (z)j}(z)(x,\xi )$$
is holomorphic in $`zW`$ as an element of $`C^{\mathrm{}}(W\times U\times \mathrm{I}\mathrm{R}^n,\mathrm{End}V)`$ with $`k^{\mathrm{th}}`$ $`z`$-derivative
$$\sigma _{(N)}^{(k)}(z)(x,\xi ):=_z^k(\sigma _{(N)}(z)(x,\xi ))\mathrm{S}^{\alpha (z)N+ฯต}(U,V)$$
(1.16)
for any $`ฯต>0`$.
A family $`zA(z)`$ of log-classical $`\psi \mathrm{do}`$s on $`C^{\mathrm{}}(M,E)`$ parametrized by a domain $`W`$ is holomorphic if in each local trivialisation of $`E`$ one has
$$A(z)=\mathrm{Op}(\sigma _{A(z)})+R(z)$$
with $`\sigma _{A(z)}`$ a holomorphic family of log-polyhomogeneous symbols and $`R(z)`$ a smoothing operator with Schwartz kernel $`R(z,x,y)C^{\mathrm{}}(W\times X\times X,\mathrm{End}(V))`$ holomorphic in $`z`$.
There are, of course, other ways to express these conditions; for example, in terms of the truncated kernel $`K^{(N)}(z)(x,y):=_{T_x^{}U}e^{i\xi (xy)}\sigma _{(N)}(z)(x,\xi )๐\overline{}\xi `$ with large $`N`$, and its derivatives $`_z^kK^{(N)}(z)(x,y)`$, used in the case $`k=0`$ in \[KV\] to compute the pole of $`\text{Tr }(A(z))`$ at $`z_0P`$. When dealing with the full Laurent expansion the essential requirement is that a balance be preserved between the Taylor expansion in $`z`$, in terms of the growth rates of the $`z`$-derivatives of the symbol, and the asymptotic symbol expansion in $`\xi `$.
###### Proposition 1.15.
If $`\sigma (z)(x,\xi )\text{CS}^{\alpha (z),m}(U,V)`$ is a holomorphic family of log-classical symbols, then so is each derivative
$$\sigma ^{(k)}(z)(x,\xi ):=_z^k\left(\sigma (z)(x,\xi )\right)\text{CS}^{\alpha (z),m+k}(U,V).$$
(1.17)
Precisely, $`\sigma ^{(k)}(z)(x,\xi )`$ has asymptotic expansion
$$\sigma ^{(k)}(z)(x,\xi )\underset{j0}{}\sigma ^{(k)}(z)_{\alpha (z)j}(x,\xi )$$
(1.18)
where as elements of $`_{l=0}^{m+k}\text{CS}^{\alpha (z)j,l}(U,V)`$
$$\sigma ^{(k)}(z)_{\alpha (z)j}(x,\xi )=_z^k\left(\sigma (z)_{\alpha (z)j}(x,\xi )\right).$$
(1.19)
That is,
$$\left(_z^k\sigma (z)\right)_{\alpha (z)j}(x,\xi )=_z^k\left(\sigma (z)_{\alpha (z)j}(x,\xi )\right).$$
(1.20)
###### Proof.
We have to show that
$$_z^k\left(\sigma (z)(x,\xi )\right)\underset{j0}{}_z^k\left(\sigma (z)_{\alpha (z)j}(x,\xi )\right)$$
(1.21)
where the summands are log-polyhomogeneous of the asserted order. First, the estimate
$$_z^k\left(\sigma (z)(x,\xi )\right)\underset{j=0}{\overset{N1}{}}_z^k\left(\sigma (z)_{\alpha (z)j}(x,\xi )\right)\mathrm{S}^{\alpha (z)N+ฯต}(U,V)$$
any $`ฯต>0`$, needed for (1.21) to hold is equation (1.16) of the definition. It remains to examine the form of the summands in $`_{j=0}^{N1}_z^k\left(\sigma (z)_{\alpha (z)j}(x,\xi )\right).`$ Taking differences of remainders $`\sigma _{(N)}(z)(x,\xi )`$ implies that each term $`\sigma (z)_{\alpha (z)j}(x,\xi )`$ is holomorphic. In order to compute $`_z\left(\sigma (z)_{\alpha (z)j}(x,\xi )\right)`$ one must compute the derivative of each of its homogeneous components; for $`|\xi |1`$ and any $`l\{0,\mathrm{},m\}`$
$`_z\left(\sigma _{\alpha (z)j,l}(z)(x,\xi )\right)`$ $`=`$ $`_z\left(|\xi |^{\alpha (z)j}\sigma _{\alpha (z)j,l}(z)(x,{\displaystyle \frac{\xi }{|\xi |}})\right)`$
$`=`$ $`\left(\alpha ^{}(z)|\xi |^{\alpha (z)j}\sigma _{\alpha (z)j,l}(z)(x,{\displaystyle \frac{\xi }{|\xi |}})\right)\mathrm{log}|\xi |`$
$`+|\xi |^{\alpha (z)j}_z\left(\sigma _{\alpha (z)j,l}(z)(x,{\displaystyle \frac{\xi }{|\xi |}})\right).`$
Since $`\sigma _{\alpha (z)j,l}(z)(x,\xi |\xi |^1)`$ is a symbol of constant order zero, so is its $`z`$-derivative. Hence
$$_z\left(\sigma (z)_{\alpha (z)j,l}(x,\xi )\right)=\alpha ^{}(z)\sigma (z)_{\alpha (z)j,l}(x,\xi )\mathrm{log}[\xi ]+p_{\alpha (z)j,l}(z)(x,\xi )$$
(1.22)
where $`\sigma _{\alpha (z)j,l}(z),p_{\alpha (z)j,l}(z)\text{CS}^{\alpha (z)j}(U)`$ are homogeneous in $`\xi `$ of order $`\alpha (z)j`$. Hence $`_z\left(\sigma (z)_{\alpha (z)j}\right)\text{CS}^{\alpha (z)j,m+1}(U)`$. Iterating (1.22), $`_z^k\left(\sigma _{\alpha (z)j}(z)(x,\xi )\right)`$ is thus seen to be a polynomial in $`\mathrm{log}[\xi ]`$ of the form
$$(\alpha ^{}(z))^k\sigma _{\alpha (z)j,k}(z)(x,\xi )\mathrm{log}^{k+m}[\xi ]+\mathrm{}+|\xi |^{\alpha (z)j}_z^k(\sigma _{\alpha (z)j,l}(z)(x,\frac{\xi }{|\xi |}))\mathrm{log}^0[\xi ]$$
with each coefficient homogeneous of order $`\alpha (z)j`$. This completes the proof. โ
Thus, taking derivatives adds more logarithmic terms to each term $`\sigma (z)_{\alpha (z)j}(x,\xi )`$, increasing the log-degree, but the order is unchanged. Specifically, $`\sigma ^{(k)}(z)_{\alpha (z)j}`$ takes the form
$$\sigma ^{(k)}(z)_{\alpha (z)j}(x,\xi )=\underset{l=0}{\overset{m+k}{}}\sigma ^{(k)}(z)_{\alpha (z)j,l}(x,\xi )\mathrm{log}^l[\xi ],$$
(1.23)
where the terms $`\sigma ^{(k)}(z)_{\alpha (z)j,l}(x,\xi )`$ are positively homogeneous in $`\xi `$ of degree $`\alpha (z)j`$ for $`|\xi |1`$ and can be computed explicitly from the lower order derivatives of $`\sigma (z)_{\alpha (z)j,m}(x,\xi )`$. The following more precise inductive formulae will be needed in what follows.
###### Lemma 1.16.
Let $`\sigma (z)(x,\xi )\text{CS}(U,V)`$ be a holomorphic family of classical symbols. Then for $`|\xi |1`$
$`\sigma _{\alpha (z)j,k+1}^{(k+1)}(z)(x,\xi )`$ $`=`$ $`\alpha ^{}(z)\sigma _{\alpha (z)j,k}^{(k)}(z)(x,\xi ),`$
$`\sigma _{\alpha (z)j,l}^{(k+1)}(z)(x,\xi )`$ $`=`$ $`\alpha ^{}(z)\sigma _{\alpha (z)j,l1}^{(k)}(z)(x,\xi )`$
$`+|\xi |^{\alpha (z)j}_z(\sigma _{\alpha (z)j,l}^{(k)}(z)(x,\xi /|\xi |)),1lk,`$
$`\sigma _{\alpha (z)j,\mathrm{\hspace{0.17em}0}}^{(k+1)}(z)(x,\xi )`$ $`=`$ $`|\xi |^{\alpha (z)j}_z(\sigma _{\alpha (z)j,0}^{(k)}(z)(x,\xi /|\xi |)).`$
###### Proof.
From the above
$$\sigma ^{(k)}(z)_{\alpha (z)j}(x,\xi )=_z^k(\sigma (z)_{\alpha (z)j}(x,\xi ))=\underset{l=0}{\overset{k}{}}\sigma ^{(k)}(z)_{\alpha (z)j,l}(x,\xi )\mathrm{log}^l[\xi ],$$
so that
$$\sigma ^{(k+1)}(z)_{\alpha (z)j}(x,\xi )=\underset{l=0}{\overset{k}{}}_z\left(\sigma ^{(k)}(z)_{\alpha (z)j,l}(x,\xi )\right)\mathrm{log}^l[\xi ].$$
(1.24)
Hence for $`|\xi |1`$
$$\underset{l=0}{\overset{k+1}{}}\sigma ^{(k+1)}(z)_{\alpha (z)j,l}(x,\xi )\mathrm{log}^l|\xi |=$$
$$\underset{r=0}{\overset{k}{}}\alpha ^{}(z)\sigma ^{(k)}(z)_{\alpha (z)j,r}(x,\xi )\mathrm{log}^{r+1}|\xi |+|\xi |^{\alpha (z)j}_z\left(\sigma ^{(k)}(z)_{\alpha (z)j,r}(x,\frac{\xi }{|\xi |})\right)\mathrm{log}^r|\xi |$$
where for the right-side we apply (1.22) to each of coefficient on the right-side of (1.24). Equating coefficients completes the proof. โ
A corresponding result on the level of operators follows in a straightforward manner:
###### Proposition 1.17.
Let $`zA(z)\text{Cl}^{\alpha (z),m}(M,E)`$ be a holomorphic family of log-polyhomogeneous $`\psi \mathrm{do}`$s. Then for any non-negative integer $`k`$, $`A^{(k)}(z_0)`$ lies in $`\text{Cl}^{\alpha (z_0),m+k}(M,E)`$.
###### Example 1.18.
For real numbers $`\alpha ,q`$ with $`q>0`$ the function $`\sigma (z)(x,\xi )=\psi (\xi )|\xi |^{\alpha qz}`$, where $`\psi `$ is a smooth cut-off function which vanishes near the origin and is equal to $`1`$ outside the unit ball, provides a holomorphic family of classical symbols; at any point $`z=z_0\mathrm{I}\mathrm{C}`$ we have $`\sigma ^{(k)}(z_0)(x,\xi )=(q)^k\psi (\xi )\mathrm{log}^k|\xi ||\xi |^{\alpha qz_0}`$ which lies in $`\text{CS}^{\alpha qz_0,k}(U)`$. More generally, if $`Q\text{Cl}^q(M,E)`$ is a classical elliptic $`\psi \mathrm{do}`$ of order $`q>0`$ with principal angle $`\theta `$, then one has for each $`z`$ the complex power $`Q_\theta ^z\text{Cl}^{qz}(M,E)`$ \[Se1\] represented in a local coordinate chart $`U`$ by a classical symbol $`๐ช(z)(x,\xi )\text{CS}^{qz}(U,V)`$. Let $`A\text{Cl}^\alpha (M,E)`$ be a coefficient classical $`\psi \mathrm{do}`$ represented in $`U`$ by $`๐(x,\xi )\text{CS}^\alpha (U)`$. Then $`\sigma _{AQ_\theta ^z}(x,\xi )\text{CS}^{\alpha qz}(U,V)`$ is a holomorphic family of symbols parametrized by $`W=\mathrm{I}\mathrm{C}`$ whose convergent Taylor expansion in $`C^{\mathrm{}}(\mathrm{I}\mathrm{C}\times U,V)`$ around each $`z_0\mathrm{I}\mathrm{C}`$ is from \[O1\] Lemma(2.1) given by
$$(\sigma _{AQ_\theta ^z})_{\alpha qzj}(x,\xi )=\underset{k=0}{\overset{\mathrm{}}{}}\underset{l=0}{\overset{k}{}}(1)^k\left(๐\mathrm{log}^k(๐ช)๐ช(z_0)\right)_{\alpha qz_0j,l}(x,\xi )\mathrm{log}^l|\xi |\frac{(zz_0)^k}{k!},$$
where $``$ denotes the usual mod($`\mathrm{S}^{\mathrm{}})`$ symbol product, $`๐ช:=๐ช(1)`$ and $`\mathrm{log}^k(๐ช)(x,\xi ):=(\mathrm{log}(๐ช)\mathrm{}\mathrm{log}(๐ช))(x,\xi )\text{CS}^{0,k}(U,V)`$ with $`k`$ factors.
### 1.4. A Laurent expansion for finite part integrals of holomorphic symbols
The following theorem computes the Laurent expansion for finite part integrals of holomorphic families of classical symbols of order $`\alpha (z)`$ in terms of local canonical and residue densities. This extends Proposition 3.4 in \[KV\], and results of \[Gu, Wo\], where the pole, the first coefficient in the expansion, was identified as the residue trace. The proof uses the property that each term of the Taylor series of a holomorphic family of classical symbols has an asymptotic symbol expansion, allowing the Laurent expansion of $`\sigma (z)(x,\xi )๐\overline{}\xi `$ to be computed through Lemma 1.8. Notice that although the Taylor expansion in the $`C^{\mathrm{}}`$ topology gives no control over the symbol as $`|\xi |\mathrm{}`$, (1.15), (1.16) impose what is needed to ensure integrability requirements.
###### Definition 1.19.
A holomorphic function $`\alpha :W`$ defined on a domain $`W`$ is said to be non-critical on
$$P:=\alpha ^1(\mathrm{Z}\mathrm{Z}[n,+\mathrm{}[)W$$
if $`\alpha ^{}(z_0)0`$ at each $`z_0P`$.
###### Theorem 1.20.
(1) Let $`U`$ be an open subset of $`\mathrm{I}\mathrm{R}^n`$. Let $`z\sigma (z)\text{CS}^{\alpha (z)}(U,V)`$ be a holomorphic family of classical symbols parametrized by a domain $`W\mathrm{I}\mathrm{C}`$ such that the order function $`\alpha `$ is non-critical on $`P`$. Then for each $`xU`$ the map $`z_{T_x^{}U}\sigma (z)(x,\xi )๐\overline{}\xi `$ is a meromorphic function on $`W`$ with poles located in $`P`$. The poles are at most simple and for $`z`$ near $`z_0P`$ one has
$`{\displaystyle _{T_x^{}U}}\sigma (z)(x,\xi )๐\overline{}\xi ={\displaystyle \frac{1}{\alpha ^{}(z_0)}}{\displaystyle _{S_x^{}U}}\sigma (z_0)_n(x,\xi )๐\overline{}_S\xi {\displaystyle \frac{1}{(zz_0)}}`$
$`+\left({\displaystyle _{T_x^{}U}}\sigma (z_0)(x,\xi )๐\overline{}\xi {\displaystyle \frac{1}{\alpha ^{}(z_0)}}{\displaystyle _{S_x^{}U}}\sigma ^{}(z_0)_{n,0}(x,\xi )๐\overline{}_S\xi \right)`$
$`+{\displaystyle \frac{\alpha ^{\prime \prime }(z_0)}{2\alpha ^{}(z_0)^2}}{\displaystyle _{S_x^{}U}}\sigma (z_0)_n(x,\xi )๐\overline{}_S\xi `$
$`+{\displaystyle \underset{k=1}{\overset{K}{}}}\left({\displaystyle _{T_x^{}U}}\sigma ^{(k)}(z_0)(x,\xi )๐\overline{}\xi {\displaystyle _{S_x^{}U}}_k(\sigma (z_0),\mathrm{},\sigma ^{(k+1)}(z_0))_{n,0}(x,\xi )๐\overline{}_S\xi \right){\displaystyle \frac{(zz_0)^k}{k!}}`$
$`+o\left((zz_0)^K\right),`$ (1.25)
where
$$_k(\sigma (z_0),\mathrm{},\sigma ^{(k+1)}(z_0))=\underset{j=0}{\overset{k+1}{}}\frac{p_{k+1j}}{\alpha ^{^{}}\left(z_0\right)^{k+2j}}\sigma ^{(j)}(z_0)\text{CS}^{\alpha (z_0),k+1}(U,V),$$
(1.26)
and $`p_{k+1j}`$ is an explicitly computable polynomial of degree $`k+1j`$ in $`\alpha ^{^{}}(z_0),\mathrm{},\alpha ^{(k+1)}(z_0)`$. Furthermore, the coefficient of $`\frac{(zz_0)^k}{k!}`$ in (1.20) is equal to $`\mathrm{fp}_{z=z_0}_{T_x^{}U}\sigma ^{(k)}(z).`$ If $`\alpha `$ is a linear function $`\alpha (z)=qz+b`$ with $`q0`$ then (1.20) reduces to
$`{\displaystyle _{T_x^{}U}}\sigma (z)(x,\xi )d\overline{}\xi ={\displaystyle \frac{1}{q}}{\displaystyle _{S_x^{}U}}\sigma (z_0)_n(x,\xi )d\overline{}_S\xi {\displaystyle \frac{1}{(zz_0)}}`$
$`+\left({\displaystyle _{T_x^{}U}}\sigma (z_0)(x,\xi )๐\overline{}\xi {\displaystyle \frac{1}{q}}{\displaystyle _{S_x^{}U}}\sigma ^{}(z_0)_{n,0}(x,\xi )๐\overline{}_S\xi \right)`$
$`+{\displaystyle \underset{k=1}{\overset{K}{}}}\left({\displaystyle _{T_x^{}U}}\sigma ^{(k)}(z_0)(x,\xi )๐\overline{}\xi {\displaystyle \frac{1}{q(k+1)}}{\displaystyle _{S_x^{}U}}\sigma ^{(k+1)}(z_0)_{n,0}(x,\xi )๐\overline{}_S\xi \right){\displaystyle \frac{(zz_0)^k}{k!}}`$
$`+o\left((zz_0)^K\right).`$ (1.27)
If $`z_0W`$ but $`z_0P`$, then $`_{T_x^{}U}\sigma (z)(x,\xi )๐\overline{}\xi `$ is holomorphic at $`z=z_0`$ and (1.20) then simplifies to the Taylor expansion
$$_{T_x^{}U}\sigma (z)(x,\xi )๐\overline{}\xi =_{T_x^{}U}\sigma (z_0)(x,\xi )๐\overline{}\xi +\underset{k=1}{\overset{K}{}}_{T_x^{}U}\sigma ^{(k)}(z_0)(x,\xi )๐\overline{}\xi \frac{(zz_0)^k}{k!}$$
$$+o\left((zz_0)^K\right).$$
(1.28)
(2) For any holomorphic family $`zA(z)\text{Cl}^{\alpha (z)}(M,E)`$ of classical $`\psi \mathrm{do}`$s parametrized by a domain $`W\mathrm{I}\mathrm{C}`$, such that order function $`\alpha `$ is non-critical on $`P`$, the map $`z\mathrm{TR}(A(z)):=_M๐x_{T_x^{}M}\mathrm{tr}_x\left(\sigma _{A(z)}(x,\xi )\right)๐\overline{}\xi `$ is a meromorphic function on $`W`$ with poles located in $`P`$. The poles are at most simple and for $`z`$ near $`z_0P`$
$`\mathrm{TR}(A(z))`$ $`=`$ $`{\displaystyle \frac{1}{\alpha ^{}(z_0)}}\mathrm{res}(A(z_0)){\displaystyle \frac{1}{(zz_0)}}`$ (1.29)
$`+{\displaystyle _M}๐x\left(\text{TR}_x(A(z_0)){\displaystyle \frac{1}{\alpha ^{}(z_0)}}\mathrm{res}_{x,0}(A^{}(z_0))\right)+{\displaystyle \frac{\alpha ^{\prime \prime }(z_0)}{2\alpha ^{}(z_0)^2}}\mathrm{res}\left(A(z_0)\right)`$
$`+{\displaystyle \underset{k=1}{\overset{K}{}}}{\displaystyle _M}๐x\left(\text{TR}_x(A^{(k)}(z_0))\text{res}_{x,0}\left(_k(\sigma _{A(z_0)},\mathrm{},\sigma _{A^{(k+1)}(z_0)})\right)\right){\displaystyle \frac{(zz_0)^k}{k!}}`$
$`+o\left((zz_0)^K\right).`$
Furthermore, the coefficient of $`\frac{(zz_0)^k}{k!}`$ in (1.29) is equal to $`\mathrm{fp}_{z=z_0}\mathrm{TR}(A^{(k)}(z)).`$ If $`A(z)`$ has order $`\alpha (z)=qz+b`$ with $`q0`$ then
$`\mathrm{TR}(A(z))`$ $`=`$ $`{\displaystyle \frac{1}{q}}\mathrm{res}(A(z_0)){\displaystyle \frac{1}{(zz_0)}}`$ (1.30)
$`+{\displaystyle _M}๐x\left(\text{TR}_x(A(z_0)){\displaystyle \frac{1}{q}}\mathrm{res}_{x,0}(A^{}(z_0))\right)`$
$`+{\displaystyle \underset{k=1}{\overset{K}{}}}{\displaystyle _M}๐x\left(\text{TR}_x(A^{(k)}(z_0)){\displaystyle \frac{\text{res}_{x,0}\left(\sigma _A^{(k+1)}(z_0)\right)}{q(k+1)}}\right){\displaystyle \frac{(zz_0)^k}{k!}}`$
$`+o\left((zz_0)^K\right).`$
If $`z_0W`$ but $`z_0P`$, then $`\mathrm{TR}(A(z))`$ is holomorphic at $`z=z_0`$ and (1.29) then simplifies to the Taylor expansion
$$\mathrm{TR}(A(z))=\mathrm{TR}(A(z_0))+\underset{k=1}{\overset{K}{}}\mathrm{TR}(A^{(k)}(z_0))\frac{(zz_0)^k}{k!}+o\left((zz_0)^K\right).$$
(1.31)
###### Remark 1.21.
Since $`\alpha `$ is non-critical on $`P`$, we have from Proposition 1.17 and equation (1.22) that the operators $`A^{(k)}(z_0)\text{Cl}^{\alpha (z_0),k}(M,E)`$ in equation (1.29) are not classical for $`k1`$.
###### Remark 1.22.
At a point $`z_0P`$, $`\alpha ^{}(z_0)0`$; writing $`\alpha (z)=\alpha (z_0)+\alpha ^{}(z_0)(zz_0)+o(zz_0)`$ we find that $`\alpha `$ is injective in a neighborhood of $`z_0`$. As a consequence, $`\mathrm{Z}\mathrm{Z}`$ being countable, so is the set of poles $`P=\alpha ^1(\mathrm{Z}\mathrm{Z}[n,+\mathrm{}))W`$ countable.
###### Remark 1.23.
Setting $`\alpha (z)=\frac{z}{1+\lambda z}`$ with $`\lambda \mathrm{I}\mathrm{R}^{}`$ for $`z\mathrm{I}\mathrm{C}\backslash \{\lambda ^1\}`$ gives rise to an additional finite part $`\frac{\alpha ^{\prime \prime }(0)}{2\alpha ^{}(0)^2}_{S_x^{}U}\sigma (0)_n(x,\xi )๐\overline{}_S\xi =\lambda _{S_x^{}U}\sigma (0)_n(x,\xi )๐\overline{}_S\xi `$ just as a rescaling $`Re^\lambda R`$ in the finite part integrals gives rise to the extra term $`\lambda _{S_x^{}U}\sigma (0)_n(x,\xi )๐\overline{}_S\xi `$ (see (1.10) with $`k=0`$ and $`\mu =e^\lambda `$).
###### Proof.
Since the orders $`\alpha (z)`$ define a holomorphic map at each point of $`P`$, for any $`z_0P`$ there is a ball $`B(z_0,r)W\mathrm{I}\mathrm{C}`$ centered at $`z_0W`$ with radius $`r>0`$ such that $`\left(B(z_0,r)\backslash \{z_0\}\right)P=\varphi `$. In particular, for all $`zB(z_0,r)\backslash \{z_0\}`$, the symbols $`\sigma (z)`$ have non-integer order. As a consequence, outside the set $`P`$, the finite part integral $`_{T_x^{}U}\sigma (z)(x,\xi )๐\overline{}\xi `$ is defined without ambiguity and $`_{T_x^{}U}\sigma (z)(x,\xi )๐\overline{}\xi ๐x`$ defines a global density on $`M`$.
Since $`z_0P`$, there is some $`j_0\mathrm{I}\mathrm{N}\{0\}`$ such that $`\alpha (z_0)+nj_0=0`$. On the other hand, for $`zB(z_0,r)\backslash \{z_0\}`$ we have $`\alpha (z)+nj0`$ and $`N>\mathrm{Re}(\alpha (z))+n1`$ can be chosen uniformly to ensure that $`\sigma _{(N)}(z)\mathrm{S}^{<n}(U,V)`$. Hence for $`zB(z_0,r)\backslash \{z_0\}`$ equation (1.9) yields (with $`k=0`$)
$$_{T_x^{}U}\sigma (z)(x,\xi )๐\overline{}\xi $$
$`=`$ $`{\displaystyle \underset{j=0}{\overset{N}{}}}{\displaystyle _{B_x^{}(0,1)}}\sigma (z)_{\alpha (z)j}(x,\xi )๐\overline{}\xi +{\displaystyle _{T_x^{}U}}\sigma _{(N)}(z)(x,\xi )๐\overline{}\xi `$ (1.32)
$`{\displaystyle \underset{j=0}{\overset{N}{}}}{\displaystyle \frac{1}{\alpha (z)+nj}}{\displaystyle _{S_x^{}U}}\sigma (z)_{\alpha (z)j}(x,\xi )๐\overline{}_S\xi `$
$`=`$ $`{\displaystyle \underset{j=0}{\overset{N}{}}}{\displaystyle _{B_x^{}(0,1)}}\sigma (z)_{\alpha (z)j}(x,\xi )๐\overline{}\xi +{\displaystyle _{T_x^{}U}}\sigma _{(N)}(z)(x,\xi )๐\overline{}\xi `$
$`{\displaystyle \underset{j=0,jj_0}{\overset{N}{}}}{\displaystyle \frac{1}{\alpha (z)+nj}}{\displaystyle _{S_x^{}U}}\sigma (z)_{\alpha (z)j}(x,\xi )๐\overline{}_S\xi `$
$`{\displaystyle \frac{1}{\alpha (z)\alpha (z_0)}}{\displaystyle _{S_x^{}U}}\sigma (z)_{\alpha (z)j_0}(x,\xi )๐\overline{}_S\xi `$
where, in view of the growth conditions (1.15) and (1.16), it is not hard to see that each of the integrals on the right-side of (1.32) is holomorphic in $`z`$. Since $`\sigma _{\alpha (z)j}(z)(x,\xi )`$ is a holomorphic family of classical symbols, there is a Taylor expansion (1.13)
$$\sigma (z)_{\alpha (z)j}(x,\xi )=\underset{k=0}{\overset{\mathrm{}}{}}\sigma ^{(k)}(z_0)_{\alpha (z_0)j}(x,\xi )\frac{(zz_0)^k}{k!}$$
(1.33)
with coefficients in $`\text{CS}^{\alpha (z_0)j,k}(U)`$
$$\sigma ^{(k)}(z_0)_{\alpha (z_0)j}(x,\xi ):=_z^k\left(\sigma (z)_{\alpha (z)j}\right)|_{z=z_0}=\left(_z^k\sigma (z)\right)_{\alpha (z)j}|_{z=z_0},$$
(1.34)
where the first equality is by definition while the second equality is equation (1.20), and likewise there is a Taylor expansion of the remainder $`\sigma _{(N)}(z)(x,\xi )`$ with coefficients
$$_z^k\left(\sigma _{(N)}(z)(x,\xi )\right)|_{z=z_0}=\left(_z^k\sigma (z)\right)_{(N)}(x,\xi )|_{z=z_0},$$
(1.35)
where again the equality is consequent on equations (1.18) and (1.20). For any non-negative integer $`K`$ we may therefore rewrite the first two lines of (1.32) as a polynomial $`_{k=0}^Ka_k\frac{(zz_0)^k}{k!}`$ plus an error term of order $`o((zz_0)^K)`$ with
$`a_k`$ $`=`$ $`{\displaystyle \underset{j=0}{\overset{N}{}}}{\displaystyle _{B_x^{}(0,1)}}\left(_z^k\sigma (z)\right)_{\alpha (z)j}(x,\xi )|_{z=z_0}d\overline{}\xi +{\displaystyle _{T_x^{}U}}\left(_z^k\sigma (z)\right)_{(N)}(x,\xi )|_{z=z_0}d\overline{}\xi `$ (1.36)
$`{\displaystyle \underset{j=0,jj_0}{\overset{N}{}}}_z^k|_{z=z_0}\left({\displaystyle \frac{1}{\alpha (z)+nj}}{\displaystyle _{S_x^{}U}}\sigma (z)_{\alpha (z)j}(x,\xi )๐\overline{}_S\xi \right).`$
Here, since $`jj_0`$, we use the fact that each factor in the terms of the final summation of (1.36) are holomorphic in a neighbourhood of $`z_0`$ (including at $`z=z_0`$). On the other hand, from (1.18)
$$\sigma ^{(k)}(z_0)(x,\xi ):=_z^k\sigma (z)(x,\xi )|_{z=z_0}\underset{j0}{}\left(_z^k\sigma (z)\right)_{\alpha (z)j}(x,\xi )|_{z=z_0}$$
while we know from (1.17) that $`\sigma ^{(k)}(z)\text{CS}^{\alpha (z),k}(U)`$. Hence (1.9) may be applied to see that
$$_{T_x^{}U}\sigma ^{(k)}(z_0)(x,\xi )๐\overline{}\xi $$
$`=`$ $`{\displaystyle \underset{j=0}{\overset{N}{}}}{\displaystyle _{B_x^{}(0,1)}}\left(_z^k\sigma (z)\right)_{\alpha (z)j}(x,\xi )|_{z=z_0}d\overline{}\xi +{\displaystyle _{T_x^{}U}}\left(_z^k\sigma (z)\right)_{(N)}(x,\xi )|_{z=z_0}d\overline{}\xi `$ (1.37)
$`+`$ $`{\displaystyle \underset{j=0,jj_0}{\overset{N}{}}}{\displaystyle \underset{l=0}{\overset{k}{}}}{\displaystyle \frac{(1)^{l+1}l!}{(\alpha (z_0)j+n)^{l+1}}}{\displaystyle _{S_x^{}U}}\left(_z^k\sigma (z)\right)_{\alpha (z)j,l}(x,\xi )|_{z=z_0}d\overline{}_S\xi `$
From the following lemma we conclude that the expressions in (1.36) and (1.37) are equal.
###### Lemma 1.24.
For $`jj_0`$ one has in a neighbourhood of $`z_0`$
$$_z^k\left(\frac{1}{\alpha (z)+nj}_{S_x^{}U}\sigma (z)_{\alpha (z)j}(x,\xi )๐\overline{}_S\xi \right)$$
$$=\underset{l=0}{\overset{k}{}}\frac{(1)^{l+1}l!}{(\alpha (z)j+n)^{l+1}}_{S_x^{}U}\left(_z^k\sigma (z)\right)_{\alpha (z)j,l}(x,\xi )๐\overline{}_S\xi .$$
(1.38)
###### Proof.
We choose $`z`$ in a neighbourhood of $`z_0`$ such that each of the factors on both sides of (1.38) are holomorphic. The equality holds trivially for $`k=0`$. For clarity we check the case $`k=1`$ before proceeding to the general inductive step. For $`k=1`$ the left-side of (1.38) is equal to
$$\frac{\alpha ^{}(z)}{(\alpha (z)j+n)^2}_{S_x^{}U}\sigma (z)_{\alpha (z)j}(x,\xi )๐\overline{}_S\xi $$
$$\frac{1}{\alpha (z)j+n}_{S_x^{}U}_z\left(\sigma (z)_{\alpha (z)j}\right)(x,\xi )d\overline{}_S\xi .$$
(1.39)
From (1.18) and (1.22), for $`|\xi |1`$
$$\left(_z\sigma (z)\right)_{\alpha (z)j}(x,\xi )=\alpha ^{}(z)\sigma (z)_{\alpha (z)j}(x,\xi )\mathrm{log}|\xi |+p_{\alpha (z)j}(z)(x,\xi )$$
and hence $`\left(_z\sigma (z)\right)_{\alpha (z)j,\mathrm{\hspace{0.17em}1}}(x,\xi )=\alpha ^{}(z)\sigma (z)_{\alpha (z)j}(x,\xi )`$ for $`|\xi |1`$. The expression in (1.39) is therefore equal to
$$\frac{1}{(\alpha (z)j+n)^2}_{S_x^{}U}\left(_z\sigma (z)\right)_{\alpha (z)j,\mathrm{\hspace{0.17em}1}}(x,\xi )๐\overline{}_S\xi $$
$$\frac{1}{\alpha (z)j+n}_{S_x^{}U}_z\left(\sigma (z)_{\alpha (z)j}\right)(x,\xi )d\overline{}_S\xi $$
which is the right-side of (1.38) for $`k=1`$.
Assume now that (1.38) holds for some arbitrary fixed $`k0`$. Then the left-side of (1.38) for $`k+1`$ is equal to
$$_z\left(_z^k\left(\frac{1}{\alpha (z)+nj}_{S_x^{}U}\sigma (z)_{\alpha (z)j}(x,\xi )๐\overline{}_S\xi \right)\right)$$
$`=`$ $`_z\left({\displaystyle \underset{l=0}{\overset{k}{}}}{\displaystyle \frac{(1)^{l+1}l!}{(\alpha (z)j+n)^{l+1}}}{\displaystyle _{S_x^{}U}}\left(_z^k\sigma (z)\right)_{\alpha (z)j,l}(x,\xi )๐\overline{}_S\xi \right)`$ (1.40)
$`=`$ $`{\displaystyle \underset{l=0}{\overset{k}{}}}{\displaystyle \frac{(1)^l(l+1)!\alpha ^{}(z)}{(\alpha (z)j+n)^{l+2}}}{\displaystyle _{S_x^{}U}}\sigma ^{(k)}(z)_{\alpha (z)j,l}(x,\xi )๐\overline{}_S\xi `$
$`+{\displaystyle \underset{l=0}{\overset{k}{}}}{\displaystyle \frac{(1)^{l+1}l!}{(\alpha (z)j+n)^{l+1}}}{\displaystyle _{S_x^{}U}}_z\left(\sigma ^{(k)}(z)_{\alpha (z)j,l}(x,\xi )\right)d\overline{}_S\xi ,`$
where for the second equality we use the property that both of the factors in each summand on the right-side of (1.38) are holomorphic near $`z_0`$, and in the notation of (1.23)
$$\left(_z^k\sigma (z)\right)_{\alpha (z)j}(x,\xi )=\underset{r=0}{\overset{k}{}}\sigma ^{(k)}(z)_{\alpha (z)j,r}(x,\xi )\mathrm{log}^r[\xi ].$$
In that notation the right-side of (1.38) for $`k`$ replaced by $`k+1`$ reads
$$\underset{l=0}{\overset{k+1}{}}\frac{(1)^{l+1}l!}{(\alpha (z)j+n)^{l+1}}_{S_x^{}U}\sigma ^{(k+1)}(z)_{\alpha (z)j,l}(x,\xi )๐\overline{}_S\xi ,$$
(1.41)
while on the (co-)sphere $`S_x^{}U`$ where $`|\xi |=1`$ the identities of Lemma 1.16 become
$`\sigma _{\alpha (z)j,k+1}^{(k+1)}(z)(x,\xi )`$ $`=`$ $`\alpha ^{}(z)\sigma _{\alpha (z)j,k}^{(k)}(z)(x,\xi ),`$
$`\sigma _{\alpha (z)j,l}^{(k+1)}(z)(x,\xi )`$ $`=`$ $`\alpha ^{}(z)\sigma _{\alpha (z)j,l1}^{(k)}(z)(x,\xi )+_z(\sigma _{\alpha (z)j,l}^{(k)}(z)(x,\xi )),1lk,`$
$`\sigma _{\alpha (z)j,\mathrm{\hspace{0.17em}0}}^{(k+1)}(z)(x,\xi )`$ $`=`$ $`_z(\sigma _{\alpha (z)j,0}^{(k)}(z)(x,\xi )).`$
Substitution of these identities in (1.41) immediately shows (1.41) to be equal to (1.40). This completes the proof of Lemma 1.24. โ
Returning to the proof of Theorem 1.20, from (1.36) and (1.37) and Lemma 1.24 we now have
$$a_k=_{T_x^{}U}\sigma ^{(k)}(z_0)(x,\xi )๐\overline{}\xi ,$$
and so the first two lines of (1.32) may be replaced by $`_{k=0}^K_{T_x^{}U}\sigma ^{(k)}(z_0)(x,\xi )๐\overline{}\xi \frac{(zz_0)^k}{k!}+o((zz_0)^K)`$. Hence (1.32) becomes
$`{\displaystyle _{T_x^{}U}}\sigma (z)(x,\xi )๐\overline{}\xi `$ $`=`$ $`{\displaystyle \underset{k=0}{\overset{K}{}}}{\displaystyle _{T_x^{}U}}\sigma ^{(k)}(z_0)(x,\xi )๐\overline{}\xi {\displaystyle \frac{(zz_0)^k}{k!}}+o((zz_0)^K)`$ (1.42)
$`{\displaystyle \frac{1}{\alpha (z)\alpha (z_0)}}{\displaystyle _{S_x^{}U}}\sigma (z)_n(x,\xi )๐\overline{}_S\xi .`$
To expand the sphere integral term in (1.42), since $`\alpha `$ is holomorphic we have in a neighbourhood of each $`z_0P`$ a Taylor expansion
$$\alpha (z)\alpha (z_0)=\underset{l=1}{\overset{L}{}}\frac{\alpha ^{(l)}(z_0)}{l!}(zz_0)^l+o(zz_0)^L$$
and hence since $`\alpha ^{}(z_0)0`$ an expansion
$$\frac{1}{\alpha (z)\alpha (z_0)}=\frac{1}{\alpha ^{}(z_0)(zz_0)}.\frac{1}{1+_{l=1}^L\frac{\alpha ^{(l+1)}(z_0)}{\alpha ^{}(z_0)}\frac{(zz_0)^l}{(l+1)!}+o(zz_0)^L}$$
$$=\frac{1}{\alpha ^{}(z_0)}.\frac{1}{(zz_0)}\frac{\alpha ^{\prime \prime }(z_0)}{2\alpha ^{}(z_0)^2}+\underset{j=1}{\overset{J}{}}\beta _j(z_0)(zz_0)^j+o(zz_0)^J$$
(1.43)
with $`\beta _j(z_0)`$ an explicitly computable rational function in $`\alpha ^{(k)}(z_0),\mathrm{\hspace{0.17em}1}kj+1`$ with denominator an integer power of $`\alpha ^{}(z_0)`$. On the other hand, since $`\alpha (z_0)j_0=n`$, the expansion (1.33) for $`j=j_0`$ becomes
$$\sigma (z)_{\alpha (z)j_0}(x,\xi )=\underset{k=0}{\overset{\mathrm{}}{}}\underset{l=0}{\overset{k}{}}(\sigma ^{(k)}(z_0))_{n,l}(x,\xi )\mathrm{log}^l[\xi ]\frac{(zz_0)^k}{k!}.$$
(1.44)
Since $`(\sigma ^{(k)}(z_0))_{n,l}(x,\xi )\mathrm{log}^l|\xi |=0`$ for $`l1`$ on $`S_x^{}U`$, we find from the expansions (1.43) and (1.44)
$$\frac{1}{\alpha (z)\alpha (z_0)}_{S_x^{}U}\sigma (z)_{\alpha (z)j_0}(x,\xi )๐\overline{}_S\xi =\frac{1}{\alpha ^{}(z_0)}._{S_x^{}U}\left(\sigma ^{}(z_0)\right)_{n,0}(x,\xi )๐\overline{}_S\xi \frac{1}{(zz_0)}$$
$$\underset{k=0}{\overset{K}{}}_{S_x^{}U}_k(\sigma (z_0),\sigma ^{}(z_0),\mathrm{},\sigma ^{(k+1)}(z_0))_{n,0}(x,\xi )๐\overline{}_S\xi \frac{(zz_0)^k}{k!}$$
$$+o\left((zz_0)^K\right),$$
(1.45)
where $`_k(\sigma (z_0),\sigma ^{}(z_0),\mathrm{},\sigma ^{(k+1)}(z_0))`$ is readily seen to have the form in (1.26). In particular, the explicit formulae given for the first two terms in (1.43) lead to the formula
$$_0(\sigma (z_0),\sigma ^{}(z_0))(x,\xi )=$$
$$\frac{1}{\alpha ^{}(z_0)}_{S_x^{}U}\sigma ^{}(z_0)_{n,0}(x,\xi )๐\overline{}_S\xi \frac{\alpha ^{\prime \prime }(z_0)}{2\alpha ^{}(z_0)^2}_{S_x^{}U}\sigma (z_0)_n(x,\xi )๐\overline{}_S\xi $$
which with the contribution from the $`k=0`$ finite-part integral on the right-side of (1.42) gives the stated constant term in the expansion (1.20). The next term up, for example, is
$$_1(\sigma (z_0),\sigma ^{}(z_0),\sigma ^{\prime \prime }(z_0))(x,\xi )=$$
$$\frac{1}{\alpha ^{}(z_0)}_{S_x^{}U}\sigma ^{\prime \prime }(z_0)_{n,0}(x,\xi )๐\overline{}_S\xi \frac{\alpha ^{\prime \prime }(z_0)}{2\alpha ^{}(z_0)^2}_{S_x^{}U}\sigma ^{}(z_0)_{n,0}(x,\xi )๐\overline{}_S\xi $$
$$+\frac{3\alpha ^{\prime \prime }(z_0)^22\alpha ^{\prime \prime \prime }(z_0)\alpha ^{}(z_0)}{12\alpha ^{}(z_0)^3}_{S_x^{}U}\sigma (z_0)_n(x,\xi )๐\overline{}_S\xi .$$
When $`\alpha (z)=qz+b`$ with $`q0`$ the right-side of (1.43) is $`\frac{1}{q(zz_0)}`$ and so from (1.33) one then has $`_k(\sigma (z_0),\sigma ^{}(z_0),\mathrm{},\sigma ^{(k+1)}(z_0))=\frac{\sigma ^{(k+1)}(z_0)}{q(k+1)!}`$ and so (1.27) follows.
If $`z_0P`$ then $`\alpha (z)\mathrm{I}\mathrm{C}\backslash \mathrm{Z}\mathrm{Z}`$ and so the log-polyhomogeneous symbols $`_k`$ in (1.45) then have non-integer order and hence have no component of degree $`n`$, and therefore vanish. Likewise the pole in(1.45) vanishes and so (1.20) simplifies, in this case, to (1.28). Alternatively, this can be seen in a simpler more direct way by using the linearity of the finite-part integral over log-polyhomogeneous symbols of non-integer order applied to the Taylor expansion of the symbol at $`z_0`$. (Indeed, in this case the term $`j=j_0`$ in (1.32) does not need to be treated separately from the sum in the previous line and (1.36) holds by linearity, from which Lemma 1.24 may then be inferred and now including the case $`j=j_0`$.)
This shows the first part of the theorem.
For the second part we use a partition of unity $`\{(U_i,\varphi _i)|iJ\}`$ such that for $`i,jJ`$ there is an $`l_{ij}J`$ with $`\text{supp}(\varphi _i)\text{supp}(\varphi _j)U_{ij}:=U_{l_{ij}}`$. We suppose trivialisations of $`\pi :EM`$ over each open set $`U_i`$. Then, with $`U_{ij}`$ identified with an open subset of $`\mathrm{I}\mathrm{R}^n`$, one has $`A(z)=_{i,j}\varphi _iA(z)\varphi _j`$ where $`\varphi _iA(z)\varphi _j=\mathrm{Op}(\sigma _{(ij)}(z))`$ is the localization of $`A`$ over $`U_{ij}`$ with amplitude
$$\sigma _{(ij)}(z)(x,y,\xi )\text{CS}^{\alpha (z)}(U_{ij}\times U_{ij},V)$$
a local holomorphic family of symbols in $`(x,y)`$ form. Each finite-part integral $`_{T_xU_{ij}}\sigma _{(ij)}(z)(x,x,\xi )๐\overline{}\xi `$ is well defined outside $`P`$, since $`A(z)`$ has non-integer order for those values of $`z`$. Using the linearity there of the canonical trace functional it follows that for $`zP`$
$$\mathrm{TR}(A(z))=\underset{i,j}{}_{U_{ij}}_{T_xU_{ij}}\mathrm{tr}\left(\sigma _{(ij)}(z)(x,x,\xi )\right)๐\overline{}\xi ๐x,$$
where $`\mathrm{tr}`$ is the trace on $`\mathrm{End}(V)`$, allowing (1.20) to be applied to each of the summands defined over the trivialising charts. Each locally defined coefficient in the Laurent expansion is seen by holomorphic continuation to define a global density on $`M`$ in the way explained in Proposition 1.25. The first part of the theorem therefore yields that $`\mathrm{TR}(A(z))`$ is meromorphic with simple poles in $`P`$ and since
$$\sigma _{A(z_0)}^{(k)}=\sigma _{A^{(k)}(z_0)}$$
(1.46)
the identity (1.29) now follows from the formula (1.20) applied to each localization.
The fact that the coefficients of $`\frac{(zz_0)^k}{k!}`$ in the Laurent expansions of the meromorphic maps $`z_{T_x^{}U}\sigma (z)`$ and $`z\mathrm{TR}(A(z))`$ correspond to the finite part at $`z=z_0`$ of their derivative at order $`k`$ follows from the general property for a meromorphic function $`f`$ on an open set $`W\mathrm{I}\mathrm{C}`$ with Laurent expansion around $`z_0`$ given by $`f(z)=_{j=1}^J\frac{b_j}{(zz_0)^j}+_{k=0}^Ka_k\frac{(zz_0)^k}{k!}+o((zz_0)^K)`$ that
$$\mathrm{fp}_{z=z_0}f^{(k)}(z)=a_k.$$
(1.47)
Combined with the equality $`_z^k\mathrm{TR}(A(z))=\mathrm{TR}(A^{(k)}(z))`$ valid for $`zP`$ we reach the conclusion.
Since the formulas (1.30), (1.31) now follow from (1.27) and (1.28), this ends the proof of the theorem. โ
In passing from the local formula (1.20) to the global formula (1.29) in the proof of Theorem 1.20 we have implicitly used the following fact, yielding the Laurent coefficients to be global densities on $`M`$ which can be integrated.
###### Proposition 1.25.
Let $`c_k(x)`$ denote the coefficient of $`\frac{(zz_0)^k}{k!}`$ in the Laurent expansion (1.20). Then $`c_k(x)dx`$ is defined independently of the choice of local coordinates on $`M`$.
###### Proof.
By formula (1.47), the coefficient $`c_k(x)`$ of $`\frac{(zz_0)^k}{k!}`$ in the Laurent expansion (1.20) with $`\sigma (z)(x,)=\sigma _{A(z)}(x,)`$ is identified with the finite part at $`z_0`$ of the $`k`$-th derivative of the map
$$zI_{A(z)}(x):=_{T_x^{}U}\sigma _{A(z)}(x,\xi )๐\overline{}\xi $$
i.e. $`c_k(x)=\mathrm{fp}_{z=z_0}I_{A^{(k)}(z)}(x).`$ For $`zP`$ the property (1.12) holds for the finite part integral $`I_{A(z)}(x)`$ as well as for the finite part integrals $`I_{A^{(k)}(z)}(x)`$ since the order of $`A^{(k)}`$ differs from that of $`A(z)`$ by an integer.
The map $`zI_{A^{(k)}(z)}(x)`$ has a Laurent expansion $`I_{A^{(k)}(z)}(x)=_{j=1}^{k+1}\frac{b_j(x)}{(zz_0)^j}+_{k=0}^Kc_k(x)\frac{(zz_0)^k}{k!}+o((zz_0)^k)`$ and $`(zz_0)^{k+1}I_{A^{(k)}(z)}(x)`$ can be extended to a holomorphic function in a small ball centered at $`z_0`$ with value $`b_{k+1}(x)`$ at $`z_0`$. Since property (1.12) holds for $`I_{A^{(k)}(z)}(x)`$ outside $`z_0`$ in this ball, it holds for the holomorphic extension on the whole ball and hence for $`b_{k+1}(x)`$. Using (1.12), we deduce that $`b_{k+1}(x)dx`$ is defined independently of the choice of local coordinates on $`M`$ and so is the difference $`\left(I_{A^{(k)}(z)}(x)\frac{b_{k+1}(x)}{(zz_0)^{k+1}}\right)dx`$ for any $`z`$ outside $`z_0`$ in a small ball centered at $`z_0`$. Iterating this argument, one shows recursively on the integer $`1Jk`$ that $`\left(I_{A^{(k)}(z)}(x)_{j=1}^{k+1J}\frac{b_j(x)}{(zz_0)^j}\right)dx`$ is defined independently of the choice of local coordinates on $`M`$ in a small ball centered at $`z_0`$. Consequently, the finite part $`\left(\mathrm{fp}_{z=z_0}I_{A^{(k)}(z)}(x)\right)dx`$ at $`z_0`$ is also defined independently of the choice of local coordinates. Since this finite part coincides with $`k!c_k(x)`$, we have that $`c_k(x)dx`$ is defined independently of the choice of local coordinates on $`M`$. โ
Examining the singular and constant terms in the expansions of Theorem 1.20 we have the following corollaries.
First, the singular term yields the known identification of the residue trace with complex residue of the canonical trace, derived in \[Gu\], \[Wo\], \[KV\]. With the assumptions of Theorem 1.20:
###### Corollary 1.26.
The map $`z_{T_x^{}U}\sigma (z)(x,\xi )๐\overline{}\xi `$ is meromorphic with at most a simple pole at $`z_0P`$ with complex residue
$$\mathrm{Res}_{z=z_0}_{T_x^{}U}\sigma (z)(x,\xi )๐\overline{}\xi =\frac{1}{\alpha ^{}(z_0)}_{S_x^{}U}\sigma (z_0)_n(x,\xi )๐\overline{}_S\xi .$$
(1.48)
For the holomorphic family $`zA(z)`$ of $`\psi \mathrm{do}`$s parametrized by $`W`$, the form $`\frac{1}{\alpha ^{}(z_0)}_{S_x^{}U}\left(\sigma _{A(z_0)}\right)_n(x,\xi )๐\overline{}_S\xi ๐x`$ defines a global density on the manifold $`M`$ and the map $`z\mathrm{TR}(A(z)):=_M๐x\text{TR}_x(A(z))`$ is a meromorphic function with at most a simple pole at $`z_0P`$ with complex residue
$$\mathrm{Res}_{z=z_0}\mathrm{TR}(A(z))=\frac{1}{\alpha ^{}(z_0)}\mathrm{res}\left(A(z_0)\right).$$
(1.49)
Thus, consequent to Proposition 1.25, one infers here the global existence of the residue density for integer order operators from the existence of the canonical trace density for non-integer order operators and holomorphicity.
On the other hand, the constant term provides a โdefect formulaโ for finite part integrals.
With the assumptions of Theorem 1.20:
###### Theorem 1.27.
For a holomorphic family of symbols $`z\sigma (z)\text{CS}(U,V)`$ parametrized by a domain $`W\mathrm{I}\mathrm{C}`$ and for any $`xU`$,
$`\mathrm{fp}_{z=z_0}{\displaystyle _{T_x^{}U}}\sigma (z)(x,\xi )๐\overline{}\xi `$ $`=`$ $`{\displaystyle _{T_x^{}U}}\sigma (z_0)(x,\xi )๐\overline{}\xi {\displaystyle \frac{1}{\alpha ^{}(z_0)}}{\displaystyle _{S_x^{}U}}\sigma ^{}(z_0)_{n,0}(x,\xi )๐\overline{}_S\xi .`$ (1.50)
$`+{\displaystyle \frac{\alpha ^{\prime \prime }(z_0)}{2\alpha ^{}(z_0)^2}}{\displaystyle _{S_x^{}U}}\sigma (z_0)_n(x,\xi )๐\overline{}_S\xi .`$
For the holomorphic family $`zA(z)\text{Cl}(M,E)`$ of $`\psi \mathrm{do}`$s parametrized by $`W\mathrm{I}\mathrm{C}`$,
$`\mathrm{fp}_{z=z_0}\mathrm{TR}(A(z))`$ $`=`$ $`{\displaystyle _M}๐x\left(\text{TR}_x(A(z_0)){\displaystyle \frac{1}{\alpha ^{}(z_0)}}\text{res}_{x,0}(A^{}(z_0))\right)`$ (1.51)
$`+{\displaystyle \frac{\alpha ^{\prime \prime }(z_0)}{2\alpha ^{}(z_0)^2}}\mathrm{res}\left(A(z_0)\right)`$
###### Remark 1.28.
Since $`\alpha `$ is non-critical on $`P`$, from Proposition 1.17 if $`z_0P`$ the operator $`A^{}(z_0)\text{Cl}^{\alpha (z_0),1}(M,E)`$ in equation (1.51) is not classical.
###### Remark 1.29.
If $`\mathrm{res}_{x,0}(A(z_0))=0`$ then $`\mathrm{fp}_{z=z_0}\text{TR}_x(A(z))=lim_{zz_0}\text{TR}_x(A(z))`$. If this holds for all $`xM`$, then $`\mathrm{TR}(A(z))`$ is holomorphic at $`z_0`$ and $`\mathrm{fp}_{z=z_0}\mathrm{TR}(A(z))=lim_{zz_0}\mathrm{TR}(A(z))`$.
One therefore has the following statement on the existence of densities associated to the local canonical trace.
###### Theorem 1.30.
With the assumptions of Theorem 1.20, for a holomorphic family $`zA(z)\text{Cl}(M,E)`$ parametrized by a domain $`W\mathrm{I}\mathrm{C}`$, and irrespective of the order $`\alpha (z_0)`$ of $`A(z_0)`$
$$\left(\text{TR}_x(A(z_0))\frac{1}{\alpha ^{}(z_0)}\text{res}_{x,0}(A^{}(z_0))\right)dx$$
(1.52)
defines a global density on $`M`$ which integrates on $`M`$ to $`\mathrm{fp}_{z=z_0}\mathrm{TR}(A(z))`$. If $`\alpha (z_0)\mathrm{Z}\mathrm{Z}`$ then (1.52) reduces to the canonical trace density on non-integer order classical $`\psi \mathrm{do}`$s of \[KV\].
Though this follows on the general grounds of Proposition 1.25, we have, for completeness, given a direct proof of Theorem 1.30 in Appendix A. This specializes to give the previously known existence of the canonical trace on non-integer order $`\psi \mathrm{do}`$s, recalled in Section (1.2).
With the assumptions of Theorem 1.20:
###### Theorem 1.31.
Let $`zA(z)\text{Cl}(M,E)`$ be a holomorphic family of classical $`\psi \mathrm{do}`$s parametrised by $`W\mathrm{I}\mathrm{C}`$ and let $`z_0W`$. If either
$$\text{TR}_x(A(z_0))dx=\left(_{T_x^{}M}\mathrm{tr}_x(\sigma _{A(z_0)})(x,\xi )๐\xi \right)dx$$
or
$$\mathrm{res}_x(A^{}(z_0))dx:=_{S_x^{}M}\mathrm{tr}_x\left(\left(\sigma _{A^{}(z_0)}\right)_n(x,\xi )\right)d_S\xi ๐x$$
defines a global density on $`M`$, then $`\mathrm{TR}\left(A(z_0)\right)`$ and $`\mathrm{res}\left(A^{}(z_0)\right)=_M\mathrm{res}_{x,0}(A^{}(z_0))๐x`$ are both well defined. The following defect formula then holds
$$\mathrm{fp}_{z=z_0}\mathrm{TR}(A(z))=\mathrm{TR}(A(z_0))\frac{1}{\alpha ^{}(z_0)}\mathrm{res}\left(A^{}(z_0)\right).$$
(1.53)
This holds in the following cases:
(i) If $`A(z_0)\text{Cl}^{\alpha (z_0),\mathrm{\hspace{0.17em}0}}(M,E)`$ satisfies one of the cases (1), (2) or (3) of Proposition 1.10 then $`\mathrm{TR}(A(z_0))`$ is defined and (1.53) holds. In case (1) this reduces to
$$\mathrm{fp}_{z=z_0}\mathrm{TR}(A(z))=\mathrm{TR}(A(z_0)).$$
(1.54)
(ii) If $`\mathrm{res}_{x,0}\left(A^{}(z_0)\right)=0`$ for all $`xM`$ then $`\mathrm{TR}(A(z))`$ is holomorphic at $`z_0W`$, so that $`\mathrm{fp}_{z=z_0}\mathrm{TR}(A(z))=lim_{zz_0}\mathrm{TR}(A(z))`$, and (1.54) holds.
(iii) If $`A(z_0)`$ is a differential operator, and more generally whenever $`\text{TR}_x(A(z_0))=0`$ for all $`xM`$, (1.53) reduces to
$$\mathrm{fp}_{z=z_0}\mathrm{TR}(A(z))=\frac{1}{\alpha ^{}(z_0)}\mathrm{res}\left(A^{}(z_0)\right).$$
(1.55)
###### Remark 1.32.
(1.53) can hold with both summands on the right-side of the equation non-zero. See Example 2.8.
###### Proof.
The first statement is consequent to Theorem 1.30. Since $`\text{TR}_x(A(z_0))dx`$ then defines a global density the transformation rule for finite part integrals in Proposition 1.9 implies that
$$_{S_x^{}M}\mathrm{tr}_x\left(\sigma _{A(z_0)}\right)_{n,\mathrm{\hspace{0.17em}0}}(x,\xi )\mathrm{log}|C^1\xi |d\overline{}_S\xi =0CGL_n(\mathrm{I}\mathrm{C})$$
and hence (taking $`C=\lambda I`$, $`\lambda \mathrm{I}\mathrm{C}`$) that
$$\mathrm{res}_{x,0}(A(z_0)):=_{S_x^{}M}๐\overline{}_S\xi \mathrm{tr}_x\left(\sigma _{A(z_0)}\right)_{n,\mathrm{\hspace{0.17em}0}}(x,\xi )=0.$$
Equation (1.53) now follows from (1.51). Parts (i), (ii), (iii) are now obvious in view of Proposition 1.10 and Proposition 1.12 and the vanishing of the residue trace on non-integer order operators and on differential operators. โ
## 2. Application to the Complex Powers
An operator $`Q\text{Ell}(M,E)`$ of positive order is called admissible if there is a proper subsector of $``$ with vertex 0 which contains the spectrum of the leading symbol $`\sigma _L(Q)`$ of $`Q`$. Then there is a half line $`L_\theta =\{re^{i\theta },r>0\}`$ (a spectral cut) with vertex $`0`$ and determined by an Agmon angle $`\theta `$ which does not intersect the spectrum of $`Q`$. Let $`\text{Ell}_{\mathrm{ord}>0}^{adm}(M,E)`$ denote the subset of admissible operators in $`\text{Ell}(M,E)`$ with positive order.
Let $`Q\text{Ell}_{\mathrm{ord}>0}^{adm}(M,E)`$ with spectral cut $`L_\theta `$. For Re $`z<0`$, the complex power $`Q_\theta ^z`$ of $`Q`$ is a bounded operator on any space $`H^s(M,E)`$ of sections of $`E`$ of Sobolev class $`H^s`$ defined by the contour integral:
$$Q_\theta ^z=\frac{i}{2\pi }_{C_\theta }\lambda ^z(Q\lambda I)^1๐\lambda $$
(2.1)
where $`C_\theta =C_{1,\theta ,r}C_{2,\theta ,r}C_{3,\theta ,r}`$. Here $`r`$ is a sufficiently small positive number and $`C_{1,\theta ,r}=\{\lambda =|\lambda |e^{i\theta }|+\mathrm{}>|\lambda |r\}`$, $`C_{2,\theta ,r}=\{\lambda =re^{i\varphi }|\theta \varphi \theta 2\pi \}`$ and $`C_{3,\theta ,r}=\{\lambda =|\lambda |e^{i(\theta 2\pi )}|r|\lambda |<+\mathrm{}\}`$. Here $`\lambda ^z=\mathrm{exp}(z\mathrm{log}\lambda )`$ where $`\mathrm{log}\lambda =\mathrm{log}|\lambda |+i\theta `$ on $`C_{1,\theta ,r}`$ and $`\mathrm{log}\lambda =\mathrm{log}|\lambda |+i(\theta 2\pi )`$ on $`C_{3,\theta ,r}`$.
For $`k\mathrm{I}\mathrm{N}`$ the complex power $`Q^z`$ is then extended to the half plane Re $`z<k`$ via the relation \[Se1\]
$$Q^kQ_\theta ^{zk}=Q_\theta ^z.$$
The definition of a complex power depends in general on the choice of $`\theta `$ and yields for any $`z\mathrm{I}\mathrm{C}`$ an elliptic operator $`Q_\theta ^z`$ of order $`z\mathrm{ord}(Q)`$. In spite of this $`\theta `$-dependence, we may occasionally omit it in order to simplify notations.
###### Remark 2.1.
For $`z=0`$
$$Q_\theta ^0=I\mathrm{\Pi }_Q$$
where $`\mathrm{\Pi }_Q`$ is the smoothing operator projection
$$\mathrm{\Pi }_Q=\frac{i}{2\pi }_{C_0}(Q\lambda I)^1๐\lambda $$
with $`C_0`$ a contour containing the origin but no other element of $`\mathrm{spec}(Q)`$, with range the generalized kernel $`\{\psi C^{\mathrm{}}(M,E)|Q^N\psi =0\mathrm{for}\mathrm{some}N\mathrm{I}\mathrm{N}\}`$ of $`Q`$. (See \[Bu\], \[Wo\], presented recently in \[Po\]).
Let $`Q\text{Ell}_{\mathrm{ord}>0}^{adm}(M,E)`$ be of order $`q`$ with spectral cut $`L_\theta `$. For arbitrary $`k\mathrm{Z}\mathrm{Z}`$, the map $`zQ_\theta ^z`$ defines a holomorphic function from $`\{z\mathrm{I}\mathrm{C},\mathrm{Re}z<k\}`$ to the space $`\left(H^s(M,E)H^{skq}(M,E)\right)`$ of bounded linear maps and we can set
$$\mathrm{log}_\theta Q:=\left[\frac{}{z}Q_\theta ^z\right]_{z=0}.$$
From (1.22), in a local trivialisation $`E_{|_U}U\times V`$ of $`E`$ over an open set $`U`$ of $`M`$ the symbol of $`\mathrm{log}_\theta Q`$ reads $`\sigma _{\mathrm{log}_\theta Q}(x,\xi )=\mathrm{ord}(Q)\mathrm{log}|\xi |\mathrm{Id}+\rho (x,\xi )`$ with $`\rho \text{Cl}^0(U,V)`$, and so $`\mathrm{log}_\theta Q\text{Cl}^{0,1}(M,E)`$ has order zero and log degree one. The logarithmic dependence is slight, for $`P,Q\text{Ell}_{\mathrm{ord}>0}^{adm}(M,E)`$, of non zero order $`p,q`$ respectively and admitting spectral cuts $`L_\theta `$ and $`L_\varphi `$ we have $`\frac{\mathrm{log}_\theta P}{p}\frac{\mathrm{log}_\varphi Q}{q}\text{Cl}^0(M,E).`$ More generally, higher derivatives of the complex powers have symbols with polynomial powers of $`\mathrm{log}|\xi |`$ and it follows from Proposition 1.17 that
$$\mathrm{log}_\theta ^kQ:=\left[\frac{^k}{z^k}Q_\theta ^z\right]_{z=0}\text{Cl}^{0,k}(M,E).$$
(2.2)
Theorem 1.20 leads to the following Laurent expansion.
###### Theorem 2.2.
Let $`Q\text{Ell}_{ord>0}^{adm}(M,E)`$ with spectral cut $`\theta `$ and of order $`q`$ and let $`A\text{Cl}^\alpha (M,E)`$. On the half plane $`\mathrm{Re}(z)>\frac{\alpha +n}{q}`$ the local Schwartz kernel $`K_{AQ_\theta ^z}(x,y)`$ of $`AQ_\theta ^z`$ is well defined and holomorphic and the restriction to the diagonal $`K_{AQ_\theta ^z}(x,x)dx=_{T_x^{}M}\sigma _{AQ_\theta ^z}(x,\xi )๐\overline{}\xi ๐x`$ defines a global density, an element of $`C^{\mathrm{}}(M,\mathrm{End}(E))`$. There is a meromorphic extension of $`K_{AQ_\theta ^z}(x,y)`$ to all $`z\mathrm{I}\mathrm{C}`$
$$K_{AQ_\theta ^z}(x,x)|^{\mathrm{mer}}:=_{T_x^{}M}\sigma _{AQ_\theta ^z}(x,\xi )d\overline{}\xi $$
(2.3)
with at most simple poles, each of which is located in $`P:=\{\frac{\alpha j}{q}|j[n,\mathrm{}[\mathrm{Z}\mathrm{Z}\}`$. For any $`xM`$, we have for $`z`$ near $`\frac{\alpha j}{q}P`$
$`K_{AQ_\theta ^z}(x,x)|^{\mathrm{mer}}={\displaystyle \frac{1}{q}}{\displaystyle _{S_x^{}M}}\left(\sigma _{AQ_\theta ^{(j\alpha )/q}}\right)_n(x,\xi )d\overline{}_S\xi {\displaystyle \frac{1}{(z\frac{\alpha j}{q})}}`$ (2.4)
$`+`$ $`{\displaystyle \underset{k=0}{\overset{K}{}}}{\displaystyle \frac{(1)^k}{k!}}\left(z{\displaystyle \frac{\alpha j}{q}}\right)^k`$
$`\times \left({\displaystyle _{T_x^{}U}}\sigma _{AQ_\theta ^{(j\alpha )/q}\mathrm{log}_\theta ^kQ}(x,\xi )๐\overline{}\xi {\displaystyle \frac{1}{q(k+1)}}{\displaystyle _{S_x^{}M}}\left(\sigma _{AQ_\theta ^{(j\alpha )/q}\mathrm{log}_\theta ^{k+1}Q}\right)_{n,0}(x,\xi )๐\overline{}_S\xi \right)`$
$`+`$ $`o\left(\left(z{\displaystyle \frac{\alpha j}{q}}\right)^K\right).`$
It follows that the map $`z\mathrm{TR}(AQ_\theta ^z):=_M\mathrm{tr}_x\left(K_{AQ_\theta ^z}(x)|^{\mathrm{mer}}\right)`$ is a meromorphic function with no more than simple poles located in $`P`$, and for $`z`$ near $`\frac{\alpha j}{q}P`$
$`\mathrm{TR}(AQ_\theta ^z)={\displaystyle \frac{1}{q}}\mathrm{res}(AQ_\theta ^{\frac{j\alpha }{q}}){\displaystyle \frac{1}{(z\frac{\alpha j}{q})}}`$
$`+{\displaystyle \underset{k=0}{\overset{K}{}}}{\displaystyle \frac{(1)^k}{k!}}\left(z{\displaystyle \frac{\alpha j}{q}}\right)^k`$
$`\times {\displaystyle _M}dx(\text{TR}_x(AQ_\theta ^{\frac{j\alpha }{q}}\mathrm{log}_\theta ^kQ){\displaystyle \frac{1}{q(k+1)}}\mathrm{res}_{x,0}(AQ_\theta ^{\frac{j\alpha }{q}}\mathrm{log}_\theta ^{k+1}Q))`$
$`+o\left(\left(z{\displaystyle \frac{\alpha j}{q}}\right)^K\right).`$ (2.5)
If $`z_0P`$ then $`\mathrm{TR}(AQ_\theta ^z)`$ is holomorphic at $`z_0`$ and for $`z`$ in a small enough neighbourhood of $`z_0`$
$$\mathrm{TR}(AQ_\theta ^z)=\underset{k=0}{\overset{K}{}}\frac{(1)^k}{k!}\mathrm{TR}(AQ_\theta ^{z_0}\mathrm{log}_\theta ^kQ)\frac{(zz_0)^k}{k!}+o\left(\left(zz_0\right)^K\right).$$
(2.6)
###### Proof.
Since $`\sigma (z):=\sigma _{AQ_\theta ^z}`$ has order $`\alpha (z)=\alpha qz`$, (1.27) of Theorem 1.20 can be applied to equation (2.3). Using (2.2) and Example 1.18, this yields (2.4). Applying the fibrewise trace $`\mathrm{tr}_x`$ and integrating over $`M`$ yields equation (2.5). Equation (2.6), to which (2.5) reduces when $`\alpha \mathrm{Z}\mathrm{Z}`$, as the operators inside the local residue traces then have non-integer order, follows from (1.31); that $`\text{TR}_x(AQ_\theta ^{z_0}\mathrm{log}_\theta ^kQ)dx`$ defines a global density on $`M`$ in this case is known from \[Le\]. โ
Because of the identity with the generalized zeta-function
$$\zeta _\theta (A,Q,z)=\mathrm{TR}\left(AQ_\theta ^z\right)$$
the expansion (2.5) is of particular interest near $`z=0`$, owing to the role of the Laurent coefficients there in geometric analysis.
###### Theorem 2.3.
If $`\mathrm{ord}(A)=\alpha [n,\mathrm{}[`$ then $`0P`$ and one then has near $`z=0`$
$`\zeta _\theta (A,Q,z)={\displaystyle \frac{1}{q}}\mathrm{res}(A){\displaystyle \frac{1}{z}}`$
$`+{\displaystyle _M}๐x\left(\text{TR}_x(A){\displaystyle \frac{1}{q}}\mathrm{res}_{x,0}(A\mathrm{log}_\theta Q)\right)\mathrm{tr}(A\mathrm{\Pi }_Q)`$
$`+{\displaystyle \underset{k=1}{\overset{K}{}}}(1)^k{\displaystyle \frac{z^k}{k!}}`$
$`\times {\displaystyle _M}dx(\text{TR}_x(A\mathrm{log}_\theta ^kQ){\displaystyle \frac{1}{q(k+1)}}\mathrm{res}_{x,0}(A\mathrm{log}_\theta ^{k+1}Q))\mathrm{tr}(A\mathrm{log}_\theta ^kQ\mathrm{\Pi }_Q)`$
$`+o(z^K).`$ (2.7)
If $`\alpha [n,\mathrm{}[`$ then $`\zeta _\theta (A,Q,z)`$ is holomorphic at zero and one has for $`z`$ near zero
$$\zeta _\theta (A,Q,z)=\underset{k=0}{\overset{K}{}}(1)^k\left(\mathrm{TR}(A\mathrm{log}_\theta ^kQ)\mathrm{tr}(A\mathrm{log}_\theta ^kQ\mathrm{\Pi }_Q)\right)\frac{z^k}{k!}+o(z^K).$$
(2.8)
###### Remark 2.4.
The formula (2.8) can also be deduced from exact formulas for the case $`\alpha \mathrm{Z}\mathrm{Z}`$ in \[Gr1\] Sect(3). All formulas presuppose the existence shown in \[KV\], \[Le\] of the canonical trace for non-integer order $`\psi \mathrm{do}`$s with log-polyhomogeneous symbol.
###### Proof.
The assumption $`\alpha [n,\mathrm{}[`$ means that $`(\alpha j_0)/q=0`$ for some $`j_0[n,\mathrm{}[`$. Hence (2.7) is almost obvious from (2.5); the subtle point is to take care to replace $`Q_\theta ^{\frac{\alpha j}{q}}=Q_\theta ^0`$ by $`I\mathrm{\Pi }_Q`$, see Remark 2.1. Since the spectral projection $`\mathrm{\Pi }_Q`$ is a smoothing operator the term $`\text{TR}_x(A\mathrm{log}_\theta ^kQ\mathrm{\Pi }_Q)dx`$ is an ordinary integral valued density and globally defined, yielding the term $`\mathrm{tr}(A\mathrm{log}_\theta ^kQ\mathrm{\Pi }_Q)`$. The formula (2.8) for $`A`$ of non-integer order (to which (2.7) reduces in this case) is immediate from (2.6). โ
We denote the coefficient of $`(z\frac{\alpha j}{q})^k/k!`$ in the Laurent expansion of the generalized zeta function at $`\frac{\alpha j}{q}P`$ by $`\zeta _\theta ^{(k)}(A,Q,\frac{\alpha j}{q})`$. In the case $`k=0`$, we use the simpler convention of writing the constant term $`\zeta _\theta ^{(0)}(A,Q,\frac{\alpha j}{q}):=\mathrm{fp}_{z=\frac{\alpha j}{q}}\zeta _\theta (A,Q,z)`$ as $`\zeta _\theta (A,Q,\frac{\alpha j}{q})`$. When $`A=I`$ write $`\zeta _\theta (Q,\frac{\alpha j}{q}):=\zeta _\theta (I,Q,\frac{\alpha j}{q}).`$
###### Corollary 2.5.
For any operator $`A\text{Cl}(M,E)`$,
$$\zeta _\theta (A,Q,0)=_M๐x\left(\text{TR}_x(A)\frac{1}{q}\mathrm{res}_{x,0}(A\mathrm{log}_\theta Q)\right)\mathrm{tr}(A\mathrm{\Pi }_Q).$$
(2.9)
More generally, for any non-negative integer $`k`$
$`\zeta _\theta ^{(k)}(A,Q,0)`$ $`=`$ $`(1)^k{\displaystyle _M}๐x\left(\text{TR}_x(A\mathrm{log}_\theta ^kQ){\displaystyle \frac{1}{q(k+1)}}\mathrm{res}_{x,0}(A\mathrm{log}_\theta ^{k+1}Q)\right)`$ (2.10)
$`+(1)^{k+1}\mathrm{tr}\left(A\mathrm{log}_\theta ^kQ\mathrm{\Pi }_Q\right).`$
If $`A`$ has integer order $`\alpha [n,\mathrm{})`$ then
$$\zeta _\theta (A,Q,\frac{\alpha j}{q})=_M๐x\left(\text{TR}_x(AQ_\theta ^{\frac{\alpha j}{q}})\frac{1}{q}\mathrm{res}_{x,0}(AQ_\theta ^{\frac{\alpha j}{q}}\mathrm{log}_\theta Q)\right)$$
Applied to the complex powers, the general statement on the existence of densities associated to the canonical and residue traces of Theorem 1.30 now states that independently of the order of $`A\text{Cl}(M,E)`$,
$$\left(\text{TR}_x(A)\frac{1}{q}\mathrm{res}_{x,0}(A\mathrm{log}_\theta Q)\right)dx$$
always defines a global density on $`M`$.
If $`A`$ has non-integer order this reduces to the KV canonical trace density and (by (2.8)) the identity (2.10) loses its residue defect term and one then has the known formula (cf. \[Gr1\] Cor. (3.8))
$$\zeta _\theta ^{(k)}(A,Q,0)=(1)^k\mathrm{TR}(A\mathrm{log}_\theta ^kQ)(1)^k\mathrm{tr}(A\mathrm{log}_\theta ^kQ\mathrm{\Pi }_Q).$$
(2.11)
Applying Theorem 1.31 to the zeta function at $`z=0`$ yields the following refinement of (2.9).
###### Theorem 2.6.
Let $`Q\text{Ell}_{ord>0}^{adm}(M,E)`$ be a classical $`\psi \mathrm{do}`$ with spectral cut $`\theta `$ and of order $`q`$ and let $`A\text{Cl}^\alpha (M,E)`$ be a classical $`\psi \mathrm{do}`$ of order $`\alpha `$. If either $`\text{TR}_x(A)dx`$ or $`\mathrm{res}_x(A\mathrm{log}_\theta Q)dx`$ defines a global density on $`M`$, then $`\zeta _\theta (A,Q,z)`$ is holomorphic at $`z=0`$, $`\mathrm{TR}(A)`$ and $`\mathrm{res}(A\mathrm{log}_\theta Q)`$ both exist, and one has
$$\zeta _\theta (A,Q,0)=\mathrm{TR}(A)\frac{1}{q}\mathrm{res}(A\mathrm{log}_\theta Q)\mathrm{tr}(A\mathrm{\Pi }_Q).$$
(2.12)
###### Proof.
If $`\text{TR}_x(A)dx`$ defines a global density then $`\mathrm{res}(A)`$ vanishes, as accounted for in the proof of Theorem 1.31, and so $`\zeta _\theta (A,Q,z)`$ is holomorphic at $`z=0`$. The formula is obvious from (1.53). โ
Notice that the assumptions of Theorem 2.6 also force $`\mathrm{res}(A)=0`$.
The situation of Theorem 2.6 can be seen to hold for certain combinations of even-even and even-odd $`\psi \mathrm{do}`$s. First, it holds in the following circumstances.
###### Corollary 2.7.
(i) If $`A`$ satisfies one of the cases (1), (2) or (3) of Proposition 1.10 then $`\mathrm{TR}(A)`$ is defined and (2.12) holds. In case (1) this reduces to
$$\zeta _\theta (A,Q,0)=\mathrm{TR}(A)\mathrm{tr}(A\mathrm{\Pi }_Q).$$
(2.13)
If $`Q`$ is an even-even operator and has even order, then (2.13) also holds when $`A`$ satisfies case (2) (assumes $`M`$ is odd-dimensional) or (3) (assumes $`M`$ is even-dimensional) of Proposition 1.10. These facts are known from \[Gr1\].
(ii) If $`A`$ is a differential operator, and more generally whenever $`\text{TR}_x(A)=0`$ for all $`xM`$, (2.12) reduces to
$$\zeta _\theta (A,Q,0)=\frac{1}{q}\mathrm{res}(A\mathrm{log}_\theta Q)\mathrm{tr}(A\mathrm{\Pi }_Q).$$
(2.14)
###### Proof.
Part (ii) follows from Proposition 1.12. For part (i), it is clear that (2.13) holds when $`\mathrm{res}_x(A\mathrm{log}_\theta Q)=0`$ for each $`xM`$. This is evident for case (1) operators. If $`A`$ satisfies case (2) (resp. case (3)) of Proposition 1.10 and if $`Q`$ is even-even and of even order, then it is not hard to see that $`\sigma _{A\mathrm{log}_\theta Q}(x,\xi )`$ is also even-even (resp. even-odd) and hence $`(\sigma _{A\mathrm{log}_\theta Q})_{n,0}(x,\xi )`$ vanishes when integrated over the $`n1`$ sphere. โ
###### Example 2.8.
To see that (2.12) may hold with all three terms non-zero, take $`A=D+S`$ with $`D`$ a differential operator and $`S`$ a smoothing operator, and let $`Q\text{Ell}_{ord>0}^{adm}(M,E)`$. Then $`\mathrm{TR}(A)=\mathrm{tr}(S)`$ and $`\mathrm{res}(A\mathrm{log}_\theta Q)=\mathrm{res}(D\mathrm{log}_\theta Q)`$ both exist (note Corollary 2.7 (ii)) and are non-zero in general. For example, if $`Q=D\text{Ell}_{ord>0}^{adm}(M,E)`$ is invertible one has $`\mathrm{res}(D\mathrm{log}_\theta D)=\zeta _\theta (D,1)`$.
###### Remark 2.9.
In Corollary 2.7 (i), if $`Q`$ has odd-order then (2.12) may hold with all three terms non-zero due to dependence on the choice of the spectral cut. The distinct behaviour for odd-order $`Q`$ was kindly pointed out to the authors by Gerd Grubb.
###### Remark 2.10.
Using Theorem 1.31 similar facts to those in Corollary 2.7 can be seen to hold for the $`\zeta _\theta ^{(k)}(A,Q,\frac{\alpha j}{q})`$, see also \[Gr1\] Sect.3. The regularity of $`\zeta _\theta (A,Q,z)`$ at $`z=0`$ in (ii) is proved in \[GS\]. When $`A=I`$ the identity (2.14) was shown in \[Sc\]. On the other hand, when $`Q`$ is a differential operator and taking $`A=Q^m`$ in (2.14) gives
$$\zeta _\theta (Q,m)=\frac{1}{q}\mathrm{res}\left(Q^m\mathrm{log}_\theta Q\right)\mathrm{tr}(Q^m\mathrm{\Pi }_Q),$$
(2.15)
which was obtained in the case when $`Q`$ is positive and invertible by other methods in \[Lo\]. Note that for sufficiently large $`m`$ one has $`\mathrm{tr}(Q^m\mathrm{\Pi }_Q)=0`$.
Looking at the next term up in the Laurent expansion, around $`z=0`$ the zeta function $`\zeta _\theta (Q,z)=\mathrm{TR}(Q_\theta ^z)`$ is holomorphic and hence the $`\zeta `$-determinant
$$\mathrm{det}_{\zeta ,\theta }Q=\mathrm{exp}(\zeta _\theta ^{}(Q,0)),$$
is defined, where $`\zeta _\theta ^{}(Q,0)=_z\zeta _\theta (Q,z))_{|_{z=0}}`$.
###### Theorem 2.11.
One has
$$\mathrm{log}\mathrm{det}_{\zeta ,\theta }(Q)=_M๐x\left(\text{TR}_x\left(\mathrm{log}_\theta Q\right)\frac{1}{2q}\text{res}_{x,0}\left(\mathrm{log}_\theta ^2Q\right)\right)\mathrm{tr}(\mathrm{log}_\theta Q\mathrm{\Pi }_Q).$$
(2.16)
If $`M`$ is odd-dimensional and $`Q`$ is an even-even operator and has even order then one has (as known from \[O2\],\[Gr1\] Sect.3, see also \[KV\] Sect. 4)
$$\mathrm{log}\mathrm{det}_{\zeta ,\theta }(Q)=\mathrm{TR}\left(\mathrm{log}_\theta Q\right)\mathrm{tr}(\mathrm{log}_\theta Q\mathrm{\Pi }_Q),$$
(2.17)
where $`\mathrm{TR}\left(\mathrm{log}_\theta Q\right)=_M\text{TR}_x\left(\mathrm{log}_\theta Q\right)๐x,`$
###### Proof.
Examining the coefficient of $`z`$ in the Laurent expansion (2.7) immediately yields (2.16). If $`Q`$ is even-even and of even order then the classical component of the local symbol of $`\mathrm{log}_\theta ^2Q\text{Cl}^{0,2}(M,E)`$ also has even-even parity. Hence the local residue integral of the term of homogeneity $`n`$ then vanishes, $`\text{TR}_x\left(\mathrm{log}_\theta Q\right)dx`$ defines a global density on $`M`$, and (2.16) reduces to (2.17). โ
### 2.1. The canonical trace on commutators and the residue trace on logarithms
The canonical trace $`\mathrm{TR}`$ is not defined on a commutator of classical $`\psi \mathrm{do}`$s which has integer order. Rather the following property holds.
###### Theorem 2.12.
Let $`Q\text{Ell}_{ord>0}^{adm}(M,E)`$ be of order $`q`$ and with spectral cut $`\theta `$, and let $`A\text{Cl}^\alpha (M,E)`$, $`B\text{Cl}^\beta (M,E)`$ for any $`\alpha ,\beta `$. Then
$$\left(\text{TR}_x\left([A,B]\right)\frac{1}{q}\mathrm{res}_{x,0}\left([A,B\mathrm{log}_\theta Q]\right)\right)dx$$
defines a global density on $`M`$ and one has
$$_M๐x\left(\text{TR}_x\left([A,B]\right)\frac{1}{q}\mathrm{res}_{x,0}\left([A,B\mathrm{log}_\theta Q]\right)\right)=0$$
(2.18)
independently of the choice of $`Q`$.
###### Proof.
Using the vanishing of $`\mathrm{TR}`$ in Proposition 1.13 (1), for $`z0`$ sufficiently close to $`0`$ we have
$$\mathrm{TR}\left([A,BQ_\theta ^z]\right)=0.$$
(2.19)
Hence the function $`z\mathrm{TR}\left([A,BQ_\theta ^z]\right)`$ also vanishes identically for such non-zero $`z`$. But from (2.7), $`z\mathrm{TR}\left([A,BQ_\theta ^z]\right)`$ extends holomorphically to include $`z=0`$. By equation (2.19) this analytically continued function must also vanish at $`z=0`$. It follows that $`\mathrm{TR}\left([A,BQ_\theta ^z]\right)`$ is holomorphic near $`z=0`$ and so (2.7) implies $`\mathrm{fp}_{z=0}\mathrm{TR}\left([A,BQ_\theta ^z]\right)=lim_{z0}\mathrm{TR}\left([A,BQ_\theta ^z]\right)=0.`$ Applying Proposition 1.31 to $`A(z)=[A,BQ_\theta ^z]`$ with $`z_0=0`$ we have by Theorem 1.27
$`0`$ $`=`$ $`\mathrm{fp}_{z=0}\mathrm{TR}\left([A,BQ_\theta ^z]\right)`$
$`=`$ $`{\displaystyle _M}๐x\left(\text{TR}_x\left([A,B(I\mathrm{\Pi }_Q)]\right)+{\displaystyle \frac{1}{q}}\text{res}_{x,0}\left([A,B\mathrm{log}_\theta Q]\right)\right)`$
which is equation (2.18), since $`\mathrm{TR}([A,B\mathrm{\Pi }_Q)])=\mathrm{tr}([A,B\mathrm{\Pi }_Q)])=0.`$
###### Corollary 2.13.
Let $`Q\text{Ell}_{ord>0}^{adm}(M,E)`$ be of order $`q`$ and with spectral cut $`\theta `$, and let $`A\text{Cl}^\alpha (M,E)`$, $`B\text{Cl}^\beta (M,E)`$. Then in cases (1), (2) and (3) of Proposition 1.13 the form $`\mathrm{res}_x\left([A,B\mathrm{log}_\theta Q]\right)dx`$ determines a global density on $`M`$ and one has
$$\mathrm{res}\left([A,B\mathrm{log}_\theta Q]\right)=0$$
independently of the choice of $`Q`$.
###### Remark 2.14.
The independence from $`Q`$ can also be seen for the residue trace term directly; given $`Q_1,Q_2\text{Ell}_{ord>0}^{adm}(M,E)`$ of order $`q_1`$ and $`q_2`$ respectively with common spectral cut $`\theta `$, the difference
$$\left(\frac{1}{q_1}\mathrm{res}_{x,0}\left([A,B\mathrm{log}_\theta Q_1]\right)\frac{1}{q_2}\mathrm{res}_{x,0}\left([A,B\mathrm{log}_\theta Q_2]\right)\right)dx$$
defines a global density which integrates to
$$\mathrm{res}\left([A,B\left(\frac{\mathrm{log}_\theta Q_1}{q_1}\frac{\mathrm{log}_\theta Q_2}{q_2}\right)]\right)=0$$
since $`\frac{\mathrm{log}_\theta Q_1}{q_1}\frac{\mathrm{log}_\theta Q_2}{q_2}`$ is a classical $`\psi \mathrm{do}`$.
A useful consequence of Theorem 2.12 and Proposition 1.12 is:
###### Corollary 2.15.
Let $`Q\text{Ell}_{ord>0}^{adm}(M,E)`$ of order $`q`$ and with spectral cut $`\theta `$ and let $`A,B\text{Cl}(M,E)`$. Whenever $`\mathrm{TR}\left([A,B]\right)=_M๐x\mathrm{TR}_x\left([A,B]\right)`$ is well defined then $`\mathrm{res}_{x,0}\left([A,B\mathrm{log}_\theta Q]\right)dx`$ is globally defined and one then has
$$\mathrm{res}\left([A,B\mathrm{log}_\theta Q]\right)=q\mathrm{TR}\left([A,B]\right).$$
(2.20)
In particular, if $`[A,B]`$ is a differential operator then $`\mathrm{res}_{x,0}\left([A,B\mathrm{log}_\theta Q]\right)dx`$ is globally defined and one has
$$\mathrm{res}\left([A,B\mathrm{log}_\theta Q]\right)=0.$$
In that case, whenever $`\mathrm{res}_{x,0}\left(AB\mathrm{log}_\theta Q\right)dx`$ defines a global density, then so does $`\mathrm{res}_{x,0}\left(B\mathrm{log}_\theta QA\right)dx`$ and
$$\mathrm{res}\left(B\mathrm{log}_\theta QA\right)=\mathrm{res}\left(AB\mathrm{log}_\theta Q\right).$$
In particular, since $`\mathrm{res}\left(\mathrm{log}_\theta Q\right)`$ exists \[O1\], for any invertible $`A\text{Cl}(M,E)`$
$$\mathrm{res}\left(A^1\mathrm{log}_\theta QA\right)=\mathrm{res}\left(\mathrm{log}_\theta Q\right).$$
(2.21)
###### Remark 2.16.
This proposition partially generalizes the fact \[O1\] that $`\mathrm{res}_{x,0}\left([A,\mathrm{log}_\theta Q]\right)dx`$ for $`A`$ a classical $`\psi \mathrm{do}`$ defines a global density and $`\mathrm{res}\left([A,\mathrm{log}_\theta Q]\right)=0`$, which when $`A`$ is a differential operator follows from the corollary applied to $`B=I`$.
On the other hand, the well known (\[MN\], \[O1\], \[CDMP\], \[Gr2\]) trace defect formula
$$\zeta _\theta ([A,B],Q,0)=\frac{1}{q}\mathrm{res}\left(A[B,\mathrm{log}_\theta Q]\right).$$
(2.22)
follows easily by applying the same argument as in the proof of Theorem 2.12 to $`C(z)=A[B,Q^z]`$. From (2.18) and (2.22) we infer:
###### Corollary 2.17.
For classical $`\psi \mathrm{do}`$s $`A`$ and $`B`$
$$\frac{1}{q}\mathrm{res}\left(A[B,\mathrm{log}_\theta Q]\right)=_M๐x\left(\text{TR}_x([A,B])\frac{1}{q}\mathrm{res}_x\left([A,B]\mathrm{log}_\theta Q\right)\right).$$
In cases (1), (2) and (3) of Proposition 1.13 the form $`\mathrm{res}_x\left([A,B]\mathrm{log}_\theta Q\right)dx`$ determines a global density on $`M`$ and one has
$$\mathrm{res}\left(A[B,\mathrm{log}_\theta Q]\right)=\mathrm{res}\left([A,B]\mathrm{log}_\theta Q\right).$$
While from Proposition 2.12 we conclude:
###### Corollary 2.18.
The density $`\text{res}_x\left([A,B\mathrm{log}_\theta Q][A,B]\mathrm{log}_\theta Q\right)dx`$ is globally defined on $`M`$ for classical $`\psi \mathrm{do}`$s $`A`$ and $`B`$ and one has
$$\mathrm{res}\left(A[B,\mathrm{log}_\theta Q]\right)=\mathrm{res}\left([A,B]\mathrm{log}_\theta Q[A,B\mathrm{log}_\theta Q]\right)$$
###### Proof.
$`{\displaystyle \frac{1}{q}}\mathrm{res}\left(A[B,\mathrm{log}_\theta Q]\right)`$ $`=`$ $`{\displaystyle _M}๐x\left(\text{TR}_x([A,B]){\displaystyle \frac{1}{q}}\mathrm{res}_x([A,B]\mathrm{log}_\theta Q)\right)`$
$`=`$ $`{\displaystyle _M}dx(\text{TR}_x([A,B]){\displaystyle \frac{1}{q}}\mathrm{res}_x\left([A,B\mathrm{log}_\theta Q]\right)`$
$`+{\displaystyle \frac{1}{q}}\mathrm{res}_x([A,B]\mathrm{log}_\theta Q[A,B\mathrm{log}_\theta Q]))`$
$`=`$ $`{\displaystyle \frac{1}{q}}{\displaystyle _M}๐x\mathrm{res}_x\left([A,B]\mathrm{log}_\theta Q[A,B\mathrm{log}_\theta Q]\right)`$
$`=`$ $`{\displaystyle \frac{1}{q}}\mathrm{res}\left([A,B]\mathrm{log}_\theta Q[A,B\mathrm{log}_\theta Q]\right)`$
We point out that Corollary 2.19 and (2.22) imply the following local index formulae.
###### Corollary 2.19.
Let $`A`$ be an elliptic $`\psi \mathrm{do}`$ with parametrix $`B`$. Let $`Q\text{Ell}_{ord>0}^{adm}(M,E)`$ be of order $`q`$ and with spectral cut $`\theta `$. Then, independently of the choice of $`Q`$,
$$\mathrm{res}\left([A,B\mathrm{log}_\theta Q]\right)=\mathrm{res}\left(A[B,\mathrm{log}_\theta Q]\right)$$
(2.23)
and are equal to $`q\text{index }(A)`$.
###### Proof.
In this case $`\text{index }(A)=\mathrm{tr}([A,B])`$ and since $`[A,B]`$ is smoothing equal to $`\mathrm{TR}\left([A,B]\right)`$. The first equality thus follows from (2.20). Since $`AB=I+S`$ where $`S`$ is a smoothing operator, and since $`\text{res}_{x,0}(S\mathrm{log}_\theta Q)`$ is therefore equal to zero, the second equality also follows. โ
## Appendix A: Proof of the density formula
The purpose here is to give a direct elementary proof of Theorem 1.30, which for the family $`zA(z)\text{Cl}(M,E)`$ parametrized by a domain $`W\mathrm{I}\mathrm{C}`$ states that irrespective of the order $`\alpha (z_0)`$ of $`A(z_0)`$
$$\left(\text{TR}_x(A)\frac{1}{\alpha ^{}}\text{res}_{x,0}(A^{})\right)dx$$
(2.24)
defines a global density on $`M`$. Here, we have written $`A=A(z_0)`$, $`A^{}=A^{}(z_0):=d/dz|_{z=z_0}(A(z))`$, and $`\alpha ^{}=\alpha ^{}(z_0)`$.
From previous works \[KV\] it is known that $`\text{TR}_x(A(z_0))dx`$ defines a global density on $`M`$ when $`\alpha (z_0)`$ is not integer valued; this follows immediately from (2.24) and Proposition 1.17.
The method of proof uses a generalization of the method used in \[O1\] to show that the residue density is globally defined for any classical $`\psi \mathrm{do}`$, and the method in \[Le\] used to show that the canonical density is globally defined for classical $`\psi \mathrm{do}`$s of non-integer order. We will take $`A`$ to be scalar valued for notational brevity, but the proof works in the same way for endomorphism valued operators; indeed it works equally for the pre-tracial density $`(_{T_xM}\sigma _A(x,\xi )๐\overline{}\xi \frac{1}{\alpha ^{}}_{S_x^{}M}(\sigma _A^{})_{n,0}(x,\xi )๐\overline{}_S\xi )dx.`$
First, we have a lemma, generalizing Lemma C.1 in \[O1\].
###### Lemma 2.20.
Let $`f(\xi )`$ be a smooth function on $`^n`$ which is homogeneous of degree $`n`$ for $`|\xi |1`$ and let $`T`$ be an invertible linear map on $`^n`$. Then for $`s`$ and any non-negative integer $`k`$
$$_{|\eta |=1}f(T\eta )|T\eta |^s\mathrm{log}^k|T\eta |d\overline{}_S\eta =\frac{(1)^k}{|\mathrm{det}T|}_{|\xi |=1}f(\xi )|T^1\xi |^s\mathrm{log}^k|T^1\xi |d\overline{}_S\xi .$$
Specifically, one has
$$_{|\eta |=1}f(T\eta )\mathrm{log}|T\eta |d\overline{}_S\eta =\frac{1}{|\mathrm{det}T|}_{|\xi |=1}f(\xi )\mathrm{log}|T^1\xi |d\overline{}_S\xi .$$
(2.25)
$$_{|\eta |=1}f(T\eta )๐\overline{}_S\eta =\frac{1}{|\mathrm{det}T|}_{|\xi |=1}f(\xi )๐\overline{}_S\xi .$$
(2.26)
###### Proof.
It is enough to prove this for $`k=0`$, differentiation with respect to $`s`$ yields the general formula. We have, using the linearity of $`T`$,
$`{\displaystyle _{1|\eta |2}}f(T\eta )|T\eta |^s๐\eta `$ $`=`$ $`{\displaystyle _{|\eta |=1}}{\displaystyle _{1r2}}f(rT\eta )r^s|T\eta |^sr^{n1}๐r๐\overline{}_S\eta `$
$`=`$ $`\left({\displaystyle \frac{2^s1}{s}}\right){\displaystyle _{|\eta |=1}}f(T\eta )|T\eta |^s๐\overline{}_S\eta .`$
On the other hand, changing variable,
$`{\displaystyle _{1|\eta |2}}f(T\eta )|T\eta |^s๐\eta `$ $`=`$ $`{\displaystyle \frac{1}{|\mathrm{det}T|}}{\displaystyle _{1|T^1\eta |2}}f(\eta )|\eta |^s๐\eta `$
$`=`$ $`{\displaystyle \frac{1}{|\mathrm{det}T|}}{\displaystyle _{|\eta |=1}}{\displaystyle _{1/|T^1\eta |r2/|T^1\eta |}}f(r\eta )r^s|\eta |^sr^{n1}๐r๐\overline{}_S\eta `$
$`=`$ $`{\displaystyle \frac{1}{|\mathrm{det}T|}}\left({\displaystyle \frac{2^s1}{s}}\right){\displaystyle _{|\eta |=1}}f(\eta )|T^1\eta |^s๐\overline{}_S\eta .`$
Consider now a local chart on $`M`$ defined by a diffeomorphism $`x:\mathrm{\Omega }U`$ from an open subset $`\mathrm{\Omega }`$ of $`M`$ to an open subset $`U`$ of $`^n`$. For $`p\mathrm{\Omega }`$ we then have the local coordinate $`x(p)^n`$. Let $`\kappa :UV`$ be a diffeomorphism to a second open subset $`V`$ of $`^n`$. Then $`y(p)=\kappa (x(p))`$ is also a local coordinate for $`\mathrm{\Omega }`$.
Let $`a(x(p),\xi )=\stackrel{~}{a}(x(p),x(p),\xi )`$ where $`\stackrel{~}{a}(x(p),y(p),\xi )`$ denotes the local amplitude of $`A`$ in $`x`$-coordinates, and likewise let $`b(y(p),\xi )`$ denote the amplitude along the diagonal in $`y`$-coordinates. From \[Ho\] with $`T(p):=(D\kappa _{x(p)})^t`$ we have
$`\mathrm{TR}_{y(p)}(A)dy(p)`$ $`:=`$ $`{\displaystyle _^n}b(y(p),\xi )๐\overline{}\xi ๐y(p)`$
$`=`$ $`{\displaystyle _^n}a(x(p),T(p)\xi )๐\overline{}\xi ๐y(p).`$
According to the transformation rule in Proposition 1.9, for $`f\text{CS}(V)`$ and $`T`$ an invertible linear map on $`^n`$
$$_^nf(T\xi )๐\overline{}\xi =\frac{1}{|\mathrm{det}T|}\left(_^nf(\xi )๐\overline{}\xi _{|\xi |=1}f(\xi )_{(n)}\mathrm{log}|T^1\xi |d\overline{}\xi \right)$$
with $`f(\xi )_{(n)}`$ the homogeneous component of $`f`$ of degree $`n`$. Hence
$`\mathrm{TR}_{y(p)}(A)dy(p)`$
$`={\displaystyle \frac{1}{|\mathrm{det}T(p)|}}\left({\displaystyle _^n}a(x(p),\xi )๐\overline{}\xi ๐y(p){\displaystyle _{|\xi |=1}}a(x(p),\xi )_{(n)}\mathrm{log}|T(p)^1\xi |d\overline{}\xi dy(p)\right)`$
$`={\displaystyle _^n}a(x(p),\xi )๐\overline{}\xi ๐x(p){\displaystyle _{|\xi |=1}}a(x(p),\xi )_{(n)}\mathrm{log}|T(p)^1\xi |d\overline{}\xi dx(p)`$
$`=\mathrm{TR}_{x(p)}(A)dx(p){\displaystyle _{|\xi |=1}}a(x(p),\xi )_{(n)}\mathrm{log}|T(p)^1\xi |d\overline{}\xi dx(p).`$ (2.27)
We turn now to the other component of (2.24) given in $`y`$-coordinates by
$$\frac{1}{\alpha ^{}}_{|\xi |=1}b^{}(y(p),\xi )_{(n)}๐\overline{}_S\xi ๐y(p),$$
where $`b^{}(y(p),\xi )=d/dz|_{z=z_0}(\sigma _A(z)(y(p),\xi ))`$ is the symbol derivative in $`y`$-coordinates and where $`b^{}(y(p),\xi )_{(n)}`$ denotes its log-homogeneous (cf. (1.2)) component of degree $`n`$. From \[Ho\] we have the asymptotic formula
$$b^{}(y(p),\xi )\underset{|\mu |0}{}_\xi ^\mu a^{}(x(p),T(p)\xi )\mathrm{\Psi }_\mu (x,\xi )$$
(2.28)
with $`\mathrm{\Psi }_\mu (x,\xi )`$ polynomial in $`\xi `$ of degree of at most $`|\alpha |/2`$. To begin with, suppose that $`a(x(p),\xi )`$ is homogeneous in $`\xi `$ of degree $`n`$. Then from (1.22) for $`|\eta |1`$
$$a^{}(x(p),\eta )=\alpha ^{}a(x(p),\eta )\mathrm{log}|\eta |+p_n(x(p),\eta )$$
(2.29)
with $`p_n(x(p),\eta )`$ positively homogeneous in $`\eta `$ of degree $`n`$, and $`a^{}(x(p),\eta )=a^{}(x(p),\eta )_{(n)}`$. Thus, if $`a(x(p),\xi )`$ is homogeneous in $`\xi `$ of degree $`n`$, by (2.28) and (2.29)
$`{\displaystyle \frac{1}{\alpha ^{}}}{\displaystyle _{|\xi |=1}}b^{}(y(p),\xi )_{(n)}๐\overline{}_S\xi ๐y(p)`$ $`=`$ $`{\displaystyle \frac{1}{\alpha ^{}}}{\displaystyle _{|\xi |=1}}a^{}(x(p),T(p)\xi )๐\overline{}_S\xi ๐y(p)`$ (2.30)
$`=`$ $`{\displaystyle \frac{1}{\alpha ^{}}}{\displaystyle _{|\xi |=1}}\alpha ^{}a(x(p),T(p)\xi )\mathrm{log}|T(p)\xi |d\overline{}_S\xi dy(p)`$
$`{\displaystyle \frac{1}{\alpha ^{}}}{\displaystyle _{|\xi |=1}}p_n(x(p),T(p)\xi )๐\overline{}_S\xi ๐y(p)`$
$`=`$ $`{\displaystyle _{|\xi |=1}}a(x(p),T(p)\xi )\mathrm{log}|T(p)\xi |d\overline{}_S\xi dy(p)`$
$`{\displaystyle \frac{1}{\alpha ^{}}}{\displaystyle _{|\xi |=1}}p_n(x(p),T(p)\xi )๐\overline{}_S\xi ๐y(p).`$
Using equations (2.25) and (2.26) of Lemma 2.20, (2.30) becomes
$`{\displaystyle \frac{1}{\alpha ^{}}}\mathrm{res}_{y(p),0}(A^{})dy(p)`$ $`={\displaystyle \frac{1}{|\mathrm{det}T(p)|}}{\displaystyle _{|\xi |=1}}a(x(p),\xi )\mathrm{log}|T(p)^1\xi |d\overline{}_S\xi dy(p)`$ (2.31)
$`{\displaystyle \frac{1}{\alpha ^{}}}{\displaystyle \frac{1}{|\mathrm{det}T(p)|}}{\displaystyle _{|\xi |=1}}p_n(x(p),\xi )๐\overline{}_S\xi ๐y(p)`$
$`={\displaystyle _{|\xi |=1}}a(x(p),\xi )\mathrm{log}|T(p)^1\xi |d\overline{}_S\xi dx(p)`$
$`{\displaystyle \frac{1}{\alpha ^{}}}{\displaystyle _{|\xi |=1}}p_n(x(p),\xi )๐\overline{}_S\xi ๐x(p)`$
$`={\displaystyle _{|\xi |=1}}a(x(p),\xi )\mathrm{log}|T(p)^1\xi |d\overline{}_S\xi dx(p)`$
$`{\displaystyle \frac{1}{\alpha ^{}}}\mathrm{res}_{x(p),0}(A^{})dx(p),`$
where the final equality follows from (2.29). Adding (2.27) and (2.31) we have when $`a(x(p),\xi )`$ is homogeneous in $`\xi `$ of degree $`n`$
$$\left(\mathrm{TR}_{y(p)}(A)\frac{1}{\alpha ^{}(z_0)}\mathrm{res}_{y(p),0}(A^{})\right)dy(p)=\left(\mathrm{TR}_{x(p)}(A)\frac{1}{\alpha ^{}(z_0)}\mathrm{res}_{x(p),0}(A^{})\right)dx(p),$$
(2.32)
proving the invariance of (2.24) in this case.
Next suppose that $`a(x(p),\xi )`$ is homogeneous in $`\xi `$ of degree $`\alpha >n`$. Then from (2.28) and since we can commute the $`z`$ and $`\mu `$ derivatives
$$b^{}(y(p),\xi )_{(n)}=\underset{|\mu |\alpha +n}{}\frac{d}{dz}|_{z=z_0}_\xi ^\mu \left(a(z)(x(p),T(p)\xi )\right)\mathrm{\Psi }_{\mu ,n}(x,\xi ).$$
where $`\mathrm{\Psi }_{\mu ,n}(x,\xi )`$ is a polynomial in $`\xi `$ of degree $`|\mu |n\alpha `$. Hence
$`{\displaystyle _{|\xi |=1}}b^{}(y(p),\xi )_{(n)}๐\overline{}_S\xi ๐y(p)`$
$`={\displaystyle \underset{|\mu |\alpha +n}{}}{\displaystyle \frac{d}{dz}}|_{z=z_0}{\displaystyle _{|\xi |=1}}_\xi ^\mu \left(a(z)(x(p),T(p)\xi )\right)\mathrm{\Psi }_{\mu ,n}(x,\xi )d\overline{}_S\xi dy(p)`$
$`=0.`$
The final equality follows using the integration by parts property in Lemma C1 of \[O1\], which states that if $`g(\xi )`$ and $`h(\xi )`$ are homogeneous in $`\xi `$ of degrees $`\gamma ,\delta `$ where $`\gamma +\delta =1n`$, then
$$_{|\xi |=1}(_{\xi _j}g(\xi ))h(\xi )๐\overline{}_S\xi =_{|\xi |=1}g(\xi )_{\xi _j}h(\xi )d\overline{}_S\xi ,$$
along with the fact that $`\mathrm{\Psi }_{\mu ,n}(x,\xi )`$ polynomial in $`\xi `$ of degree $`|\mu |n\alpha `$.
This completes the proof that (2.24) is a density independent of coordinates.
## Appendix B: Proof of Lemma 1.6 and Lemma 1.8
For a fixed $`N\mathrm{I}\mathrm{N}`$ chosen large enough such that $`\mathrm{Re}(\alpha )N1<n`$, we write $`\sigma (x,\xi )=_{j=0}^{K_N}\sigma _{\alpha j}(x,\xi )+\sigma _{(N)}(x,\xi )`$ and split the integral accordingly as
$$_{B_x^{}(0,R)}\sigma (x,\xi )๐\overline{}\xi =\underset{j=0}{\overset{N}{}}_{B_x^{}(0,R)}\sigma _{\alpha j}(x,\xi )๐\overline{}\xi +_{B_x^{}(0,R)}\sigma _{(N)}(x,\xi )๐\overline{}\xi .$$
Since $`\mathrm{Re}(\alpha )N1<n`$, $`\sigma _{(N)}`$ lies in $`L^1(T_x^{}U)`$ and the integral $`_{B_x^{}(0,R)}\sigma _{(N)}(x,\xi )๐\overline{}\xi `$ converges when $`R\mathrm{}`$ to $`_{T_x^{}U}\sigma _{(N)}(x,\xi )๐\overline{}\xi `$. On the other hand, for any $`jN`$
$`{\displaystyle _{B_x^{}(0,R)}}\sigma _{\alpha j}={\displaystyle _{B_x^{}(0,1)}}\sigma _{\alpha j}+{\displaystyle _{D_x^{}(1,R)}}\sigma _{\alpha j}.`$ (2.33)
Here $`D_x^{}(1,R)=B_x^{}(0,R)\backslash B_x^{}(0,1)`$. The first integral on the r.h.s. converges and since $`\sigma _{\alpha j}(x,\xi )_{l=0}^k\sigma _{\alpha j,l}(x,\xi )\mathrm{log}^l[\xi ],`$ the second integral reads:
$$_{D_x^{}(1,R)}\sigma _{\alpha j}(x,\xi )๐\overline{}\xi =\underset{l=0}{\overset{k}{}}_1^Rr^{\alpha j+n1}\mathrm{log}^lrdr_{S_x^{}U}\sigma _{\alpha j,l}(x,\omega )๐\omega .$$
Hence the following asymptotic behaviour:
$`{\displaystyle _{D_x^{}(1,R)}}๐\overline{}\xi \sigma _{\alpha j}(x,\xi )_R\mathrm{}`$
$`{\displaystyle \underset{l=0}{\overset{k}{}}}{\displaystyle \frac{\mathrm{log}^{l+1}R}{l+1}}{\displaystyle _{S_x^{}U}}\sigma _{\alpha j,l}(x,\omega )๐\overline{}_S\xi ={\displaystyle \underset{l=0}{\overset{k}{}}}{\displaystyle \frac{\mathrm{log}^{l+1}R}{l+1}}{\displaystyle _{S_x^{}U}}\sigma _{n,l}(x,\xi )๐\overline{}_S\xi \mathrm{if}\alpha j=n`$
$`{\displaystyle _{D_x^{}(1,R)}}๐\overline{}\xi \sigma _{\alpha j}(x,\xi )`$ $`_R\mathrm{}`$ $`{\displaystyle \underset{l=0}{\overset{k}{}}}({\displaystyle \underset{i=0}{\overset{l}{}}}{\displaystyle \frac{(1)^{i+1}\frac{l!}{(li)!}\mathrm{log}^iR}{(\alpha j+n)^i}}R^{\alpha j+n}{\displaystyle _{S_x^{}U}}\sigma _{\alpha j,l}(x,\xi )d\overline{}_S\xi `$
$`+`$ $`(1)^ll!{\displaystyle \frac{R^{\alpha j+n}}{(\alpha j+n)^{l+1}}}{\displaystyle _{S_x^{}U}}\sigma _{\alpha j,l}(x,\xi )๐\overline{}_S\xi `$
$`+`$ $`{\displaystyle \frac{(1)^{l+1}l!}{(\alpha j+n)^{l+1}}}{\displaystyle _{S_x^{}U}}\sigma _{\alpha j,l}(x,\xi )d\overline{}_S\xi )\mathrm{if}\alpha jn.`$
Putting together these asymptotic expansions yield the statements with
$$C_x(\sigma )=_{T_x^{}U}\sigma _{(N)}+\underset{j=0}{\overset{N}{}}_{B_x^{}(0,1)}\sigma _{a_j}+\underset{j=0,a_j+n0}{\overset{N}{}}\underset{l=0}{\overset{L}{}}\frac{(1)^{l+1}l!}{(a_j+n)^{l+1}}_{S_x^{}U}\sigma _{a_j,l}.$$
The $`\mu `$-dependence follows from
$`\mathrm{log}^{l+1}(\mu R)`$ $`=`$ $`\mathrm{log}^{l+1}R\left(1+{\displaystyle \frac{\mathrm{log}\mu }{\mathrm{log}R}}\right)^{l+1}`$
$`_R\mathrm{}`$ $`\mathrm{log}^{l+1}R{\displaystyle \underset{k=0}{\overset{l+1}{}}}C_{l+1}^k\left({\displaystyle \frac{\mathrm{log}\mu }{\mathrm{log}R}}\right)^k.`$
The logarithmic terms $`_{l=0}^k\frac{1}{l+1}_{S_x^{}U}\sigma _{n,l}(x,\xi )๐\overline{}_S\xi \mathrm{log}^{l+1}(\mu R)`$ therefore contribute to the finite part by $`_{l=0}^k\frac{\mathrm{log}^{l+1}\mu }{l+1}_{S_x^{}U}\sigma _{n,l}(x,\xi )๐\overline{}_S\xi `$ as claimed in the lemma.
Laboratoire de Mathรฉmatiques, Complexe des Cรฉzeaux, Universitรฉ Blaise Pascal, 63 177 Aubiรจre Cedex F. E-mail: sylvie.paycha@math.univ-bpclermont.fr
Department of Mathematics, Kingโs College London. E-mail: sgs@mth.kcl.ac.uk
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# Effect of spin fluctuations on Tc from density-functional theory for superconductors.
## 1 Introduction
Since the discovery of superconductivity many theories have been born to explain this phenomenon and calculate observables. First papers about the role of spin fluctuations by Doniach and Engelsberg Doniach and Izuyama et al. Japan were published in sixties. Till today, fluctuations have been introduced to many-body and phenomenological models and a very popular semiempirical theory proposed by Eliashberg Eliashberg .
The goal of this work is to include the spin fluctuations into the density functional theory for superconductors which, in principle, anables to calculate all material properties, also in the superconducting state, from first principles. The framework of the SCDFT was set up by Oliveira, Gross and Kohn OGK in 1988. Recently, the SCDFT gap equation has been solved numerically for simple metals Martin-OGK ; Miguel and MgB<sub>2</sub> Massidda .
As for the critical temperatures, it is known for a long time, that spin fluctuations decrease considerably T<sub>c</sub> of some superconductors Berk ; Winter . In our previous work for niobium under pressure our-Nb , we solved the gap equation of the Eliashberg theory Eliashberg with and without spin fluctuations and the SCDFT gap equation only with the Coulomb and phonon interactions. We found that the effect of paramagnons decreased T<sub>c</sub> obtained from the Eliashberg theory by 3-4 K, however, an approximate treatment of the Coulomb interactions by a simple constant, $`\mu ^{}`$, led to a large disagreement of the theoretical results with the experimetal data Struzhkin . In contrast to the Eliashberg theory, the SCDFT scheme is parameter free, but the critical temperature calculated without spin fluctuations for Nb at ambient pressure our-Nb was about 3.7 K higher than the experimental T<sub>c</sub>.
In this work, we follow the derivations of the SCDFT gap equation given in a number of PhD theses<sup>1</sup><sup>1</sup>1available at URL: www.physik.fu-berlin.de/$``$ag-gross SK ; ML ; MM , and we include the spin fluctuations. The paramagnon spectral function is calculated within the random phase approximation (RPA) with the assumption of the homogeneous electron gas, similarly to the work by Berk and Schrieffer Berk done for the Eliashberg theory. We solve the obtained gap equation for a few simple metals and update our previous results for niobium under pressure.
In the following Sections, we introduce the SCDFT gap equation and the construction of the exchange-correlation functional, $`F_{xc}`$, by collecting the most important building blocks of the theory given by its authors OGK and first developers SK ; ML ; MM ; SK-dec ; KC . These Sections are: II. SCDFT gap equation, III. Exchange-correlation functional, and IV. Coulomb interaction and phonons in $`F_{xc}`$. Above Sections are written using the notation according to Parks Parks ; Kazumi and Vonsovsky Vonsovsky . This notation is in some points, such as Nambu Greenโs function and the selfenergy, different than the notation previously used for the SCDFT SK ; ML ; MM . We introduce the spin fluctuations in Sections: V. Paramagnons in $`F_{xc}`$ and VI. Gap equation with paramagnons and implementation details. We report obtained critical temperatures in Section VII, and we summarize in Section VIII.
## 2 SCDFT gap equation
In this Section, we wish to guide the reader, step by step, to the gap equation which will be solved at the end of this work to calculate the critical temperatures. We start by bringing the fundaments of the SCDFT OGK and the main approximations, such as the decoupling of band energies and the superconducting gap and a linearization of the gap equation close to $`T_c`$, which were assumed for numerical convenience SK ; ML ; MM ; SK-dec . We believe that these approximations do not cause any significant difference in the calculated critical temperatures.
Turning to details of the SCDFT, in order to obtain the gap equation one needs to follow the points below:
1. The grand-canonical Hamiltonian for a superconductor reads
$`\widehat{H}_{v,\mathrm{\Delta }}={\displaystyle \underset{\sigma }{}}{\displaystyle d^3r\widehat{\psi }_\sigma ^{}(๐ซ)\left[\frac{^2}{2}+v(๐ซ)\mu \right]\widehat{\psi }_\sigma (๐ซ)}`$
$`+{\displaystyle \frac{1}{2}}{\displaystyle d^3rd^3r^{}\widehat{\psi }_\sigma ^{}(๐ซ)\widehat{\psi }_\sigma ^{}(๐ซ^{})\frac{1}{|๐ซ๐ซ^{}|}\widehat{\psi }_\sigma (๐ซ^{})\widehat{\psi }_\sigma (๐ซ)}`$
$`{\displaystyle d^3r_1d^3r_1^{}d^3r_2d^3r_2^{}\widehat{\psi }_{}^{}(๐ซ_1^{})\widehat{\psi }_{}^{}(๐ซ_1)}`$
$`\times w(๐ซ_1^{},๐ซ_1,๐ซ_2,๐ซ_2^{})\widehat{\psi }_{}(๐ซ_2)\widehat{\psi }_{}(๐ซ_2^{})`$
$`[{\displaystyle }d^3rd^3r^{}\mathrm{\Delta }^{}(๐ซ,๐ซ^{})\widehat{\psi }_{}(๐ซ)\widehat{\psi }_{}(๐ซ^{})+H.c.],`$ (1)
where $`v(๐ซ)`$ and $`\mathrm{\Delta }(๐ซ,๐ซ^{})`$ are an external potential and an anomalous pair potential respectively. The pairing interaction $`w`$ in the particular BCS case satisfies $`w(๐ซ_1^{},๐ซ_1,๐ซ_2,๐ซ_2^{})=w(๐ซ_1^{}๐ซ_1,๐ซ_2๐ซ_2^{})`$. The normal and anomalous densities, $`n(๐ซ)`$ and $`\chi (๐ซ,๐ซ^{})`$, are defined as
$`n(๐ซ)`$ $`=`$ $`{\displaystyle \underset{\sigma }{}}\widehat{\psi }_\sigma ^{}(๐ซ)\widehat{\psi }_\sigma (๐ซ),`$ (2)
$`\chi (๐ซ,๐ซ^{})`$ $`=`$ $`\widehat{\psi }_{}(๐ซ)\widehat{\psi }_{}(๐ซ^{}).`$ (3)
2. The Hohenberg-Kohn theorem for superconductors says that, at each temperature $`\theta =1/\beta `$, the normal and anomalous densities, $`n(๐ซ)`$ and $`\chi (๐ซ,๐ซ^{})`$, determine uniquely the density operator $`\widehat{\rho }=e^{\beta \widehat{H}_{v,\mathrm{\Delta }}}/Tre^{\beta \widehat{H}_{v,\mathrm{\Delta }}}`$ which minimizes the thermodynamic potential, $`\mathrm{\Omega }_{v,\mathrm{\Delta }}[\widehat{\rho }]`$, given by
$$\mathrm{\Omega }_{v,\mathrm{\Delta }}[\widehat{\rho }]=Tr\{\widehat{\rho }\widehat{H}_{v,\mathrm{\Delta }}+\theta \widehat{\rho }ln\widehat{\rho }\}.$$
(4)
3. Furthermore, the thermodynamic potential can be expressed in terms of the densities and potentials by involving a universal functional of the densities, $`F[n,\chi ]`$, as follows
$`\mathrm{\Omega }_{v,\mathrm{\Delta }}[n,\chi ]`$ $`=`$ $`F[n,\chi ]+{\displaystyle d^3rv(๐ซ)n(๐ซ)}`$ (5)
$``$ $`{\displaystyle }d^3rd^3r^{}[\mathrm{\Delta }^{}(๐ซ,๐ซ^{})\chi (๐ซ,๐ซ^{})+H.c.].`$
4. The universal functional contains the exchange-correlation (xc) free-energy functional, $`F_{xc}[n,\chi ]`$, as below
$`F[n,\chi ]`$ $`=`$ $`T_s[n,\chi ]\theta S_s[n,\chi ]\mu N`$ (6)
$`+`$ $`{\displaystyle \frac{1}{2}}{\displaystyle d^3rd^3r^{}\frac{n(๐ซ)n(๐ซ^{})}{|๐ซ๐ซ^{}|}}`$
$``$ $`{\displaystyle d^3r_1d^3r_1^{}d^3r_2d^3r_2^{}\chi ^{}(๐ซ_1,๐ซ_1^{})}`$
$`\times `$ $`w(๐ซ_1^{},๐ซ_1,๐ซ_2,๐ซ_2^{})\chi ^{}(๐ซ_2,๐ซ_2^{})+F_{xc}[n,\chi ],`$
where $`T_s[n,\chi ]`$ and $`S_s[n,\chi ]`$ are the kinetic energy and the entropy of a noninteracting system with the noninteracting potentials, $`v_s`$ and $`\mathrm{\Delta }_s`$, such that the densities $`n`$ and $`\chi `$ are equal to those of the noninteracting system. In the above formula, $`\mu `$ is the chemical potential.
5. The noninteracting grand-canonical Hamiltonian can be written in terms of the noninteracting densities and potentials as
$`\widehat{H}_s`$ $`=`$ $`{\displaystyle \underset{\sigma }{}}{\displaystyle d^3r\widehat{\psi }_\sigma ^{}(๐ซ)\left[\frac{^2}{2}+v_s(๐ซ)\mu \right]\widehat{\psi }_\sigma (๐ซ)}`$ (7)
$``$ $`[{\displaystyle }d^3rd^3r^{}\mathrm{\Delta }_s^{}(๐ซ,๐ซ^{})\widehat{\psi }_{}(๐ซ)\widehat{\psi }_{}(๐ซ^{})+H.c.].`$
6. The diagonalization of the noninteracting Hamiltonian, $`\widehat{H}_s`$, using the Bogoliubov transformation leads to the Kohn-Sham-Bogoliubov-de Gennes (KS-BdG) equations
$`\left[{\displaystyle \frac{^2}{2}}+v_s(๐ซ)\mu \right]u_i(๐ซ)`$ $`+`$ $`{\displaystyle d^3r^{}\mathrm{\Delta }_s(๐ซ,๐ซ^{})v_i(๐ซ^{})}`$ (8)
$`=`$ $`E_iu_i(๐ซ),`$
$`\left[{\displaystyle \frac{^2}{2}}+v_s(๐ซ)\mu \right]v_i(๐ซ)`$ $`+`$ $`{\displaystyle d^3r^{}\mathrm{\Delta }_s^{}(๐ซ,๐ซ^{})u_i(๐ซ^{})}`$ (9)
$`=`$ $`E_iv_i(๐ซ),`$
with $`v_i(๐ซ)`$ and $`u_i(๐ซ)`$ being the pair creation and anihilation amplitudes respectively.
7. The noninteracting potentials, $`v_s`$ and $`\mathrm{\Delta }_s`$, consist of the external potentials, $`v_0`$ and $`\mathrm{\Delta }_0`$, and Hartree potentials, and the exchange-correlation potentials, $`v_{xc}`$ and $`\mathrm{\Delta }_{xc}`$, as follows
$`v_s[n,\chi ](๐ซ)`$ $`=`$ $`v_0(๐ซ)+{\displaystyle d^3r^{}\frac{n(๐ซ^{})}{|๐ซ๐ซ^{}|}}`$ (10)
$`+`$ $`v_{xc}[n,\chi ](๐ซ),`$
$`\mathrm{\Delta }_s[n,\chi ](๐ซ,๐ซ^{})`$ $`=`$ $`\mathrm{\Delta }_0(๐ซ,๐ซ^{})+{\displaystyle d^3r^{}\frac{\chi (๐ซ,๐ซ^{})}{|๐ซ๐ซ^{}|}}`$ (11)
$`+`$ $`\mathrm{\Delta }_{xc}[n,\chi ](๐ซ,๐ซ^{}).`$
The external pairing potential has been introduced in order to break the symmetry, thus, in the calculations $`\mathrm{\Delta }_0(๐ซ,๐ซ^{})0`$ in Eq. (11).
8. The exchange-correlation potentials, $`v_{xc}`$ and $`\mathrm{\Delta }_{xc}`$, are defined as the derivatives of the $`xc`$ functional, $`F_{xc}[n,\chi ]`$, with respect to the densities, $`n`$ and $`\chi `$, correspondingly as below
$`v_{xc}[n,\chi ](๐ซ)`$ $`=`$ $`{\displaystyle \frac{\delta F_{xc}[n,\chi ]}{\delta n(๐ซ)}},`$ (12)
$`\mathrm{\Delta }_{xc}[n,\chi ](๐ซ,๐ซ^{})`$ $`=`$ $`{\displaystyle \frac{\delta F_{xc}[n,\chi ]}{\delta \chi ^{}(๐ซ,๐ซ^{})}}.`$ (13)
9. The densities, $`n`$ and $`\chi `$ are defined as functions of the amplitudes $`u_i(๐ซ)`$ and $`v_i(๐ซ)`$ as
$`n(๐ซ)`$ $`=`$ $`2{\displaystyle \underset{i}{}}[|u_i(๐ซ)|^2f_{\beta ,i}+|v_i(๐ซ)|^2(1f_{\beta ,i})],`$ (14)
$`\chi (๐ซ,๐ซ^{})`$ $`=`$ $`{\displaystyle \underset{i}{}}[v_i^{}(๐ซ^{})u_i(๐ซ)(1f_{\beta ,i})v_i^{}(๐ซ)u_i(๐ซ^{})f_{\beta ,i}],`$
with the Fermi distribution function $`f_{\beta ,i}=1+exp(\beta E_i)`$.
At this point, one could guess the densities, $`n`$ and $`\chi `$, and find the potentials, $`v_{xc}`$ and $`\mathrm{\Delta }_{xc}`$, and solve the KS-BdG equations, and find new densities etc. Further for practical reasons, as we already mentioned at the begin of this Section, one can make two approximations which we will discuss now.
1. The energy scales for the electronic energies and the superconducting energy gap differ by orders of magnitude. Therefore, the KS-BdG equations can be decoupled into the Kohn-Sham equation and the gap equation. This approximation was introduced to the SCDFT in Ref. SK-dec .
It holds within the decoupling approximation that:
1. the amplitudes $`u_i(๐ซ)`$ and $`v_i(๐ซ)`$ can be written in a form
$`u_i(๐ซ)u_i\phi _i(๐ซ);`$ $`v_i(๐ซ)v_i\phi _i(๐ซ),`$ (16)
2. the eigenvalues in Eqs. (8) and (9) are defined by
$$E_i=\pm \sqrt{\xi _i^2+|\mathrm{\Delta }_i}$$
(17)
where $`\xi _i=\epsilon _i\mu `$
3. the coefficients $`u_i`$ and $`v_i`$ are given by
$`u_i`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}sgn(E_i)e^{i\varphi _i}\sqrt{1+{\displaystyle \frac{\xi _i}{E_i}}},`$ (18)
$`v_i`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\sqrt{1{\displaystyle \frac{\xi _i}{E_i}}},`$ (19)
and the phase factor $`\varphi _i`$ is defined by
$$e^{i\varphi _i}=\frac{\mathrm{\Delta }_i}{|\mathrm{\Delta }_i|},$$
(20)
4. the matrix elements $`\mathrm{\Delta }_i`$ are defined as
$$\mathrm{\Delta }_i=d^3rd^3r^{}\phi _i^{}(๐ซ)\mathrm{\Delta }_s(๐ซ,๐ซ^{})\phi _i(๐ซ^{}),$$
(21)
5. and the normal and anomalous densities read respectively
$`n(๐ซ)={\displaystyle \underset{i}{}}\left(1{\displaystyle \frac{\xi _i}{E_i}}\right)\mathrm{tanh}\left({\displaystyle \frac{\beta E_i}{2}}\right)|\phi _i(๐ซ)|^2,`$ (22)
$`\chi (๐ซ,๐ซ^{})={\displaystyle \frac{1}{2}}{\displaystyle \underset{i}{}}{\displaystyle \frac{\mathrm{\Delta }_i}{E_i}}\mathrm{tanh}\left({\displaystyle \frac{\beta E_i}{2}}\right)\phi _i(๐ซ)\phi _i(๐ซ^{}).`$ (23)
2. The decoupling of the two energy scales yields a transformation of the KS-BdG equations into the ordinary Kohn-Sham equation
$$\left[\frac{^2}{2}+v_s[n,\chi ](๐ซ)\mu \right]\phi _i(๐ซ)=ฯต_i\phi _i(๐ซ),$$
(24)
and the gap equation
$$\mathrm{\Delta }_i=\mathrm{\Delta }_{Hxci}[\mu ,\mathrm{\Delta }_i].$$
(25)
The Eq. (25) stems from including Eqs. (22) and (23) into Eq. (11), and using the potential given by formula (11) in Eq. (21).
3. In vicinity of $`T_c`$, the gap function is vanishing, therefore, it can be linearized in $`\mathrm{\Delta }_i`$.
The above twelve steps lead to the gap equation which can be expressed in the form
$`\mathrm{\Delta }_i`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{j}{}}M_{Hxc,ij}[\mu ]{\displaystyle \frac{tanh(\frac{\beta }{2}\xi _j)}{\xi _j}}\mathrm{\Delta }_j,`$ (26)
$`M_{Hxc,ij}[\mu ]`$ $`=`$ $`{\displaystyle \frac{\delta \mathrm{\Delta }_{Hxc,i}}{\delta \chi _j}},`$ (27)
where $`\mathrm{\Delta }_{Hxc,i}`$ is defined by Eq. (13).
In other way, Eq. (27) can be written as
$$\mathrm{\Delta }_i=Z_i[\mu ]\mathrm{\Delta }_i\frac{1}{2}\underset{j}{}K_{ij}[\mu ]\frac{tanh(\frac{\beta }{2}\xi _j)}{\xi _j}\mathrm{\Delta }_j.$$
(28)
$`K_{ij}`$ and $`Z_i`$ are the functionals only of the chemical potential in the case when the gap equation is linearized. The above gap equation will be solved later in this work. The explicit form of the kernel $`K_{ij}`$ and the norm $`Z_i`$ will be given in Section 6.
Since the gap function (25) contains the exchange-correlation part defined by Eq. (13), we will focus on the construction of the exchange-correlation free-energy functional, $`F_{xc}`$, in the following Section.
## 3 Exchange-correlation functional, $`F_{xc}[n,\chi ]`$
The derivation of the exchange-correlation energy $`F_{xc}`$, by making use of the perturbative expansion of the thermodynamic potential, was given in Ref. SK . For the purpose of inclusion the spin interactions, we will briefly draw a skeleton of this derivation here.
First, one can notice from Eqs. (5) and (6) that
$`F_{xc}`$ $`=`$ $`\mathrm{\Omega }\mathrm{\Omega }_s+{\displaystyle d^3r[v_H(๐ซ)+v_{xc}(๐ซ)]n(๐ซ)}`$ (29)
$``$ $`{\displaystyle d^3rd^3r^{}[\mathrm{\Delta }_{xc}^{}(๐ซ,๐ซ^{})\chi (๐ซ,๐ซ^{})+\mathrm{\Delta }_{xc}(๐ซ,๐ซ^{})\chi ^{}(๐ซ,๐ซ^{})]}`$
$``$ $`{\displaystyle \frac{1}{2}}{\displaystyle d^3rd^3r^{}\frac{n(๐ซ)n(๐ซ^{})}{|๐ซ๐ซ^{}|}}.`$
Then, we take the coupling constant integration formula which reads
$$\mathrm{\Omega }\mathrm{\Omega }_s=_0^1\frac{d\lambda }{\lambda }\lambda \widehat{H_1},$$
(30)
where $`\lambda `$ is the coupling constant, and the perturbation Hamiltonian $`\widehat{H_1}`$ satisfies $`\widehat{H}=\widehat{H_s}+\lambda \widehat{H_1}`$ with the interacting and noninteracting Hamiltonians, $`\widehat{H}`$ and $`\widehat{H_s}`$, respectively. The Hamiltonian $`\widehat{H_1}`$ contains the difference between the exact Coulomb interaction and the exchange-correlation potentials, the electron-phonon interaction, the electron-paramagnon interaction etc.
The average in Eq. (30) has to be taken with the density operator $`\widehat{\rho }_\lambda =e^{\beta H_\lambda }/Z_\lambda `$.
Before an explicit evaluation of the coupling constant integration formula (30), we write here a definition of the Nambu Greenโs function
$`\overline{G}_{\sigma \sigma ^{}}(๐ซ\tau ,๐ซ^{}\tau ^{})=`$
$`\left(\begin{array}{cc}G_{\sigma \sigma ^{}}(๐ซ\tau ,๐ซ^{}\tau ^{})& F_{\sigma \sigma ^{}}(๐ซ\tau ,๐ซ^{}\tau ^{})\\ F_{\sigma \sigma ^{}}^{}(๐ซ\tau ,๐ซ^{}\tau ^{})& G_{\sigma ^{}\sigma }(๐ซ^{}\tau ^{},๐ซ\tau )\end{array}\right),`$ (33)
which is a 2$`\times `$2-matrix of the normal and anomalous single particle Greenโs functions, $`G_{\sigma \sigma ^{}}`$ and $`F_{\sigma \sigma ^{}}`$, given respectively by
$`G_{\sigma \sigma ^{}}(๐ซ\tau ,๐ซ^{}\tau ^{})`$ $`=`$ $`\widehat{T}\widehat{\psi }_\sigma (๐ซ\tau )\widehat{\psi }_\sigma ^{}^{}(๐ซ^{}\tau ^{}),`$ (34)
$`F_{\sigma \sigma ^{}}(๐ซ\tau ,๐ซ^{}\tau ^{})`$ $`=`$ $`\widehat{T}\widehat{\psi }_\sigma (๐ซ\tau )\widehat{\psi }_\sigma ^{}(๐ซ^{}\tau ^{}),`$ (35)
$`F_{\sigma \sigma ^{}}^{}(๐ซ\tau ,๐ซ^{}\tau ^{})`$ $`=`$ $`\widehat{T}\widehat{\psi }_\sigma ^{}(๐ซ\tau )\widehat{\psi }_\sigma ^{}^{}(๐ซ^{}\tau ^{}).`$ (36)
The detailed derivation of $`\widehat{H}_1`$ is given in Refs. SK ; ML ; MM . This derivation starts from the equations of motion for the field operator, $`\widehat{\psi }_\sigma `$, and for the noninteracting Greenโs function, $`\overline{G}_{\sigma \sigma ^{}}^s`$, which are as follows
$`{\displaystyle \frac{}{\tau }}\widehat{\psi }_\sigma (๐ซ\tau )`$ $`=`$ $`e^{\widehat{H}\tau }[\widehat{H},\widehat{\psi }_\sigma (๐ซ)]e^{\widehat{H}\tau },`$ (37)
$`\widehat{}\overline{G}_{\sigma \sigma ^{}}^s(๐ซ\tau ,๐ซ^{}\tau ^{})`$ $`=`$ $`\delta _{\sigma \sigma ^{}}\delta (๐ซ๐ซ^{})\delta (\tau \tau ^{}),`$ (38)
with the Kohn-Sham Hamiltonian for the normal state, $`\widehat{h}_s`$, and the operator $`\widehat{}`$ given respectively by
$`\widehat{h}_s(๐ซ)`$ $`=`$ $`{\displaystyle \frac{^2}{2}}+v_s(๐ซ)\mu ,`$ (39)
$`\widehat{}`$ $`=`$ $`\left(\begin{array}{cc}\frac{}{\tau }+\widehat{h}_s(๐ซ)& \widehat{\mathrm{\Delta }}_s(๐ซ)\\ \widehat{\mathrm{\Delta }}_s^{}(๐ซ)& \frac{}{\tau }\widehat{h}_s(๐ซ)\end{array}\right).`$ (42)
The operator $`\widehat{\mathrm{\Delta }}_s(๐ซ)`$ is defined as
$$\widehat{\mathrm{\Delta }}_s(๐ซ)f(๐ซ)=d^3r^{}\widehat{\mathrm{\Delta }}_s(๐ซ,๐ซ^{})f(๐ซ^{}).$$
(43)
In order to complete the derivation, one also needs to make use of the Dysonโs equation
$`\overline{G}_{\sigma ^{}\sigma }(๐ซ\tau ,๐ซ^{}\tau ^{})`$ $`=`$ $`\overline{G}_{\sigma ^{}\sigma }^s(๐ซ\tau ,๐ซ^{}\tau ^{})`$ (44)
$`+`$ $`{\displaystyle \underset{\sigma \sigma ^{}}{}}{\displaystyle d^3r_1d^3r_2๐\tau _1๐\tau _2\overline{G}_{\sigma \sigma _1}^s(๐ซ\tau ,๐ซ_1\tau _1)}`$
$`\times `$ $`\overline{\mathrm{\Sigma }}(๐ซ_1\tau _1,๐ซ_2\tau _2)\overline{G}_{\sigma _2\sigma ^{}}(๐ซ_2\tau _2,๐ซ^{}\tau ^{}),`$
with $`\overline{\mathrm{\Sigma }}`$ being the self-energy.
The above building blocks make us to arrive, after some algebra, at the relation
$`\mathrm{\Omega }\mathrm{\Omega }_s`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _0^1}{\displaystyle \frac{d\lambda }{\lambda }}\{{\displaystyle \underset{\sigma \sigma ^{}}{}}{\displaystyle }d^3rd^3r^{}`$ (45)
$`\times `$ $`{\displaystyle ๐\tau ^{}[\overline{\mathrm{\Sigma }}_{\sigma \sigma ^{}}^\lambda (๐ซ\tau ,๐ซ^{}\tau ^{})\overline{G}_{\sigma ^{}\sigma }^\lambda (๐ซ^{}\tau ^{},๐ซ\tau ^+)]_{11}}`$
$``$ $`\lambda {\displaystyle d^3r[v_H(๐ซ)+v_{xc}(๐ซ)]n^\lambda (๐ซ)}`$
$`+`$ $`2\lambda {\displaystyle }d^3rd^3r^{}\mathrm{\Delta }_{xc}^{}(๐ซ,๐ซ^{})\chi ^\lambda (๐ซ,๐ซ^{})\},`$
which we can plug into the Eq. (29) for the exchange-correlation functional, $`F_{xc}[n,\chi ]`$.
As for the first-order selfenergy, $`\overline{\mathrm{\Sigma }}_{\sigma \sigma ^{}}`$, for the nonmagnetic systems with the potential $`v(๐ซ\tau ,๐ซ^{}\tau ^{})`$, this energy is defined as
$$\overline{\mathrm{\Sigma }}(๐ซ\tau ,๐ซ^{}\tau ^{})=v(๐ซ\tau ,๐ซ^{}\tau ^{})\tau _3\overline{G}(๐ซ\tau ,๐ซ^{}\tau ^{})\tau _3,$$
(46)
and $`\tau _3`$ is one of the Pauli matrices:
$`\tau _1=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right),`$ $`\tau _2=\left(\begin{array}{cc}0& i\\ i& 0\end{array}\right),`$ (51)
$`\tau _3=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right),`$ $`\tau _0=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right).`$ (56)
For the magnetic systems, the matrix $`\tau _3`$ in the each vertex of Feynman diagrams for the selfenergy with the Coulomb and phonon interactions has to be replaced with the matrix $`\tau _0\tau _3`$.
In this Section, we sketched main steps to be done for finding a general form of the $`F_{xc}[n,\chi ]`$ functional for a superconductor. The final formula involves the selfenergy which will be evaluated in detail for the Coulomb and electron-phonon interactions in the next Section and for the paramagnons in Section 5.
## 4 Coulomb and electron-phonon interactions in $`F_{xc}[n,\chi ]`$
The derivation of $`F_{xc}`$ for the Coulomb and phonon interactions is given in detail in Refs. Martin-OGK ; ML . Here, we report this derivation starting with the interactions in the selfenergy (in Eq. (46)) defined by
$`v^{el}(๐ซ,๐ซ^{})`$ $`=`$ $`{\displaystyle \frac{1}{|๐ซ๐ซ^{}|}},`$ (57)
$`v^{ph}(๐ซ\tau ,๐ซ^{}\tau ^{})`$ $`=`$ $`V_{\lambda ๐ช}(๐ซ)D_{\lambda ๐ช}(\tau \tau ^{})V_{\lambda ๐ช}(๐ซ^{}),`$ (58)
where $`V_{\lambda ๐ช}`$ is the electron-phonon interaction vertex and $`D_{\lambda ๐ช}`$ is the phonon Greenโs function defined as
$$D_{\lambda ๐ช}(\tau ,\tau ^{})=\widehat{T}\widehat{\mathrm{\Phi }}_{\lambda ๐ช}(\tau )\widehat{\mathrm{\Phi }}_{\lambda ๐ช}^{}(\tau ^{}),$$
(59)
with $`\widehat{\mathrm{\Phi }}_{\lambda ๐ช}=b_{\lambda ,๐ช}+b_{\lambda ,๐ช}^{}`$, and $`b_{\lambda ,๐ช}^{}`$ ($`b_{\lambda ,๐ช}`$) being the phonon creation (anihilation) operators.
Let us have a look now at the expression (29) for $`F_{xc}`$ and the definitions of the Nambu Greenโs function and selfenergy given by Eqs. (33) and (46) respectively. The (1,1)-element of the ($`\overline{\mathrm{\Sigma }}\overline{G}`$)-matrix, present in the formula (45) and enterring Eq. (29), is proportional to
$$G_{}G_{}F_{}F_{}^{}=G_{}G_{}+F_{}F_{}^{}$$
(60)
and the corresponding terms with the opposite spins. The above terms appear for both the Coulomb and electron-phonon interactions, and later will lead to the opposite signums in the kernel $`K_{ij}`$ and the norm $`Z_i`$ of the gap equation. Just mentioned difference in signum, in the first order terms of the total energy with the normal and anomalous Greenโs functions, stems from the factor of (-1) which one has to associate with the each loop of anomalous Greenโs functions.
In order to evaluate further $`F_{xc}`$, we bring here the explicite expressions for the noninteracting propagators. The formulas given below were derived from the definitions (34-36) assumming the decoupling approximation, i.e. Eqs. (16); the Kohn-Sham orbitals $`\phi _๐ค(๐ซ)`$ were chosen to those of a homogeneous gas ($`w_n`$ are the odd Matsubara frequencies)
$`G_{\sigma \sigma ^{}}^s(๐ค,w_n)=\delta _{\sigma ,\sigma ^{}}`$
$`\times \left[{\displaystyle \frac{|u_๐ค|^2}{i\omega _nE_๐ค}}+{\displaystyle \frac{|v_๐ค|^2}{i\omega _n+E_๐ค}}\right],`$ (61)
$`F_{\sigma \sigma ^{}}^s(๐ค;w_n)=\delta _{\sigma ,\sigma ^{}}sgn(\sigma ^{})`$
$`\times u_๐คv_๐ค^{}\left({\displaystyle \frac{1}{i\omega _n+E_๐ค}}{\displaystyle \frac{1}{i\omega _nE_๐ค}}\right),`$ (62)
$`F_{\sigma \sigma ^{}}^s(๐ค;w_n)=\delta _{\sigma ,\sigma ^{}}sgn(\sigma )`$
$`\times u_๐ค^{}v_๐ค\left({\displaystyle \frac{1}{i\omega _n+E_๐ค}}{\displaystyle \frac{1}{i\omega _nE_๐ค}}\right).`$ (63)
Now, we will combine Eqs. (29) and (45), for the $`F_{xc}`$ and $`\mathrm{\Omega }\mathrm{\Omega }_s`$ respectively, with a definition of the Nambu Greenโs function, Eq. (33), and an expression for the selfenergy, Eq. (46). As for the noniteracting Greenโs functions, we use those obtained within the decoupling approximation, i.e. (61-63). This way, one arrives to formulas for the xc energy, stemming from the normal and anomalous loops, which we write here. The โnormalโ and โanomalousโ terms of $`F_{xc}`$ for the electronic contributions, $`F_{xc}^{el,1}`$ and $`F_{xc}^{el,2}`$, are as follows
$`F_{xc}^{el,1}`$ $`=`$ $`{\displaystyle \frac{1}{4}}{\displaystyle \underset{\mathrm{๐ค๐ค}^{}}{}}\left(1{\displaystyle \frac{\xi _๐ค}{E_๐ค}}\right)v(๐ค,๐ค^{})\left(1{\displaystyle \frac{\xi _๐ค^{}}{E_๐ค^{}}}\right)`$ (64)
$`\times \mathrm{tanh}\left({\displaystyle \frac{\beta }{2}}E_๐ค\right)\mathrm{tanh}\left({\displaystyle \frac{\beta }{2}}E_๐ค^{}\right),`$
$`F_{xc}^{el,2}`$ $`=`$ $`{\displaystyle \frac{1}{4}}{\displaystyle \underset{\mathrm{๐ค๐ค}^{}}{}}v(๐ค,๐ค^{}){\displaystyle \frac{\mathrm{\Delta }_๐ค}{E_๐ค}}{\displaystyle \frac{\mathrm{\Delta }_๐ค^{}}{E_๐ค^{}}}`$ (65)
$`\times \mathrm{tanh}\left({\displaystyle \frac{\beta }{2}}E_๐ค\right)\mathrm{tanh}\left({\displaystyle \frac{\beta }{2}}E_๐ค^{}\right),`$
and the electron-phonon terms, with the normal and anomalous loops, $`F_{xc}^{ph,1}`$ and $`F_{xc}^{ph,2}`$i, respectively are given below
$`F_{xc}^{ph,1}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{\mathrm{๐ค๐ค}^{}}{}}{\displaystyle ๐\mathrm{\Omega }\alpha ^2F(\mathrm{\Omega })}`$ (66)
$`\times [(1+{\displaystyle \frac{\xi _๐ค\xi _๐ค^{}}{E_๐คE_๐ค^{}}})I(E_๐ค,E_๐ค^{},\mathrm{\Omega })`$
$`+(1{\displaystyle \frac{\xi _๐ค\xi _๐ค^{}}{E_๐คE_๐ค^{}}})I(E_๐ค,E_๐ค^{},\mathrm{\Omega })],`$
$`F_{xc}^{ph,2}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{\mathrm{๐ค๐ค}^{}}{}}{\displaystyle ๐\mathrm{\Omega }\alpha ^2F(\mathrm{\Omega })\frac{\mathrm{\Delta }_๐ค\mathrm{\Delta }_๐ค^{}^{}}{E_๐คE_๐ค^{}}}`$ (67)
$`\times [I(E_๐ค,E_๐ค^{},\mathrm{\Omega })I(E_๐ค,E_๐ค^{},\mathrm{\Omega })].`$
The function $`I(E_๐ค,E_๐ค^{},\mathrm{\Omega })`$ is defined as
$`I(E_๐ค,E_๐ค^{},\mathrm{\Omega })`$ $`=`$ $`{\displaystyle \frac{1}{\beta ^2}}{\displaystyle \underset{\omega _1\omega _2}{}}{\displaystyle \frac{1}{i\omega _1E_๐ค}}{\displaystyle \frac{1}{i\omega _2E_๐ค^{}}}`$ (68)
$`\times {\displaystyle \frac{2\mathrm{\Omega }}{(\omega _1\omega _2)^2+\mathrm{\Omega }^2}}.`$
For the completeness, we give the definitions:
$`v(๐ค,๐ค^{})`$ $`=`$ $`{\displaystyle d^3rd^3r^{}\phi _๐ค^{}(๐ซ)\phi _๐ค(๐ซ^{})}`$ (69)
$`\times {\displaystyle \frac{1}{|๐ซ๐ซ^{}|}}\phi _๐ค^{}^{}(๐ซ)\phi _๐ค^{}(๐ซ^{}),`$
$`g_{๐ค,๐ค+๐ช}^{\lambda ๐ช}`$ $`=`$ $`{\displaystyle d^3r\phi _๐ค^{}(๐ซ)V_{\lambda ๐ช}\phi _{๐ค+๐ช}(๐ซ)},`$ (70)
$`\alpha ^2F(\mathrm{\Omega })`$ $`=`$ $`{\displaystyle \frac{1}{N(\epsilon _F)}}{\displaystyle \underset{\lambda ๐ช}{}}{\displaystyle \underset{๐ค}{}}|g_{๐ค,๐ค+๐ช}^{\lambda ๐ช}|^2\delta (\mathrm{\Omega }\omega _{\lambda ๐ช})`$ (71)
$`\times \delta (\epsilon _๐ค\epsilon _F)\delta (\epsilon _{๐ค+๐ช}\epsilon _F),`$
where $`\omega _{\lambda ๐ช}`$ is the phonon frequency and $`N(\epsilon _F)`$ is the density of states.
Using the formulas (64-68), one is ready to derive the exchange-correlation potential defined by Eq. (13). This derivation can be performed with the help of the chain rule as follows
$`\mathrm{\Delta }_{xc,i}`$ $`=`$ $`{\displaystyle \frac{\delta F_{xc}}{\delta \mu }}{\displaystyle \frac{\delta \mu }{\delta \chi _i^{}}}{\displaystyle \underset{j}{}}[{\displaystyle \frac{\delta F_{xc}}{\delta |\mathrm{\Delta }_j|^2}}{\displaystyle \frac{\delta |\mathrm{\Delta }_j|^2}{\delta \chi _i^{}}}`$ (72)
$`+{\displaystyle \frac{\delta F_{xc}}{\delta (\varphi _j)}}{\displaystyle \frac{\delta (\varphi _j)}{\delta \chi _i^{}}}].`$
Further evaluation of the above expression is given in detail in Refs. Martin-OGK ; SK . In this work, we give the final formula for $`\mathrm{\Delta }_{xc,i}`$ which involves the phonon and paramagnon spectral functions and can be implemented in a straightforward way. We will give the details of implementation in Section 6.
At this point, we arrived to the explicit expressions for $`F_{xc}`$ with the electronic and phononic parameters such as: the chemical potential $`\mu `$, the density of states $`N(\epsilon _F)`$, the single particle energies $`\epsilon _๐ค`$, and the Eliashberg function $`\alpha ^2F(\mathrm{\Omega })`$. Now, we are ready to introduce the spin fluctuations into the discussed formalism, and we will this in the following Section.
## 5 Paramagnons in $`F_{xc}[n,\chi ]`$
We will introduce the transverse spin-fluctuations to the total energy within the SCDFT. For the simplicity, we will assume the singlet pairing and the $`s`$-wave symmetry of the gap function. The extension to triplet superconductors could be done following the work by Capelle et al. KC ; triplet . For the case of magnetic superconductors, one should take also into account a correction for the Zeeman effect, i.e. the spin gap. As for the paring potentials with the higher angular-momentum, one cannot average spherically the angular part of the interaction in the RPA formula for the paramagnon susceptibility, which formula will be used later in this Section.
Here, we start with the Nambu Greenโs function for the superconductors with magnetic interactions included into the description. This matrix is now 4$`\times `$4 dimentional and reads
$$\overline{G}(๐ซ\tau ,๐ซ^{}\tau ^{})=\widehat{T}\widehat{\mathrm{\Psi }}^{}(๐ซ,\tau )\widehat{\mathrm{\Psi }}(๐ซ^{},\tau ^{}),$$
(73)
with the 4-component field operators (the notation has been chosen according to Maki in Ref. Kazumi and $`x`$ denotes the vector (r,$`\tau `$))
$`\widehat{\mathrm{\Psi }}(x)=\left(\begin{array}{c}\widehat{\psi }_{}(x)\\ \widehat{\psi }_{}(x)\\ \widehat{\psi }_{}^{}(x)\\ \widehat{\psi }_{}^{}(x)\end{array}\right),`$ $`\widehat{\mathrm{\Psi }}^{}(x)=\left(\widehat{\psi }_{}^{}(x)\widehat{\psi }_{}^{}(x)\widehat{\psi }_{}(x)\widehat{\psi }_{}(x)\right).`$ (78)
The first-order selfenergy with the spin dependent interaction $`v^{\mu \nu }`$, where $`\mu `$ and $`\nu `$ denote the cartesian components of the spin orientations of two interacting electrons, is given by
$`\overline{\mathrm{\Sigma }}(๐ซ\tau ,๐ซ^{}\tau ^{})`$ $`=`$ $`v^{\mu \nu }(๐ซ\tau ,๐ซ^{}\tau ^{})\widehat{\alpha }_\mu \overline{G}(๐ซ\tau ,๐ซ^{}\tau ^{})\widehat{\alpha }_\nu ,`$ (80)
$`v^{\mu \nu }(๐ซ\tau ,๐ซ^{}\tau ^{})`$ $`=`$ $`I_{ex}(๐ซ)D^{\mu \nu }(\tau \tau ^{})I_{ex}(๐ซ^{}).`$ (81)
The quantity $`I_{ex}`$ is the spin exchange interaction, and $`D^{\mu \nu }`$ is the spin Greenโs function. The matrix $`\widehat{\alpha }_\mu `$ is defined as
$$\widehat{\alpha }_\mu =\left(\begin{array}{cc}\sigma _\mu & 0\\ 0& \sigma _\mu ^{tr}\end{array}\right),$$
(82)
where $`\sigma _\mu ^{tr}`$ denotes a matrix transposed to the Pauli matrix $`\sigma _\mu `$ (see Ref. Vonsovsky ).
For the transverse spin fluctuations, the $`\alpha `$-matrix, given by formula (82), involves the Pauli matrices $`\sigma ^+`$ and $`\sigma ^{}`$ defined as $`\sigma ^\pm =\frac{1}{2}(\sigma _x\pm i\sigma _y)`$; explicitely
$`\sigma ^+=\left(\begin{array}{cc}0& 1\\ 0& 0\end{array}\right),`$ $`\sigma ^{}=\left(\begin{array}{cc}0& 0\\ 1& 0\end{array}\right).`$ (87)
Evaluation of the selfenergy with paramagnons, according to Eqs. (80-82), yields a very sparse 4$`\times `$4-matrix which reads
$`\overline{\mathrm{\Sigma }}(๐ซ\tau ,๐ซ^{}\tau ^{})=v^+(๐ซ\tau ,๐ซ^{}\tau ^{})`$
$`\times \left(\begin{array}{cccc}G_{}(๐ซ\tau ,๐ซ^{}\tau ^{})& 0& 0& F_{}(๐ซ\tau ,๐ซ^{}\tau ^{})\\ 0& 0& 0& 0\\ 0& 0& 0& 0\\ F_{}^{}(๐ซ\tau ,๐ซ^{}\tau ^{})& 0& 0& G_{}^{}(๐ซ\tau ,๐ซ^{}\tau ^{})\end{array}\right),`$ (92)
(93)
where $`G_{}^{}=G_{}`$.
Now, if we go back to the previous section and look again at the (1,1)-element of the ($`\overline{\mathrm{\Sigma }}\overline{G}`$)-matrix, we will remind to us that for the Coulomb and electron-phonon interactions, the total energy is proportional to the expression (60). For the magnetic interactions, however, for which the Nambu Greenโs function has been defined by Eq. (73) and the selfenergy has been given by Eq. (93), the total energy is proportional to
$$G_{}G_{}F_{}F_{}^{}=G_{}G_{}F_{}F_{}^{}.$$
(94)
The above expression differs from relation (60) by signum in front of the anomalous Greenโs functions. This difference will show up in the kernel $`K_{ij}`$ and the norm $`Z_i`$ of the gap equation such that, both the phonon and paramagnon spectral functions enter the kernel with different signum (originating from the anomalous loop of Greenโs functions) and the norm with the same signum (originating from the normal loop).
To proceed further with the evaluation of the xc-free energy, $`F_{xc}`$, we write explicitely the spin-fluctuation Greenโs function, $`D^{\mu \nu }(\tau \tau ^{})`$, used in Eq. (81). In the case of paramagnons, $`D^{\mu \nu }(\tau \tau ^{})`$ is the transverse spin susceptibility, $`\chi ^+`$, defined as
$$\chi ^+(๐ซ๐ซ^{},\tau \tau ^{})=\widehat{T}\widehat{S}^{}(๐ซ,\tau )\widehat{S}^+(๐ซ^{},\tau ^{}),$$
(95)
with the operators increasing and lowerring spin which are defined respectively as
$`\widehat{S}^+(๐ซ,\tau )`$ $`=`$ $`\widehat{\psi }_{}^{}(๐ซ,\tau )\widehat{\psi }_{}(๐ซ,\tau ),`$ (96)
$`\widehat{S}^{}(๐ซ,\tau )`$ $`=`$ $`\widehat{\psi }_{}^{}(๐ซ,\tau )\widehat{\psi }_{}(๐ซ,\tau ).`$ (97)
For the conduction band, we can use a model of the homogeneous electron gas with the fluctuations treated on the level of the random phase approximation. The Fourier transform of the RPA-โdressedโ paramagnon propagator is
$$\chi ^+(๐ช,\nu _n)=\frac{\chi ^0(๐ช,\nu _n)}{1I_{ex}\chi ^0(๐ช,\nu _n)},$$
(98)
with the Pauli susceptibility $`\chi ^0`$ and the even Matsubara frequencies $`\nu _n`$.
It is convenient to introduce the spectral representation
$`\chi ^+(๐ช,\nu _n)`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{d\mathrm{\Omega }}{\pi }}D^0(\mathrm{\Omega },\nu _n)\mathrm{}m\chi ^+(๐ช,\mathrm{\Omega }),`$
$`D^0(\mathrm{\Omega },\nu _n)`$ $`=`$ $`{\displaystyle \frac{2\mathrm{\Omega }}{\nu _n^2+\mathrm{\Omega }^2}},`$ (100)
and the momentum averaged paramagnon spectral function
$`P(\mathrm{\Omega })`$ $`=`$ $`N(\epsilon _F){\displaystyle _0^{2k_F}}๐q{\displaystyle \frac{q}{2k_F^{\mathrm{\hspace{0.33em}\hspace{0.33em}2}}}}`$ (101)
$`\times |I(q)|^2\left[{\displaystyle \frac{1}{\pi }}\mathrm{}m\chi ^+(q,\mathrm{\Omega })\right].`$
We assume that the interaction function, $`I(q)`$, is the momentum independent quantity $`I_{ex}`$, which can be calculated in a way given for instance in Ref. our-Nb .
Therefore, for the systems with the electron-paramagnon interactions, the exchange-correlation free energy is given by
$`F_{xc}^{sf,1}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{\mathrm{๐ค๐ค}^{}}{}}{\displaystyle ๐\mathrm{\Omega }P(\mathrm{\Omega })}`$ (102)
$`\times [(1+{\displaystyle \frac{\xi _๐ค\xi _๐ค^{}}{E_๐คE_๐ค^{}}})I(E_๐ค,E_๐ค^{},\mathrm{\Omega })`$
$`+(1{\displaystyle \frac{\xi _๐ค\xi _๐ค^{}}{E_๐คE_๐ค^{}}})I(E_๐ค,E_๐ค^{},\mathrm{\Omega })],`$
$`F_{xc}^{sf,2}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{\mathrm{๐ค๐ค}^{}}{}}{\displaystyle ๐\mathrm{\Omega }P(\mathrm{\Omega })\frac{\mathrm{\Delta }_๐ค\mathrm{\Delta }_๐ค^{}^{}}{E_๐คE_๐ค^{}}}`$ (103)
$`\times [I(E_๐ค,E_๐ค^{},\mathrm{\Omega })I(E_๐ค,E_๐ค^{},\mathrm{\Omega })],`$
where the function $`I(E_๐ค,E_๐ค^{},\mathrm{\Omega })`$ is defined by Eq. (68). The explicit formula for the paramagnon spectral function, $`P(\mathrm{\Omega })`$, within the RPA is given for instance in Refs. Berk ; Vonsovsky ; our-Nb .
## 6 Gap equation with paramagnons and implementation details
At this point, when we have completed the derivation of all components of the exchange-correlation free energy: the Coulomb part - Eqs. (64-65), the phonon part - Eqs. (66-67), and the spin-fluctuation part - Eqs. (102-103), we can write explicitely the gap equation given by Eqs. (26-28).
The $`M_{ij}`$-matrix of the linearized equation (26) is the following function of the kernel $`K_{ij}`$ and the norm $`Z_i`$
$$M_{ij}=\frac{1}{2}\frac{K_{ij}\left[\mathrm{\Delta }=0\right]}{1Z_i\left[\mathrm{\Delta }=0\right]}.$$
(104)
The nondiagonal part of the $`M_{ij}`$-matrix is given by
$$K_{ij}=K_{ij}^{el}+K_{ij}^{ph+sf},$$
(105)
where the electronic part is defined by
$`K_{ij}^{el}`$ $`=`$ $`w_{ij},`$ (106)
$`w_{ij}`$ $`=`$ $`{\displaystyle \frac{2\pi }{k_ik_j}}log\left({\displaystyle \frac{(k_i+k_j)^2+k_{TF}^2}{(k_ik_j)^2+k_{TF}^2}}\right).`$ (107)
The Coulomb interaction $`w_{ij}`$ has been spherically averaged over the angular coordinates since, as we said before, we assumed the $`s`$-wave pairing. The electron correlations are taken into account by the Thomas-Fermi screening constant, $`k_{TF}`$, and $`k_i`$ is an absolute value of the reciprocal vector.
The electron-phonon and -paramagnon interaction diagonal part of the $`M_{ij}`$-matrix is given by
$`K_{ij}^{ph+sf}`$ $`=`$ $`{\displaystyle \frac{2}{tanh(\beta \xi _i/2)tanh(\beta \xi _j/2)}}`$ (108)
$`\times {\displaystyle }d\mathrm{\Omega }[\alpha ^2F(\mathrm{\Omega })P(\mathrm{\Omega })]`$
$`\times \left[I(\xi _i,\xi _j,\mathrm{\Omega })I(\xi _i,\xi _j,\mathrm{\Omega })\right].`$
The diagonal part of the $`M_{ij}`$-matrix is
$$Z_i=Z_i^{el}+Z_i^{ph+sf},$$
(109)
where the purely electronic part is
$`Z_i^{el}`$ $`=`$ $`{\displaystyle \frac{1}{2\xi _i}}\{{\displaystyle \underset{j}{}}w_{ij}[1tanh(\beta \xi _j/2)]`$ (110)
$`{\displaystyle \frac{_{jk}\frac{\beta w_{jk}/2}{cosh^2(\beta \xi _j/2)}\left[1tanh(\beta \xi _k/2)\right]}{_k\frac{\beta /2}{cosh^2(\beta \xi _k/2)}}}\},`$
and the phononic and paramagnon part is
$`Z_i^{ph+sf}`$ $`=`$ $`{\displaystyle \frac{4\pi }{tanh(\beta \xi _i/2)}}{\displaystyle \frac{1}{\beta }}{\displaystyle ๐\mathrm{\Omega }\left[\alpha ^2F(\mathrm{\Omega })+P(\mathrm{\Omega })\right]}`$ (111)
$`\times {\displaystyle \underset{\omega _2}{}}\omega _2sgn(\omega _2)[Z_{i,sym}^{ph+sf}+Z_{i,asym}^{ph+sf}],`$
$`Z_{i,sym}^{ph+sf}`$ $`=`$ $`\left[n_\beta (\mathrm{\Omega })+f_\beta (\xi _i)\right]{\displaystyle \frac{2(\xi _i+\mathrm{\Omega })}{\left[\omega _2^2+(\xi _i+\mathrm{\Omega })^2\right]^2}}`$ (112)
$`+\left[n_\beta (\mathrm{\Omega })+f_\beta (\xi _i)\right]{\displaystyle \frac{2(\xi _i\mathrm{\Omega })}{\left[\omega _2^2+(\xi _i\mathrm{\Omega })^2\right]^2}}.`$
Functions $`f_\beta `$ and $`n_\beta `$ are the Fermi-Dirac and Bose-Einstein distribution functions respectively.
For the electronic part of the norm, i.e. $`Z_i^{el}`$, we used the zero temperature approximation given in Refs. MM ; our-Nb . This approximation can be justified by the fact that the critical temperatures of simple metals, which we calculate in this work, are very low. The above simplification is done for sake of a numerical convenience since there are many singularities in the formula (110).
The subscripts โsymโ and โasymโ mean the symmetric and antisymmetric part of $`Z_i`$ with respect to the electron-phonon coupling elements $`g_{๐ค,๐ค+๐ช}`$. The electron-paramagnon interaction constant, $`I(๐ช)`$, has been also averaged in q leading to $`I_{ex}`$. The antisymmetric part $`Z_{i,asym}^{ph+sf}`$ is ommited in our calculations according to the reasons discussed in Refs. Martin-OGK ; MM and in our previous work our-Nb . Therefore, we do not give the expression for $`Z_{i,asym}^{ph+sf}`$ in this work.
## 7 Critical temperatures of simple metals
In the following two subsections, we report the critical temperatures obtained by solving the SCDFT gap equation with spin fluctuations included. We compare these results with the results without spin fluctuations and results from the Eliashberg theory. First, we calculate parameters of the gap equation for several simple metals: V, Mo, Ta, and Pd (fcc and bcc) at ambient pressure. At the end, we complete our previous results for Nb under pressure our-Nb reporting $`T_c`$ obtained within the SCDFT with the paramagnons included.
The electronic structures, the densities of states (DOS) and the electron-phonon coupling constants and the phonon and magnon spectral functions for studied metals were calculated within the local density approximation (LDA). We used the pseudopotential plane wave codes pwscf PWscf and espresso espresso . The phonons and electron-phonon couplings were obtained from the density functional perturbation theory (DFPT) Review . Since the calculation of the spectral function $`\alpha ^2F`$ is very time consuming, we used the ultrasoft pseudopotentials (US PPs) Vanderbilt . The kinetic energy cut-offs for the wavefunctions and densities were 45 Ry and 270 Ry respectively in order to reproduce well all features of the phonon dispersions especially for the low frequency phonons (see Ref. our-Nb ). The metallic broadening at the Fermi energy Paxton was assumed at 0.03 Ry. We used the Monkhorst-Pack mesh mesh of (64,64,64)-points for the DOS calculations, (16,16,16)-points for the self-consistent calculation of the electron-phonon-coupling matrix elements for the each phonon, the mesh of (8,8,8)-points to fit the phonon dispersions, and the fit from (16,16,16) into (64,64,64) mesh-points to perform the integrands with the double-delta function present in the definition of the electron-phonon coupling constant, $`\lambda ^{ph}`$, and the spectral function, $`\alpha ^2F(\omega )`$.
The spin-exchange interaction contants, $`I_{ex}`$, for metals at ambient pressure were taken from the work by Singalas et al. Singalas , and further we used them for the calculation of the spectral functions, $`P(\omega )`$, and the electron-paramagnon coupling constant, $`\lambda ^{sf}`$. For niobium under pressure, we used $`I_{ex}`$ and $`P(\omega )`$ calculated in our previous work our-Nb .
All electronic parameters and the phonon and magnon spectral functions were assumed to be the same for the normal and superconducting state. The accuracy of functions $`\alpha ^2F(\omega )`$ and $`P(\omega )`$ is very important for an exact estimation of the critical temperature. The electron-phonon spectral function, very time consumming for calculations, contains all the specific information about the studied system. In contrast to $`\alpha ^2F(\omega )`$, the approximation which we used for the paramagnon spectral function, to avoid calculation of this quantity from the time-dependent density functional theory, is insuficient. We made the assumption of the homogeneous electron gas for the spin susceptibility and the only spin-dependent quantity which we calculated specifically for a given metal was the exchange constant. The calculation of this constant, i.e. $`I_{ex}`$, is very difficult and obtained results have a large error due to their very small values and necessity to calculate a response function to small magnetizations applied to the system. Therefore, as we will see below, the obtained critical temperatures are not always very close to the experimental ones. Further development should be directed into more accurate calculation of the spectral functions, especially $`P(\omega )`$.
### 7.1 Transition metals at ambient pressure
In TABLE 1, we report the critical temperatures and parameters which enter the gap equation calculated by means of the Eliashberg theory and the SCDFT for a few simple metals: vanadium, molibdenium and tantallum in bcc lattice structure and palladium in fcc and bcc structures. Our calculated densities of states, $`N(\epsilon _F)`$, and electron-phonon coupling constants, $`\lambda ^{ph}`$, are in a good agreement with previous calculations by Savrasov et al. Sav-lambda . The Eliashberg functions calculated within the DFPT are presented in FIG. 1. The Coulomb parameter, $`\mu ^{}`$, was obtained from the Bennemen-Garland formula mu ; our-Nb , which employs the density of states. The spin exchange constant, $`I_{ex}`$, taken from Ref. Singalas , has been used to obtain the paramagnon spectral function, $`P(\omega )`$, which we draw in FIG. 2.
As for the critical temperatures, for tantallum, the SCDFT result is in a very small relative error, defined in TABLE 1, of 3$`\%`$ with respect to the experimental data expTc . While, the Eliashberg result with spin fluctuations included is in the error of 81$`\%`$. For molibdenium, $`T_c`$ from the Eliashberg gap equation is smaller than the experimental one, even without the paramagnon effect. But the absolute error of all calculated temperatures for Mo is smaller than 1 K. Palladium in both structures fcc and bcc is nonsuperconducting and the SCDFT reproduces well this result. In contrast to the SCDFT result, from the Eliashberg theory we obtained superconductivity for Pd in the bcc structure with a very small $`T_c`$.
Usually, the critical temperatures from the SCDFT are lower than temperatures from the Eliashberg theory. In some cases, however, the SCDFT temperatures are higher. This situation is for vanadium and molibdenium. Especially for vanadium, $`T_c`$ from the SCDFT gap equation is about 2 K higher than the experimental data expTc , even after inclusion of spin fluctuations. This fact may indicate that, either the spin exchange constant, $`I_{ex}`$, was underestimated, or a contribution of the asymmetric part of the phononic term in the SCDFT gap equation is quite large. As we know from results reported in Refs. Martin-OGK ; our-Nb , if we neglect the asymmetric part in the electron-phonon-coupling matrix elements by taking the $`\alpha ^2F(\omega )`$ avaraged at the Fermi level, the critical temperatures are higher (see the discussion in Section 6). The last approximation, however, has to be done if we do not evaluate formulas with the $`g_{๐ค,๐ค+๐ช}`$ elements explicitely.
In general, the critical temperatures obtained from the SCDFT are in a good agreement with the measured temperatures expTc , and the effect of paramagnons improves the result considerable for many simple metals.
### 7.2 Niobium under pressure
In TABLE 2, we present critical temperatures and parameters of the gap equation for niobium at eight pressures in the range from -17 GPa up to 80 GPa. The spin exchange constants, $`I_{ex}`$, have been calculated from first principles in Ref. our-Nb , and the electron-phonon and electron-magnon spectral functions for Nb have been presented also in that work.
Here, we complete our previous results by reporting the effect of paramagnons on $`T_c`$ calculated from the SCDFT. After the inclusion of spin fluctuations, the critical temperatures obtained from the SCDFT are closer to the experimental $`T_c`$โs for pressures in the range of 0-40 GPa, i.e. pressures between two anomalies measured by Struzhkin et al. Struzhkin . The dependence of the measured critical temperature as a function of pressure is no longer reproduced by our calculations when we take into account paramagnons. At ambient pressure and for higher pressures, paramagnons seem to make the theoretical result worse. The above effect, could be explained by making the observation that, in every case where the exchange constant $`I_{ex}`$ is large, the theoretical temperature underestimates the measured temperature, and vice versa, for the smallest $`I_{ex}`$ the critical temperature obtained from the SCDFT is the highest and the error is positive.
Concluding this Subsection, the implementation of paramagnons to the SCDFT generally makes calculated critical temperatures closer to the experimental ones. But our calculated exchange constants, $`I_{ex}`$, are not sufficiently accurate. This fact gives a direction for the future development.
## 8 Summary
In the present work, we included the transverse spin fluctuations to the density functional theory for superconductors. The SCDFT is presented from its foundations, through the decoupling approximation, the gap equation and details of implementation. We assumed singlet and $`s`$-wave pairing potential; The extension to triplet superconductors could be done following the work by Capelle et al. KC ; triplet . The electron-phonon couplings and the electron-paramagnon couplings were averaged at the Fermi energy, therefore the asymmetric part of the functional with respect to the electron-phonon matrix elements and to the spin-exchange interaction constants were ommited. Through the whole work, we kept the notation to be consistent with Parks Parks ; Kazumi and Vonsovsky Vonsovsky .
Paramagnons and phonons in the superconducting state were assumed to be the same like in the normal state. The Eliashberg spectral function has been calculated within the density functional perturbation theory and it is fully material specific. Paramagnons, in contrast, have been obtained from the random phase approximation for the homogeneous electron gas and only the spin exchange constants were calculated from the electronic structure.
We reported the critical temperatures obtained from the SCDFT and the Eliashberg linearized gap equation with and without spin fluctuations for a few simple metals: V, Mo, Ta, Pd at ambient pressure and Nb at several pressures up to 80 GPa. Some discrepancies between the temperatures calculated from the SCDFT and the measured temperatures are due to the fact that it is quite difficult to obtain the accurate spin-exchange constants and/or to the fact that the spectral functions have been averaged at the Fermi level. Netherveless, the results show that inclusion of paramagnons improves the critical temperatures obtained from both methods, the SCDFT and the Eliashberg theory. The critical temperatures obtained from the parameter-free SCDFT are in most cases closer to the experimental data than the results obtained from the Eliashberg theory.
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# Energy density bounds for black strings
## I Introduction
One of most important and interesting properties of black strings is the classical instability which was discovered by Gregory and Laflamme GL . However, at this stage, we have not obtain the clear answers for what is the final state of the instability. However, the positivity of the energy of the spacetimes, that is, whether there is the lower bound of the energy of the spacetimes, or not, might give us the information about the stability of them or the final states. For an example, the Schwarzschild black hole with negative mass is unstable under the metric perturbation, and if the positivity of the energy holds, then the spacetimes with negative mass would have naked singularity somewhere. A sufficient condition for the stability of spacetimes is that its energy saturates a lower bound on the energy of all field configurations satisfying the same boundary conditions. It is known that some charged extreme black hole solutions saturate a lower bound such a bound, which is the Bogomolโ nyi-type bound GH .
In a class of four dimensional asymptotically flat spacetimes, the trial to show the positivity of the ADM energy had been made of by a lot of authors. However, they proved it in the only special cases. Finally, Shoen and Yau SY succeeded in giving a complete proof on the positivity of the ADM energy in asymptotically flat spacetimes, which can be applied whether there is a black hole or not. The more elegant and simpler proof was given by Witten and Nester W ; N , who used the spinorial method, when there is no horizon. The introduced spinor is the solution of Dirac-Witten equation and is asymptotically constant at spatial infinity. Gibbons et. al. G extended Wittenโs arguments to black hole spacetimes containing an apparent horizon without assuming anything about the interior region of a black hole. They integrated the Dirac-Witten equation over a sapcelike hypersurface with the boundaries the spatial inifinity and the horizon and applied Stokesโs theorem there. The term over the apparent horizon as an inner boundary arise in addition to the surface term over the spatial infinity, which coincides with the ADM 4-momentum. Since this spinor boundary term over the apparent horizon becomes proportional to the expansion of null geodesics on the apparent horizon, it vanishes. As a result, the discussion of Gibbons et. al. becomes equal to that of the original Wittenโs and the ADM mass also becomes positive in the presence of a black hole.
The ADM mass in a $`p`$-brane spacetime or a black string will diverge because the spatial $`p`$ directions or the direction of the black sting are infinite in extent and it is given by an integral over the spatial infinity (with its topology $`R^p\times S^{D(p+2)}`$) which encloses the $`p`$-brane. Therefore, in such a spacetime the ADM mass is not well-defined conserved charge in general except for the cases where they are compactified. The proof of positivity of such compactified cases have been already given SIT . Through this paper, we consider the conserved charges from asymptotic Killing-Yano tensors, which was introduced by Kastor and Traschen KT . This charge is called Y-ADM charge and gives us good physical interpretation of the mass density rather than that of a total mass of the spacetime. In the previous work YADM , we investigated the positivity of Y-ADM mass density in $`1`$ brane spacetimes, which are assume to be transverse asymptotically flat to the $`1`$ brane. As results, we could establish the positivity in the two special cases ; (a) conformally flat cases or algebraically special spacetimes and (b) cases where there exists translational Killing vector field along the string.
In this paper, we will investigate the positivity of the Y-ADM mass density in transversely asymptotically flat black string spacetimes containing an apparent horizon, using Witten-Nesterโs spinorial technique. Due to the presence of an horizon, the spinor boundary term over the cross section of the horizon and the transverse space arises, as is seen in the proof of Gibbons et. al. G . If this boundary term vanishes, we can establish the positivity of the Y-ADM mass density in the above cases (a) and (b) from the results in the previous studies YADM . We will show the sufficient conditions for the horizon surface term to vanish, which is the cases where (A) the apparent horizon becomes null. Therefore, if at least, the black string spacetimes satisfy the conditions either (a), or (b) and (A), then the positivity of Y-ADM mass density is assured.
This paper is organized as follows. In section II, we will show that the spinorial boundary term over the cross section of a black string and the transverse space becomes proportional to the expansion of the cross section rather than the horizon, unlike the proof of Gibbons et. al. G . In section III, we give the sufficient conditions for this surface term to vanish on the apparent horizon. In section LABEL:sec:sum, we will summarize the results.
## II Positivity bounds of Y-ADM mass density
We assume that the black string spacetime $`(M,g_{ab})`$ containing an apparent horizon is $`D(5)`$ dimensions and transverse asymptotically flat, which means asymptotically flat in the direction transverse to the black string but always not in the direction parallel to it. We also assume that the spacetime is foliated by time slices $`V_\mathrm{t}`$ with normal timelike vector field $`t^a`$, and that each timeslice $`V_\mathrm{t}`$ is foliated by the submanifolds $`V_{\mathrm{tx}}`$ binormal to vectors $`t^a`$ and $`x^a`$, where $`x^a`$ is the vector field pallarel to the black string, satisfying the normalization $`x^ax_a=1`$ and $`t^ax_a=0`$. Let us assume the codimension two surface $`V_{\mathrm{tx}}`$ which is the transverse space to the black sting) is foliated by the codimension three surface $`V_{\mathrm{txr}}`$ normal to the transverse radial directional vector field $`r^a`$ satisfying the $`r^ar_a=1`$ and $`r^at_a=r^ax_a=0`$. At transverse spatial infinity we take $`t^a=(/x^0)^a,x^a=(/x^1)^a`$. In this spacetime, there exists the asymptotic Killing-Yano tensor $`f^{(01)}=dx^0dx^1`$, by which we can define the Y-ADM mass density as the surface integral over the transverse spatial infinity $`S_{\mathrm{}}`$ KT ; YADM , as seen in FIG.1.
The full spacetime metric can be written in the form
$`g_{ab}`$ $`=`$ $`t_at_b+h_{ab}`$ (1)
$`=`$ $`t_at_b+x_ax_b+q_{ab}`$
$`=`$ $`t_at_b+x_ax_b+r_ar_b+r_{ab},`$
where $`h_{ab}`$, $`q_{ab}`$ and $`r_{ab}`$ are the induced metrics on the submanifolds $`V_\mathrm{t}`$, $`V_{\mathrm{tx}}`$ and $`V_{\mathrm{txr}}`$, respectively. In this paper, we use the quantities $`\gamma ^{\widehat{0}}=\gamma ^at_a,\gamma ^{\widehat{1}}=\gamma ^ax_a,\gamma ^{\widehat{2}}=\gamma ^ar_a`$.
To prove the positivity of the ADM energy, Witten W used the spinorial method, which consists of the two steps. As the first step, he showed that the surface integral of so-called Nester 2 form $`E^{ab}=\psi ^{}\gamma ^{\widehat{0}}\gamma ^{abc}_c\psi `$ over the spatial infinity (codimension two surface) coincides with the ADM mass (exactly speaking, ADM $`4`$-momentum). The following step is to transform this surface integral into the volume integral, using Stokesโ theorem, and to relate it to Eintein tensor equal to the energy momentum tensor satisfying the dominant energy condition. To discuss the positivity of Y-ADM mass density, let us introduce the Nester $`3`$-form YADM defined as
$`B^{abc}=\psi ^{}\gamma ^{\widehat{0}\widehat{1}}\gamma ^{abcd}_d\psi `$ (2)
in analogy with the Nester $`2`$ form in the proof of the positive energy theorem W , where $`\psi `$ is asymptotically constant spinor in the transverse direction and satisfies the Dirac-Witten equation $`q_b^a\gamma ^b_a\psi =0`$. However, in the reference YADM , we dealt with the transversely asymptotically flat $`1`$ brane spacetimes with no horizon since the in general, in the case with a horizon there exists singularity in the interior region and so the discussion in the reference YADM cannot be applied to the spacetimes with horizons, such as, black string. In the latter case, we need add the surface integral over the cross section $``$ of horizon and transverse space as the inner boundary to the surface integral over the transverse spatial infinity. To explain this, let us begin with the volume integral over the transverse space $`V_{\mathrm{tx}}`$ (See the equation (19) in the reference YADM ),
$`{\displaystyle \frac{1}{8\pi }}{\displaystyle _{V_{tx}}}dS_{bc}(_aB^{abc}+_aB^{abc})={\displaystyle \frac{1}{8\pi }}{\displaystyle _{V_{\mathrm{tx}}}}[(G_{ab}t^at^b`$ (3)
$``$ $`R_{abcd}x^ax^cq^{bd})\psi ^{}\psi +2(_a\psi ^{})q^{ab}(_b\psi )].`$
If we choose the cross section of the horizon and the transverse space as an inner boundary in analogy with the reference G , then as the result of applying Stokesโ theorem to the left hand side in the above eqaution, we obtain the following eqaution,
$`{\displaystyle \frac{1}{8\pi }}{\displaystyle _S_{\mathrm{}}}๐S_{abc}(B^{abc}+B^{abc})`$ $``$ $`{\displaystyle \frac{1}{8\pi }}{\displaystyle _{}}๐S_{abc}(B^{abc}+B^{abc})`$ (4)
$`=`$ $`{\displaystyle \frac{1}{8\pi }}{\displaystyle _{V_{\mathrm{tx}}}}[(G_{ab}t^at^bR_{abcd}x^ax^cq^{bd})\psi ^{}\psi +2(_a\psi ^{})q^{ab}(_b\psi )],`$
where the first term in the left hand side coincides with the Y-ADM mass density $``$, following the discussion in the reference YADM . On the other hand, the second term in the left hand side appears due to the presence of the horizon $``$, which implies the cross section of the horizon and the transverse space. The right hand side becomes positive in at least, two special cases, (1) conformally flat cases or algebraically special spacetimes and (2) cases where there exists translational Killing vector field along the string, as shown in YADM .
The purpose of this paper is to estimate the surface term on the horizon $``$ and to investigate when this surface term vanishes. This term is
$`{\displaystyle \frac{1}{8\pi }}{\displaystyle _{}}๐S_{abc}(B^{abc}+B^{abc})={\displaystyle \frac{1}{8\pi }}{\displaystyle _{}}๐S(B^{\widehat{0}\widehat{1}\widehat{2}}+B^{\widehat{0}\widehat{1}\widehat{2}}),`$ (5)
where $`B^{\widehat{0}\widehat{1}\widehat{2}}=\psi ^{}\gamma ^{\widehat{0}}\gamma ^{\widehat{1}}\gamma ^{\widehat{0}\widehat{1}\widehat{2}}\gamma ^{\widehat{A}}_{\widehat{A}}\psi =\psi ^{}\gamma ^{\widehat{2}}\gamma ^{\widehat{A}}_{\widehat{A}}\psi `$ and $`dS`$ is the volume element on the cross section of the horizon and the transverse space.
We have the following relation between the projection into $`V_t`$, $`V_{tx}`$ and $`V_{txr}`$ of the $`D`$ dimensional spinor covariant derivatives $`_a`$ and the intrinsic $`(D1`$) dimensional, $`(D2)`$ dimensional, and $`(D3)`$ dimensional spinor covariant derivatives $`{}_{}{}^{(D1)}_{\widehat{A}}^{},{}_{}{}^{(D2)}_{\widehat{A}}^{}`$, and $`{}_{}{}^{(D3)}_{\widehat{A}}^{}`$,
$`_{\widehat{A}}\psi `$ $`=`$ $`{}_{}{}^{(D1)}_{\widehat{A}}^{}\psi +{\displaystyle \frac{1}{2}}K_{A\widehat{i}}\gamma ^{\widehat{i}}\gamma ^{\widehat{0}}\psi `$ (6)
$`=`$ $`{}_{}{}^{(D2)}_{\widehat{A}}^{}\psi +{\displaystyle \frac{1}{2}}J_{\widehat{A}\widehat{n}}\gamma ^{\widehat{n}}\gamma ^{\widehat{1}}\psi +{\displaystyle \frac{1}{2}}K_{A\widehat{i}}\gamma ^{\widehat{i}}\gamma ^{\widehat{0}}\psi `$
$`=`$ $`{}_{}{}^{(D3)}_{\widehat{A}}^{}\psi +{\displaystyle \frac{1}{2}}k_{\widehat{A}\widehat{B}}\gamma ^{\widehat{B}}\gamma ^{\widehat{2}}\psi +{\displaystyle \frac{1}{2}}J_{\widehat{A}\widehat{n}}\gamma ^{\widehat{n}}\gamma ^{\widehat{1}}\psi `$
$`+{\displaystyle \frac{1}{2}}K_{A\widehat{i}}\gamma ^{\widehat{i}}\gamma ^{\widehat{0}}\psi ,`$
respectively, where $`K_{ab}`$, $`J_{ab}`$ and $`k_{ab}`$ are the extrinsic curvatures of the submanifolds $`V_\mathrm{t}`$, $`V_{\mathrm{tx}}`$ and $`V_{\mathrm{txr}}`$, respectively. The indexes $`\widehat{i},\widehat{n}`$ and $`\widehat{A}`$ run $`\widehat{i}=\widehat{1},\mathrm{},\widehat{D}1`$, $`\widehat{n}=\widehat{2},\mathrm{},D1`$ and $`\widehat{A}=\widehat{3},\mathrm{},\widehat{,}D1`$. Therefore, the real part of $`B^{abc}t_{[a}x_br_{c]}=B^{\widehat{0}\widehat{1}\widehat{2}}`$ becomes
$`B^{\widehat{0}\widehat{1}\widehat{2}}+B^{\widehat{0}\widehat{1}\widehat{2}}=\psi ^{}\gamma ^{\widehat{2}}\gamma ^{\widehat{A}}{}_{}{}^{(D3)}_{\widehat{A}}^{}\psi +K_{\widehat{A}\widehat{2}}\psi ^{}\gamma ^{\widehat{0}}\gamma ^{\widehat{A}}\psi +(k_{\widehat{A}}^{\widehat{A}}K_{\widehat{1}}^{\widehat{1}}K_{\widehat{2}}^{\widehat{2}})\psi ^{}\psi +K_{\widehat{i}}^{\widehat{i}}\psi ^{}\gamma ^{\widehat{2}}\gamma ^{\widehat{0}}\psi +\mathrm{c}.\mathrm{c}`$ (7)
Let us impose the boundary condition with respect to the spinor on the cross section $``$ of the horizon as follows G ,
$`\gamma ^{\widehat{2}}\gamma ^{\widehat{0}}\psi =\psi ,`$ (8)
where we used the fact that the eigenvalues of the Hermitian matrix $`\gamma ^{\widehat{2}}\gamma ^{\widehat{0}}`$ are $`\pm 1`$ and therefore, in the above, we restrict the freedom of the spinor to the half. Under this condition, we find that the first term in the left hand side in the equation (7) vanishes, since
$`\psi ^{}\gamma ^{\widehat{2}}\gamma ^{\widehat{A}}{}_{}{}^{(D3)}_{\widehat{A}}^{}\psi `$ $`=`$ $`(\gamma ^{\widehat{2}}\gamma ^{\widehat{0}}\psi )^{}\gamma ^{\widehat{2}}\gamma ^{\widehat{A}}{}_{}{}^{(D3)}_{\widehat{A}}^{}\psi `$ (9)
$`=`$ $`\psi ^{}\gamma ^{\widehat{0}}\gamma ^{\widehat{2}}\gamma ^{\widehat{2}}\gamma ^{\widehat{A}}{}_{}{}^{(D3)}_{\widehat{A}}^{}\psi `$
$`=`$ $`\psi ^{}\gamma ^{\widehat{2}}\gamma ^{\widehat{A}}{}_{}{}^{(D3)}_{\widehat{A}}^{}(\gamma ^{\widehat{2}}\gamma ^{\widehat{0}}\psi )`$
$`=`$ $`\psi ^{}\gamma ^{\widehat{2}}\gamma ^{\widehat{A}}{}_{}{}^{(D3)}_{\widehat{A}}^{}\psi `$
$`=`$ $`0.`$
Similarly, the second term in the left hand side in the equation (7) also vanishes, because $`\gamma ^{\widehat{0}}\gamma ^{\widehat{A}}`$ anticommute with $`\gamma ^{\widehat{2}}\gamma ^{\widehat{0}}`$, that is,
$`\psi ^{}\gamma ^{\widehat{0}}\gamma ^{\widehat{A}}\psi =\psi ^{}\gamma ^{\widehat{0}}\gamma ^{\widehat{A}}\psi =0.`$ (10)
Then we obtain the simple expression with respect to the suface integral over the horizon $``$ as follows,
$`{\displaystyle \frac{1}{8\pi }}{\displaystyle _{}}๐S_{abc}(B^{abc}+B^{abc})`$ (11)
$`=`$ $`{\displaystyle \frac{1}{8\pi }}{\displaystyle _{}}[\psi ^{}(k_{\widehat{A}}^{\widehat{A}}K_{\widehat{1}}^{\widehat{1}}K_{\widehat{2}}^{\widehat{2}}+K_{\widehat{i}}^{\widehat{i}})\psi ],`$
which we find that this integrand coincides with the expansion $`\theta _+`$ of the cross section and black string and the transverse space, since we can observe that the above term coincides with $`\theta _+`$ defined as
$`\theta _+`$ $`:=`$ $`r^{ab}_al_b`$ (12)
$`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}(K_{\widehat{i}}^{\widehat{i}}K_{\widehat{1}}^{\widehat{1}}K_{\widehat{2}}^{\widehat{2}}+k_{\widehat{A}}^{\widehat{A}}),`$
where $`l^a=\frac{1}{\sqrt{2}}(t^a+r^a)`$ is outgoing null vector. Therefore, we observe that when this expansion of the cross section $``$ of the horizon and the transverse space $`V_{\mathrm{tx}}`$, the surface term on the cross section of the horizon vanishes. We should note that unlike the proof in the references G ; SIT , in general, this term does not vanishes on the apparent horizon, since $`\theta _+`$ is not expansion of horizon but that of the cross section of the horizon and the transverse space.
## III Suface term over a horizon
As mentioned in the previous section, the surface integral over the cross section of the horizon and the transverse space generally does not vanish on an apparent horizon unlike the proof of Gibbons et. al G . As seen in the below, a sufficient for this surface term to vanish is the cases where the apparent horizon becomes null surface. In this section, we will show that the above surface term vanishes on an apparent horizon in the above conditions.
### III.1 Cases where an apparent horizon is null
Following Haywardโs discussion Hay , an apparent horizon is null if and only if null energy condition hold, the shear and the normal energy density vanishes Hay (exactly speaking, this is a necessary and sufficient condition for a so-called trapping horizon defined by the closure of a $`D1`$ dimensional surface foliated by marginal surfaces on which $`\mathrm{\Theta }_+=0,\mathrm{\Theta }_{}0`$ and $`_{}\mathrm{\Theta }_+0`$, where $`_{}`$ denotes the Lie derivative along the ingoing null direction and $`\mathrm{\Theta }_{}`$ is the expansion of the ingoing null geodesics, becomes null surface. In this paper, the apparent horizon means this trapping horizon. ). We can show that if an apparent horizon is null surface, the above surface term vanishes. The expansion of out going null geodesics on the apparent horizon is defied as
$`\mathrm{\Theta }_+=s^{ab}_al_b=x^ax^b_al_b+\theta _+=0`$ (13)
On the other hand, the shear of the horizon can be expressed in the form
$`\sigma _{ab}`$ $`=`$ $`s_{(a}^cs_{b)}^d_al_b{\displaystyle \frac{1}{D2}}\mathrm{\Theta }_+s_{ab}`$ (14)
$`=`$ $`(x^cx^d_al_b)x_ax_b+(x^c_cl_d)x_{(a}q_{b)}^d`$
$`+q_{(a}^cq_{b)}^d_cl_d`$
$`=0,`$
whrere $`s_{ab}=g_{ab}+t_at_br_ar_b`$ and the equation (13) was used. Therefore, since $`\sigma _{xx}:=\sigma _{ab}x^ax^b=x^cx^d_cl_d=0`$, we can show that on the apparent horizon $`\theta _+=0`$ in combination with the equation (13).
We will mention when the physical situations where an apparent horizon becomes null would realize. In stationary black string spacetimes, the apparent horizon would coincide with the Killing horizon of the black string. However, even in non-stationary cases, such situation would occur, if there exists the apparent horizon and there is no radiation across and along the horizon. Here we should note that in non-stationary cases, since the apparent horizon is generally inside the horizon, we have to replace $``$ with the cross section $`๐`$ of this apparent horizon $`A`$ and the transverse as the inter boundary in the equation (3) in order that the inner boundary vanishes.
## IV Summary
Finally, in combination with the previous work, we will mention the sufficient conditions for the positivity of the Y-ADM mass density to hold black strings whose boundary condition satisfie transverse asymptotically flatness.
The Y-ADM mass density becomes positive if a transversely asymptotically flat black string spacetime which contains an apparent horizon and satisfies the dominant energy condition follows the below conditions ;
(1) either conformally flat cases or algebraically special spacetimes, or there is a translationary Killing field along the black string, i.e. the black string is uniform and (2) the apparent horizon becomes null.
## Acknowledgements
We thank Associated Professor K. Nakao and Professor H. Ishihara for useful comments. We would like to thank Associated Professor T. Shiromizu for continuous encouragement.
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# INTRODUCTION TO DARK ENERGY AND DARK MATTER
## 1 Plan of Introduction to Dark Energy
* What observations and theoretical assumptions underly dark energy (DE)?
* If general relativity (GR) holds at all length scales, the most conservative assumption, then DE follows from the supernovae Type 1A (SNe1A) or, independently, from the Cosmic Microwave Background (CMB) combined with Large Scale Structure (LSS).
* Should we seriously query GR at large distance scales?
## 2 Einstein-Friedmann Equation
The Einstein equations relate geometry on the Left-Hand-Side (LHS) to the distribution of mass-energy on the Right-Hand-Side (RHS)
$$G_{\mu \nu }=8\pi GT_{\mu \nu }$$
(1)
We hesitate to change the LHS but it is really checked with precision only at Solar System (SS) scales. At cosmological length scales, we may consider using a modification such as higher-dimensional gravity.
On the RHS, if we include only luminous and dark matter it is insufficient (keeping the LHS intact) and there is needed a further term which could be a cosmological constant or, more generally, dark energy.
## 3 Observational Issues
How can we constrain DE?
* Measurement of the expansion history H(t)
* The time-dependence of the equation of state w(t)
* Looking for any clustering property of DE. No evidence for this presently.
* How does DE couple to Dark Matter (DM)? This is related to the question of clustering.
* Local tests of GR and the equivalence principle, though the extrapolation from the SS to the Universe is some 13-15 orders of magnitude comparable to the extrapolation from the weak scale to the GUT scale in particle phenomenology. The usual prior is a desert hypothesis.
## 4 $`\mathrm{\Lambda }`$ as DE: Why $`10^{122}`$ (Planck Mass)<sup>4</sup>?
We know from the Lamb Shift and Casimir Effect in quantum electrodynamics that vacuum fluctuations are real effects.
If we calculate the value of $`\mathrm{\Lambda }`$, it will naively be ultra-violet (UV) quartically divergent. The most natural UV cut-off in GR is the Planck mass $`10^{19}GeV`$ whereupon
$$\mathrm{\Lambda }(10^{19}GeV)^4=(10^{28}eV)^4=10^{112}(eV)^4$$
(2)
If we use, instead, the weak scale $`100GeV`$ as our UV cut-off, we arrive at
$$\mathrm{\Lambda }(100GeV)^4=(10^{11}eV)^4=10^{44}(eV)^4$$
(3)
The observed value for $`\mathrm{\Lambda }`$, by contrast, is approximately
$$\mathrm{\Lambda }(3\times 10^3eV)^410^{10}(eV)^4$$
(4)
## 5 Coincidence Problem
As if the fine-tuning problem for $`\mathrm{\Lambda }`$ were not enough, there is a second problem with $`\mathrm{\Lambda }`$, the coincidence problem. Let us define $`\mathrm{\Omega }_\mathrm{\Lambda }=\rho _\mathrm{\Lambda }/\rho _C`$ as the fraction of the critical density $`\rho _C`$.
The present value is $`\mathrm{\Omega }_\mathrm{\Lambda }0.7`$ but it scales, since $`\rho _\mathrm{\Lambda }`$ is constant and assuming $`\mathrm{\Omega }_{TOT}=1`$, like $`\rho _C^1(1+Z)^3`$ so at a redshift $`Z>10`$ it was $`\mathrm{\Omega }_\mathrm{\Lambda }<0.001`$ while for a future redshift $`Z<0.9`$ one has $`\mathrm{\Omega }_\mathrm{\Lambda }>0.999`$.
If we plot $`\mathrm{\Omega }_\mathrm{\Lambda }`$ versus log R over cosmic history from $`60<log_{10}R<+60`$, it appears like a step function changing from zero to one abruptly around the present era. Even more dramatic is a plot of $`d\mathrm{\Omega }_\mathrm{\Lambda }/dR`$ which approximates a Dirac delta function and the coincidence problem is then why we live right in the middle of the spike of the delta function.
If the dark energy had appeared earlier it would have interfered with structure formation: if later, we would still be unaware of it.
## 6 The Quintessence Possibility
One parametrization of the dark energy can be made using a dynamical scalar field, now generically called quintessence.
### 6.1 Scaling potentials
Examples are:
$$Ve^{\lambda \mathrm{\Phi }}$$
(5)
as in $`^{\mathrm{?},\mathrm{?}}`$
$$V((\mathrm{\Phi }A)^2+C)e^{\lambda \mathrm{\Phi }}$$
(6)
as in $`^\mathrm{?}`$.
### 6.2 Tracker Potentials
Examples are
$$V\mathrm{\Phi }^\alpha $$
(7)
as in $`^\mathrm{?}`$,
$$Vexp\left(\frac{M}{Q}1\right)$$
(8)
as in $`^\mathrm{?}`$.
### 6.3 Approaches to the Coincidence Problem
We may assume that our universe sees periodic epochs of acceleration$`^\mathrm{?}`$ with potential
$$VM^4e^{\lambda \mathrm{\Phi }}(1+A\mathrm{sin}a\mathrm{\Phi })$$
(9)
Another possibility is that it is important that our epoch is close to the matter/radiation equality time. This may be incorporated by having a non-minimal coupling to matter$`^\mathrm{?}`$, to gravity$`^\mathrm{?}`$ or in a k-essence theory with a non-trivial kinetic term in the lagrangian$`^\mathrm{?}`$.
## 7 Dark Energy with Equation of State $`w=p/\rho <1`$
Present data on SNe1A, CMB and LSS are consistent with w=-1 as for a cosmological constant.
Since the possibility that $`w<1`$ is still allowed$`^\mathrm{?}`$, I shall spend a disproportionate amount of time on it because, if it persisted, it could well signal new physics.
One interpretation of dark energy comes from string theory, closed strings on a toroidal cosmology$`^\mathrm{?}`$. This leads generically to $`w<1`$ $`^\mathrm{?}`$.
In general, without dark energy (as in most cosmology texts pre-1998), the destiny of the Universe was tied to geometry in a simple manner: the Universe will expand forever if it is open or flat; it will stop expanding and contract to a Big Crunch if it is closed.
With Dark Energy, this connection between geometry and destiny is lost and the future fate depends entirely on how the presently-dominant dark energy will evolve.
This question is studied in $`^{\mathrm{?},\mathrm{?},\mathrm{?}}`$. If $`w<1`$ and is time-independent, the scale factor diverges at a finite future time \- the Big Rip. Generally, this will be at least as far in the future as the Big Bang was in the past.
Such a cosmology may have philosophical appeal? There is more symmetry between past and future.
If one allows a time-dependent w(t), there are two other possible fates:
(i)An infinite-lifetime universe where dark energy dominates at all future times.
(ii)A disappearing dark energy where the Universe becomes (again) matter dominated.
The case $`w<1`$ gives rise to some exceptionally interesting puzzles for theoretical physics.
There is the question of violation of the weak energy condition universally assumed in general relativity. This means there are inertial frames where the energy density is negative signaling vacuum instability$`^{\mathrm{?},\mathrm{?}}`$.
Let us make three assumptions, any or all of which may be incorrect, just so that we may say something more: that (i) There is a stable vacuum with $`\mathrm{\Lambda }=0`$; (ii) The dark energy decays to it by a 1st-order phase transition; (iii) There is some, albeit feeble, interaction between dark energy and the electromagnetic field.
Then one can use old arguments$`^\mathrm{?}`$ to investigate nucleation. The result is that$`^\mathrm{?}`$ even with the tiniest coupling of dark energy to the electromagnetic field the dark energy would have spontaneously decayed long ago unless the appropriate bubble radius is at least galactic in size.
In this model, because the energy density of the DE is so small compared to e.g. the energy density in a common macroscopic magnetic field of, say, 10T the 1st order phase transition can be adequately suppressed only by decoupling the DE completely from all but gravitational forces or by arguing that a collision would need to be between galaxies or larger objects to be effected. Certainly, no terrestrial experiment can be influenced: for one contrary suggestion of a Josephson junction experiment which might well be justified for other reasons, see e.g. $`^\mathrm{?}`$.
Of course, this is only a toy model but the general conclusion is probably correct - that there can be no microscopic effect of the dark energy.
This makes the DE very difficult or impossible to investigate except through astronomical observations.
## 8 Dark Energy and Neutrinos
It has been pointed out by many theorists that the density of the dark energy $`(10^3eV)^4`$ is suggestive of the neutrino mass.
Very interesting attempts to strengthen such a connection have been made $`^{\mathrm{?},\mathrm{?}}`$. Such mass-varying neutrino models seek to make a direct identification of the DE density with neutrino mass$`^{\mathrm{?},\mathrm{?}}`$ itself.
## 9 Precision Experiment
We know well of the precision experiments to test Newtonโs Law of Gravity down to a distance of 100 microns and below.
One originator of such ideas suggests$`^\mathrm{?}`$ a different precision test, of the Earth-Moon distance, to a similar accuracy of 100 microns, presumably the distance between the centers of mass. A particular modification of gravity$`^\mathrm{?}`$ might have a tiny effect on our lunar system. Clearly if this experiment can be achieved, the present accuracy being at the level of centimeters, it would be an impressive achievement.
## 10 Conclusions on Dark Energy
* The theoretical community has yet to come up with a definitive proposal to explain the dark energy.
* The nature of the dark energy is so profound for cosmology and particle physics that we desperately need more SNe1A observations from important proposed experiments e.g. SNAP (for which NASA funding has sadly been suspended for 5 years as a result of prioritizing sending humans to Mars!), as well as complementary observational constraints on the CMB from e.g. the Planck mission.
* The equation of state will be decisive. If w=-1, itโs a cosmological constant with its fine-tuning and coincidence problems. If $`w>1`$ quintessence will receive a shot in the arm.
* If the data would settle down to a value $`w<1`$ we could be at the dawn of a revolution in theory with general relativity at the largest distance scales called into question.
## 11 Introduction to Dark Matter
.
Existence of darl matter is supported by disparate cosmological measurements.
Values of energy and matter densities at the present time, determined by: the temperature fluctuations in the CMB data; distance-luminosity for supernovae type 1A; distribution of galaxies on large scales (LSS); abundance of light elements (BBN).
In terms of the critical density $`\mathrm{\Omega }`$ for the various components is found to be as follows (taking $`h^2=0.5,h=0.707`$).
* Relativistic particles, radiation e.g. the CMB photons. Only $`\mathrm{\Omega }_\gamma =5.934\pm 0.008\times 10^5`$.
* $`\mathrm{\Omega }_\mathrm{\Lambda }=0.72\pm 0.08`$ in a smoothly distributed dark energy.
* $`\mathrm{\Omega }_M=0.27\pm 0.016`$ in non-relativistic particles (Matter) of which
$`\mathrm{\Omega }_b=0.0448\pm 0.0018`$ in baryons (protons and neutrons)
$`\mathrm{\Omega }_{HDM}<0.0152(95\%CL)`$ in non-baryonic hot dark matter.
$`\mathrm{\Omega }_{CDM}=0.223\pm 0.016`$ in non-baryonic cold dark matter.
The excess of total matter density ($`0.27`$) over baryonic mass density ($`0.0224`$) constitutes the evidence for non-baryonic dark matter.
No known elementary particle can account for the non-baryonic dark matter.
One obvious candidate, the neutrinos, are so light they constitute hot dark matter and contribute to the $`\mathrm{\Omega }_{HDM}<0.0152`$.
Many hypothetical particles have been proposed for the CDM. Some come from extensions of the standard model, most notably the axion and the lightest supersymmetric particle.
Other possibilities include Wimpzillas, solitons, self-interacting dark matter, Kaluza-Klein dark matter, etc.
The class of non-baryonic dark matter candidate of greatest interest are the Weakly Interacting Massive Particles (WIMPs). Therefore I shall focus on their detection and the claims to have discovered them.
## 12 WIMPs and their detection
WIMPs are appealing because of the simple mechanism by which they can achieve the appropriate present cosmic density. In the early universe they were in thermal and chemical equilibrium with the rest of matter and radiation. With the expansion of the Universe, their reactions (including annihilation) slowed down and decoupled from the rest of the world leaving a constant number of WIMPs expanding with the Universe.
The correct present density is obtained for WIMPs with couplings of order the weak interactions and masses in the 1 GeV - 1 TeV range. The neutralino$`^\mathrm{?}`$ is the most popular example although in any extension of the standard model one typical seeks a WIMP candidate, e.g. the nark is a WIMP candidate$`^\mathrm{?}`$ in the sark model$`^\mathrm{?}`$.
Detection can be direct or indirect.
Direct signals are from collisions with nuclei in a detector. A very sensitive low-background detector (bolometer) records the amount of energy deposited by WIMPs and (in the future) the direction of motion of the struck nucleus.
Indirect signals come from WIMP reactions in planets, stars or galaxies. The most common reaction is WIMP annihilation with anti-WIMP. Out of this annihilation come $`\nu `$, $`e^+`$, $`\overline{p}`$ and high-energy $`\gamma `$. Such annihilations occur at a detectable rate where the anti-WIMPs are concentrated e.g. in the center of the Sun, the center of the Earth and in galactic centers including the Milky Way. Neutrino telescopes, gamma-ray telescopes and cosmic-ray detectors can be used in these indirect searches.
Now we look at three claims for seeing WIMPs.
## 13 HEAT positron detection
Two separate balloon flights with different detectors have seen$`^\mathrm{?}`$ more cosmic ray positrons above $`7`$ GeV than predicted in models for cosmic ray propagation in the galaxy. Wimp annihilation can be invoked to explain this excess.
The extra positrons can be fitted by a assuming neutralino annihilation. The best fit requires WIMP mass of 238 GeV.
The positron spectrum lacks any discriminating feature which clearly singles out WIMP annihilation.
## 14 $`\gamma `$-Rays from Galactic Bulge
Gamma rays from WIMP annihilation offer a characteristic signature in the spectrum: a gamma-ray line. Each photon will carry an energy equal to the WIMP mass, 10GeV to 100 TeV. No competing process is know that could produce such a line.
No line has been detected yet. The estimates for GLAST (launch scheduled 2006) are encouraging.
Another suggestion$`^\mathrm{?}`$ has been that the 511 keV gamma excess from the galactic bulge arises from positrons associated with unexpectedly light WIMP annihilation. The necessary WIMP mass is in the region between 1 MeV and 100 MeV. There are other explanations for the 511 KeV line including primordial black holes as dark matter$`^\mathrm{?}`$.
## 15 DAMA Annular Modulation
Because the Earthโs motion changes the relative speed of the Earth and WIMPs the WIMP detection rate varies$`^\mathrm{?}`$ and repeats itself every year. The maximum occurs in June for the canonical halo model with Maxwellian velocity distribution.
The DAMA group has claimed$`^\mathrm{?}`$ to have detected such annual modulation in their NaI data. No alternative explanation of the DAMA data has been forthcoming.
No other direct detection of such a WIMP signal has been made by any other group but there are differences between the targets used as well as the nuclear spin thereof. So comparison between experiments requires some theoretical assumptions.
Nevertheless, it does appear that CDMS data$`^\mathrm{?}`$ completely or almost (?) excludes the DAMA claim.
Future detectors will measure the direction of motion of the recoil nucleus and enable a more clearcut WIMP signature.
## 16 Conclusions on DARK MATTER
How to be Sure of WIMP Detection?
We require features that can be due to WIMPs and nothing else.
*
* (i) Gamma-ray annihilation from WIMP annihilation should show a gamma-line in correspondence with the WIMP mass.
* (ii)Annual modulation should show the correct periodicity both in rate and, in future, directional dependence.
Compatible indirect (i) and direct (ii) detection could provide compelling evidence for WIMPs.
Better would be production in a collider consistent with cosmological detection!
## Acknowledgments
We thank Tran Thanh Van for the invitation to La Thuile. This work was supported in part by the US Department of Energy under Grant No. DE-FG02-97ER-41036.
## References
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# Numerical simulation of 2D steady granular flows in rotating drum: On surface flows rheology
## I Introduction
Granular media present numbers of interesting and unusual behaviours: They can flow as liquids, but, under some circumstances, they can jam and resist to external shear stress without deforming. Understanding rheology of granular systems has thus developed along two major themes: The Rapid flows \- gaseous-like - regime where grains interact through binary collisions, are generally described in the framework of the kinetic theory Savage81\_jfm ; Jenkins83\_jfm ; Lun86\_am ; The slow flow \- solid-like - regime where grain inertia is negligible is most commonly described using the tools of soil mechanics and plasticity theory nedderman92\_book . In between these two regimes there exists a dense flow \- liquid-like - regime where grain inertia becomes important but contacts between grains are kept. Rheology of this last regime has been widely investigated experimentally, numerically and theoretically (see Gdrmidi04\_epje for a review), but still remains far from being understood. Several models have been proposed recently to describe dense granular flows by accounting for non-local effects Mills99\_epl ; Andreotti01\_pre ; Jenkins02\_pof ; Bonamy03\_epl ; Rajchenbach03\_prl , by adapting kinetic theory Savage98\_jfm ; Bocquet02\_pre ; Mohan02\_jfm , by modelling dense flows as partially fluidized flows Aranson02\_pre ; Josserand04\_epje or by considering them as quasi-static flows where the mean motion results from transient fractures modelled as self activated process Pouliquen96\_pre ; Debregeas00\_epl ; Pouliquen01\_acs ; Lemaitre02\_prl , but, to our knowledge, none of them succeed to account for all the features experimentally observed.
The most spectacular manifestation of this solid/liquid duality occurs during an avalanche when a thin layer of grains starts to roll at the surface of the packing, most of the grains remaining apparently static. The global evolution of such surface flows can be captured by models derived from non-linear physics Bouchaud94\_jpi ; Boutreux98\_pre ; Aradian98\_pre or fluid mechanics savage89\_jfm ; Douady99\_epjb ; Khakhar97\_pof ; Bonamy02\_pof . However, some experimental results remain unexplained: For instance, experimental velocity profiles measured in two-dimensional (2D) flows Rajchenbach00\_ap ; Bonamy01\_phd ; Bonamy03\_gm or three dimensional (3D) flows Bonamy02\_pof ; Orpe01\_pre ; Jain02\_pof ; Felix02\_phd clearly exhibit the selection of a constant velocity gradient within the flowing layer while momentum balance implies that the shear stress increases linearly with depth. This observation is incompatible with any local and one-to-one stress/strain constitutive relations. Recent experiments Bonamy02\_prl have provided evidence of โjammedโ aggregates embedded in the avalanche. These โsolidโ clusters are found to be power-law distributed without any characteristic length-scales, and may explain the failure of present models. But a clear understanding of the avalanche rheology is still missing.
The purpose of this paper is to investigate the surface flows rheology through numerical simulations of 2D โminimalโ granular systems made of cohesionless weakly polydisperse cylinders confined in a slowly rotating drum. Those allow us to track the evolution of quantities like stress that are not accessible in real experiments. Moreover, they allow to get rid of artefacts such as the friction of beads on the lateral boundaries of the drum that may confuse the interpretation of an experiment. The numerical simulation were performed using contact dynamic methods Jean92\_conf ; Moreau94\_ejmas based on a fully implicit resolution of the contact forces, without any resort to regularization schemes. At a given step of evolution, all the kinematic constraints within the packing are simultaneously taken into account together with the equations of motions to determine all the contact forces in the packing. This allows to deal properly with nonlocal momentum transfers implied in multiple collisions, contrary to Molecular Dynamics schemes traditionally used that reduce the system evolution to a succession of binary collisions.
The simulation scheme and the description of the simulated systems are detailed in Sec. II. In Sec. III, we report comprehensive analysis of volume fraction and velocity (translational and rotational) profiles at the center of the drum. They are compared with experimental data available in the literature. Stress analysis and implications on the rheology of free surface flows are discussed in Sec. IV. and Section V respectively. Sec. VI is focussed on the analysis of both the translational and rotational velocity fluctuations. Finally, Sec. VII summarizes our findings.
## II Simulation Methodology
For this study, we simulate granular systems similar to those investigated experimentally in Rajchenbach00\_ap ; Bonamy01\_phd ; Bonamy02\_prl ; Bonamy03\_gm . We model a 2D rotating drum of diameter $`D_0`$ equal to $`450\mathrm{mm}`$ half-filled with $`7183`$ rigid disks of density $`\rho _0=2.7\mathrm{g}.\mathrm{cm}^3`$ and diameter uniformly distributed between $`3\mathrm{mm}`$ and $`3.6\mathrm{mm}`$. This weak polydispersity prevents any 2D ordering effect that may induce nongeneric effects. Normal restitution coefficient between two disks (respectively between disks and drum) is set to $`0.46`$ (respectively $`0.46`$) and the friction coefficient to $`0.4`$ (respectively $`0.95`$). Normal restitution coefficients and disk/disk friction coefficient were chosen to mimic the experimental flows of aluminium beads investigated in Bonamy01\_phd ; Bonamy03\_gm . The drum/disk friction coefficient was set close to $`1`$ to prevent sliding at the drum boundary.
Numerical simulations dedicated to evolution of granular media can be based either on explicit Drake91\_jfm ; Cundall79\_g ; Kishino88\_conf ; Silbert01\_pre or implicit Jean92\_conf ; Moreau94\_ejmas ; DeSaxe91\_jmsm method. One of the drawbacks of explicit models is to reduce non-local momentum transfers implied in multiple collisions to a succession of binary collisions. Moreover, numerical instabilities can occur in granular flows. They are corrected either by introducing some artificial viscosity or by reducing the size of the time step. The Non Smooth Contact Dynamics method used here is implicit. It provides a nonsmooth formulation of the body s impenetrability condition, the collision rules and the dry Coulomb friction law. The method is extensively described in Jean99\_cmame , and briefly explained below.
Firstly, equations of motion are written for a collection of rigid bodies and discretized by a time integrator Moreau99\_cmame . The interaction problem is then solved at contact level (local level) rather than at particle level (global level) as commonly performed in explicit methods. In other words, equations are written in term of relative velocities $`๐ฎ_\alpha `$ and local impulsions $`๐ซ_\alpha `$ defined at each contact point indexed by $`\alpha `$. The impenetrability condition evoked previously means that particles candidates for contact should not cross the boundaries of antagonistโs bodies. We consider also that contacting bodies do not attract each other, i.e. that the reaction force is positive, and vanishes when the contact vanishes. This can be summarized in the following so-called velocity Signorini Condition:
$$u_n0r_n0u_n.r_n=0,$$
(1)
where the index $`n`$ denotes the normal component of the various quantities (index $`\alpha `$ is omitted). Let us note that this philosophy is different from what is used in explicit methods, where normal forces are usually proportional to the penetration between two particles. The dry frictional law is the Coulombโs one for which the basic features are: The friction force lies in the Coulombโs cone ($`r_t\mu r_n`$, $`\mu `$ friction coefficient), and if the sliding relative velocity is not equal to zero, its direction is opposed to the friction force ($`r_t=\mu r_n`$). This summarized in the following relation:
$$r_t\mu r_nu_t0r_t=\mu r_n\frac{u_t}{u_t}$$
(2)
For rigid bodies we also need to adopt a collision law because the velocity Signorini condition does not give enough information. We adopt the Newton restitution law, $`u_n^+=e_nu_n^{}`$, realistic for collection of disks. The reader can refer to Moreau88\_conf for more explanations about collision laws. Time discretization of equations of motion - where the global contact forces are the only missing quantities to determine the motion of each bead - lead to the following scheme:
$$\{\begin{array}{c}๐๐ซ๐ฎ=๐\hfill \\ law_\alpha [๐ฎ_\alpha ,๐ซ_\alpha ]=.true.,\alpha =1,n_c\hfill \end{array}$$
(3)
where $`๐ฎ`$ and $`๐ซ`$ denotes the vectors containing the relative velocity and the mean contact impulse for all the contact points respectively. The matrix $`๐`$ is the Delassus operator Delassus17\_asens that contains all the local informations (local frames and contact points) as well as the informations related to the contacts connectivity. The right hand side of first line in Eq. 3 represents the free relative velocity calculated by only taking into account the external forces. The operator $`law_\alpha `$ encodes the friction-contact law which should be satisfied by each component of couple $`(๐ฎ_\alpha ,๐ซ_\alpha )`$; $`n_c`$ denotes the number of contact. Systems of Eq. 3 can be solved by a classical non linear Gau$`\beta `$-Seidel algorithm Jean99\_cmame or a Conjugate Projected Gradient one Renouf04\_came . This two algorithms benefit from parallel versions Renouf04\_jcam ; Renouf04\_reef which show their efficiency in the simulations of large systems. Information from this local level, the contact level, is transfered to the global level, the grain level and the configuration of the system is updated.
The procedure to achieve a numerical experiment is the following: All the disks are placed in an immobile drum; Once the packing is stabilized, a constant rotation speed $`\mathrm{\Omega }`$ (ranging from $`2\mathrm{rpm}`$ to $`15\mathrm{rpm}`$) is imposed to the drum; After one round, a steady continuous surface flow is reached (This has been checked by looking at the time evolution of the total kinetic energy within the packing over the next round); One starts then to capture 400 snapshots equally distributed over a rotation of the drum. The time-step is set to $`6.10^3\mathrm{s}`$. The number of time-steps necessary to achieve an experiment ranges from $`4.10^3`$ to $`10^4`$ depending on the rotating speed. All simulations have been performed with LMGC90 software Dubois03\_conf . On a SGI Origin 3800 with 16 processors, about $`20\mathrm{h}`$ are required to achieve one of these simulations.
A typical snapshot of the simulated granular packing is shown in Fig. 1. For each bead of each of the $`400`$ frames within a given numerical experiment, one records the position $`๐ฑ`$ of its center of mass, the โinstantaneousโ velocity $`๐`$ of this center of mass measured over a time window $`\delta _t=6.10^3\mathrm{s}`$ and its angular velocity $`w`$. Voronoรฏ tessellation was used to define the local volume fraction associated to each bead (see e.g. Bonamy02\_prl ). The components of contact stress tensor $`\sigma `$ associated to each bead $`i`$ are computed as Radjai96\_prl ; Radjai98\_prl :
$$\sigma _{\alpha \beta }=\frac{1}{2V_i}\underset{ji}{}x_{ji}^\alpha F_{ji}^\beta ,\alpha ,\beta \{1,2\},$$
(4)
where $`V_i`$ is the volume of the Voronoรฏ polyedra associated to the bead $`i`$, $`๐
_{ji}`$ the contact force between $`i`$ and $`j`$, and $`๐ฑ_{ji}=๐ฑ_j๐ฑ_i`$. In all the following, these quantities are nondimensionalized: Calling $`g`$ the gravity constant and $`d`$ the mean disk diameter, distances, time, velocities and stresses are given in units of $`d`$, $`\sqrt{d/g}`$, $`\sqrt{gd}`$ and $`\rho _0gd`$ respectively. In this paper, we concentrate on the continuum scale by looking at profiles of the time and space averaged quantities. Statistical analysis of these quantities at the grain scale will be presented in a separated paper.
In rotating drum geometries, the surface flow is not fully developed. The frame of study should now be chosen appropriately. One thus define the frame $`\mathrm{}`$ rotating with the drum that coincides with the reference frame $`\mathrm{}_0=(๐_x,๐_z)`$ fixed in the laboratory, so that $`๐_x`$ (resp. $`๐_z`$) is parallel (resp. perpendicular) to the free surface (see Fig. 1) defAngle . In the frame $`\mathrm{}`$, the flow can be considered as quasi-homogeneous at the center of the drum, e.g. within the elementary slice $`\mathrm{\Sigma }`$ (see Fig. 1) 20 beads diameter wide, parallel to $`๐_z`$ located at $`x=0`$. This slice is divided into layers of one mean bead diameter wide parallel to the flow. The given value of a given continuum quantity $`\overline{a}(z)`$ (volume fraction, velocity, stressโฆ) at depth $`z`$ is then defined as the average of the corresponding quantity defined at the grain scale over all the beads in all the 400 frames of the sequence whose center of mass is inside the layer.
## III Kinematic analysis
### III.1 Volume fraction profile
Let us first focus on volume fraction profiles within the packing. Figure 2 displays the volume fraction profile measured for $`\mathrm{\Omega }=6\mathrm{rpm}`$. To check the homogeneity of the flow with regard to $`\nu `$, the elementary slice $`\mathrm{\Sigma }`$ was translated of an increment of 5 bead diameters in both positive $`x`$ and negative $`x`$. The volume fraction profile is found to be invariant under infinitesimal translation along $`\stackrel{}{e}_x`$. At the free surface, $`\nu `$ drops quickly to zero within a small zone of thickness around three/four beads diameter independent of $`\mathrm{\Omega }`$. In all the following, the free surface boundary is set at the lower boundary of this small region (mixed line in Fig. 2), defined at the point where $`\nu `$ becomes larger than $`0.7`$. At the drum boundary, $`\nu `$ jumps also to a much smaller value within a small region about two/three beads diameters thick, which should be attributed to the presence of the smooth drum boundary. Apart fr21.8om these two narrow regions, the volume fraction $`\nu `$ is almost constant within the drum, around the random close packing (RCP) value $`\nu ^{RCP}0.82`$. A closer look (Inset of Fig. 2) suggests that $`\nu `$ is constant within the static phase, and decreases weakly within the flowing layer as defined from the velocity profile in next section. Such behaviour is expected since granular systems should dilate before being able to deform. However, this decreasing is very small and compressible effects can thus be neglected with regard to momentum balance, even if they may significantly alter the local flow rheology Bonamy02\_prl .
### III.2 Velocity profiles
As expected for a quasi-homogeneous flow, the normal component $`v_z`$ of the velocity was found to be negligible compared to the tangential component $`v_x`$ at any depth $`z`$. Figure 3 depicts both the streamwise velocity profile $`v_x(z)`$ (Fig. 3a) and the shear rate profile $`_zv_x(z)`$ (Fig. 3b) for $`\mathrm{\Omega }=6\mathrm{rpm}`$. Both these profiles were found to be invariant under infinitesimal translation along $`\stackrel{}{e}_x`$. Two phases can clearly be observed: A flowing layer exhibiting a linear velocity profile and a static phase experiencing creep motion where both $`v_x`$, and $`v_x/z`$ decay exponentially with depth (see Inset of Fig. 3b). Such shapes are very similar to the ones observed experimentally in 2D flows Rajchenbach00\_ap ; Bonamy01\_phd ; Bonamy03\_gm as well as in 3D flows Bonamy02\_pof ; Orpe01\_pre ; Jain02\_pof ; Felix02\_phd . The interface between the two phases can then be defined by extrapolating the linear velocity profile of the flowing phase to zero (see Fig. 3a). The flowing layer thickness $`H`$ can then deduced.
Velocity profiles measured for $`\mathrm{\Omega }=6\mathrm{rpm}`$ at five different locations $`x`$ are represented in Fig. 4. At these locations, $`_xH`$ is no more equal to zero. The width of the elementary slice $`\mathrm{\Sigma }`$ has thus been decreased to two beads diameters in order to minimize this drift. The shape of the velocity profile remains the same in these locations, with a clear linear velocity profile within the flowing layer and an exponentially decaying velocity within the static phase. Both the characteristic decay length $`\lambda `$ of the exponential creep within the static phase and the constant velocity gradient $`\dot{\gamma }_0`$ within the flowing layer are observed to be independent of the precise location $`x`$ for a given value of $`\mathrm{\Omega }`$.
The โnaturalโ control parameter in our experiment is the rotating speed $`\mathrm{\Omega }`$. However, comparisons between experiments in heap geometry and rotating drum geometry Gdrmidi04\_epje suggest that the main control parameter for the surface flow is the non-dimensionalized flow rate $`Q`$, defined as:
$$Q=_{z=R_0}^0\nu (z)v_x(z)\mathrm{dz}$$
(5)
Its variation as well as the one of the flowing layer thickness $`H`$ and the mean slope $`\theta `$ with respect to $`\mathrm{\Omega }`$ are reported in Tab.1.
Velocity profiles obtained in the center of the drum for various $`Q`$ are represented in Fig. 5a. Apart from the flowing layer thickness $`H`$, the streamwise velocity profile at the center of the drum is characterized by two parameters, namely the characteristic decay length $`\lambda `$ of the exponential creep within the static phase and the constant shear rate $`\dot{\gamma }_0`$ within the flowing layer. Their evolution with respect to the flow rate $`Q`$ is reported in Fig 5b,c. Within the errorbars, $`\lambda `$ is independent of $`Q`$, of the order of $`3\pm 0.3`$. This behaviour is similar to what is reported in both 3D heap flows experiments Komatsu01\_prl and 3D rotating drum experiments Bonamy01\_phd ; Bonamy02\_pof ; CourrechdePont05\_prl , where $`\lambda `$ was found to be $`\lambda 1.4`$ and $`\lambda 2.5`$ respectively.
In our numerical simulation, $`\dot{\gamma }_0`$ exhibits a weak dependence with $`Q`$ (see Fig. 5c). It ranges typically from $`0.1`$ to $`0.25`$ when $`Q`$ is made vary from $`8`$ to $`58`$. This dependency is compatible with the ones observed experimentally in 2D rotating drum by Rajchenbach Rajchenbach00\_ap , who proposed that $`\dot{\gamma }_0`$ scales as $`\dot{\gamma }_0(\mathrm{sin}\theta \mathrm{sin}\mathrm{\Phi })^{1/2}/\mathrm{cos}^{1/2}\mathrm{\Phi }`$, where $`\mathrm{\Phi }`$ refers to the Coulomb friction angle. Value of this angle can be estimated from the variation of the mean flow angle with respect to $`Q`$ (see Fig. 6 and next section) and was found to be $`\mathrm{\Phi }=17.4^{}`$. Inset of Fig. 5c shows that the scaling proposed by Rajchenbach is compatible with our results. It is worth to mention that $`\dot{\gamma }_0`$ was found to be constant, around 0.5, independent of $`Q`$ in 3D experiments in Hele-Shaw drums Bonamy02\_pof ; Gdrmidi04\_epje . This strongly suggests some non-trivial effect of either the lateral confinement or the flow dimension on the profile within the flowing layer. This will be explored in future 3D simulations of rotating drums.
### III.3 Flowing layer thickness and mean flow angle
The thickness of the flowing layer $`H`$ is plotted as a function of the flow rate $`Q`$ in Fig. 6a. As observed experimentally Bonamy02\_pof ; Gdrmidi04\_epje , $`H`$ scales as $`\sqrt{Q}`$, which is expected since the shear rate varies weakly within the flowing layer.
The mean flow angle $`\theta `$ can then be assimilated to an effective friction coefficient $`\mu _{eff}=\mathrm{tan}\theta `$ between the flowing layer and the static phase Douady99\_epjb . Its evolution with respect to the flow rate $`Q`$ is represented in Fig. 6. The effective friction coefficient is found to increase with $`Q`$. Similar increasing was observed experimentally, - at a much larger scale -. It was attributed to wall effects Bonamy02\_pof ; Gdrmidi04\_epje since this dependency was observed to be weaker when the drum thickness is increased Bonamy01\_phd ; Gdrmidi04\_epje . No such wall effects can be invoked in the present study. In other words, part of the increase of the effective friction with flow rate cannot be induced by wall friction contrary to what was suggested in Bonamy02\_pof ; Gdrmidi04\_epje and should be found in the granular flow rheology.
### III.4 Angular velocity profiles
A typical mean angular velocity profile $`\omega (z)`$ has been represented in Fig. 7. It is interesting to plot $`\omega `$ with respect to the vorticity $`\times ๐ฏ=_zv_x`$ (see Inset of Fig. 7). There is a clear relationship between these two quantities: $`\omega =\frac{1}{2}\times ๐ฏ`$ in the whole packing independently of $`Q`$. This relationship is analogue to the one obtained in classical hydrodynamics where the mean rotating speed of the particles is equal to half the vorticity. Such relationship was observed in Molecular Dynamics simulations of dilute granular flows Campbell85\_jam ; Lun91\_jfm , but expected to fail at higher volume fraction Campbell86\_aa ; Campbell86\_jfm . In this latter case, grains were expected to organize in layers the grains of which rotate in the same direction. This would decrease the mean angular velocity of the grains, and $`\omega `$ would be smaller than $`\frac{1}{2}\times ๐ฏ`$. In our numerical simulation, such behaviour is not observed, which suggests that the grains spins do not organize in layer despite the high density of the flow.
## IV Stress analysis
### IV.1 Stress tensor profile
Let us now look at stress profiles - that cannot be measured experimentally. The stress tensor $`\sigma `$ is the sum of three contributions: $`\sigma =\sigma ^c+\sigma ^k+\sigma ^r`$ where $`\sigma ^c`$, $`\sigma ^k`$ and $`\sigma ^r`$ refer to the contact, kinetic and rotational components of the stress tensor respectively. In our dense free surface flows, $`\sigma ^k`$ and $`\sigma ^r`$ are found to be negligible with regard to $`\sigma ^c`$. One can thus assume that $`\sigma \sigma ^c`$.
Components of the contact stress tensor associated to each bead has been computed for each snapshot of each numerical experiment (see Sec.II). The profile of the continuum value of each component of the contact stress tensor - and consequently the total stress tensor - $`\sigma _{\alpha \beta }(z)`$ is then defined over the elementary slice $`\mathrm{\Sigma }`$ located in the center of the drum (see Fig. 1) according to the same procedure used to calculate velocity and volume fraction profile. The tensor $`\sigma `$ is found to be symmetric, i.e. $`\sigma _{xz}=\sigma _{zx}`$. For 2D surface flows, it is thus defined by three independent components $`\sigma _{xx}`$, $`\sigma _{xz}`$ and $`\sigma zz`$. Typical profile of these components with respect to the depth $`z`$ at the center of the drum are represented in Fig. 8.
Shapes of these profiles are quite surprising. In the rotating frame $`\mathrm{}`$, velocity and volume fraction profiles were found to be invariant along infinitesimal translation along $`๐_x`$ within the elementary slice $`\mathrm{\Sigma }`$. In other words, for $`x=0`$, one gets $`\nu /x๐ฏ/x0`$. More generally, it is commonly assumed that at the center of the drum, the $`x`$ derivative of the stress tensor vanishes Gdrmidi04\_epje ; Rajchenbach03\_prl ; Rajchenbach00\_ap ; Bonamy01\_phd ; Bonamy02\_prl . The Cauchy equations would then read:
$$\begin{array}{cc}(a)\hfill & \frac{\sigma _{xz}}{z}=\nu \mathrm{sin}\theta \hfill \\ (b)\hfill & \frac{\sigma _{zz}}{z}=\nu \mathrm{cos}\theta +\nu \mathrm{\Omega }v_x+\nu \mathrm{\Omega }^2z\hfill \end{array}$$
(6)
where $`\theta `$ denotes the mean flow angle. The second right-handed term of Eq. 6b is the Coriolis term. This term is maximum at the free surface where it reaches $`15\%`$ of the first right-handed gravity term e.g. for $`\mathrm{\Omega }=6\mathrm{rpm}`$. the last right-handed term of Eq. 6b is the centrifugal term. This term is maximum at the drum boundary where it reaches $`1\%`$ of the first right-handed gravity term e.g. for $`\mathrm{\Omega }=6\mathrm{rpm}`$. Finally, inertial effects can be neglected and the Cauchy equations for pure steady homogenous flows would come down to:
$$\begin{array}{cc}(a)\hfill & \frac{\sigma _{xz}}{z}=\nu \mathrm{sin}\theta \hfill \\ (b)\hfill & \frac{\sigma _{zz}}{z}=\nu \mathrm{cos}\theta \hfill \end{array}$$
(7)
and, since the volume fraction $`\nu `$ is almost constant, close to the random close packing value $`\nu ^{RCP}=0.82`$:
$$\begin{array}{cc}(a)\sigma _{xz}(z)=z\nu ^{RCP}\mathrm{sin}\theta \hfill & \\ (b)\sigma _{zz}(z)=z\nu ^{RCP}\mathrm{cos}\theta \hfill & \end{array}$$
(8)
These predictions were compared to the measured profiles (Fig. 8). The measured profile $`\sigma _{zz}`$ fits well with Eq. 8b. However, $`\sigma _{xz}`$ departs from Eq. 8a within the static phase. To understand this discrepancy, one looks at the gradient of the stress tensor (see Fig. 9). The $`x`$-derivative of the various components were calculated by translating the elementary slice $`\mathrm{\Sigma }`$ of an increment $`\delta x=5`$ from one side to the other of the reference position $`x=0`$. We checked that the obtained values do not depend on $`\delta x`$. Both $`\sigma _{zz}/x`$ and $`\sigma _{xz}/x`$ vanish within $`\mathrm{\Sigma }`$ at the center of the drum. However $`\sigma _{xx}/x`$ does not. In other words, steady surface flows in rotating drums cannot be considered as quasi-homogenous even at the center of the drum. The Cauchy equations should then read:
$$\begin{array}{cc}(a)\hfill & \frac{\sigma _{xx}}{x}+\frac{\sigma _{xz}}{z}=\nu \mathrm{sin}\theta \hfill \\ (b)\hfill & \frac{\sigma _{zz}}{z}=\nu \mathrm{cos}\theta \hfill \end{array}$$
(9)
This may explain the slight discrepancies observed between homogenous steady heap surface flows and steady surface flows in rotating drum (see e.g. Bonamy03\_epl for related discussions).
## V Constitutive laws
### V.1 Inertial number $`I`$
It was recently suggested Dacruz04\_proc ; Dacruz04\_phd ; Gdrmidi04\_epje ; Jop05\_jfm that the shear state of a dense granular flow can be characterized through a dimensionless number $`I`$, referred to as the inertial number, defined as:
$$I=\frac{_zv_x}{\sqrt{\sigma _{zz}}}$$
(10)
This parameter can be regarded as the ratio between the typical time of deformation $`1/_zv_x`$ and the typical time of confinement $`1/\sqrt{\sigma _{zz}}`$ Gdrmidi04\_epje .
A typical profile of the inertial number $`I`$ is plotted in Fig. 10a. This non-dimmensionalized parameter was shown to be the relevant parameter to account for the transition from the quasi-static regime to the dense inertial regime in plane shear configuration, annular shear and inclined plane configuration Dacruz04\_proc ; Dacruz04\_phd ; Gdrmidi04\_epje . Therefore, it is natural to consider $`I`$ as the relevant parameter to describe the transition from the quasi-static phase and the flowing layer in the surface flow geometry. To check this assumption, we determine the value $`I_{th}`$ of the inertial number at the interface between the static phase/flowing layer interface - defined by extrapolating the linear velocity profile of the flowing phase to zero (see Fig. 3)- for all our numerical experiments carried out at various $`\mathrm{\Omega }`$. Variations of $`I_{th}`$ as a function of $`Q`$ is represented in Fig. 10b. This threshold is found to be constant, equal to:
$$I_{th}\mathrm{1.8.10}^2$$
(11)
which provides a rather strong argument to consider this non-dimensionalized parameter as the relevant one to describe surface flows.
### V.2 Rheology
Now that a relevant parameter describing the local shear state of the flow has been proposed, one can discuss in more detail the flow rheology. As a first guess, it is tempting to consider local constitutive laws relating the components of the stress tensor to $`I`$ through a one-to-one relation. In this case, dimensional analysis leads to:
$$\sigma _{xz}/\sigma _{zz}=\mu (I),\sigma _{xx}/\sigma _{zz}=k(I)$$
(12)
Typical variations of the effective friction coefficient $`\mu `$ as a function of the inertial number $`I`$ are plotted on Fig. 11a. A semilogarithmic representation (see inset of Fig. 11a) shows that data collected for different flow rates $`Q`$ collapse relatively well within the scaling:
$$\mu =a+b\mathrm{log}I$$
(13)
with $`a0.35`$ and $`b0.013`$ when $`I`$ ranges from $`10^4`$ to $`10^1`$. A departure from this scaling is observed when $`I`$ becomes smaller than $`10^4`$. In this latter case, $`\mu `$ decreases more rapidly with $`I`$. It is worth to note that the scaling given by Eq.13 is quantitatively similar to the one observed in the incline plane geometry Dacruz04\_phd , which suggests that both free surface flow and flows down to a rough incline plane may be described through the same constitutive laws. Relating $`\mu `$ and $`I`$ through a local constitutive law seems thus to be relevant.
Figure 11b shows the variations of $`k=\sigma _{xx}/\sigma _{zz}`$ as a function of $`I`$. In the flowing layer i.e. when $`I`$ exceed $`I_{th}`$, $`k1`$. The non monotonic behaviour observed in the static phase is much more suprising: The parameter $`k`$ starts from a value lower than $`1`$ at the drum boundary $`k(I0)0.8`$, increases and reaches a maximum for $`I10^3`$ where $`k(I10^3)1.2`$ and finally decreases for increasing $`I`$ and tends to $`1`$ within the flowing layer. Such observation is very different from the $`k=1`$ behaviour observed in the whole materials in both annular shear and incline geometry Silbert01\_pre ; Dacruz04\_phd ; Gdrmidi04\_epje .
While the profile $`\{\mu (z)\}`$ is observed to be invariant along infinitesimal translation, the profile $`\{k(z)\}`$ is not (Fig. 12). The $`x`$-derivative of $`k`$ is found to be almost constant $`k/x0.05`$ within the whole packing. In other words, the flow cannot be considered as homogeneous at the center of the drum as regard with the parameter $`k`$. Furthermore, while the curves $`\mu (I)`$ collected for different flow rates $`Q`$ collapse fairly well, the curves $`k(I)`$ do not. This strongly suggest that the non-local effects implied e.g. by the existence of multi-scale rigid clusters embedded in the flow Bonamy02\_prl should be found in the constitutive law $`k(I)`$ rather than in $`\mu (I)`$ contrary to what was suggested in Bonamy02\_prl ; Bonamy03\_epl ; Bonamy03\_gm ; Gdrmidi04\_epje .
## VI Fluctuation analysis
Let us now analyse the fluctuations $`\delta v`$ and $`\delta \omega `$ of the velocity and the vorticity respectively. Calling $`๐(๐ฑ,t)`$ the โinstantaneousโ velocity of a bead located at the position $`๐ฑ`$ within the elementary slice $`\mathrm{\Sigma }`$ at a given time $`t`$, the fluctuating part of the velocity $`\delta ๐(๐ฑ,t)`$ is defined as $`\delta ๐(๐ฑ,t)=๐(๐ฑ,t)v_x(z)๐_x`$ where $`v_x(z)`$ denotes the average velocity at the depth $`z`$ (see Fig. 3a). Profiles of velocity fluctuation $`\delta v^2(z)`$ are them computed by dividing $`\mathrm{\Sigma }`$ into layers of one mean bead diameter wide, and averaging $`\delta c^2`$ over all the beads of the 400 frames whose center of mass is inside the corresponding layer. Same procedure is applied to determine the profiles of angular velocities fluctuations. In our athermal granular systems, the only time scale is provided by the velocity gradient $`_zv_x`$. Therefore, we looked at profiles of $`\{\delta v/_zv_x\}`$ and $`\{\delta \omega /_zv_x\}`$ rather than direct profiles of $`\{\delta v\}`$ and $`\{\delta \omega \}`$.
Figure 13 displays both translational velocity fluctuation profile (Fig. 13a) and angular fluctuation (Fig. 13b) nondimensionalized by the shear rate $`_zv_x`$. In both cases, the nondimensionalized fluctuations are found to be constant within the flowing layer i.e. :
$$\begin{array}{cc}\frac{\delta v}{_zv_x}2.65\hfill & \mathrm{for}zH\mathrm{or}II_{th}\hfill \\ \frac{\delta \omega }{_zv_x}3.35\hfill & \mathrm{for}zH\mathrm{or}II_{th}\hfill \end{array}$$
(14)
In the static phase, both $`\delta v/_zv_x`$ and $`\delta \omega /_zv_x`$ are found to increase with the distance from the static/flowing interface. Figure 14 plots both $`\delta v/_zv_x`$ (Fig. 14a) and $`\delta \omega /_zv_x`$ (Fig. 14b) as a function of the inertial number $`I`$. It evidences the existence of two different scalings within the static phase, namely:
$$\begin{array}{cc}(a)\frac{\delta v}{_zv_x}I^{1/2}\hfill & \mathrm{for}II_{th}\hfill \\ (b)\frac{\delta \omega }{_zv_x}I^{1/3}\hfill & \mathrm{for}II_{th}\hfill \end{array}$$
(15)
Such scaling are very similar to the one observed in the shear geometry Dacruz04\_phd . The importance of these fluctuations with regards to the typical rate of deformation $`_zv_x`$ (up to 40), as well as the scaling given by Eq. 15a exhibited within the static phase are compatible with the picture presented in Gdrmidi04\_epje to describe quasi-static flow: The average grains motion is made of a succession of very slow motions when the particle climbs over the next one, and a rapid motion when it is pushed back into the next hole by the confining picture.
## VII Concluding discussion
Rheologies of 2D dense granular flows were investigated through Non Smooth Contact Dynamics simulations of steady surface flows in a rotating drum. Profiles of the different continuum quantities were measured at the center of the drum where the flow is non-accelerating. Volume fraction $`\nu `$ is found to be almost constant, around the Random Close Packing value $`\nu ^{RCP}0.82`$ within the whole packing, except for a tiny dilation (few percents) within the flowing layer, as expected from dilatancy effects. As observed experimentally Rajchenbach00\_ap ; Bonamy01\_phd ; Bonamy03\_gm ; Bonamy02\_pof ; Orpe01\_pre ; Jain02\_pof ; Felix02\_phd , the streamwise velocity profile $`\{v_x(z)\}`$ is found to be linear within the flowing layer, and to decrease exponentially within the static phase. Mean profile of the angular velocity was also measured at the center of the drum and was shown to be equal to half of the vorticity in the whole packing.
In a second step, profiles of the three independent component of the stress tensor were measured at the center of the drum. Quite surprisingly, the flow is found to be non-homogeneous at the center of the drum with regard to one of this component, namely $`\sigma _{xx}`$. In other words, $`_x\sigma _{xx}`$ does not vanish whereas $`_x\nu `$, $`_x๐ฏ`$, $`_x\sigma _{zz}`$ and $`_x\sigma _{xz}`$ vanish.
The inertial number $`I`$ \- defined as the ratio between inertial solicitations and confinement solicitations was determined. This number is shown to be the relevant one to investigate quantitatively the rheology of the surface flows. The transition from the static to the flowing phase is found to occur to a fixed value $`I_{th}`$ of $`I`$,independently of the flow rate $`Q`$. Constitutive laws relating the components of the stress tensor to $`I`$ were determined. The effective friction$`\mu =\sigma _{xz}/\sigma _{zz}`$ is found to increase logarithmically with $`I`$, independently of the flow rate $`Q`$. This relation is found to match quantitatively the one observed in rough incline geometry. On the other hand, the ratio $`k=\sigma _{xx}/\sigma _{zz}`$ is found to be be significantly different from $`k=1`$ in contrast to what was observed in plane shear, annular shear, and rough incline geometry Silbert01\_pre ; Dacruz04\_phd . To be more precise, $`k`$ if found to vary non monotonically with $`I`$. Moreover, $`_xk`$ is found not to vanish contrary to the x-derivative of the other continuum quantities. It is worth to note that $`k=1`$ together with a univocal relation between $`\mu `$ and $`I`$ would have naturally implied a Bagnold velocity profile, as observed in rough incline geometry, but not in the present free surface flow geometry. In other words, the selection of the velocity profile resides more in the function $`k(I,\mathrm{})`$ than in $`\mu (I)`$.
Dependencies of $`\{k(I)\}`$ with Q, as well as the fact that $`_xk`$ does not vanish lead us to conjecture that the ratio $`k`$ encodes the structure of the percolated network of grains in extended contact with each others - referred to as the arches network. In this scenario, the structure of this network - and therefore the ratio $`k`$ \- is related to the global geometry of the packing as well as to the orientation of the flow. This picture is broadly consistent with nonlocal models based on the coexistence of particle chains and fluidlike materials Mills99\_epl ; Bonamy03\_epl . However, a more detailed study is needed to verify this and understand how $`k`$ can be related to the global structure of the force network.
Finally, both velocity $`\delta v`$ and angular velocity $`\delta \omega `$ fluctuations were analysed. These quantities non-dimensionalized by the shear rate $`_zv_x`$ were found to be constant - independent of the flow rate - within the flowing layer thickness. In the static phase, both $`\delta v/_zv_x`$ and $`\delta v/_zv_x`$ were found to decrease as different power-laws with $`I`$. This behaviour is consistent with the idea of an intermittent dynamics generated from the succession of rapid rearrangements and slow displacement Gdrmidi04\_epje . This change of behaviour at the static/flowing interface is broadly consistent with recent measurements of Orpe and Khakhar Orpe04\_prl revealing a sharp transition in the rms velocity distribution at this interface. Understanding what set precisely the scaling laws require precise statistical analysis of beads velocities at the grain scale. This represents interesting topic for a future investigation.
###### Acknowledgements.
This work is supported by the CINES (Centre Informatique National de lโEnseignement Supรฉrieur), Montpellier-France, project mgc257.
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# Quantum squeezing and entanglement in a two-mode Bose-Einstein condensate with time-dependent Josephson-like coupling
## I Introduction
Bose-Einstein condensates (BECs) provide a useful system for investigating matter wave quantum squeezing and entanglement. Squeezing in BECs has already been investigated previously from a number of different perspectivessqueeze1 ; squeeze2 ; squeeze3 ; squeeze4 ; squeeze5 . Squeezed states are quantum states for which no classical analog existsWallsNature . The definition of quantum mechanical squeezing, or the reduction of quantum fluctuations below the standard quantum limit (SQL), derives directly from the Heisenberg Uncertainty Principle: For two arbitrary operators $`\widehat{A}`$ and $`\widehat{B}`$ which obey the commutation relation $`[\widehat{A},\widehat{B}]=\widehat{C}`$, quantum squeezing exists when one of the variables satisfies the relationwallszoller
$$[\mathrm{\Delta }\widehat{A}]^2<\frac{1}{2}|\widehat{C}|.$$
(1)
Although the origin of quantum squeezing in many experiments is enhanced quantum correlations, it must be noted that the inequality (1) is not a sufficient proof of enhanced quantum correlations between the particles, as discussed by Kitagawa and Ueda in the context of spin squeezingueda . A measure of the amount of spin squeezing was defined asueda ; wineland
$$\xi =\frac{2[\mathrm{\Delta }J_{}]^2}{J},$$
(2)
where $`[\mathrm{\Delta }J_{}]^2`$ is the variance in the direction orthogonal to the total spin vector for an $`N`$ spin-$`\frac{1}{2}`$ particle system. This definition takes into account the effect of quantum correlations as it derives from the fact that a minimum uncertainty state for $`N`$ elementary spin-$`\frac{1}{2}`$ particles has spin $`J/2`$ ($`J=N/2`$) equally distributed over any two orthogonal components normal to the mean spin direction. This follows from the fact that, in the absence of quantum correlations, the total variance in the normal direction is simply given by the sum of the variances of the individual elementary spins. Any state with uncertainty less than the minimum uncertainty state, i.e. spin squeezed due to enhanced quantum correlations between the particles, therefore gives $`\xi <1`$.
Quantum entanglement, which is closely related to quantum squeezing, has recently generated a lot of activity amongst researchers, owing to its significant role in the studies of quantum information theory. A large number of theoretical studies exist on generating entangled states using BECs in a variety of different physical settings. Many of these involve two-mode BECs (TBECs) in which atoms in two different hyperfine states are entangledtwocomp1 ; Micheli ; kennedy ; twocomp2 ; twocomp3 . In simple terms, a quantum state is said to be entangled when the state cannot be written as a simple product state, the most prominent example being the famous Bell-state for a two particle system: $`|\psi =(|1|0+|0|1)/\sqrt{2}`$, where โ1โ and โ0โ represent spin โupโ and โdownโ states respectively. From the inseparability criterion for the $`N`$-particle density matrix, a parameter for the entanglement between atoms in TBECs has been derivedtwocomp1 ; Micheli , which has a form similar to Eq. (2):
$$\xi _๐ง^2=\frac{N[\mathrm{\Delta }๐ง_1\widehat{J}]^2}{๐ง_2\widehat{J}^2+๐ง_3\widehat{J}^2}$$
(3)
where $`๐ง_i`$, $`i=1,2,3`$ are unit orthogonal vectors, and $`\widehat{J}`$ is the well-known Schwinger angular momentum operator representing a many particle system. It is easy to see that Eq. (3) may be viewed as a generalization of Eq. (2).
On the other hand, it has also been argued that quantum entanglement in a TBEC is more meaningful when one considers the system as a bipartite system of two modes, analogous to considering entanglement of electromagnetic modes, as opposed to the photons themselves, as there is no definitive measure for entanglement between three or more subsystemsMilburn . In addition, the two modes of a TBEC are clearly distinguishable and experimentally accessible subsystems. In this case, the standard measure of entanglement is the von Neumann entropy of the reduced density operator of either of the subsystems:
$$E(\rho _A)=\mathrm{Tr}\rho _A\mathrm{log}_2(\rho _A),$$
(4)
where $`\rho _A`$ is the reduced density operator for subsystem $`A`$ defined by the partial trace over the other subsystem $`B`$, $`\rho _A=\mathrm{Tr}_B\rho `$. It can be easily shown that for a general bipartite system of TBECs written in the Fock basis $`|\psi (t)=_mc_m(t)|Nm_A|m_B`$, $`\rho _A=_{m=0}^N|c_m(t)|^2|mm|`$ and hence the dynamically evolving bipartite entanglement in TBECs can be parameterized by
$$E(t)=\underset{m=0}{\overset{N}{}}|c_m(t)|^2\mathrm{log}_2|c_m(t)|^2.$$
(5)
In this paper, we calculate analytically dynamical evolution of quantum squeezing as well as quantum entanglement for both types of system decompositions of a TBEC: $`N`$ particle subsystems based on Eq. (3) and the subsystems composed of the two modes based on Eq. (5). We shall consider a TBEC with a continuous time-dependent Raman coupling between the two states, with and without nonlinear interactions. A physical realization of such a system already exists; the Josephson-like coupling between the two levels is provided by an external laser in the Rb TBEC of Ref. myatt , while a state-dependent magnetic field gradient may be applied to induce Josephson tunneling in the Na spinor system of Ref. stenger . We consider a system with an adiabatically time-varying coupling with an off-resonance atom-light interaction. The TBEC is itself an interesting and important system to study: not only does it display effects similar to that of Josephson junction in superconductors, such as the collapse and revival phenomenon in macroscopic scale, it is also the simplest Bose-Hubbard model with a two site lattice potential for which complete description of entanglement can be given.
The central result of this paper is that we use an exact solution to the time-dependent Schrรถdinger equation to find analytical expressions for the dynamical evolution of squeezing and entanglement. In particular, we explicitly identify the effect of the nonlinear interactions on squeezing and entanglement. Since one of the comments of Ref. Milburn with regard to the atom-atom entanglement parameterized by Eq. (3) was that of practicality, we shall restrict our calculations to two different possible initial states for a TBEC, the Dicke state and the phase state. We consider atom-atom entanglement for which clear physical measurement is in principle possible, namely those involving the measurement of the mean and variance of the atom number difference and the relative phase between the two species.
The paper is organized as follows: In Sec. II, we describe, using the Schwinger notation for the angular momentum operators, the TBEC and the exact solution to the time-dependent Schrรถdinger Equation. In Sec. III, we present our main results: the dynamics of the TBEC system with a time-dependent Josephson-like coupling between the two species and the evolution of quantum squeezing and entanglement for the two types of subsystem decompositions for various values of laser couplings, with and without the nonlinear interaction. We conclude in Sec. IV.
## II Formalism
### II.1 Hamiltonian
We consider atomic BECs in two different hyperfine states trapped in a single trap, with a time-varying Raman coupling between the two levels given by a spatially uniform electromagnetic field. The Hamiltonian for this system, with the annihilation operators for the two distinct states denoted $`\widehat{a}`$ and $`\widehat{b}`$ under the two mode approximation is the following:
$`\widehat{H}`$ $`=`$ $`\widehat{H}_a+\widehat{H}_b+\widehat{H}_{\mathrm{int}}+\widehat{H}_{\mathrm{las}}`$ (6)
$`\widehat{H}_a`$ $`=`$ $`\omega _a\widehat{a}^{}\widehat{a}+{\displaystyle \frac{U_a}{2}}\widehat{a}^{}\widehat{a}^{}\widehat{a}\widehat{a}`$ (7)
$`\widehat{H}_b`$ $`=`$ $`\omega _b\widehat{b}^{}\widehat{b}+{\displaystyle \frac{U_b}{2}}\widehat{b}^{}\widehat{b}^{}\widehat{b}\widehat{b}`$ (8)
$`\widehat{H}_{\mathrm{int}}`$ $`=`$ $`{\displaystyle \frac{U_{ab}}{2}}\widehat{a}^{}\widehat{a}\widehat{b}^{}\widehat{b}`$ (9)
$`\widehat{H}_{\mathrm{las}}`$ $`=`$ $`\mathrm{\Omega }(t)(\widehat{a}^{}\widehat{b}e^{i\phi (t)}+\widehat{b}^{}\widehat{a}e^{i\phi (t)}).`$ (10)
$`\widehat{H}_a`$ and $`\widehat{H}_b`$ describe the two condensates undergoing self-interaction while $`\widehat{H}_{\mathrm{int}}`$ and $`\widehat{H}_{\mathrm{las}}`$ describe the condensates interacting with each other via collisions and laser-induced interactions respectively. $`\widehat{H}_{\mathrm{las}}`$ describes a time-dependent coupling, rendering an overall time-dependent Hamiltonian for the system. The Hamiltonian may be rewritten by employing the Schwinger notation for the angular momentum operators, namely, $`\widehat{J}_x=\frac{1}{2}(\widehat{a}^{}\widehat{b}+\widehat{b}^{}\widehat{a})`$, $`\widehat{J}_y=\frac{1}{2i}(\widehat{a}^{}\widehat{b}\widehat{b}^{}\widehat{a})`$, and $`\widehat{J}_z=\frac{1}{2}(\widehat{a}^{}\widehat{a}\widehat{b}^{}\widehat{b})`$ with the Casimir invariant $`J^2=\frac{\widehat{N}}{2}(\frac{\widehat{N}}{2}+1)`$ where $`\widehat{N}=\widehat{a}^{}\widehat{a}+\widehat{b}^{}\widehat{b}`$ is the total number operator and is a conserved quantity. Physically, eigenvalues of the operator $`\widehat{J}_z`$ represents the difference in the number of atoms in different hyperfine levels, while $`\widehat{J}_x`$ and $`\widehat{J}_y`$ takes on the meaning of relative phase between the two species. The $`z`$ component angular momentum eigenstates, known also as the Dicke statesbloch ; barnettradmore , are written as $`|j,m`$ where $`m=j,\mathrm{},j`$ with $`\widehat{J}_z|j,m=m|j,m`$. Here, $`j=N/2`$ and is the quantum number of angular momentum. The raising and lowering operators are defined in the usual way as $`\widehat{J}_\pm =\widehat{J}_x\pm i\widehat{J}_y`$ such that $`\widehat{J}_\pm |j,m=\sqrt{(jm)(j\pm m+1)}|j,m\pm 1`$. In terms of the angular momentum operators, the Hamiltonian takes the formchoi ; chen :
$$\widehat{H}(t)=\omega _0\widehat{J}_z+q\widehat{J}_z^2+\mathrm{\Omega }(t)\left[\widehat{J}_+e^{i\phi (t)}+\widehat{J}_{}e^{i\phi (t)}\right],$$
(11)
where $`\omega _0=\omega _a\omega _b+(N1)(U_aU_b)/2`$, $`q=(U_a+U_bU_{ab})/2`$. It is noted that $`U_a(U_b)`$ or $`U_{ab}`$ may, in principle, be tuned via Feshbach resonance through the application of an external magnetic fieldfeshbach ; the factors $`\omega _0`$ and $`q`$ are consequently adjustable parameters.
### II.2 Solution to the time-dependent Schrรถdinger Equation
An exact solution to the Schrรถdinger Equation
$$i\mathrm{}\frac{d}{dt}|\psi (t)=\widehat{H}(t)|\psi (t)$$
(12)
with the time-dependent Hamiltonian $`\widehat{H}(t)`$ of Eq. (11) can be given in terms of a time evolution operator $`\widehat{U}(t)`$, $`|\psi (t)=\widehat{U}(t)|\psi (0)`$ chen :
$$U(t)=\widehat{R}^{}(t)e^{iH^{}(t)}\widehat{R}(0),$$
(13)
where $`\widehat{R}`$ is a time dependent unitary transformation defined as
$$\widehat{R}(t)=\mathrm{exp}\left[\frac{\lambda }{2}(\widehat{J}_{}e^{i\phi (t)}\widehat{J}_+e^{i\phi (t)})\right],$$
(14)
and
$$\widehat{H}^{}(t)=\widehat{R}\widehat{H}(t)\widehat{R}^{}i\widehat{R}\frac{}{t}\widehat{R}^{}.$$
(15)
The parameter $`\lambda `$ in Eq. (14) is an auxiliary parameter which may be chosen to simplify the transformed Hamiltonian, $`\widehat{H}^{}`$. It can be shown that $`\widehat{R}`$ generates a gauge transformation under which the time-dependent Schrรถdinger Equation is covariantliang :
$$i\mathrm{}\frac{d}{dt}|\psi ^{}(t)=\widehat{H}^{}(t)|\psi ^{}(t),$$
(16)
where the transformed state $`|\psi ^{}(t)=\widehat{R}|\psi (t)`$. For our case, $`\lambda `$ is chosen such that the Hamiltonian is diagonal in the $`\widehat{J}_z`$ representation. With an additional assumption of adiabaticity conditions $`d\phi /dt0`$ and $`d\lambda /dt0`$, and also assuming that the two-photon transition terms proportional to $`\widehat{J}_+^2e^{i2\phi (t)}`$ and $`\widehat{J}_{}^2e^{i2\phi (t)}`$ can be neglected tunneling , one obtains for the transformed Hamiltonian
$$\widehat{H}^{}(t)=\sqrt{\omega _0^2+4\mathrm{\Omega }^2(t)}\widehat{J}_z\frac{q}{2}\widehat{J}_z^2,$$
(17)
where $`\lambda `$ is chosen such that
$$\mathrm{tan}\lambda =\frac{2\mathrm{\Omega }(t)}{\omega _0}.$$
(18)
The time evolution of any observable $`\widehat{A}`$ is then given in the Heisenberg picture as $`\widehat{A}(t)\psi (0)|U^{}(t)\widehat{A}U(t)|\psi (0)`$, where $`U(t)`$ is defined in Eqs. (13-15) and $`|\psi (0)`$ is the initial quantum state of the system. We note here that by ignoring the two-photon transitions for the purpose of obtaining analytical solutions, the effect of โtwo-axis counter-twistingโueda is not included in our effective Hamiltonian. For concreteness, we shall consider in this paper the case $`\phi (t)=\mathrm{\Delta }t`$ where $`\mathrm{\Delta }`$ gives the detuning of the laser from the $`|A|B`$ transition between the two species. Also, we shall consider the case $`U_a=U_b>U_{ab}`$ and identical trapping potentials for the two species, corresponding to the $`|F=1,M_F=\pm 1`$ hyperfine states of Na trapped in an optical dipole traptwocomp1 . This implies $`q>0`$ as well as $`\omega _0=0`$ i.e. $`\lambda =\pi /2`$ in Eq. (14).
### II.3 Initial quantum state of TBEC
We consider in this paper an $`SU(2)`$ atomic coherent state or a coherent spin state (CSS), $`|\theta ,\varphi `$bloch ; barnettradmore as the initial quantum state for a TBEC. A CSS is a minimum uncertainty state that describes a system with well-defined relative phase between the two species, and provides a good description of TBECs under suitable experimental conditionssavage . It is defined mathematically by applying the rotation operator on the extreme Dicke state $`|j,j`$ or $`|j,j`$. The definition that we use in this paper is:
$$|\theta ,\varphi =\mathrm{exp}\left[\frac{\theta }{2}(\widehat{J}_{}e^{i\varphi }\widehat{J}_+e^{i\varphi })\right]|j,j.$$
(19)
A CSS $`|\theta ,\varphi `$ is therefore an eigenstate of the spin component in the $`(\theta ,\varphi )`$ direction $`\widehat{J}_{\theta ,\varphi }\widehat{J}_x\mathrm{sin}\theta \mathrm{cos}\varphi +\widehat{J}_y\mathrm{sin}\theta \mathrm{sin}\varphi +\widehat{J}_z\mathrm{cos}\varphi `$ with eigenvalue $`J`$ where $`\theta `$ and $`\varphi `$ denote polar and azimuthal angles. It can be shown, following from Eq. (19), that CSS may be written as a superposition state:
$`|\theta ,\varphi `$ $`=`$ $`{\displaystyle \underset{m=j}{\overset{j}{}}}_m^j(\theta ,\varphi )|j,m,`$ (20)
where
$`_m^j(\theta ,\varphi )`$ $`=`$ $`(C_{j+m}^{2j})^{1/2}\mathrm{cos}^{j+m}\left({\displaystyle \frac{\theta }{2}}\right)\mathrm{sin}^{jm}\left({\displaystyle \frac{\theta }{2}}\right)`$ (21)
$`\times e^{i(jm)\varphi }`$
and $`C_m^n`$ denotes the combination, $`C_m^n=n!/[(nm)!m!]`$. Further extensive discussions on the properties of the CSS may be found in Refs. bloch ; barnettradmore . In this paper, we shall consider initial CSSโs that may be produced experimentally, $`|\theta =0,\varphi =0`$, and $`|\theta =\pi /2,\varphi =0`$. The CSS $`|\theta =0,\varphi =0`$ is simply a Dicke state $`|j,j`$ in which all atoms are found in one hyperfine level, and the CSS $`|\theta =\pi /2,\varphi =0`$ is an eigenstate of $`\widehat{J}_x`$ which is a state with a well-defined phase difference $`\varphi =0`$, and is a phase state of a two-mode boson systemCastin . States with different values of $`\theta `$, although simple to deal with mathematically, are difficult to produce and observe experimentally. For the same reason, we shall restrict our consideration of squeezing to those in the $`x`$, $`y`$, and $`z`$ directions since the variances in the relative phase and atom number are physical quantities clearly measurable without the ambiguity associated with the measurement of variances in other directions.
## III Results
In this section, we present various results for the two initial states for a TBEC, the Dicke state and the phase state. The coefficient $`q`$ provides the strength of the scattering interaction between bosons, and its magnitude in relation to the tunnelling coupling $`\mathrm{\Omega }`$ (as we consider $`\omega _0=0`$ in this paper) is an important parameter in determining distinct coupling regimesleggett . Of the three coupling regimes, namely, the Rabi ($`q/\mathrm{\Omega }1/N`$), the Josephson ($`1/Nq/\mathrm{\Omega }N`$) and the Fock ($`q/\mathrm{\Omega }N`$) regimes, it is clear that the regimes of more physical interest are the Rabi and the Josephson regimes, since, in the Fock regime the tunnelling coupling is overwhelmed by the nonlinear interaction term, resulting in no new effects of particular interest. This has been confirmed numericallytonel .
In the Rabi regime, the tunnelling coupling dominates with $`q/\mathrm{\Omega }`$ near zero. We shall therefore consider the limiting case of $`q=0`$ i.e. the linear regime which not only illustrates the essential physics of the Rabi regime but also demonstrate clearly the effect of nonlinear interaction on squeezing and entanglement. In addition, we shall consider the lower end of the Josephson regime, as it was found that important features rapidly disappear as one nears the Fock regime. The two couplings considered in this paper are therefore $`q=0`$ and $`q=q_j`$ with $`q_j/\mathrm{\Omega }=3/N`$ for the Josephson regime, where we choose $`N`$ to be 400. We also have an additional parameter in our model not considered in some of the previous work, that of a detuning $`\mathrm{\Delta }`$. As will be shown below, this enters mainly as a phase shift which affects the results in a nontrivial manner. To illustrate the effect of this parameter we shall consider detuning of the same magnitude as the $`q`$ for the Josephson regime, ($`\mathrm{\Delta }=q_j`$) as well as a relatively large value of detuning ($`\mathrm{\Delta }/\mathrm{\Omega }=10q_j`$). In particular, the results in the linear ($`q=0`$) regime will help us see the role of $`\mathrm{\Delta }`$ with respect to $`q`$. The result for $`\mathrm{\Delta }=0`$ can be easily interpolated from these two values of detunings so will not be explicitly presented.
### III.1 Dynamics of the macroscopic spin vector
We first consider how the quantum state of a TBEC represented by the fictitious spin vector evolves in time by calculating the expectation values $`\widehat{J}_\alpha (t)`$ where $`\alpha `$ stands for $`x`$, $`y`$, and $`z`$. This later gives us insight into how the squeezing and entanglement in a TBEC dynamically evolves. We evaluate the mean values of spin components based on Eq. (13),
$$\widehat{J}_\alpha =\theta ,\varphi |\widehat{R}^{}(0)e^{i\widehat{H}^{}}\widehat{R}(t)\widehat{J}_\alpha \widehat{R}^{}(t)e^{i\widehat{H}^{}}\widehat{R}(0)|\theta ,\varphi .$$
(22)
In the first instance, this may be calculated by using the commutation relation amongst the spin-$`J`$ components, namely $`[\widehat{J}_z,\widehat{J}_\pm ]=\pm \widehat{J}_\pm `$ and $`[\widehat{J}_+,\widehat{J}_{}]=2\widehat{J}_z`$, and employing the Baker-Hausdorff theorembarnettradmore . Useful identities are given in Eqs. (23-25) with $`\phi (t)\mathrm{\Delta }t`$:
$`\widehat{R}\widehat{J}_z\widehat{R}^{}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(\widehat{J}_+e^{i\mathrm{\Delta }t}+\widehat{J}_{}e^{i\mathrm{\Delta }t}),`$ (23)
$`\widehat{R}\widehat{J}_+\widehat{R}^{}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(\widehat{J}_+\widehat{J}_{}e^{2i\mathrm{\Delta }t})+\widehat{J}_ze^{i\mathrm{\Delta }t},`$ (24)
$`\widehat{R}\widehat{J}_{}\widehat{R}^{}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(\widehat{J}_{}\widehat{J}_+e^{2i\mathrm{\Delta }t})+\widehat{J}_ze^{i\mathrm{\Delta }t}.`$ (25)
In principle, it is possible to repeatedly apply the Baker-Hausdorff theorem to obtain a general expression for $`\widehat{J}_\alpha (t)\widehat{U}^{}(t)\widehat{J}_\alpha \widehat{U}(t)`$; for the special case of $`q=0`$, the expressions may, in fact, be simplified further, giving
$$\widehat{J}_z(t)=\frac{N}{2}\mathrm{cos}\theta \mathrm{cos}[2\mathrm{\Omega }t+\mathrm{\Delta }t+\varphi ],$$
(26)
and
$`\widehat{J}_\pm (t)`$ $`=`$ $`{\displaystyle \frac{N}{2}}e^{i\mathrm{\Delta }t}[\mathrm{cos}\left({\displaystyle \frac{\theta \pi /2}{2}}\right)`$ (27)
$`i\mathrm{sin}\left({\displaystyle \frac{\theta \pi /2}{2}}\right)\mathrm{sin}(2\mathrm{\Omega }t+\mathrm{\Delta }t+\varphi )].`$
It is found however that the expressions quickly become unwieldy especially for $`q0`$. The way around the problem is to note an important and useful point that, for the Hamiltonian under consideration, the transformation operator $`\widehat{R}`$ happens to have the same form as the generator of $`SU(2)`$ atomic coherent states\[Eq. (19)\] if one writes $`\phi (t)=\varphi `$. In particular, the effect of $`\widehat{R}`$ is to rotate the CSS such that $`\widehat{R}(\theta ^{},\varphi ^{})|\theta ,\varphi |\theta +\theta ^{},\varphi \varphi ^{}`$, i.e.
$$\widehat{R}(t)|\theta ,\varphi =\underset{m=j}{\overset{j}{}}_m^j(\theta \pi /2,\varphi \mathrm{\Delta }t)|j,m,$$
(28)
where $`_m^j(\theta ,\varphi )`$ is defined in Eq. (21). Using the well-known relation for the raising and lowering operators $`\widehat{J}_\pm `$, and the fact that $`F(\widehat{J}_z)|j,m=F(m)|j,m`$ where $`F`$ denotes some analytic function, one can obtain, after some algebra, an analytical expression for $`\widehat{J}_z(t)`$ as a summation:
$`\widehat{J}_z(t)`$ $`=`$ $`{\displaystyle \underset{m=N/2}{\overset{N/21}{}}}๐_1(\theta ,m)`$ (29)
$`\times \mathrm{cos}\left[2\mathrm{\Omega }tq\left(m+{\displaystyle \frac{1}{2}}\right)t+\mathrm{\Delta }t+\varphi \right],`$
where we have defined
$`๐_\alpha (\theta ,m)`$ $`=`$ $`C_{N/2+m+1}^N\left({\displaystyle \frac{N}{2}}+m+1\right)`$ (30)
$`\times \mathrm{cos}^{2N}\left({\displaystyle \frac{\theta \pi /2}{2}}\right)`$
$`\times \mathrm{tan}^{N2m\alpha }\left({\displaystyle \frac{\theta \pi /2}{2}}\right).`$
Following somewhat more involved but identical steps as above, we obtain:
$`\widehat{J}_\pm (t)`$ $`=`$ $`{\displaystyle \underset{m=N/2}{\overset{N/21}{}}}e^{i(\mathrm{\Delta }t\pi /2)}๐_1(\theta ,m)`$ (31)
$`\times \mathrm{sin}\left[2\mathrm{\Omega }tq\left(m+{\displaystyle \frac{1}{2}}\right)t+\mathrm{\Delta }t+\varphi \right]`$
$`+{\displaystyle \underset{m=N/2}{\overset{N/2}{}}}e^{i\mathrm{\Delta }t}๐_0(\theta ,m)\left({\displaystyle \frac{m}{jm}}\right).`$
from which $`\widehat{J}_x(t)`$ and $`\widehat{J}_y(t)`$ may be deduced: $`\widehat{J}_x=\frac{1}{2}(\widehat{J}_++\widehat{J}_{})`$ and $`\widehat{J}_y=\frac{1}{2i}(\widehat{J}_+\widehat{J}_{})`$.
We first present in Figs. 1 and 2, the expectation values $`\widehat{J}_\alpha (t)`$, $`\alpha =x,y,z`$ scaled to the total number of atoms $`N`$ for $`q=0`$, and $`q=q_j`$ respectively for the detunings $`\mathrm{\Delta }=q_j`$ and $`\mathrm{\Delta }=10q_j`$. It is noted that no significant qualitative difference is observed in these scaled amplitudes when one uses different number of atoms. In all the figures in this paper, the solid line represents the result for the initial CSS $`|\theta ,\varphi =0`$ with $`\theta =0`$ i.e. the Dicke state, while the dashed line corresponds to the initial phase state, $`\theta =\pi /2`$. In Fig. 1, it is seen that, for the initial state $`|\theta =0,\varphi =0`$, the $`y`$ and $`z`$ components of spin, $`\widehat{J}_y`$ and $`\widehat{J}_z`$, undergo oscillatory evolution with a $`\pi /2`$ phase shift. This implies that the macroscopic spin vector may be visualized as if undergoing a circular motion in the $`y`$-$`z`$ plane. Superposed on this motion is the gradual increase of the oscillation in the $`x`$ component, $`\widehat{J}_x`$ with the amplitude of the $`y`$ component being proportionally reduced with a $`\pi /2`$ phase shift. This implies that the spin vector initially undergoing a circular motion in the $`y`$-$`z`$ plane rotates around the $`z`$-axis with increasing amplitude. This motion gets more pronounced for the larger $`\mathrm{\Delta }`$ in that the spin vector undergoing circular motion in the $`y`$-$`z`$ plane rotates around the $`z`$-axis relatively quickly into a circular motion in the $`x`$-$`z`$ plane. This continues on to return to the $`y`$-$`z`$ plane, repeating this pattern over time. For the initial state $`|\theta =\pi /2,\varphi =0`$, there is no oscillation in the $`y`$-$`z`$ plane, and only a slow rotation through the $`x`$-$`y`$ plane is observed. This state may be identified as a โself-trappingโ state, as there is no transfer of populations during the evolution. With higher $`\mathrm{\Delta }`$, the frequency of this rotation is clearly increased.
For the case of $`q=q_j`$, Fig. 2 shows the dephasing or โcollapseโ of oscillations for the initial state $`|\theta =0,\varphi =0`$. The three spin components collapse to give average spin of almost zero. As will be seen below, this has a significant impact on the variance, and consequently on squeezing and entanglement in a TBEC. For the self-trapping initial state $`|\theta =\pi /2,\varphi =0`$, there is no change in the behavior with $`q=q_j`$. As to be expected, higher values of nonlinearity $`q`$ results in shorter time scales for the collapse and with higher number of atoms, no significant qualitative difference is observed except that the collapse happens somewhat faster. The long time simulation is provided in Fig. 3, demonstrating clear revivals.
The zero mean spin due to collapse indicates that the quantum state of a TBEC evolves from a CSS which is similar to the usual coherent state in quantum optics to a state which is similar to the number eigenstate with an equal number of atoms in each mode. The state that the CSS evolves into is, however, not exactly a self-trapping state as it clearly revives back into a CSS; it may be viewed as a quantum state with a small spread around the mean self-trapping state. For lack of better terminology we shall refer to this collapsed quantum state as a โquasi-self-trapping stateโ in this paper.
### III.2 Quantum squeezing and fluctuations
In the language of spin squeezing, any reduction of the quantum noise for $`q=0`$ corresponds to an effective squeezing due to the rotation of the coordinate axes, while squeezing observed for $`q=q_j`$ comes from a complex combination of coordinate rotation on top of the interatomic collision effect of the one-axis twisting term proportional to $`\widehat{J}_z^2`$ueda . In order to calculate the amount of squeezing in the system, we first calculate the variances in the three components of the macroscopic spin vector, i.e. $`[\mathrm{\Delta }\widehat{J}_\alpha (t)]^2=\widehat{J}_\alpha ^2(t)\widehat{J}_\alpha (t)^2`$, $`\alpha =x,y,z`$ which, as mentioned above are experimentally accessible quantities. The unitary property of $`\widehat{R}`$ i.e. $`\widehat{R}^{}\widehat{R}=๐`$ is used to evaluate terms of the form $`\widehat{R}\widehat{J}_\alpha \widehat{J}_\beta \widehat{R}^{}\widehat{R}\widehat{J}_\alpha \widehat{R}^{}\widehat{R}\widehat{J}_\beta \widehat{R}^{}`$ along with the identities Eqs. (23-25).
It can be shown that $`\widehat{J}_z^2(t)`$ is given by
$`\widehat{J}_z^2(t)`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{m=N/2}{\overset{N/22}{}}}๐_2(\theta ,m)\left({\displaystyle \frac{N}{2}}m1\right)`$ (32)
$`\times \mathrm{cos}[4\mathrm{\Omega }t2q(m+1)t+2\mathrm{\Delta }t+2\varphi ]`$
$`+{\displaystyle \frac{1}{4}}{\displaystyle \underset{m=N/2}{\overset{N/21}{}}}๐_2(\theta ,m)\left({\displaystyle \frac{N}{2}}m\right)`$
$`+๐_0(\theta ,m)\left({\displaystyle \frac{N}{2}}+m+1\right),`$
so that the variance in the $`z`$ direction, $`[\mathrm{\Delta }\widehat{J}_z(t)]^2`$, is simply given by subtracting the square of Eq. (29) from Eq. (32). On the other hand, $`\widehat{J}_x^2=\frac{1}{4}[\widehat{J}_+^2+\widehat{J}_{}^2+\widehat{J}_+\widehat{J}_{}+\widehat{J}_{}\widehat{J}_+]`$ and $`\widehat{J}_y^2=\frac{1}{4}[\widehat{J}_+\widehat{J}_{}+\widehat{J}_{}\widehat{J}_+\widehat{J}_{}^2\widehat{J}_+^2]`$. The required expectation values of $`\widehat{J}_\pm ^2`$ and $`\widehat{J}_\pm \widehat{J}_{}`$ are given by:
$`\widehat{J}_+^2(t)`$ $`=`$ $`{\displaystyle \underset{m=N/2}{\overset{N/22}{}}}{\displaystyle \frac{e^{2i\mathrm{\Delta }t}}{2}}๐_2(\theta ,m)\left({\displaystyle \frac{N}{2}}m1\right)\mathrm{cos}[4\mathrm{\Omega }t2q(m+1)t+2\mathrm{\Delta }t+2\varphi ]`$ (33)
$`{\displaystyle \underset{m=N/2}{\overset{N/21}{}}}{\displaystyle \frac{e^{2i\mathrm{\Delta }t}}{4}}\{๐_2(\theta ,m)({\displaystyle \frac{N}{2}}m)+๐_0(\theta ,m)({\displaystyle \frac{N}{2}}+m+1)`$
$`+(2m+1)i\mathrm{sin}[2\mathrm{\Omega }tq(m+{\displaystyle \frac{1}{2}})t+\mathrm{\Delta }t+\varphi ]\}+{\displaystyle }_{m=N/2}^{N/2}e^{2i\mathrm{\Delta }t}๐_0(\theta ,m)({\displaystyle \frac{m^2}{jm}}),`$
$`\widehat{J}_+(t)\widehat{J}_{}(t)`$ $`=`$ $`{\displaystyle \underset{m=N/2}{\overset{N/22}{}}}{\displaystyle \frac{1}{2}}๐_2(\theta ,m)\left({\displaystyle \frac{N}{2}}m1\right)\mathrm{cos}[4\mathrm{\Omega }t2q(m+1)t+2\mathrm{\Delta }t+2\varphi ]`$ (34)
$`+{\displaystyle \underset{m=N/2}{\overset{N/21}{}}}{\displaystyle \frac{1}{4}}\{๐_2(\theta ,m)({\displaystyle \frac{N}{2}}m)+๐_0(\theta ,m)({\displaystyle \frac{N}{2}}+m+1)`$
$`4๐_1(\theta ,m)\mathrm{cos}[2\mathrm{\Omega }tq(m+{\displaystyle \frac{1}{2}})t+\mathrm{\Delta }t+\varphi ]\}+{\displaystyle }_{m=N/2}^{N/2}๐_0(\theta ,m)({\displaystyle \frac{m^2}{jm}}),`$
with $`\widehat{J}_{}^2(t)\widehat{J}_+^2(t)^{}`$ where asterisk (\*) denotes complex conjugate, and $`\widehat{J}_{}(t)\widehat{J}_+(t)=\widehat{J}_+(t)\widehat{J}_{}(t)2\widehat{J}_z(t)`$, as to be expected from the commutation relation $`[\widehat{J}_+,\widehat{J}_{}]=2\widehat{J}_z`$. From these expressions and those for $`\widehat{J}_x(t)`$ and $`\widehat{J}_y(t)`$ obtained above, the variances $`[\mathrm{\Delta }\widehat{J}_x(t)]^2`$ and $`[\mathrm{\Delta }\widehat{J}_y(t)]^2`$ may be calculated.
We find that, for $`q=0`$, the variances $`[\mathrm{\Delta }\widehat{J}_\alpha ]^2`$, $`\alpha =x,y,z`$, show an oscillatory behavior which is bounded above by $`N/4`$ \[Fig. 4\]. This is an expected result for a CSS, which is known to have variances in the standard quantum limit (SQL) of $`J/2`$. One key observation regarding Fig. 4 is that the variances do go below the SQL. This kind of squeezing which is due to the rotation of the coordinate axisueda naturally occurs since a spin vector is an eigenstate of spin in one direction with zero variance in that direction. As the spin vector traverses the phase space due to dynamical evolution, the error ellipsoid follows the path of the spin vector in such a way that the minor axis of the ellipsoid periodically lines up with the $`x`$, $`y`$ or $`z`$ axis, resulting in the reduction of quantum fluctuations in that direction.
For small $`\mathrm{\Delta }`$, the variance in the $`z`$ direction show oscillatory reduction with $`\pi /2`$ phase shift with that of the $`y`$ component. The variance in the $`x`$ direction is seen as more or less maintaining its value near the SQL. For larger $`\mathrm{\Delta }`$, the uncertainty in the $`x`$ direction is found oscillating with increasing amplitude, which is consistent with the error ellipsoid following the motion of the rotating spin vector discussed above.
For the $`q=q_j`$ case, the variances in the $`x`$, $`y`$, and $`z`$ directions give significantly different behavior, as shown in Fig. 5. The unexpected feature is the large variance equivalent to the maximum relative uncertainty (standard deviation) of order $`\pm 35\%`$ per measurement. Mathematically, this can be understood as the corollary of the collapsing mean spin components; the collapsing spin vector implies that the the variance $`[\mathrm{\Delta }\widehat{J}_\alpha ]^2=\widehat{J}_\alpha ^2\widehat{J}_\alpha ^2`$ is dominated by the $`\widehat{J}_\alpha ^2`$ term of the order $`N^2`$ as it cannot be cancelled by the $`\widehat{J}_\alpha ^2`$ term. Careful analysis of the variance for the $`q=0`$ case discussed above reveals that the terms of the form $`\widehat{J}_\alpha ^2`$ are almost exactly cancelled by the $`\widehat{J}_\alpha ^2`$ which are also of the order $`N^2`$ to maintain the variances of the order of the SQL, $`N/4`$. This type of cancellation cannot happen as the spin vector collapses with $`q=q_j`$. Physically, the increased uncertainty can be attributed to the quantum state evolving away from the minimum uncertainty state of CSS into a quasi-self-trapping state with much higher quantum fluctuations. In order to display periodic behavior in the variances, we plot in Fig. 6 the variances over the identical time period as in Fig. 3. The result is in agreement with the purely numerical result of Tonel et al.tonel . It is found that, as the revivals occur, the variances do go below the SQL periodically. Although it is not clearly depicted due to scaling, the results for the initial phase state $`|\theta =\pi /2,\varphi =0`$ are identical to that for the $`q=0`$ case i.e. the variance stays at 1 for the atom number difference, and oscillates between $`0`$ and the SQL for the relative phase.
### III.3 Dynamical evolution of quantum entanglement in a TBEC
As mentioned in Ref. Milburn , the amount of entanglement depends on the way a system is partitioned into subsystems. In this subsection we calculate the dynamical evolution of quantum entanglement in two different subsystem decomposition of the TBEC: entanglement between the particles and that between the two modes.
#### III.3.1 Entanglement between the particles
The atom-atom entanglement is parameterized by Eq. (3). Equation (3) is useful since, regardless of the actual spin direction at a given time, one need only to identify three unit orthogonal directions $`๐ง_i`$ to check whether the system has become entangled due to system dynamics. If oneโs goal is to identify maximum entanglement possible between the atoms, it is necessary to scan through all possible $`๐ง_i`$โs within the unit sphere at each point in time. However, it should be noted that keeping to a fixed direction may help in constructing an actual experimental scheme to measure $`\xi _\alpha ^2(t)`$. In particular, we consider in this paper experimentally meaningful variances and amplitudes of the macroscopic spin vector in the $`x`$, $`y`$, and $`z`$ directions. We therefore consider
$$\xi _\alpha ^2(t)=\frac{N[\mathrm{\Delta }\widehat{J}_\alpha (t)]^2}{\widehat{J}_\beta (t)^2+\widehat{J}_\gamma (t)^2}$$
(35)
where $`\alpha ,\beta ,\gamma `$ cycle through $`x,y,z`$.
For $`q=0`$, it was found that $`\xi _\alpha ^2(t)1`$, $`\alpha =x,y,z`$, at all times for both initial states i.e. not entangled despite the fact that, as already seen in Fig. 4, squeezing due to the rotation of the coordinate axis does occur. This shows that the observed squeezing is not accompanied by an increase in quantum correlations between the particles. On the other hand, for $`q=q_j`$, it is found that $`\xi _\alpha ^2`$ is maintained at 1 for the initial phase state as to be expected, while for the initial Dicke state, the very quickly increasing quantum fluctuations and the collapsing spin vector due to dephasing $`\xi _\alpha ^2(t)`$ large. Physically, this is due to the fact that the CSS is evolving into a quasi-self-trapping quantum state with low atom-atom correlations. A correct balance between the rapidly increasing variance and the collapsing spin vector is clearly needed in order to maintain entanglement with $`\xi _\alpha ^2(t)<1`$. From Fig. 5, it is seen that the variance in $`\widehat{J}_x`$ increases most slowly amongst the variances, while Fig. 2 indicates the spin vector $`\widehat{J}_y+\widehat{J}_z`$ oscillating with a non-zero amplitude, showing the most promise in finding entanglement in this spin direction. In Fig. 7 we show $`\xi _x^2(t)`$ for $`q=q_j`$, for $`\mathrm{\Delta }=q_j`$ and $`\mathrm{\Delta }=10q_j`$. Differently from other results, it was found that with $`\mathrm{\Delta }=0`$, $`\xi _x^2(t)1`$ for all times, indicating the crucial role $`\mathrm{\Delta }`$ plays in this system. For the initial Dicke state the parameter quickly becomes large, clearly indicating that quantum correlations between the atoms are destroyed rapidly. However, it is seen that the parameter $`\xi _x^2(t)`$ does dip below 1 for a brief time period demonstrating atom-atom entanglement.
Somewhat counter-intuitively, one may extend the duration over which the atoms are entangled by reducing the nonlinearity constant $`q`$; it is clear that with a smaller $`q`$, the variance in $`\widehat{J}_x`$ increases more gently while the collapse of the spin vector occurs over a longer time scale, extending the duration over which $`\xi _x^2<1`$. This is presented in Fig. 8 for $`q=q_j/10`$. As before, the increase in $`\mathrm{\Delta }`$ results in the increase in oscillation frequency. This indicates a potential for the quantum control of entanglement properties of a TBEC via externally adjustable parameters.
#### III.3.2 Entanglement between the two modes
Writing the quantum state of a TBEC in the form
$`|\mathrm{\Psi }(t)`$ $`=`$ $`{\displaystyle \underset{m=0}{\overset{N}{}}}c_m(t)|N|Nm`$ (36)
$`=`$ $`{\displaystyle \underset{m=j}{\overset{j}{}}}c_m(t)|j,m`$ (37)
the well-known von Neumann entropy that shows degree of entanglement can be writtenMilburn :
$$E(t)=\frac{1}{\mathrm{log}_2(N+1)}\underset{m=j}{\overset{j}{}}|c_m(t)|^2\mathrm{log}_2|c_m(t)|^2$$
(38)
where the normalization factor $`\mathrm{log}_2(N+1)^1`$ was included so that $`0<E<1`$. The expansion coefficients are given by
$`c_m(t)`$ $`=`$ $`m|\widehat{U}(t)|\psi (0)`$ (39)
$`=`$ $`m|\widehat{R}^{}(t)e^{iH^{}t}\widehat{R}(0)|\theta ,\varphi ,`$ (40)
where the rotation operator $`\widehat{R}`$, reduced Hamiltonian $`H^{}`$ and the CSS $`|\theta ,\varphi `$ are as defined above. Using Eq. (28), and inserting the completeness relation $`_m^{}|m^{}m^{}|=๐`$ one can write
$`c_m(t)`$ $`=`$ $`{\displaystyle \underset{m^{}}{}}_{m,m^{}}^{(\lambda )}(t)_m^{}^j(\theta \pi /2,\varphi \mathrm{\Delta }t)`$ (41)
$`\times e^{i[2m^{}\mathrm{\Omega }tqm^2t/2]},`$
where
$`_{m,m^{}}^{(\lambda )}(t)`$ $`=`$ $`m|\widehat{R}^{}(t)|m^{}`$ (42)
represents the matrix element of the rotation operator. This can be evaluated by first applying the disentangling theorem on the rotation operator $`\widehat{R}`$bloch :
$`\widehat{R}(\lambda ,\mathrm{\Delta }t)`$ $`=`$ $`\mathrm{exp}\left[{\displaystyle \frac{\lambda }{2}}(\widehat{J}_{}e^{i\mathrm{\Delta }t}\widehat{J}_+e^{i\mathrm{\Delta }t})\right]`$ (43)
$`=`$ $`e^{\tau \widehat{J}_+}e^{\mathrm{ln}(1+|\tau |^2)\widehat{J}_z}e^{\tau ^{}\widehat{J}_{}},`$ (44)
where $`\tau =e^{i\mathrm{\Delta }t}\mathrm{tan}(\frac{\lambda }{2})`$. Taking into account the unitary nature of the rotation operator$`\widehat{R}^{}(\theta ,\varphi )=\widehat{R}(\theta ,\varphi )`$ and expanding the exponential operators containing $`\widehat{J}_+`$ and $`\widehat{J}_{}`$ one can write
$`_{m,m^{}}^{(\lambda )}(t)`$ $`=`$ $`{\displaystyle \underset{n,n^{}}{}}{\displaystyle \frac{(\tau )^n}{n!}}{\displaystyle \frac{\tau ^n^{}}{n^{}!}}(1+|\tau |^2)^{m^{}n^{}}`$ (45)
$`\times m|\widehat{J}_+^n\widehat{J}_{}^n^{}|m^{}.`$
The final expression for the matrix element, which was obtained by using the ladder operator nature of $`\widehat{J}_\pm `$ and found to be identical to the irreducible representation of full rotation group bloch is:
$`_{m,m^{}}^{(\lambda )}(t)`$ $`=`$ $`{\displaystyle \underset{n=0}{\overset{2j}{}}}{\displaystyle \frac{(1)^{m^{}m+n}}{(n+m^{}m)!n!}}{\displaystyle \frac{(jm+n)!}{(j+mn)!}}`$ (46)
$`\times \left[{\displaystyle \frac{(j+m)!(j+m^{})!}{(jm)!(jm^{})!}}\right]^{1/2}e^{i(mm^{})\mathrm{\Delta }t}`$
$`\times \mathrm{sin}^{2n+m^{}m}\left({\displaystyle \frac{\lambda }{2}}\right)\mathrm{cos}^{m^{}m}\left({\displaystyle \frac{\lambda }{2}}\right).`$
The normalized entanglement parameter $`E(t)`$ is plotted in Fig. 9, for the four combinations of $`q`$ and $`\mathrm{\Delta }`$: $`q=0`$, $`q=q_j`$ and $`\mathrm{\Delta }=q_j`$, $`\mathrm{\Delta }=10q_j`$. It is seen that the nonlinearity is crucial for achieving high degree of entanglement. The effect of $`\mathrm{\Delta }`$ was, as expected, to increase the frequency of oscillations. To generate Fig. 9, we used Eqs. (38), (41), and (46) with $`N=40`$ due to the computational limitations in numerically calculating factorials of large numbers. However, it was found that, with non-zero $`q`$, as long as the ratio of the various parameters such as $`q/N`$ is kept the same, there are no discernible changes in the plot as a function of $`N`$. Our result agrees with the numerical result obtained by Tonel et al.tonel with higher number of atoms.
It was found that, for the initial phase state $`|\theta =\pi /2,\varphi =0`$, the entanglement parameter does depend on the number of atoms, as the normalized coefficients $`|c_m(t)|^2`$ are given in this case by the binomial distribution, namely, $`|c_m(t)|^2C_{j+m}^{2j}`$. It is also notable that the coefficients are no longer time dependent. As the number of particles $`N`$ and hence $`j`$ increases, the broad binomial distribution approximates narrower and narrower Gaussian as a function of $`m`$, and hence the state becomes less and less entangled. For example, for $`N=400`$, it is found that the scaled $`E=0.62`$. In Fig. 10, we plot the entanglement parameter as a function of the number of atoms $`N`$ for the initial phase state. It is seen that the value tends towards 0.6 as the number is increased up to 1000 atoms. Such consistent behavior exhibited for the bipartite entanglement suggests a possibility to use this property to create a macroscopic matter-wave state with known entanglement.
## IV Conclusion
We have studied the dynamical evolution of a TBEC in the presence of an adiabatically varying, off-resonant atom-light coupling. The macroscopic, fictitious spin vector was found to undergo a rather complex motion, and further complications in the form of collapses and revivals were found with nonzero nonlinearity. As the main result of this paper, we used an exact solution to the time-dependent Schrรถdinger equation to calculate analytically the amount of quantum mechanical squeezing and entanglement in a TBEC under various conditions. In particular, we considered entanglement generated between the atoms and that between the modes. For the case with the nonlinearity turned off, although the variance was found to be reduced below the SQL (i.e. squeezed), the system never demonstrated atom-atom entanglement, and the entanglement between the two modes remained well below its maximal value. With the nonlinearity turned on, it was found that the atom-atom entanglement can be generated initially, and then become unentangled rapidly. On the other hand, the entanglement between the two modes was found to more or less maintain maximal values throughout, except for some fluctuations. The dynamics of entanglement was found to be controllable via various parameters present in this system, namely the nonlinearity $`q`$ and the detuning $`\mathrm{\Delta }`$.
Potentially useful results of this work are the squeezing below the SQL in the atom number difference and relative phase which may help reduce the projection noise in spectroscopic and interferometric applications as discussed by Wineland et al.wineland , and the consistently high degree of entanglement between the two modes with the nonlinearity turned on, which could be useful in the context of quantum information science. In addition, the dynamically stable features of the self-trapping state and the time-independent entanglement between the two modes which is also relatively insensitive to the changes in the atom number (for large $`N`$) could prove useful in designing a robust, macroscopic bipartite quantum state with known entanglement. Future work could involve studying quantum control methods for maintaining optimal squeezing and entanglement in a TBEC, for such applications as precision matter-wave interferometry using entangled BECsinterferometry .
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# Coupled-mode theory for periodic side-coupled microcavity and photonic crystal structures
## I Introduction
In the past several years the linear and nonlinear properties of side-coupled waveguiding structures have attracted the attention of many researchers LittlePaper <sup>-</sup>WaksPaper . These structures consist of one or more waveguiding elements in which forward and backward propagating waves are indirectly coupled to each other via one or more mediating resonant cavities. Perhaps the most common proposals for realizing these structures involve photonic crystal (PC) waveguides with defect modes slightly displaced from the waveguiding region (Fig. 1a, left) HausPaper1 <sup>,</sup>YarivPaper1 , or micro-ring resonator structures in which two channel waveguides are side-coupled to micro-ring resonators (Fig. 1a, right) InvitedPaper . In the PC structure the forward and backward propagating modes within the waveguide are coupled via the defect; for the micro-ring structure, the forward going mode in the lower (upper) channel waveguide is coupled, via the micro-ring, to the backward going mode in the upper (lower) channel. The linear and nonlinear properties of both types of structures have been studied HausPaper1 <sup>,</sup>YarivPaper1 <sup>-</sup>SoljacicPaper .
The electromagnetic properties of these structures can be accurately determined in great detail using numerically intensive methods such as finite-difference time-domain (FDTD) simulations FDTDbook . An analysis in terms of Wannier functions can substantially reduce computation time for the PC structure KurtPaper1 , but the numerical problem remains daunting. In particular, full FDTD calculations of the micro-ring structures have to date been confined to two-dimensional analogs of the actual structures of interest FDTDbook . Furthermore, direct numerical simulation, while valuable for design purposes, offers little insight into the physics of the structures. Consequently, semi-analytical techniques, such as the scattering-matrix approach of S. Fan et al.HausPaper1 and Yong Xu et al.YarivPaper1 , have been proposed. Using these techniques the optical properties of side-coupled structures can be understood in terms of the interactions between a small number of modes.
In this paper we concentrate our attention on periodic, side coupled structures (Fig. 1b). Our primary objective is to derive coupled mode equations (CME) that describe pulse propagation in such structures. Coupled mode theory has long been used as an effective design tool for grating structures where forward and backward propagating waves are directly coupled via an index grating YarivQEBook . In directly coupled structures, it is well known that a Bragg gap opens in the dispersion relation of the structure when the phase accumulated in one round trip through a period of the grating is an integer multiple of $`2\pi `$, so that the slight reflections that are incurred due to the grating are coherently enhanced. Structures possessing a Bragg gap have found a variety of uses, such as dispersion compensation EggletonPaper and wavelength division multiplexing GilesPaper . In the side-coupled structure the Bragg feedback mechanism, and hence the Bragg gap, does exist, although it is now mediated by the coupling cavity. However, there is also a second type of gap: a resonator gap, which is associated with the resonance frequencies - and therefore the geometry - of the mediating cavity. For the micro-ring resonator structure the interpretation of this gap is straightforward: when the phase accumulated in a round-trip through the micro-ring resonator is an integer multiple of $`2\pi `$, then the coupling between the forward and backward going waves is resonantly enhanced. Of these two gaps, the resonator gap is perhaps the more important, because it exhibits a deep transmission dip seen even in a structure with only one unit cell.
Because side-coupled structures exhibit both Bragg and resonator gaps, it is to be expected that a CME description of optical pulse propagation will be more complicated than in Bragg gratings. The CME for Bragg gratings involve two fields (forward and backward going) interacting via a coupling coefficient. For side-coupled structures, the most interesting situation is when a resonator gap lies near one of the Bragg gaps, and we show in this paper that the relevant CME then involves three fields: a cavity field and forward and backward going fields.
We derive our CME using a phenomenological Hamiltonian approach, which distills the essential physical interactions of the structure, and hence provides a simple physical picture of optical interactions. We build the fields in our CME as Fourier superpositions of the modes in the Hamiltonian. Hence, our CME are derived for infinite, periodic structures in which the coupling to each cavity is the same. Nevertheless, we show that our CME can be generalized to describe finite, apodized structures, in which the coupling (but not the period) varies from cavity to cavity. Therefore, the CME can be used to describe finite structures with only a small number of cavities. Indeed, the general Hamiltonian approach we advocate can be applied even to structures with only one or two cavities, if the formalism we introduce in Sec. II is extended to a discrete number of (not necessarily identical) cavities. In both discrete and periodic scenarios, the Hamiltonian approach exhibits the similarities of the optical dynamics of these artificially structured materials to more traditional problems in solid state physics. As well, it allows for an easy quantization of the description to address the quantum optics of these structures. We plan to turn to this, as well as the direct derivation of our phenomenological Hamiltonian from the underlying electrodynamics, in future publications.
The present paper is organized as follows. In Sec. II we describe the Hamiltonian model for a system with a single microresonator, investigate the transmission/reflection spectrum of the structure, and indicate how the parameters in our phenomenological Hamiltonian can be set from more common models of cavity resonators. In Sec. III we discuss how the Hamiltonian can be used to model a periodic waveguide-resonator structure. We then discuss methods of reducing the number of fields and interactions in our Hamiltonian while retaining the basic physics. In Sec. IV we derive the coupled mode equations in terms of effective fields built as Fourier superpositions of the modes in the Hamiltonian of Sec. III, and we show how to modify these CME to describe finite, apodized structures. In Sec. V we conclude.
## II Hamiltonian model and transmission for a single cavity structure
In this section we construct a Hamiltonian model for a structure in which forward and backward propagating waves are indirectly coupled to each other via a cavity centred at $`z=z_0`$. We will focus on classical optics here, but because its easy generalization to quantum optics is one of the strengths of this approach, we adopt a quantum notation and, for the classical Poisson bracket $`\{..,..\}`$, we write $`(i\mathrm{})^1[..,..]`$; we also use to indicate complex conjugation. We will also often speak of operators rather than variables, especially when it makes the physics more clear. For example, we introduce $`a_k^{}`$ and $`c_k^{}`$ as creation operators for photons propagating with wavenumber $`k`$ in the forward and backward direction respectively. Because $`k>0`$ ($`k<0`$) indicates that the photons are propagating in the forward (backward) direction, $`a_k^{}`$ exists for $`k>0`$ and $`c_k^{}`$ for $`k<0`$. For a given $`k`$, the energy in these fields is $`\mathrm{}\omega _ka_k^{}a_k`$ and $`\mathrm{}\omega _kc_k^{}c_k`$, with $`\omega _k=c\left|k\right|/n`$, where $`c`$ is the speed of light in a vacuum, and $`n`$ is a constant effective index, equal for the forward and backward propagating waves. By ignoring the frequency dependence of $`n`$ we are neglecting the underlying material dispersion within the waveguides; we discuss the validity of this approximation after equation (5) below. To describe light in the cavity, we define a creation operator $`b^{}`$, and identify the energy in the field as $`\mathrm{}\omega _0b^{}b`$, where $`\omega _0`$ is the resonant frequency of the cavity. For the micro-ring resonator structure of Fig. 1a (right), the $`a_k^{}`$ and $`c_k^{}`$ could represent creation operators for light propagating in the forward direction in the lower waveguide and the backward direction in the upper waveguide, while $`b^{}`$ could represent the field circulating in the counter-clockwise direction in the micro-ring resonator. Our notation implies that the two waveguides have a common mode index $`n`$, but this could easily be generalized. For the PC structure of Fig. 1a (left), the $`a_k^{}`$ and $`c_k^{}`$ would represent creation operators for light propagating in the forward and backward direction in a waveguide mode of the PC waveguide, and $`b^{}`$ would represent the creation operator for the field inside the single mode defect. Regardless of their interpretation, the operators satisfy the commutation relations
$`[a_k,a_k^{}^{}]`$ $`=`$ $`\delta \left(kk^{}\right),`$
$`[c_k,c_k^{}^{}]`$ $`=`$ $`\delta \left(kk^{}\right),`$
$`[b,b^{}]`$ $`=`$ $`1,`$ (1)
with all other commutation relations vanishing. Assuming that no light couples directly between the propagating modes governed by $`a_k^{}`$ and $`c_k^{}`$, but that light can couple from these modes to the cavity, we use the following model Hamiltonian for the system HausPaper1 <sup>,</sup>YarivPaper1 :
$$H=H_o+H_{coupling},$$
(2)
where
$`H_o`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}๐k\mathrm{}\omega _ka_k^{}a_k+{\displaystyle _{\mathrm{}}^0}๐k\mathrm{}\omega _kc_k^{}c_k+\mathrm{}\omega _ob^{}b,`$
$`H_{coupling}`$ $`=`$ $`\mathrm{}{\displaystyle _0^{\mathrm{}}}\xi _k\left[a_k^{}be^{ikz_0}+b^{}a_ke^{ikz_0}\right]๐k`$
$`\left(1\right)^q\mathrm{}{\displaystyle _{\mathrm{}}^0}\xi _k\left[c_k^{}be^{ikz_0}+b^{}c_ke^{ikz_0}\right]๐k.`$
The quantities $`\xi _k`$ and $`\left(1\right)^q\xi _k`$ characterize the strength of the coupling between cavity field and waveguide fields, propagating in the forward and backward direction; $`q`$ is an integer that depends on the symmetry of the cavity mode YarivPaper1 . Note that except for the factor $`(1)^q`$ our notation implies that the coupling to forward and backward propagating waveguide modes is identical. In the micro-ring structure, for example, this means that we assume equal coupling to the two waveguides; generalization of this is straightforward, but for simplicity we will not do it here. The time evolution of the operators is given by the Heisenberg equations of motion
$$i\mathrm{}\frac{dO}{dt}=[O,H],$$
(5)
where $`O`$ is any operator.
In writing down (2), (II) and (LABEL:2\_Ham\_eqn3) we have implicitly assumed that the cavity supports only one mode, with resonant frequency $`\omega _0`$, and that the waveguides guide light in only a single spatial mode profile. Strictly speaking, of course, neither of these assumptions is valid. In general, cavities support more than one mode, oscillating at one or more resonance frequencies, and for sufficiently high frequencies a waveguide will support multiple transverse modes. However, we are primarily interested in the physics of these structures for frequencies at or near a specific resonant frequency $`\omega _0`$. We then assume that within this frequency range only one resonance of the cavity exists or, alternatively, that only a single mode of a multi-mode cavity is excited, and that the waveguides of the structure are single mode. Furthermore, we assume that the underlying material or modal dispersion of the structure is negligible within the frequency range of interest. For our purposes, the inclusion of material dispersion would lead to quantitative, but not qualitative changes.
In Appendix $`1`$ we show that our Hamiltonian formulation leads to a Lorentzian transmission and reflection across the cavity for frequencies in the vicinity of $`\omega _0`$:
$`t\left(\omega \right)`$ $``$ $`{\displaystyle \frac{i\mathrm{\Delta }}{\gamma i\mathrm{\Delta }}},`$ (6)
$`r\left(\omega \right)`$ $``$ $`\left(1\right)^q\left({\displaystyle \frac{\gamma }{\gamma i\mathrm{\Delta }}}\right),`$ (7)
where $`\gamma =2\pi n\xi _{\stackrel{~}{\omega }_0}^2/c`$, and $`\xi _{\stackrel{~}{\omega }_o}`$ is the coupling coefficient between the cavity and waveguides evaluated at $`k=\stackrel{~}{\omega }_0n\omega _0/c`$, and where $`\mathrm{\Delta }=\left(\omega \omega _0\alpha \left(\omega \right)\right)`$ characterizes the detuning from the renormalized resonance frequency $`\omega _0+\alpha \left(\omega \right)`$ . An expression for the quantity $`\alpha \left(\omega \right)`$ is given in Appendix $`1`$. For our structures of interest $`\alpha \left(\omega \right)`$ is sufficiently small that $`\omega \omega _0\alpha \left(\omega \right)\omega \omega _0`$ to a good approximation.
The transmission and reflection coefficients in (6),(7) are of precisely the form that follows from simple transfer matrix models of resonant cavities or ring resonators YarivPaper1 <sup>,</sup>InvitedPaper . In the latter structure, for example, the coupling of the cavity to the waveguides is described by self-coupling and cross-coupling coefficients $`\sigma `$ and $`\kappa `$ respectively, which in a simple case (where the coupling is assumed to occur at the point of smallest separation) are real and satisfy $`\sigma ^2+\kappa ^2=1`$. Comparing the transmission and reflection coefficients found there with (6),(7), we find that they become equivalent if we put
$$\gamma =\frac{c}{2\pi \overline{n}R}\left(\frac{1\sigma ^2}{\sigma ^2}\right)$$
(8)
where $`\overline{n}`$ and $`R`$ are the effective index and radius of the resonator respectively. Thus if a given resonator is parameterized by $`\sigma `$ and $`\kappa `$, as well of course by the resonance frequency $`\omega _0`$, then relation (8) allows one to determine the effective coupling coefficient $`\xi _{\stackrel{~}{\omega _0}}`$ and thus set what will be, as we will see, the crucial elements in the phenomenological Hamiltonian (2). The appropriate values of $`\sigma `$and $`\kappa `$ for a single resonator could be determined by experiment, or directly calculated from the underlying channel and resonator geometries, as discussed by Waks and Vuckovic WaksPaper .
A typical spectrum for a single cavity structure is shown in Fig. 2. On resonance, the reflection induced by the cavity reaches 100% (albeit only for a single wavelength), and remains significant as long as the detuning, $`\mathrm{\Delta }`$, is on the order of $`\gamma `$. The width of the spectrum is dictated by $`\gamma `$, and the larger the coupling to the cavity, the broader the resonance. In physical terms, this means that as the waveguides are brought closer to the cavity of Fig. $`1a`$, the resonance width increases.
## III Hamiltonian for a periodic structure
We now generalize the single-cavity Hamiltonian to describe a periodic structure, in which the forward and backward propagating modes are coupled to an infinite series of periodically spaced cavities (Fig. 1b). We assume that the resonators are not directly coupled to each other, although of course they do couple indirectly via the waveguides. Generalizing the Hamiltonian (2) to include the periodic sequence of resonators, we write
$`H`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}๐k\mathrm{}\omega _ka_k^{}a_k+{\displaystyle _{\mathrm{}}^0}๐k\mathrm{}\omega _kc_k^{}c_k+{\displaystyle \underset{l}{}}\mathrm{}\omega _ob_l^{}b_l`$
$`\mathrm{}{\displaystyle \underset{l}{}}{\displaystyle _0^{\mathrm{}}}๐k\xi _k\left[b_l^{}a_ke^{ikz_l}+a_k^{}b_le^{ikz_l}\right]`$
$`\left(1\right)^q\mathrm{}{\displaystyle \underset{l}{}}{\displaystyle _{\mathrm{}}^0}\xi _k๐k\left[b_l^{}c_ke^{ikz_l}+c_k^{}b_le^{ikz_l}\right],`$
where $`a_k^{}`$ $`\left(c_k^{}\right)`$ are again the creation operators for light propagating the forward (backward) direction. The main difference between (III) and (2) is that we have now included a countably infinite number of resonators, each with the same resonance frequency, $`\omega _0`$, and associated with the creation operator $`b_l^{}`$, where $`l`$ indexes the resonator. The resonators are evenly spaced at $`z_l=l\mathrm{\Lambda }`$, which gives a fundamental reciprocal lattice vector $`G_0=2\pi /\mathrm{\Lambda }`$. The Hamiltonian (III) can be re-written as
$`H`$ $`=`$ $`{\displaystyle \underset{G}{}}{\displaystyle _{B.Z.}}๐k\mathrm{}\omega _{k+G}a_{k+G}^{}a_{k+G}+{\displaystyle \underset{G}{}}{\displaystyle _{B.Z.}}๐k\mathrm{}\omega _{kG}c_{kG}^{}c_{kG}+{\displaystyle \underset{l}{}}\mathrm{}\omega _ob_l^{}b_l`$ (10)
$`\mathrm{}{\displaystyle \underset{l}{}}{\displaystyle \underset{G}{}}{\displaystyle _{B.Z.}}๐k\xi _{k+G}\left[b_l^{}a_{k+G}e^{i\left(k+G\right)z_l}+a_{k+G}^{}b_le^{i\left(k+G\right)z_l}\right]`$
$`\left(1\right)^q\mathrm{}{\displaystyle \underset{l}{}}{\displaystyle \underset{G}{}}{\displaystyle _{B.Z.}}๐k\xi _{k+G}\left[b_l^{}c_{kG}e^{i\left(kG\right)z_l}+c_{kG}^{}b_le^{i\left(kG\right)z_l}\right],`$
where $`_G`$ represents the summation over an infinite number of positive reciprocal lattice vectors (with $`G=0,G_0,2G_0,\mathrm{}`$), and where in the integrations we restrict the wavenumber $`k`$ to the first Brillouin zone ($`G_0/2<kG_0/2`$); We sum only over the positive reciprocal lattice vectors so that $`a_{k+G}^{}`$ and $`c_{kG}^{}`$ retain their association with forward and backward propagation modes respectively. The operators satisfy commutation relations
$`[a_{k+G},a_{k^{}+G^{}}^{}]`$ $`=`$ $`\delta \left(kk^{}\right)\delta _{G,G^{}},`$
$`[c_{kG},c_{k^{}G^{}}^{}]`$ $`=`$ $`\delta \left(kk^{}\right)\delta _{G,G^{}},`$
$`[b_l,b_l^{}^{}]`$ $`=`$ $`\delta _{l,l^{}},`$ (11)
with all other commutators vanishing; the first two of these follow immediately from (1). Because the system is periodic, we can identify a countably infinite set of Bragg frequencies in (10). These are the frequencies $`\omega _{k\pm G}`$ evaluated at $`k=0`$ or $`G_0/2`$. Hence, since $`\omega _{k\pm G}=c\left|k\pm G\right|/n`$ for ring resonator structures, the $`M^{th}`$ Bragg frequency occurs at $`\omega _b^{(M)}=M\left(cG_0/2n\right)`$ (with $`M0`$ an integer).
To simplify (10), we introduce the collective operator
$$b_k=\sqrt{\frac{\mathrm{\Lambda }}{2\pi }}\underset{l}{}b_le^{ikz_l},$$
(12)
where $`k`$ is now a continuous variable that ranges over the first Brillouin zone. In Appendix 2 we introduce this operator by first considering only excitations of the resonators periodic over a length $`L=N\mathrm{\Lambda }`$, and then taking $`N\mathrm{}`$. We find in that limit
$$\underset{l}{}\mathrm{}\omega _0b_l^{}b_l_{B.Z.}๐k\mathrm{}\omega _ob_k^{}b_k,$$
and that
$$[b_k,b_k^{}^{}]=\delta \left(kk^{}\right)$$
for $`k`$ and $`k^{}`$ in the first Brillouin zone, with all other commutators vanishing. In terms of this collective operator the Hamiltonian (10) becomes
$`H`$ $`=`$ $`{\displaystyle \underset{G}{}}{\displaystyle _{B.Z.}}๐k\mathrm{}\omega _{k+G}a_{k+G}^{}a_{k+G}+{\displaystyle \underset{G}{}}{\displaystyle _{B.Z.}}๐k\mathrm{}\omega _{kG}c_{kG}^{}c_{kG}+{\displaystyle _{B.Z.}}๐k\mathrm{}\omega _ob_k^{}b_k`$ (13)
$`\mathrm{}{\displaystyle \underset{G}{}}{\displaystyle _{B.Z.}}๐k\mathrm{\Xi }_{+k+G}\left[b_k^{}a_{k+G}+a_{k+G}^{}b_k\right]\left(1\right)^q\mathrm{}{\displaystyle \underset{G}{}}{\displaystyle _{B.Z.}}๐k\mathrm{\Xi }_{k+G}\left[b_k^{}c_{kG}+c_{kG}^{}b_k\right],`$
where $`\mathrm{\Xi }_{\pm k\pm G}`$ $`\sqrt{\frac{2\pi }{\mathrm{\Lambda }}}\xi _{\pm k\pm G}`$. In Table I we give typical values for parameters characterizing side-coupled structures, and we use them in our sample calculations below. There and for the rest of this paper we assume that the coupling $`\mathrm{\Xi }_{\pm k\pm G^{}}`$ is approximately constant at wavevectors corresponding to frequencies within our region of interest, and take $`\mathrm{\Xi }_{\pm k\pm G^{}}\mathrm{\Xi }`$. This approximation is reasonable if the $`G`$ of interest satisfy $`G\mathrm{\Delta }k`$, where $`\mathrm{\Delta }k`$ is the range over which the $`\xi _k`$ varies significantly. We can expect $`\mathrm{\Delta }k2\pi /(1\mu m)`$ for the structures of interest (see Appendix 1), and since $`G`$ is at most a few times $`G_0=2\pi /\mathrm{\Lambda }`$ ($`=2\pi /(32\mu m)`$ from Table I), this inequality is indeed satisfied.
The dispersion relation of the system can be determined by traditional transfer matrix methods, using (6), (7) for the transmission and reflection coefficients of a single resonator. However, to see the connection with the coupled mode equations we will derive, we consider determining the dispersion relation directly from the Hamiltonian (13), by applying the Heisenberg equation of motion to generate equations for the time derivatives of $`a_{k+G}`$, $`c_{kG}`$ and $`b_k`$. Assuming harmonic time dependence $`e^{i\omega t}`$ for the operators, we determine an expression for $`\omega `$ as a complicated function of the countably infinite set of $`\omega _{\pm \left|k\right|\pm G}`$, and the discrete value $`\omega _0`$. Alternately (and equivalently) we can exhibit the Hamiltonian in a matrix form (13)
$$H=\mathrm{}_{BZ}๐k๐_k^{}๐ต_k๐_k,$$
(14)
where
$$๐_k^{}=(a_{k+G_0}^{},a_{k+2G_0}^{},\mathrm{},c_{kG_0}^{},c_{k2G_0}^{},\mathrm{},b_k^{}),$$
(15)
and $`๐ต_k`$ contains all of the interactions between the $`a_{k+G}`$, $`c_{kG}`$ and $`b_k`$. Then, by diagonalizing the (infinite-dimensional) matrix $`๐ต_k`$ we can in principle determine the dispersion relation of the structure. In Fig. 3 we consider a typical uncoupled (in the limit where $`\mathrm{\Xi }=0`$) and coupled dispersion relation for the structure. The dotted line shows the uncoupled dispersion relation, and the solid line shows the dispersion relation of the coupled system, as determined by the transfer matrix approach.
If one of the Bragg frequencies is close to the resonant frequency $`\omega _0`$, then we show below that a truncation of the matrix V<sub>k</sub> to three terms is a good approximation. The restricted Hamiltonian that results is
$$H\mathrm{}_{B.Z.}๐k\left[\begin{array}{ccc}a_{k+G^{}}^{}\hfill & c_{kG^{}}^{}\hfill & b_k^{}\hfill \end{array}\right]\left[\begin{array}{ccc}\omega _{k+G^{}}\hfill & 0\hfill & \mathrm{\Xi }\hfill \\ 0\hfill & \omega _{kG^{}}\hfill & \left(1\right)^q\mathrm{\Xi }\hfill \\ \mathrm{\Xi }\hfill & \left(1\right)^q\mathrm{\Xi }\hfill & \omega _0\hfill \end{array}\right]\left[\begin{array}{c}a_{k+G^{}}\hfill \\ c_{kG^{}}\hfill \\ b_k\hfill \end{array}\right].$$
(16)
where $`G^{}`$ is the reciprocal lattice vector associated with the forward (backward) band that has $`\omega _{k+G^{}}`$ $`(\omega _{kG^{}})`$ closest to $`\omega _0`$. Here we have assumed that the resonant frequency is very close to a Bragg frequency with its associated gap at the Brillouin zone centre, and so $`\omega _G^{}=\omega _G^{}\omega _b`$, where $`\omega _b`$ is the Bragg frequency closest to the resonance frequency<sup>8</sup>. We refer to eqn. 16 as the โthree mode model.โ Its validity near a resonance frequency for any particular structure can be formally investigated by including the omitted terms in a multiple scales analysis, or by simply comparing the dispersion relation following from eqn. 16 with a full solution of the dispersion relation using a transfer matrix approach. This is done in Fig. 4, using the parameters in Table I as was done in Fig. 3. In Fig. 4 we also plot the imaginary part of $`k`$ within the gaps. Note that the exact solution and that from the three mode model are in good agreement for the frequency range shown in Fig. 4. Such agreement fails at other Bragg frequencies that are further from the resonant gap, of course, since the three mode model (eqn. 16) only contains the physics of the Bragg gap closest to $`\omega _0`$. It is to frequencies near $`\omega _0`$ that we henceforth restrict ourselves.
## IV Coupled-mode equations in the three-mode model
In this section we derive a set of coupled-mode equations which describe pulse propagation in the periodic structure, based on the three-mode Hamiltonian (16). We then demonstrate that although these coupled mode equations are derived for an infinite periodic system with equal coupling at each resonator, they can, with only slight modifications, be used to describe finite systems with varying coupling at each resonator. We start by defining effective fields in terms of the amplitudes $`a_{k+G^{}}`$, $`c_{kG^{}}`$ and $`b_k`$:
$`g_+(z,t)`$ $`=`$ $`{\displaystyle _{B.Z.}}{\displaystyle \frac{dk}{\sqrt{2\pi }}}a_{k+G^{}}e^{ikz},`$
$`g_{}(z,t)`$ $`=`$ $`{\displaystyle _{B.Z.}}{\displaystyle \frac{dk}{\sqrt{2\pi }}}c_{kG^{}}e^{ikz},`$
$`b(z,t)`$ $`=`$ $`{\displaystyle _{B.Z.}}{\displaystyle \frac{dk}{\sqrt{2\pi }}}b_ke^{ikz}.`$ (17)
where $`G^{}`$ indexes the reciprocal lattice vector that is retained within the three mode approximation. These fields can be interpreted as a forward propagating field, a backward propagating field, and the field distribution in the resonators respectively. Using the definitions in (17), the effective fields satisfy the equal time commutation relations,
$`[g_\pm (z,t),g_\pm ^{}(z^{},t)]`$ $`=`$ $`\widehat{\delta }\left(zz^{}\right)`$
$`[b(z,t),b^{}(z^{},t)]`$ $`=`$ $`\widehat{\delta }\left(zz^{}\right),`$ (18)
with all other commutation relations vanishing. The function $`\widehat{\delta }(zz^{})`$ is an effective delta function such that $`_{\mathrm{}}^{\mathrm{}}f(z)\widehat{\delta }(zz^{})๐z=f(z^{})`$ when the function $`f(z)`$ has its wavenumber restricted to the first Brillouin zone of the system. In terms of the effective fields, the Hamiltonian in (16) becomes
$`H`$ $`=`$ $`\mathrm{}\omega _b{\displaystyle ๐zg_+g_+^{}}+i{\displaystyle \frac{\mathrm{}c}{2n}}{\displaystyle ๐z\left(\frac{g_+^{}}{z}g_+g_+^{}\frac{g_+}{z}\right)}`$ (19)
$`+`$ $`\mathrm{}\omega _b{\displaystyle ๐zg_{}g_{}^{}}i{\displaystyle \frac{\mathrm{}c}{2n}}{\displaystyle ๐z\left(\frac{g_{}^{}}{z}g_{}g_{}^{}\frac{g_{}}{z}\right)}`$
$`+`$ $`\mathrm{}\omega _0{\displaystyle }dzbb^{}\mathrm{}\mathrm{\Xi }{\displaystyle }dz(b^{}g_++c.c.)`$
$``$ $`(1)^q\mathrm{}\mathrm{\Xi }{\displaystyle }dz(b^{}g_{}+c.c.)`$
where $`\omega _b`$ denotes the Bragg frequency centered at the Brillouin zone center and closest to $`\omega _0`$ Bragg\_definition . Using the Heisenberg equations of motion for the effective fields, we obtain the coupled equations
$`\left({\displaystyle \frac{}{t}}+{\displaystyle \frac{c}{n}}{\displaystyle \frac{}{z}}\right)g_+(z,t)`$ $`=`$ $`i\omega _bg_+(z,t)+i\mathrm{\Xi }b(z,t),`$
$`\left({\displaystyle \frac{}{t}}{\displaystyle \frac{c}{n}}{\displaystyle \frac{}{z}}\right)g_{}(z,t)`$ $`=`$ $`i\omega _bg_{}(z,t)+i\left(1\right)^q\mathrm{\Xi }b(z,t),`$
$`{\displaystyle \frac{}{t}}b(z,t)=i\omega _ob(z,t)`$ $`+`$ $`i\mathrm{\Xi }g_+(z,t)+i\left(1\right)^q\mathrm{\Xi }g_{}(z,t).`$
One can obtain the dispersion relation directly from (IV) by assuming that each field is a plane wave $`e^{ikzi\omega t}`$, with $`k`$ restricted to the first Brillouin zone. The results are equivalent to those in Fig.4, obtained by diagonalizing (16).
Although the CME (IV) were derived assuming an infinite medium, they can be used to describe a structure where the coupling constant $`\mathrm{\Xi }`$ varies slowly over a distance on the order of the spacing between the resonators. A multiple scale analysis can be used to identify this limit and corrections to it. A more striking inhomogeneous structure is one beginning with a region where there are no resonators, followed by a length $`L`$ over which resonators are placed with an equal spacing and equal coupling to the channel(s), followed by a region where again there are no resonators. A simple model for such a region would be to use the equations (IV), but replacing $`\mathrm{\Xi }`$ with a position dependent coupling constant $`\left[\theta \left(z\right)\theta \left(zL\right)\right]\mathrm{\Xi },`$ where $`\theta `$ is the usual step function. It can be easily seen that this model formally violates our assumptions. Consider, for example, fields with a stationary time dependence, so $`g_+(z,t)=g_+\left(z\right)e^{i\overline{\omega }t}`$, and similarily for all other fields. Then the first equation gives
$$i\overline{\omega }g_+\left(z\right)+\frac{c}{n}\frac{}{z}g_+\left(z\right)=i\omega _bg_+\left(z\right)+i\left[\theta \left(z\right)\theta \left(zL\right)\right]\mathrm{\Xi }b(z,t),$$
(21)
where in fact the factor $`\left[\theta \left(z\right)\theta \left(zL\right)\right]`$ could be omitted, since the third of (IV) together with the position dependent coupling constant guarantees that $`b\left(z\right)`$ will only be nonzero in the region between $`z=0`$ and $`z=L`$. Note however that at $`z=0`$ and $`z=L`$ the equation (21) leads to a discontinuous $`g_+\left(z\right)/z`$ if it is assumed that $`g_+(z)`$ is everywhere continuous. This violates, of course, the assumption that fields such as $`g_+(z,t)`$ are of the form (17).
Despite such a formal violation of our assumptions, this simple model in fact gives a good description of the optical response of a finite structure. To see this, consider first the fields $`g_\pm (z,t)`$ within the structure. It is clear from (IV) that for a supposed frequency $`\overline{\omega }`$ there are two Bloch wavenumbers, which equivalently follow from (16); they are given by $`k\left(\overline{\omega }\right)=\pm \overline{k}`$, where
$$\overline{k}=\frac{n}{c}\sqrt{\frac{\left(\mathrm{\Delta }_o\mathrm{\Delta }_1\mathrm{\Xi }^2\right)^2\mathrm{\Xi }^4}{\mathrm{\Delta }_0^2}}.$$
(22)
In the equation above $`\mathrm{\Delta }_0=\left(\overline{\omega }\omega _0\right)`$ is the detuning from the resonance frequency and $`\mathrm{\Delta }_1=\left(\overline{\omega }\omega _b\right)`$ is the detuning from the Bragg frequency that lies closest to $`\omega _0`$. As a result, one can write the forward and backward propagating effective fields, $`g_\pm (z,t)`$, as
$`g_\pm (z,t)`$ $`=`$ $`g_\pm \left(z\right)e^{i\overline{\omega }t}`$
$`g_\pm \left(z\right)`$ $`=`$ $`g_\pm ^{\left(1\right)}e^{i\overline{k}z}+g_\pm ^{\left(2\right)}e^{i\overline{k}z},`$ (23)
Once $`g_+^{(1)}`$ and $`g_+^{(2)}`$ are set, $`g_{}^{(1)}`$ and $`g_{}^{(1)}`$ are determined by the dispersion relation, or equivalently (IV). Hence there are only two independent constants. Outside the structure ($`\mathrm{\Xi }=0`$) there are also two independent constants in each of the regions $`z<0`$ and $`z>L`$, but the solution of (IV) is simpler. There it takes the form
$`g_\pm (z,t)`$ $`=`$ $`g_\pm (z)e^{i\overline{\omega }t}`$
$`g_+(z)`$ $`=`$ $`g_+e^{iqz}`$
$`g_{}(z)`$ $`=`$ $`g_{}e^{iqz},`$
where $`g_+`$, $`g_{}`$ are independent and $`q=\overline{\omega }n/c`$. For $`z<0`$ we denote the constants by $`g_+^<`$ and $`g_{}^<`$, and for $`z>L`$ we denote them by $`g_+^>`$ and $`g_{}^>`$. Now we consider the boundary condition at $`z=L`$, and note that since no field is incident from $`z>L`$, we have $`g_{}^>=0`$; an incident field is specified by $`g_+^<`$. Our independent unknowns are then $`g_{}^<,g_+^>`$, and the constants $`g_+^{(1)}`$ and $`g_+^{(2)}`$ that specify the field in the structure. We solve for these four unknowns by requiring the continuity of $`g_\pm (z)`$ at $`z=0`$ and $`z=L`$. The resulting transmittance of the structure can be written as
$$T\left(\omega \right)=\left|\frac{g_+^>e^{iqL}}{g_+^{(1)}+g_+^{(2)}}\right|^2$$
(24)
with
$`g_+^{\left(1\right)}`$ $`=`$ $`{\displaystyle \frac{e^{i\overline{k}L}}{2}}\left[1+{\displaystyle \frac{\mathrm{\Xi }^2}{\overline{k}\mathrm{\Delta }_o}}\left({\displaystyle \frac{\mathrm{\Delta }_o\mathrm{\Delta }_1}{\mathrm{\Xi }^2}}1\right)\right]g_+^>e^{iqL},`$
$`g_+^{\left(2\right)}`$ $`=`$ $`{\displaystyle \frac{e^{i\overline{k}L}}{2}}\left[1{\displaystyle \frac{\mathrm{\Xi }^2}{\overline{k}\mathrm{\Delta }_o}}\left({\displaystyle \frac{\mathrm{\Delta }_o\mathrm{\Delta }_1}{\mathrm{\Xi }^2}}1\right)\right]g_+^>e^{iqL}.`$
In Fig. 5 we compare the transmission spectrum of a two channel micro-ring resonator structure with 30 cavities, calculated both using the transfer matrix technique,<sup>7</sup> and using the coupled mode equation result eqn. 24. Again we adopt the parameters of Table I. Generally there is good qualitative agreement, with the main features of the spectrum well described by the coupled mode equation result (24), although as noted above it is being applied beyond its strict range of applicability. An extension of this approach leads to the use of the CME (20) to treat a finite structure where the coupling constant $`\mathrm{\Xi }`$ varies from one resonator to the next. To describe this we simply allow $`\mathrm{\Xi }`$ in (20) to adopt a $`z`$-dependence,
$`\left({\displaystyle \frac{}{t}}+{\displaystyle \frac{c}{n}}{\displaystyle \frac{}{z}}\right)g_+(z,t)`$ $`=`$ $`i\omega _bg_+(z,t)+i\mathrm{\Xi }\left(z\right)b(z,t),`$
$`\left({\displaystyle \frac{}{t}}{\displaystyle \frac{c}{n}}{\displaystyle \frac{}{z}}\right)g_{}(z,t)`$ $`=`$ $`i\omega _bg_{}(z,t)+i\left(1\right)^q\mathrm{\Xi }\left(z\right)b(z,t),`$
$`{\displaystyle \frac{}{t}}b(z,t)=i\omega _0b(z,t)`$ $`+`$ $`i\mathrm{\Xi }\left(z\right)g_+(z,t)+i\left(1\right)^q\mathrm{\Xi }\left(z\right)g_{}(z,t).`$ (25)
In Fig. 6 we plot the transmission spectrum for a 5 cavity structure apodized such that the cavities (from left to right) are characterized by coupling constants $`(\sigma _1,\mathrm{}\sigma _5)=(0.993,0.986,0.98,0.986,0.993)`$, corresponding to $`(\mathrm{\Xi }_1\mathrm{\Lambda }n/\pi c,..\mathrm{\Xi }_5\mathrm{\Lambda }n/\pi c)=(0.0208,0.0287,0.0351,0.0287,0.0208)`$. The transfer matrix results is presented, as well as a very simple application of the CME (25) using a piecewise uniform function to represent $`\mathrm{\Xi }`$, where in the $`n^{th}`$ unit well we set $`\mathrm{\Xi }=\mathrm{\Xi }_n`$. Again there is good qualitative agreement, although the CME are being applied beyond their strict range of applicability. Besides the difference between the CME and transfer matrix results with respect to the Fabry-Perot type oscillations, as seen in Fig. 5, here the CME solution also consistently overestimates the transmission on the high-frequency side of the stop gap. This can be traced back to the effects on the band curvature induced by the next highest Bragg gap, which are implicitly included in the transfer matrix solution but not in the CME calculation.
Finally, we note that while at least three coupled mode equations are necessary to describe the kind of structures we consider here if we deal with both their space and time dependence, if we instead restrict ourselves to a stationary time dependence, $`g_\pm (z,t)=g_\pm \left(z\right)e^{i\overline{\omega }t}`$ and $`b(z,t)=b(z)e^{i\overline{\omega }t}`$, then in fact we can eliminate the variable $`b(z,t)`$ and construct coupled mode equations involving only $`g_+(z,t)`$ and $`g_{}(z,t)`$. They are
$`{\displaystyle \frac{}{z}}g_+\left(z\right)`$ $`=`$ $`i\nu \left(\omega \right)g_+\left(z\right)+i\left(1\right)^q\mu \left(\omega \right)g_{}\left(z\right),`$
$`{\displaystyle \frac{}{z}}g_{}\left(z\right)`$ $`=`$ $`i\nu \left(\omega \right)g_{}\left(z\right)i\left(1\right)^q\mu \left(\omega \right)g_+\left(z\right),`$
where
$`\nu \left(\omega \right)`$ $`=`$ $`{\displaystyle \frac{n}{c}}\left[{\displaystyle \frac{\mathrm{\Xi }^2}{\left(\omega _0\omega \right)}}\left(\omega _b\omega \right)\right],`$
$`\mu \left(\omega \right)`$ $`=`$ $`{\displaystyle \frac{n}{c}}{\displaystyle \frac{\mathrm{\Xi }^2}{\left(\omega _0\omega \right)}},`$ (27)
These equations are valid for $`\omega \omega _0`$. It is well-known that a photonic band gap opens in the dispersion relation described by these equations when $`\left|\mu \left(\omega \right)\right|\left|\nu \left(\omega \right)\right|`$, ProgInOptics and that the width of the gap is larger for larger values of $`\left|\mu \left(\omega \right)\right|`$. Consequently we see from these equations an analytic confirmation of features that our dispersion relation display. Within our three mode model, one edge of the resonator gap occurs at $`\omega \omega _0`$ (in which case $`\nu `$ and $`\mu `$ both diverge equally quickly and are hence equal in the limit as $`\omega `$ approaches $`\omega _0`$), and one edge of the Bragg gap occurs at $`\omega \omega _b`$, because then the second term in the expression for $`\nu \left(\omega \right)`$ vanishes, and $`\nu \left(\omega _b\right)=\mu \left(\omega _b\right)`$.
## V Conclusion
We have presented a phenomenological Hamiltonian description of light propagation in side-coupled resonators. This formulation is appealing in its simplicity, since it captures the basic physics of the structures via a set of readily understandably parameters. The most interesting special case is perhaps where a resonator gap is close to a Bragg gap, and at frequencies close to these gaps a three mode model gives a good description of the dynamics of a periodic structure of resonators. Coupled mode equations based on these captures the dispersion relation even deep within the gaps, and a naive extension of these equations to describe finite structures, although not within the strict range of applicability of the model, gives a good qualitative description.
A hallmark of the kind of approach we have taken here is the connection of theoretically calculated or experimentally observed parameters, such as the coupling coefficient $`\sigma `$, to the parameters that appear in our phenomenological Hamiltonian. Such a strategy is particularly amenable to the description of quantum and nonlinear optical effects. The Hamiltonian description leads to straightforward quantization, of course, and appropriate nonlinear terms can easily be added to the Hamiltonian. In a previous study by Grimshaw et al.<sup>22</sup>, it was shown that three nonlinear coupled mode equations support stationary solitary wave solution in the presence of Kerr nonlinearity. Numerical studies have indicated that soliton-like waves exist in resonator structures. In future work we plan to apply the approach we have detailed here to study such field excitations, where a Hamiltonian framework provides the ability to characterize conserved quantities in terms of the symmetries of the nonlinear field theory.
## VI Acknowledgments
This project was partly funded by the Natural Science and Engineering Research Council (NSERC) of Canada. Philip Chak acknowledges financial support from Photonic Research Ontario and an Ontario Graduate Scholarship.
## VII Appendix 1
In this appendix we use the Hamiltonian (2) to determine the transmission properties of a single-cavity structure. These transmission properties have been intensively studied using various methods such as finite difference, time domain simulationsFDTDbook , and scattering matrix techniquesHausPaper1 YarivPaper1 , and it is well-known that a Lorentzian function gives an excellent approximation to the response of the structure. Here we show that our Hamiltonian also leads to a Lorentzian spectrum. To discuss transmission and reflection, we assume that there is a time-dependent source, $`u\left(t\right)`$, coupled to the forward propagating modes at $`z_s<z_0`$. We therefore modify the Hamiltonian (2) to include a source term:
$$H=H_o+H_{coupling}+H_{source},$$
(28)
with
$$H_{source}=\mathrm{}_0^{\mathrm{}}\left[a_k^{}u\left(t\right)e^{ikz_s}+a_ku^{}\left(t\right)e^{ikz_s}\right]๐k,$$
(29)
where $`e^{ikz_s}`$ accounts for the fact that the light is generated at $`z=z_s`$. Using the Hamiltonian (28) and the commutation relations (1) in the Heisenberg equations of motion (5) we find
$`a_k(t)`$ $`=`$ $`i\xi _k{\displaystyle _{\mathrm{}}^t}b\left(t^{}\right)e^{i\omega _k\left(tt^{}\right)}e^{ikz_o}๐t^{}`$
$`+`$ $`i{\displaystyle _{\mathrm{}}^t}u\left(t^{}\right)e^{i\omega _k\left(tt^{}\right)}e^{ikz_s}๐t^{},`$
$`c_k(t)`$ $`=`$ $`i\left(1\right)^q\xi _k{\displaystyle _{\mathrm{}}^t}b\left(t^{}\right)e^{i\omega _k\left(tt^{}\right)}e^{ikz_o}๐t^{},`$
$`{\displaystyle \frac{db(t)}{dt}}`$ $`=`$ $`i\omega _ob(t)+i{\displaystyle _0^{\mathrm{}}}\xi _ka_k(t)e^{ikz_o}๐k`$ (30)
$`+`$ $`i\left(1\right)^q{\displaystyle _{\mathrm{}}^0}\xi _kc_k(t)e^{ikz_o}๐k.`$
where we have formally integrated the Heisenberg equations for $`da_k/dt`$ and $`dc_k/dt`$, so that both $`a_k\left(t\right)`$ and $`c_k\left(t\right)`$ are expressed entirely in terms of $`b\left(t\right)`$ and $`u\left(t\right)`$. Using the expressions for $`a_k\left(t\right)`$ and $`c_k\left(t\right)`$ in the equation for $`db/dt`$, and expanding $`b\left(t\right)`$ and $`u\left(t\right)`$ in terms of Fourier components,
$`b\left(t\right)`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}b\left(\omega \right)e^{i\omega t}๐\omega ,`$
$`u\left(t\right)`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}u\left(\omega \right)e^{i\omega t}๐\omega ,`$ (31)
we obtain
$$b\left(\omega \right)=\left[\frac{2\pi n\xi _{\stackrel{~}{\omega }}/c}{2\pi n\xi _{\stackrel{~}{\omega }}^2/ci\mathrm{\Delta }}\right]u\left(\omega \right)e^{i\stackrel{~}{\omega }\left(z_oz_s\right)},$$
(32)
where $`\mathrm{\Delta }=\left(\omega \omega _o+\alpha \left(\omega \right)\right)`$ and $`\stackrel{~}{\omega }=\omega n/c`$, with
$`\alpha \left(\omega \right)=2{\displaystyle _0^{\mathrm{}}}\mathrm{}\left({\displaystyle \frac{\xi _k^2}{\frac{c}{n}k\omega }}\right)๐k`$ (33)
describing the small shift in the resonance frequency of the cavity due to the presence of the waveguide. To estimate the effect of $`\alpha (\omega )`$, we assume $`\xi _k`$ takes a gaussian form in k space with a peak centered at $`k=\stackrel{~}{\omega _0}`$. We take the width of the gaussian profile to be about $`1\mu m^1`$, associated with a typical length over which the coupling between the waveguide and resonator is significant. Using this approximate form for $`\xi _k`$ in the expression for $`\alpha (\omega )`$ and numerically evaluating the integral, we have verified that $`\alpha (\omega )`$ is much smaller than the resonance frequency $`\omega _0`$ for structures of interest. Note that in (32) we have switched our notation for wavenumber from $`k`$ to $`\stackrel{~}{\omega }=n\omega /c=\left|k\right|`$ to stress that we are now considering the frequency response of the structure. To determine the transmission and reflection spectrum of the structure we define a set of effective fields
$`f_+(z,t)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2\pi }}}{\displaystyle _0^{\mathrm{}}}๐ka_k\left(t\right)e^{ikz},`$
$`f_{}(z,t)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2\pi }}}{\displaystyle _{\mathrm{}}^0}๐kc_k\left(t\right)e^{ikz}.`$ (34)
We then substitute the values (30) for $`a_k\left(t\right)`$ and $`c_k\left(t\right)`$ in the effective fields (34), and use the Fourier transforms (31) of $`b\left(t\right)`$ and $`u\left(t\right)`$ to simplify the integrals. We are specifically interested in the following two quantities
$`\underset{z\mathrm{}}{lim}f_+(z,t)`$ $`=`$ $`{\displaystyle \frac{i}{c}}{\displaystyle _0^{\mathrm{}}}\left[{\displaystyle \frac{i\mathrm{\Delta }}{2\pi n\xi _{\stackrel{~}{\omega }}^2/ci\mathrm{\Delta }}}\right]u\left(\omega \right)e^{ik\left(zz_s\right)}e^{i\omega t}๐\omega ,`$ (35)
$`\underset{z\mathrm{}}{lim}f_{}(z,t)`$ $`=`$ $`{\displaystyle \frac{i}{c}}\left(1\right)^q{\displaystyle _0^{\mathrm{}}}\left[{\displaystyle \frac{2\pi n\xi _{\stackrel{~}{\omega }}^2/c}{2\pi n\xi _{\stackrel{~}{\omega }}^2/ci\mathrm{\Delta }}}\right]u\left(\omega \right)e^{ik\left(z+z_s\right)}e^{i2\stackrel{~}{\omega }z_o}e^{i\omega t}๐\omega .`$
Note that in the absence of coupling we would have
$`\underset{z\mathrm{}}{lim}f_+(z,t)`$ $`=`$ $`{\displaystyle \frac{i}{c}}{\displaystyle _0^{\mathrm{}}}u\left(\omega \right)e^{ik\left(zz_s\right)}e^{i\omega t}๐\omega ,`$ (36)
$`\underset{z\mathrm{}}{lim}f_{}(z,t)`$ $`=`$ $`0`$
The first (second) of the expressions in (35) is the transmitted (reflected) field built as a superposition of the Fourier components of the source term, $`u\left(\omega \right)`$. We can therefore define the transmission and reflection coefficients as
$`t\left(\omega \right)`$ $`=`$ $`{\displaystyle \frac{i\mathrm{\Delta }}{2\pi n\xi _{\stackrel{~}{\omega }}^2/ci\mathrm{\Delta }}},`$
$`r\left(\omega \right)`$ $`=`$ $`\left(1\right)^q{\displaystyle \frac{2\pi n\xi _{\stackrel{~}{\omega }}^2/c}{2\pi n\xi _{\stackrel{~}{\omega }}^2/ci\mathrm{\Delta }}}.`$
From these coefficients, it is clear that the cavity affects the transmission/reflection of the structure when the detuning, $`\mathrm{\Delta }`$, is on the order of $`2\pi n\xi _{\stackrel{~}{\omega }}^2/c`$. In the limit of very weak coupling โ that is, when the value of $`2\pi n\xi _{\stackrel{~}{\omega }}^2/c`$ is approximately constant over a frequency range centered at $`\omega _0`$ and spanning several multiples of $`2\pi n\xi _{\stackrel{~}{\omega }}^2/c`$, then the transmission and reflection are well approximated by a Lorentzian lineshape
$`t\left(\omega \right)`$ $``$ $`{\displaystyle \frac{i\mathrm{\Delta }}{\gamma i\mathrm{\Delta }}},`$ (37)
$`r\left(\omega \right)`$ $``$ $`\left(1\right)^q\left({\displaystyle \frac{\gamma }{\gamma i\mathrm{\Delta }}}\right),`$ (38)
where $`\gamma 2\pi n\xi _{\stackrel{~}{\omega }_o}/c`$. This condition yields $`\gamma c\mathrm{\Delta }k/2n`$; for our assumed $`\mathrm{\Delta }k2\pi /(1\mu m)`$ this gives the requirement $`\gamma 300ps^1`$, which is met by typical values of $`\gamma `$ (see equation (8) and Table I).
## VIII Appendix 2
In this appendix we build the continuous collective operator $`b_k`$ (12) that applies for an infinite system of discrete resonators by first considering only excitations that are periodic over a length $`L=N\mathrm{\Lambda }`$, and then passing to the limit $`N\mathrm{}`$. In the periodic case there are still an infinite number of resonators, but only $`N`$ of the $`b_l`$ are independent. Assuming $`N`$ is even, we can take them to be
$$l=\frac{N}{2}+1,\frac{N}{2}+2,\mathrm{},\frac{N}{2}1,\frac{N}{2}.$$
(39)
We denote this range by $`R`$. For an $`l`$ outside $`R_l`$, we have $`b_l=b_{lpN}`$ where $`p`$ is an integer such that $`lpN`$ is within the range (39). If we now introduce discrete wavevectors $`k_m=2\pi m/L`$, where
$$m=\frac{N}{2}+1,\frac{N}{2}+2,\mathrm{},\frac{N}{2}1,\frac{N}{2},$$
(40)
(that is, $`mR`$) we can introduce Fourier amplitudes $`\overline{b}_m`$ according to
$$\overline{b}_m\frac{1}{\sqrt{N}}\underset{lR}{}b_le^{ik_mz_l},$$
(41)
where $`z_l=l\mathrm{\Lambda }`$. We then find immediately that
$$b_l=\frac{1}{\sqrt{N}}\underset{mR}{}\overline{b}_me^{ik_mz_l},$$
and that
$$\underset{lR}{}b_l^{}b_l=\underset{mR}{}\overline{b}_m^{}\overline{b}_m,$$
(42)
while
$$[\overline{b}_m,\overline{b}_m^{}^{}]=\delta _{mm^{}},$$
for example, so
$$\underset{m^{}}{}[\overline{b}_m,\overline{b}_m^{}^{}]=1$$
or
$$\frac{2\pi }{L}\underset{m^{}}{}[\sqrt{\frac{L}{2\pi }}\overline{b}_m,\sqrt{\frac{L}{2\pi }}\overline{b}_m^{}^{}]=1,$$
(43)
a form that we will presently find useful.
We now consider letting $`N\mathrm{}`$, with $`L\mathrm{}`$ such that $`\mathrm{\Lambda }`$ is fixed. Then the range $`R`$ approaches all the integers from $`\mathrm{}`$to $`+\mathrm{}`$, while $`k_m`$ become more closely spaced and approach a dense distribution of points ranging from $`\pi /\mathrm{\Lambda }`$ to $`\pi /\mathrm{\Lambda }`$; this is the first Brillouin zone, and we denote it by $`B.Z.`$ In the usual way, then, we take
$$\frac{2\pi }{L}\underset{m^{}}{}_{B.Z.}๐k^{},$$
(44)
and, if we introduce $`b_k`$ such that
$$\sqrt{\frac{L}{2\pi }}\overline{b}_mb_k,$$
(45)
where the $`k`$ in $`b_k`$ is first identified with $`k_m`$ but then allowed to vary continuously as $`N\mathrm{}`$, from (43) we have
$$_{B.Z.}๐k^{}[b_k,b_k^{}^{}]=1,$$
and so we can identify
$$[b_k,b_k^{}^{}]=\delta (kk^{}),$$
for $`k`$ and $`k^{}`$ within $`B.Z.`$ In this limit, using (44,45), we find
$$\underset{l}{}b_l^{}b_l_{B.Z.}๐kb_k^{}b_k$$
from (42), where the integer $`l`$ now ranges from $`\mathrm{}`$ to $`\mathrm{}`$, and we recover (12) from (41).
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# Hubble Space Telescope ultraviolet spectroscopy of blazars: emission lines properties and black hole masses
## 1 Introduction
Highly Polarized Quasars (HPQ, also referred to as Flat Spectrum Radio Quasars) and, occasionally, BL Lacertae objects, collectively known as blazars, exhibit broad emission lines superimposed on their optical and ultraviolet (UV) continua (Netzer et al. 1994; Scarpa, Falomo, & Pian 1995; Vermeulen et al. 1995; Corbett et al. 1996; Scarpa & Falomo 1997; Koratkar et al. 1998; Pian et al. 2002; DโElia, Padovani, & Landt 2003). Emission lines play an important role in the energetics of blazars: some models of multiwavelength blazar emission (Dermer & Schlickeiser 1993; Sikora et al. 1994; Ghisellini & Madau 1996) predict that the broad line region (BLR) photons are Compton upscattered to X- and gamma-ray energies by the relativistic particles composing the jet plasma, and form luminous high energy spectral components, which often dominate the overall blazar output (Mattox et al. 1997; Bloom et al. 1997; Wehrle et al. 1998; Hartman et al. 2001; Ballo et al. 2002; Pian et al. 2002). The role of broad line emission in shaping the spectrum of different classes of blazars is however not fully assessed (Fossati et al. 1998; Ghisellini et al. 1998; Padovani et al. 2003).
More in general, the characteristics of the BLR of blazars may help in investigating the interplay between the accretion disk and the relativistic jet, which is more prominent in blazars than in normal QSO and Seyferts (Celotti et al. 1997; Maraschi & Tavecchio 2003; DโElia et al. 2003; Wang, Luo, & Ho 2004). AGN broad emission lines may be used also to estimate the masses of the compact objects residing in the nuclear centers, most likely supermassive black holes (BH), by exploiting their dynamical effect on the line-emitting gas clouds. The application of the virial theorem requires that the size of the BLR is determined either with direct methods (reverberation mapping technique, Peterson & Wandel 2000; Kaspi et al. 2000; Peterson et al. 2004), or with indirect arguments (Kaspi et al. 2000; Vestergaard 2002). In the latter case the uncertainties are obviously larger, and the methods must be tested carefully. Since this has important consequences on the evolution and demographics of AGNs, it is crucial to accomplish these measurements both for low and high redshift sources. This approach has been adopted in the estimate of BH masses of large samples of quasars based on optical emission lines (Woo & Urry 2002; McLure & Dunlop 2004). The advantage of using UV, rather than optical emission lines is that the former correspond to a higher ionization state and are therefore presumably more representative of the dynamics close to the central massive object.
In this paper we present the analysis of the broad and intense UV emission lines of 16 blazars observed by the Hubble Space Telescope (HST) and the Faint Object Spectrograph (FOS). Previous studies of AGN UV spectra have been carried out by Bechtold et al. (2002), who focussed on the absorption systems, by Kuraszkiewicz et al. (2002) and Kuraszkiewicz et al. (2004), who present a complete HST FOS atlas of emission line parameters of AGNs, and by Evans and Koratkar (2004), who recalibrated the pre-COSTAR AGN spectra. We concentrate here on the UV spectra of blazars. Based on the radiative and kinematic properties of the broad emission line region (BLR), we measure the luminosities of their BLRs and derive estimates of the BLR sizes and of the central BH masses.
We adopt the โconcordant cosmologyโ, $`\mathrm{\Omega }_m`$ = 0.3, $`\mathrm{\Omega }_\mathrm{\Lambda }`$ = 0.7, and assume $`H_0`$ = 72 km s<sup>-1</sup> Mpc<sup>-1</sup> (Spergel et al. 2003). Luminosities reported by other authors and used in this paper have been transformed into this cosmology.
## 2 Sample selection and data analysis
We have retrieved from the HST archive<sup>1</sup><sup>1</sup>1using MAST, the Multi-mission Archive at STScI, see http://archive.stsci.edu. all pre- and post-COSTAR FOS grating spectra of sources previously classified as blazars (Wall & Peacock 1985; Impey & Tapia 1988; Impey & Tapia 1990; Impey et al. 1991; Stickel et al. 1991; Stocke et al. 1991; Padovani & Urry 1992; Wills et al. 1992; Perlman et al. 1996). We also included in our final list PKS 1229-021, which, despite having low polarization (Wills et al. 1992), is considered a blazar because of significant emission at MeV-GeV frequencies (Hartman et al. 1999), and 3C 273, which has intermediate properties between those of blazars (strong radio emission, superluminal motion, gamma-ray emission, jet emission dominance at hard X-ray energies, e.g., Haardt et al. 1998; Grandi & Palumbo 2004), and those of Seyfert galaxies (broad emission lines, big blue bump). We selected spectra taken with high-resolution gratings in the UV region (G130H, G190H, G270H, G400H).
This search yielded 24 objects with measurable spectra. Spectra taken with the same grating within one day were averaged to increase the signal-to-noise ratio. For this work we have considered only the 16 sources with significant (larger than 3$`\sigma `$) emission line detections. These are reported in Table 1.
For 6 sources there are also low-resolution grating (G160L) observations in the archive, obtained nearly simultaneously to the high-resolution spectra (i.e., within one day). Since the line parameters and spectral indices derived from the former are not significantly different from those measured in the high-resolution spectra, we have neglected these spectra.
Although the considered objects do not represent a complete sample, they form a sizeable dataset to investigate the UV line properties of blazars.
After applying a correction for the Galactic absorption using the maps of Schlegel, Finkbeiner and Davis (1998) and the extinction curve of Cardelli, Clayton and Mathis (1989), we measured the equivalent widths (EWs), the intensities and the full width at half maximum (FWHM) values of the emission lines fitting a linear local continuum on each side of the line (see Table 2). The EW uncertainties were estimated by assuming 2$`\sigma `$ variations of the local continuum.
For each object we have combined the spectra taken quasi-simultaneously (within 1 day) with different gratings, excluded the regions affected by emission or absorption features, binned the signal in 20-50 ร
wavelength intervals and fitted the continuum with a power-law. To account for calibration uncertainties of the data, we added a 5% systematic error to the statistical errors. The derived power-law spectral indices and flux normalizations are given in Table 1.
A similar analysis of the continuum and line properties of these objects has been presented in Kuraszkiewicz et al. (2002) and Kuraszkiewicz et al. (2004).
## 3 Results
The redshifts of our objects (see Table 1) range between $`z=0.158`$ and $`z=1.404`$, with an average value of $`<z>=0.84\pm 0.31`$. Thus, the lines typically detected in our FOS spectra are Ly$`\beta `$, Ly$`\alpha `$, C IV $`\lambda `$1549, C III\] $`\lambda `$1909, Si IV $`\lambda `$1400 and in some cases Mg II $`\lambda `$2798. We report in Table 1 the spectral indices and normalizations of the UV continua and in Table 2 the EW and the FWHM values of the emission lines.
In 3 cases (3C 273, 3C 345, and 3C 454.3) observations at more than one epoch are available. Variations of line and continuum emission are observed, with maximum amplitudes of factors of $``$2 and 7, respectively. However the observations are too limited and sparse to allow a meaningful assessment of correlated line and continuum variability.
In section 3.1 we describe the average properties of the emission lines of blazars in the UV spectral region; in section 3.2 we use the line and continuum properties of the UV spectra to estimate the masses of the central BHs.
### 3.1 Emission line properties
In order to produce a representative high signal-to-noise ratio UV spectrum of blazars, we have combined all the UV spectra in our dataset. Each spectrum was first reduced to rest-frame and then normalized to its average continuum flux. The resulting composite blazar spectrum, normalized to unity at the reference wavelength of 1500 ร
, is shown in Fig. 1. The EW, relative intensities, and FWHM of the emission lines of the composite spectrum are given in Table 3.
The composite spectrum of blazars is similar to that of normal QSO. The line ratios of blazars and normal AGNs are also not significantly different. This is illustrated in Fig. 2, where we report the luminosities of the Ly$`\alpha `$ and C IV $`\lambda `$1549 lines of blazars, compared with those of a list of radio-loud quasars (RLQ) observed by HST FOS (Wills et al. 1995). Note that in this list of RLQ there are 8 objects in common with our sample of blazars; therefore, for the purpose of the comparison, these 8 sources have been considered as blazars and have been excluded from the RLQ list. We also compare with the RLQ 3C 390 (1845+79, $`z=0.056`$), which has hybrid properties, i.e., it has substantial polarization (1.3%, Impey et al. 1991), but is lobe-dominated at radio wavelengths (Ghisellini et al. 1993). Figure 2 shows that the intensity ratio of Ly$`\alpha `$ vs the C IV $`\lambda `$1549 in blazars is consistent with that exhibited by normal RLQ.
More in general, we have compared in Table 3 the average intensity ratios of the lines we detect in our composite blazar spectrum with those reported by other authors for larger samples of QSO or RLQ (Francis et al. 1991; Zheng et al. 1997; Telfer et al. 2002). Except for the C III\] $`\lambda `$1909 line, which appears somewhat underluminous in blazars, there is good overall agreement. This comparison suggests that the structure and physical state of the BLR in blazars and normal RLQ are indistinguishable. This is also confirmed by the comparison of the Ly$`\alpha `$ line luminosities. In Fig. 3 we report the continuum luminosity at 1350 ร
as a function of the Ly$`\alpha `$ luminosity for blazars and RLQ. With the exception of 3C 390, which is at relatively low redshift, the blazars and RLQ have a similar range of Ly$`\alpha `$ luminosities. The averages of the logarithmic distributions, in erg s<sup>-1</sup>, are $`<\mathrm{logL}(\mathrm{Ly}\alpha )>=44.55\pm 0.11`$ and $`<\mathrm{logL}(\mathrm{Ly}\alpha )>=44.72\pm 0.12`$ for blazars and RLQ, respectively (the uncertainties represent the errors associated with the averages, i.e. the standard deviations divided by the square root of the number of objects). However, due to relativistic beaming, blazars have more luminous continua than RLQ, i.e., blazar lines have smaller EW (see Fig. 3). While RLQ emission lines, for any continuum luminosity, have EW between 100 ร
and 1000 ร
, part of the blazars exhibit line EW between 10 ร
and 100 ร
, and these have the most strongly boosted continua. From comparison with the RLQ, we have estimated the luminosity enhancement due to beaming.
Relativistic aberration affects the non-thermal synchrotron luminosity and depends on the fourth power (for a jet geometry) of the relativistic Doppler factor $`\delta `$ (= \[$`\mathrm{\Gamma }(1\beta cos\theta )]^1`$). RLQ are thought to be the parent population of blazars: their jets are directed away from the line of sight, so that their luminosities are only weakly affected by relativistic beaming. Therefore, we have estimated the beaming amplification for the blazars by assuming that their continuum luminosities should exhibit a dependence on Ly$`\alpha `$ line luminosities similar to that of RLQ (Fig. 3). We fitted the RLQ line (in erg s<sup>-1</sup>) and continuum (in erg s<sup>-1</sup> ร
<sup>-1</sup>) luminosities to a power-law and obtained the dependence:
$$L(1350\mathrm{\AA })=0.46\times 10^3L(Ly\alpha )^{1.02},$$
(1)
which has a scatter of 0.2dex in L(1350 ร
). We will use this relationship in Section 3.2.1 to correct the continuum luminosities of blazars for the beaming.
### 3.2 Black hole masses
Under the assumption that the dominant mechanism responsible for the width of the broad emission lines is the gravitational potential of the central supermassive BH, and that the line widths reflect the Keplerian velocities of the line-emitting material in a virialized system (Wandel et al. 1999; McLure & Dunlop 2001), the BH mass M<sub>BH</sub> is given by:
$$M_{BH}=G^1v^2R_{\mathrm{BLR}}$$
(2)
where $`v`$ is the velocity of the gas gravitationally bound to the central BH, $`R_{\mathrm{BLR}}`$ is the size of the BLR, and G is the gravitational constant. The velocity $`v`$ can be obtained directly from the FWHM of the broad emission lines ($`v=f\times v_{FWHM}`$), where $`f`$ is a factor that depends on the geometry and kinematics of the BLR (e.g., McLure & Dunlop 2002; Vestergaard 2002).
#### 3.2.1 Size of the BLR
The most reliable method to derive $`R_{\mathrm{BLR}}`$ is through the reverberation mapping technique (e.g. Peterson et al. 2004, and references therein). This uses the time lag of the emission line light curve with respect to the continuum light curve to determine the light crossing size of the BLR in AGNs. However, this method requires intensive monitoring of the UV continuum and of the lines and can be used only for a limited number of objects (e.g., Korista et al. 1995; Onken et al. 2002), including one of our sources, 3C 273 Paltani & Tรผrler 2005). An alternative way to estimate $`R_{\mathrm{BLR}}`$ is to use the empirical relationship found between $`R_{\mathrm{BLR}}`$ and the optical continuum luminosity (Kaspi et al. 2000).
We have derived this relationship in the UV for a sample of 15 PG quasars and 10 Seyfert 1 galaxies having BLR radii determined via reverberation mapping in the optical (Kaspi et al. 2000) and measured UV spectra (Vestergaard 2002). One of the PG quasars is 3C 273, which is also a member of our blazar sample. For this object, we corrected the continuum luminosity for the beaming effect (see Eq. 1). For one of the Seyferts, NGC 4151, the UV continuum of which varies with high amplitude, we re-measured the continuum luminosity at 1350 ร
from the average IUE spectrum. These quantities are reported in Fig. 4. We fitted the BLR radii and the luminosities at 1350 ร
(rest frame) with a power-law and obtained the following relationship:
$$R_{BLR}=(22.4\pm 0.8)\left(\frac{\lambda L_\lambda (1350\mathrm{\AA })}{10^{44}\mathrm{erg}\mathrm{s}^1}\right)^{0.61\pm 0.02}\mathrm{lt}\mathrm{days}$$
(3)
If we exclude 3C 273 from the fit, the result is unchanged (the power-law index is $`0.60\pm 0.02`$ in this case). We note that the slope in Eq. 3 is consistent with that determined in the same wavelength range by Kaspi et al. (2005, index $`0.56\pm 0.05`$), it is slightly flatter than that found in the optical by Kaspi et al. (2000, index $`0.70\pm 0.03`$) and slightly steeper than that found by McLure and Jarvis (2002, index $`0.50\pm 0.02`$) at 3000 ร
, based on a very similar sample of PG quasars and Seyfert galaxies.
Since the relationship between the BLR radius and continuum luminosity is supposed to be valid in the case of a thermal continuum (see also discussion in Paltani and Tรผrler 2005), we must correct the blazar UV continuum luminosities for the effect of relativistic beaming.
For the blazars with continuum luminosity exceeding the power-law dependence between the RLQ line and continuum luminosities (Eq. 1), we adopted the continuum luminosities computed with Eq. 1 at the corresponding line luminosity, and derived the BLR radii through Eq. 3. These are reported in Table 4. The correction of the continuum luminosity is relevant (i.e., larger than $``$3 times the scatter) for 4 objects (see Fig. 3). We note that our estimate of the BLR radius of 3C 273 is consistent with that reported by Paltani and Tรผrler (2005).
An alternative, independent method for evaluating $`R_{\mathrm{BLR}}`$ consists in coupling the luminosity of the BLR with the information carried by the multiwavelength spectrum of the blazar. Since blazars, among all AGNs, are the only ones with a spectrum extending to gamma-rays, this method is specific for the blazar class of AGNs.
Ten of our blazars have multiwavelength energy distributions which have been fitted with synchrotron and inverse Compton radiation components (Ghisellini et al. 1998). The latter component dominates at the X- and gamma-ray energies and originates from the scattering of relativistic electrons off both synchrotron photons (internal to the jet) and external radiation fields. These include broad line photons, the density of which, $`U_{ext}`$, is thus estimated through the multiwavelength spectral fit.
Following the procedure adopted by Celotti, Padovani and Ghisellini (1997), we reconstructed the total luminosity of the BLR for each of our sources by using the intensities of the observed UV emission lines and by assuming for the unobserved lines the line ratios of an average quasar spectrum (Francis et al. 1991). These derived BLR luminosities are reported in Table 4.
From the fitted densities of the external photons $`U_{ext}`$ and from the observed BLR luminosities, the size of the BLR, $`R_{BLR}`$, can be derived according to:
$$R_{\mathrm{BLR}}=\sqrt{\frac{L_{\mathrm{BLR}}}{4\pi cU_{\mathrm{ext}}\delta ^2}}$$
(4)
where $`\delta `$ is the relativistic Doppler factor required by the multiwavelength modeling. The BLR radii computed with Eq. 4 are reported in Table 4. We have identified this second method as โspectral energy distribution (SED) methodโ, in order to distinguish it from the one based on the empirical determination of $`R_{BLR}`$ from the continuum luminosity. No clear correlation is found between the BLR radii determined with the two methods. One probable explanation of the discrepancy is that the radiation density of the BLR is generally smaller than the parameter $`U_{ext}`$ obtained with multiwavelength fits. This parameter includes not only the BLR photons, but also additional contributions, such as photons coming directly from the accretion disk, or produced by the dusty torus, or by larger regions of the jet (Ghisellini priv. comm.). Thus, the SED method may underestimate the BLR sizes (and therefore the BH masses) in some cases. Moreover, the uncertainties associated with the SED fit parameters are large. Therefore, although we had proposed the โSEDโ method for BH mass determination in an individual source (PKS 0537โ441, Pian et al. 2002), it appears that this method cannot be generalized.
#### 3.2.2 Mass estimates
In order to evaluate the BH masses of our objects we have used Eqs. 2 and 3, adopting a standard value of $`f=\sqrt{3}/2`$ for the kinematic factor, corresponding to an isotropic distribution of the BLR clouds (Wandel 1999; Kaspi et al. 2000; Vestergaard 2002), that yields virial BH masses consistent with those derived from the M<sub>BH</sub>-L<sub>bulge</sub> relationship (Labita, Falomo, & Treves, in prep.). After setting $`v_{FWHM}`$ to suitable units we obtain the relation:
$$M_{BH}=3.26\times 10^6\left(\frac{\lambda L_\lambda (1350\mathrm{\AA })}{10^{44}\mathrm{erg}\mathrm{s}^1}\right)^{0.61}\left(\frac{v_{\mathrm{FWHM}}}{10^3\mathrm{km}\mathrm{s}^1}\right)^2M_{}$$
(5)
Vestergaard (2002) also derives a formula for the central BH mass, based on the continuum measurements at the rest frame wavelength of 1350 ร
of a sample of 26 AGN with BLR radii determined by reverberation mapping. However, this relationship was calibrated against the BH masses determined from optical measurements, and not directly from the BLR size, as we do.
By using Eq. 5 with the beaming-corrected L(1350 ร
) luminosities and $`v_{\mathrm{FWHM}}`$ estimated from our spectra, we have computed the central BH masses. For consistency with Vestergaard (2002) we have used the FWHM of C IV $`\lambda `$1549. These masses are reported in Column 5 of Table 4, and in Fig. 5, where they are compared with those computed with our Eq. 5 for a sample of PG quasars, for which Vestergaard (2002) reports UV luminosities and C IV $`\lambda `$1549 line FWHM values. The blazar BH masses are statistically consistent with those of quasars: the averages of the logarithmic distributions, in solar masses, are $`<\mathrm{logM}_{\mathrm{BH}}>=8.31\pm 0.10`$ and $`<\mathrm{logM}_{\mathrm{BH}}>=8.42\pm 0.08`$ for blazars and RLQ, respectively (the uncertainties represent the errors associated with the averages, i.e. the standard deviations divided by the square root of the number of objects).
For comparison, we have also computed the blazar BH masses using Vestergaardโs relationship (Eq. 8 of Vestergaard 2002), with our measured blazar luminosities and FWHM values of the C IV $`\lambda `$1549 line. These are systematically lower than the corresponding ones determined using Eq. 5, by a factor 1.4-2. We have also compared our estimated BH masses with those determined for the same blazars by other authors: a number of our objects are in common with the samples of Liang and Liu (2003), Woo and Urry (2002), and Wang et al. (2004). Their mass estimates have been reported in the last 3 columns of Table 4: our masses are generally smaller. This may be partially due to our correction of the continuum luminosity for the beaming effect. In particular, this must be the case for two of the four objects in our sample with the most strongly beamed continuum, 3C 279 and 0537-441.
We do not find a correlation of the BH mass with redshift for our blazar sample.
## 4 Summary and Conclusions
We have studied the properties of the UV emission lines of blazars, mostly from single epoch HST FOS spectra, and found that the average blazar UV spectrum is similar to that of RLQ. This is the sum of a thermal and non-thermal component. Our targets are mainly HPQ and Low-Frequency Peaked BL Lacs (Padovani & Giommi 1995; Fossati et al. 1998) where the emission of the non-thermal synchrotron component peaks at optical/IR frequencies. Therefore a large relative contribution from the thermal accretion disk is expected in the UV region (e.g., Bregman et al. 1986). This is clearly evident from the spectral energy distribution of 3C 273 (Ulrich et al. 1980; Courvoisier 1998) and may be significant in 3C 279 (Pian et al. 1999).
With the aim of estimating central BH masses of blazars, we have assumed Keplerian conditions in the BLR gas motion and have evaluated the BLR size using the results of a fit of UV luminosities and BLR radii of a sample of QSOs having BLR sizes determined via reverberation mapping in the optical. We have derived a relationship between $`R_{BLR}`$ and luminosity in the UV domain (1350 ร
at rest frame), which exhibits a slope consistent with that of Kaspi et al. (2005), although slightly flatter and steeper than proposed in the optical and near-UV by Kaspi et al. (2000) and McLure & Jarvis (2002), respectively.
For those 10 blazars having multiwavelength spectral fits we have also applied an independent method of BLR size determination, based on the combination of the observed $`L_{BLR}`$ and fitted external radiation density (โSEDโ method). We have not found a clear correlation between the BLR sizes obtained with the two methods, with the largest deviations observed in the sense of a deficit of the โSEDโ radii with respect to those obtained with the empirical $`R_{BLR}\lambda L_\lambda `$ relationship. We conclude that the SED method yields a BLR size inconsistent with that derived from the continuum luminosity.
Our estimated BH masses have an average of $`(2.8\pm 2.0)\times 10^8M_{}`$ (the quoted uncertainty is the standard deviation) and are comparable with those of lower redshift blazars, estimated with different methods (Barth, Ho, Sargent 2003; Falomo, Carangelo, & Treves 2003).
The distribution of our blazar BH masses computed with Eq. 5 is consistent with the distribution of the PG quasar masses computed with the same equation. These results suggest that the differences between radio powerful sources and radio-weak ones are not due to the mass of the BHs residing at their centers. However, the validity of this conclusion at the intermediate/high redshifts must be corroborated by the analysis of wider samples of homogeneous datasets. Moreover, further intensive spectroscopic monitoring of the brightest blazars at optical and UV wavelengths is required, in order to construct well sampled continuum and emission line light curves for the application of the reverberation mapping technique.
## Acknowledgments
We thank R. Bohlin for assistance with HST data analysis, G. Ghisellini for helpful discussion, and the referee, A. Koratkar, for a constructive report. We acknowledge the use of the SIMBAD and NED databases, publicly available online. This work was partially supported by the Italian Space Agency under the contract I/R/056/02.
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# A comparative analysis of empirical calibrators for nebular metallicity
## 1 Introduction
HII regions, from diffuse HII regions in the Galaxy to Extragalactic HII Regions (GEHR) as well as HII Galaxies, have been for many years, the main source of information about metallicity in distant galaxies. Their bright emission line spectra are visible in all kinds of objects where there are recent episodes of star formation. The analysis of these nebular spectra constitutes the best method, if not the only one, for the determination of the chemical abundances of elements such as: He, N, O, Ne, Ar, S having optical emission lines corresponding to different ionization states. An accurate knowledge of these abundances is essential for a complete understanding of the evolution of stars and stellar systems and has allowed some light to be shed on several questions concerning the chemical evolution of galaxies in the Local Universe. They are now becoming even more relevant with the regard to the younger Universe.
Recombination lines yield the most accurate abundances because of their weak dependence on nebular temperature. In fact, helium abundances can be derived to an accuracy better than 5 %. Unfortunately, most of the observed emission lines in ionised nebulae are collisionally excited and their intensities depend exponentially on temperature. In principle, this temperature can be determined from appropriate line ratios, the most widely used being that of \[OIII\]$`\lambda 4363`$ ร
/$`(\lambda 4959`$ ร
$`+\lambda 5007`$ ร
) although, recently, the improved sensitivity of new detectors in larger telescopes allows for the measurement of other auroral lines, which are less temperature sensitive (Kinkel & Rosa 1994; Castellanos, Dรญaz & Terlevich 2002; hereinafter CDT02).
All these ratios involve the detection and measurement of one intrinsically weak line which in many objects is too faint to be observed. This is the case for regions with high metal content โ where the efficient cooling exerted by metallic ions renders weak lines undetectable โ, HII regions in distant galaxies and objects with low surface brightness. In these cases, empirical methods based on the intensities of strong, easily observable, optical lines have been developed and are nowadays widely used.
The so called โempirical methodsโ are based on the cooling properties of ionised nebulae which ultimately translate into a relationship between emission line intensities and oxygen abundance. In fact, when the cooling is dominated by oxygen, the electron temperature depends inversely on oxygen abundance. Since the intensities of collisionally excited lines depend exponentially on electron temperature, a relation is expected to exist between these intensities and oxygen abundances.
According to Pagel et al. (1979), under the assumptions that a) the nebula is ionisation bounded, b) the region can be represented by small clumps of gas with a given electron density surrounded by a much less dense material, so that the degree of ionisation is proportional to ($`ฯต^2n_eQ_H)^{1/3}`$ where $`n_e`$ is the clump electron density, $`Q_H`$ is the number of hydrogen ionising photons and $`ฯต`$ is the filling factor and c) the cooling is fixed by oxygen abundance, we can consider that the emission line spectrum of the nebula will depend on: the energy distribution of the ionising radiation field, the ionisation parameter and the oxygen abundance . Therefore, if a single relation between the chosen calibrator and the oxygen abundance is sought, further assumptions are needed implying that either the hardness of the radiation field or the degree of ionisation or both depend on oxygen abundance.
Following these ideas, several abundance calibrators have been proposed involving different emission line ratios: among others, \[OIII\] $`\lambda `$ 5007 ร
/H$`\beta `$ (Jensen, Strom & Strom 1976; \[OIII\] $`\lambda `$ 5007 ร
/ \[NII\] $`\lambda `$ 6584 ร
(Alloin et al. 1979) and (\[OII\] $`\lambda `$ 3727 ร
+ \[OIII\] $`\lambda `$ 5007 ร
)/ H$`\beta `$ (R<sub>23</sub>; Pagel et al. 1979) . The advantages and drawbacks of the different calibrators have been discussed by several authors (see Pagel, Edmunds & Smith 1980; Kennicutt & Garnett 1996; Kewly & Dopita 2002). Although abundances derived through the use of these calibrations are recognised to suffer from considerable uncertainties, they are still believed by many authors to trace large-scale trends in galaxies. Empirical methods have been used to derive abundances in objects as different as dwarf irregular galaxies (e.g. Skillman, Kennicutt & Hodge 1989), individual HII regions in spiral galaxies (e.g. Oey & Kennicutt 1993), low surface brightness galaxies (McGaugh 1994), nuclear starbursts (e.g. Storchi-Bergmann, Calzetti & Kinney 1994) and even active galactic nuclei (Storchi-Bergmann et al. 1998). They have also been employed in the derivation of abundance distributions in the discs of spiral galaxies (e.g. Belley & Roy 1992; van Zee et al. 1998) and emission line galaxies at intermediate redshift (e.g. Kobulnicky & Kewley 2004)
In this work we perform a comparative analysis of the principal empirical calibrations of abundances which are based on the intensities of the nebular lines of oxygen, nitrogen and sulphur, visible in the optical and far red spectral regions. All these calibrations present a considerable scatter, usually larger than that associated with observational errors and probably related to the assumptions mentioned above. The aim of our work is to understand the reasons for this scatter and, whenever possible, to find ways of improving the empirical derivation of abundances.
In order to do that, we have compiled a large sample of emission line objects (HII galaxies, GEHR and diffuse HII regions in the Galaxy and the Local Group) with a direct determination of the total oxygen abundance through the measurement of the auroral lines of \[OIII\]$`\lambda `$4363 ร
, \[OII\]$`\lambda `$ 7327 ร
or \[SIII\]$`\lambda `$6312 ร
and we have constructed a complete sequence of photo-ionisation models, covering the main physical properties of these objects.
In the next section we describe the sample of objects and the process of determination of the oxygen abundance. In Section 3 , we summarize the main properties of the photoionisation models used for our analysis which is then presented in Section 4. Finally, Section 5 summarizes our results and the main conclusions reached.
## 2 Sample of objects and abundance derivation
Our sample is comprised of a combination of different emission line objects ionised by young massive stars: diffuse HII regions in the Galaxy and the Local Group (DHR), Giant Extragalactic HII regions (GEHR) and HII galaxies (HIIG) and therefore does not include planetary nebulae or objects with non-thermal activity. For all of them direct determinations of electron temperature exist thus allowing the derivation of the oxygen abundance which we have taken as the observational metallicity indicator. The sample includes the objects analysed in (Dรญaz & Pรฉrez-Montero 2000, hereinafter: DPM00) with the addition of low excitation GEHRs from CDT02; GEHRs in M101 (Kennicutt, Bresolin & Garnett 2003) and M51 (Garnett, Kennicutt & Bresolin 2004); GEHRs in galaxies in the Sculptor Group (Skillman, Cรดtรฉ & Miller 2003) and the Virgo cluster (Vรญlchez & Iglesias-Pรกramo 2004); and HII galaxies from the works of Guseva et al. (2000), Popescu & Hopp (2000)and Kniazev et al. (2001). Data from these latter objects have been complemented with information in the spectral range between 7000 ร
and 1 micron from Pรฉrez-Montero & Dรญaz (2003; hereinafter PMD03, 12 HII galaxies) and Garnett (1992; 13 objects), including the \[SIII\] strong emission lines.
The sources for the line intensities, together with the number and class of the collected objects, are summarized in Table 1. The total sample comprises 367 objects with lines in the optical part of the spectrum, 282 of them with \[NII\] data, and 126 with near IR \[SIII\] data.
The physical conditions โ electron temperature, electron density and oxygen abundance โ for the whole sample have been recalculated using the same procedures as in PMD03, based on the fivel-level statistical equilibrium model in the task TEMDEN contained in the software package IRAF (De Robertis, Dufour & Hunt 1987; Shaw & Dufour 1995). The atomic coefficients used are the same as in PMD03 and are referenced in Table 4 of that work. Electron densities were determined from the \[SII\] $`\lambda `$ 6717ร
/ $`\lambda `$ 6731ร
line ratio. Electron temperatures have been calculated from the \[OIII\] ($`\lambda `$ 4959ร
+$`\lambda `$5007ร
)/ $`\lambda `$ 4363ร
line ratio for all but 13 objects of the sample for which the \[SIII\] ($`\lambda `$ 9069ร
+$`\lambda `$9532ร
)/$`\lambda `$ 6312ร
line ratio has been used isnstead. These latter objects are of low excitation and lie on the high metallicity range (for example, CDT1 in NGC1232 (Castellanos et al. 2002) or S5 in M101 (Kinkel & Rosa, 1994)). For them, an empirical relation between \[OIII\] and \[SIII\] electron temperatures has been used:
$$t([OIII])=0.95t([SIII])+0.08$$
based on the grids of photo-ionisation models described in the next section and, differing slightly from the empirical relation found by Garnett (1992), due mainly to the introduction of the new atomic coefficients for S<sup>2+</sup> from Tayal & Gupta (1999).
Regarding \[OII\] temperatures, for 81 objects of the sample it has been possible to derive its value from the \[OII\]($`\lambda `$ 3726ร
+$`\lambda `$3729ร
) /$`\lambda `$ 7325ร
line ratio. <sup>3</sup><sup>3</sup>3The \[OII\] $`\lambda `$7319ร
+$`\lambda `$7330ร
lines can have a contribution by direct recombination which increases with temperature. Using the calculated \[OIII\] electron temperatures, we have estimated these contributions to be less than 4 % in all cases and therefore we have not corrected for this effect. For the rest of the objects in the sample we have resorted to the model relations between t(\[OII\]) and t(\[OIII\]) found in PMD03 that take explicitly into account the dependence of t(\[OII\]) on electron density. This can affect the deduced abundances of $`O^+/H^+`$ by non-negligible factors, larger in all cases than the reported errors.
Figure 1, shows a comparison of the total oxygen abundances derived as described above with the values published in the original sources. The deviations from the 1:1 relation arise mostly in the high metallicity range as a result of the dependence of t(\[OII\]) on density which affects the calculated O<sup>+</sup>/H<sup>+</sup> abundances. This can be better seen in Figure 2 where the abundance differences are plotted as a function of the O<sup>+</sup>/O<sup>++</sup> ionic fraction.
The oxygen abundances of the sample objects cover the range 0.02Z (IZw18; Skillman & Kennicutt 1993) to 1.82Z (CDT1 in NGC1232; Castellanos et al. 2002)<sup>4</sup><sup>4</sup>4A solar value of 12+log(O/H)= 8.69 (Allende-Prieto, Lambert & Asplund 2001) is assumed through this paper. . The re-calculated \[SII\] electron densities, \[OII\] and \[OIII\] electron temperatures and oxygen abundances are listed in Table 2. The quoted uncertainties have been derived from the emission line flux errors as published in the corresponding references. In the upper panel of Figure 1, the values of these errors (half-error bars) are plotted as a function of oxygen abundance. It can be seen that the errors are almost constant, with an average value of about $`\pm `$0.07 dex, up to 12+log(O/H)$``$ 8.1 and thereafter increase with metallicity up to $`\pm `$ 0.5 dex at the highest derived abundances.
The complete table will be available in electronic form at CDS (Centre de Donnรฉes astronomiques de Strasbourg), via anonymous ftp to cdsarc.u-strasbg.fr (130.79.128.5), (http://cdsweb.u-strasbg.fr) , or at http://pollux.ft.uam.es/enrique/Table1/. Only an example is given here.
At any rate, it should be recalled that the determination of the gaseous chemical abundances is usually accomplished by combining results from photoionization models and observed emission line intensity ratios. Even when the electron temperature can be determined with good accuracy, there are several major unsolved problems that severely limit the confidence of present results including: (1) the effect of temperature structure in multiple-zone models (PMD03); (2) the presence of temperature fluctuations across a given nebula (Peimbert, 2003); (3) collisional and density effects on ion temperatures (Luridiana, Peimbert, & Leitherer 1999); (4) the presence of neutral zones affecting the determination of ionization correction factors (ICFโs)(Peimbert, Peimbert & Luridiana 2002); (5) the ionization structure which is not adequately reproduced by current models (PMD03); (6) the possible photon leakage that affects the low ionization lines formed in the outer parts of the ionized regions (Castellanos, Dรญaz, & Tenorio-Tagle 2002). The first three effects can introduce uncertainties with respect to the derived oxygen abundances of about 0.2, 0.3 and 0.4 dex respectively, depending on excitation. The uncertainties introduced by the latter effects have not yet been quantified.
## 3 Photo-ionisation models
In order to identify and understand the possible sources of scatter in the different empirical calibrations, we have calculated a set of photo-ionisation models covering the physical conditions of the observed objects. No attempt has been made however to perform recalibrations using these models.
The photoionisation models have been computed with the most recent version of the photo-ionisation code CLOUDY 96 (Ferland 2002).
Each modelled HII region has been assumed to be spherically symmetric, with the ionised emitting gas taken to be of constant density (10 and 100 particles per cm<sup>3</sup>), located at a distance very large compared to its thickness and therefore allowing the approximation of plane-parallel geometry. The gas is ionised by a single massive star whose spectral energy distribution (SED) is represented by a CoStar NLTE stellar atmosphere (Schaerer & de Koter, 1997) with effective temperature between 35000 K and 50000 K . The impact on the analysis of the general trends shown by empirical calibrations using other SEDs is negligible and, in practice, only slightly affects the absolute effective temperature scale. Ionisation parameters ($`U`$) between 10<sup>-2</sup> and 10<sup>-3</sup>, which is the range corresponding to the ionisation degree shown by the sample objects, have been chosen. Finally, the solar chemical abundances used are those given in Table 3.
We have computed models with values of this solar abundance multiplied by factors: 1.7, 0.85, 0.34, 0.17, 0.08 and 0.04 corresponding to 1, 0.5, 0.2, 0.1, 0.05 and 0.025 times the solar Grevesse & Sauval (1998) value ($`\mathrm{log}(O/H)`$=-3.08). The refractory elements: Fe, Mg, Al, Ca, Na and Ni have been depleted by a factor of 10 and Si by a factor of 2 (Garnett et al. 1995), to take into account the presence of dust grains. In the case of nitrogen, we have considered, for the models with $`U`$ = 10<sup>-2.5</sup>, another set of abundances with a value of (N/O) 0.5 dex lower than the solar value, close to the values found in low metallicity nebulae.
For the sake of clarity, only the models for $`n_e`$=100 cm<sup>-3</sup>, T<sub>eff</sub> = 35000 K and 50000 K, and $`U`$ = 10<sup>-2.0</sup> and 10<sup>-3.0</sup> are shown in the figures.
## 4 Empirical abundance parameters
### 4.1 The R<sub>23</sub> (O<sub>23</sub>) parameter
The R<sub>23</sub> parameter, that here we have preferred to rename as O<sub>23</sub> in order to differentiate it from the analogous parameter based on sulphur emission lines, was defined as:
$$O_{23}\frac{\mathrm{I}(3727\AA )+I(4959\AA )+I(5007\AA )}{\mathrm{I}(H\beta )}$$
by Pagel et al. (1979). Its relation with oxygen abundance for the objects of the compiled sample can be seen in the left panel of Figure 3. The relation is double-valued. This is due to the efficiency of oxygen as a cooling agent thus decreasing the strength of the oxygen emission lines at high metallicities. At low metallicities however, the cooling is mainly exerted by hydrogen and the oxygen line strengths increase with metallicity. The value of logO<sub>23</sub> reaches a maximum of about 1.2 at an oxygen abundance of 12+log(O/H) $``$ 8.0.
Three different regions can be distinguished in the plot: a lower branch in which O<sub>23</sub> increases with incresing abundance, an upper branch in which the opposite occurs and a turnover region. The two branches can, in principle, be fitted by regression lines with positive and negative slope respectively providing a low to moderate uncertainty in the determination of the metallicity. In the turnover region with log O<sub>23</sub> $``$ 0.8 and 12+log(O/H) $``$ 8.0, although the precise values are difficult to assess, objects showing the same value of O<sub>23</sub> can have oxygen abundances that differ by almost an order of magnitude. It should be noted that a large proportion of the data lie on top of this ill defined zone (up to 40% of the total number of objects and even more in the case of HII galaxies) where the abundance determination can be very uncertain.
Another characteristic of the calibration which is readily apparent is the existence of a scatter larger than accounted for by observational errors. This scatter is related to the fact that, in general, ionised regions do not constitute a single parameter family, hence different geometries of the emitting gas (ionisation parameter) and different ionising radiation temperatures can affect the values of the abundance parameter O<sub>23</sub>. This can be seen in the right panel of Figure 3, where data are shown together with different model sequences . Error bars have been omitted for the sake of clarity. O<sub>23</sub> is seen to depend on both ionisation parameter and stellar effective temperature. The dependence on ionisation parameter is more evident at low metallicities, while the dependence on effective temperature is important in all metallicity regimes. This double dependence makes of O<sub>23</sub> a rather unsuitable abundance parameter. In fact, at an oxygen abundance of 12+log(O/H) = 7.8, logO<sub>23</sub> can vary between 0.5 and 0.9.
Different assumptions about the effects of metallicity on either nebular structure or ionising temperature have been used by different authors in order to define a sequence of models that would eventually allow the calibration of the upper branch where observational data is very scarce. From analyses of HII region data, McCall, Rybski & Shields (1985) concluded that the stellar ionising temperature varied with metallicity while the filling factor remained constant, whereas Dopita & Evans (1986) concluded the opposite: that ionising temperature was constant while U varied with oxygen abundance. These two different assumptions led to calibrations yielding abundances that differ by more than a factor of two.
Theoretical stellar evolution models point to a relation between stellar metallicity and effective temperature in the sense that, for a given mass, stars of higher metallicities show lower effective temperatures. This fact led McGaugh (1991;McG91) to produce a new calibration based on more realistic theoretical models in which the ionisation is provided by stellar clusters of different metallicities. According to his models, in the upper branch, O<sub>23</sub> is relatively insensitive to both ionising temperature and U and the models converge to a single sequence. In the lower branch however, O<sub>23</sub> is mostly dependent on U, as has already been shown by Skillman (1989; S89) and additional information about this parameter is needed in order to apply the empirical method with greater confidence.
Some authors have used the \[OII\]/\[OIII\] ratio as an ionisation parameter indicator to obtain this additional information (e.g. Kobulnicky et al. 1999). The hardness of the ionising radiation however also affects this ratio in a significant way. In fact, at a given value of U, the \[OII\]/\[OIII\] ratio is lower for higher stellar effective temperatures as a result of the increase of the ionisation of O<sup>+</sup> to O<sup>++</sup>. These effects can be seen in Figure 4 where $`\mathrm{log}O_{23}`$ is plotted as a function of log (\[OII\]/\[OIII\]) (left panel) where it can be seen that the \[OII\]/\[OIII\] ratios corresponding to models with the same ionisation parameter and different stellar ionising temperature widely differ.
Just the opposite happens in the case of \[SII\]/\[SIII\], another line ratio used as an ionisation parameter indicator (e.g Dรญaz et al. 1991). At constant U, \[SII\]/\[SIII\] increases somewhat with increasing stellar effective temperature as more S<sup>++</sup> is converted to S<sup>3+</sup>, although in this latter case the effect is important only for the highest ionisation parameters (U$``$ 10<sup>-2</sup>), thus making the \[SII\]/\[SIII\] ratio a more useful ionisation parameter diagnostic. Figure 4 (right panel) shows that \[SII\]/\[SIII\] depends mostly on the ionisation parameter and is rather insensitive to the stellar effective temperature.
Regarding observational data, no clear relation is found between $`\mathrm{log}O_{23}`$ and log (\[SII\]/\[SIII\]), implying that oxygen abundance and the ionisation parameter are not correlated. On the other hand, a definite trend between $`\mathrm{log}O_{23}`$ and \[OII\]/\[OIII\] is clearly seen which can be explained by the expected dependence of stellar effective temperature and metallicity (see also Kewley & Dopita 2002).
Finally, in the turnover region, O<sub>23</sub> is sensitive to both ionisation parameter and ionising temperature and is almost insensitive to oxygen abundance.
It should be taken into account that McGaugh models use zero age star clusters. The situation becomes much more complicated when the evolution of these clusters is taken into account; the evolution of massive stars is fast and metallicity dependent and the cluster ionising temperature may not be a monotonicaly decreasing function of age once WR stars begin to appear (Garcรญa-Vargas, Bressan & Dรญaz 1995; Stasiลska & Leitherer 1996).
Figures 5 and 6, show the residuals of several published calibrations of the O<sub>23</sub> parameter in both lower and upper branches as well as in the intermediate region. The calibrations of S89 for the lower branch and Zaritsky, Kennicutt & Huchra (1994; hereinafter ZKH94) for the upper one, involve only the O<sub>23</sub> parameter. In contrast, the calibrations of McG91 and Pilyugin (2000, 2001a; hereinafter Pil00 and Pil01a, respectively) take also into account the dependence on ionisation parameter and effective temperature via the (\[OII\]/\[OIII\]) ratio (Kobulnicky et al. 1999) and the parameter, P, defined as the quotient of \[OIII\] and (\[OII\]+\[OIII\]) by Pil00.
For the lower branch, that we have considered as corresponding to 12+log(O/H) $`<`$ 8.0, the best fit is found for the calibration by Pil00, with a mean value for the metallicity 0.03 dex higher than the mean value of those directly derived. Its uncertainty in this regime, understood as the standard deviation of the residuals, is $`\pm `$0.15 dex which should be compared to the $`\pm `$0.07 dex average error in the direct oxygen abundance determination (see Figure 1). The calibrations by Skillman and McGaugh present also the same dispersion, although the mean values of the deduced abundances are 0.06 dex lower and 0.21 dex higher respectively than the mean value of those directly derived. A slight trend in the sense of abundances being more underestimated as the metallicity increases is found in the S89 calibration. This trend is probably introduced by the slight dependence of O<sub>23</sub> on ionisation parameter at high effective temperature that as shown by HII galaxy data in left panel of Figure 4.
In fact, in this metallicity range, he $`O_{23}`$ values predicted by photoionisation models for different values of the ionisation parameter and stellar effective temperature produce a scatter in the $`O_{23}`$ versus 12+log(O/H) relationship larger than shown by observational data. This probably indicates that the objects compiled to perform the calibrations, mainly HII galaxies, show very similar properties, ฤฑ.e. they show a very restricted range of ionisation parameters and ionising temperatures. It should be noted however, that not all HII galaxies share those properties and that, in particular, those that lack detectable electron temperature sensitive lines and therefore are the probable targets of the empirical calibrations, show ionisation parameters which are significantely lower. These include Luminous Compact Blue Galaxies (LCBGs) (Hoyos & Dรญaz 2005).
For the upper branch, in the range 12+log(O/H) $``$ 8.4 the best fit is found with McG91 calibration, although it predicts abundances 0.08 dex larger than the mean value, with a dispersion of 0.19 dex. Taking into account that the average error in the direct oxygen abundance determination in this regime is about $`\pm `$0.20, this calibration provides a good estimate of the oxygen abundance for metallic objects. The calibration by Pil01a underestimates abundances by 0.14 dex on average with a dispersion of 0.21 dex. Contrary to what is found for the lower branch, the only calibration not taking into account the dependences on log $`U`$ and T<sub>eff</sub> (ZKH94) shows the largest dispersion (0.27 dex) which probably implies that, in this case, the calibrating sample objects do not share ionising properties.
Finally, for the intermediate region, in the range 8.0 $``$ 12+log(O/H) $`<`$ 8.4, we have evaluated both the calibrations for the lower and upper branches. In this case it can be observed that the lower branch calibrations underestimate the metallicity and the upper branch calibrations overestimate it with residuals that increase with increasing and decreasing metallicity respectively. In this regime it is virtually impossible to choose any reliable calibration of O<sub>23</sub>.
There are not many ways to improve on the O<sub>23</sub> abundance parameter calibration, since HII regions and HII galaxies are ionised by young star clusters and, as these clusters evolve, their ionisation parameters and ionising temperatures change in ways that are not easy to parametrize. A more promissing approach is the search for other potentially useful abundance parameters, some of which are examined below.
### 4.2 Parameters involving \[NII\]
The N2 parameter was defined as:
$$N2\mathrm{log}\frac{I(6584\mathrm{\AA })}{I(H\alpha )}$$
by Denicolรณ, Terlevich & Terlevich (2002, hereinafter DTT02), although it was used before as an empirical estimator by Storchi-Bergmann et al. (1994) and Van Zee et al. (1998). The relation between N2 and the logarithmic oxygen abundance is shown in the left panel of Figure 7 for all the objects in the sample for which nitrogen data exists. The N2 parameter has several advantages: first of all, contrary to O<sub>23</sub>, the relationship between N2 and oxygen abundance is single-valued and secondly, since the emission lines on which it is based are very close in wavelength, the N2 parameter is almost free of uncertainties introduced by reddening corrections or flux calibrations. The dash-dotted line in the plot corresponds to the relation found by DTT02, for their sample objects:
$$12+\mathrm{log}(O/H)=9.12+0.73\mathrm{N2}$$
This line represents a reasonable fit to the data, but shows a large scatter at all metallicities. Most of the scatter is shown by GEHR data, while HII galaxies define a much narrower relation. Most GEHR data fall below the line. This is not surprising since many of the metal rich GEHR used by DTT02 in their calibration had oxygen abundances derived from O<sub>23</sub> and directly derived abundances for metal rich HII regions tend to be lower than those derived from this parameter (Castellanos et al. 2002; Garnett et al. 2004).
Again, the main reason for the dispersion is probably related to the different ionisation parameters and stellar effective temperatures of GEHR since the N2 parameter depends on both as can be seen in the right panel of Figure 7. The \[NII\] lines become weaker as the excitation degree and/or the ionising temperature increase. N2 reaches a maximum of about -0.5 for models with low effective temperature (35000 K) while the lowest values of the parameter are found in models with high effective temperature (50000 K) and high ionisation parameter (logU=-2.0). An additional source of scatter is related to the possibly different N/O relative abundances. To try to quantify this effect we have added a set of photo-ionisation models with a value of log(N/O) 0.5 dex lower than our solar assumed value (see Table 3) and log $`U=10^{2.5}`$. This model sequence can be seen also in Figure 7 (right panel). A lower N/O ratio mimics a higher ionisation parameter. Model sequences of high ionising temperatures and constant N/O ratio seem to reproduce adequately HII galaxy data while a sequence of models with N/O incresing with oxygen abundance would seem more adequate for GEHR data.
The residuals of the fit from the DTT02 calibration against the directly determined oxygen abundances are represented in Figure 8. This empirical calibrator works reasonably well in the turnover region of the logO<sub>23</sub> versus 12+log(O/H) plot, although the dispersion reaches a value of 0.27 dex for the reasons given above. This dispersion is of 0.25 dex for HII galaxies. For the rest of the sample it reaches 0.3 dex with the mean value 0.1 dex higher thus implying that, in these cases, oxygen abundances could be overestimated.
Other empirical parameters involving the \[NII\] lines are the \[NII\]$`\lambda \lambda `$6548,6584 ร
/\[OIII\]$`\lambda \lambda `$4959,5007 ร
ratio, first proposed by Alloin et al. (1979) and recently revindicated by Pettini & Pagel (2004), and the \[NII\]$`\lambda \lambda `$6548,6584 ร
/\[OII\]$`\lambda \lambda `$3727,3729 ร
and \[NII\]$`\lambda \lambda `$ 6548,6584 ร
/\[SII\]$`\lambda \lambda `$ 6716,6731 ร
ratios suggested by Dopita & Evans (1986) and Kewley & Dopita (2002) as metallicity calibrators in the high abundance regime. In Figure 9, 12+log(O/H) is represented as a function of the \[NII\]/\[OIII\] parameter (top panel), the \[NII\]/\[OII\] parameter (middle panel) and the \[NII\]/\[SII\] parameter (bottom panel) for the compiled sample of objects. ยฟFrom the direct comparison with observational data, it can be seen that these three parameters are valid only for a metallicity higher than 12+log(O/H) $``$ 7.8 and with a scatter similar to that found for the N2 parameter. This scatter could again be related to the different objects presenting different N/O ratios. In fact, a clear segregation is found between HII galaxy and GEHR data, more evident in the upper panel, which is probably related to the GEHR showing higher values of N/O and hence N/S ratios. Figure 10 (upper panel) shows that a tight relation exists between log(\[NII\]/\[OII\]) and log(N<sup>+</sup>/O<sup>+</sup>) that is, in turn, a very good indicator of log(N/O). The fit of a regresion line to the data produces the relation:
$$\mathrm{log}\left(\frac{N}{O}\right)=1.144\mathrm{log}\left(\frac{[NII]}{[OII]}\right)0.232$$
The uncertainty involved in the determination of the N/O ratio from \[NII\]/\[OII\] for the whole sample, represented by the standard deviation of the residuals, is 0.14 dex, but decreases to only 0.08 for HII galaxies (see Figure 10, lower panel).
### 4.3 The S<sub>23(4)</sub> parameter
The S<sub>23</sub> parameter was defined by Vรญlchez & Esteban (1996) as:
$$S_{23}\frac{\mathrm{I}(6717\AA )+I(6731\AA )+I(9069\AA )+I(9532\AA )}{\mathrm{I}(H\beta )}$$
using the \[SII\] and \[SIII\] lines analogous to those of \[OII\] and \[OIII\] in the O<sub>23</sub> parameter. It was proposed by Christensen, Petersen & Gammerlgaard (1997) as a sulphur abundance indicator and more recently by DPM00 as an oxygen abundance indicator due to the characteristics evident in Figure 11 (left panel): firstly its single-valued behaviour up to solar metallicities, and secondly its lower dependence on the other functional parameters. The fact that S<sub>23</sub> presents a lesser dependence on effective temperature and ionisation parameter than $`O_{23}`$, seems to be confirmed by photo-ionisation models (see Figure 11, right panel). There is a dependence of S<sub>23</sub> on log$`U`$ and T<sub>eff</sub> but it is weaker than for O<sub>23</sub> (see Figure 3 for a comparison).
ยฟFrom the observational point of view, the S<sub>23</sub> parameter has two important advantages: firstly the sulphur lines remain intense even for the highest metallicity objects and secondly it is relatively independent of reddening, since the lines of both \[SII\] and \[SIII\] can be measured relative to nearby hydrogen recombination lines. On the negative side, \[SIII\] lines shift out of the far red spectral region for redshifts higher than 0.1.
Using the newly added observational data we have improved the DPM00 relation to:
$$12+log(O/H)=8.15+1.85\mathrm{log}S_{23}+0.58(\mathrm{log}S_{23})^2$$
whose residuals for the complete sample relative to the directly determined oxygen abundances are represented as a function of oxygen abundance in the upper panel of Figure 16 . The dispersion is approximately equal to 0.2 dex in all the abundance ranges although it decreases to 0.10 dex for the HII galaxy sample. The relation is not linear. Values of S<sub>23</sub> lower than expected are found for higher excitation nebulae having low metallicity probably due to the presence of \[SIV\] in non-negligible amounts, as seems to be indicated by the position on the diagram of IZw18, the least metallic object. Unfortunately, despite recent observations of a sample of HII galaxies in the near IR spectral range (PMD03), there is a considerable lack of data on objects of low metallicity whose inclusion would definitely improve the calibration.
Oey & Shields (2000) have defined a new parameter, S<sub>234</sub>,
$$S_{234}\frac{\mathrm{I}(6725\AA )+I(9069\AA )+I(9532\AA )+I(10.5\mu )}{\mathrm{I}(H\beta )}$$
which takes into account the contribution of \[SIV\] through its emission line at 10.52 $`\mu `$. The contribution of \[SIV\] is expected to be relevant only in objects with a high degree of ionisation (Dรญaz et al. 1991) and therefore the use of S<sub>234</sub> almost eliminates the dependence on ionisation parameter found for S<sub>23</sub> (Kennicutt et al. 2000). In fact, photoinisation models indicate that S<sub>23</sub> is only slightly dependent on ionisation parameter but shows a non negligible dependence on effective temperature, which becomes more evident at high metallicities (see Figure 13).
Unfortunately, the sample of objects for which the \[SIV\] line in the mid infrared is measured is very poor. Using the available data for these objects we have confirmed that the contribution of this line to S<sub>234</sub> can be rather large. We have found very little \[SIV\] data for HII regions, GEHR and HII galaxies: the Orion nebula (Lester et al. 1979), Mrk 209 (Nollenberg et al. 2002) and a sample of objects in the Magellanic Clouds (Vermeij et al. 2002). Data on these objects are plotted in Figure 12 (S<sub>23</sub> in the left panel and S<sub>234</sub> in the right panel). Any improvement in the abundance calibration is difficult to quantify given the scarcity of data. At any rate, since no observations of the $`\lambda `$ 10.5 $`\mu `$ line exist for most objects, it would have to be calculated from photo-ionisation models which would make S<sub>234</sub> a semi-empirical parameter.
### 4.4 The S<sub>23</sub>/O<sub>23</sub> parameter
One fact that becomes evident from the examination of the different abundance parameters discussed above is that the validity of each one of them seems to be restricted to a given metallicity range. This means that it is necessary to have some a priori knowledge about the metallicity of an object or a sample of objects in order to choose the appropriate abundance indicator.
Traditional ways of doing this include an examination of the \[NII\]$`\lambda `$ 6584 ร
/H$`\alpha `$ ratio that can discriminate between objects with 12+log(O/H) higher or lower than about 8.0 (S89). More recently, the values of the O<sub>23</sub> and S<sub>23</sub> taken together have been used by DPM00 to discriminate between subsolar and oversolar abundances. However, in some cases, the interest is focused on the comparison of global abundances trends shown by different objects or abundance distributions over a wide range of metallicities. In those cases it would be desirable, to obtain abundances by means of the same calibrator so that comparisons are meaningful. This calibrator should be valid for the whole metallicity range.
The study of metallicity gradients over galaxy discs is one of the issues what could be improved in this way. For example, there are many different conclusions about the value of the oxygen abundance distribution in the well studied galaxy M101 . Different authors (Zaritsky, 1992; Scowen et al. 1992; Vila-Costas & Edmunds, 1992) have pointed to an increase in the slope of the gradient in the inner regions whereas other authors (Henry & Howard, 1995; Pilyugin, 2001b) obtain an exponential law throughout the whole disc. Kennicutt & Garnett (1996) have shown how the use of one calibration of O<sub>23</sub> or another leads to different conclusions.
In Figure 14 we represent the gradient of some of the parameters studied here as a function of the galactocentric distance to the center of M101 (data from Kennicutt & Garnett, 1996). In the upper panel the O<sub>23</sub> parameter is seen to increase with increasing galactocentric radius up to a value of 0.3 R<sub>0</sub> and then it remains almost constant. In the middle panel, the S<sub>23</sub> parameter shows the same behaviour as O<sub>23</sub> in the central regions of the disc but decreases with increasing galactocentric radius from 0.3 R<sub>0</sub> onwards. These two trends taken together point to the central disc regions of M101 ( R $`<`$ 0.3 R<sub>0</sub> ) being over-solar and thus lying on the upper branch of the O<sub>23</sub> and S<sub>23</sub> parameters. The outer regions of the disc would have under-solar abundances and lie on the lower branch of the S<sub>23</sub> calibration. Most of the regions in this regime ( R $`>`$ 0.3 R<sub>0</sub> ) lie on the turnover region of the O<sub>23</sub> calibration and therefore show an almost constant value of this parameter. In the lower panel of Figure 14, we can see that a combination of the two parameters, S<sub>23</sub>/O<sub>23</sub>, shows a continously decreasing trend through the disc of M101.
Using all the objects of our sample with measurements of the \[OII\], \[OIII\], \[SII\] and \[SIII\] lines and a direct determination of the metallicity, we have calibrated for the first time this new parameter (Figure 15 left panel). As can be seen, the relation remains single-valued for the whole range of metallicity, though it is non-linear, due probably to the contribution of \[SIV\] for high excitation regions of low metallicity. The addition of high metallicity objects is not possible due to the lack of data with he necessary auroral lines, but the fact that the parameter increases towards the inner parts of the disc of M101 (see Figure 14) suggests that it does not undergo any turnover at high metallicites. Of course, this parameter keeps the same sources of uncertainty as its two progenitors. In Figure 15 (right), where we compare the observational data with results from our photo-ionisation models, it can be seen that the value of S<sub>23</sub>/O<sub>23</sub> increases for low degrees of ionisation and lower effective temperatures.
Using the compiled data, we propose the following relation to derive oxygen abundances in the whole metallicity range in the absence of auroral emission line data:
$$12+\mathrm{log}(O/H)=9.09+1.03\mathrm{log}\left(\frac{S_{23}}{O_{23}}\right)0.23\left[\mathrm{log}\left(\frac{S_{23}}{O_{23}}\right)\right]^2$$
This relation is plotted along with data in the left panel of Figure 15. The residuals of this fit from the values of 12+log(O/H) deduced from the direct method, are plotted in the lower panel of Figure 16 and show a dispersion of 0.27 dex, comparable to that found for the N2 parameter.
The upper panel of Figure 17 shows the oxygen abundance gradient in M101. Solid circles correspond to oxygen abundances derived using the S<sub>23</sub>/O<sub>23</sub> parameter from the emission line data of Kennicutt & Garnett (1996); open circles correspond to abundances derived by the direct method (Kennicutt et al. 2003). The agreement between directly determined and empirically derived abundances is excellent.
Unfortunately, sulphur line intensity data is scarce and therefore it is difficult to assess the suitability of the S<sub>23</sub>/O<sub>23</sub> parameter as an abundance gradient indicator. We have found only two more galaxies with data covering a substantial part of the disc: M33 (data from Kwitter & Aller 1981 and Vรญlchez et al. 1988) and NGC300 (data from Deharveng et al. 1988 and Christensen et al. 1997). Their abundance gradients are shown in the middle and lower panels of Figure 17 respectively. Again, solid circles correspond to oxygen abundances derived using the S<sub>23</sub>/O<sub>23</sub> parameter and open circles correspond to abundances derived by the direct method. We have also added oxygen abundance data derived from the spectroscopic analysis of early B-type supergiant stars (Monteverde et al. 1988; Urbaneja et al. 2003). Although the discs of these two galaxies are not as well sampled as that of M101, the agreement between empirically and directly derived abundances is good and the agreement between nebular and stellar abundances is encouraging. Regarding the shape of the gradients, in all three cases an increase in slope for the central galactic regions is apparent, although in the cases of M33 and NGC300 nothing conclusive can be said. We think that more observations and more work along these lines would greatly help to derive the true abundance distributions across the discs of galaxies.
## 5 Summary and conclusions
In this work we have revised the different proposed oxygen abundance calibrations using a large compiled sample of observations comprising the emission lines of \[OII\], \[OIII\], \[SII\], \[SIII\] and \[NII\] for objects with oxygen abundances derived by the direct method. The data has been compared with results from a set of photo-ionisation models in order to seek an explanation to the sources of scatter in the calibrations. The direct calibration of any parameter with model results alone would lead to non quantifiable systematic errors and we consider strictly empirical calibrations to be much more reliable.
In Table 4 we summarize the main properties of the parameters studied, including their metallicity range of validity, and the uncertainty obtained for each calibration, understood as the standard deviation of the residuals of the deduced oxygen abundance and that derived through the direct method. The O<sub>23</sub> parameter is the most widely used due to the important role of oxygen in the cooling of the ionised gas and because the oxygen emission lines remain in the optical-far red part of the spectrum until redshift $``$ 1. This parameter presents a double-valued relation with metallicity and in the lower and upper branch (in the latter case, taking into account the strong dependence on effective temperature) the uncertainty remains below 0.2 dex. Nevertheless, the uncertainty in the turnover region (for 12+log(O/H) between 8.0 and 8.4) may reach almost an order of magnitude. No ways of improving this calibration further have been found in this work.
The best alternative to this parameter in this metallicity regime is S<sub>23</sub>. In objects where it is possible to observe the near-IR \[SIII\] lines (up to redshift $``$ 0.1), the metallicity can be deduced with less than 0.2 dex dispersion up to oxygen abundances 12+log(O/H) $``$ 8.9. The S<sub>23</sub> parameter also offers the advantage of being relatively independent of reddening. The contribution of the \[SIV\] emission line in the mid-IR is relevant only for high excitation objects and can be taken into account by means of the S<sub>234</sub> parameter. Unfortunately, there is, at the moment, very little data to calibrate it empirically, and avalaible photo-ionisation models fail to correctly reproduce the ionisation structure for sulphur (PMD03).
The N2 parameter, also reddening independent, is a good alternative for distant objects (up to $`z`$ $``$ 0.5) with intermediate metallicity ( 8.0 $``$ 12+log(O/H) $``$ 8.4). Nevertheless, this parameter suffers from uncertainties due to its dependence on ionisation parameter and N/O ratio. Besides, the calibration done by DTT02 works better for HII galaxies data but its application to higher metallicity regions carries a higher uncertainty and possibly overestimates the derived oxygen abundances.
Regarding the \[NII\]/\[OIII\], \[NII\]/\[OII\] and \[NII\]/\[SII\] parameters, these are only correlated with oxygen abundances at moderate to high metallicities and show an uncertainty similar to that of N2. However, as the main source of uncertainty for the calibrations involving the nitrogen lines is probably related to the N/O ratio, these calibrations could be improved with the use of \[NII\]/\[OII\] as a N/O calibrator which shows a dispersion of only 0.10 dex for HII galaxies.
Finally we have used the compiled sample of objects to produce for the first time a calibration for the S<sub>23</sub>/O<sub>23</sub> parameter, which could be useful to study variations over a wide rage of metallicities, as is the case for the discs of galaxies. This new parameter includes non-negligible uncertanties, inherited from its two predecesors, O<sub>23</sub> and S<sub>23</sub> but the results for the disc of M101 are encouraging.
## Acknowledgements
We would like to thank M. Castellanos, R. Terlevich, E. Terlevich, C. Esteban, E. Pรฉrez and D. Valls-Gabaud for very interesting discussions and suggestions and an anonimous referee for a careful revision of the manuscript. We would also like to acknowledge the thorough revision of the English done by Michael Taylor.
This work has been partially supported by DGICYT projects AYA-2000-0973 and AYA-2004-08260-C03-03.
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# Two pion mediated scalar isoscalar NN interaction in the nuclear medium
## 1 Introduction
The determination of the binding energy of nuclei starting from realistic $`NN`$ potentials is one of the subjects which has received permanent attention from the early days when the Brueckner-Bethe-Goldstone (BBG) equation introduced methods to overcome the strong repulsion of the nuclear forces at short distances.
At present, several many body techniques compete to accurately determine the binding energy of nuclear matter starting from the realistic NN potentials. One of them is the correlated basis functions (CBF) , which follows the line of the Hypernetted Chain Approach (FHNC) . Another one follows the traditional BBG approach , and, although costly numerically, the variational Monte Carlo method (MC) has allowed to make, in principle, exact calculations, although limited to nuclei with small or medium value of A . Methods like the selfconsistent treatment of the nucleon selfenergy have also introduced new advances in the field .
The need for three body forces has also been emphasized and the present status is that it is difficult to be quantitative on the strength of these forces, and usually they are parameterized in order to adjust the precise value of the binding energy . It has also been noted in that short range correlations play an important role when considering these three body forces.
A common feature of these approaches is that they start from the realistic nucleon nucleon interaction, obtained from fits to NN data and deuteron data. They use hence the free NN interaction as input. One of the important ingredients of this interaction is the one pion exchange (OPE). However, from detailed studies of the pion nuclear interaction it is well known that the pion properties in the nuclear medium are sizably renormalized .
There is also the question of the intermediate range attraction, which is basic in the binding of nuclei. Models for this interaction would contain $`\sigma `$ exchange, uncorrelated two pion exchange and omega exchange . In as much as the pion properties are changed in the medium, so should the two pion exchange be modified. Medium effects in the two-pion exchange have been investigated in early works like in Ref. restricting themselves in this case to a subset of two pion exchange diagrams with no $`\mathrm{\Delta }`$-isobar intermediate states, by including Pauli blocking in the intermediate nucleons.
The medium modification of the two pion exchange got a new impulse after the models of the $`\pi \pi `$ interaction in the medium showed large modifications , later on softened by the introduction of chiral constraints . The implementation of this medium modified $`\pi \pi `$ interaction in the correlated two pion exchange $`NN`$ potential increased appreciably the $`NN`$ attraction in nuclear matter. This was partially reduced by the consideration of the chiral constraints in Refs. . In these latter references the importance of short range correlations which modify the $`\pi `$ nucleus selfenergy was already discussed. The further use of medium modified vector meson masses led to improvements in the nuclear matter saturation curve.
A new perspective into this problem has been made possible by studies of meson meson interaction within chiral unitary approaches which allow to improve the description of the correlated two pion exchange $`NN`$ interaction , as an alternative to the conventional $`\sigma `$-meson exchange interaction. This picture of the $`\sigma `$ exchange was mandatory after extensive studies showing that the $`\sigma `$ is not a genuine resonance, made up of $`q\overline{q}`$ but just the manifestation of a pion - pion resonance state created by the interaction of the pions, what is called a dynamically generated resonance. This shows up naturally within the context of chiral unitary approaches which use the input of the chiral Lagrangians for the meson meson interaction and extends chiral perturbation theory to implement exactly unitarity in coupled channels . This means the $`\sigma `$ exchange inside a nuclear medium will also be modified as direct consequence of the change of the pion properties.
Our aim in the present paper is to start from this new picture for the $`\sigma `$ exchange, use also the standard approach for the uncorrelated two pion exchange and modify these in the nuclear medium to see what changes one finds from these sources. Further improvements come from the consideration of short range correlations not only in the pion selfenergy but also in the vertex functions appearing in the model. The changes obtained are moderate, thanks to the simultaneous consideration of these nuclear short range effects in the calculation. In the absence of these, the renormalization of the $`NN`$ interaction is huge. Yet, even the moderate results obtained are large enough to motivate further calculations of the nuclear binding and other properties of matter.
The paper is organized as follows. In Sec. II we provide those elements of the chiral Lagrangian which are relevant for the present calculations and briefly discuss peculiarities of the finite baryon density. In Sec. III we consider the modification of the one pion exchange $`NN`$ force in the nuclear medium and in Sec. IV we discuss the propagation of two pions in the nuclear matter. Sec. V and VI are devoted to the in-medium two pion exchange in the scalar-isoscalar channel, both, correlated and uncorrelated. The technical details are relegated to the Appendix.
## 2 Effective Lagrangian
In this section we will briefly specify those elements of the effective chiral Lagrangian in the meson-baryon sector which are relevant for the subsequent calculations.
The effective chiral Lagrangian is written as the sum of a purely mesonic Lagrangian $`_M`$ and the baryonic Lagrangian $`_B`$
$$_{eff}=_M+_B$$
(1)
Both are organized in a derivative and quark mass expansion. The lowest order mesonic Lagrangian $`_2`$ is given by
$$_2=\frac{f_\pi ^2}{4}^\mu U^{}_\mu U+\chi U^{}+U\chi ^{}$$
(2)
and contains the most general low-energy interactions of the pseudo-scalar meson octet. In Eq. (2) the symbol $`\mathrm{}`$ indicates the trace in flavor space, the Goldstone fields are collected in a unitary matrix $`U`$, $`f_\pi 93`$ MeV is the pseudoscalar decay constant and the leading symmetry-breaking term $`\chi `$ is linear in the quark masses. For $`SU\left(2\right)`$ and in the isospin limit $`\chi =\text{diag}(m_\pi ^2,m_\pi ^2)`$. The lowest order baryon octet - meson octet Lagrangian reads
$$_B^{\left(1\right)}=\overline{B}\left(i\gamma ^\mu D_\mu m_B\right)B+\frac{D}{2}\overline{B}\gamma ^\mu \gamma _5\{u_\mu ,B\}+\frac{F}{2}\overline{B}\gamma ^\mu \gamma _5[u_\mu ,B]$$
(3)
where the brackets $`\left[\mathrm{}\right]`$ and $`\left\{\mathrm{}\right\}`$ denote commutators and anti-commutators, respectively. The covariant derivative of the $`SU\left(3\right)`$ baryon matrix $`B`$ is defined as
$$D_\mu B=_\mu B+[\mathrm{\Gamma }_\mu ,B]$$
(4)
In the absence of external field Eqs. (3) and (4) involve other quantities
$$u=\sqrt{U},u^\mu =iu^{}^\mu Uu^{},\mathrm{\Gamma }_\mu =\frac{1}{2}\left(u^{}_\mu u+u_\mu u^{}\right)$$
(5)
The $`SU\left(3\right)`$ axial vector coupling constants are determined by neutron and hyperon $`\beta `$-decay. One finds $`F0.51`$, $`D0.76`$ and the axial coupling constant is $`g_A=F+D1.27`$. In the $`SU\left(2\right)`$ limit the Lagrangian simplifies to
$$_N^{\left(1\right)}=\overline{\psi }\left(i\gamma ^\mu D_\mu m_N\right)\psi +\frac{D+F}{2}\overline{\psi }\gamma ^\mu \gamma _5u_\mu \psi $$
(6)
where $`\psi `$ is a two component Dirac field $`\psi =(p,n)^T`$.
In the pion-nucleon sector the chiral Lagrangian (6) constrains all possible interactions of the pion fields with fermions at the lowest chiral order we are working. For instance, using the exponential parameterization of the unitary matrix $`U`$
$$U=\mathrm{exp}\left[i\frac{\mathrm{\Phi }}{f_\pi }\right],\mathrm{\Phi }=๐๐
$$
(7)
where $`๐`$ are the Pauli operators and expanding $`u`$ and $`u^{}`$ we obtain the pion-nucleon couplings, including up to three pion fields
$$u_\mu =\frac{_\mu \mathrm{\Phi }}{f_\pi }+\frac{1}{24f_\pi ^3}\left(_\mu \mathrm{\Phi }\mathrm{\Phi }^22\mathrm{\Phi }_\mu \mathrm{\Phi }\mathrm{\Phi }+\mathrm{\Phi }^2_\mu \mathrm{\Phi }\right)+๐ช\left(\mathrm{\Phi }^5\right)$$
(8)
If we supplement Eq. (8) with the lowest order interaction of pions as provided by Eq. (2)
$$_2=\frac{1}{48f_\pi ^2}\left(_\mu \mathrm{\Phi }\mathrm{\Phi }\mathrm{\Phi }^\mu \mathrm{\Phi }\right)^2+m_\pi ^2\mathrm{\Phi }^4$$
(9)
we arrive to the set of Feynman graphs shown in Fig. 1. Here, in addition to the standard $`\pi NN`$ vertex, Fig. 1a, the chiral perturbation theory generates the contact term of Fig. 1b (see Appendix A) and the pion pole term of Fig. 1c. These last two diagrams are the basic element in the description of $`\pi N\pi \pi N`$ reaction near the threshold . The appearance of the pole terms and contact $`3\pi NN`$ interactions at the same order of the chiral expansion is crucial for the in-medium calculations where due to the partial cancellations the physical amplitudes become independent of the parameterization of $`U`$ matrix even in the presence of the nuclear background in accord with the equivalence theorem . One can see explicitely, that the contact term (b) cancels exactly the off shell part of the pion pole term (c) coming from the $`\left(q^2m_\pi ^2\right)`$ part of the $`\pi \pi \pi \pi `$ vertex . We shall also see that the off shell part of the $`\pi \pi \pi \pi `$ amplitude cancel exactly with other terms when we perform the calculation of the $`NN`$ interaction in the nuclear medium.
## 3 One-pion exchange at finite density
We start with the one-pion exchange $`NN`$ potential (OPEP). The typical diagrams modifying it are shown in Fig. 2 where the propagation of exchanged pions is distorted by interactions with nucleons forming the Fermi sea. The fermionic bubbles describe the decay of the pion in $`ph`$ and $`\mathrm{\Delta }h`$ states and take into account the conventional nuclear matter polarization effects. All these diagrams are responsible for the interaction of two nucleons in the particle-particle ladder channel with the in-medium virtual excitations. In Fig. 2 the first graph is the vacuum contribution. The second and third diagrams correspond to the $`ph`$ RPA series and the last diagrams accounts for the excitation of $`2p2h`$ states and the contribution of the $`S`$-wave optical potential, respectively. The $`p`$-wave pion self energy is given by
$$\mathrm{\Pi }(k,\rho )=\left(\frac{D+F}{2f_\pi }\right)^2๐^2๐ฐ(k,\rho )\left[1\left(\frac{D+F}{2f_\pi }\right)^2g^{}๐ฐ(k,\rho )\right]^1$$
(10)
where $`g^{}=0.6`$ is the Landau-Migdal parameter , $`\rho `$ is the nuclear matter density and $`๐ฐ(k,\rho )=๐ฐ^d(k,\rho )+๐ฐ^c(k,\rho )`$ is the Lindhard function accounting for the direct and crossed contributions of $`ph`$ and $`\mathrm{\Delta }h`$ excitations with the normalization of the appendix of Ref. .
The OPEP in the momentum space takes the form
$`iV_{OPEP}(๐,\rho )`$ $`=`$ $`iW(๐,\rho )\widehat{๐}_i\widehat{๐}_j\sigma _1^i\sigma _2^j๐_1๐_2`$ (11)
where we have defined
$$W(๐,\rho )=\left(\frac{D+F}{2f_\pi }\right)^2๐^2F^2\left(๐\right)\stackrel{~}{D}_\pi (0,๐)$$
(12)
where $`\stackrel{~}{D}_\pi (q_0,๐)`$ is the pion propagator in the medium
$$i\stackrel{~}{D}_\pi \left(k\right)=\frac{i}{k^2m_\pi ^2\mathrm{\Pi }(k,\rho )+i0^+}$$
(13)
and $`F\left(๐\right)`$ stands for a monopole form factor $`\mathrm{\Lambda }^2/\left(\mathrm{\Lambda }^2+๐^2\right)`$ with the cut off scale $`\mathrm{\Lambda }=1`$ GeV, $`๐`$ is a momentum transfer and $`\widehat{๐}=๐/\left|๐\right|`$. Note, that OPEP depends on the real part of the polarization operator only, since for $`q=(0,๐)`$ one has $`\text{Im}\mathrm{\Pi }(0,๐,\rho )=0`$. The well-known vacuum $`NN`$ amplitude is recovered in Eq. (12) by setting $`\mathrm{\Pi }=0`$ or in the limit $`\rho =0`$.
Our results for $`W(๐,\rho )`$ are presented in Fig. 3 (left). The standard vacuum behavior is shown by the solid curve. The dashed curve represents the modified OPEP at normal nuclear matter density were we observe an additional strong increase of the strength associated with the attractive excitation of $`ph`$ and $`\mathrm{\Delta }h`$ collective states, with $`ph`$ playing a dominant role. Their individual contributions are shown by dot-dot-dashed and dot-dashed curves, respectively.
The analytic properties of the in-medium pion propagator which we use here may be verified by using the dispersion representation for the Green function
$`i\stackrel{~}{D}_\pi (k_0,๐)=i{\displaystyle \frac{2}{\pi }}{\displaystyle _0^{\mathrm{}}}๐xx{\displaystyle \frac{\text{Im}\stackrel{~}{D}_\pi (x,๐)}{k_0^2x^2+i0^+}}`$ (14)
In this case the OPEP of Eq. (12) can be written in terms of the absorptive part of the pion propagator only
$`W(๐,\rho )={\displaystyle \frac{2}{\pi }}\left({\displaystyle \frac{D+F}{2f_\pi }}\right)^2๐^2F^2\left(๐\right){\displaystyle \underset{0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{dx}{x}}\left[\text{Im}\stackrel{~}{D}_\pi (x,๐)\right]`$ (15)
The causality requires that both equations must produce the same result. Indeed as one can see in Fig. 3 the curves calculated with dispersive and absorptive parts of the pion propagator are practically indistinguishable. We would like to note that a strong modification of the OPEP at finite baryonic density observed here is not new and was predicted long time ago by Migdal . There it was also shown that in-medium modified OPEP helps to explain the unnatural parity states in finite nuclei, for instance, the shift of $`0^{}`$ state in closed shell nuclei.
In the right panel of Fig. 3 we show our combined plot for a few densities $`\rho _0,\rho _0/2`$ and $`\rho _0/4`$. Here, we also show the sensitivity of our results to the value of the Landau-Migdal parameter $`g^{}`$. We find that the increase of $`g^{}`$ from 0.6 to 0.7 makes the OPEP softer. This fact suggests that the proper treatment of the $`NN`$ short range correlations is important for understanding the in-medium properties of the OPEP.
At this point we would like to mention that in a realistic calculation one will have to add strong repulsive forces at short distances. This can be done in a straightforward way using any of many body schemes discussed in the introduction. The correlations of this part of the interaction would effectively modulate the $`\pi `$ exchange interaction, introducing the correlation parameter $`g^{}`$ . The denominator in Eq. (13) takes into account this effect between $`p`$-wave bubbles in the diagrams of Fig. 2, but not between the external nucleon and the contiguous bubble. To account for this we make the separation between the longitudinal and transverse parts of the pion effective interaction
$$\left(\frac{D+F}{2f_\pi }\right)^2F\left(๐\right)^2\frac{q_iq_j}{q_0^2๐^2m_\pi ^2+i0^+}๐ฑ_l\left(q\right)\widehat{q}_i\widehat{q}_j+๐ฑ_t\left(q\right)\left(\delta _{ij}\widehat{q}_i\widehat{q}_j\right)$$
(16)
where $`\widehat{q}_i`$ is the Cartesian component of the unit vector $`\widehat{๐}=๐/\left|๐\right|`$ and
$$๐ฑ_l\left(q\right)=\left(\frac{D+F}{2f_\pi }\right)^2\left[\frac{๐^2}{q_0^2๐^2m_\pi ^2+i0^+}+g^{}\right]F\left(๐\right)^2$$
(17)
$$๐ฑ_t\left(q\right)=\left(\frac{D+F}{2f_\pi }\right)^2g^{}F\left(๐\right)^2$$
(18)
When we perform the sum of diagrams in Fig. 2 then we get
$`V_{OPEP}(q,\rho )`$ $`=`$ $`\left[{\displaystyle \frac{๐ฑ_l\left(q\right)}{1๐ฐ(q,\rho )๐ฑ_l\left(q\right)}}\widehat{q}_i\widehat{q}_j+{\displaystyle \frac{๐ฑ_t\left(q\right)}{1๐ฐ(q,\rho )๐ฑ_t\left(q\right)}}\left(\delta _{ij}\widehat{q}_i\widehat{q}_j\right)\right]\sigma _1^i\sigma _2^j๐_1๐_2`$ (19)
$`=`$ $`W_0(q,\rho )๐_1๐_2๐_1๐_2+W_2(q,\rho )\left(๐_1\widehat{๐}๐_2\widehat{๐}{\displaystyle \frac{1}{3}}๐_1๐_2\right)๐_1๐_2`$
$`W_0(q,\rho )`$ $`=`$ $`{\displaystyle \frac{1}{3}}\left[{\displaystyle \frac{๐ฑ_l\left(q\right)}{1๐ฐ(q,\rho )๐ฑ_l\left(q\right)}}+{\displaystyle \frac{2๐ฑ_t\left(q\right)}{1๐ฐ(q,\rho )๐ฑ_t\left(q\right)}}\right]`$ (20)
$`W_2(q,\rho )`$ $`=`$ $`{\displaystyle \frac{๐ฑ_l\left(q\right)}{1๐ฐ(q,\rho )๐ฑ_l\left(q\right)}}{\displaystyle \frac{๐ฑ_t\left(q\right)}{1๐ฐ(q,\rho )๐ฑ_t\left(q\right)}}`$ (21)
where we have explicitely separated the central and the tensor parts of the interaction. In Fig. 4 we show the results for the central $`W_0`$ (left panel) and tensor $`W_2`$ (right panel) parts (omitting the spin-isospin operators) as a function of the baryonic density. Note that in the limit $`\rho 0`$, $`W_2=W`$ where $`W`$ is given by Eq. 12.
## 4 Two pions in the medium
In this section we turn to the dynamics of the two pion system in the nuclear medium. Our main interest here is the propagation of two $`S`$-wave pion pairs. As it is well known in the $`S`$-wave scattering the use of the proper unitarization schemes lead to the generation of the $`\sigma `$-meson. Later on we will use this result for the in-medium $`NN`$ interaction mediated by exchange of two correlated pions in the scalar-isoscalar $`\sigma `$-meson channel.
For the $`\pi _a\pi _b\pi _c\pi _d`$ scattering process, defined by the Cartesian isospin indices $`a,\mathrm{}`$, the use of the standard $`\chi `$PT procedure in expanding the $`_2`$ of Eq. (2) to order $`๐ช\left(๐
^4\right)`$ results in the tree level contact interaction
$`iV_{\pi \pi }^{abcd}=\delta _{ab}\delta _{cd}A\left(s\right)+\delta _{ac}\delta _{bd}A\left(t\right)+\delta _{ad}\delta _{bc}A\left(u\right),`$ (22)
where
$$A\left(s\right)=\frac{i}{f_\pi ^2}\left(sM_\pi ^2\frac{1}{3}\underset{i=a,b,c,d}{}\mathrm{\Lambda }_i\right)+๐ช\left(q^4\right),$$
(23)
and $`\mathrm{\Lambda }_i=k_i^2M_\pi ^2`$ is the off-shell part of the invariant $`\pi \pi `$ amplitude. At this order of the pion field expansion the isoscalar $`S`$-wave $`\pi \pi `$ partial amplitude ($`L=0`$) is obtained from the standard decomposition
$$V_{\pi \pi }^{L,I=0}=\frac{1}{2}\frac{1}{\left(\sqrt{2}\right)^\alpha }_1^1d\mathrm{cos}\theta P_L\left(\mathrm{cos}\theta \right)V_{\pi \pi }^{I=0}\left(\theta \right)$$
(24)
where $`P_L\left(\mathrm{cos}\theta \right)`$ are the Legendre polynomials and $`\left(\sqrt{2}\right)^\alpha `$ accounts for the statistical factor occurring in states with identical particles: $`\alpha =2`$ for $`\pi \pi \pi \pi `$ in the unitary normalization of the states . The tree level scalar-isoscalar $`\pi \pi `$ scattering amplitude is
$$V_{\pi \pi }^{L=I=0}=\frac{1}{f_\pi ^2}\left(s\frac{M_\pi ^2}{2}\frac{1}{3}\underset{i}{}\mathrm{\Lambda }_i\right).$$
(25)
In Eq. (25) the off shell part depends on choice of $`U`$ and is equal to zero for on mass shell pions.
Following Ref. and using the Bethe-Salpeter equation we unitarize the $`S`$-wave $`\pi \pi `$ scattering amplitude (see Fig. 5)
$$V_{\pi \pi }^{L=I=0}T_{\pi \pi }^{L=I=0}=\left[f_\pi ^2\left(s\frac{M_\pi ^2}{2}\right)^1G_{\pi \pi }\left(s\right)\right]^1,$$
(26)
where $`G_{\pi \pi }\left(s\right)`$ is a scalar two-pion loop function
$$G_{\pi \pi }\left(s\right)=i\frac{d^4k}{\left(2\pi \right)^4}\frac{1}{\left[\left(Pk\right)^2M_\pi ^2+i0^+\right]\left(k^2M_\pi ^2+i0^+\right)}$$
(27)
where $`P^\mu =(P_0,๐ท)`$ and $`s=P^2`$. The $`G_{\pi \pi }\left(s\right)`$ function is analytic with a cut along the positive real axis starting at the $`\pi \pi `$ threshold. Note that, Eq. (26) contains a pole in the second Riemann sheet corresponding to the $`\sigma `$-meson with mass and width $`Mi\mathrm{\Gamma }/2450i221`$ MeV.
In the nuclear medium the $`S`$-wave $`\pi \pi `$ scattering amplitude and therefore the $`\sigma `$-meson get renormalized and explicit calculations were done in Refs. . The diagrams at one loop level are shown in Fig. 6. For instance, the amplitudes corresponding to the insertion of fermion bubbles in the upper meson line are given by
$`V_\text{I}^{up}\left(s\right)`$ $`=`$ $`i\left({\displaystyle \frac{D+F}{2f_\pi }}\right)^2{\displaystyle \frac{d^4k}{\left(2\pi \right)^4}๐^2๐ฐ\left(k\right)D_\pi ^2\left(k\right)D_\pi \left(Pk\right)\left[V_{on}\left(s\right)+\frac{1}{3f_\pi ^2}\underset{i}{}\mathrm{\Lambda }_i\right]^2}`$ (28)
$`V_{\text{II}}^{up}\left(s\right)`$ $`=`$ $`{\displaystyle \frac{i2}{3f_\pi ^2}}\left({\displaystyle \frac{D+F}{2f_\pi }}\right)^2{\displaystyle \frac{d^4k}{\left(2\pi \right)^4}๐^2๐ฐ\left(k\right)D_\pi \left(k\right)D_\pi \left(Pk\right)\left[V_{on}\left(s\right)+\frac{1}{3f_\pi ^2}\underset{i}{}\mathrm{\Lambda }_i\right]}`$ (29)
$`V_{\text{III}}^{up}\left(s\right)`$ $`=`$ $`{\displaystyle \frac{i}{9f_\pi ^4}}\left({\displaystyle \frac{D+F}{2f_\pi }}\right)^2{\displaystyle \frac{d^4k}{\left(2\pi \right)^4}๐^2๐ฐ\left(k\right)D_\pi \left(Pk\right)}`$ (30)
As was shown in Refs. , in the center of mass frame of the two pions the off shell part ($`\mathrm{\Lambda }`$ terms) in $`V_I^{up}`$ cancels exactly the $`V_{II}^{up}`$ and $`V_{III}^{up}`$ terms. And one is left only with the diagrams of type (I) but with the on shell $`\pi \pi `$ amplitude
$`V^{up}\left(s\right)`$ $`=`$ $`i\left({\displaystyle \frac{D+F}{2f_\pi }}\right)^2V_{on}^2\left(s\right){\displaystyle \frac{d^4k}{\left(2\pi \right)^4}๐^2๐ฐ\left(k\right)D_\pi ^2\left(k\right)D_\pi \left(Pk\right)}`$ (31)
It is straightforward to iterate the $`ph`$ and $`\mathrm{\Delta }h`$ excitations in Fig. 6 (I) and the loop function, $`T\left(s\right)`$, is given by
$$T\left(s\right)=V_{on}^2\left(s\right)\stackrel{~}{G}_{\pi \pi }\left(s\right)$$
(32)
where $`\stackrel{~}{G}_{\pi \pi }\left(s\right)`$ is in-medium modified scalar loop integral
$$\stackrel{~}{G}_{\pi \pi }\left(s\right)=i\frac{d^4k}{\left(2\pi \right)^4}\stackrel{~}{D}_\pi \left(k\right)\stackrel{~}{D}_\pi \left(Pk\right)$$
(33)
Using the spectral representation for the in-medium pion propagators, Eq. (13), we get for the loop function
$`\stackrel{~}{G}_{\pi \pi }\left(P\right)={\displaystyle \frac{2}{\pi ^2}}{\displaystyle \frac{d๐}{\left(2\pi \right)^3}\underset{0}{\overset{\mathrm{}}{}}๐x\left[\text{Im}\stackrel{~}{D}_\pi (x,๐)\right]\underset{x}{\overset{\mathrm{}}{}}๐y\left[\text{Im}\stackrel{~}{D}_\pi (yx,๐ท๐)\right]}`$
$`\times {\displaystyle \frac{y}{\left(P_0+y\right)\left(P_0y+i0^+\right)}}`$ (34)
We refer to Ref. where different aspects of $`S`$-wave $`\pi \pi `$ scattering in the nuclear medium are discussed and also the behavior of the $`\sigma `$-meson mass and width at finite baryonic density is addressed. But here, we would like to illustrate the impact of the nuclear medium on the $`\pi \pi `$ system. For that consider the imaginary part of the loop function in the $`\pi \pi `$ center of mass frame $`P^\mu =(P_0,0)`$ with $`P_0=\sqrt{s}`$. This situation is relevant for the in-medium $`\pi \pi `$ scattering and contains the proper information about the dynamics of the pole position of $`\sigma `$ at finite density. Our results for the imaginary part of the scalar loop function for several densities are shown in Fig. 7(left panel). The solid curve correspond to the vacuum loop function which can be obtained from Eq. (4) by substituting the imaginary part of in-medium pion propagators by their vacuum expressions
$$\text{Im}G_{\pi \pi }\left(\sqrt{s}\right)=\frac{1}{16\pi }\sqrt{1\frac{4m_\pi ^2}{s}}$$
(35)
which is proportional to the two-body phase space of two particles (pions). As one can see in Fig. 7 (left) the effect of the medium is remarkable due to the increase of available for pions phase space because of additional pion decay branches like $`ph`$, $`\mathrm{\Delta }h`$ and $`2p2h`$.
The kinematics relevant for the $`NN`$ force, where two nucleons interact by exchange of mesons is defined by the moving reference frame where $`P^\mu =(0,๐)`$. In this case the $`\pi \pi `$ interaction in the medium has to be generalized from the results in since there $`P_00`$ but $`๐=0`$. We repeat all the steps that led to cancellations in and find again the same cancellations as before but with some remnant $`๐^2`$ dependent terms vanishing in the limit $`๐^2=0`$. To evaluate these terms we simplify the calculation assuming $`๐`$ relatively small (this is fine for momenta below the Fermi momentum). Concretely, we assume
$$1)๐^2/k_{max}^21,$$
where $`k_{max}`$ is the cut off in the three momentum (of the order of $`1`$ GeV in ) that one uses to regularize the $`G_{\pi \pi }`$ function. We also assume
$$2\left)๐^2D_\pi \right(k)1$$
which implies $`k_{max}^2m_\pi ^2`$ as it is the case. And 3) we expand $`D_\pi \left(k\right)`$ in terms of $`D_\pi \left(Pk\right)`$ and vice-versa to relate different terms. After all this is done we find that the corrections can be taken into account by means of the change in the Lindhard function
$$๐ฐ\left(k\right)๐ฐ\left(k\right)\left[1+\frac{๐^2}{3}\left(s\frac{m_\pi ^2}{2}\right)^1+\frac{๐^4}{3}\left(s\frac{m_\pi ^2}{2}\right)^2\right]$$
(36)
The expression in brackets in Eq. (36) is $`1`$ for very small $`q`$ and also $`1`$ for $`๐^2m_\pi ^2/2`$. Hence with very good approximation we can take the bracket equal to unity and thus there are no other corrections to be done to the result of except the obvious one of changing $`s๐^2`$. For these value of $`s`$ the $`G_{\pi \pi }`$ is only real, contrary to the case studied in . Results for $`G_{\pi \pi }\left(๐^2\right)`$ can be seen in Fig. 7 (right panel) for different nuclear densities $`\rho _0`$, $`\rho _0/2`$ and $`\rho _0/4`$. We can see that the corrections are sizable particularly at small values of $`\left|๐\right|`$.
## 5 In-medium renormalization of the correlated two pion ($`\sigma `$-meson) exchange
The correlated two pion exchange (CrTPE) in the scalar-isoscalar channel, the equivalent to a $`\sigma `$ exchange in meson exchange models , or correlated two pion exchange in the dispersion relations in was studied in within the context of chiral Lagrangians. One starts from the diagrams of Fig. 8, where the $`\pi \pi \pi \pi `$ scattering shows the off shell ambiguities. To avoid these ambiguities with isoscalar exchange it was stated in Ref. that one must include the subset of diagrams of Fig. 9 to find cancellations of the off shell $`\pi \pi `$ isoscalar amplitude. This statement was rigorously verified in Ref. , where it was shown that the consideration of these subset of chiral diagrams, Fig. 9, including the contact $`3\pi NN`$ interactions, results in the cancellation of the off-shell part of the $`\pi \pi `$ amplitude, and the on-shell part of the $`\pi \pi `$ amplitude can be factorized out from the loop integrals. One step forward was given in , where iteration of the $`\pi \pi `$ interaction, through the Bethe-Salpeter equation, was done by means of which a simple analytical expression was obtained for the correlated two pion exchange in the scalar-isoscalar channel
$$V_\sigma \left(t\right)=6V^2\left(t\right)\left[f_\pi ^2\left(t\frac{M_\pi ^2}{2}\right)^1G_{\pi \pi }\left(t\right)\right]^1$$
(37)
where $`t=๐^2`$ in the $`NN`$ c.m. frame. The vertex function $`V\left(t\right)=V_N\left(t\right)+V_\mathrm{\Delta }\left(t\right)`$ for the triangle loop with two mesons and one baryon propagator, including $`N`$ and $`\mathrm{\Delta }`$ intermediate states, is evaluated in using a cut off in $`\left|๐\right|`$ of about 1 GeV. Note that, the bracket in Eq. (37) contains a pole in the $`s`$-channel corresponding to the $`\sigma `$-meson. This restores the relation to the $`\sigma `$-meson exchange, which now enters the formalism as a dynamical resonance in the $`\pi \pi `$ system (see also related discussions in Refs. ).
The diagrams responsible for the renormalization of CrTPE in the nuclear medium are shown in Fig. 10. There, in analogy to what was done for the $`\pi \pi `$ interaction we include $`\pi `$ selfenergy corrections as well as vertex corrections. We find that the cancellation of the off shell part of the vacuum $`\sigma `$-exchange discussed in is also exact at finite baryonic density but for zero momentum transfer only (see Appendix B). The results of the derivation for $`๐0`$ can be summarized as follows:
1) The expressions for the triangle vertex functions $`\stackrel{~}{V}_N`$ and $`\stackrel{~}{V}_\mathrm{\Delta }`$ of Ref. are obtained in the same way replacing the two free pion propagators by the renormalized ones.
2) The Lindhard function entering the pion selfenergy is changed to
$$๐ฐ\left(k\right)๐ฐ\left(k\right)\left[1\frac{1}{6}\frac{๐^2}{๐^2+m_\pi ^2/2}\right]$$
(38)
to account for the corrections obtained at $`๐0`$. These corrections are negligible for small $`\left|๐\right|`$ and for $`\left|๐\right|>m_\pi `$ of the order of $`15\%`$, hence moderate in all cases.
3) The final expression for the potential $`V_\sigma `$ is given by Eq. (37), which accounts for the $`\pi \pi `$ rescattering, by substituting $`G_{\pi \pi }\left(t\right)`$ by $`\stackrel{~}{G}_{\pi \pi }\left(t\right)`$ of Eq. (4) and taking the expression for the in-medium vertex function $`V\left(t\right)=\stackrel{~}{V}_N\left(t\right)+\stackrel{~}{V}_\mathrm{\Delta }\left(t\right)`$ where
$`\stackrel{~}{V}_N\left(t\right)`$ $`=`$ $`{\displaystyle \frac{2\kappa _N}{\pi ^2}}{\displaystyle \frac{d^3๐}{\left(2\pi \right)^3}\frac{M_N}{E\left(๐\right)}\left(๐^2+๐๐\right)\underset{0}{\overset{\mathrm{}}{}}๐x\left[\text{Im}\stackrel{~}{D}_\pi (x,๐)\right]}`$ (39)
$`\times {\displaystyle \underset{0}{\overset{\mathrm{}}{}}}dy{\displaystyle \frac{x+y+E\left(๐\right)M_N}{\left(x+y\right)\left(x+E\left(๐\right)M_N\right)\left(y+E\left(๐\right)M_N\right)}}\left[\text{Im}\stackrel{~}{D}_\pi (y,๐+๐)\right]`$
$`\stackrel{~}{V}_\mathrm{\Delta }\left(t\right)`$ $`=`$ $`{\displaystyle \frac{2\kappa _\mathrm{\Delta }}{\pi ^2}}{\displaystyle \frac{4}{9}}{\displaystyle \frac{d^3๐}{\left(2\pi \right)^3}\frac{M_\mathrm{\Delta }}{E_\mathrm{\Delta }\left(๐\right)}\left(๐^2+๐๐\right)\underset{0}{\overset{\mathrm{}}{}}๐x\left[\text{Im}\stackrel{~}{D}_\pi (x,๐)\right]}`$ (40)
$`\times {\displaystyle \underset{0}{\overset{\mathrm{}}{}}}dy{\displaystyle \frac{x+y+E_\mathrm{\Delta }\left(๐\right)M_N}{\left(x+y\right)\left(x+E_\mathrm{\Delta }\left(๐\right)M_N\right)\left(y+E_\mathrm{\Delta }\left(๐\right)M_N\right)}}\left[\text{Im}\stackrel{~}{D}_\pi (y,๐+๐)\right]`$
The coupling constants $`\kappa _n`$ are defined by
$$\kappa _N=\left(\frac{D+F}{2f_\pi }\right)^2,\kappa _\mathrm{\Delta }=\left(\frac{3}{\sqrt{2}}\frac{D+F}{2f_\pi }\right)^2$$
(41)
There is still one more correction to be done in order to account for short range correlations. So far we have replaced the free pion propagator by the renormalized one of Eq. (13). However, since the interaction between the $`ph`$ bubble and the external nucleons is affected by correlations one should take $`๐ฑ_l`$ (see Eq. (17)) instead of $`\left[\left(D+F\right)^2/2f_\pi \right]^2D_\pi `$. Thus we would get the series
$`D_\pi \left(k\right)+D_\pi \left(k\right)๐ฐ\left(k\right)๐ฑ_l\left(k\right)+D_\pi \left(k\right)๐ฐ\left(k\right)๐ฑ_l\left(k\right)๐ฐ\left(k\right)๐ฑ_l\left(k\right)+\mathrm{}={\displaystyle \frac{D_\pi \left(k\right)}{1๐ฐ\left(k\right)๐ฑ_l\left(k\right)}}`$
$`={\displaystyle \frac{\stackrel{~}{D}_\pi \left(k\right)}{1\left({\displaystyle \frac{D+F}{2f_\pi }}\right)^2g^{}F^2\left(k\right)๐ฐ\left(k\right)}}`$ (42)
where $`D_\pi `$ and $`\stackrel{~}{D}_\pi `$ are the free and dressed pion propagators, respectively. Hence, the expressions for $`\stackrel{~}{V}_N`$ and $`\stackrel{~}{V}_\mathrm{\Delta }`$ get modified by including inside the $`d^3๐`$ integral the factors
$$\frac{1}{1\left({\displaystyle \frac{D+F}{2f_\pi }}\right)^2g^{}F^2\left(k\right)๐ฐ\left(k\right)}\times \frac{1}{1\left({\displaystyle \frac{D+F}{2f_\pi }}\right)^2g^{}F^2\left(k+q\right)๐ฐ\left(k+q\right)}$$
(43)
Our results for $`\sigma `$-meson exchange in the momentum space are shown if Fig. 11 for both cases, no short range correlations (left) and with short range correlations (right).
## 6 Uncorrelated two pion exchange
In this section we consider another sort of intermediate distance contributions to the $`NN`$ force generated by the uncorrelated two pion exchange. The material presented here for the vacuum $`NN`$ scattering is standard and we merely generalize it to the nuclear medium.
In the perturbative expansion of the $`NN`$ force we must to take into account the planar and crossed box diagrams shown in Figs. 12 and 13, respectively. We will discuss the scalar-isoscalar part of this contributions only. The contribution of the isovector exchange is small and can be found, for instance, in Ref.
In vacuum, the expression for the planar box diagrams, Fig. 12, including the nucleon pole and $`N\mathrm{\Delta }`$ and $`\mathrm{\Delta }\mathrm{\Delta }`$ intermediate states reads
$$iV_{NN}^{\left(P\right)}=\underset{n,m}{}\kappa _n\kappa _m\frac{d^3๐}{\left(2\pi \right)^3}\left[๐บ_n^{\left(1\right)}\left(๐+๐\right)\right]\left[๐บ_n^{\left(1\right)}๐\right]^{}\left[๐บ_m^{\left(2\right)}\left(๐+๐\right)\right]\left[๐บ_m^{\left(2\right)}๐\right]^{}\left\{_{nm}^{\left(P\right)}(๐,๐)I_{nm}^{\left(P\right)}\right\}$$
(44)
In Eq. (44) the sum over $`n,m=N,\mathrm{\Delta }`$ is assumed and the spin transition operators are $`๐บ_N=๐`$, $`๐บ_\mathrm{\Delta }=๐บ`$. The coupling constants $`\kappa _n`$ are defined in Eq. (41). The isospin factors are given by $`I_{nm}^{\left(P\right)}`$ for which we find
$$I_{NN}^{\left(P\right)}=32๐_1๐_2,I_{N\mathrm{\Delta }}^{\left(P\right)}=2+\frac{2}{3}๐_1๐_2,I_{\mathrm{\Delta }\mathrm{\Delta }}^{\left(P\right)}=\frac{4}{3}\frac{2}{9}๐_1๐_2$$
(45)
The function $`_{nm}^{\left(P\right)}`$ in Eq. (44) contains the integration over the time-like component of the four vector $`k`$.
$`_{nm}^{\left(P\right)}(๐,๐)`$ $`=`$ $`{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \frac{dk_0}{2\pi }}D_\pi \left(k+q\right)D_\pi \left(k\right)S_n\left(P_1k\right)S_m\left(P_2+k\right)`$ (46)
where $`D_\pi `$ and $`S_i`$ is the pion and nonrelativistic baryon propagators, respectively. $`S_i`$ is given by
$$iS_i\left(P\right)=\frac{M_i}{\sqrt{๐ท^2+M_i^2}}\frac{i}{P_0\sqrt{๐ท^2+M_i^2}+i0^+}$$
(47)
This integration can be carried out explicitely
$`_{nm}^{\left(P\right)}(\omega _1,\omega _2)`$ $`=`$ $`{\displaystyle \frac{i}{2}}[(\omega _1^2+\omega _2^2+3\omega _1\omega _2+[E_nE][E_mE])(E_n+E_m2E)`$ (48)
$`+(\omega _1+\omega _2)(2\omega _1\omega _2+[E_n+E_m2E]^2)]`$
$`\times {\displaystyle \frac{1}{\omega _1\omega _2\left(\omega _1+\omega _2\right)}}\left\{{\displaystyle \frac{1}{E_n+E_m2Ei0^+}}\right\}`$
$`\times {\displaystyle \frac{M_n}{E_n\left[\omega _1+E_nE\right]\left[\omega _2+E_nE\right]}}`$
$`\times {\displaystyle \frac{M_m}{E_m\left[\omega _1+E_mE\right]\left[\omega _2+E_mE\right]}}`$
Recall that $`q=(0,๐)`$, $`P_1=(E,๐ท)`$, $`P_2=(E,๐ท)`$ in $`NN`$ c.m. frame and in Eq. (48)
$$\omega _1=\sqrt{๐^2+M_\pi ^2},\omega _2=\sqrt{\left(๐+๐\right)^2+M_\pi ^2},E_n=\sqrt{\left(๐ท๐\right)^2+M_n^2},E_m=\sqrt{\left(๐ท๐\right)^2+M_m^2},$$
(49)
are the on-shell energies of intermediate pions and baryons with c.m. energy of the initial nucleons $`E=\sqrt{๐ท^2+m}`$.
In the nuclear matter $`\stackrel{~}{}`$, and the expression for $`\stackrel{~}{}`$ is obtained by using the dispersion representation for the in-medium pion propagator
$`\stackrel{~}{}_{nm}^{\left(P\right)}(๐,๐,\rho )`$ $`=`$ $`{\displaystyle \frac{1}{\pi ^2}}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}๐x^2{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}๐y^2_{nm}^{\left(P\right)}(x,y)\text{Im}\stackrel{~}{D}_\pi (x,๐,\rho )\text{Im}\stackrel{~}{D}_\pi (y,๐+๐,\rho )`$ (50)
where $`_{nm}^{\left(P\right)}(x,y)`$ is given by Eq. (48) with substitution $`\omega _1x`$ and $`\omega _2y`$.
For the crossed box diagrams, Fig. 13, we have
$$iV_{NN}^{\left(C\right)}=\underset{n,m}{}\kappa _n\kappa _m\frac{d^3๐}{\left(2\pi \right)^3}\left[๐บ_n^{\left(1\right)}\left(๐+๐\right)\right]\left[๐บ_n^{\left(1\right)}๐\right]^{}\left[๐บ_m^{\left(2\right)}๐\right]\left[๐บ_m^{\left(2\right)}\left(๐+๐\right)\right]^{}\left\{_{nm}^{\left(C\right)}(๐,๐)I_{nm}^{\left(C\right)}\right\}$$
(51)
The isospin factors are the following
$$I_{NN}^{\left(C\right)}=3+2๐_1๐_2,I_{N\mathrm{\Delta }}^{\left(C\right)}=2\frac{2}{3}๐_1๐_2,I_{\mathrm{\Delta }\mathrm{\Delta }}^{\left(C\right)}=\frac{4}{3}+\frac{2}{9}๐_1๐_2$$
(52)
and the vacuum expression for $`_{nm}^{\left(C\right)}(๐,๐)`$ is given by
$`_{nm}^{\left(C\right)}(๐,๐)`$ $`=`$ $`{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \frac{dk_0}{2\pi }}D_\pi \left(k+q\right)D_\pi \left(k\right)S_n\left(P_1k\right)S_m\left(P_2kq\right)`$ (53)
The integration can be done analytically and our result reads
$`_{nm}^{\left(C\right)}(\omega _1,\omega _2)`$ $`=`$ $`{\displaystyle \frac{i}{2}}\left[\omega _1^2+\omega _2^2+\omega _1\omega _2+\left(\omega _1+\omega _2\right)\left(\stackrel{~}{E}_n+E_m2E\right)+\left(\stackrel{~}{E}_nE\right)\left(E_mE\right)\right]`$ (54)
$`\times {\displaystyle \frac{1}{\omega _1\omega _2\left(\omega _1+\omega _2\right)}}`$
$`\times {\displaystyle \frac{M_n}{\stackrel{~}{E}_n\left[\omega _1+\stackrel{~}{E}_nE\right]\left[\omega _2+\stackrel{~}{E}_nE\right]}}`$
$`\times {\displaystyle \frac{M_m}{E_m\left[\omega _1+E_mE\right]\left[\omega _2+E_mE\right]}}`$
here $`\omega _1,\omega _2`$, $`E`$ and $`E_m`$ are defined in Eq. (49) and
$$\stackrel{~}{E}_n=\sqrt{\left(๐ท+๐+๐\right)^2+M_n^2}$$
(55)
The corresponding expression for the in-medium crossed box diagrams takes the form
$$\stackrel{~}{}_{nm}^{\left(C\right)}(๐,๐,\rho )=\frac{1}{\pi ^2}\underset{0}{\overset{\mathrm{}}{}}๐x^2\underset{0}{\overset{\mathrm{}}{}}๐y^2_{nm}^{\left(C\right)}(x,y)\text{Im}\stackrel{~}{D}_\pi (x,๐,\rho )\text{Im}\stackrel{~}{D}_\pi (y,๐+๐,\rho )$$
(56)
For the generic case of non-vanishing initial momenta the analytical structure of Eq. (48) is driven by the term in figure brackets
$$\frac{1}{E_n+E_m2Ei0^+}$$
(57)
Here we have to pay a special attention to the case where two intermediate $`NN`$ states $`n=m=N`$ appear, because in time ordered perturbation theory these diagrams are generated by iterations of the OPEP in a Lippmann-Schwinger equation (LSE). Considering the $`NN`$ intermediate state only, the iterated TPE can be easily identified and comes from the nucleon pole in Eq. (46) (in the lower half of the complex plane) corresponding to $`k_0=E_NEi0^+`$ with $`E=\sqrt{๐ท^2+M_N^2}`$
$$_{NN}^{\left(P\right)Npole}(\omega _1,\omega _2)=\frac{i}{2}\left(\frac{M_N}{E_N}\right)^2\frac{1}{EE_N+i0^+}\left[\frac{1}{\left(EE_N\right)^2\omega _1^2}\frac{1}{\left(EE_N\right)^2\omega _2^2}\right]$$
(58)
Expanding Eq. (58) in powers of $`1/M_N`$
$$_{NN}^{\left(P\right)Npole}(\omega _1,\omega _2)i\frac{1}{\omega _1^2\omega _2^2}\frac{M_N}{๐ท^2\left(๐ท๐\right)^2+i0^+}+๐ช\left(1/M_N\right)$$
(59)
and inserting the leading order result in Eq. (44) one can get the second order term in the non-relativistic Lippmann-Schwinger equation.
It is instructive to derive the contributions of the two remaining poles (in lower half-plane) from the pion propagators
$$_{NN}^{\left(P\right)\pi poles}(๐,๐)=\frac{i}{2\omega _1\omega _2}\frac{\omega _1^2+\omega _2^2+\omega _1\omega _2\left(EE_N\right)^2}{\left(\omega _1+\omega _2\right)\left(\omega _1^2\left(EE_N\right)^2\right)\left(\omega _2^2\left(EE_N\right)^2\right)}\frac{M_N^2}{E_N^2}$$
(60)
After the expansion of this result in powers of $`1/M_N`$ we get
$$_{NN}^{\left(P\right)\pi poles}(๐,๐)\frac{i}{2}\frac{\omega _1^2+\omega _2^2+\omega _1\omega _2}{\omega _1^3\omega _2^3\left(\omega _1+\omega _2\right)}+๐ช\left(1/M_N\right)$$
(61)
In this limit Eq. (61) cancels exactly the corresponding crossed box diagram with two intermediate nucleons in the isoscalar channel. Indeed, considering the crossed box diagram with two intermediate nucleons, the leading term of $`_{NN}^{\left(C\right)}(๐,๐)`$ in the $`1/M_N`$ expansion is given by
$$_{NN}^{\left(C\right)}(๐,๐)\frac{i}{2}\frac{\omega _1^2+\omega _2^2+\omega _1\omega _2}{\omega _1^3\omega _2^3\left(\omega _1+\omega _2\right)}+๐ช\left(1/M_N\right)$$
(62)
Note that for the isovector $`\pi \pi `$ exchange because of the different sign of the $`๐_1๐_2`$ term in $`I_{NN}^{\left(P\right)}`$, $`I_{NN}^{\left(C\right)}`$ in Eqs. (45) and (52) these two contributions would add. The cancellation discussed above hold also in the nuclear medium.
The result of the vacuum scalar-isoscalar $`NN`$ force generated by the planar and crossed box diagrams, with the nucleon pole diagrams excluded, is shown in the left panel of Fig. 14 by the solid curve. It is in agreement with Ref. where it was shown that the consistent use of the cut off regularization in both the correlated two pion exchange, and uncorrelated two pion exchange, together with the contribution of a repulsive $`\omega `$-exchange lead to a scalar-isoscalar potential in good agreement with the Argonne potential in the whole range of relevant distances. The corresponding results for the nuclear matter are shown in Fig. 14 (left) for three densities $`\rho _0`$, $`\rho _0/2`$ and $`\rho _0/4`$. Qualitatively the behavior is similar to the vacuum case but with a strong enhancement toward the small momentum transfer. As we have seen this feature is generic in present calculations. Uncorrelated $`\pi \pi `$ exchange becomes extremely attractive at intermediate distances.
The spin sum over intermediate baryon states of the $`๐บ๐`$ operator in Eqs. (44) and (51) gives in the scalar channel
$$\beta \left[\left(๐+๐\right)๐\right]^2$$
(63)
where $`\beta =1`$ $`\left(NN\right)`$, $`\beta =2/3`$ $`\left(N\mathrm{\Delta }\right)`$ and $`\beta =4/9`$ $`\left(\mathrm{\Delta }\mathrm{\Delta }\right)`$. Now we again wish to take into account the short range correlations
$`\left({\displaystyle \frac{D+F}{2f_\pi }}\right)^4F^2\left(๐\right)F^2\left(๐+๐\right)D_\pi \left(๐+๐\right)D_\pi \left(๐\right)\left[\left(๐+๐\right)๐\right]^2`$
$`\stackrel{~}{๐ฒ}_l\left(k+q\right)\stackrel{~}{๐ฒ}_l\left(k\right)\left[\left(\widehat{๐+๐}\right)\widehat{๐}\right]^2`$
$`+\left[\stackrel{~}{๐ฒ}_l\left(k+q\right)\stackrel{~}{๐ฒ}_t\left(k\right)+\stackrel{~}{๐ฒ}_l\left(k\right)\stackrel{~}{๐ฒ}_t\left(k+q\right)\right]\left\{1\left[\left(\widehat{๐+๐}\right)\widehat{๐}\right]^2\right\}`$
$`+\stackrel{~}{๐ฒ}_t\left(k+q\right)\stackrel{~}{๐ฒ}_t\left(k\right)\left\{1+\left[\left(\widehat{๐+๐}\right)\widehat{๐}\right]^2\right\}`$
where
$$\stackrel{~}{๐ฒ}_i\left(k\right)=\frac{๐ฑ_i\left(k\right)}{1๐ฐ\left(k\right)๐ฑ_i\left(k\right)}$$
(65)
where $`๐ฐ\left(k\right)`$ is the Lindhard function and $`๐ฑ_l`$ and $`๐ฑ_t`$ are given by Eq. (17) and (18), respectively.
## 7 Results and discussions
By comparing the results in Fig. 3 and Fig. 4 we observe that the effect of correlations has been essential and reduces drastically the medium effects found in Fig. 3 without corrections. We observe that the medium corrections weaken the strength of the central part of the OPEP. On the other hand the effect of the medium corrections on the tensor part are more moderate.
The $`\sigma `$ exchange presents similar features as we can see in Fig. 11. There we see that the medium effects in the absence of short range correlations are rather large and increase the strength of the interaction in about a factor two at $`\rho =\rho _0`$. However, as soon as the correlations are taken into account the medium modifications are reduced drastically and at $`\rho =\rho _0`$ just reduce the strength in about 25 percent.
The medium effects in the uncorrelated two pion exchange are also relatively small, with respect to the size of the interaction in the vacuum, see Fig. 14 (right). Once again the consideration of the correlations has been essential and reducing the size of the medium effects. However, the fact that the strength of this interaction is bigger than that of the correlated two pion exchange makes the absolute correction relevant and of the same size as the corrections discussed before.
Since we have taken also the effect of correlations in the free part of the interaction and any realistic calculation of the binding energy of matter will explicitly account for these correlations, the use of our modified in-medium potentials would lead to double-counting if any of these methods is used. For this purpose the relevant results from this work should be the differences between the in-medium potential and the free one. In this case we start already with one bubble and the correlations account for the repulsion between the external leg and the one in the bubble. The two external legs still have to be correlated and this will be done with the use, for instance, of the ordinary Bethe Golstone equation.
After this discussion of the medium effects in the different terms, we show the differences in momentum space between the medium and free parts of the different terms in Fig. 15 (left). The results are shown at $`\rho =\rho _0`$ and we see that these corrections are moderate, but they could have a relevant role in the binding of matter. In order to have a qualitative idea of the relevance of these corrections to the potentials we rewrite them in coordinate space and show these results in Fig. 15 (right). The effects seem sizeable at short distances, but this will be irrelevant in any realistic calculation of binding energies since the consideration of the short range correlations between the external legs will make this interaction inoperative. More interesting is the strength of the corrections around 1 fm, and there we see that all them are relatively small, of the order of 20 MeV or less, the biggest one being the central part of the OPE. However, given the size of the empirical scalar isoscalar attraction which is of the order of 20 MeV at intermediate densities, the corrections found here are not negligible. It would be thus interesting to see the effects of the results obtained here in observables like the binding of nuclear matter and other properties, which we hope to stimulate with the results obtained here.
## 8 Conclusions
In this paper we studied the modification of the one pion exchange, as well as two pion exchange potential inside a nuclear medium. For this purpose we separated the two pion exchange into an uncorrelated two pion exchange and the correlated two pion exchange. We study both in the scalar isoscalar channel, which is by large the most important one generated by this interaction. The correlated two pion exchange gave rise to the equivalent of the $`\sigma `$ exchange in other models and we studied the medium modification to it. On the other hand, for the uncorrelated two pion exchange we have followed a traditional approach in which only terms with at least one intermediate $`\mathrm{\Delta }`$ state are considered.
One of the important findings here was the effect played by the NN short range correlations which drastically moderated the medium corrections to the potential. In the absence of these, the corrections where unrealistically large. Yet, even if relatively moderate, the medium corrections found in this paper are sizeable enough to have relevant repercussions in the binding and other properties of nuclear matter and we would like to encourage calculations in this direction.
## Acknowledgments
We would like to acknowledge A. Ramos for a critical reading of the manuscript and useful suggestions. This work is partly supported by DGICYT contract number BFM2003-00856, and the E.U. EURIDICE network contract no. HPRN-CT-2002-00311. This research is part of the EU Integrated Infrastructure Initiative Hadron Physics Project under contract number RII3-CT-2004-506078.
## Appendix A Feynman rules for vertices
The set of Feynman diagrams shown in Fig. 16 appear in the construction of the $`\pi N\pi \pi N`$ transition amplitude. The corresponding $`3\pi NN`$ vertex functions are given by
$`iL_a`$ $`=`$ $`{\displaystyle \frac{D+F}{2}}{\displaystyle \frac{2\sqrt{2}}{12f_\pi ^3}}\left[๐\left(๐_1+2๐_1๐_2\right)\right]`$ (66)
$`iL_b`$ $`=`$ $`{\displaystyle \frac{D+F}{2}}{\displaystyle \frac{2\sqrt{2}}{12f_\pi ^3}}\left[๐\left(๐_2+2๐_2๐_1\right)\right]`$ (67)
$`iL_c`$ $`=`$ $`\pm {\displaystyle \frac{D+F}{2}}{\displaystyle \frac{2}{12f_\pi ^3}}\left[๐\left(2๐_1+๐_1+๐_2\right)\right]`$ (68)
$`iL_d`$ $`=`$ $`{\displaystyle \frac{D+F}{2}}{\displaystyle \frac{2\sqrt{2}}{12f_\pi ^3}}\left[๐\left(2๐_1+๐_1+๐_2\right)\right]`$ (69)
$`iL_e`$ $`=`$ $`{\displaystyle \frac{D+F}{2}}{\displaystyle \frac{2\sqrt{2}}{12f_\pi ^3}}\left[๐\left(2๐_2+๐_1+๐_2\right)\right]`$ (70)
$`iL_f`$ $`=`$ $`0`$ (71)
## Appendix B Correlated exchange in the medium
In this Appendix we demonstrate the cancellation of the off shell part of the correlated two pion exchange in the scalar isoscalar channel. For simplicity we consider the nucleon intermediate states only. The generalization to $`N\mathrm{\Delta }`$ and $`\mathrm{\Delta }\mathrm{\Delta }`$ is straightforward. The diagrams are shown in Fig. 10 and corresponding amplitudes are given by
$`T_a`$ $`=`$ $`T_a^{on}+T_a^{\left(1\right)}+T_a^{\left(2\right)}`$ (72)
$`iT_a^{on}`$ $`=`$ $`6\left({\displaystyle \frac{D+F}{2f_\pi }}\right)^4V_\pi ^{on}\left(t\right)\left[V_{N\pi \pi }\left(t\right)\right]`$ (73)
$`\times {\displaystyle }{\displaystyle \frac{d^4k}{\left(2\pi \right)^4}}\left[๐(๐+๐)\right]\left[๐๐\right]\left[๐^2๐ฐ\left(k\right)\right]S_F(P_1k)D^2\left(k\right)D(k+q)`$
$`iT_a^{\left(1\right)}`$ $`=`$ $`2\left({\displaystyle \frac{D+F}{2f_\pi }}\right)^4{\displaystyle \frac{1}{f_\pi ^2}}\left[V_{N\pi \pi }\left(t\right)\right]`$ (74)
$`\times {\displaystyle }{\displaystyle \frac{d^4k}{\left(2\pi \right)^4}}\left[๐(๐+๐)\right]\left[๐๐\right]\left[๐^2๐ฐ\left(k\right)\right]S_F(P_1k)D\left(k\right)D(k+q)`$
$`iT_a^{\left(2\right)}`$ $`=`$ $`2\left({\displaystyle \frac{D+F}{2f_\pi }}\right)^4{\displaystyle \frac{1}{f_\pi ^2}}\left[V_{N\pi \pi }\left(t\right)\right]`$ (75)
$`\times {\displaystyle }{\displaystyle \frac{d^4k}{\left(2\pi \right)^4}}\left[๐(๐+๐)\right]\left[๐๐\right]\left[๐^2๐ฐ\left(k\right)\right]S_F(P_1k)D^2\left(k\right)`$
$`T_b`$ $`=`$ $`T_b^{on}+T_b^{\left(1\right)}+T_b^{\left(2\right)}`$ (76)
$`iT_b^{on}`$ $`=`$ $`6\left({\displaystyle \frac{D+F}{2f_\pi }}\right)^4V_\pi ^{on}\left(t\right)\left[V_{N\pi \pi }\left(t\right)\right]`$ (77)
$`\times {\displaystyle }{\displaystyle \frac{d^4k}{\left(2\pi \right)^4}}\left[๐(๐+๐)\right]\left[๐๐\right]\left[(๐+๐)^2๐ฐ(k+q)\right]S_F(P_1k)D\left(k\right)D^2(k+q)`$
$`iT_b^{\left(1\right)}`$ $`=`$ $`2\left({\displaystyle \frac{D+F}{2f_\pi }}\right)^4{\displaystyle \frac{1}{f_\pi ^2}}\left[V_{N\pi \pi }\left(t\right)\right]`$ (78)
$`\times {\displaystyle }{\displaystyle \frac{d^4k}{\left(2\pi \right)^4}}\left[๐(๐+๐)\right]\left[๐๐\right]\left[(๐+๐)^2๐ฐ(k+q)\right]S_F(P_1k)D\left(k\right)D(k+q)`$
$`iT_b^{\left(2\right)}`$ $`=`$ $`2\left({\displaystyle \frac{D+F}{2f_\pi }}\right)^4{\displaystyle \frac{1}{f_\pi ^2}}\left[V_{N\pi \pi }\left(t\right)\right]`$
$`\times {\displaystyle }{\displaystyle \frac{d^4k}{\left(2\pi \right)^4}}\left[๐(๐+๐)\right]\left[๐๐\right]\left[(๐+๐)^2๐ฐ(k+q)\right]S_F(P_1k)D^2(k+q)`$
$`iT_c`$ $`=`$ $`iT_a\left(kl,qq,P_1P_2\right)`$
$`iT_d`$ $`=`$ $`iT_b\left(kl,qq,P_1P_2\right)`$ (79)
where $`V_{N\pi \pi }\left(t\right)`$ is the triangle loop integral
$$V_{N\pi \pi }\left(t\right)=i\left(\frac{D+F}{2f_\pi }\right)^2\frac{d^4l}{\left(2\pi \right)^4}\left[๐\left(๐๐\right)\right]\left[๐๐\right]S_F\left(P_2l\right)D\left(lq\right)D\left(l\right)$$
(80)
and the on-mass-shell $`\pi \pi `$ scattering amplitude is
$$V_\pi ^{on}\left(t\right)=\frac{1}{f_\pi ^2}\left(t\frac{M_\pi ^2}{2}\right)$$
(81)
$`iT_e`$ $`=`$ $`\left({\displaystyle \frac{D+F}{2f_\pi }}\right)^4{\displaystyle \frac{1}{f_\pi ^2}}\left[V_{N\pi \pi }\left(t\right)\right]`$ (82)
$`\times {\displaystyle }{\displaystyle \frac{d^4k}{\left(2\pi \right)^4}}\left[๐(๐+๐)\right]\left[๐๐\right]\left[(3๐๐+2๐^2)๐ฐ\left(k\right)\right]S_F(P_1k)D\left(k\right)D(k+q)`$
$`iT_f`$ $`=`$ $`\left({\displaystyle \frac{D+F}{2f_\pi }}\right)^4{\displaystyle \frac{1}{f_\pi ^2}}\left[V_{N\pi \pi }\left(t\right)\right]`$
$`\times {\displaystyle }{\displaystyle \frac{d^4k}{\left(2\pi \right)^4}}\left[๐(๐+๐)\right]\left[๐๐\right]\left[(๐^22๐^2๐๐)๐ฐ(k+q)\right]S_F(P_1k)D\left(k\right)D(k+q)`$
$`iT_g`$ $`=`$ $`iT_e\left(kl,qq,P_1P_2\right)`$ (84)
$`iT_h`$ $`=`$ $`iT_f\left(kl,qq,P_1P_2\right)`$ (85)
$`iT_i`$ $`=`$ $`\left({\displaystyle \frac{D+F}{2f_\pi }}\right)^4{\displaystyle \frac{1}{f_\pi ^2}}\left[V_{N\pi \pi }\left(t\right)\right]{\displaystyle \frac{d^4k}{\left(2\pi \right)^4}\left[๐\left(๐2๐\right)\right]\left[๐๐\right]\left[๐^2๐ฐ\left(k\right)\right]S_F\left(P_1k\right)D^2\left(k\right)}`$ (86)
$`iT_j`$ $`=`$ $`\left({\displaystyle \frac{D+F}{2f_\pi }}\right)^4{\displaystyle \frac{1}{f_\pi ^2}}\left[V_{N\pi \pi }\left(t\right)\right]`$ (87)
$`\times {\displaystyle }{\displaystyle \frac{d^4k}{\left(2\pi \right)^4}}\left[๐๐\right]\left[๐(๐+2๐)\right]\left[๐^2๐ฐ\left(k\right)\right]S_F(P_1+qk)D^2\left(k\right)`$
$`iT_k`$ $`=`$ $`iT_i\left(kl,qq,P_1P_2\right)`$ (88)
$`iT_l`$ $`=`$ $`iT_j\left(kl,qq,P_1P_2\right)`$ (89)
In the present case $`q=(0,๐)`$ and for $`๐0`$ the amplitude $`T_a^{\left(2\right)}`$ cancels exactly $`T_i`$. The same cancelation is found for $`T_b^{\left(2\right)}`$ and $`T_j`$. The sum of $`T_e`$ and $`T_a^{\left(1\right)}`$ is given by
$`iT_a^{\left(1\right)}iT_e`$ $`=`$ $`3\left({\displaystyle \frac{D+F}{2f_\pi }}\right)^4{\displaystyle \frac{1}{f_\pi ^2}}\left[V_{N\pi \pi }\left(t\right)\right]`$ (90)
$`\times {\displaystyle }{\displaystyle \frac{d^4k}{\left(2\pi \right)^4}}\left[๐(๐+๐)\right]\left[๐๐\right]\left[๐๐๐ฐ\left(k\right)\right]S_F(P_1k)D\left(k\right)D(k+q)`$
The sum of $`T_f`$ and $`T_b^{\left(1\right)}`$ takes the form
$`iT_b^{\left(1\right)}iT_f`$ $`=`$ $`3\left({\displaystyle \frac{D+F}{2f_\pi }}\right)^4{\displaystyle \frac{1}{f_\pi ^2}}\left[V_{N\pi \pi }\left(t\right)\right]`$
$`\times {\displaystyle }{\displaystyle \frac{d^4k}{\left(2\pi \right)^4}}\left[๐(๐+๐)\right]\left[๐๐\right]\left[(๐+๐)๐๐ฐ(k+q)\right]S_F(P_1k)D\left(k\right)D(k+q)`$
From this the sum of four diagrams is given by
$`iT_a^{\left(1\right)}iT_eiT_b^{\left(1\right)}iT_f=3\left({\displaystyle \frac{D+F}{2f_\pi }}\right)^4{\displaystyle \frac{1}{f_\pi ^2}}\left[V_{N\pi \pi }\left(t\right)\right]`$ (92)
$`\times {\displaystyle }{\displaystyle \frac{d^4k}{\left(2\pi \right)^4}}\left[๐(๐+๐)\right]\left[๐๐\right][๐^2๐ฐ(k+q)+๐๐\{๐ฐ(k+q)๐ฐ\left(k\right)\}]S_F(P_1k)D\left(k\right)D(k+q)`$
Alternatively, the spin flip parts of Eqs. (90) and (B) can be canceled if we change in Eq. (B) $`๐\left(๐+๐\right)`$ and using again the fact that $`q=(0,๐)`$ and the properties of the Lindhard function $`๐ฐ(p_0,๐)=๐ฐ(p_0,๐)`$ we get
$`iT_b^{\left(1\right)}iT_f`$ $`=`$ $`3\left({\displaystyle \frac{D+F}{2f_\pi }}\right)^4{\displaystyle \frac{1}{f_\pi ^2}}\left[V_{N\pi \pi }\left(t\right)\right]`$
$`\times {\displaystyle }{\displaystyle \frac{d^4k}{\left(2\pi \right)^4}}\left[๐๐\right]\left[๐(๐+๐)\right]\left[๐๐๐ฐ\left(k\right)\right]S_F(P_1k)D\left(k\right)D(k+q)`$
and
$`iT_a^{\left(1\right)}iT_eiT_b^{\left(1\right)}iT_f`$ $`=`$ $`6\left({\displaystyle \frac{D+F}{2f_\pi }}\right)^4{\displaystyle \frac{1}{f_\pi ^2}}\left[V_{N\pi \pi }\left(t\right)\right]`$
$`\times {\displaystyle }{\displaystyle \frac{d^4k}{\left(2\pi \right)^4}}\left[(๐^2+๐๐)๐๐\right]๐ฐ\left(k\right)S_F(P_1k)D\left(k\right)D(k+q)`$
Finally
$`iT_c^{\left(1\right)}iT_giT_d^{\left(1\right)}iT_h`$ $`=`$ $`6\left({\displaystyle \frac{D+F}{2f_\pi }}\right)^4{\displaystyle \frac{1}{f_\pi ^2}}\left[V_{N\pi \pi }\left(t\right)\right]`$
$`\times {\displaystyle }{\displaystyle \frac{d^4l}{\left(2\pi \right)^4}}\left[(๐๐๐^2)๐๐\right]๐ฐ\left(l\right)S_F(P_2l)D\left(l\right)D(lq)`$
As one can see in the limit $`๐0`$ we get an exact cancellation of the off mass shell part of the $`\pi \pi `$ amplitude and only $`T_a^{on}`$ and $`T_b^{on}`$ are left.
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# First Results on {_ฮ}ยนยฒC production at DAฮฆNE
## 1 Introduction
Even though the first hypernucleus was identified more than fifty years ago , Hypernuclear Physics was systematically studied only in the last decade, in spite of its great interest and discovery potential for nuclear structure, strong and weak interactions and possible quark effects in nuclei. The most recent experiments were performed at AGS (Brookhaven) and at the 12 GeV PS (KEK) , and hypernuclei production was based on the strangeness exchange $`(K^{},\pi ^{})`$ reaction on nuclear targets with $`K^{}`$ in flight and at rest, or on the associated production $`(\pi ^+,K^+)`$ one.
This experimental scenario led to the idea of performing hypernuclear physics experiments with a dedicated detector (FINUDA) using a source of $`K^{}`$ different from traditional hadron facilities; that is, the $`\varphi `$-factory DA$`\mathrm{\Phi }`$NE at the Frascati National Laboratories of I.N.F.N., Italy .
FINUDA (acronym for โFIsica NUcleare a DA$`\mathrm{\Phi }`$NEโ) can be considered an experiment of third generation in hypernuclear physics. The original design of the FINUDA apparatus and, in particular, the large angle covered for the detection of charged and neutral particles emitted after the formation and decay of hypernuclei, allows for the simultaneous measurement of observables like excitation energy spectra, lifetimes and partial decay widths for mesonic and non-mesonic decay, with high statistics and good energy resolution (better than 1 MeV). Furthermore, these observables can be measured for different targets at the same time, thus reducing systematic errors in comparing properties of different hypernuclei.
The first FINUDA data taking at DA$`\mathrm{\Phi }`$NE started in December 2003 and was successfully concluded in March 2004. In the following the first results from the experiment will be reported, which fully confirm the expected capability of FINUDA to perform high quality hypernuclear physics at the DA$`\mathrm{\Phi }`$NE collider.
## 2 The FINUDA Experiment at DA$`\mathrm{\Phi }`$NE
DA$`\mathrm{\Phi }`$NE (Double Annular $`\mathrm{\Phi }`$-factory for Nice Experiments) consists of two rings, one for electrons and the other for positrons, that overlap in two straight sections where the beams collide. The energy of each beam is 510 MeV in order to produce the $`\varphi `$(1020) meson in the collisions.
At the luminosity $``$=10<sup>32</sup>cm<sup>-2</sup>s<sup>-1</sup>, the $`\varphi `$ meson is produced at a rate $`4.4\times 10^2`$ s<sup>-1</sup>. The $`\varphi K^+K^{}`$ branching ratio is $`49\%`$ and therefore, since the $`\varphi `$ is produced almost at rest, DA$`\mathrm{\Phi }`$NE is a source of $`2.2\times 10^2`$ ($`K^+K^{}`$) pairs/s, collinear, background free and of very low energy ($`16`$ MeV). The low energy of the kaons is the key-feature for performing hypernuclear physics experiments at the DA$`\mathrm{\Phi }`$NE $`\varphi `$-factory.
The main idea of FINUDA is to slow down to rest the negative kaons from the $`\varphi K^+K^{}`$ decay in thin solid targets, so as to study the following formation and decay of hypernuclei produced by the strangeness exchange reaction:
$$K_{stop}^{}+{}_{}{}^{A}Z{}_{\mathrm{\Lambda }}{}^{A}Z+\pi ^{}$$
(1)
where $`{}_{}{}^{A}Z`$ indicates a target nucleus and $`{}_{\mathrm{\Lambda }}{}^{A}Z`$ the produced hypernucleus. The method of producing hypernuclei via reaction (1) was the standard one with emulsions or bubble chambers in the sixties. A first attempt to use reaction (1) even in a counter experiment was done in 1973; a <sup>12</sup>C target was employed and the overall energy resolution was $``$ 6 MeV FWHM . A substantial experimental effort with a dedicated apparatus and on several targets was then performed at KEK in the late eighties ; however, the instrumental resolution did not exceed 2.4 MeV FWHM.
The use of $`K^{}`$โs from a $`\varphi `$-factory to produce hypernuclei has several advantages when compared to the extracted $`K^{}`$ beams or intense $`\pi ^+`$ beams . First of all, the low-energy and almost monochromatic $`K^{}`$ emitted from $`\varphi `$ decay can be efficiently stopped in thin targets (0.2 g/cm<sup>2</sup>). At hadron machines, extracted $`K^{}`$ beams require thick targets (some g/cm<sup>2</sup>), in order to obtain sufficient event rates. In addition, the uncertainty on the interaction point and the energy straggling of the emitted particles impair the achievable resolution. The same problem occurs in hypernuclear spectroscopy performed via the more efficient ($`\pi ^+,K^+`$) reaction . In FINUDA, furthermore, the use of thin targets along with the low-mass of the spectrometer tracking system permits the detection of charged particles other than pions (mainly $`p`$โs and $`d`$โs) with a solid angle similar to that of pions, and a threshold as low as $`100`$ MeV/$`c`$ for protons and $`200`$ MeV/$`c`$ for deuterons. Finally, the cylindrical symmetry of the interaction region allowed for the construction of a spectrometer of cylindrical shape with a large solid angle which, for the detection of the $`\pi ^{}`$โs coming from reaction (1), is larger than $`\pi `$ sr, therefore much bigger than those available at fixed target machines, typically $``$ 100 msr. Such an acceptance, along with the excellent performances of DA$`\mathrm{\Phi }`$NE, enables the detection of hypernuclei with a rate of about 80 hypernuclei/hour at $``$=10<sup>32</sup> cm<sup>-2</sup>s<sup>-1</sup> (with a 10<sup>-3</sup> capture rate).
Fig. 1 shows a global view of the apparatus. The layers of the tracker are contained inside a superconducting solenoid, which provides a highly homogeneous (within 2% inside the tracking volume) magnetic field of 1.0 T over a cylindrical volume of 146 cm radius and 211 cm length.
Three main regions can be distinguished inside the FINUDA apparatus.
* The interaction/target region is shown schematically in fig. 2a). Here, the highly ionizing $`(K^+,K^{})`$ pairs are detected by a barrel of 12 thin scintillator slabs (dubbed TOFINO for short), surrounding the beam pipe, with a time resolution of $`\sigma 250`$ ps. The TOFINO barrel is surrounded by an octagonal array of silicon microstrips (ISIM) featuring a spatial resolution $`\sigma 30`$ $`\mu `$m and an energy resolution on $`\mathrm{\Delta }E/\mathrm{\Delta }x`$ for the kaons from $`\varphi `$ decay of 20% . Thin solid target modules are positioned at a distance of a few millimeters on the external side of each element of the octagon. The task of the ISIM detector is the determination of the interaction points of the $`(K^+,K^{})`$ pairs in the thin targets.
* The external tracking device consists of four different layers of position sensitive detectors. It is arranged in cylindrical symmetry and is immersed in a He atmosphere to reduce the effects of the multiple Coulomb scattering. The trajectories of charged particles coming from the targets and crossing the tracking system are measured by: (i) a first array of ten double-sided silicon microstrip modules (OSIM) placed close to the target elements (see fig. 2a); (ii) two arrays of eight planar low-mass drift chambers (LMDC) filled with a (70%He-30%C<sub>4</sub>H<sub>10</sub>) mixture, featuring a spatial resolution $`\sigma _{\rho \varphi }150`$ $`\mu `$m and $`\sigma _z1.0`$ cm ; (iii) a straw tube detector, composed by six layers of longitudinal and stereo tubes, which provide a spatial resolution $`\sigma _{\rho \varphi }150`$ $`\mu `$m and $`\sigma _z500`$ $`\mu `$m . The straw tubes are positioned at 1.1 m from the beams interaction point. With the magnetic field set at 1.0 T, the design momentum resolution of the spectrometer, for 270 MeV/c $`\pi ^{}`$โs, is $`\mathrm{\Delta }p/p`$=0.4% FWHM. It corresponds to an energy resolution on hypernuclear spectra better than 1.0 MeV. On the other hand, the energy resolution, for the 80 MeV protons emitted in the hypernuclear non-mesonic decay, is 1.6 MeV FWHM.
* The external time of flight barrel (TOFONE) is composed by 72 scintillator slabs, 10 cm thick and 255 cm long, and provides signals for the first level trigger and for the measurement of the time-of-flight of the charged particles, with a time resolution $`\sigma `$ 350 ps. Moreover, it allows for the detection of neutrons following hypernucleus decays with an efficiency of $``$10%, an angular acceptance of 70% and an energy resolution of 8 MeV FWHM for neutrons of 80 MeV .
Further details concerning the design and performances of the FINUDA apparatus can be found in Refs. .
An important feature of the FINUDA apparatus (see fig. 2a) is the possibility to host eight different targets close to the interaction region; therefore, the possibility of obtaining data on different hypernuclei at the same time. For the starting run, the following targets were selected : two <sup>6</sup>Li (isotopically enriched to 90%), one <sup>7</sup>Li (natural isotopic abundance), three <sup>12</sup>C, one <sup>27</sup>Al and one <sup>51</sup>V. Physical motivations for the performed choice are described in Refs. .
## 3 Data taking and apparatus performances
Many experimental tests were performed during the data taking in order to monitor the machine performance as well as the calibrations of the spectrometer. Hereafter the most relevant ones are listed:
1. the luminosity of the DA$`\mathrm{\Phi }`$NE collider was continuously evaluated by means of the Bhabha scattering events, and was in agreement (within $`10\%`$) with the values provided by the machine. The top luminosity reached during the run was 0.7$`\times 10^{32}`$ cm<sup>-2</sup>s<sup>-1</sup>, with a daily integrated luminosity of about 4 pb<sup>-1</sup>;
2. the profile of the interaction region was also continuously monitored by FINUDA, and used to control the collider;
3. the energy of the colliding beams was measured on-line via Bhabha scattering and through the reconstruction of the $`K_S\pi ^+\pi ^{}`$ invariant mass, where the $`K_S`$โs are due to the $`\varphi K_SK_L`$ decay.
The trigger selecting hypernucleus formation events requires two fired back-to-back TOFINO slabs, with signal amplitude above an energy threshold accounting for the high ionization of slow kaons, and a fast coincidence on the TOFONE barrel . This allows $`(K^+,K^{})`$ pairs, together with a fast particle crossing the spectrometer and hitting the external scintillator barrel, to be selected against the physical background coming from the other $`\varphi `$ decays or against fake events generated by the accelerator electromagnetic background.
The reconstruction procedure of the $`\varphi `$ formation point and of the kaon directions and momenta at vertex uses the kaon interaction points in the ISIM modules, identified through their high stopping power. The procedure is based on a two helix algorithm which accounts for the kinematics of the $`\varphi `$ decay, the average value of the $`\varphi `$ mass, the crossing angle (12.5 mrad) of the $`e^+e^{}`$ beams, measured by using Bhabha events, and the geometry of the vertex region. The stopping points of the kaons in the targets are computed by a tracking procedure based on the GEANE package , which performs a numerical integration of the trajectory starting from the $`\varphi `$ formation point and the kaon direction and momenta and accounting for the geometrical structure and the material composition of the FINUDA interaction region.
The beam crossing angle determines a small total momentum of the $`\varphi `$ (boost: 12.3 MeV/$`c`$) directed towards the positive $`x`$ side. This boost adds to the 127 MeV/$`c`$ average momentum of the kaons from the $`\varphi `$ decay introducing a left-right asymmetry clearly visible in Fig. 2b), which shows the scatter-plot of the reconstructed $`y`$ vs $`x`$ coordinates of the $`K^{}`$ stopping points. The distribution of points on the outer octagon represents the positions of the eight targets, where most of the $`K^{}`$ stop ($`75\%`$ of all $`K^{}`$ interactions in the apparatus). A partial accumulation of points also occurs on the left-side ISIM modules (10%). The events corresponding to the $`K^{}`$โs stopping in the ISIM modules provide an additional sample in a supplementary silicon target. The remaining density of points partially depicts TOFINO.
Hypernuclear events are selected by the simultaneous presence of $`K^+`$ and $`K^{}`$ particles. The $`K^+`$โs enable the $`K^{}`$ tagging and moreover offer the possibility to perform an accurate and continuous in-beam calibration of FINUDA. The positive kaons, stopping in the target array, decay at rest with a mean life of 12.4 ns. The two main two-body decays $`K^+\mu ^+\nu _\mu `$ (B.R.=63.51%) and $`K^+\pi ^+\pi ^0`$ (21.16%) are a source of monochromatic particles fully crossing the spectrometer, with momenta 235.5 MeV/c for the $`\mu ^+`$ and and 205.1 MeV/c for the $`\pi ^+`$, respectively. The absolute scale of the momenta was determined with a precision better than 200 keV/$`c`$, even in the simplified hypothesis (applied in the analysis presented here) of a constant magnetic field of 1.0 T, directed along the $`z`$ axis, in the whole tracking volume. This precision can be assumed as the systematic error on the measurement of the particlesโ momenta in the range between 200 and 300 MeV/$`c`$.
For the present analysis only high quality tracks were selected. Such tracks are emitted in the forward hemisphere, with respect to the direction of the $`K^+`$, and cross a minimum amount of materials inside the spectrometer. Fig. 3 shows the momentum distribution of the positive tracks coming from the stopped $`K^+`$. The two peaks at 236 MeV/c and 205 MeV/c correspond to the previously mentioned decays. The tails on the left of the two peaks are due to different contributions, the biggest part played by instrumental effects due to the momentum loss of particle crossing the edges of the chambers and their supports. Moreover, in this region two additional $`K^+`$ decay channels open: the $`K_{e3}^+`$ mode (B.R.=4.8%), giving a continuum spectrum of positrons (which cannot be distinguished from $`\mu ^+`$โs) ending at 228 MeV/$`c`$, and the $`K_{\mu 3}^+`$ one (3.2%), which gives again a continuum spectrum with end point at 215 MeV/$`c`$. By analyzing these different contributions to the peaks shape one can conclude that the asymmetry affects, overall, the gaussian line shape at the level of about 4%. This peak asymmetry was however not considered in the fit of the spectra described in Sec. 4, since other error sources were overwhelming.
From the width of the $`\mu ^+`$ peak the present momentum resolution of the apparatus can be estimated to be $`\mathrm{\Delta }p/p`$=0.6% FWHM, which corresponds to 1.29 MeV FWHM for the hypernuclear levels in agreement with the results of the hypernuclear spectra reported in the next section. We expect that the momentum resolution of the spectrometer should improve to the design value of 0.4% FWHM once the final detector calibration and alignment will be performed, and the mapped magnetic field will be inserted in the reconstruction and fitting procedure.
## 4 Results on $`{}_{\mathrm{\Lambda }}{}^{}{}_{}{}^{12}`$C spectroscopy. Discussion and conclusions.
In order to evaluate the capabilities of FINUDA to yield relevant spectroscopic parameters, the analysis started from <sup>12</sup>C targets. We recall that for $`{}_{\mathrm{\Lambda }}{}^{}{}_{}{}^{12}`$C an excitation spectrum with a 1.45 MeV FWHM resolution was recently obtained at KEK using the $`(\pi ^+,K^+)`$ reaction at 1.05 GeV/$`c`$ by the E369 Collaboration .
The spectra out of only two of the three available <sup>12</sup>C targets were added since the third one showed a slight systematic energy displacement, of about 0.5 MeV. The reason of this is under study, and therefore for the current analysis these data are not included. The requirement of high quality tracks (long tracks crossing the whole spectrometer, with a hit on each tracking detector, i.e OSIM, LMDCโs and straw tubes) reduced the analysed data to about the 40% of the whole available sample of events with vertex coming from a <sup>12</sup>C target.
The raw momentum spectrum of the $`\pi ^{}`$ coming from the analysed <sup>12</sup>C targets is shown in Fig. 4. Different processes produce $`\pi ^{}`$ after $`K^{}`$ absorption and reproduce well the experimental spectra :
a) quasi-free $`\mathrm{\Sigma }^+,\mathrm{\Sigma }^0`$ and $`\mathrm{\Lambda }`$ production: $`K^{}p\mathrm{\Sigma }^+\pi ^{}`$, $`K^{}n\mathrm{\Sigma }^0\pi ^{}`$, $`K^{}n\mathrm{\Lambda }\pi ^{}`$;
b) quasi-free $`\mathrm{\Lambda }`$ decay: $`\mathrm{\Lambda }p\pi ^{}`$;
c) quasi-free $`\mathrm{\Sigma }^{}`$ production: $`K^{}p\mathrm{\Sigma }^{}\pi ^+`$, followed by $`\mathrm{\Sigma }^{}n\pi ^{}`$;
d) two nucleon $`K^{}`$ absorption: $`K^{}(NN)\mathrm{\Sigma }^{}N`$, followed by $`\mathrm{\Sigma }^{}n\pi ^{}`$.
All the mentioned reactions were simulated in detail in the FINUDA Monte Carlo program. The simulated events were reconstructed by the same program used for the real events, with the same selection criteria, in order to accurately take into account the geometrical acceptance and the reconstruction efficiency of the apparatus. In particular, the size of the spectrometer and the value of the magnetic field determine an acceptance momentum cut of about 180 MeV/c for four-hits tracks, which excludes most of the reactions producing low energy $`\pi ^{}`$โs. However, in the momentum region where the bound states of <sup>12</sup>C are expected (beyond $`260`$ MeV/$`c`$), only process d) is contributing. We remark that both processes c) and d) are due to $`\mathrm{\Sigma }^{}`$ decay in flight, but the $`\pi ^{}`$ distribution from the process c) is peaked at 190 MeV/$`c`$, and goes to zero beyond 260 MeV/$`c`$. The dashed line in Fig. 4 represents the contribution due to process d), normalized to the number of entries in the $`(275รท320)`$ MeV/$`c`$ momentum region, beyond the physical region for the production of $`\mathrm{\Lambda }`$-hypernuclei via reaction (1).
In order to obtain the $`\mathrm{\Lambda }`$ binding energy distribution the d) process is subtracted from the $`\pi ^{}`$ momentum distribution, and the momenta are converted into binding energies ($`B_\mathrm{\Lambda }`$). The two prominent peaks, as can be seen in Figs. 5a) and b), at $`B_\mathrm{\Lambda }`$ around 11 MeV (ground state) and 0 MeV, were already observed in previous experiments and interpreted as $`(\nu p_{\frac{1}{2}}^1,\mathrm{\Lambda }s)`$ and $`(\nu p_{\frac{3}{2}}^1,\mathrm{\Lambda }p)`$ ($`\nu `$= nucleon). The experimental energy resolution was determined by fitting the $`B_\mathrm{\Lambda }`$ 11 MeV peak with a gaussian curve ($`\chi ^2/d.o.f.=1.71`$), and amounts to 1.29 MeV FWHM. The ground state of $`{}_{\mathrm{\Lambda }}{}^{}{}_{}{}^{12}`$C is assumed to be a single state. Indeed, it is known that it consists of a $`(1^{},\mathrm{\hspace{0.33em}2}^+)`$ doublet, but theoretical calculations predict splittings of 70 keV , 80 keV and 140 keV between them, one order of magnitude smaller than the present instrumental resolution. The peak at about 0. MeV has a more complicated structure, and we tried to disentagle different contributions in the analysis described in the following.
The experimental spectrum closely resembles the one from E369 experiment . This is expected, as the production of hypernuclear states is, in first approximation, determined by the momentum transferred to $`\mathrm{\Lambda }`$โs, which is grossly comparable for both experiments ($`250`$ MeV/$`c`$ for FINUDA, $`350`$ MeV/$`c`$ for E369). The $`100`$ MeV/$`c`$ difference may account for the different yield of the two main peaks.
The absolute values of the capture rates for the different peaks could be obtained in a simple way by the method of the $`K^{}`$ tagging. Indeed, in the events where the $`K^+`$ is seen to decay in the $`K_{\mu 2}`$ and $`K_{\pi 2}`$ decay modes with the produced $`\mu ^+`$ or $`\pi ^+`$ crossing the spectrometer and hitting the TOFONE barrel, we are sure the trigger condition on the prompt TOFONE coincidence has been satisfied by the charged products of the $`K^+`$ decay.
Hence, in these events triggered by the decay products of the $`K^+`$, the interactions of the corresponding $`K^{}`$ in the targets are observed without any trigger bias. Using this subsample of events, the number of $`K^{}`$ stopping in the targets can be counted directly and the number of $`\pi ^{}`$ produced by the $`K^{}`$ interactions can be accurately determined by only correcting for the apparatus acceptance for $`\pi ^{}`$ of selected momentum and for detector efficiency. The acceptance is calculated using the FINUDA Monte Carlo and the detector efficiency is determined by calibration data.
The value obtained for the $`{}_{\mathrm{\Lambda }}{}^{}{}_{}{}^{12}`$C ground state formation is $`(1.01\pm 0.11_{stat}\pm 0.10_{sys})\times 10^3`$/(stopped $`K^{}`$). It agrees very well with the value $`(0.98\pm 0.12)\times 10^3`$/(stopped $`K^{})`$ measured at KEK ; we recall that the first generation CERN experiment reported the value $`(2\pm 1)\times 10^4`$/(stopped $`K^{}`$) .
In between the two main peaks, there are also indications of other states produced with weaker strength. In order to reproduce, at least qualitatively, this spectrum six gaussian functions were used, centered at the $`B_\mathrm{\Lambda }`$ values reported in Ref. ; the widths were fixed, for all of them, to $`\sigma =0.55`$ MeV, corresponding to the experimental resolution. The abscissa scale is affected only by a scale error of $`\pm 80`$ keV. The result of this fit is shown in Fig. 5a).
The spectrum is not well reproduced, the resulting reduced $`\chi ^2/d.o.f.`$ is 3.8 (for 64 d.o.f.), and in particular the region $`10\mathrm{MeV}<B_\mathrm{\Lambda }<5\mathrm{MeV}`$ is poorly fitted. The capture rates for these different contributions normalized to the ground state capture rate are reported in the second column of the upper part of Table 1. A better $`\chi ^2/d.o.f.`$= 2.3 is obtained by adding a further contribution, and leaving the positions of the seven levels free (57 d.o.f.). Their values are reported in the second column of the lower part of Tab. 1. The capture rates for these different contributions are again normalized to the capture rate for the $`{}_{\mathrm{\Lambda }}{}^{}{}_{}{}^{12}`$C ground state formation. The result of the fit is shown in Fig. 5b). A contribution from the quasi-free $`\mathrm{\Lambda }`$-production, starting from $`B_\mathrm{\Lambda }=0`$ and properly smeared by taking into account the instrumental resolution, was included in both fits.
The peaks $`\mathrm{\#}2`$ and $`\mathrm{\#}3`$ can be attributed to the <sup>11</sup>C core excited states at 2.00 and 4.80 MeV. The excitation of these states was expected in several theoretical calculations ; their energies may be sensitive to the $`\mathrm{\Lambda }`$-$`N`$ interaction matrix elements. However, the peak $`\mathrm{\#}`$4 and a newly observed peak $`\mathrm{\#}`$5 are not explained with such a simple way. Excluding from the fit the peak $`\mathrm{\#}`$5 the value for $`\chi ^2/d.o.f.`$ worsened to 3.3. There exist several positive-parity excited states of the <sup>11</sup>C core in this energy region which could contribute the these structures . It can be noticed that the integrated strength for the excitation of all these weakly excited states compared to that of the two main peaks is more than twice larger than the one reported by E369.
The sum of the capture rates for the $`B_\mathrm{\Lambda }=0.27`$ MeV and $`B_\mathrm{\Lambda }=2.1`$ MeV states is $`(2.59\pm 0.19_{stat})\times 10^3`$/(stopped $`K^{})`$, and agrees with the KEK result $`(2.3\pm 0.3)\times 10^3`$/(stopped $`K^{})`$ , in which the contributions for the two states were not resolved. The CERN experiment reports $`(3\pm 1)\times 10^4`$/(stopped $`K^{})`$. In the present analysis these states are indeed resolved, though, inevitably, strongly correlated in our fit. It is however remarkable that the relative intensities for the contributions at $`B_\mathrm{\Lambda }=0.27`$ MeV and $`B_\mathrm{\Lambda }=2.1`$ MeV are close to the values found by Dalitz et al. in an emulsion experiment.
Theoretical calculations for the ground state formation quote the values $`0.33\times 10^3`$/(stopped $`K^{}`$) , $`0.23\times 10^3`$/(stopped $`K^{}`$) and $`0.12\times 10^3`$/(stopped $`K^{}`$) . Analogous calculations for the capture rate leading to states in which the $`\mathrm{\Lambda }`$ is in a $`p`$ state quote , respectively, $`0.96\times 10^3`$/(stopped $`K^{}`$) according to the theoretical prediction of Ref. , and $`0.59\times 10^3`$/(stopped $`K^{}`$) following Ref. . As general remark it may be noticed that our measured values are larger by factors $`(3รท6)`$ as compared with theoretical predictions. Finally, the pattern of the relative strength for the excited-core states is also significantly larger than the theoretical calculation reported in Ref. .
In conclusion, the method of producing hypernuclei stopping in thin nuclear targets the low energy $`K^{}`$ from $`\varphi `$ decay at a $`\varphi `$-factory was proved to work, and may be used to perform accurate measurements on many hypernuclei observables. A first analysis allowed to achieve this purpose and to obtain also some interesting new physical information.
## 5 Acknowledgements
We are greatly indebted to Dr. S. Bertolucci, former Director of LNF, for his continuous encouragement and help. Dr. P. Raimondi and the DA$`\mathrm{\Phi }`$NE machine staff are warmly acknowledged for their very skillful handling of the DA$`\mathrm{\Phi }`$NE collider.
The excellent and qualified contribution of the whole FINUDA technical staff, at all stages of the experiment setting-up, is deeply appreciated.
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# Mass loss at the lowest stellar massesBased on observations obtained at the European Southern Observatory using the Very Large Telescope in Cerro Paranal, Chile (observing runs 67.C-0549(B), 69.B-0126(A), 71.C-0429(C) and 71.C-0429(D)).
## 1 Introduction
In the last decade evidence has been found that the mass accretion model that simultaneously explains the mass accretion and mass loss features observed in a great variety of sources, from active galaxies to low-mass pre-main sequence stars (classical T Tauri stars), can be also applied at very low stellar masses and, possibly, below the substellar limit (e.g., Fernรกndez & Comerรณn Fernandez01.1 (2001), Muzerolle et al. Muzerolle03.1 (2003), Comerรณn et al. Comeron03.1 (2003), Natta et al. Natta04.1 (2004), Barrado y Navascuรฉs & Jayawardhana Barrado04.2 (2004)). In fact, in the framework of low mass star formation, the mass accretion - mass loss proportionality proposed by Cabrit et al. (Cabrit90.1 (1990)) seems to hold below the substellar limit. Natta et al. (Natta04.1 (2004)) have studied the accretion properties of very low mass objects, more than doubling the number of substellar objects for which the mass accretion rate, แน<sub>acc</sub>, is known. They confirm the trend of lower แน<sub>acc</sub> for lower M, although with a large spread, possibly due to an age effect. This trend has been recently confirmed for the entire substellar domain, down to nearly the deuterium-burning limit (Mohanty et al. Mohanty05.1 (2005)). Little information can be found for mass loss rates, but the estimation made by Comerรณn et al. (Comeron03.1 (2003)) for the very low mass star LS-RCrA 1 falls in the expected range of values, according to the mentioned proportionality.
In order to confirm that the classical T Tauri star paradigm also applies at the very low mass regime, it is necessary to find out whether these objects can be the exciting sources of jets or outflows. Here we report the discovery of a jet emanating from a M5 star (Par-Lup3-4) and we confirm previous evidence (Fernรกndez & Comerรณn Fernandez01.1 (2001), Barrado y Navascuรฉs et al. Barrado04.1 (2004)) that the M6.5, or later, star LS RCrA-1 is the exciting source of an outflow. Both stars were discovered in the course of H<sub>ฮฑ</sub> surveys carried out in the Lupus 3 and R CrA regions, respectively (Comerรณn et al. Comeron03.1 (2003); Fernรกndez & Comerรณn Fernandez01.1 (2001)). Their strong emission at permitted lines indicates a strong mass accretion process, while the forbidden emissions make them very good candidates in the search for jets or outflows. These two stars also share an unexpected property: they are quite underluminous (Fernรกndez & Comerรณn Fernandez01.1 (2001); Comerรณn et al. Comeron03.1 (2003)). When compared to other young objects of similar spectral types, Par-Lup3-4 and LS-RCrA 1 happen to be fainter by almost 4 mag and 1.8 mag, respectively. These low luminosities put both stars on the 50 Myr isochrone on the HR diagram. This age is much older than those estimated for the other members of their young associations, which are well below 10 Myr. The results reported here suggest that neither of these objects is obscured by an edge-on disk, thus constraining possible explanations of the underluminosity.
## 2 Observations
All the observations reported in this paper were carried out with the VLT (Cerro Paranal, Chile) in service mode. The narrow-band \[SII\] and H<sub>ฮฑ</sub> imaging observations of Par-Lup3-4 took place on the night of May 2, 2003. They are integrations of 950 s each using the visible imaging and low-resolution spectrograph FORS1. We have measured a seeing of 0$`\stackrel{}{.}`$6 for the \[SII\] image and of 0$`\stackrel{}{.}`$7 for the H<sub>ฮฑ</sub> one.
High resolution ($`R`$ = 57,000) spectroscopy of Par-Lup3-4 was carried out over 5 nights in 2003, from July 4 to 30, using UVES, the Ultraviolet and Visual Echelle Spectrograph (Kaufer et al. Kaufer03.1 (2003)). Each observation consisted of two consecutive spectra covering the range $`\lambda \lambda `$ 3300 to 6800 ร
with exposure times of 1512 s. The slit width was always 1$`\stackrel{}{.}`$2. Since imaging and spectroscopy were scheduled in the same period, recognition of the Par-Lup3-4 jet in the FORS1 images was possible only after the UVES observations had been obtained, thus preventing the selection of a position angle (PA) of the UVES slit matching the direction of the jet. Instead, the parallactic angle was set so as to minimize losses due to atmospheric differential refraction, thus resulting in our observations probing a range of PAs. The log of spectroscopic observations is presented in Table 1. The third column contains the seeing values measured with the DIMM (Differential Image Motion Monitor) during the time of the exposures; such values are stored in the ESO Observatories Ambient Conditions Database<sup>1</sup><sup>1</sup>1http://archive.eso.org/asm/ambient-server. The range of parallactic angles covered, as well as the average one for each pair of exposures, are also listed.
Narrow-band imaging observations of LS-RCrA 1 using the same setup as for Par-Lup3-4 were carried out on June 2, 2003. High resolution spectroscopic observations of LS-RCrA 1 were carried out from June 3 to July 4, 2003. The instrumental setup was the same as for Par-Lup3-4.
The spectra were reduced and analyzed using IRAF<sup>2</sup><sup>2</sup>2IRAF is distributed by the National Optical Astronomy Observatories, which are operated by the Association of Universities for Research in Astronomy, Inc. (AURA), under cooperative agreement with the National Science Foundation, USA., paying special attention to the small spatially extended structure present in some of the spectra. All velocities are referred to the local standard of rest (LRS).
## 3 Results
### 3.1 Par-Lup3-4
Figure 1 shows three of the echellograms obtained for Par-Lup3-4. Each of them is the average of two consecutive exposures. The echellogram taken on July 27 (central panel) shows spatially unresolved \[SII\] emission lines, as expected from a point source. The data from the other two nights show emission from the surroundings of the star. The fact that on each of these two nights the redshifted emission (bottom of the lines) originates at opposite sides of the star is due to the different parallactic angles; all the emission comes, spatially, from the same side of the star. The observed emission can be traced up to 3$`\stackrel{}{.}`$6 from the star. Hints of very faint blueshifted emission coming from the other side of the star can be seen on the echellogram from July 4.
The extended emission is clearly detected on the \[SII\] narrow band image of Par-Lup3-4 (see Fig. 2), where a distinct knot is seen at a PA = 129$`\stackrel{}{.}`$7 and at a distance of 1$`\stackrel{}{.}`$3 from the star, corresponding to a 260 AU projected distance from the star assuming a distance of 200 pc for the Lupus region; see discussion in Comerรณn et al. (Comeron03.1 (2003)). Fainter jet-like emission can be traced further away up to 4$`\stackrel{}{.}`$2 from the star (840 AU), as well as in the opposite direction reaching up to 2$`\stackrel{}{.}`$0 (400 AU). The orientation of these features implies that the slit was oriented along the jet in our spectra of July 4, as seen on the right panel of Fig. 2, where the range of slit PAs covered by the echelle observations are plotted, giving the average PAs as a reference. The extended emission is, nevertheless, difficult to see on the narrow band H<sub>ฮฑ</sub> image, most probably due to the dominance of the H<sub>ฮฑ</sub> emission from the star. This extended emission is not detected on the H<sub>ฮฑ</sub> echellogram taken along the orientation of the jet.
The \[SII\] lines, presented in Fig. 3, show a double-peaked profile, in which the relative intensities of both peaks change. Since variability with time scales of weeks is not expected at such a distance from the exciting source, changes are likely to be due to the different PA of the slit probing different parts of the extended emission. The brighter and redshifted component is associated with the emission knot seen in the \[SII\] images to the south-east of Par-Lup3-4.
The \[NII\] lines are in emission in some of the spectra, but they are absent in others. We have found no correlation between the intensity of these lines and the slit PA, but the intensity seems to correlate with the seeing at the time of the observations: the better the seeing, the more intense the lines are. The spectra with better signal to noise show a double-peaked profile. This is a clear signature of a bipolar jet, since \[NII\] emission only originates in the high-velocity component of the bipolar outflows (Hirth et al. Hirth97.1 (1997)). The \[OI\] $`\lambda `$ 6300 and the \[OI\] $`\lambda `$ 6363 emission lines, on the contrary, present single peak profiles.
The LSR velocities measured for the \[OI\], \[NII\], and \[SII\] lines are listed in Table 2. These velocities correspond to the central position of each line or line component, in the case of double-peaked profiles.
We have also detected emission from the HeI $`\lambda `$5876 permitted line. Since this line is too noisy in the individual spectra, we have added up all of them and we have measured an equivalent width of 3.7$`{}_{}{}^{+0.8}{}_{1.7}{}^{}`$ ร
. An equivalent width of 1.6 ร
was already reported for the HeI $`\lambda `$6678 emission of this star by Comerรณn et al. (Comeron03.1 (2003)).
H<sub>ฮฑ</sub> emission, coming from an object located at 4$`\stackrel{}{.}`$2 from Par-Lup3-4, has been detected in the two spectra taken at PA 24 and 42 (average PA of 33). Several reasons let us discard the possibility that it is an artifact, i.e. light reflection: the little resemblance to the H<sub>ฮฑ</sub> profile of Par-Lup3-4; the lack of any bright feature on the spectrum (the brightest one amounts to less than 1000 counts); the fact that its position does not change from one spectrum to the next one, but it gets fainter, as if the object slowly moves out of the slit; and the fact that the H<sub>ฮฑ</sub> profile of the closest bright star (RX J1608.9 $``$3905, located at 1โ) is not in emission, but filled. From the position of this spectrum on the echellogram we estimate a PA of $``$30 or $``$210 for the unknown object<sup>3</sup><sup>3</sup>3The uncertainty in the PA is due to the fact that the spectrograph is allowed to turn 360.. The line that connects Par-Lup3-4 and this object subtends, thus, an angle of $``$80 with the jet. An outflow direction perpendicular to a binary axis has been reported for V536 Aql (Mundt & Eislรถffel, Mundt98.1 (1998)); the distance between both components being 0$`\stackrel{}{.}`$52 ($``$ 120 AU). The H<sub>ฮฑ</sub> emission of the new object, presented in Fig. 4, shows a wide, double-peaked profile with extended wings covering more than 300 km$``$s<sup>-1</sup>. The two peaks are centered at $``$-10 km$``$s<sup>-1</sup> and $``$90 km$``$s<sup>-1</sup>, the red peak being the faintest one. No visible counterpart has been found on R, I, and z band images up to limiting magnitudes of R=24 mag, I=22 mag, and z=21 mag (Comerรณn et al. in preparation). No infrared counterpart has been found brighter than J=20.8 mag, H=17.5 mag and K$`{}_{S}{}^{}`$20.8 mag in images obtained at the New Technology Telescope (NTT), at La Silla, with SOFI in June 2001 (J band) and July 2002 (H band), and with ISAAC at the VLT (H<sub>2</sub> filter) in August 2003. The line profile, as well as the velocities of the two intensity peaks, resembles the H<sub>ฮฑ</sub> emission of pre-main sequence objects (Fernรกndez et al. Fernandez95.1 (1995); Reipurth et al. Reipurth96.1 (1996)). If this is the case, and taking into account the fact that the object is more than $``$8 mag fainter than the young M5 stars of this star forming region, its position on the luminosity vs. age diagram of Burrows et al. (Burrows97.1 (1997), see their Fig.7) would fall on the regime of the young, very low mass brown dwarfs and planetary mass objects. Its H<sub>ฮฑ</sub> emission could be due to mass accretion or to a flare, like the one that has been recently observed with UVES on the old M9 dwarf DENIS 104814.7-395606.1 (Fuhrmeister & Schmitt Fuhrmeister04.1 (2004)). The flare option is, nevertheless, less plausible, because the H<sub>ฮฑ</sub> intensity does not change much from the first spectrum to the second and both were taken over a time interval of 50 minutes. Mass accretion, on the other hand, is known to play a very important role in the formation of brown dwarfs and, possibly, of Jovian planets (Quillen & Trilling Quillen98.1 (1998)). A conservative value of the full width at 10% of the H<sub>ฮฑ</sub> peak profile is in the range of 250$``$300 km$``$s<sup>-1</sup>, which supports the hypothesis of mass accretion (see Sec. 4.1).
Deep H<sub>2</sub> (2.12 $`\mu `$m) images of Par-Lup3-4 taken with ISAAC in August 2003 show a point source located at 1$`\stackrel{}{.}`$2 from it, at a PA$``$63ยฐ. If it has no emission lines in this band, the source is about 6 mag fainter than Par-Lup3-4 in K<sub>S</sub>, reaching K$`{}_{S}{}^{}`$19.6 mag. If physically related to Par-Lup3-4 the object could be in the planetary mass regime (Burrows et al. Burrows97.1 (1997)). On the echellograms obtained on July 4, 2003, we are not able to detect H<sub>ฮฑ</sub> emission at this position, perhaps because it is dominated by the emission from Par-Lup3-4, due to a seeing $``$1โณ.
### 3.2 LS-RCrA 1
The echellograms show no hint of extended emission close to the star and therefore there is no dependence on the PA of the slit. Fig. 5 shows, for this reason, the average forbidden lines obtained for LS-RCrA 1.
Unlike in the case of Par-Lup3-4, no double-peaked profile is identified in any of the forbidden lines. However, asymmetric profiles are clearly seen. We like to note that such types of profiles are more common among classical T Tauri stars than the double-peaked profiles seen in Par-Lup3-4 (see Hirth et al. Hirth97.1 (1997)). The asymmetry of the \[OI\] lines seems to be due to the absorption of the redshifted emission; while the \[SII\] lines show a bump on the red wing, centered at $``$ 50 km s<sup>-1</sup>. The \[NII\] lines present quite symmetric profiles, but both of them show faint emission at $``$ 50 km s<sup>-1</sup>. All the forbidden lines have a FWHM of about 40 km$``$s<sup>-1</sup>.
We detect emission from both \[NII\] $`\lambda `$6583 and \[NII\] $`\lambda `$6548, two lines that have been reported to form only a high velocity component (HVC) and not a low velocity component (LVC) in classical T Tauri stars (Hirth et al. Hirth97.1 (1997)).
## 4 Discussion
### 4.1 Par-Lup3-4
As expected from its low resolution spectrum (Comerรณn et al. Comeron03.1 (2003)), the observed forbidden lines of Par-Lup3-4 are characteristic of the shocked gas usually observed in jets from pre-main sequence stars. The precise physical characteristics of Par-Lup3-4 are difficult to determine due to its anomalously low luminosity, which prevents its comparison with theoretical evolutionary tracks in the temperature-luminosity diagram. Nevertheless, the late-type spectrum of the central object places Par-Lup3-4 among the least massive objects known to excite a jet. Other very low mass objects have been also reported to power outflows. Froebrich et al. (Froebrich03.1 (2003)) found VLA 1623 to be the lowest mass star among a sample of Class 0 sources powering outflows; they estimate that it will reach a mass of 0.07 M. Very low masses are also expected for BKLT J162658-241836 and WLY 2-36, the likely exciting sources of Herbig-Haro objects in the $`\rho `$ Ophiuchi embedded cluster (Gรณmez et al. Gomez03.1 (2003)).
The fact that only one side of the jet is clearly detected on the visible images is not strange among young stars. One-sided jets, with a very faint counterjet, have been reported for several of them (e.g., DG Tau, Solf & Bรถhm Solf93.1 (1993)). Hirth et al. (Hirth94.1 (1994)) discuss asymmetries in bipolar jets from young stars, which can be related to the source itself or to its immediate environment.
The lack of photospheric features on the weak continuum of the high resolution spectra prevents us from computing the LSR velocity of the star. Nevertheless, an estimation of this velocity can be obtained from the CO observations carried out by Gahm et al. (Gahm93.1 (1993)) towards the Lupus 3 cloud, which gave an average value between 5 and 6 km$``$s<sup>-1</sup>. More recently, Hara et al. (Hara99.1 (1999)) measured the LSR velocities of a C<sup>18</sup>O core located at less than 1โ from the position of Par-Lup3-4; the beam size of the telescope was 2$`\stackrel{}{.}`$6. They measured a velocity of 4.13 $`\pm `$1.2 km$``$s<sup>-1</sup>. Since C<sup>18</sup>O traces only the dense parts of the clouds, unlike <sup>12</sup>CO, the C<sup>18</sup>O velocity relates more to young stars formed very recently. The average LSR velocities that we measure for the jet of Par-Lup3-4, 4.7$`\pm `$1.7 km$``$s<sup>-1</sup> for \[SII\] $`\lambda `$6716 and 2.7$`\pm `$1.5 km$``$s<sup>-1</sup> for \[SII\] $`\lambda `$6731, match the values obtained from radio observations, thus confirming that the red- and blue-shifted peaks come from the jet and the counterjet, respectively.
Gaussian fits to the double-peaked \[SII\] lines show an average difference between peaks of $`41`$ km$``$s<sup>-1</sup>. If we assume a jet velocity of 150 km$``$s<sup>-1</sup> perpendicular to the plane of the disk, we get a disk tilt of 8ยฐ with respect to the plane of the sky, if there is a symmetric distribution of velocities in the jet. For a velocity of about 100 km$``$s<sup>-1</sup> (the measured velocity width at the base of the individual line components of the \[SII\] emissions) the corresponding inclination is 12ยฐ. Flared disks, for which the ratio of the disk scale height H to the radial distance R increases with R (see Hartmann Hartmann98.1 (1998)), can hide the star more easily than flat disks. To date only small samples are available for the study of the frequency of flared disks among brown dwarfs; nevertheless, this frequency does not seem to be high. From the study of the disks around 12 brown dwarfs, Natta et al. (Natta02.1 (2002)) concluded that nine of them might have flat disks, in spite of the strong bias of their sample against objects with flat disks. Mohanty et al. (Mohanty04.1 (2004)), on the other side, found strong evidence of flared disks for two brown dwarfs out of a sample of three and Sterzik et al. (Sterzik04.1 (2004)) reported a flare disk geometry for ChaH<sub>ฮฑ</sub>1.
The H<sub>ฮฑ</sub> emission of three T Tauri stars with edge-on disks, where the contributions arising from the surface and its closest vicinity are blocked from direct view, has been recently studied by Appenzeller et al. (Appenzeller05.1 (2005)). In all cases the line has a narrow profile, with a full width half maximum (FWHM) below 100 km$``$s<sup>-1</sup>. Appenzeller et al. interpret it as the H<sub>ฮฑ</sub> contribution to the outflows. This result strongly supports the conclusion of White & Basri (White03.1 (2003)), confirmed by Natta et al. (Natta04.1 (2004)), that H<sub>ฮฑ</sub> emission can only be undoubtedly attributed to the mass accretion process, if the full width of the emission profile at 10% of the maximum intensity (hereafter 10% width) is above<sup>4</sup><sup>4</sup>4For a Gaussian function the 10% width is 1.8226 times the FWHM. Note, however, that stellar line profiles are not always Gaussian. 270 km$``$s<sup>-1</sup>. Jayawardhana et al. (Jayawardhana03.1 (2003)) suggest, nevertheless, that for some accreting objects the 10% width could be as low as 200 km$``$s<sup>-1</sup>. The wide and complex H<sub>ฮฑ</sub> profile that we have observed for Par-Lup3-4, with a 10% width in the range from 340 to 400 km$``$s<sup>-1</sup>, should come, then, from the accretion related regions, which lie very close to or on the stellar surface, strongly supporting the non-edge disk hypothesis. We also detect in our spectra HeI $`\lambda `$5876 emission, which is usually interpreted as being formed very close to the stellar surface. However, Appenzeller et al. (Appenzeller05.1 (2005)) have detected it also in the spectrum of their sample of stars with edge-on disks, thus suggesting that HeI emission can also be produced far from the surface.
The flux ratios of the observed forbidden lines inform about the physical characteristics of the jet. Bacciotti & Eislรถffel (Bacciotti99.1 (1999)) have developed a technique which allows one to determine the local ionization fraction, the electron density and electron temperature using these ratios. The low signal to noise along the jet prevents us from carrying a detailed spatial study of these line ratios, but we can get average values for the whole jet. The observed line ratios are
$$\frac{[SII]\lambda 6716}{[SII]\lambda 6731}=0.64\pm 0.04\{\begin{array}{c}\text{blue component }0.74\pm 0.11\hfill \\ \text{red component }0.58\pm 0.06\hfill \end{array}$$
$$\frac{[SII](\lambda 6716+\lambda 6731)}{[OI](\lambda 6300+\lambda 6363)}=0.37\pm 0.04$$
$$\frac{[OI](\lambda 6300+\lambda 6363)}{[NII](\lambda 6548+\lambda 6583)}14.8$$
From comparison to other jets for which these lines have been measured (Bacciotti & Eislรถffel 1999), we can conclude that we deal with a low excitation jet.
The widths of the \[OI\] lines are very similar to those of the \[SII\] lines, in contrast to what has been observed for some classical T Tauri stars (Hirth et al. Hirth94.1 (1994)). Hirth et al. showed that the \[OI\] lines form at smaller distances from the star than the \[SII\] lines. The fact that the velocities involved in the formation of both sets of lines seem to be similar could suggest that the collimation mechanism is already working very close to the star.
### 4.2 LS-RCrA 1
Like for Par-Lup3-4, the non-detection of the photospheric continuum in the high resolution spectra of LS-RCrA 1 prevents us from computing the LSR velocity of the star. Barrado y Navascuรฉs et al. (Barrado04.1 (2004)) have measured a radial velocity of 2$`\pm 3`$ km$``$s<sup>-1</sup>. This velocity falls slightly outside of the range of velocities that Neuhรคuser et al. (Neuhaeuser00.1 (2000)) have measured for 12 earlier type T Tauri stars in the CrA complex (0 to -5 km$``$s<sup>-1</sup>), but it contains within its error bars the LSR H<sub>2</sub>CO velocity reported by Loren (Loren79.1 (1979)) at a position located at less than 40โ from the star and measured with a beam size of 2$`\stackrel{}{.}`$3; for this position Loren reported a velocity of 6.1$`\pm `$0.6 km$``$s<sup>-1</sup>. More recently, Vilas-Boas et al. (VilasBoas00.1 (2000)) observed C<sup>18</sup>O emission from a condensation located at less than 4โ from LS-RCrA 1, using a beam size of about 1$`\stackrel{}{.}`$5, and they obtained a LSR velocity of 5.65$`\pm `$1.18 km$``$s<sup>-1</sup>.
Taking either 6 or 2 km$``$s<sup>-1</sup> as the stellar LSR velocity, all the measured forbidden lines are blueshifted with respect to the star (see Table 2). This may be interpreted as meaning that only the blueshifted component of the outflow is seen from our vantage point, with the redshifted one probably occulted by a circumstellar disk, as observed in most classical T Tauri stars. This hypothesis is supported by the asymmetric line profile of the \[OI\] lines, which display an extended blue wing but miss the red wing. These results, together with the fact that the forbidden lines are known to form at different distances from the star, strongly argue in favor of a disk seen at a geometry markedly different from edge-on. As in the case of Par-Lup3-4, the 10% width of the H<sub>ฮฑ</sub> emission, with values in the range from 265 to 300 km$``$s<sup>-1</sup>, indicates that H<sub>ฮฑ</sub> is dominated by the accretion component formed near the surface of the star.
Barrado y Navascuรฉs et al. (Barrado04.1 (2004)) favour the hypothesis of an edge-on disk, as an explanation for the puzzling aspects of LS-RCrA 1, namely the lack of near-infrared (NIR) excess combined with accretion, the unusually prominent outflow signatures without high-velocity components or asymmetries, the very broad H<sub>ฮฑ</sub>, and the sub-luminosity. Nevertheless, we think that there are also explanations for all these features in the framework of the non edge-on disk hypothesis. No NIR excesses are expected for very low mass stars and brown dwarfs, as has been modeled by Natta & Testi (Natta01.1 (2001)) and has been confirmed by Barrado y Navascuรฉs et al. (Barrado04.3 (2004)), except for very few objects; such an excess is predicted for wavelengths longer than $``$3$`\mu `$m. Our high resolution spectra show asymmetries on the forbidden line profiles. The very broad H<sub>ฮฑ</sub> profile indicates an unimpeded view to the close proximity of the stellar surface, where the largest velocities of the H<sub>ฮฑ</sub>-emitting gas arise (Appenzeller et al. Appenzeller05.1 (2005)). Concerning the sub-luminosity, we still support the hypothesis suggested by Fernรกndez & Comerรณn (Fernandez01.1 (2001)) that, at these low masses, strong mass accretion might have an important effect on the position of the star on the HR diagram.
The central velocities measured for the forbidden lines differ notably. The lower velocities are found for the \[OI\] lines, while the \[NII\] lines give the highest values (see Table 2). The \[OI\] $`\lambda `$ 6300 line has a critical density higher than that of \[NII\] 6583, and roughly 100 times that of the \[SII\] lines (Hartigan et al. Hartigan95.1 (1995)). Hartigan et al. analyzed a sub-sample of four stars for which they detect emission from three lines, \[OI\] $`\lambda `$ 5577, \[OI\] $`\lambda `$ 6300, and \[SII\] $`\lambda `$ 6731, that have very different critical densities. The low velocity component of the three emission lines show a correlation in which the lower velocities correspond to the line with highest critical density (\[OI\] $`\lambda `$5577); a similar behaviour is observed for LS-RCrA 1. They found this correlation to be consistent with acceleration in a disk wind, but also with an origin in an accretion column, because in either case the lines with higher critical density (like \[OI\]) form closer to the disk than the lines with low critical density, and the flow accelerates as it rises from the disk. Hirth et al. (Hirth97.1 (1997)) found that for a sample of 12 T Tauri stars, located at $``$120-140 pc, the centroid of the \[OI\] 6300 emission is located at an average distance of 0$`\stackrel{}{.}`$2 from the star, whereas that of the \[SII\] $`\lambda `$6731 and \[NII\] $`\lambda `$6583 lines are factors of 3 and 3.5 times further away, respectively.
The bumps observed at about 50 km$``$s<sup>-1</sup> on the \[SII\] and \[NII\] lines may be due to either matter ejected closer to our line of sight or to a faster knot. Similar bumps have been observed, e.g., for the T Tauri stars DK Tau, GG Tau ad IP Tau (Hartigan et al. Hartigan95.1 (1995)).
## 5 Conclusions
We report the discovery of a jet emanating from the very low mass star Par-Lup3-4 (M5), and we confirm previous evidence of strong mass loss from another very low mass star, LS-RCrA 1 (M6.5 or later), most probably in the form of a jet or disk wind.
The line ratios of the forbidden lines of the jet of Par-Lup3-4 point to a low excitation jet. The double-peaked \[SII\] emission, centered on the LSR velocity of the Lupus 3 cloud in this region, allows us to set lower limits for the jet inclination: angles below 8ยฐ, with respect to the plane of the sky, would imply unlikely velocities above 150 km$``$s<sup>-1</sup>. With such inclination only a very flared disk would hide the star. The large 10% width is attributed to the accretion related regions, which lie very close to or on the stellar surface, suggesting that the large H<sub>ฮฑ</sub> equivalent width measured for this object is not due to the selective blocking of the central object by an edge-on disk.
H<sub>ฮฑ</sub> emission, coming for an object located at 4$`\stackrel{}{.}`$2 from Par-Lup3-4, has been detected in the two spectra taken at PA $``$ 33. The line profile, different from that of Par-Lup3-4, resembles that of other pre-main sequence objects. Upper limits for its brightness at visible and near infrared wavelengths suggest that, if associated with the Lupus 3 star forming region, it could be a young, very low mass brown dwarf.
All the forbidden lines that we have measured for LS-RCrA 1 are blueshifted with respect to the LSR velocity of the star. The emission from the receding part of the jet seems to be hidden by a non edge-on disk; a hypothesis that is supported by the fact that we detect H<sub>ฮฑ</sub> emission coming from the accretion related regions located close to the surface of the star. The velocities of the \[OI\] and \[SII\] forbidden emission lines are ordered inversely with their respective critical densities. This has been interpreted by Hartigan et al. (Hartigan95.1 (1995)), for more massive classical T Tauri stars, as acceleration from the most dense regions, close to the star, what would be consistent with acceleration in a disk wind, but also with an origin in an accretion column.
If both Par-Lup3-4 and LS-RCrA 1 have no edge-on disk, an alternative explanation is required in order to explain their unusual low luminosities. Strong accretion has been suggested to modify the position of classical T Tauri stars on the HR diagram (Hartmann et al. Hartmann97.1 (1997), Siess et al. Siess97.1 (1993)). Extending the modelling of the accretion effects towards the lowest stellar and substellar masses may indicate whether or not this is a viable explanation for the observed properties of these objects.
###### Acknowledgements.
We acknowledge the ESO staff who carried out the service mode observations, the ESO User Support Group for their valuable assistance in the preparation of our observations, and the ESO Data Flow Operations Group for the preparation of our data package. A. Kaufer, S. DโOdorico and L. Kaper, authors of the UVES User Manual, are also acknowledged, as well as those who prepared the UVES web pages. Fruitful discussions with Eike Guenther, Jochen Eislรถffel, Jens Woitas, Enrique Pรฉrez, Reinhard Mundt and Ferdinando Patat were very helpful, as were comments from David Barrado y Navascuรฉs and from David Butler. MF acknowledges ESO and the Thรผringer Landessternwarte (Germany) for their hospitality. She received support from the Deutsches Zentrum fรผr Luft- und Raumfahrt (DLR), Fรถrderkennzeichen 50 OR 0401, and from the Spanish grant AYA2004-05395. This work has made use of the Digitized Sky Surveys, produced at the Space Telescope Science Institute under U.S. Government grant NAG W-2166, and of the NASA/IPAC Infrared Science Archive, which is operated by the Jet Propulsion Laboratory, California Institute of Technology, under contract with the National Aeronautics and Space Administration.
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# Free Isotropic-Nematic Interfaces in Fluids of Charged Platelike Colloids
## I Introduction
Platelike colloidal particles play a decisive, constitutive role in processes like agriculture, construction, oil drilling, or coating. This wide range of applicability mirrors a very rich phase behavior of suspensions of platelets, including liquid crystalline phases, sol-gel transitions, and flocculation, depending on numerous material parameters, like size, shape, or charge of the particles, as well as on effective, solvent mediated interactions which can be tuned, e.g., by the choice of the solvent, salt concentration, or $`p\mathrm{H}`$-value. Whereas bulk properties of suspensions of charged platelets have been investigated for decades Davi05 , free interfaces between coexisting fluid phases in such systems have not yet been studied. Here we focus on this latter issue by proposing and studying a density-functional theory of such inhomogeneous multicomponent systems of charged anisotropic particles.
On the experimental side, characterization of bulk phases have been conducted for several model systems like natural clay Lang38 ; Brit69 , laponite Mour95 ; Gabr96 ; Pign97 ; Kroo98 ; Mour98 ; Saun99 ; Levi00 ; Nico00 ; Nico01 ; Pori01 ; Pori03 , sterically stabilized gibbsite Kooi98 ; Kooi01 ; Beek03 ; Beek04 , or nickel(II)hydroxide Brow98 ; Brow99 using methods like polarized light analysis Gabr96 ; Kooi98 ; Kooi01 ; Beek03 ; Beek04 , light scattering Kroo98 ; Nico00 ; Nico01 , small-angle scattering with neutrons or x-rays Mour95 ; Pign97 ; Kroo98 ; Saun99 ; Brow98 ; Brow99 , rheological measurements Mour95 ; Mour98 ; Levi00 , or NMR Pori01 ; Pori03 . Since coexistence between bulk phases of charged gibbsite platelets has already been observed experimentally Beek03 ; Beek04 , we expect that the spatially varying structural properties between them are also experimentally accessible using, e.g., scattering Diet95 or optical Bain98 techniques.
The theoretical description of suspensions of charged platelets is rather complicated due to long-ranged, anisotropic interactions and many different length scales. Interfaces and surfaces in such systems add further difficulties induced by the partial loss of translational symmetry. Under these circumstances it is advisable to start with simplified models. In the presented one, the particles are modeled as hard cuboids with pointlike charges concentrated in their center. Furthermore, the platelet orientations are restricted to three mutually perpendicular directions, which is commonly known as the Zwanzig model Zwan63 . Finally, in order to gain computational advantages, distances between particles are not measured by the usual Euclidean norm but by the so-called supremum norm.
Density-functional theory Evan79 is a very effective method to investigate inhomogeneous fluid systems. It has recently been applied to describe suspensions of platelets with pure hard-core interactions near interfaces and surfaces Harn02a ; Harn02b ; Bier04 . Here we study platelets interacting via a hard-core plus a Coulomb potential; further interactions like dispersion forces are not considered, which corresponds to a suspension in which the indices of refraction between solvent and solute are matched. The density functional is constructed by functional integration of the two-particle density with respect to the interaction potential Evan79 , which is the analogue to a Debye-charging process Deby23 . The two-particle density can be obtained, e.g., by interaction site model calculations Harn01 ; Harn02c or expressed in terms of the potential of mean force which may be approximated by effective pair potentials Soga91 ; Rowa00 ; Triz02 . For reasons of computational advantages, here we choose an extension of the even simpler Debye-Hรผckel pair distribution function Deby23 in which the Debye screening factor is replaced by a spatially varying quantity.
In view of these approximations, this theoretical model can be expected to be only qualitatively correct. The most subtle point is the choice of the pair distribution function, which requires further improvements in order to increase the quantitative reliability. However, the other parts of the formalism are expected to remain valid and thus provide the basis for future efforts. In this sense, the following sections present a generic formalism for free interfaces in fluids of charged platelike colloids, implemented exemplarily for the above-mentioned extended Debye-Hรผckel pair distribution function.
The text is structured as follows. In Sec. II the general formalism with a detailed derivation of the actual density functional is presented. Section III is devoted to bulk phase diagrams of the model. The structures of free interfaces between an isotropic and a nematic bulk phase are calculated in Sec. IV. Section V discusses the current approach and summarizes our results.
## II General formalism
### II.1 Definitions
We consider a ternary mixture of charged hard square cuboids with their edges required to be parallel to the Cartesian axes (Zwanzig model Zwan63 ) dissolved in a dielectric solvent (e.g., water) with dielectric constant $`\epsilon `$. The solvent is treated as a continuum. For simplicity, the charges are fixed, monodisperse, and concentrated in the centers of the particles. The particles of the first component, representing the macroions $`\mathrm{M}`$, have size $`D_\mathrm{M}\times D_\mathrm{M}\times L_\mathrm{M}`$, $`D_\mathrm{M}L_\mathrm{M}`$, and charge $`Q_\mathrm{M}0`$. Within the Zwanzig approximation, macroions can take three different orientations, denoted as $`\mathrm{M}_\mathrm{x}`$, $`\mathrm{M}_\mathrm{y}`$, or $`\mathrm{M}_\mathrm{z}`$ corresponding to whether the $`L_\mathrm{M}`$-edges are parallel to the $`x`$-, $`y`$-, or $`z`$-axis, respectively (see Fig. 1). The second component consists of salt anions $`\mathrm{A}`$ modeled as cubes of side length $`D_\mathrm{A}:=D_\mathrm{S}`$ and charge $`Q_\mathrm{A}:=Q_\mathrm{S}<0`$ (see Fig. 1). Finally, the third component consists of salt cations $`\mathrm{C}`$ and counterions guaranteeing overall charge neutrality. They are also described by cubes with the same side length $`D_\mathrm{C}:=D_\mathrm{S}`$ but opposite charge $`Q_\mathrm{C}:=Q_\mathrm{S}>0`$ (see Fig. 1).
We denote as $`\varrho _i(๐ซ),i\{\mathrm{M}_\mathrm{x},\mathrm{M}_\mathrm{y},\mathrm{M}_\mathrm{z},\mathrm{A},\mathrm{C}\}`$, the number density at point $`๐ซ`$ of the centers of macroions with orientation $`\mathrm{M}_{\mathrm{x},\mathrm{y},\mathrm{z}}`$, anions, and cations, respectively. Note that the position $`๐ซV^3`$, with $`V`$ denoting the system volume, is a *continuous* variable in contrast to the orientation of macroions, which varies within a *discrete* set. As an abbreviation, we introduce $`\underset{ยฏ}{\varrho }(๐ซ):=(\varrho _{\mathrm{M}_\mathrm{x}}(๐ซ),\mathrm{},\varrho _\mathrm{C}(๐ซ))`$.
The system under consideration is coupled to two particle reservoirs: One supplies neutralized macroions (chemical formula $`\mathrm{C}_k\mathrm{M},k:={\displaystyle \frac{Q_\mathrm{M}}{Q_\mathrm{S}}}`$) and the other neutral salt (chemical formula $`\mathrm{CA}`$); $`\mu _{\mathrm{C}_k\mathrm{M}}`$ and $`\mu _{\mathrm{CA}}`$ denote the corresponding chemical potentials. Upon entering the solvent, these molecules dissociate:
$`\mathrm{C}_k\mathrm{M}`$ $``$ $`k\mathrm{C}^{Q_\mathrm{C}}+\mathrm{M}^{Q_\mathrm{M}},`$
$`\mathrm{CA}`$ $``$ $`\mathrm{C}^{Q_\mathrm{C}}+\mathrm{A}^{Q_\mathrm{A}}.`$ (1)
These equilibrium chemical reactions lead to the following relations between the reservoir chemical potentials ($`\mu _{\mathrm{C}_k\mathrm{M}}`$ and $`\mu _{\mathrm{CA}}`$) and the particle chemical potentials ($`\mu _i,i\{\mathrm{M}_\mathrm{x},\mathrm{M}_\mathrm{y},\mathrm{M}_\mathrm{z},\mathrm{A},\mathrm{C}\},\mu _{\mathrm{M}_\mathrm{x}}=\mu _{\mathrm{M}_\mathrm{y}}=\mu _{\mathrm{M}_\mathrm{z}}`$):
$`\mu _{\mathrm{C}_k\mathrm{M}}`$ $`=`$ $`k\mu _\mathrm{C}+\mu _{\mathrm{M}_{\mathrm{x},\mathrm{y},\mathrm{z}}}`$
$`\mu _{\mathrm{CA}}`$ $`=`$ $`\mu _\mathrm{C}+\mu _\mathrm{A}.`$ (2)
### II.2 Density-functional theory
The configurations of this system are characterized by the set of number density profiles $`\underset{ยฏ}{\varrho }`$. The equilibrium states minimize the grand canonical density functional Evan79 ; Units
$$\mathrm{\Omega }[\underset{ยฏ}{\varrho }]=\underset{i}{}\underset{V}{}\mathrm{d}^3r\varrho _i(๐ซ)\left(\mathrm{ln}\left(\varrho _i(๐ซ)\right)1\mu _i^{}\right)+F^{\mathrm{ex}}[\underset{ยฏ}{\varrho }],$$
(3)
where $`F^{\mathrm{ex}}`$ is the free energy in excess over the ideal gas contribution. Here, the reduced particle chemical potentials $`\mu _i^{}:=\mu _i\mathrm{ln}\left(\mathrm{\Lambda }_i^3\right)`$ with the thermal de Broglie wavelength $`\mathrm{\Lambda }_i`$ ($`\mathrm{\Lambda }_{\mathrm{M}_\mathrm{x}}=\mathrm{\Lambda }_{\mathrm{M}_\mathrm{y}}=\mathrm{\Lambda }_{\mathrm{M}_\mathrm{z}}`$) for particles of class $`i`$ have been introduced. With the reduced reservoir chemical potentials
$`\mu _{\mathrm{C}_k\mathrm{M}}^{}`$ $`:=`$ $`\mu _{\mathrm{C}_k\mathrm{M}}\left(\mathrm{ln}\left(\mathrm{\Lambda }_{\mathrm{M}_{\mathrm{x},\mathrm{y},\mathrm{z}}}^3\right)+k\mathrm{ln}\left(\mathrm{\Lambda }_\mathrm{C}^3\right)\right)`$
$`\mu _{\mathrm{CA}}^{}`$ $`:=`$ $`\mu _{\mathrm{CA}}\left(\mathrm{ln}\left(\mathrm{\Lambda }_\mathrm{A}^3\right)+\mathrm{ln}\left(\mathrm{\Lambda }_\mathrm{C}^3\right)\right),`$ (4)
Eq. $`(\text{2})`$ takes the form
$`\mu _{\mathrm{C}_k\mathrm{M}}^{}`$ $`=`$ $`k\mu _\mathrm{C}^{}+\mu _{\mathrm{M}_{\mathrm{x},\mathrm{y},\mathrm{z}}}^{}`$
$`\mu _{\mathrm{CA}}^{}`$ $`=`$ $`\mu _\mathrm{C}^{}+\mu _\mathrm{A}^{}.`$ (5)
For given reservoir chemical potentials $`\mu _{\mathrm{C}_k\mathrm{M}}^{}`$ and $`\mu _{\mathrm{CA}}^{}`$, the particle chemical potentials $`\mu _i^{}`$ are fixed by Eq. $`(\text{5})`$ and the constraint of global charge neutrality:
$$\underset{V}{}\mathrm{d}^3r\underset{i}{}Q_i\varrho _i(๐ซ)=0.$$
(6)
The Euler-Lagrange equations of the minimization problem read
$$\frac{\delta \mathrm{\Omega }}{\delta \varrho _i(๐ซ)}=\mathrm{ln}\left(\varrho _i(๐ซ)\right)\mu _i^{}c_i(๐ซ)=0$$
(7)
with the one-particle direct correlation function
$$c_i(๐ซ):=\frac{\delta F^{\mathrm{ex}}}{\delta \varrho _i(๐ซ)}.$$
(8)
If $`\underset{ยฏ}{\varrho }=\underset{ยฏ}{\varrho }^{\mathrm{eq}}`$ minimizes the density functional in Eq. $`(\text{3})`$, the grand potential $`\mathrm{\Omega }(T,V,\mu _{\mathrm{C}_k\mathrm{M}}^{},\mu _{\mathrm{CA}}^{})=p(T,V,\mu _{\mathrm{C}_k\mathrm{M}}^{},\mu _{\mathrm{CA}}^{})V`$ with the osmotic pressure $`p(T,V,\mu _{\mathrm{C}_k\mathrm{M}}^{},\mu _{\mathrm{CA}}^{})`$ equals $`\mathrm{\Omega }[\underset{ยฏ}{\varrho }^{\mathrm{eq}}]`$. Phase coexistence corresponds to different states with equal values of the pressure $`p`$, the chemical potential $`\mu _{\mathrm{C}_k\mathrm{M}}^{}`$ of the reservoir of neutralized platelets, and the chemical potential $`\mu _{\mathrm{CA}}^{}`$ of the salt reservoir. In particular, coexistence does *not* imply equal values of the particle chemical potentials $`\mu _i^{}`$, $`i\{\mathrm{M}_\mathrm{x},\mathrm{M}_\mathrm{y},\mathrm{M}_\mathrm{z},\mathrm{A},\mathrm{C}\}`$. Rather, coexisting bulk phases give rise to a Donnan potential maintaining different chemical potentials $`\mu _i^{}`$.
### II.3 Excess free energy
The above considerations are valid for any interaction between the particles. In this subsection, our choice of the model and the resulting excess free energy $`F^{\mathrm{ex}}`$ are specified.
The interaction energy $`U_{ij}(๐ซ,๐ซ^{})`$ of a particle of class $`i`$ at position $`๐ซ`$ with a particle of class $`j`$ at position $`๐ซ^{}`$ comprises a hard-core potential $`U_{ij}^\mathrm{h}(๐ซ,๐ซ^{})`$, which prevents the particles from overlapping, and a contribution $`U_{ij}^\mathrm{c}(๐ซ,๐ซ^{})`$ due to the charges: $`U=U^\mathrm{h}+U^\mathrm{c}`$. As stated in the Introduction (Sec. I), we do not consider dispersion forces.
The interactions between the charges are approximated as
$$U_{ij}^\mathrm{c}(๐ซ,๐ซ^{}):=\frac{Q_iQ_j}{๐ซ๐ซ^{}_{\mathrm{}}},$$
(9)
where the usual Euclidean norm $`๐ซ_2=\sqrt{x^2+y^2+z^2}`$ is replaced by the supremum norm $`๐ซ_{\mathrm{}}=\mathrm{max}(|x|,|y|,|z|)`$ because of computational advantages. Since these two norms are equivalent, i.e., $`๐ซ_{\mathrm{}}๐ซ_2\sqrt{3}๐ซ_{\mathrm{}}`$, we do not expect that the results change qualitatively due to this approximation. Furthermore, $`_{\mathrm{}}`$-spheres are cubes with their edges parallel to the Cartesian axes; therefore, the supremum norm is the most natural and adapted norm in the context of a Zwanzig model for cuboids.
As described in Ref. Evan79 , the exact relation
$$\frac{\delta F^{\mathrm{ex}}}{\delta U_{ij}(๐ซ,๐ซ^{})}=\frac{1}{2}\varrho _i(๐ซ)\varrho _j(๐ซ^{})g_{ij}(๐ซ,๐ซ^{})$$
(10)
with the pair distribution function $`g`$ corresponding to the pair potential $`U`$ can be functionally integrated along the path $`U^{(\eta )}:=U^\mathrm{h}+\eta U^\mathrm{c},\eta [0,1]`$, which yields
$`F^{\mathrm{ex}}[\underset{ยฏ}{\varrho }]`$ $`=`$ $`F^{\mathrm{ex},\mathrm{h}}[\underset{ยฏ}{\varrho }]+{\displaystyle \frac{1}{2}}{\displaystyle \underset{ij}{}}{\displaystyle \underset{V}{}}\mathrm{d}^3r{\displaystyle \underset{V}{}}\mathrm{d}^3r^{}`$ (11)
$`\varrho _i(๐ซ)\varrho _j(๐ซ^{})U_{ij}^\mathrm{c}(๐ซ,๐ซ^{}){\displaystyle \underset{0}{\overset{1}{}}}d\eta g_{ij}^{(\eta )}(๐ซ,๐ซ^{}),`$
where $`F^{\mathrm{ex},\mathrm{h}}`$ is the excess free energy corresponding to the pure hard-core potential $`U^\mathrm{h}`$, and $`g^{(\eta )}`$ denotes the (inhomogeneous) pair distribution function for the pair potential $`U^{(\eta )}`$. $`F^{\mathrm{ex},\mathrm{h}}`$ is chosen as the fundamental measure functional introduced by Cuesta and Martรญnez-Ratรณn Cues97a ; Cues97b :
$$F^{\mathrm{ex},\mathrm{h}}[\underset{ยฏ}{\varrho }]:=\underset{V}{}\mathrm{d}^3r\mathrm{\Phi }\left(\underset{ยฏ}{n}(๐ซ)\right)$$
(12)
with the weighted densities
$$n_\alpha (๐ซ)=\underset{i}{}\underset{V}{}\mathrm{d}^3r^{}\omega _{\alpha ,i}(๐ซ๐ซ^{})\varrho _i(๐ซ^{})$$
(13)
for $`\alpha \{0,1x,1y,1z,2x,2y,2z,3\}`$ and the excess free energy density
$`\mathrm{\Phi }(\underset{ยฏ}{n})`$ $`=`$ $`n_0\mathrm{ln}(1n_3)+{\displaystyle \underset{q\{x,y,z\}}{}}{\displaystyle \frac{n_{1q}n_{2q}}{1n_3}}+{\displaystyle \frac{n_{2x}n_{2y}n_{2z}}{(1n_3)^2}}.`$
Due to Eq. $`(\text{9})`$, the $`\eta `$-integration in Eq. $`(\text{11})`$ may be interpreted as a Debye charging process Deby23 ; McQu00 . This motivates to approximate $`g^{(\eta )}`$ by an expression similar to the pair distribution function of the Debye-Hรผckel theory Deby23 ; Hans86 ; Lee97 ; McQu00 :
$`g_{ij}^{(\eta )}(๐ซ,๐ซ^{})`$ $`:=\mathrm{exp}(U_{ij}^\mathrm{h}(๐ซ,๐ซ^{}))\mathrm{max}[0,1U_{ij}^\mathrm{c}(๐ซ,๐ซ^{})`$ (15)
$`\times \eta \mathrm{exp}(\sqrt{\eta }\kappa (๐ซ,๐ซ^{};[\underset{ยฏ}{\varrho }])๐ซ๐ซ^{}_{\mathrm{}})].`$
This pair distribution function has non-negative values, it vanishes within the hard-core, and it approaches unity at infinitely large distances. Again, here the Euclidean norm $`_2`$ has been replaced by the supremum norm $`_{\mathrm{}}`$. The factor $`\sqrt{\eta }`$ in Eq. $`(\text{15})`$ is introduced because in a Debye charging process all charges $`Q_i`$ are replaced by $`\sqrt{\eta }Q_i`$. Furthermore, the charges $`Q_i`$ have to be interpreted as effective charges in order to reproduce the actual effective interactions between charged particles within Debye-Hรผckel theory Aubo03 .
The screening factor $`\kappa `$ in Eq. $`(\text{15})`$ is chosen as
$$\kappa (๐ซ,๐ซ^{};[\underset{ยฏ}{\varrho }]):=\frac{1}{2}\left(\stackrel{~}{\kappa }(๐ซ;[\underset{ยฏ}{\varrho }])+\stackrel{~}{\kappa }(๐ซ^{};[\underset{ยฏ}{\varrho }])\right)$$
(16)
where only anions and cations contribute to the screening Warr00 :
$$\stackrel{~}{\kappa }(๐ซ;[\underset{ยฏ}{\varrho }]):=\sqrt{4\pi Q_\mathrm{S}^2(\varrho _\mathrm{A}(๐ซ)+\varrho _\mathrm{C}(๐ซ))}.$$
(17)
Our analysis rendered that, for a spatially constant screening factor $`\kappa `$, the above model does not yield stable interfacial profiles. This led us to introduce the spatially varying expression in Eq. $`(\text{16})`$. Alternative expressions for inhomogeneous screening factors are known from the theory of electrolytes: In Refs. Onsa34 ; Lee97 non-symmetric screening factors are provided whereas in Ref. Groh98 $`\kappa `$ is calculated from the mean salt density. We prefer the definition in Eqs. $`(\text{16})`$ and $`(\text{17})`$ because it is symmetric and the screening is determined by the salt concentration *at* the actually investigated positions.
The expansion of $`\mathrm{\Omega }[\underset{ยฏ}{\varrho }=\underset{ยฏ}{\varrho }^{(\mathrm{hom})}+\delta \underset{ยฏ}{\varrho }]`$ around a spatially homogeneous state $`\underset{ยฏ}{\varrho }^{(\mathrm{hom})}`$ in powers of perturbations $`\delta \underset{ยฏ}{\varrho }`$ shows that the spatially homogeneous state $`\underset{ยฏ}{\varrho }^{(\mathrm{hom})}`$ is unstable with respect to spatial variations if the macroion charge $`|Q_\mathrm{M}|`$ is sufficiently large because the second order term can become negative. Thus, the choice for $`g^{(\eta )}`$ in Eq. $`(\text{15})`$ leads to spatially *inhomogeneous* bulk phases if the macroion charge $`|Q_\mathrm{M}|`$ is larger than some threshold value. Here, we restrict ourselves to the case of spatially *homogeneous* bulk phases, i.e., only sufficiently small macroion charges are considered.
With Eq. $`(\text{15})`$, the innermost integral in Eq. $`(\text{11})`$ can be evaluated leading to an expression
$`{\displaystyle \underset{0}{\overset{1}{}}}d\eta g_{ij}^{(\eta )}(๐ซ,๐ซ^{})`$
$`=`$ $`\mathrm{exp}\left(U_{ij}^\mathrm{h}(๐ซ,๐ซ^{})\right)\left(1+G_{ij}(\kappa (๐ซ,๐ซ^{};[\underset{ยฏ}{\varrho }]),๐ซ๐ซ^{}_{\mathrm{}})\right)`$
with functions
$$G_{ij}(\kappa ,s):=\underset{0}{\overset{1}{}}d\eta \mathrm{min}[1,U_{ij}^\mathrm{c}(s)\eta \mathrm{exp}\left(\sqrt{\eta }\kappa s\right)]$$
(19)
which decay for $`s\mathrm{}`$ as
$`G_{ij}(\kappa ,s)`$ $``$ $`{\displaystyle \underset{0}{\overset{1}{}}}d\eta U_{ij}^\mathrm{c}(s)\eta \mathrm{exp}\left(\sqrt{\eta }\kappa s\right)`$ (20)
$`=`$ $`{\displaystyle \frac{Q_iQ_j}{s}}{\displaystyle \underset{0}{\overset{1}{}}}d\eta \eta \mathrm{exp}\left(\sqrt{\eta }\kappa s\right)`$
$`=`$ $`{\displaystyle \frac{Q_iQ_j}{s}}{\displaystyle \underset{0}{\overset{1}{}}}d\zeta \mathrm{\hspace{0.33em}2}\zeta ^3\mathrm{exp}\left(\zeta \kappa s\right)`$
$`=`$ $`{\displaystyle \frac{2Q_iQ_j}{\kappa ^4s^5}}\gamma (4,\kappa s)`$
$``$ $`{\displaystyle \frac{12Q_iQ_j}{\kappa ^4s^5}},`$
where $`\gamma `$ denotes the incomplete gamma function Grad80 ; Abra72 . Therefore, the integrand in Eq. $`(\text{11})`$ vanishes at small distances $`๐ซ๐ซ^{}_{\mathrm{}}`$ due to $`U_{ij}^\mathrm{h}`$, whereas it decays as $`๐ซ๐ซ^{}_{\mathrm{}}^1`$ for $`๐ซ๐ซ^{}_{\mathrm{}}\mathrm{}`$. In order to isolate the $`๐ซ๐ซ^{}_{\mathrm{}}^1`$ asymptotics, we add and subtract unity on the right-hand side of Eq. $`(\text{II.3})`$, which, after insertion into Eq. $`(\text{11})`$, leads to the following decomposition of the excess free energy
$$F^{\mathrm{ex}}=F^{\mathrm{ex},\mathrm{h}}+F_{\mathrm{el}}^{\mathrm{ex},\mathrm{c}}+F_{\mathrm{corr}}^{\mathrm{ex},\mathrm{c}}$$
(21)
with the electrostatic part
$`F_{\mathrm{el}}^{\mathrm{ex},\mathrm{c}}[\underset{ยฏ}{\varrho }]`$ $`:=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{ij}{}}{\displaystyle \underset{V}{}}\mathrm{d}^3r{\displaystyle \underset{V}{}}\mathrm{d}^3r^{}\varrho _i(๐ซ)\varrho _j(๐ซ^{})U_{ij}^\mathrm{c}(๐ซ,๐ซ^{})`$ (22)
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{V}{}}\mathrm{d}^3r{\displaystyle \underset{V}{}}\mathrm{d}^3r^{}{\displaystyle \frac{\varrho ^\mathrm{Q}(๐ซ)\varrho ^\mathrm{Q}(๐ซ^{})}{๐ซ๐ซ^{}_{\mathrm{}}}}`$
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{V}{}}\mathrm{d}^3r\varrho ^\mathrm{Q}(๐ซ)\psi (๐ซ)`$
and the correlation part
$`F_{\mathrm{corr}}^{\mathrm{ex},\mathrm{c}}[\underset{ยฏ}{\varrho }]`$ $`:=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{ij}{}}{\displaystyle \underset{V}{}}\mathrm{d}^3r{\displaystyle \underset{V}{}}\mathrm{d}^3r^{}\varrho _i(๐ซ)\varrho _j(๐ซ^{})U_{ij}^\mathrm{c}(๐ซ,๐ซ^{})[\mathrm{exp}(U_{ij}^\mathrm{h}(๐ซ,๐ซ^{}))1`$ (23)
$`+\mathrm{exp}(U_{ij}^\mathrm{h}(๐ซ,๐ซ^{}))G_{ij}(\kappa (๐ซ,๐ซ^{};[\underset{ยฏ}{\varrho }]),๐ซ๐ซ^{}_{\mathrm{}})].`$
Here, the local charge density
$$\varrho ^\mathrm{Q}(๐ซ):=\underset{i}{}Q_i\varrho _i(๐ซ)$$
(24)
and the electrostatic potential
$$\psi (๐ซ):=\underset{V}{}\mathrm{d}^3r^{}\frac{\varrho ^\mathrm{Q}(๐ซ^{})}{๐ซ๐ซ^{}_{\mathrm{}}}$$
(25)
have been introduced. Note that, although the integrands in Eqs. $`(\text{22})`$ and $`(\text{23})`$ are undefined for $`๐ซ=๐ซ^{}`$, the *three-dimensional* integrals exist due to the $`๐ซ๐ซ^{}_{\mathrm{}}^1`$ asymptotics for $`๐ซ๐ซ^{}_{\mathrm{}}0`$. Since the integrand in Eq. $`(\text{23})`$ decays as $`๐ซ๐ซ^{}_{\mathrm{}}^6`$ for $`๐ซ๐ซ^{}_{\mathrm{}}\mathrm{}`$, $`F_{\mathrm{corr}}^{\mathrm{ex},\mathrm{c}}`$ is well-defined for all finite system volumes $`V`$ and the thermodynamic limit of the ratio $`\frac{1}{V}F_{\mathrm{corr}}^{\mathrm{ex},\mathrm{c}}`$ exists. According to the last expression in Eq. $`(\text{22})`$, the same statements are true for the electrostatic contribution $`F_{\mathrm{el}}^{\mathrm{ex},\mathrm{c}}`$ provided the electrostatic potential $`\psi `$ is well-defined. For locally charge neutral systems ($`\varrho ^\mathrm{Q}=0`$), e.g., for bulk phases, the latter holds because of Eq. $`(\text{25})`$. In the next subsection it is shown that $`\psi `$ can also be calculated in systems with only lateral translational symmetry.
### II.4 Planar geometry
By imposing suitable boundary conditions, we consider only systems with translational symmetry in the lateral $`x`$ and $`y`$ directions. Hence, in the absence of spontaneous symmetry breaking, all densities $`\varrho _i`$ depend at most on the $`z`$ coordinate. Since the thermodynamic limit of globally charge neutral systems of Coulomb interacting hard particles exists Lieb72 , i.e., the bulk free energy density depends neither on the shape nor on the boundaries of the system volume $`V`$, the following system volumes of size $`2L\times 2L\times L`$ are considered in the limit $`L\mathrm{}`$:
$$V(L):=A(L)\times [\frac{L}{2},\frac{L}{2}],$$
(26)
where $`A(L)`$ is a square in the $`x`$-$`y`$ plane of side length $`2L`$ with periodic boundary conditions.
The electrostatic potential in Eq. $`(\text{25})`$ can be expressed as
$$\psi (z)=\underset{\frac{L}{2}}{\overset{\frac{L}{2}}{}}dz^{}\varrho ^\mathrm{Q}(z^{})\underset{A(L)}{}\mathrm{d}^2a^{}\frac{1}{(๐^{},zz^{})_{\mathrm{}}},$$
(27)
where $`๐^{}`$ denotes a two-dimensional vector in the $`x`$-$`y`$ plane. The inner integral in Eq. $`(\text{27})`$ leads to
$`{\displaystyle \underset{A(L)}{}}\mathrm{d}^2a^{}{\displaystyle \frac{1}{(๐^{},zz^{})_{\mathrm{}}}}`$ (28)
$`=`$ $`{\displaystyle \underset{0}{\overset{|xx^{}|}{}}}da^{}\mathrm{\hspace{0.33em}8}a^{}{\displaystyle \frac{1}{|xx^{}|}}+{\displaystyle \underset{|xx^{}|}{\overset{L}{}}}da^{}\mathrm{\hspace{0.33em}8}a^{}{\displaystyle \frac{1}{a^{}}}`$
$`=`$ $`4|xx^{}|+8L.`$
In conjunction with the global charge neutrality constraint of Eq. $`(\text{6})`$, Eq. $`(\text{27})`$ reduces to
$$\psi (z)=4\underset{\frac{L}{2}}{\overset{\frac{L}{2}}{}}dz^{}\varrho ^\mathrm{Q}(z^{})|zz^{}|.$$
(29)
By differentiating twice, one finds that $`\psi `$ fulfills the Poisson equation: $`\psi ^{\prime \prime }=8\varrho ^\mathrm{Q}`$. Furthermore, by making use of global charge neutrality (Eq. $`(\text{6})`$), one finds
$`\psi \left({\displaystyle \frac{L}{2}}\right)`$ $`=`$ $`4{\displaystyle \underset{\frac{L}{2}}{\overset{\frac{L}{2}}{}}}dz^{}\varrho ^\mathrm{Q}(z^{})\left(z^{}+{\displaystyle \frac{L}{2}}\right)`$ (30)
$`=`$ $`4{\displaystyle \underset{\frac{L}{2}}{\overset{\frac{L}{2}}{}}}dz^{}\varrho ^\mathrm{Q}(z^{})z^{}`$
$`=`$ $`4{\displaystyle \underset{\frac{L}{2}}{\overset{\frac{L}{2}}{}}}dz^{}\varrho ^\mathrm{Q}(z^{})\left({\displaystyle \frac{L}{2}}z^{}\right)`$
$`=`$ $`\psi \left({\displaystyle \frac{L}{2}}\right).`$
Thus, the density functional in Eq. $`(\text{3})`$ takes the final form
$`\mathrm{\Omega }[\underset{ยฏ}{\varrho }]`$ $`=`$ $`4L^2{\displaystyle \underset{i}{}}{\displaystyle \underset{\frac{L}{2}}{\overset{\frac{L}{2}}{}}}\mathrm{d}z\varrho _i(z)(\mathrm{ln}\left(\varrho _i(z)\right)1\mu _i^{}`$ (31)
$`+{\displaystyle \frac{1}{2}}Q_i\psi (z))+F^{\mathrm{ex},\mathrm{h}}[\underset{ยฏ}{\varrho }]+F^{\mathrm{ex},\mathrm{c}}_{\mathrm{corr}}[\underset{ยฏ}{\varrho }]`$
which has to be minimized under the constraint of global charge neutrality (Eq. $`(\text{6})`$). This is carried out by numerically solving the Euler-Lagrange equations (Eq. $`(\text{7})`$) with a Picard-iteration scheme on a one-dimensional grid.
In the isotropic and nematic bulk fluid, the densities $`\underset{ยฏ}{\varrho }`$ are spatially constant. In this case, the Euler-Lagrange equations Eq. $`(\text{7})`$ comprise five coupled equations:
$$\mathrm{ln}(\varrho _i^{(\mathrm{bulk})})\mu _i^{(\mathrm{bulk})}c_i^{\mathrm{h},(\mathrm{bulk})}c_{\mathrm{corr},i}^{\mathrm{c},(\mathrm{bulk})}=0,$$
(32)
where the chemical potentials $`\mu _i^{(\mathrm{bulk})}`$ fulfill Eq. $`(\text{5})`$ and the local charge neutrality condition ($`_iQ_i\varrho _i^{(\mathrm{bulk})}=0`$). The electrostatic contribution $`c_{\mathrm{el},i}^{\mathrm{c},(\mathrm{bulk})}=Q_i\psi ^{(\mathrm{bulk})}`$ is absent in Eq. $`(\text{32})`$ because the electrostatic potential $`\psi `$ vanishes in locally charge neutral systems (see Eq. $`(\text{25})`$).
For determining the number density profiles at free interfaces between coexisting bulk phases $`_1`$ and $`_2`$, the Euler-Lagrange equations
$$\mathrm{ln}\left(\varrho _i(z)\right)\mu _i^{}+Q_i\psi (z)c_i^\mathrm{h}(z)c_{\mathrm{corr},i}^\mathrm{c}(z)=0$$
(33)
are to be solved with the boundary conditions
$$\varrho _i\left(z=\frac{L}{2}\right)=\varrho _i^{(_1)},\varrho _i\left(z=\frac{L}{2}\right)=\varrho _i^{(_2)}.$$
(34)
In order that for $`z=\frac{L}{2}`$ and $`z=\frac{L}{2}`$ Eq. $`(\text{33})`$ reduces to Eq. $`(\text{32})`$ for $`_1`$ and $`_2`$, respectively, one has the requirements
$`\mu _i^{}+Q_i\psi \left({\displaystyle \frac{L}{2}}\right)`$ $`=`$ $`\mu _i^{(_1)},`$
$`\mu _i^{}+Q_i\psi \left({\displaystyle \frac{L}{2}}\right)`$ $`=`$ $`\mu _i^{(_2)}.`$ (35)
Using Eq. $`(\text{30})`$, one readily concludes
$$\mu _i^{}=\frac{1}{2}\left(\mu _i^{(_1)}+\mu _i^{(_2)}\right)$$
(36)
and
$$\psi _D:=\psi \left(\frac{L}{2}\right)\psi \left(\frac{L}{2}\right)=\frac{1}{Q_i}\left(\mu _i^{(_1)}\mu _i^{(_2)}\right).$$
(37)
$`\psi _D`$ is known as *Donnan potential* Adam73 between the two bulk phases $`_1`$ and $`_2`$. It maintains a density gradient of the mobile particles at the interface between two coexisting bulk phases. Its definition given above is unique, i.e., the rightmost expression is in fact independent of $`i`$ due to Eq. $`(\text{5})`$, e.g.,
$`{\displaystyle \frac{\mu _\mathrm{A}^{(_1)}\mu _\mathrm{A}^{(_2)}}{Q_\mathrm{A}}}`$ $`=`$ $`{\displaystyle \frac{\left(\mu _{\mathrm{CA}}^{}\mu _\mathrm{C}^{(_1)}\right)\left(\mu _{\mathrm{CA}}^{}\mu _\mathrm{C}^{(_2)}\right)}{Q_\mathrm{C}}}`$ (38)
$`=`$ $`{\displaystyle \frac{\mu _\mathrm{C}^{(_1)}\mu _\mathrm{C}^{(_2)}}{Q_\mathrm{C}}}.`$
## III Bulk fluid
As a first step in the investigation of the density functional developed in the last section, bulk phase diagrams are determined for various macroion charges $`Q_\mathrm{M}`$ by solving the bulk Euler-Lagrange equations $`(\text{32})`$.
The macroion and the salt number densities are given by $`\varrho _\mathrm{M}:=\varrho _{\mathrm{M}_\mathrm{x}}+\varrho _{\mathrm{M}_\mathrm{y}}+\varrho _{\mathrm{M}_\mathrm{z}}`$ and $`\varrho _\mathrm{S}:=\varrho _\mathrm{A}`$, respectively. In order to detect the formation of liquid crystalline phases of the macroions, the equilibrium nematic order parameter for the director oriented relative to the $`z`$-direction,
$$s_\mathrm{M}:=\frac{3}{2}\frac{\varrho _{\mathrm{M}_\mathrm{z}}}{\varrho _\mathrm{M}}\frac{1}{2}[\frac{1}{2},1],$$
(39)
and the equilibrium biaxial order parameter,
$$q_\mathrm{M}:=\frac{\varrho _{\mathrm{M}_\mathrm{x}}\varrho _{\mathrm{M}_\mathrm{y}}}{\varrho _\mathrm{M}},$$
(40)
have been determined. The definition of $`s_\mathrm{M}`$ agrees with the well-known scalar liquid-crystal order parameter $`S=P_2(\mathrm{cos}\vartheta )=\frac{3}{2}(\mathrm{cos}\vartheta )^2\frac{1}{2}`$ because within the Zwanzig model only macroion orientations $`\mathrm{M}_\mathrm{z}`$ parallel ($`\mathrm{cos}\vartheta =1`$) and $`\mathrm{M}_{\mathrm{x},\mathrm{y}}`$ perpendicular ($`\mathrm{cos}\vartheta =0`$) to the $`z`$-axis are possible. $`s_\mathrm{M}`$ vanishes in an isotropic phase ($`\varrho _{\mathrm{M}_\mathrm{x}}=\varrho _{\mathrm{M}_\mathrm{y}}=\varrho _{\mathrm{M}_\mathrm{z}}`$), whereas it is positive in a nematic phase with director parallel to the $`z`$-axis ($`\varrho _{\mathrm{M}_\mathrm{z}}>\varrho _{\mathrm{M}_\mathrm{x}},\varrho _{\mathrm{M}_\mathrm{y}}`$). A discrimination of the orientation $`\mathrm{M}_\mathrm{z}`$ leads to negative values of $`s_\mathrm{M}`$.
It turned out that the biaxial order parameter $`q_\mathrm{M}`$ vanishes throughout the whole inspected range of reduced chemical potentials $`\mu _{\mathrm{C}_k\mathrm{M}}^{}`$ and $`\mu _{\mathrm{CA}}^{}`$, whereas the nematic order parameter $`s_\mathrm{M}`$ indicates either an isotropic fluid ($`s_\mathrm{M}=0`$) or a nematic fluid ($`s_\mathrm{M}>0`$).
Figure 2
shows phase diagrams for the parameters Units (compare Fig. 1) $`D_\mathrm{M}=20\mathrm{}_B14\mathrm{nm}`$, $`L_\mathrm{M}=\mathrm{}_B0.72\mathrm{nm}`$, $`D_\mathrm{S}=\mathrm{}_B0.72\mathrm{nm}`$, $`Q_\mathrm{S}=e`$ with $`Q_\mathrm{M}=0`$, $`Q_\mathrm{M}=0.25Q_\mathrm{S}`$, $`Q_\mathrm{M}=0.5Q_\mathrm{S}`$, $`Q_\mathrm{M}=0.75Q_\mathrm{S}`$, and $`Q_\mathrm{M}=Q_\mathrm{S}`$ in terms of the macroion packing fraction $`\eta _\mathrm{M}=\varrho _\mathrm{M}D_\mathrm{M}^2L_\mathrm{M}`$ and the salt density $`\varrho _\mathrm{S}`$.
One isotropic phase ($`I`$) and one nematic phase ($`N`$) are found separated by first-order phase transitions. Whereas for coexisting phases $`\eta _\mathrm{M}`$ is always smaller in the isotropic phase than in the nematic phase, $`\varrho _\mathrm{S}`$ of coexisting phases is higher in the isotropic and lower in the nematic phase. A similar displacement of salt from regions of large concentrations of charged macroions is known as Donnan effect Adam73 . Whereas the original Donnan effect has been discovered in systems subdivided by membranes which are impermeable for macroions, here the density difference of the macroions occurs due to two coexisting bulk phases. As for the case of membrane equilibrium, here a Donnan potential $`\psi _D`$ (Eq. $`(\text{37})`$) maintains the density gradients between the coexisting phases. Figure 3
shows its dependence on the macroion charge $`Q_\mathrm{M}`$ and the salt density $`\varrho _\mathrm{S}`$. $`\psi _D`$ decreases with increasing salt density $`\varrho _\mathrm{S}`$. This tendency is intuitively expected as the Donnan effect becomes more pronounced with increasing macroion charge whereas increasing the salt density gives rise to a stronger screening of the macroion charge. For fixed salt density $`\varrho _\mathrm{S}`$ well below $`0.1\mathrm{mM}`$, $`\psi _D`$ decreases with increasing $`|Q_\mathrm{M}|`$, whereas this behavior is reversed for fixed salt density $`\varrho _\mathrm{S}`$ above $`0.1\mathrm{mM}`$.
Upon an increase of the macroion charge $`|Q_\mathrm{M}|`$, the isotropic and the nematic binodals are shifted to larger values of the macroion packing fraction $`\eta _\mathrm{M}`$. This may be qualitatively understood by introducing the notion of an effective shape, which, in the present case, for macroions is given by a hard core surrounded by a soft $`_{\mathrm{}}`$-sphere, i.e., a cube, with its linear extension proportional to $`Q_\mathrm{M}^2`$ due to the pairwise Coulomb repulsion. For small macroion charges, the effective shape is still platelike whereas for highly charged colloids, the effective shape tends towards a cube leading to a shift of the two-phase region to larger macroion packing fractions.
For fixed macroion charge as well as particle shape and increasing salt density $`\varrho _\mathrm{S}`$, the isotropic-nematic binodals in Fig. 2 bend towards smaller macroion packing fractions. This behavior is expected intuitively, because high ionic strength causes strong screening which in turn leads to effectively quasi-hard platelets (see Fig. 2(a)).
## IV Free interfaces
Based on the bulk properties provided in the previous section, we are now able to calculate the density profiles at the free interfaces between the coexisting isotropic and nematic phases by solving the spatially varying Euler-Lagrange equations (Eq. $`(\text{33})`$). The density and order parameter profiles corresponding to the parameters used in the previous section (see also Figs. 2 and 3) and to nematic bulk salt density $`\varrho _\mathrm{S}^{(N)}=2.210^5\widehat{=}0.1\mathrm{mM}`$ Units are depicted in Figs. 4
and 5,
respectively. The interface position $`z=0`$ is chosen such that $`\varrho _\mathrm{M}(0)=\frac{1}{2}\left(\varrho _\mathrm{M}^{(\mathrm{I})}+\varrho _\mathrm{M}^{(\mathrm{N})}\right)`$.
First, and most important, the formalism described in Sec. II renders stable free interfaces between coexisting bulk phases. This should be regarded as an accomplishment which can be traced back to using the spatially varying screening factor $`\kappa `$ introduced in Eq. $`(\text{16})`$; trials with spatially constant $`\kappa `$ were not successful.
For a given density profile $`\varrho _\mathrm{M}(z)`$ (Fig. 4), the corresponding interface width $`\zeta `$ is defined as the spatial distance between the loci, where the tangent at the density profile at position $`z=0`$ reaches the values of the nematic bulk density $`\varrho _\mathrm{M}^{(\mathrm{N})}`$ and the isotropic bulk density $`\varrho _\mathrm{M}^{(\mathrm{I})}`$, respectively. This interface width $`\zeta `$ decreases monotonically with increasing macroion charge from $`\zeta =1.8D_\mathrm{M}`$ for $`Q_\mathrm{M}=0.25Q_\mathrm{S}`$ to $`\zeta =1.3D_\mathrm{M}`$ for $`Q_\mathrm{M}=Q_\mathrm{S}`$ (see Fig. 4). The Debye length $`\kappa ^1`$ decreases monotonically from $`\kappa ^1=1.5D_\mathrm{M}`$ for $`Q_\mathrm{M}=0.25Q_\mathrm{S}`$ to $`\kappa ^1=D_\mathrm{M}`$ for $`Q_\mathrm{M}=Q_\mathrm{S}`$. Finally, the bulk correlation lengths $`\xi `$ of the coexisting isotropic and nematic bulk phases, inferred from the exponential decay lengths of $`\varrho _\mathrm{M}(z)\varrho _\mathrm{M}^{(\mathrm{I},\mathrm{N})}`$, also decrease monotonically upon increasing $`|Q_\mathrm{M}|`$ and the values are by and large equal to those of $`\kappa ^1`$.
The nematic order parameter profiles $`s_\mathrm{M}(z)`$ (Fig. 5) interpolate almost monotonically between $`s_\mathrm{M}(\mathrm{})>0`$ in the nematic bulk phase ($`N`$) and $`s_\mathrm{M}(\mathrm{})=0`$ in the isotropic bulk phase ($`I`$). Note that $`s_\mathrm{M}`$ has been defined for a director in $`z`$-direction, i.e., platelets on the nematic side ($`z<0`$) are preferably oriented parallel to the free interface. At a fixed position on the nematic side ($`z<0`$), $`s_\mathrm{M}(z)`$ decreases with increasing macroion charge $`|Q_\mathrm{M}|`$. This behavior is consistent with the picture of an increasingly isotropic effective shape, introduced in the previous section.
The charge density profiles $`\varrho ^Q(z)`$ (Eq. $`(\text{24})`$) displayed in Fig. 6
show deviations from local charge neutrality within the interfacial region $`4D_\mathrm{M}z4D_\mathrm{M}`$. A negative charge density occurs on the nematic side ($`N`$) and a positive charge density on the isotropic side ($`I`$). Such a local charging is necessary for the appearance of the non-vanishing Donnan potentials $`\psi _D`$ shown in Fig. 3. The full electrostatic potential profiles $`\psi (z)`$ are depicted in Fig. 7.
They increase monotonically from the macroion-rich nematic phase $`N`$ to the macroion-poor isotropic phase $`I`$, maintaining the density gradients occurring in the interface region. The potential difference $`\psi (\mathrm{})\psi (\mathrm{})`$ equals the Donnan potential $`\psi _D`$ (see Eq. $`(\text{37})`$ and Fig. 3).
The interfacial tensions $`\gamma `$ of the interfaces shown in Fig. 4 decrease monotonically from $`\gamma =1.3310^5\widehat{=}108\mathrm{nN}\mathrm{m}^1`$ for $`Q_\mathrm{M}=0.25Q_\mathrm{S}`$ to $`\gamma =610^7\widehat{=}5\mathrm{nN}\mathrm{m}^1`$ for $`Q_\mathrm{M}=Q_\mathrm{S}`$ Units , which are comparable to experimental findings for laponite suspensions Kooi01 . The corresponding wetting parameters $`\omega =(4\pi \gamma \xi ^2)^1`$ Schi90 are in the range $`6\mathrm{}340`$. If these values for $`\gamma `$ are indeed so small, the free isotropic-nematic interfaces are expected to be strongly affected by capillary wavelike fluctuations which are not captured by the present theory.
Calculating bulk phase diagrams and density profiles for macroion charges $`Q_\mathrm{M}1.25Q_\mathrm{M}`$ along the lines described above lead to unphysical results like, e.g., negative interfacial tensions. The reason for this phenomenon is that the bulk phases determined in Sec. III are assumed to be spatially homogeneous, whereas it can be shown that the equilibrium bulk states are spatially inhomogeneous for sufficiently large macroion charges (see Subsec. II.3).
## V Discussion and Summary
The numerical results for the bulk systems and the free interfaces presented in Secs. III and IV, respectively, are in good qualitative agreement with intuitive expectations. Within the density functional theory described in Sec. II one gains access not only to interfacial density profiles at free interfaces between coexisting bulk phases but also to local charge densities and electrostatic potential profiles.
Since here we have been interested in spatially homogeneous (isotropic or nematic) bulk phases, only very small macroion charges ($`|Q_\mathrm{M}|e`$) have been considered. Determining phase coexistence for larger platelet charges is computationally more demanding because spatially inhomogeneous bulk phases are involved. Note that, according to Subsec. II.3, charges have to be interpreted as effective charges.
A further difficulty related to the choice of the pair distribution functions $`g_{ij}^{(\eta )}`$ used in Eq. $`(\text{11})`$ may appear if, as for Eq. $`(\text{15})`$, the effective shape of the macroions becomes more and more isotropic upon increasing the macroion charge: In this case, the two-phase region between isotropic and anisotropic phase is shifted to unrealistically large packing fractions. Therefore, the pair distribution functions must be chosen properly in order to yield *anisotropic* effective macroion shapes up to large macroion charges. Unfortunately, deriving *analytical* expressions for pair distribution functions of platelike particles with inhomogeneous charge distributions is still a big challenge.
In conclusion, the density-functional theory of charged platelike particles developed here can be regarded as a first step to understand qualitatively free interfaces between isotropic and nematic bulk phases of suspensions of charged platelets and salt. Here, we have constructed a density functional for charged platelike particles and applied it to a ternary mixture of platelike macroions and salt ions (Fig. 1) in the bulk and at free interfaces between coexisting isotropic and nematic phases. For sufficiently small macroion charges, the bulk phase diagrams in terms of densities (Fig. 2) exhibit one isotropic phase and one nematic phase. For increasing macroion charge and fixed salt densities, the two-phase coexistence region is shifted to larger macroion packing fractions. For fixed macroion charge and increasing salt density, the limit of quasi-hard platelets is approached. The Donnan potential between coexisting phases (Fig. 3) can be expressed in terms of the particle chemical potentials gained from bulk structure calculations. Density and nematic order parameter profiles at free interfaces between isotropic and nematic phases at coexistence (Figs. 4 and 5) show non-monotonic behavior. The value of the nematic order parameter in the nematic bulk phase decreases upon increasing the macroion charge. The width of the interface and the bulk correlation lengths are approximately given by the Debye length. The interfacial tension decreases upon increasing the macroion charges. Electrically charged layers form at the free interface (Fig. 6). The corresponding electrostatic potential profiles (Fig. 7) exhibit monotonic behavior. Investigations of spatially inhomogeneous bulk phases are necessary in order to apply the theory to larger macroion charges. Improvements of the present theory call for more accurate analytical expressions for the pair distribution function between charged platelike particles.
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# The spectral-type/luminosity and the spectral-type/satellite-density relations in the 2dFGRS
## 1 Introduction
The existence of a relation between galaxy type and environment dates back to pioneering work by Hubble& Humason (1931) and was first quantified by Davis & Geller (1976) and Dressler (1980). Dressler found that, in nearby clusters, the fraction of elliptical galaxies increases and the fraction of spirals decreases with increasing density. Whitmore & Gilmore (1991) and Whitmore et al. (1993) suggested that the correlation between morphology and cluster-centric-radius is tighter than the correlation between morphology and density. Deep surveys such as the 2dFGRS and the SDSS have confirmed that the star-formation level in galaxies decreases at large galaxy density (Lewis et al. (2002); Gomez et al. (2003); Goto et al. (2002, 2003); Balogh et al. 2004a ; Christlein & Zabludoff (2005)) and that a threshold is reached at low densities ($``$ 1 gal Mpc<sup>-2</sup>, M<sub>B</sub> $``$ โ19), below which no further increase in star-formation is observed. These results provide evidence that, at least at low redshift, a strong correlation exists between the characteristics of the stellar component of a galaxy and its surrounding environment, with luminosity and color being the galaxy properties most strongly correlated with environment (Kauffmann et al. (2004); Blanton et al. 2005a ). But the exact dependence of the morphology/density relation on density and on luminosity is still a matter of debate. It is not clear yet whether the relation extends to galaxies in systems less dense than clusters (i.e. to the majority of galaxies). Postman & Geller (1984) and Maia & da Costa (1990) claim that the morphology/density relation certainly extends to groups, whereas Whitmore (1995) does not confirm their result. Tran et al. (2001) and Helsdon & Ponman (2003) support the existence of the morphology/density relation in X-ray bright groups. Dominguez et al. (2002) show that the relation is observed in very massive (optically selected) groups only, whereas Kelm & Focardi (2004a) report that the frequency of early-type galaxies is larger in Compact Groups than among isolated galaxies. Mateus & Sodrรฉ (2004) and Gerken et al. (2004) provide evidence that, even outside clusters, star formation properties are affected in all ranges of density. Tanaka et al. (2004) find that only faint galaxies show a break in star formation and morphology at a critical local density. Also the role of luminosity on the morphology/density relation has still to be fully disentangled. Galaxy clustering appears to depend on luminosity for luminous galaxies and on color for low luminous ones (Norberg et al. 2002a ; Hogg et al. (2003); Balogh et al. 2004b ; Berlind et al. (2004); Zehavi et al. (2004)). For blue galaxies the relation between environment and luminosity is typically weak, whereas for red galaxies clustering is likely a non-monotonic function of luminosity, peaking at both high and low luminosities.
In this paper we investigate the relation linking the galaxy spectral type mix with both luminosity and local environment. Specifically we explore 1) how the spectral-type/luminosity relation varies as a function of environment and 2) how the spectral-type/satellite-density relation varies as a function of luminosity. We use data from the 2dF to select 10 different volume limited samples, covering a wide luminosity range (โ22.5 $``$ M<sub>B</sub> -5 $`\mathrm{log}`$ $`h_{75}`$ $``$ โ17.0). We evaluate for each galaxy the neighbour density on the characteristic scale of galaxy groups ($``$ 1 $`h_{75}^1`$ Mpc) which further corresponds to the present day typical virial radius of halos. Galaxy properties are expected to correlate most strongly with densities evaluated on this scale, also from a theoretical standpoint (Blanton et al. 2005b ; Kauffmann et al. (2004); Berlind et al. (2004)).
At variance with previous analysis we compute neighbour density applying a maximum magnitude difference criterion and count neighbours over a 2 magnitude interval. The adopted range in magnitude reduces the number of galaxies that have no neighbours, on a group scale, to $``$15% and associates most galaxies (2/3) with a number of neighbours (1$``$neigh$``$8) that matches the typical observed environment of galaxies in groups. We also limit neighbour computation to galaxies that are fainter than (or equally luminous to) the galaxy itself. Usually, when computing density in volume-limited samples, no distinction is made between brighter and fainter neighbours. This implies that, within the same volume-limited sample, the density definition depends on luminosity: the environment of luminous galaxies is defined by fainter neighbours whereas the environment of low luminous galaxies is defined by brighter neighbours. But, obviously, the impact of a brighter or a fainter companion on a galaxy is different. Less massive companions have likely been, or will be, accreted by the galaxy halo, whereas more massive companions will likely accrete the galaxy and destroy its halo.
Our density definition is luminosity-independent. For luminous galaxies, the environment on the 1 h$`{}_{}{}^{1}{}_{75}{}^{}`$ Mpc scale is likely to correspond to the density of satellites that have been captured by the galaxy halo. Conversely, for low luminous galaxies, it likely corresponds to the clustering of small halos among themselves, or, in the case of a galaxy swallowed up in the halo of a bigger companion, for the richness of satellites within this large halo.
In $`\mathrm{\S }`$2 we present the sample, in $`\mathrm{\S }`$3 we discuss the link between density distribution and luminosity. In $`\mathrm{\S }`$4 and $`\mathrm{\S }`$5 we explore the dependence on density of the spectral-type/luminosity relation. In $`\mathrm{\S }`$6 we investigate the dependence on luminosity of the spectral-type/density relation. The summary and conclusions are given in $`\mathrm{\S }`$7. We assume $`\mathrm{\Omega }_M=0.3`$, $`\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$, and $`h_{75}`$ = H<sub>o</sub> /(75 km s<sup>-1</sup> Mpc<sup>-1</sup>) = 1.
## 2 The sample
The sample we use for the present analysis is selected from the 2dFGRS (Colless et al. (2001, 2003)). The 2dF covers $``$ 1800 square degrees and is complete for galaxies down to an extinction-corrected limit of $`b_J`$ = 19.45. It provides redshifts, in the range 0 $``$ $`z`$ $``$ 0.3, for 221,496 galaxies selected from the APM catalogue (Norberg et al. 2002b ), which is 90-95% complete (Maddox et al. (1990)). Because saturation effects and stellar contamination cannot be ignored for bright galaxies, we exclude from the sample galaxies brighter than b<sub>j</sub> = 16.
Each 2dF galaxy spectrum is typed on the basis of the relative strength of its first two principal components (for details on the PCA see Folkes et al. 1999), which are the emission and absorption components within the spectrum. The parameter $`\eta `$ (Madgwick et al. (2002)) is the linear combination of these two components. Qualitatively $`\eta `$ is an indicator of the ratio of the present to the past star-formation activity of each galaxy, but it is reliable only for z $``$ 0.15. Clusters are dominated by galaxies with the lowest $`\eta `$ values, whereas the field contains a much larger proportion of galaxies with higher ($`\eta >`$0) values. The median $`\eta `$ correlates with morphological classes, (low $`\eta `$ are typically early type galaxies, high $`\eta `$ late type galaxies) although there is a large scatter in the $`\eta `$ values of spectra that lie within a given morphological class.
As in Madgwick et al. (2002) we divide the $`\eta `$ scale into 4 intervals:
$$\eta <1.4Type\mathrm{\hspace{0.17em}1}$$
$$1.4\eta <1.1Type\mathrm{\hspace{0.17em}2}$$
$$1.1\eta <3.5Type\mathrm{\hspace{0.17em}3}$$
$$\eta 3.5Type\mathrm{\hspace{0.17em}4}$$
We group together Type 3 and Type 4 galaxies, the latter being rare. Throughout the paper, Type 1, Type 2 and Type 3 + Type 4 galaxies are named passive, quiet-SF and active-SF respectively. We keep quiet-SF and active-SF galaxies separate, in order to investigate any dependence of SF triggering processes on specific density characteristics.
For each 2dF galaxy (random fields excluded) with z $``$ 0.15 we have automatically identified neighbours within 1 $`h_{75}^1`$ Mpc projected distance and $`\pm `$ 1000 km s<sup>-1</sup> depth. We count as neighbours all galaxies fainter than the galaxy itself that satisfy a maximum magnitude difference criterion (โ2 $``$ $`M_{gal}M_{neigh}`$ $``$0). We reject from the sample all galaxies whose 2 magnitude fainter companions would fail the 2dF selection criteria. The minimum fiber separation of the 2dF survey ($``$ 30 โณ) tends to reduce the number of close neighbours of galaxies, a bias that might affect passive galaxies more severely than SF galaxies. However, the bias is likely marginal as this separation corresponds to less than one-tenth of the explored distance even for the highest redshift galaxies.
The final sample includes $``$14,000 galaxies in the redshift range 0.0156 $``$ z $``$ 0.15, and absolute magnitude range โ22.5 $``$ M<sub>B</sub> -5 $`\mathrm{log}`$ $`h_{75}`$ $``$ โ17.0. Absolute magnitudes are computed adopting the k-correction as in Magdwick et al. (2002), which varies with galaxy spectral-type. We split the sample into 10 different volume-limited subsamples, covering a 1 magnitude range each and overlapping by 0.5 magnitude. Since the k-corrections are class dependent, the z<sub>min</sub> and z<sub>max</sub> values corresponding to a given absolute magnitude range are also class dependent. Hence, the volumes defining the samples for two different spectral classes, for the same bin in absolute magnitude, will not exactly coincide (see also Norberg et al. 2002a.).
Table 1 lists the spectral-type composition in each volume limited sample. The gap between galaxies in the faintest and the brightest volume limited samples corresponds to a factor $``$150 in luminosity.
## 3 The dependence on luminosity of the satellite-density distribution
We assign a local density (number of fainter neighbours within 1 $`h_{75}^1`$ Mpc, $`\pm `$ 1000 km<sup>-1</sup> depth and a 2-$`\mathrm{\Delta }`$Mag range) to all galaxies in our sample. We then define four distinct density regimes characterized by different number of neighbours:
$$neigh=0$$
$$1neigh2$$
$$3neigh8$$
$$neigh>8$$
Table 2 lists the number of galaxies per spectral-type in the four defined density ranges and for each volume limited sample. The density parameterization is such that most galaxies ($``$ 2/3) are in the two central bins, which exhibit the typical galaxy density of groups. Obviously, computing densities on the galaxy group scale does not correspond to selecting a sample of 2dF groups (Eke et al. 2004a ; Eke et al. 2004b ; Merchรกn & Zandivarez (2002)).
The relation linking luminosity and fainter neighbour density is shown in Fig. 1. Distributions of passive, quiet-SF, active-SF and all-type galaxies, normalized to the total number of galaxies of a given type, are shown, for the 10 volume-limited samples. The last column of Fig.1 shows that in a composite (all-type) population, the number of fainter neighbours associated with galaxies is a weak function of luminosity. This recalls the result by Zehavi et al. (2002), showing that all-type galaxy subsamples in 3 distinct absolute magnitude ranges have real-space correlation functions that are parallel power-laws.
If we assume that for luminous galaxies, the environment on a 1 $`h_{75}^1`$ Mpc scale essentially stands for the density of satellites that have been captured by the galaxy halo whereas for faint galaxies, it stands for the number of neighbour galaxies still in their own small halo, then Fig. 1 actually indicates that the distribution of satellites surrounding luminous central galaxies within large halos, and the distribution of fainter companions surrounding low luminous galaxies, are almost self similar. This implies that without information on the luminosity of the galaxies, the neighbour density distribution of galaxies (on 1 $`h_{75}^1`$ Mpc scale) cannot be used to discriminate between massive group-size halos and associations of galaxies in distinct small-size halos.
Figure 1 also clearly shows that at all luminosities, passive and SF galaxies exhibit different distributions and that the excess of companions surrounding passive galaxies is not limited to luminous galaxies (Norberg et al. 2002a ; Hogg et al. (2003); Berlind et al. (2004); Blanton et al. 2005a ) but is instead a general characteristic of passive galaxies. The all-type galaxy distribution reflects the passive population at the bright end, and the star-forming galaxy population at the faint end. Our assumption that luminous galaxies are central galaxies within group-size halos and faint galaxies are central galaxies within small size halos is therefore further consistent with the expectation that the SFR of a galaxy is a decreasing function of its halo mass .
However SF galaxies are found among bright galaxies and passive galaxies among faint ones. How can we explain their existence? We will assume that luminous SF galaxies are hosted in small mass halos; they may exhibit several neighbours, but, at variance with luminous passive galaxies, neighbours are not embedded within the galaxy halo. As a consequence optically selected passive dominated groups are predicted to be systematically more massive than optically selected SF dominated groups (Kelm & Focardi 2004a ; Kelm & Focardi 2004b ; Mulchaey et al. (2003)). Similarly, we explain the existence of low luminous passive galaxies assuming that they are satellites embedded within the halo of a large (group-size) system. Actually, the large fraction of faint passive galaxies with $`>`$8 neighbours among faint galaxies (see Table 2) suggests that this population is tracing a large potential well (Norberg et al. 2002a ; Hogg et al. (2003); Berlind et al. (2004); Zehavi et al. (2004); Jing & Borner (2004)), with the galaxy and all of its fainter neighbours having been accreted by a massive system.
The stronger clustering of passive galaxies relative to SF galaxies, on the group scale, appears to arise from two distinct contributions. At the luminous end, it is due to an excess of satellites surrounding central galaxies inside large halos. At the faint end it is due to an excess of satellites that are strongly correlated among themselves.
## 4 The dependence on satellite-density of the spectral-type/luminosity relation
Figure 2 shows the fractions of passive, quiet-SF and active-SF galaxies in the 10 volume-limited samples for the total (all neighbour density) population. It is a global spectral-type/luminosity plot that indicates how the fraction of passive and SF galaxies varies as a function of luminosity. It is similar to Fig. 9 in Norberg et al. (2002a), in which fractions have been derived for late-type and early-type galaxies only.
The horizontal line in Fig. 2 denotes the 40% fraction: points above this threshold mark dominant populations. Passive galaxies are โdominantโ in galaxy samples brighter than M<sub>B</sub> \- 5 $`\mathrm{log}`$ $`h_{75}`$ = -20, active-SF galaxies are โdominantโ in samples fainter than M<sub>B</sub> \- 5 $`\mathrm{log}`$ $`h_{75}`$ = โ19. Contributions from different spectral-type populations are comparable in the \[$`20รท19`$\] magnitude bin. These trends confirm that star-formation activity in the local universe definitely is a characteristic of low luminosity galaxies.
In Fig. 3 we break down the contributions of passive and SF galaxies to the Fig. 2 plot into their contributions from systems exhibiting different numbers of fainter neighbours. This allows us to explore how strongly the relative fraction of passive and SF galaxies depends on environment. If the dominance of passive galaxies at high luminosity and the dominance of active-SF galaxies at low luminosity were independent of neighbour density, we would expect all panels in Fig. 3 to be similar. This is not the case, however, differences among panels are modest: whatever the number of satellites, bright samples are dominated by passive galaxies, and faint samples by active-SF ones.
Therefore, in general, luminosity dominates over neighbour multiplicity in setting the spectral-type mix of a galaxy population. A luminous galaxy might have few or many satellites, but will likely trace a deep potential. A faint galaxy might have few or many neighbours but will likely trace a shallow potential. It is only in extreme environments that the mix set by luminosity is significantly modified: galaxies with neigh$`>`$8 have their star-formation level (a typical active-SF one) suppressed even in faint samples, whereas isolated (neigh=0) galaxies are still 40% likely to be star-forming (quiet-SF) at M<sub>B</sub> โ 5 $`\mathrm{log}`$ $`h_{75}`$$``$$`21.5`$. These trends are consistent with result discussed in $`\mathrm{\S }`$3, namely that โminorityโ population can be identified in very luminous and very faint samples whose luminosities are inaccurate tracers of their halo mass: luminous SF galaxies trace small (sub-group size) halos, whereas faint passive galaxies trace massive (group/cluster size) halos.
Figures 2 and 3 also provide evidence that fractions of active-SF and quiet-SF galaxies exhibit distinct trends with luminosity: the fraction of active-SF galaxies decreases towards increasing luminosity, while the fraction of quiet-SF galaxies is nearly independent of luminosity, except for the most luminous samples. While Fig.1 indicates that the dependence on density is the same for quiet-SF and active-SF galaxies (Madgwick et al. (2003)), Fig. 2 and Fig. 3 indicate that the dependence on luminosity is different. The data thus suggest a bimodal behaviour for galaxies with satellite-density and a โtrimodalโ behaviour (passive, quiet-SF, and active-SF) with luminosity. Bimodality in the distribution of galaxies properties has been addressed in many recent papers (Strateva et al. (2001); Hogg et al. (2002); Balogh et al. 2004b ; Berlind et al. (2004); Blanton et al. 2005a ).
## 5 The role of satelite density for passive and SF galaxies
To further explore the trend of increasing passive and decreasing SF galaxy fraction with luminosity and fainter neighbour density we also show, in Fig. 4, the fractional content of passive, quiet-SF and active-SF galaxies in different environments. The relative role of extremely dense and intermediate dense environment can be easily explored by comparing the gap between the neigh=0 and the neigh$`>`$8 lines with the the gap between the neigh=0 and the neigh=3-8 lines. Figure 4 shows that, for all 3 types, the gaps undergo a strong variation at magnitude M<sub>B</sub> \- 5 $`\mathrm{log}`$ $`h_{75}`$ $``$ โ19. Therefore, we will keep the analysis of faint and bright samples separate.
Figure 4 shows that in luminous samples the fraction of isolated passive galaxies (neigh=0) is always below the fraction of passive galaxies with neigh$`>`$8, and that the gap does not depend on luminosity. A similar specular large gap is observed for quiet-SF galaxies, whereas the gap is smaller for active-SF galaxies. A smaller, but still significant, gap is observed between the neigh=0 and the neigh=3-8 lines, in passive and quiet-SF galaxy samples. The size of the gap is, again, nearly luminosity independent. Conversely, no gap is associated with active-SF galaxies.
In faint samples (M -5 $`\mathrm{log}`$ $`h_{75}`$$``$ โ19), passive and active-SF galaxies exhibit a large gap between the neigh=0 and the neigh$`>`$8 lines, whereas no gap is associated with the neigh=0 and neigh=3-8 transition. Again, a distinct behaviour is observed for active-SF and quiet-SF galaxies, the latter being equally frequent in all environments.
In summary, Fig. 4 indicates that a continuous parameterization of neighbour multiplicity (from 0 to 1-2 to 3-8 to $`>`$8) is indeed meaningful for bright samples, where neighbours are mainly satellites, as it relates to different fractions of passive and quiet-SF galaxies. In faint samples, however, a threshold-like density parameterization appears to describe the galaxy behaviour better than a continuous one. This suggests that a continuous relation linking spectral-type with density only occurs when computing the density of satellites surrounding very luminous galaxies and implies that the spectral-type/density relation traces an enhanced correlation inside massive halos rather than an enhanced correlation between distinct halos.
The M<sub>B</sub> \- 5 $`\mathrm{log}`$ $`h_{75}`$ $`19`$ magnitude is a critical one: it corresponds to the luminosity at which the dependence on satellite-density moves from continuous to threshold-like, and also to the luminosity where fractions of active-SF galaxies become larger than the fractions of quiet-SF galaxies (see Fig. 2). Therefore it corresponds to the luminosity above which samples of passive galaxies exhibit a dependence on fainter neighbour density that is specular relative to that of quiet-SF galaxies, and below which passive galaxies are specular to active-SF galaxies. This is consistent with the finding (Norberg et al. 2002a ; Balogh et al. 2004a ) that, for low luminosity galaxies, clustering is a strong function of color, while for luminous galaxies clustering is a strong function of luminosity.
## 6 The dependence on luminosity of the spectral-type/neighbour density relation
We have shown that the spectral-type/density relation is possibly a spectral-type/satellite-density relation that traces an enhanced correlation inside single massive halos rather than enhanced correlation between distinct halos. To test this assumption directly we next examine the dependence of the spectral-type/density ($``$morphology/density) relation on luminosity. In Fig. 5 we show the spectral-type/density relation, with fractions of passive, quiet-SF, active-SF and all-SF galaxies normalized to the total number of galaxies in a given density bin, for 9 volume-limited samples. The faintest sample is not shown because it is small (see Table 2) and more affected by statistical uncertainties.
Figures 5 provides evidence that the fractional increase of passive galaxies between extreme densities is a universal characteristic of galaxies that does not depend on luminosity. What depends on luminosity is the fractional increase (decrease) of passive (SF) galaxies between neigh=0 and a group-like density (neigh=3-8). The increase (decrease) is not observed in samples fainter than $``$ โ19. This supports our claim that the spectral-type/density relation is actually a process linked to the accretion of satellites by large massive halos, and not to enhanced correlation between distinct halos.
Figure 4 and 5 also indicate that the efficient formation of faint passive galaxies is a threshold process that only occurs in systems where a galaxy has a very large number of neighbours; it does not proceed gradually with neighbour density. Clearly this suggest that mechanisms acting only in clusters (ram pressure and stripping) are more efficient than those acting also in groups (galaxy interactions) in generating faint passive galaxies.
Finally Fig. 5 indicates that the neighbour density range corresponding to the intersection between early-type and late-type galaxies moves towards richer regions as the luminosity of the samples decreases. At M<sub>B</sub> \- 5 $`\mathrm{log}`$ $`h_{75}`$ = -22 passive galaxies appear more numerous than SF galaxies even at the lowest densities (neigh=0). At M<sub>B</sub> \- 5 $`\mathrm{log}`$ $`h_{75}`$ = -20 equipartition is reached at group-like densities (neigh=3-8). At M<sub>B</sub> \- 5 $`\mathrm{log}`$ $`h_{75}`$ = -18 equipartition is never reached; passive galaxies are no more than one-third of the population even in the densest environment (neigh$`>`$8). This confirms that, on the group scale, luminosity generally dominates over environment in setting the spectral type mix of a population. The result is complementary to the finding (Norberg et al. 2002a ) that luminosity, and not type, is the dominant factor in determining how the clustering strength of the whole galaxy population varies with luminosity. However, because in luminous (faint) samples galaxies are mainly passive (active-SF), the correlation with type is strong at both the high and low luminosity end.
Given the correlation between spectral-type and morphology (Madgwick et al. (2002)), Fig.5 also likely illustrates a strong dependence of the morphology/density relation (Davis & Geller (1976); Dressler (1980); Postman & Geller (1984); Maia & da Costa (1990); Helsdon & Ponman (2003)) on luminosity. Nevertheless, to prove this dependence for the morphology/density relation might be difficult as the correlation between environment and stellar age (color and spectral-type) appears stronger than the correlation between environment and morphology (Willmer et al. (1998); Kauffmann et al. (2004); Blanton et al. 2005a ).
## 7 Summary and conclusions
In this paper we have investigated relations linking the spectral-type properties of 2dF galaxies to their luminosity and local neighbour density characteristics. We have assigned a local density to galaxies computing the number of neighbours within a 1 $`h_{75}^1`$ Mpc projected distance and $`\pm `$1000 km s<sup>-1</sup> depth. Our approach differs from previous analyses dealing with the same issue, because we have computed densities counting only fainter neighbours and applying a maximum magnitude difference criterion (โ2 $``$ M<sub>gal</sub> \- M<sub>neigh</sub> $``$0). This implies that, at least for luminous galaxies, neighbours likely trace the density of satellites that have been captured by the galaxy halo.
We have shown that the local density distribution for the all-type galaxy sample is approximately luminosity-independent over the whole explored luminosity range. This indicates that, at least on a 1$`h_{75}^1`$ Mpc scale, the number of neighbours associated with a galaxy is very similar, whether it reflects the number of satellites accreted by a luminous galaxy halo or the number of neighbours of low luminosity galaxies still in their original small halos.
We have also found that that the excess of fainter companions surrounding passive galaxies is not limited to luminous galaxies (Norberg et al. 2002a ; Hogg et al. (2003); Berlind et al. (2004); Blanton et al. 2005a ) but is instead a general characteristic of passive galaxies. The stronger clustering of passive galaxies relative to SF galaxies, on the group scale, appears to arise from two distinct contributions. At the luminous end, it is due to an excess of satellites surrounding central galaxies inside large halos. At the faint end it is due to an excess of satellites that are strongly correlated among themselves.
We have shown that the global spectral-type/luminosity relation (Fig.2) is not significantly altered in subsamples exhibiting different satellite densities. Whatever the environment, passive galaxies (Type 1) numerically dominate in luminous samples, and active-SF galaxies (Type 3+4) numerically dominate in faint ones. In contrast, the relative content in quiet-SF galaxies shows a weak dependence on luminosity. Only galaxies in extreme environments exhibit significant departures from these general trends: in dense environments (neigh$`>`$8) a significant fraction of passive galaxies is observed even among faint galaxies, whereas among isolated galaxies (neigh=0), quiet-SF (Type 2) galaxies still represent a 40% fraction of the luminous population. We suggest that these โminorityโ populations, identified in very luminous and very faint samples, are poor tracers of halo mass: luminous SF galaxies are actually tracers of small (sub-group scale) halos, whereas faint passive galaxies are tracer of massive (group/cluster scale) halos.
Our analysis provides evidence for the existence of a global spectral-type/satellite-density relation, with the fraction of passive galaxies steadily growing (and the fraction of quiet-SF galaxies steadily decreasing) when moving from an isolated galaxy sample to galaxies with cluster-like neighbour density. But we have also shown that this relation only holds in luminous samples; in faint samples the variation in the fractional content of passive (SF) between the neigh=0 and the intermediate dense (neigh=3-8) environments is not observed: the dependence on environment becomes threshold-like, and very dense environments are required to observe a variation in the spectral-type mix. This suggests that the morphology/density relation is likely a morphology/satellite-density relation, that traces enhanced correlation inside single massive halos rather than enhanced correlation between distinct halos.
###### Acknowledgements.
We thank A. Berlind, A. Biviano, R. De Propris, T. Goto and C.N.A. Willmer for comments and suggestions. We are also indebted to the anonymous referee whose comments and criticism greatly improved the scientific content of this paper. This work was supported by MIUR, BK acknowledges a fellowship from Bologna University.
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# Critical Examinations of QSO Redshift Periodicities and Associations with Galaxies in Sloan Digital Sky Survey Data
## 1 Introduction
The debate on whether QSOs are ejected from the nuclei of low-redshift galaxies with a periodic non-cosmological โintrinsicโ redshift has been going on for many years. Some evidence has been claimed to suggest such an intrinsic redshift hypothesis, in which QSOs have redshifts that are much larger than their parent galaxies and the excess of redshift is assumed to represent an always redshifted intrinsic component (Burbidge & Burbidge, 1967; Arp et al., 1990; Karlsson, 1990; Chu et al., 1998; Burbidge & Napier, 2001; Bell, 2004b; Arp et al., 2005).
Two models have been discussed in the literature which predict exact values for the preferred redshifts. One of these is the Karlsson formula which suggests a periodicity of $`\mathrm{}\mathrm{log}(1+z_{eff})=0.089`$ with peaks lying at 0.061, 0.30, 0.60, 0.96, 1.14, 1.96 and so on (Karlsson 1977, 1990; Arp et al. 1990, 2005; Burbidge $`\&`$ Napier 2001, 2003), where $`z_{eff}`$ is the redshift of the QSO measured relative to the nearby galaxy, called effective redshift, which is defined as:
$$1+z_{eff}=(1+z_Q)/(1+z_G)$$
(1)
where $`z_Q`$ is the observed quasar redshift and $`z_G`$ is the redshift of the associated galaxy which is assumed to be the ejecting galaxy. To explain such a periodicity, they claimed that quasars are ejected by active galaxies and the putative parent galaxies are generally much brighter than their quasar off-springs (Arp et al., 2005). As claimed by Burbidge $`\&`$ Napier (2001, 2003), the typical projected association separation is about 200 kpc.
Another model, namely decreasing intrinsic redshift model (DIR model), was proposed by Bell (2004), where the QSO intrinsic redshift equation is given by the relation:
$$z_{iQ}=z_f(NM_N)$$
(2)
where $`z_f=0.62\pm 0.01`$ is the intrinsic redshift constant, $`N`$ is an integer, and $`M_N`$ varies with $`N`$ and is a function of a second quantum number $`n`$. In the DIR model, galaxies are produced continuously through the entire age of the universe, and QSOs are assumed to be ejected from the nuclei of active galaxies and represent the first very short lived stage ($`10^710^8`$ yr) in the evolution of galaxies (Bell, 2004b).
The above an intrinsic redshift hypothesis, if true, will have far-reaching consequences for cosmology and the nature of QSOs. Most of those previous studies on the Karlsson formula used rather small samples (except for Arp et al. 2005), and have been suspected that the claimed peaks were due to artifacts associated with selection effects (Basu, 2005). To avoid such a heterogeneous selection manner as well as personal prejudice, Hawkins et al. (2002) tested the periodicity in $`\mathrm{log}(1+z_{qso})`$ with 2dF redshift survey data with 67291 nearby galaxies and 10410 QSOs; it was found that there is no periodicity in $`\mathrm{log}(1+z_{qso})`$. However, Napier $`\&`$ Burbidge (2003) argued that in order to use the 2dF sample to properly test the original hypothesis, it is necessary to establish for each pair that the galaxy is at least a late-type active spiral system. Arp et al. (2005) also re-examined the 2dF sample and claimed that they found that the redshifts of brighter QSOs in the QSO density contours fit very exactly the long standing Karlsson formula and confirm the existence of preferred values in the distribution of quasar redshifts.
In an attempt to resolve these issues, we turn to the Sloan Digital Sky Survey (SDSS) (and also 2dF QSO Redshift Survey (2QZ) occasionally) to carry out this study, which have the largest homogeneous sample of data as well as the spectroscopic sub-classification of galaxies. In section 2, to test whether there is a periodicity existing in $`\mathrm{log}(1+z)`$, we construct four sets of QSO-galaxy pairs with different QSOs and galaxies, with all QSOs projected within 200 kpc from nearby galaxies at these galaxiesโ distances. QSO density contours are also presented to show that there is no periodicity in SDSS QSOs under such analysis. In section 3, we examine the relationship between high-z QSOs and nearby active galaxies to show that such QSOs are not likely to be ejected by nearby active galaxies. In section 4, we analyze the redshift distribution of QSOs in SDSS DR3 and 2QZ to show that there is no evidence for non-artificial periodicity in redshifts of QSOs, contrary to the DIR model. Discussion and conclusion are described in Section 5.
## 2 No Periodicity in $`\mathrm{log}(1+z)`$
### 2.1 The SDSS Data and Pair Selection
In this section, we use the SDSS DR1 QSO catalog (Schneider et al., 2003) and the New York University Value-Added Galaxy Catalog (NYU-VAGC) (Blanton et al., 2005). For reliability in the derived redshifts, we consider only those QSOs in the range of $`z>0.4`$, and galaxies in the range of $`0.01<z<0.2`$ with the highest plate quality labeled as โgoodโ and with no redshift warning. This quality control leaves a total of 190591 galaxies and 15747 QSOs in the sample.
Three issues need to be carefully addressed when we analyze the relation between foreground galaxies and QSOs, as well as the redshift distribution of QSOs, since due to the survey strategies and the instrumental limitations, the selections of galaxies and QSOs are not entirely independent, and the selection of QSOs in SDSS is also dependent on redshift. First, due to the mechanical constraint in SDSS that spectroscopic fibers must be separated by $`55^{\prime \prime }`$ on a given plate (Blanton et al. 2003), consequently some QSO-galaxy pairs would be missing from the spectroscopic sample. However, this issue would have little effect on the results for two reasons: (1) $`55^{\prime \prime }`$ corresponds to an angular distance of $`40`$ kpc for a galaxy at $`z=0.04`$, which is the typical value of redshift in our galaxy sample, only few pairs (about $`5\%`$) would be missing in a given separation of 200 kpc for randomly distributed QSOs and galaxies, which is also shown in Fig. 7 where the distribution of data pairs are consistent with randomly distributed pairs; (2) such fiber constraint is independent of redshifts of galaxies or QSOs, therefore its redshift distribution will not be biased, although some pairs are missing in the sample. Second, the magnitude limits of the SDSS galaxy and QSO spectroscopic surveys are quite different, i.e., $`i<19.1`$ for $`z<3`$ QSOs, $`i<20.2`$ for $`z>3`$ QSOs, $`r<17.77`$ for most sampled galaxies and $`r<19.5`$ for luminous red galaxies (Richards et al. 2002; Strauss et al. 2002), hence the magnitudes of QSOs are mostly higher than galaxies. However, since in the ejection hypothesis, the parent galaxies are generally much brighter than their QSO off-springs (Bell 2004b; Arp et al. 2005), the pair making process are not likely to be affected by the magnitude differences between QSOs and galaxies, which is also shown in Fig. 12 and Fig. 13. Moreover, the fact that the completeness of the spectroscopic selection varies with redshift (Richards et al. 2002) will consequentially affect the redshift distribution of QSOs and might cause artificial periodicities into data, as will be discussed extensively in Section 4. However, since low-z ($`z<2.5`$) QSOs which have flat and high completeness level ($`>90\%`$; Richards et al. 2002) occupied a very large fraction of all QSOs (about $`90\%`$), such selection dependence in redshift would not smear out intrinsic periodicities in QSO redshift if they do exist.
It has been suggested that quasars with bright apparent magnitude will generally be nearby and the redshifts of these quasars will require little or no correction for the periodicity effects to be manifested, and those low redshift galaxies with which such quasars appear to be preferentially associated, tend to be morphologically disturbed active galaxies (Arp et al., 2005). Therefore, to make our results more compelling, we select a sub-sample of 3724 QSOs, called bright QSOs, which have i-band magnitudes less than 18.5, and a sub-sample of 77426 galaxies, called active galaxies, which are labeled as starforming, starburst, starforming broadline or starburst broadline galaxies in the spectroscopic sub-classification. Then we construct four sets of QSO-galaxy samples by intercrossing them, and get four sets of pairs in which a QSO is projected within 200 kpc from a galaxy: 4572 pairs for QSO-nearby galaxies, 3216 pairs for QSO-active nearby galaxies, 1129 pairs for bright QSO-nearby galaxies and 791 pairs for bright QSO-active nearby galaxies. When there is more than one galaxy within the 200 kpc projected distance limit of the QSO, we take the closest galaxy in projected distance to make up the pair.
### 2.2 Analysis and Results
We make power spectrum analysis to investigate the periodicity hypothesis of Karlsson (1977). The power $`I`$ is defined as in Burbidge $`\&`$ Napier (2001):
$$I(\nu )=2R^2/N,$$
(3)
where
$$R^2=S^2+C^2,$$
(4)
with
$$S=\underset{i=1}{\overset{N}{}}w_i\mathrm{sin}(2\pi x_i/P),C=\underset{i=1}{\overset{N}{}}w_i\mathrm{cos}(2\pi x_i/P),$$
(5)
with $`\nu =1/P`$ and $`x_i=\mathrm{log}(1+z_i)`$. Here $`w_i`$ is a weighting function, and $`w_i1`$ except in section 4, as in the analysis of Burbidge $`\&`$ Napier (2001). For randomly and uniformly distributed redshifts, $`\overline{I}=2`$.
To test our code developed for this study, we first re-analyze the 290 QSOs in Karlsson and Napier $`\&`$ Burbidgeโs data sets (Karlsson, 1990; Burbidge & Napier, 2001; Napier & Burbidge, 2003), as shown in Fig. 1. Errors on $`I(P)`$ are given by using the bootstrap methods (Efron, 1979) in the following steps: (1) we take the non-zero number in each bin in the upper histogram as the expectation value of a Poisson distribution; (2) we re-sample each bin following the Poisson distribution to re-produce 1000 new sets of data, repeat the power spectral analysis on these re-samplings, and finally calculate the standard deviations in the derived values of $`I`$ at different periods $`P`$. Clearly the periodicity at around $`\mathrm{}\mathrm{log}(1+z_{eff})=0.089`$ is highly significant at above 3.5$`\sigma `$ level.
In Figs. 2-5 we show histograms of the effective redshifts of QSOs paired with galaxies and their unwindowed power spectra with standard deviations calculated in the same way as for Fig. 1. Pairs in these four figures are for QSO-nearby galaxies, QSO-active nearby galaxies, bright QSO-nearby galaxies and bright QSO-active nearby galaxies, as described in Sec. 2.1. Our results show that for these significantly larger samples than that in Fig. 1, all peaks appeared in the power spectra are consistent with Poisssonian fluctuations, i.e., there is no evidence for a periodicity at the predicted frequency in $`\mathrm{log}(1+z)`$, or at any other frequency.
### 2.3 QSO Density Contours
After the work of Hawkins et al. (2002) on 2dF data, Arp et al. (2005) argued that the predicted periodic redshifts are apparent in the brighter 2dF quasars in the QSO density contours. We therefore use SDSS DR1 data to construct the contours defined by Arp et al. (2005), where the whole region is divided into boxes $`z\times B=0.075\times 0.3`$ in the redshift/apparent magnitude plane, then the number of quasars in each box is counted. To show whether the predicted periodic redshifts are obscured by our coarse grid sizes, a contour with $`z\times B=0.05\times 0.2`$ is also presented for comparison. As shown in Fig. 6, the peak positions are consistent in the two contour plots, and there is no evidence for redshift peaks at the predicted positions.
## 3 No Strong Connection between Active Galaxies and Bellโs High-z QSOs
In Bell (2004b), a high-redshift QSO sample from SDSS and a low-redshift QSO and QSO-like object sample from Hewitt $`\&`$ Burbidge (1993) were presented. Though the dips at redshifts of 2.7 and 3.5 have been explained as being caused by the lower efficiency of the selection algorithm at these redshifts (Richards et al., 2002), Bell (2004b) nevertheless claimed that the corresponding redshift peaks at 3.1 and 3.7 in the high-z SDSS QSOs come from the intrinsic redshift broadening by Doppler ejection and Hubble flow components, which is in favor of the DIR model. Through analysis of the profiles of such peaks, Bell (2004b) derived a mean cosmological components to be $`z_c0.066`$ for the high-z sample.
In the DIR model, galaxies are produced continuously through the entire age of the universe, and QSOs are assumed to be ejected from the nuclei of active galaxies and represent the first very short lived stage ($`10^710^8`$ yr) in the evolution of galaxies. If this is true, there must be some connection between foreground active galaxies and high-z QSOs beyond gravitational lensing. Here we examine the high-redshift samples taken from the SDSS data. We test the relationship between 2691 QSOs with redshifts in 2.4$``$4.8 and 77426 nearby active galaxies with redshifts in 0.01$``$0.2 from NYU-VAGC, all of which have the highest plate quality labeled as โgoodโ and with no redshift warning. We inter-compare these two data sets to find all QSO-galaxy pairs within an angular separation corresponding to less than a given distance $`D_{separation}`$ from several kpc to 1 Mpc at the distance of the galaxy. In some cases, there is more than one galaxy within the $`D_{separation}`$ projected distance limit of the QSO; for these objects we take the closest galaxy in projected distance to make up the pair. Since it is suggested that all QSOs are born out of active galaxies and QSOs should be significantly fainter than their parent galaxies (Bell 2004b), we would not miss a considerable fraction of parent active galaxies for high-z QSOs if the DIR model is right.
The distribution of projected separation distance for all pairs is shown in Fig. 7, and the redshift distribution of active galaxies in pairs with QSOs is shown in Fig. 8. Both of them are consistent with random distributions, but totally different from the distribution from the ejection simulation with a ejection velocity of 11,000 km s<sup>-1</sup> which was given by Bell (2004b) as typical values, and the mean redshift of these galaxies is $`z0.044`$, also significantly different from Bellโs result of 0.066. Here the random distributions are obtained by moving the positions of all galaxies by 1 degree in random directions; thus these galaxies should be completely unrelated to background QSOs. The ejection simulation is done by ejecting all QSOs from randomly selected active galaxies with three ejection velocities: 11,000 km s<sup>-1</sup>, 40,000 km s<sup>-1</sup> and 80,000 km s<sup>-1</sup>, and with a uniformly distributed age of $`010^8`$ yr which is given by Bell (2004b) as a typical value.
To quantitatively show the differences between simulations and โtrueโ QSO-galaxy pairs, i.e., pairs found in the data but not necessarily physical pairs, results of chi-squared tests are given in Table 1. In $`3\sigma `$ confidence level for both distributions of projected separation distance and redshift distribution of active galaxies in pairs with QSOs, the โtrueโ QSO-galaxy pairs are consistent with random distributions, but inconsistent with ejection hypothesis with ejection velocity up to 80,000 km s<sup>-1</sup>.
## 4 No Periodicity in $`z`$
We also analyze the periodicity in redshifts of QSOs in SDSS DR3 (Schneider et al. 2005) and 2dF (Croom et al. 2004) to investigate in larger database of QSOs whether there is a periodicity of $`\mathrm{\Delta }z=0.67\pm 0.05`$, predicted by the DIR model. Six data sets are used in this section: all 46,420 QSOs in SDSS DR3 (Fig. 9), 22,497 QSOs with the highest quality flag in 2dF (Fig. 10), a high completeness (close to $`100\%`$) sub-sample containing 23,109 QSOs with Galactic-extinction corrected $`i`$-band magnitude ($`m_i`$) less than 19 and redshift less than 2 in SDSS DR3 (Richards et al. 2002) (Fig. 11(a)), and three sub-samples containing QSOs in low completeness (less than $`50\%`$) regions in SDSS DR3: 15,696 QSOs with $`m_i>19`$ and $`z<2.4`$ (Fig. 11(b)), 19,064 QSOs with $`m_i>19`$ in all redshifts (Fig. 11(c)), and 9,763 QSOs with $`z>2`$ (Fig. 11(d)). To reduce the edge effect produced by the truncated redshift distribution which has a lower redshifts cut-off due to the small space volume sampled and a higher redshifts cut-off due to the observational flux limit (see e.g. Hawkins et al. 2002), we follow Hawkins et al. to use the Hann function as a weighting in equation 5,
$$w_i=\frac{1}{2}[1\mathrm{cos}(\frac{2\pi x_i}{L})],$$
(6)
where $`L`$ is chosen to cover the range of redshifts. Here $`L=5.4`$ for the full SDSS sample, $`L=3.1`$ for the 2dF sample, and $`L=1.95,2.1,5.1`$ and 3.4 for the four SDSS sub-samples respectively.
After smoothing off the sharp edges in the lowest and highest redshifts, a periodicity around $`\mathrm{\Delta }z=0.67`$ is detected in the full sample of SDSS QSOs, as shown in Fig. 9; however a periodicity of $`\mathrm{\Delta }z=0.67\pm 0.05`$ or any other frequency is not found in the 2dF QSOs, as shown in Fig. 10. Such a difference between these two surveys is not surprising since the redshift-dependent spectroscopic completeness is relatively flat in 2dF (Croom et al. 2004), while in SDSS the spectroscopic completeness varies drastically at some redshifts (Richards et al. 2002). It is therefore improper to use all QSO redshifts in SDSS to probe any intrinsic periodicity without addressing selection bias. To further investigate whether such a periodicity around $`\mathrm{\Delta }z=0.67`$ in SDSS QSOs is spuriously produced by various incompleteness as function of redshift, we select a high-completeness sub-sample of 23,109 QSOs with $`m_i<19`$ and $`z<2`$ in SDSS DR3, and three sub-samples containing QSOs in low-completeness regions. As shown in Fig. 11, no periodicity is found in the high-completeness sample where the power spectrum is consistent with a continuously ascending curve due to the low frequency component of the redshift distribution, whereas in different low-completeness samples, strong periodicity always appears, but with different peak locations (0.88 in (b), 0.67 in (c) and 0.74 in (d)). This should be a strong indicator that the peaks in low-completeness samples are caused by different selection effects in different samples. In sum, there is no evidence for intrinsic periodicity in redshifts of QSOs.
## 5 Discussion and Conclusion
However, one might ask whether it is because we have some paired QSOs with wrong parent galaxies so that not only the effective redshifts of QSOs show no periodicity, but also high-z QSOs and nearby active galaxies show no strong connection. The wrong-pairing indeed could happen that when there is more than one galaxy within the $`D_{separation}`$ projected distance limit of the QSO and we take the closest galaxy in projected distance to make up the pair. In the following we quantitatively examine this possibility and its effect.
For the pair making process in section 2, the $`D_{separation}`$ is 200 kpc, which is less than the average projection distance between QSOs and galaxies, and the percentage of cases in which there are two or more galaxies within the projected distance is 27$`\%`$ for QSO-nearby galaxies, 19$`\%`$ for QSO-active nearby galaxies, 27$`\%`$ for bright QSO-nearby galaxies and 18$`\%`$ for bright QSO-active nearby galaxies, respectively. This means that for a majority of paired QSOs ($`>73\%`$), there is only one galaxy within the given projected distance and would not be paired incorrectly, hence the claimed periodicity should have been detected in our larger samples if they did exist.
For the pair making process in section 3, the largest $`D_{separation}`$ is 1 Mpc which is larger than the average projection distance between QSOs and galaxies ($``$400 kpc), therefore wrong-pairing may occur more frequently here. It will be even worse if the typical ejection distance is larger than the mean projection separation of QSOs and active galaxies. So would such wrong pairs result in the good agreement between ejected model and randomly generated pairs? We answer this question by making the following test. Suppose that all QSOs are ejected by randomly selected active galaxies with a given ejection velocity (11,000 km s<sup>-1</sup>, 40,000 km s<sup>-1</sup> and 80,000 km s<sup>-1</sup>), and with a uniformly distributed age of $`010^8`$ yr which is given by Bell (2004) as typical values, we get 200 sets of false pairs. As shown in Fig. 7-8 and Table 1, the distribution of such simulated ejection QSO-galaxy pairs are totally different from random distribution. We therefore conclude that the random-like distribution of QSO-active galaxy pairs could not be produced by the ejection model.
Another question is that whether we miss periodicities of QSOs by setting a lower limit of $`z=0.01`$ for galaxies and no constraint in QSO-galaxy magnitude relation. Though the lower limit of $`z=0.01`$, the same as in Hawkins et al. (2002), is chosen to have confidence in the derived angular distance, as well as avoid large projection effect of very nearby galaxies, we re-analyze SDSS DR1 QSOs and galaxies again with no redshift limits on galaxies and set magnitude constraint that all paired QSOs should be at least 5 or 3 magnitudes fainter than the paired galaxy. As shown in Fig. 11 and Fig. 12, similar to our results in section 2.2, there is no evidence for a periodicity at the predicted frequency in log$`(1+z)`$, or at any other frequency.
In summary, using samples from SDSS and 2QZ, we demonstrate that not only there is no periodicity at the predicted frequency in $`\mathrm{log}(1+z)`$ and $`z`$, or at any other frequency, but also there is no strong connection between foreground active galaxies and high redshift QSOs. These results are against the hypothesis that QSOs are ejected from active galaxies or have periodic intrinsic non-cosmological redshifts.
We thank the anonymous referee and Dr. Bell for valuable suggestions that have significantly improved this paper. This study is supported in part by the Special Funds for Major State Basic Research Projects and by the National Natural Science Foundation and the Ministry of Education of China. SNZ also acknowledges NASA for partial financial support through several research grants.
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# Kaon Electromagnetic Form Factor in the Light-Front Formalism11footnote 1To appear โPhysics of Elementary Particles and Atomic Nuclei, Vol. 36, (2005).โ
## Kaon Electromagnetic Form Factor in the Light-Front Formalism<sup>1</sup><sup>1</sup>1To appear โPhysics of Elementary Particles and Atomic Nuclei, Vol. 36, (2005).
Fabiano P. Pereira<sup>a</sup>, J.P.B.C. de Melo<sup>a,,b</sup>, T. Frederico<sup>c</sup> and Lauro Tomio<sup>a</sup>
(a) Instituto de Fรญsica Teรณrica, Universidade Estadual Paulista,
01405-900, Sรฃo Paulo, SP, Brazil
(b) Centro de Ciรชncias Exatas e Tecnolรณgicas, Universidade Cruzeiro do Sul,
08060-070, Sรฃo Paulo, Brazil
(c) Departamento de Fรญsica, ITA, Centro Tรฉcnico Aeroespacial,
12228-900, Sรฃo Josรฉ dos Campos, Brazil
## Abstract
Numerical calculations are performed and compared to the experimental data for the electromagnetic form factor of the kaon, which is extracted from both components of the electromagnetic current, $`J^+`$ and $`J^{}`$, with a pseudo-scalar coupling of the quarks to the kaon. In the case of $`J^+`$ there is no pair term contribution in the Drell-Yan frame ($`q^+=0`$). However, for $`J^{}`$, the pair term contribution is different from zero and necessary in order to preserve the rotational symmetry of the current. The free parameters are the quark masses and the regulator mass.
### 1. Introduction
The convenience of using light-front variables in QCD descriptions of hadron properties and interactions, has been established long time ago. In particular, we have a clarifying article by N.N. Bogolyubov and coworkers published in 1983 . In the subsections 3.5.3 and 3.5.4 of this article, the QCD description of the simplest composite systems (the mesons) and the corresponding form factors at high momentum transfer are discussed. The formalism is developed from a gauge invariant two-point function of the Bethe-Salpeter amplitude. The pion electromagnetic form factors are shown as example using the light-front formalism. Models for wave-functions in the light front are originally developed in . In more recent years, the use of light-front formalism has become a common procedure in QCD description of hadrons . Here, we can mention a few works that we are concerned, as , dedicated to study pseudoscalar properties of mesons, structure wave functions and quark-antiquark correlations. From such references, one can trace a more complete and detailed bibliography.
In the present communication, we report results for the kaon electromagnetic form factor that are extracted from both components of the electromagnetic current, $`J^+=J^0+J^3`$ and $`J^{}=J^0J^3`$, with a pseudo-scalar coupling of the quarks. In the case of $`J^+`$ there is no pair term contribution in the Drell-Yan frame ($`q^+=0`$). However, for the $`J^{}`$ component of the electromagnetic current, the pair term contribution is different from zero and necessary in order to preserve the rotational symmetry of the current. We note that, when considering the case of vector particles, even the $`J^+`$ electromagnetic current has contribution from pair terms to have a full covariant theory (For a more recent application of this ideas in the vector anomaly problem, see ref.\[bf05\]). In order to satisfy the angular condition for spin one particles, it is necessary to consider pair terms in the electromagnetic current $`J^+`$ . Besides the valence contribution to the $`J^{}`$ current, the pair term is necessary for both, pseudoscalar and vector particles to keep the rotational symmetry properties of the current in the light-front formalism.
## 2. Electromagnetic Current Model
In order to extract the electromagnetic form factor for the kaon, the components $`J^+`$ and $`J^{}`$ of the electromagnetic current are used. The $`J^{(\mu =\pm )}`$ components of the electromagnetic current for the kaon have contribution, from the quark ($`q`$) and the antiquark ($`\overline{q}`$), are given by
$`J_q^\mu (q^2)=`$ $`ฤฑe_qg^2N_c{\displaystyle }{\displaystyle \frac{d^4k}{(2\pi )^4}}\times `$
$`\times `$ $`\mathrm{Tr}[S(k,m_{\overline{q}})\gamma ^5S((kP^{}),m_q)\gamma ^\mu S((kP),m_q)\gamma ^5]\mathrm{\Lambda }(k,P^{})\mathrm{\Lambda }(k,P),`$
$`J_{\overline{q}}^\mu (q^2)=`$ $`q\overline{q}\text{in}J_q^\mu (q^2),`$ (1)
where the number of colors is $`N_c=3`$, $`g`$ is the coupling constant and $`e_q`$ ($`e_{\overline{q}}`$) is the quark (anti-quark) charge. We use the Breit frame, where the momentum transfer is $`q^2=(\stackrel{}{q}_{})^2`$, $`P^0=P^{\mathrm{\hspace{0.17em}0}}`$ and $`\stackrel{}{P}_{}^{}=\stackrel{}{P}_{}=\stackrel{}{q_{}}/2`$. The function $`\mathrm{\Lambda }(k,p)=N/[(pk)^2m_R^2+ฤฑฯต]`$ is used in order to regulate the divergent integral, where $`m_R`$ is the regulator mass and $`m_q`$ and $`m_{\overline{q}}`$ are, respectively, the quark and anti-quark masses. The function $`S(p)`$ is the fermion propagator:
$$S(p,m)=\frac{1}{\text{/}pm+ฤฑฯต}.$$
(2)
The light-front coordinates are defined as $`k^+=k^0+k^3,k^{}=k^0k^3,k_{}=(k^1,k^2)`$. In the following, for the calculation of the pair terms, we consider the model given in for a composite boson bound state and in the study of the Ward-Takahashi identity in the light-front formalism . The contribution of the pair term for $`J^+`$ and $`J^{}`$ components of the electromagnetic current comes from the matrix elements proportional to $`k^{}`$ in both cases (antiquark and quark on-shell).
## 3. Electromagnetic Form Factor
The most general expression for the form factor of the spin zero particles is given by:
$$P|J^\mu |P^{}=(P^{}+P)^\mu F(q^2)+(P^{}P)^\mu G(q^2)$$
(3)
In this elastic process, the form factor depends only on $`q^2`$, and $`G(q^2)=0`$ in all $`q^2`$. Here, off-shell effects are not explored. However, the off-shell effects are important and relevant in many topics for particles and nuclear physics.
In order to extract the form factor for the kaon, $`F_{K^+}(q^2)`$, we used both $`J^+`$ and $`J^{}`$ components of the electromagnetic current. One can verify that only the on-shell pole $`\overline{k}^{}=(k_{}^2+m_{\overline{q}}^2)/k^+`$ contribute to the $`k^{}`$ integration in the interval $`0<k^+<P^+`$:
$`F_{\overline{q}}^+(q^2)`$ $`=`$ $`e_{\overline{q}}{\displaystyle \frac{N^2g^2N_c}{P^+}}{\displaystyle \frac{d^2k_{}dk^+}{4\pi ^3}\frac{๐ฉ_{\overline{q}}^+\theta (k^+)\theta (P^+k^+)}{k^+(P^+k^+)^2(P^{}_{}{}^{}+k^+)^2(P^{}\overline{k}^{}\frac{f_{2,q}ฤฑฯต}{P^+k^+})}}`$ (4)
$`\times {\displaystyle \frac{1}{(P^{}\overline{k}^{}\frac{f_{3,q}ฤฑฯต}{P^+k^+})(P^{}\overline{k}^{}\frac{f_4ฤฑฯต}{P^+k^+})(P^{}\overline{k}^{}\frac{f_5ฤฑฯต}{P^+k^+})}},`$
$`F_q^+(q^2)`$ $`=`$ $`[q\overline{q}\text{in}F_{\overline{q}}^+(q^2)],`$ (5)
where $`f_{2,q}(Pk)_{}^2+m_q^2`$, $`f_{3,q}(P^{}k)_{}^2+m_q^2`$, $`f_4(Pk)_{}^2+m_R^2`$ and $`f_5(P^{}k)_{}^2+m_R^2`$ . In the numerator, $`๐ฉ_{\overline{q}}^+`$ is given by
$`๐ฉ_{\overline{q}}^+={\displaystyle \frac{1}{4}}\mathrm{Tr}[(\text{/}k+m_{\overline{q}})\gamma ^5(\text{/}k\text{/}P^{}+m_q)\gamma ^+(\text{/}k\text{/}P+m_q)\gamma ^5]|_{k_{}=\overline{k}^{}}.`$ (6)
The kaon light-front wave function of the model can be extracted from (4)and (5) as
$`\mathrm{\Phi }_Q^i(x,k_{})={\displaystyle \frac{1}{(1x)^2}}{\displaystyle \frac{N}{(m_{K^+}^2M_0^2)(m_{K^+}^2M^2(m_Q,m_R))}},`$ (7)
where $`x=k^+/P^+`$ is the momentum fraction, $`Q=\overline{q},q`$ and
$$M^2(m_Q,m_R)=\frac{k_{}^2+m_Q^2}{x}+\frac{(Pk)_{}^2+m_R^2}{1x}P_{}^2.$$
(8)
The squared free quark mass is given by $`M_0^2=M^2(m_{\overline{q}},m_q)`$. For the final wave-functions, $`\mathrm{\Phi }_{\overline{q}}^f`$ and $`\mathrm{\Phi }_q^f`$, we just need to exchange $`PP^{}`$ in (7) and (8).
The expression obtained for the electromagnetic form factor in terms of the initial $`(\mathrm{\Phi }_{\overline{q}}^i)`$ and final $`(\mathrm{\Phi }_{\overline{q}}^f)`$ wave functions is
$`F_{\overline{q}}^+\left(q^2\right)`$ $`=`$ $`e_{\overline{q}}{\displaystyle \frac{N^2g^2N_c}{4\pi ^3P^+}}{\displaystyle \frac{d^2k_{}dx}{x}๐ฉ_{\overline{q}}^+\theta \left(x\right)\theta \left(1x\right)\mathrm{\Phi }_{\overline{q}}^f(x,k_{})\mathrm{\Phi }_{\overline{q}}^i(x,k_{})},`$ (9)
$`F_q^+\left(q^2\right)`$ $`=`$ $`\left[q\overline{q}\text{in}F_{\overline{q}}^+\left(q^2\right)\right].`$ (10)
The final expression for the electromagnetic form factor obtained with $`J^+`$ is the sum of two contributions from the quark and the antiquark currents:
$$F_{K^+}^+(q^2)=F_q^+(q^2)+F_{\overline{q}}^+(q^2),$$
(11)
where the normalization is given by $`F_{K^+}^+(0)=1`$. The calculation of the kaon electromagnetic form factor in the light-front with $`J^+`$, without pair term, gives the same result as the covariant one (see Fig.1).
The contribution to the electromagnetic form factor obtained with $`J^{}`$ after the integration in $`k^{}`$ from the interval $`0<k^+<P^+`$ is given by
$`F_{\overline{q}}^{(I)}(q^2)`$ $`=`$ $`e_{\overline{q}}{\displaystyle \frac{N^2g^2N_c}{4\pi ^3P^+}}{\displaystyle \frac{d^2k_{}dx}{x}\theta (x)\theta (1x)๐ฉ_{\overline{q}}^{(I)}\mathrm{\Phi }_{\overline{q}}^f(x,k_{})\mathrm{\Phi }_{\overline{q}}^i(x,k_{})},`$ (12)
$`F_q^{(I)}(q^2)`$ $`=`$ $`[q\overline{q}\text{in}F_{\overline{q}}^{(I)}(q^2)],`$ (13)
where
$`๐ฉ_{\overline{q}}^{(I)}`$ $`=`$ $`{\displaystyle \frac{k_{}^2+m_{\overline{q}}^2}{xP^+}}\left[(m_qm_{\overline{q}})^2{\displaystyle \frac{q^2}{4}}\right]+P^+\left[2m_{\overline{q}}(m_qm_{\overline{q}})+xP^{+2}\right]`$ (14)
When using $`J^{}`$ to extract the electromagnetic form factor, besides the contribution of the interval (I), the pair term contributes to the electromagnetic form factor in the interval (II) ($`P^+<k^+<P^+`$). The pair term contribution for the form factor, as shown in , is given by $`F^{(II)}(q^2)`$.
$$F^{(II)}(q^2)=\frac{N^2g^2N_c}{P^+}\left[e_q\mathrm{\Delta }_q^{(II)}(q^2)+e_{\overline{q}}\mathrm{\Delta }_{\overline{q}}^{(II)}(q^2)\right]$$
(15)
where $`\mathrm{\Delta }_{\overline{q}}^{(II)}`$ and $`\mathrm{\Delta }_q^{(II)}`$ given below. These terms correspond to the pair contribution in the $`J^{}`$ component of the electromagnetic current, which are obtained after the integration in $`k^{}`$ and the limit $`P^+P^+`$ are performed. Then one gets the following equations for the pair terms:
$`\mathrm{\Delta }_{\overline{q}}^{(II)}=`$ $`{\displaystyle \frac{1}{4\pi ^3}}{\displaystyle \frac{d^2k_{}}{(2\pi )^4}๐ฉ_{\overline{q}}^{(II)}\underset{i=2}{\overset{5}{}}\frac{\mathrm{ln}(f_i)}{_{j=2,ij}^5(f_jf_i)}},`$ (16)
where $`f_2f_{2,q}`$, $`f_3f_{3,q}`$ and
$`๐ฉ_{\overline{q}}^{(II)}={\displaystyle \frac{1}{P^+}}\left[P^{+2}+{\displaystyle \frac{q^2}{4}}(m_qm_{\overline{q}})^2\right].`$ (17)
We obtain the corresponding quark current contribution as in the above Eqs. (16) and (17) just by replacing $`q\overline{q}`$.
In the limit $`P^+P^+`$, the pair term contribution (zero mode) is non zero and responsible for the covariance of the $`J^{}`$ component of the electromagnetic current. The sum of the contributions from the intervals (I) and (II) for $`J^{}`$ in the light-front gives the same result as in the covariant calculation .
The final expression for the electromagnetic form factor for the kaon, extracted from $`J^{}`$ is
$$F_{K^+}^{}(q^2)=\left[F_q^{(I)}(q^2)+F_{\overline{q}}^{(I)}(q^2)+F^{(II)}(q^2)\right],$$
(18)
which is normalized by the charge conservation, $`F_{K^+}^{}(0)=1`$.
## 4. Results and Conclusion
Next, we present results obtained considering the light-front formalism, as well as the covariant formalism. The parameters of the model are the constituent quark masses $`m_q=m_u=0.220`$ GeV, $`m_{\overline{q}}=m_{\overline{s}}=0.419`$ GeV, and the regulator mass $`m_R=0.946`$ GeV, which are adjusted to fit the electromagnetic radius of the kaon. With these parameters, the calculated electromagnetic radius of the kaon is $`r_{k^+}^2=0.354`$ fm<sup>2</sup>, very close to the experimental radius $`r_{k^+}^2=0.340`$ fm<sup>2</sup> .
The electromagnetic form factor is presented in Fig.1. Due to the fact that $`J^+`$ does not have light-front pair term contributions, the electromagnetic form factor results equal to the one obtained in a covariant calculation. In the case of $`J^{}`$, the light-front calculation gives results quite different from the covariant results, as shown in figure 1. After the inclusion of the pair terms, with $`J^{}`$, we observe a complete agreement between the light-front and covariant results. In conclusion, the $`J^+`$ and $`J^{}`$ components of the electromagnetic current of the kaon are obtained in the light-front and in the covariant formalisms, in a constituent quark model. In the case of $`J^{}`$, we note that the pair terms are essential to obtain a complete agreement between the covariant and the light-front results for the kaon electromagnetic form factor.
Our thanks to the Brazilian agencies FAPESP and CNPq for partial support.
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# Radio quiet neutron star 1E 1207.4-5209 as a source of gravitational waves
## I Introduction
The radio quiet neutron star 1E 1207.4-5209 (hereafter 1E1207) is at the center of the supernova remnant (SNR) PKS 1209$``$51/52. It was discovered by Helfand $`\&`$ Becker helfand84 with the Einstein Observatory. The distance to the SNR is $`d=2.1_{0.8}^{+1.8}`$ kpc Giac00 . The X-ray spectrum of the central source can be described by a thermal model which gives a distance of 2 kpc Mere96 ; Vas97 ; Zavlin98 .
The long observations devoted to 1E1207 both by Chandra and by XMM-Newton have unveiled a number of unique and somewhat contradictory characteristics that, at the moment, defy standard theoretical interpretations Bignami04 .
The characteristic age of the pulsar, 200 to 900 kyr Pavlov02 , is much larger than the estimated age of the SNR, 3 to 20 kyr Roger88 .
The values of the spin-down luminosity, $`\dot{E}1\times 10^{34}\mathrm{erg}\mathrm{s}^1`$, and conventional magnetic field (B-field), $`B3\times 10^{12}`$ G, are typical for a radio pulsars Zavlin04 . However such a B-field is significantly different from the value obtained from the cyclotron absorption lines interpreted both in terms of electrons ($`B8\times 10^{10}`$ G) as well as protons ($`B1.6\times 10^{14}`$ G) Sanwal02 ; Luca04 .
Chandra and XMM-Newton observations indicate that the pulsar is not spinning down steadily Mere02 ; Bignami03 ; Zavlin04 . Moreover, the first derivative pulse frequency varies significantly and its sign is also variable in different observations Zavlin04 .
The non-monotonic behavior of its pulse frequency, $`\mathrm{\Delta }\nu `$, is interpreted by three hypotheses: glitch, accretion and binary (with an orbital period of 0.2 to 6 yr), in which the binary hypothesis is somewhat more plausible than the other two Zavlin04 .
This letter analyzes that for binary pulsars with very small mass function, the Roemer time delay in one orbital period cannot be resolved. Such pulsars may thus be treated as โisolatedโ pulsars. However the pulse frequency and frequency derivatives of such pulsars are still affected by the orbital motion at long time scale, which causes anomalies like that of 1E1207.
Different from the orbital period of 0.2 to 6yr Zavlin04 , the one that predicted by this letter is much shorter, 0.5 to 3.3 min, therefore, 1E1207 is an ideal source of gravitational waves.
## II Orbital effect at long time scale
Roemer time delay is the propagation time across the binary orbit, which is given,
$$\frac{z}{c}=\frac{r\mathrm{sin}i}{c}\mathrm{sin}(\omega +f),$$
(1)
where $`c`$ is the speed of light, $`r`$ the distance between the focus and the pulsar, $`f`$ the true anomaly, $`\omega `$ the angular distance of the periastron from the node, and $`i`$ the orbitโs inclination. The orbital motion also causes the change of pulse frequency, $`\mathrm{\Delta }\nu `$,
$$\frac{\mathrm{\Delta }\nu }{\nu }=\frac{๐ฏ๐ง_\mathrm{p}}{c}=K[\mathrm{cos}(\omega +f)+e\mathrm{cos}\omega ],$$
(2)
where $`K2\pi a_\mathrm{p}\mathrm{sin}i/[cP_\mathrm{b}(1e^2)^{1/2}]`$ is the semi-amplitude, $`e`$, $`P_\mathrm{b}`$, $`a_\mathrm{p}`$ are eccentricity, orbital period, and pulsar semi-major axis, respectively.
Small companion mass, $`i`$ or $`P_\mathrm{b}`$ of a binary pulsar may make the Roemer time delay of Eq. (1) unmeasurable. A binary pulsar may thus be treated as an โisolatedโ pulsar. Whereas, following calculation indicates that Eq. (1) and Eq. (2) can still cause long-term effects on such โisolatedโ pulsars.
For a binary pulsar, the time received by the observer (barycentric time) is,
$$t_\mathrm{b}=t_\mathrm{p}+\frac{z}{c},$$
(3)
where $`t_\mathrm{p}`$ is the proper time of the pulsar, and $`z/c`$ is dependent on Kepler equation,
$$Ee\mathrm{sin}E=\overline{M}=\overline{n}t,$$
(4)
where $`\overline{M}`$, $`E`$ and $`\overline{n}`$ are mean anomaly, eccentric anomaly and mean angular velocity, respectively. Note that $`t`$ is the time of periastron passage, which is uniform.
For a true isolated pulsar, we have $`z/c=0`$ in Eq. (3), thus $`t_\mathrm{b}=t_\mathrm{p}`$, which means both $`t_\mathrm{b}`$ and $`t_\mathrm{p}`$ are uniform. But for a binary pulsar system, $`t_\mathrm{b}`$ is no longer uniform, whereas $`t_\mathrm{p}`$ is still uniform.
Therefore, the proper time of the pulsar, $`t_\mathrm{p}`$, can be used to replace the uniform time, $`t`$ of Eq. (4), then we have $`\overline{M}=\overline{n}t_\mathrm{p}`$.
If $`\mathrm{\Delta }\nu `$ of Eq. (2) is averaged over one orbit period by the measured time, $`t_\mathrm{b}`$, then it gives
$$\mathrm{\Delta }\nu =\frac{1}{P_\mathrm{b}}_0^{P_\mathrm{b}}\mathrm{\Delta }\nu ๐t_\mathrm{b}=\frac{1}{P_\mathrm{b}}_0^{P_\mathrm{b}}\mathrm{\Delta }\nu (dt_\mathrm{p}+\frac{\dot{z}}{c}dt_\mathrm{p})$$
$$=\frac{1}{P_\mathrm{b}}_0^{P_\mathrm{b}}\mathrm{\Delta }\nu \frac{\dot{z}}{c}๐t_\mathrm{p}=\frac{X}{P_\mathrm{b}}_0^{P_\mathrm{b}}\mathrm{\Delta }\nu \mathrm{cos}(\omega +E)\dot{E}๐t_\mathrm{p}$$
$$=\frac{XK\nu }{P_\mathrm{b}}\pi (1\frac{e^2}{4})+O(e^4),$$
(5)
where $`X`$ is the projected semi-major axis, $`Xa_\mathrm{p}\mathrm{sin}i/c`$.
In practical observation, an observer averages $`\mathrm{\Delta }\nu `$ from 0 to $`T`$ ($`TP_\mathrm{b}`$) through $`t_\mathrm{b}`$, the time received by observer, without knowing the orbital period, $`P_\mathrm{b}`$, at all. However if the pulsar measured is truly in a binary system, $`P_\mathrm{b}`$ will affect the averaged result, as given in Eq. (5), thus the averaged $`\mathrm{\Delta }\nu `$ given by the observer is
$$\mathrm{\Delta }\nu =\frac{1}{T}_0^T\mathrm{\Delta }\nu dt_\mathrm{b}=\frac{1}{T}(_0^{P_\mathrm{b}}\mathrm{\Delta }\nu dt_\mathrm{b}+\mathrm{}$$
$$+_{P_\mathrm{b}(N1)}^{P_\mathrm{b}N}\mathrm{\Delta }\nu dt_\mathrm{b}+_{P_\mathrm{b}N}^T\mathrm{\Delta }\nu dt_\mathrm{b})=\beta +o(\beta \frac{P_\mathrm{b}}{T}),$$
(6)
where $`\beta XK\nu \pi (1e^2/4)/P_\mathrm{b}`$, and $`N`$ is an integer. Eq. (6) indicates that if a pulsar is in a binary system, then $`\mathrm{\Delta }\nu `$ measured by the observer is actually contaminated by the long-term orbital effect, $`\beta `$.
## III Interpretation of four puzzles and estimation of orbital period of 1E 1207
### III.1 Puzzle 1: pulsar age vs SNR age
If the $`\mathrm{\Delta }\nu `$ of Eq. (5) and Eq. (6) (brackets $`,`$ are ignored hereafter) are unchangeable then the effect actually cannot be measured. However $`\beta `$ contains orbital elements, $`i`$, $`e`$ and $`a`$ (where $`a`$ is the semi-major axis of the orbit, $`a=a_\mathrm{p}M/M_2`$, $`M`$ and $`M_2`$ are the total mass and companion mass respectively, the mass of the pulsar is $`M_1`$) which are long-periodic terms when the Spin-Orbit coupling effect is considered. Therefore, $`\mathrm{\Delta }\nu `$ is a function of time, and the orbital effect induced $`\dot{\nu }_\mathrm{L}`$ can be found by differentiating $`\beta `$ of Eq. (6), as given in detail in the following subsection.
Thus the observational first derivative of the pulse frequency, $`\dot{\nu }_{\mathrm{obs}}`$, is given by
$$\dot{\nu }_{\mathrm{obs}}=\dot{\nu }+\dot{\nu }_L$$
(7)
where $`\dot{\nu }`$ is the intrinsic one, which is caused by magnetic dipole radiation. Thus the following relation can be obtained
$$\frac{2\dot{\nu }_{\mathrm{obs}}}{\nu _{\mathrm{obs}}}=(\frac{2\dot{\nu }}{\nu }+\frac{2\dot{\nu }_L}{\nu })\frac{\nu }{\nu _{\mathrm{obs}}}\frac{2\dot{\nu }}{\nu }+\frac{2\dot{\nu }_L}{\nu }.$$
(8)
Eq. (8) is actually
$$\frac{1}{\tau }=\frac{1}{\tau _\mathrm{p}}+\frac{2\dot{\nu }_L}{\nu }$$
(9)
where $`\tau =200900`$ kyr is the age corresponding to the contaminated spin-down (by the long-term orbital effect). In other words, when the true age of the pulsar equals the age of SNR, $`\tau _\mathrm{p}=330`$ kyr is the true characteristic age of the pulsar.
Putting $`\tau `$ and $`\tau _\mathrm{p}`$ into Eq. (9), one obtains two group solutions corresponding to maximum and minimum magnitude of $`\dot{\nu }_L`$ and $`\dot{\nu }`$, respectively, ($`\dot{\nu }_L=1.2\times 10^{11}\mathrm{Hz}\mathrm{s}^1`$, $`\dot{\nu }=1.3\times 10^{11}\mathrm{Hz}\mathrm{s}^1`$); and ($`\dot{\nu }_L=1.7\times 10^{12}\mathrm{Hz}\mathrm{s}^1`$, $`\dot{\nu }=1.9\times 10^{12}\mathrm{Hz}\mathrm{s}^1`$).
This implies that the magnitude of $`\dot{\nu }`$ and $`\dot{\nu }_L`$ are much larger than that of $`\dot{\nu }_{\mathrm{obs}}`$, since $`\dot{\nu }_L`$ and $`\dot{\nu }`$ nearly cancel each other out. Therefore the age puzzle, $`\tau _\mathrm{p}\tau `$, can also be explained.
### III.2 Puzzle 2: B-field
Section III.1 shows that the measured $`\dot{\nu }_{\mathrm{obs}}`$ may under-estimate the true intrinsic pulse frequency derivative, $`\dot{\nu }`$. Therefore, the B-field $`3\times 10^{12}`$ G inferred from $`\dot{\nu }_{\mathrm{obs}}`$ may be under-estimated also.
The two $`\dot{\nu }_\mathrm{L}`$ obtained through Eq. (9) correspond to two magnetic dipole radiation-induced $`\dot{\nu }`$, and therefore, to two B-field, $`3\times 10^{13}`$ G and $`1\times 10^{13}`$ G, respectively. It is easy for them to reconcile with the high B-field option, $`B=1.6\times 10^{14}`$ G Sanwal02 , i.e., by assuming the magnetic inclination angle, $`\alpha =11^{}`$ and $`\alpha =4^{}`$, respectively.
However, it is very difficult for these two B-fields to reconcile with another option, $`B=8\times 10^{10}`$ G Sanwal02 . Therefore, the B-fields inferred from the true intrinsic spin-down favors that 1E1207 is a magnetar.
### III.3 Puzzle 3: non-monotonic spin-down and estimation of orbital period
In the gravitational two-body problem with spin, each body precesses in the gravitational field of its companion (geodetic precession), with precession velocity of 1 Post-Newtonian order (PN) bo . The Spin-Orbit coupling causes long-periodic variations in the six orbital elements, $`i`$, $`e`$, $`a`$, $`\overline{M}`$, $`\omega `$ and $`\mathrm{\Omega }`$ (longitude of the ascending node) gong . By the definition of $`K`$ and $`X`$, $`\beta `$ of Eq. (6) can be rewritten as,
$$\mathrm{\Delta }\nu =\beta =\frac{GM\nu }{2\pi c^2a}\rho ,$$
(10)
where $`\rho \pi \mathrm{sin}^2i(M_2/M)^2(1e^2/4)/\sqrt{1e^2}`$. According to Eq. (10), $`\beta `$ contains the orbital elements, $`e`$, $`i`$ and $`a`$, which are all long-periodic terms when the Spin-Orbit effect is considered. However, the variation of $`i`$ is much smaller ($`S/L`$ times, $`S`$ and $`L`$ are the spin and orbital angular momentum respectively) than that of $`a`$ and $`e`$ gong . Thus the long-period variation of Eq. (6) can be written in a Taylor series as
$$\mathrm{\Delta }\nu =\beta =\beta _0+\dot{\beta }t+\mathrm{}=\beta _0\beta \frac{\dot{a}}{a}(1\xi )t+\mathrm{},$$
(11)
where
$$\xi \frac{(1e^2)e^2}{2(1+e^2)(1e^2/4)}+\frac{e^2}{1+e^2},$$
(12)
and
$$\frac{\dot{a}}{a}=\frac{GL(1+e^2)}{c^2a^3(1e^2)^{5/2}}(2+\frac{3M_2}{2M_1})(P_yQ_xP_xQ_y)$$
(13)
where $`P_x`$, $`P_y`$, $`Q_x`$, $`Q_y`$ are sine and cosine functions of $`\omega `$ and $`\mathrm{\Omega }`$ gong . The orbital period $`P_\mathrm{b}`$ of a few minutes corresponds to $`\dot{\omega }_{GR}10^5\mathrm{s}^1`$, which corresponds to $`\dot{a}/a10^6\mathrm{s}^1`$. Thus in the observation time span, $`\mathrm{\Delta }t10^2`$ ks, $`a`$ has changed like, one tenth of its period, which means $`\omega `$ has changed by $`\pi /5`$. This actually corresponds to a large variation amplitude in $`\mathrm{\Delta }a`$.
Define $`\delta a|\mathrm{\Delta }a/a|_{\mathrm{max}}`$, then the maximum and minimum $`a`$ of Eq. (10) are $`a_{\mathrm{max}}=a(1+\delta a)`$, and $`a_{\mathrm{min}}=a(1\delta a)`$ respectively. The discrepancy in $`\mathrm{\Delta }\nu `$ corresponds to the error bar of each observation is given,
$$\delta \nu =\frac{GM\nu \rho }{2\pi c^2a}\left(\frac{1}{1\delta a}\frac{1}{1+\delta a}\right)=\beta \frac{2\delta a}{1(\delta a)^2}.$$
(14)
The fact that the amplitude of $`\delta \nu `$ is not much larger than a few $`\mu `$Hz Zavlin04 demands that $`a>|\mathrm{\Delta }a|`$. Thus the maximum $`\delta a`$ can only be like $`\delta a=0.9`$, whereas $`\delta a=1`$ is not allowed.
From the point of view of Eq. (10), both error bars in one observation and discrepancy for different observations are dependent of the variation of $`a`$. The difference is that the discrepancy among different observations, i.e., Jan 2000, Aug 2002, corresponds to a much longer time scale, in which $`\mathrm{\Delta }\nu `$ is modulated by both $`\omega `$ and $`\mathrm{\Omega }`$ (the period of $`\mathrm{\Omega }`$ is comparable to that of $`\omega `$) for many periods. Whereas in one observation ($`10^1`$ ks โ $`10^2`$ ks), the time may be just enough for $`\mathrm{\Delta }\nu `$ to vary in a few periods of $`\omega `$, or even less than a period of $`\omega `$.
The jump of $`\mathrm{\Delta }\nu `$ between Dec 2001 and Jan 2002 Zavlin04 , can be explained by the variation of $`\omega `$ and $`\mathrm{\Omega }`$, which causes relatively sharp variation in $`a`$ and thus significant variation in $`\mathrm{\Delta }\nu `$.
The second term at the right hand side of Eq. (11) actually corresponds to $`\dot{\nu }_L`$, which is given in magnitude as
$$\dot{\nu }_L=\beta \frac{\dot{a}}{a}(1\xi ).$$
(15)
Putting the two $`\dot{\nu }_\mathrm{L}`$ obtained in Eq. (9) into Eq. (15), we have two curves $`\rho `$ vs $`P_\mathrm{b}`$ corresponding to $`\dot{\nu }_{\mathrm{L1}}`$ and $`\dot{\nu }_{\mathrm{L2}}`$ respectively, as shown in Fig. 1.
Similarly putting $`\mathrm{\Delta }\nu _1=0.18\mu `$Hz and $`\mathrm{\Delta }\nu _2=4.2\mu `$Hz into Eq. (6), we have two $`\rho `$ vs $`P_\mathrm{b}`$ curves as shown in Fig. 1, which correspond to minimum and maximum discrepancies in $`\nu `$ or error bars in different observations of Zavlin et al Zavlin04 .
The maximum orbital period, 3.3 min, is given by the cross section of $`\mathrm{\Delta }\nu _2`$ and $`\dot{\nu }_{\mathrm{L2}}`$ at $`C`$ as shown in Fig. 1. The minimum orbital period is 0.1 min corresponding to $`A`$ given by $`\mathrm{\Delta }\nu _1`$ and $`\dot{\nu }_{\mathrm{L1}}`$.
In the area ABCD of Fig. 1, a point, i.e., with $`P_\mathrm{b}=0.7`$ min and $`\rho =0.008`$ can be found. Assuming $`M_1=1.4M_{}`$, $`M_2=0.2M_{}`$ and $`e=0`$, $`\rho 0.09\mathrm{sin}^2i`$ is obtained, and by the definition of $`\rho `$, $`\mathrm{sin}i0.3`$ can be obtained. In turn $`X`$ is given, $`X2.6`$ ms, which is smaller than the time resolution of observation, 5.7 ms or 2.9 ms Zavlin04 . Therefore, the modulation induced by the orbital motion may not be detected from the timing observation. This is consistent with the fact that the side band corresponding to $`P_\mathrm{b}`$ of a few minutes has not been found in 1E1207.
Therefore, the companion of 1E1207 should have low mass and be compact enough. Because the undetected orbital modulation implies that $`X`$ must be small; and the short orbital period demands that the companion be compact object like low-mass neutron star carr03 or strange star xu03 ; xu05 . A white dwarf star companion is unlikely due to the separation of the two stars is almost equal to the radius of a white dwarf star when the orbital period is of 1 min.
### III.4 Puzzle 4: magnitude and sign of $`\dot{\nu }_{\mathrm{obs}}`$
The values of $`P_\mathrm{b}`$, $`\rho `$ and the period of $`\omega `$ corresponding to the four points ABCD of Fig. 1 are shown in Table 1.
As given by Eq. (10) and Eq. (15), both $`\mathrm{\Delta }\nu `$ and $`\dot{\nu }_\mathrm{L}`$ vary with $`a`$ which is in turn modulated by the period, $`T_\omega `$ (and $`T_\mathrm{\Omega }`$) as shown in Table 1. In the case $`P_\mathrm{b}=1`$ min, the period of $`T_\omega `$ is $`25`$ days, therefore the period of variation of $`\dot{\nu }_\mathrm{L}`$ is approximately $`25`$ days also (recall the period, $`T_\mathrm{\Omega }`$, is comparable with $`T_\omega `$).
On the other hand, the intrinsic $`\dot{\nu }`$ changes steadily, which means $`\dot{\nu }`$ and $`\dot{\nu }_\mathrm{L}`$ some times cancelling out, and some times have the same sign and enhancing, thus $`\dot{\nu }_{\mathrm{obs}}`$ can both be $`10^{14}\mathrm{s}^2`$ ($`\dot{\nu }`$ and $`\dot{\nu }_\mathrm{L}`$ cancelled out); and $`10^{11}\mathrm{s}^2`$ ($`\dot{\nu }`$ and $`\dot{\nu }_\mathrm{L}`$ enhanced). This well explains the observations of Zavlin et al Zavlin04 , which show that $`\dot{\nu }_{\mathrm{obs}}`$ can have very different magnitude, $`10^{14}\mathrm{s}^2`$ and $`10^{11}\mathrm{s}^2`$ and its sign is also changeable at different epochs.
Eq. (10) and Eq. (15) actually predict that $`\mathrm{\Delta }\nu `$ and $`\dot{\nu }_\mathrm{L}`$ can vary with periods of days, thus $`\dot{\nu }_{\mathrm{obs}}`$ can change sign in order of days, or even during one observation ($`10^2`$ ks). Comparing Table 1 with observation may extract the period $`T_\omega `$ and $`T_\mathrm{\Omega }`$ and therefore determine the orbital period $`P_\mathrm{b}`$.
## IV Discussion
Therefore, all four puzzles can be explained naturally by an ultra-compact binary system.
The best spectral model describes the continuum as the sum of two blackbody curves with $`kT=0.211\pm 0.001`$ kev, for an emitting radius $`R=2.95\pm 0.05`$ km; and $`kT=0.40\pm 0.02`$ kev ($`R=250\pm 50`$ m) Bignami03 . It is possible that these two emitting radii are from the hot spot of two stars, 1E1207 and its companion.
The characteristic amplitude of gravitational waves from a binary system is Thorne
$$h=1.4\times 10^{20}(\frac{\mu }{M_{}})(\frac{M}{M_{}})^{2/3}(\frac{P_\mathrm{b}}{1\mathrm{hr}})^{2/3}(\frac{d}{100\mathrm{pc}})^1f(e)$$
$$3\times 10^{21},$$
(16)
where $`\mu `$ is the reduced mass which equals $`\mu =0.18M_{}`$ when $`M_1=1.4M_{}`$ and $`M_2=0.2M_{}`$; $`d=2`$ kpc is the distance; $`P_\mathrm{b}=1`$ min is orbital period; and $`f(e)`$ is given by $`f(e)=(1+\frac{73}{24}e^2+\frac{37}{96}e^4)/(1e^2)^{7/2}`$, which is assumed $`f(e)1`$. In such case, $`GM/(c^2a)(v/c)^29\times 10^5`$ ($`v`$ is a characteristic orbital velocity), which means equations, Eq. (11)โEq. (15), based on the Post-Newtonian approximation are good enough to describe the dynamics of the ultra-compact binary system.
The time scale of coalescing of 1E1207 corresponding to $`P_\mathrm{b}=0.5`$ min is $`27`$ yr. In order to be consistent with the fact that it was discovered in 1984 and is still there, it is necessary that $`P_\mathrm{b}0.5`$ min. Therefore the most probable orbital period of 1E1207 is (0.5โ3.3) min.
The low wave frequency, $`10^2`$ Hz, and the extremely large wave amplitude means that 1E1207 is an ideal source for the space detector Laser Interferometer Space Antenna.
###### Acknowledgements.
I thank A. Rรผdiger and T. Kiang for very helpful comments and corrections on the presentation of this letter. I thank T.Y. Huang and X.S. Wan for useful comments on the mathematics in this letter.
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# Symmetries of modules of differential operators
## 1 Introduction
We study the space of linear differential operators acting in the space of tensor densities on $`S^1`$ as a module over the group $`\mathrm{Diff}(S^1)`$ of all diffeomorphisms of $`S^1`$. More precisely, let $`๐_{\lambda ,\mu }^k(S^1)`$ be the space of linear $`k`$-th order differential operators
$$A:_\lambda (S^1)_\mu (S^1)$$
where $`_\lambda (S^1)`$ and $`_\mu (S^1)`$ are the spaces of tensor densities of degree $`\lambda `$ and $`\mu `$ respectively. We compute the commutant of the $`\mathrm{Diff}(S^1)`$-action on $`๐_{\lambda ,\mu }^k(S^1)`$. This commutant is an associative algebra which we denote $`_{\lambda ,\mu }^k(S^1)`$ and call the algebra of symmetries.
1.1 This paper is closely related to the classical subject initiated by Veblen in his talk at IMC in 1928, namely the study of invariant operators also called natural operators (cf. ). An operator is called invariant if it commutes with the action of the group of diffeomorphisms. The main two examples are the classic de Rham differential of differential forms
$$d:\mathrm{\Omega }_k(M)\mathrm{\Omega }_{k+1}(M),$$
where $`M`$ is a smooth manifold, and the integral
$$:\mathrm{\Omega }_n(M)$$
provided $`M`$ is compact of dimension $`n`$.
Usually, one considers differential operators acting on various spaces of tensor fields on a smooth manifold. A famous theorem states that the de Rham differential is, actually, the only invariant differential operator in one argument acting on the spaces of tensor fields. This result was conjectured by Schouten and proved independently and using different approaches by Rudakov , Kirillov and Terng , for a complete historical account see .
Many classification results are available now, and it was shown that there are quite few invariant differential operators and most of them are of a great importance. For instance, bilinear invariant differential operators on tensor fields were classified by Grozman . The complete list of such operators contains well-known examples, such as the Poisson, Schouten and Nijenhuis brackets, and one exceptional bilinear third-order differential operator
$$G:_{\frac{2}{3}}(S^1)_{\frac{2}{3}}(S^1)_{\frac{5}{3}}(S^1).$$
(1.1)
Note that differential operators invariant with respect of the diffeomorphism groups can be interpreted in terms of the Lie algebras of vector fields. This viewpoint relates the subject with the Gelfand-Fuchs cohomology, see and references therein.
1.2 The main difference of our work from the classic literature is that we consider linear operators acting on differential operators (instead of tensor fields). More precisely, we classify the linear maps
$$T:๐_{\lambda ,\mu }^k(S^1)๐_{\lambda ,\mu }^k(S^1)$$
(1.2)
commuting with the $`\mathrm{Diff}(S^1)`$-action. The module of differential operators $`๐_{\lambda ,\mu }^k(S^1)`$ is not isomorphic to any module of tensor fields (but rather rasembles $`\mathrm{gl}(_\lambda )`$ or $`_\lambda ^{}_\mu `$). The problem of classification of $`\mathrm{Diff}(S^1)`$-invariant operators on $`๐_{\lambda ,\mu }^k(S^1)`$ is, therefore, different from Veblenโs problem, although similar.
A well-known example of a map (1.2) is the conjugation of differential operators. This map associates to an operator $`A`$ the adjoint operator $`A^{}`$. If $`A๐_{\lambda ,\mu }^k(S^1)`$, then $`A^{}๐_{1\mu ,1\lambda }^k(S^1)`$, so that this map defines a symmetry if and only if $`\lambda +\mu =1`$.
Let us emphasize that, unlike Rudakov-Kirillov-Terng, we consider not only differential (or local) symmetries of $`๐_{\lambda ,\mu }^k(S^1)`$ but also non-local ones. For instance, we find a version of trace which is an analog of the Adler trace (see ).
1.3 The main purpose of this paper is to show that some modules $`๐_{\lambda ,\mu }^k(S^1)`$ are particular and very interesting. It turns out that the algebra of symmetries $`_{\lambda ,\mu }^k(S^1)`$ can be quite rich depending on the values of $`\lambda `$ and $`\mu `$ as well as on $`k`$. This algebra is an important characteristic of the corresponding module which embraces those given in .
The first example which is particular (for every $`k`$) is the module $`๐_{0,1}^k(S^1)`$ of operators from the space of functions to the space of 1-forms. The algebra of symmetries in this case is always of maximal dimension.
Another interesting module is $`๐_{\frac{1}{2},\frac{3}{2}}^2(S^1)`$. It is related to the famous Virasoro algebra and also to the projective differential geometry, see . This module appears as a particular case in our classification.
Further intriguing examples of modules of differential operators are $`๐_{\frac{2}{3},\frac{5}{3}}^3(S^1)`$ and $`๐_{\frac{2}{3},\frac{5}{3}}^4(S^1)`$. These modules are related to the Grozman operator (1.1). The algebraic and geometric meaning of these modules is not known.
1.4 We will also consider symmetries of differential operators acting in the space of $`\lambda `$-densities over $``$ and compare this case with the case of $`S^1`$. The classification of the invariant differential operators (1.2) remains the same as that on $`S^1`$, except that there are no non-local symmetries.
Let us also mention that symmetries of the modules of differential operators in the multi-dimensional case have been classified in . In this case the algebra of symmetries is smaller.
1.5 This paper is organized as follows. In Section 2 we introduce the $`\mathrm{Diff}(S^1)`$-modules of differential operators and their symbols. In Section 3 we formulate the classification theorems. In Section 4 we give an explicit construction of all $`\mathrm{Diff}(S^1)`$-invariant linear maps on the modules $`๐_{\lambda ,\mu }^k(S^1)`$ in all possible cases. These operators are our main characters; some of them are known and some seem to be new. In Section 5 we calculate the associative algebras of invariant operators and identify them with some associative algebras of matrices. This gives a complete description of the symmetry algebras $`_{\lambda ,\mu }^k(S^1)`$. Finally, in Section 6 we prove that there are no other invariant operators on $`๐_{\lambda ,\mu }^k(S^1)`$ than the operators we introduce. This completes the proof of the main theorems.
## 2 The main definitions
In this section we define the space of differential operators on $`S^1`$ acting on the space of densities and the corresponding space of symbols. We pay particular attention to the action of the group of diffeomorphisms $`\mathrm{Diff}(S^1)`$ and of the Lie algebra of vector fields $`\mathrm{Vect}(S^1)`$ on these spaces.
### 2.1 Densities on $`S^1`$
Denote by $`_\lambda (S^1)`$, or $`_\lambda `$ for short, the space of $`\lambda `$-densities on $`S^1`$
$$\phi =\varphi (x)(dx)^\lambda ,$$
where $`\lambda `$ is the degree (or weight), $`x`$ is a local coordinate on $`S^1`$ and $`\varphi (x)C^{\mathrm{}}(S^1)`$. As a vector space, $`_\lambda `$ is isomorphic to $`C^{\mathrm{}}(S^1)`$.
The group $`\mathrm{Diff}(S^1)`$ naturally acts on $`_\lambda `$. If $`f\mathrm{Diff}(S^1)`$, then
$$\rho _{f^1}^\lambda :\varphi (x)(dx)^\lambda \left(f^{}\right)^\lambda \varphi (f(x))(dx)^\lambda .$$
The $`\mathrm{Diff}(S^1)`$-modules $`_\lambda `$ and $`_\mu `$ are not isomorphic unless $`\lambda =\mu `$ (cf. ).
###### Example 2.1.
The space $`_0`$ is isomorphic to $`C^{\mathrm{}}(S^1)`$, as a $`\mathrm{Diff}(S^1)`$-module; the space $`_1`$ is nothing but the space of 1-forms (volume forms); the space $`_1`$ is the space of vector fields on $`S^1`$.
The space $`_\lambda `$ can also be viewed as the space of functions on the cotangent bundle $`T^{}S^1S^1`$ (with zero-section removed) homogeneous of degree $`\lambda `$. In standard (Darboux) coordinates $`(x,\xi )`$ on $`T^{}S^1`$ one writes:
$$\varphi (x)(dx)^\lambda \varphi (x)\xi ^\lambda .$$
(2.3)
This identification commutes with the $`\mathrm{Diff}(S^1)`$-action.
### 2.2 Invariant pairing
There is a pairing $`_\lambda _{1\lambda }`$ given by
$$\varphi (x)(dx)^\lambda ,\psi (x)(dx)^{1\lambda }=_{S^1}\varphi (x)\psi (x)๐x$$
which is $`\mathrm{Diff}(S^1)`$-invariant. For instance, the space $`_{\frac{1}{2}}`$ is equipped with a scalar product; this is a natural pre-Hilbert space that is popular in geometric quantizaton.
### 2.3 Differential operators on densities
Consider the space of linear differential operators
$$A:_\lambda _\mu $$
with arbitrary $`\lambda ,\mu `$. This space will be denoted by $`๐_{\lambda ,\mu }(S^1)`$. The subspace of differential operators of order $`k`$ will be denoted by $`๐_{\lambda ,\mu }^k(S^1)`$.
Fix a (local) coordinate $`x`$, a differential operator $`A๐_{\lambda ,\mu }^k(S^1)`$ is of the form
$$A=a_k(x)\frac{d^k}{dx^k}+a_{k1}(x)\frac{d^{k1}}{dx^{k1}}+\mathrm{}+a_0(x),$$
(2.4)
where $`a_i(x)`$ are smooth functions. More precisely,
$$A(\phi )=\left(a_k(x)\frac{d^k\varphi (x)}{dx^k}+a_{k1}(x)\frac{d^{k1}\varphi (x)}{dx^{k1}}+\mathrm{}+a_0(x)\varphi (x)\right)(dx)\mu ,$$
where $`\phi =\varphi (x)(dx)^\lambda `$.
###### Example 2.2.
The space $`๐_{\lambda ,\mu }^0(S^1)`$ is nothing but $`_{\mu \lambda }`$. Any zeroth-order differential operator is the operator of multiplication by a $`(\mu \lambda )`$-density:
$$a(x)(dx)^{\mu \lambda }:\varphi (x)(dx)^\lambda a(x)\varphi (x)(dx)^\mu $$
### 2.4 $`\mathrm{Diff}(S^1)`$\- and $`\mathrm{Vect}(S^1)`$-module structure
The space $`๐_{\lambda ,\mu }(S^1)`$ is a $`\mathrm{Diff}(S^1)`$-module with respect to the action
$$\rho _f^{\lambda ,\mu }(A)=\rho _f^\mu A\rho _{f^1}^\lambda ,$$
where $`f\mathrm{Diff}(S^1)`$.
We will also consider the Lie algebra of vector fields $`\mathrm{Vect}(S^1)`$ and the natural $`\mathrm{Vect}(S^1)`$-action on $`๐_{\lambda ,\mu }^k(S^1)`$. A vector field $`X=X(x)\frac{d}{dx}`$ acts on the space of tensor densities $`_\lambda `$ by Lie derivative
$$L_X^\lambda (\phi )=\left(X(x)\varphi ^{}(x)+\lambda X^{}(x)\varphi (x)\right)(dx)^\lambda .$$
The action of $`\mathrm{Vect}(S^1)`$ on the space of differential operators is given by the commutator
$$_X^{\lambda ,\mu }(A)=L_X^\mu AAL_X^\lambda .$$
(2.5)
Note that the above formulรฆ are independent of the choice of the local coordinate $`x`$.
#### 2.4.1 Example: the module $`๐_{\lambda ,\mu }^1(S^1)`$
The space $`๐_{\lambda ,\mu }^1(S^1)`$ is split into a direct sum
$$๐_{\lambda ,\mu }^1(S^1)_{\mu \lambda 1}_{\mu \lambda }$$
(2.6)
as a $`\mathrm{Diff}(S^1)`$\- (and $`\mathrm{Vect}(S^1)`$-) module.
Indeed, every first-order differential operator $`A๐_{\lambda ,\mu }^1(S^1)`$
$$A\left(\varphi (x)(dx)^\lambda \right)=\left(a_1(x)\varphi ^{}(x)+a_0(x)\varphi (x)\right)(dx)^\mu $$
can be rewritten in the form
$$A\left(\varphi (dx)^\lambda \right)=\left(a_1\varphi ^{}+\lambda a_1^{}\varphi +(a_0\lambda a_1^{})\varphi \right)(dx)^\mu ,$$
and, finally, one obtaines an intrinsic expression
$$A(\phi )=\left(L_{a_1}\phi +(a_0\lambda a_1^{})\phi \right)(dx)^{\mu \lambda },$$
where $`a_1=a_1(x)\frac{d}{dx}`$ is understood a vector field and $`a_0(x)\lambda a_1^{}(x)`$ as a function.
Furthermore, using identification (2.3), one can write a more elegant formula for a first-order operator:
$$A=\xi ^{\mu \lambda }L_{a_1}+\mathrm{div}a_1\frac{}{\xi }+a_0.$$
Note that there are no intrinsic formulรฆ similar to the above ones in the case of modules $`๐_{\lambda ,\mu }^k(S^1)`$ with $`k2`$, and, in general, there are no splittings similar to (2.6). The geometric meaning of the modules $`๐_{\lambda ,\mu }^k(S^1)`$ was discussed in .
### 2.5 Space of symbols of differential operators
The filtration
$$๐_{\lambda ,\mu }^0(S^1)๐_{\lambda ,\mu }^1(S^1)\mathrm{}๐_{\lambda ,\mu }^k(S^1)\mathrm{}$$
is preserved by the $`\mathrm{Diff}(S^1)`$-action. The graded $`\mathrm{Diff}(S^1)`$-module $`๐ฎ_{\lambda ,\mu }(S^1)=\mathrm{gr}(๐_{\lambda ,\mu }(S^1))`$ is called the module of symbols of differential operators.
The quotient module $`๐_{\lambda ,\mu }^k(S^1)/๐_{\lambda ,\mu }^{k1}(S^1)`$ is isomorphic to the module of tensor densities $`_{\mu \lambda k}(S^1)`$; the isomorphism is provided by the principal symbol. As a $`\mathrm{Diff}(S^1)`$-module, the space of symbols depends, therefore, only on the difference
$$\delta =\mu \lambda ,$$
so that $`๐ฎ_{\lambda ,\mu }(S^1)`$ can be denoted as $`๐ฎ_\delta (S^1)`$, and finally we have
$$๐ฎ_\delta (S^1)=\underset{i=0}{\overset{\mathrm{}}{}}_{\delta i}$$
as $`\mathrm{Diff}(S^1)`$-modules.
The space of symbols $`๐ฎ_\delta (S^1)`$ can also be viewed as the space of functions on $`T^{}S^1S^1`$. Namely, any $`k`$-th order symbol $`P๐ฎ_\delta (S^1)`$ can be written in the form
$$P=a_k(x)\xi ^{k\delta }+a_{k1}(x)\xi ^{k\delta 1}+\mathrm{}+a_0(x)\xi ^\delta .$$
(2.7)
The natural lift of the action of $`\mathrm{Diff}(S^1)`$ to $`T^{}S^1`$ equips the space of functions (2.7) with a structure of $`\mathrm{Diff}(S^1)`$-module; this action coincides with the $`\mathrm{Diff}(S^1)`$-action on the space $`๐ฎ_\delta (S^1)`$.
###### Remark 2.3.
The spaces $`๐_{\lambda ,\mu }(S^1)`$ and $`๐ฎ_\delta (S^1)`$ are not isomorphic as $`\mathrm{Diff}(S^1)`$-modules. There are cohomology classes which are obstructions for existence of such an isomorphism, see .
## 3 The main results
In this section we formulate the main results of this paper and give a complete description of the algebras of symmetries $`_{\lambda ,\mu }^k(S^1)`$. We defer the proofs to Sections 4-6.
We will say that the algebra $`_{\lambda ,\mu }^k(S^1)`$ is trivial if it is generated by the identity map $`\mathrm{Id}`$, and therefore is isomorphic to $``$. Of course, we are interested in the cases when this algebra is non-trivial.
### 3.1 Introducing four algebras of matrices
We will need the following associative algebras.
1. The commutative algebra $`\mathrm{}`$ will be denoted by $`^n`$. Of course, this algebra can be represented by diagonal $`n\times n`$-matrices.
2. The algebra of (lower) triangular $`(n\times n)`$-matrices will be denoted by $`๐ฑ_n`$.
3. The commutative algebra $`๐`$ of $`(2\times 2)`$-matrices of the form
$$\left(\begin{array}{cc}a& 0\\ b& a\end{array}\right).$$
4. The algebra $`๐`$ of $`(4\times 4)`$-matrices of the form
$$\left(\begin{array}{cccc}a& 0& 0& d\\ 0& a& 0& 0\\ 0& c& b& 0\\ 0& 0& 0& b\end{array}\right)$$
It turns out that the algebras of symmetries $`_{\lambda ,\mu }^k(S^1)`$ are always direct sums of the above matrix algebras.
In Appendix we will introduce natural generators of the algebras $`๐`$, $`๐`$ and $`^n`$.
### 3.2 Stability: the case $`k5`$
We start our list of classification theorems with the โstableโ case. Namely, if $`k5`$, then the algebras $`_{\lambda ,\mu }^k(S^1)`$ do not depend on $`k`$.
###### Theorem 3.1.
For $`k5`$, the algebra $`_{\lambda ,\mu }^k(S^1)`$ is trivial for all $`(\lambda ,\mu )`$, except
1. $`_{\lambda ,\mu }^k(S^1)^2`$, for $`\{\begin{array}{cc}\lambda +\mu =1,\hfill & \lambda 0,\hfill \\ \lambda =0,\hfill & \mu 1,\mathrm{\hspace{0.17em}0}\hfill \\ \mu =1,\hfill & \lambda 1,\mathrm{\hspace{0.17em}0}\hfill \end{array}`$
2. $`_{0,0}^k(S^1)_{1,1}^k(S^1)^3`$;
3. $`_{0,1}^k(S^1)๐^2`$.
The exceptional modules $`๐_{\lambda ,\mu }^k(S^1)`$ with $`k5`$ are represented in Figure 1.
We will provide in the sequel a list of generators of the symmetry algebra for each non-trivial case, we will also give its explicit identification with the corresponding algebra of matrices.
### 3.3 Modules of differential operators of order 4
The result for the modules $`๐_{\lambda ,\mu }^k(S^1)`$ with $`k4`$ is different from the above one. (It is interesting to compare this property of differential operators with that in the case of algebraic equations.)
Consider the modules of operators of order $`k=4`$. The complete classification of symmetries in this case is given by the following result.
###### Theorem 3.2.
The algebra $`_{\lambda ,\mu }^4(S^1)`$ is trivial for all $`(\lambda ,\mu )`$ except
1. $`_{\lambda ,\mu }^4(S^1)^2,`$ for $`\{\begin{array}{cc}\lambda +\mu =1,\hfill & \lambda 0,\frac{2}{3}\hfill \\ \lambda =0,\hfill & \mu 3,\frac{5}{4},\mathrm{\hspace{0.17em}1},\mathrm{\hspace{0.17em}0}\hfill \\ \mu =1,\hfill & \lambda 1,\mathrm{\hspace{0.17em}0},\frac{1}{4},2\hfill \end{array}`$
2. $`_{\lambda ,\mu }^4(S^1)^3,`$ for $`(\lambda ,\mu )=(1,1),(0,\frac{5}{4}),(0,0),(\frac{1}{4},1),(\frac{2}{3},\frac{5}{3})`$;
3. $`_{0,3}^4(S^1)_{2,1}^4(S^1)๐`$;
4. $`_{0,1}^4(S^1)๐^2`$.
The exceptional modules $`๐_{\lambda ,\mu }^4(S^1)`$ are represented in Figure 2.
###### Remark 3.3.
We will show in Section 5.4.1 that the module $`๐_{\frac{2}{3},\frac{5}{3}}^4(S^1)`$ is, indeed, a very special one. This exceptional module is related to the Grozman operator (1.1).
### 3.4 Modules of differential operators of order 3
Symmetries of the modules of third-order operators are particularly rich.
###### Theorem 3.4.
The algebra $`_{\lambda ,\mu }^3(S^1)`$ is trivial for all $`(\lambda ,\mu )`$ except
1. $`_{\lambda ,\mu }^3(S^1)^2`$, for $`\{\begin{array}{cc}\lambda +\mu =1,\hfill & \lambda 0,\frac{1}{2},\frac{2}{3}\hfill \\ (3\lambda +1)(3\mu 4)=1,\hfill & \lambda 0,\frac{2}{3}\hfill \end{array}`$
2. $`_{\lambda ,\mu }^3(S^1)๐`$, for $`\mu \lambda =2`$, $`\lambda 0,\frac{1}{2},1`$;
3. $`_{\lambda ,\mu }^3(S^1)^3`$, for $`\{\begin{array}{cc}\lambda =0,\hfill & \mu 3,\mathrm{\hspace{0.17em}2},\mathrm{\hspace{0.17em}1}\hfill \\ \mu =1,\hfill & \lambda 0,1,2\hfill \end{array}`$
4. $`_{\lambda ,\mu }^3(S^1)๐`$, for $`(\lambda ,\mu )=(0,3),(0,2),(1,1),(2,1)`$;
5. $`_{\frac{1}{2},\frac{3}{2}}^3(S^1)๐ฑ_2`$;
6. $`_{\frac{2}{3},\frac{5}{3}}^3(S^1)^3`$;
7. $`_{0,1}^3(S^1)๐^2`$.
The exceptional modules $`๐_{\lambda ,\mu }^3(S^1)`$ are represented in Figure 3.
### 3.5 Modules of second-order differential operators
Second-order differential operators are definitely among the most popular objects of mathematics. Their invariants with respect to the action of $`\mathrm{Diff}(S^1)`$, such as monodromy or the rotation number, were thoroughly studied. Some modules of second-order differential operators on $`S^1`$ have geometric meaning, they have been related to various algebraic structures such as the Virasoro algebra and integrable systems.
The result in the second-order case is as follows.
###### Theorem 3.5.
The algebra $`_{\lambda ,\mu }^2(S^1)`$ is isomorphic to $`^2`$ for all $`(\lambda ,\mu )`$, except
1. $`_{\lambda ,\mu }^2(S^1)๐`$, for $`\{\begin{array}{cc}\mu \lambda =1,\hfill & \lambda 0\hfill \\ \mu \lambda =2,\hfill & \lambda 0,\frac{1}{2},1\hfill \end{array}`$
2. $`_{\lambda ,\mu }^2(S^1)^3`$, for $`\{\begin{array}{cc}\lambda =0,\hfill & \mu 2,\mathrm{\hspace{0.17em}1}\hfill \\ \mu =1,\hfill & \lambda 0,1\hfill \end{array}`$
3. $`_{\frac{1}{2},\frac{3}{2}}^2(S^1)๐ฑ_2`$;
4. $`_{0,2}^2(S^1)_{1,1}^2(S^1)๐`$;
5. $`_{0,1}^2(S^1)๐`$.
The exceptional modules of second-order operators are represented in Figure 4.
### 3.6 Modules of first-order differential operators
Let us finish this section with the result in the first-order case. The result is less interesting, the only particular module is $`๐_{0,1}^1(S^1)`$.
###### Theorem 3.6.
The algebra of $`_{\lambda ,\mu }^1(S^1)`$ is isomorphic to $`^2`$ for all $`\lambda ,\mu `$, except in the following cases:
1. $`_{\lambda ,\mu }^1(S^1)๐`$, for $`\mu \lambda =1`$, $`\lambda 0`$;
2. $`_{0,1}^1(S^1)๐`$.
For the sake of completeness, let us also mention the zeroth-order case: the algebra of symmetries $`_{\lambda ,\mu }^0(S^1)`$ is trivial. Indeed, the module $`๐_{\lambda ,\mu }^0(S^1)`$ is isomorphic to $`_{\mu \lambda }`$.
### 3.7 Non-compact case: differential operators on $``$
The algebras of symmetries of the modules of differential operators in the spaces of densities over $``$ can be different from that on $`S^1`$. This occurs in the โmost particularโ case $`(\lambda ,\mu )=(0,1)`$.
###### Theorem 3.7.
(i) If $`(\lambda ,\mu )(0,1)`$, then the algebra $`_{\lambda ,\mu }^k()`$ coincides with $`_{\lambda ,\mu }^k(S^1)`$.
(ii) In the exceptional case $`(\lambda ,\mu )=(0,1)`$ one has:
1. If $`k3`$, then $`_{0,1}^k()๐ฑ_2^2`$;
2. If $`k=2`$, then $`_{0,1}^k()๐ฑ_2`$;
3. If $`k=1`$, then $`_{0,1}^k()๐ฑ_2`$.
## 4 Construction of symmetries
In this section we give an explicit construction of the generators of algebras $`_{\lambda ,\mu }^k(S^1)`$ for every case where this algebra is non-trivial. We will prove in Section 6 that our list of invariant differential operators is complete.
### 4.1 The conjugation
The best known invariant map between the spaces of differential operators is the conjugation. It is a linear map
$$C:๐_{\lambda ,\mu }^k(S^1)๐_{1\mu ,1\lambda }^k(S^1)$$
that associates to each operator $`A`$ its adjoint $`A^{}`$ defined by
$$_{S^1}A^{}(\phi )\psi =_{S^1}\phi A(\psi )$$
for every $`\phi _{1\mu }`$ and $`\psi _\lambda `$. It follows that the modules $`๐_{\lambda ,\mu }^k(S^1)`$ with $`\lambda `$ and $`\mu `$ satisfying the condition
$$\lambda +\mu =1$$
have non-trivial symmetries.
The conjugation map $`C`$ is an involution; the straight line $`\lambda +\mu =1`$ will play a role of symmetry axis in the plane parameterized by $`(\lambda ,\mu )`$.
For an arbitrary local parameter $`x`$ on $`S^1`$, the conjugation map is given by the well-known formula
$$C:\underset{i=0}{\overset{k}{}}a_i(x)\frac{d^i}{dx^i}\underset{i=0}{\overset{k}{}}(1)^i\left(\frac{d}{dx}\right)^ia_i(x)$$
(4.8)
that easily follows from the definition.
###### Remark 4.1.
The expression (4.8) is independent from the choice of the parameter $`x`$. Indeed, any change of local coordinates is given by a diffeomorphism of $`S^1`$. Note that this fundamental property of coordinate independence is just a different way to express the $`\mathrm{Diff}(S^1)`$-equivariance.
### 4.2 The cases $`\lambda =0`$ and $`\mu =1`$
We will define a $`\mathrm{Diff}(S^1)`$-invariant operator
$$P_0:๐_{0,\mu }^k(S^1)๐_{0,\mu }^k(S^1).$$
Let us first consider a $`\mathrm{Diff}(S^1)`$-invariant projection $`P_0:๐_{0,\mu }^k(S^1)_\mu `$ defined by: $`AA(1),`$ where $`1_0C^{\mathrm{}}(M)`$ is a constant function on $`S^1`$. In other words,
$$P_0\left(\underset{i=0}{\overset{k}{}}a_i(x)\frac{d^i}{dx^i}\right)=a_0(x)\left(dx\right)^\mu .$$
(4.9)
Since $`_\mu ๐_{0,\mu }^k(S^1)`$, one obtains a non-trivial element of the algebra $`_{0,\mu }^k(S^1)`$.
Thanks to the conjugation map (4.8), one also has a non-trivial symmetry $`P_0^{}=CP_0C`$
$$P_0^{}:๐_{\lambda ,1}^k(S^1)๐_{\lambda ,1}^k(S^1).$$
The explicit formula follows from (4.9) and (4.8):
$$P_0^{}\left(\underset{i=0}{\overset{k}{}}a_i(x)\frac{d^i}{dx^i}\right)=\underset{i=0}{\overset{k}{}}(1)^ia_i(x)^{(i)}.$$
(4.10)
The right hand side is understood as a (scalar) differential operator from $`_\lambda `$ to $`_1`$.
### 4.3 Two additional elements of $`_{0,1}^k(S^1)`$
In the most particular case $`(\lambda ,\mu )=(0,1)`$, there are two more elements of the symmetry algebra.
* There is a non-local element of $`_{0,1}^1(S^1)`$. It is given by the expression
$$L\left(\underset{i=0}{\overset{k}{}}a_i(x)\frac{d^i}{dx^i}\right)=\left(_{S^1}a_0(x)๐x\right)d$$
(4.11)
where $`d`$ is the de Rham differential. Indeed, the projection (4.9) gives a 1-form on $`S^1`$ so that the integral above is well-defined; since $`d๐_{0,1}^1(S^1)`$, we can understand $`L`$ as a linear map $`L:๐_{0,1}^k(S^1)๐_{0,1}^k(S^1)`$ and therefore a symmetry.
###### Remark 4.2.
Equation (4.11) is an analog of the well-known Adler trace , although the latter is defined on the space of pseudodifferential operators from $`_0`$ to $`_0`$.
* There is one more element of the algebra $`_{0,1}^1(S^1)`$ given by the formula
$$P_1\left(\underset{i=0}{\overset{k}{}}a_i(x)\frac{d^i}{dx^i}\right)=\left(\underset{i=1}{\overset{k}{}}(1)^{i1}a_i(x)^{(i1)}\right)d.$$
(4.12)
It is easy to check directly that $`P_1`$ is $`\mathrm{Diff}(S^1)`$-invariant, but its intrinsic form can also be written. In the one-dimensional case, the de Rham differential $`d`$ is an element of $`๐_{0,1}^1(S^1)`$. One has a $`\mathrm{Diff}(S^1)`$-invariant operator
$$\delta :๐_{1,\mu }^k(S^1)๐_{0,\mu }^{k+1}(S^1)$$
given by right composition with the de Rham differential: $`\delta :AAd.`$ This map is a bijection between $`๐_{1,\mu }^k(S^1)`$ and $`\mathrm{Ker}P_0๐_{0,\mu }^{k+1}(S^1)`$. One has:
$$P_1=\delta P_0C\delta ^1(\mathrm{Id}P_0).$$
### 4.4 Additional elements of $`_{0,0}^k(S^1)`$ and $`_{1,1}^k(S^1)`$
The algebra $`_{0,0}^k(S^1)`$ is generated by the operator $`P_0`$ given by (4.9) and
$$S=C\delta ^1(\mathrm{Id}P_0)C\delta C.$$
(4.13)
One can check that the explicit formula for this operator is as follows:
$$S\left(\underset{i=0}{\overset{k}{}}a_i(x)\frac{d^i}{dx^i}\right)=\underset{i=0}{\overset{k1}{}}(1)^i\left(\frac{d}{dx}\right)^i\left(a_i(x)+a_{i+1}^{}(x)\right)+(1)^k\left(\frac{d}{dx}\right)^ka_k(x).$$
The algebra $`_{1,1}^k(S^1)`$ is generated by $`P_0^{}`$ given by (4.10) and the operator $`S^{}=CSC`$.
### 4.5 Symmetries and bilinear operators on tensor densities
We now give a general way to construct linear $`\mathrm{Diff}(S^1)`$-invariant differential operators on $`๐_{\lambda ,\mu }^k(S^1)`$. Assume there are two $`\mathrm{Diff}(S^1)`$-invariant differential operators:
* a bilinear differential operator $`J:_\nu _\lambda _\mu ;`$
* a linear projection $`\pi :๐_{\lambda ,\mu }^k(S^1)_\nu .`$
We define a linear map $`J\pi :๐_{\lambda ,\mu }^k(S^1)๐_{\lambda ,\mu }^k(S^1)`$ as follows
$$\left(J\pi \right)(A)()=J(\pi (A),).$$
This map is obviously $`\mathrm{Diff}(S^1)`$-invariant.
This is the way invariant differential operators on the modules $`๐_{\lambda ,\mu }^k(S^1)`$ are related to invariant bilinear differential operators on densities. We give here the complete list of bilinear operators on densities and the complete list of linear projections. We will then specify the generators of the algebras $`_{\lambda ,\mu }^k(S^1)`$ that can be obtained by the above construction.
#### 4.5.1 Bilinear invariant differential operators on tensor densities
The classification of invariant bilinear differential operators on tensor fields is due to P. Grozman . His list is particularly interesting in the one-dimensional case (see also ).
Let us recall here the complete list.
1. Every zeroth-order operator $`_\nu _\lambda _{\nu +\lambda }`$ is of the form:
$$\varphi (x)(dx)^\lambda \psi (x)(dx)^\mu c\varphi (x)\psi (x)(dx)^{\lambda +\mu },$$
where $`c`$. From now on we omit a scalar multiple $`c`$.
2. Every first order operator $`_\nu _\lambda _{\nu +\lambda +1}`$ is as follows
$$\{\varphi (x)(dx)^\nu ,\psi (x)(dx)^\lambda \}=\left(\nu \varphi (x)\psi (x)^{}\lambda \varphi (x)^{}\psi (x)\right)(dx)^{\nu +\lambda +1},$$
(4.14)
where $`x`$ is a local coordinate on $`M`$ and we identify tensor densities with functions. The operator (4.14) is nothing but the Poisson bracket on $`T^{}S^1`$ (or $`T^{}^1`$).
For every $`(\nu ,\lambda )(0,0)`$, the operator (4.14) is the only $`\mathrm{Diff}(S^1)`$\- (or $`\mathrm{Diff}(^1)`$-) invariant operator, otherwise there are two linearly independent operators: $`\varphi d(\psi )`$ and $`d(\varphi )\psi `$, where $`d`$ is the de Rham differential.
3. There exist second order operators $`_\nu _\lambda _{\nu +\lambda +2}`$ given by the compositions:
$$\begin{array}{ccc}\hfill \varphi \psi & & \{d\varphi ,\psi \}\text{for}\nu =0,\hfill \\ \hfill \varphi \psi & & \{\varphi ,d\psi \}\text{for}\lambda =0,\hfill \\ \hfill \varphi \psi & & d\{\varphi ,\psi \}\text{for}\nu +\lambda =1.\hfill \end{array}$$
(4.15)
4. Three third-order bilinear invariant differential operators $`_\nu _\lambda _{\nu +\lambda +3}`$ are also given by compositions:
$$\begin{array}{ccc}\hfill \varphi \psi & & \{d\varphi ,d\psi \}\text{for}(\nu ,\lambda )=(0,0),\hfill \\ \hfill \varphi \psi & & d\{d\varphi ,\psi \}\text{for}(\nu ,\lambda )=(0,2),\hfill \\ \hfill \varphi \psi & & d\{\varphi ,d\psi \}\text{for}(\nu ,\lambda )=(2,0).\hfill \end{array}$$
(4.16)
5. The only differential operator of order 3 which is not a composition of the operators of lesser orders is the famous Grozman operator $`G:_{\frac{2}{3}}(S^1)_{\frac{2}{3}}(S^1)_{\frac{5}{3}}(S^1)`$ already mentioned in Introduction, see (1.1). It is given by the following expression:
$$G(\varphi (x)(dx)^{\frac{2}{3}},\psi (x)(dx)^{\frac{2}{3}})=\left(2\left|\begin{array}{cc}\varphi (x)& \psi (x)\\ \varphi ^{\prime \prime \prime }(x)& \psi ^{\prime \prime \prime }(x)\end{array}\right|+3\left|\begin{array}{cc}\varphi ^{}(x)& \psi ^{}(x)\\ \varphi ^{\prime \prime }(x)& \psi ^{\prime \prime }(x)\end{array}\right|\right)(dx)^{\frac{5}{3}}.$$
(4.17)
The $`\mathrm{Diff}(S^1)`$-invariance of this operator can be easily checked directly.
###### Remark 4.3.
The operator (4.17) remains one of the most mysterious invariant differential operators. Its geometric and algebraic meaning was discussed in .
We will use the above bilinear operators to construct the symmetries, but we do not use Grozmanโs classification result in our proof.
#### 4.5.2 Invariant projections from $`๐_{\lambda ,\mu }^k(S^1)`$ to $`_\nu `$
Let us now give the list of $`\mathrm{Diff}(S^1)`$-invariant linear maps from $`๐_{\lambda ,\mu }^k(S^1)`$ to the space $`_\nu `$.
1. The well-known projection is the principal symbol map $`\sigma :๐_{\lambda ,\mu }^k(S^1)_{\mu \lambda k}.`$ given by the expression
$$\sigma \left(\underset{i=0}{\overset{k}{}}a_i(x)\frac{d^i}{dx^i}\right)=a_k(x)(dx)^{\mu \lambda k}.$$
(4.18)
The map $`\sigma `$ is obviously $`\mathrm{Diff}(S^1)`$-invariant for all $`(\lambda ,\mu )`$.
2. For all $`(\lambda ,\mu )`$, define a linear map
$$V:๐_{\lambda ,\mu }^k(S^1)_{\mu \lambda k+1}$$
as follows:
$$V(A)=\left(\alpha a_k^{}(x)+\beta a_{k1}(x)\right)(dx)^{\mu \lambda k+1}.$$
(4.19)
where
$$\alpha =\lambda k+\frac{k(k1)}{2},\beta =\mu \lambda k$$
It is easy to check that this map is $`\mathrm{Diff}(S^1)`$-invariant. The map (4.19) is a โfirst-order analogโ of the principal symbol.
###### Remark 4.4.
If $`\lambda +\mu =1`$, then this map is proportional to the principal symbol of the $`(k1)`$-th order operator $`A(1)^kA^{}`$. In other words,
$$V=\sigma (\mathrm{Id}(1)^kC)$$
if $`\lambda +\mu =1`$.
3. In the particular case,
$$\lambda =\frac{1k}{2},\mu =\frac{1+k}{2}$$
(4.20)
The map (4.19) vanishes. In this case, there are two independent projections onto $`_1`$:
$$Aa_k^{}(x)dx,Aa_{k1}(x)dx$$
which are $`\mathrm{Diff}(S^1)`$-invariant.
4. It turns out that for some special values of the parameters $`\lambda `$ and $`\mu `$, there exist second-order analogues of the operators (4.18) and (4.19).
###### Proposition 4.5.
For every $`k3`$ and $`(\lambda ,\mu )`$ satisfying the relation
$$\left(\lambda +\frac{k2}{3}\right)\left(\mu \frac{k+1}{3}\right)+\frac{1}{36}(k+1)(k2)=\mathrm{\hspace{0.17em}0},$$
(4.21)
there exists a $`\mathrm{Diff}(S^1)`$-invariant map $`W:๐_{\lambda ,\mu }^k_{\mu \lambda k+2}`$ given by
$$W(A)=\left(\alpha _2a_k^{\prime \prime }(x)+\alpha _1a_{k1}^{}(x)+\alpha _0a_{k2}(x)\right)(dx)^{\mu \lambda k+2},$$
(4.22)
where the coefficients are defined by
$$\begin{array}{ccc}\alpha _2\hfill & =\hfill & \frac{2}{3}k(k1)(k+3\lambda 2)^2\hfill \\ \alpha _1\hfill & =\hfill & 2(k1)(k+3\lambda 2)(22\lambda k)\hfill \\ \alpha _0\hfill & =\hfill & 3k^2+12\lambda k+12\lambda ^211k24\lambda +10.\hfill \end{array}$$
(4.23)
###### Proof.
Straightforward. โ
5. There exists one more invariant differential projection $`๐_{0,1}^k(S^1)_0(S^1)`$ given by the composition
$$\pi _\delta =P_0C\delta ^1(\mathrm{Id}P_0),$$
(4.24)
where $`P_0`$ is defined by (4.9) and $`C`$ is the conjugation.
6. If there is an invariant map from $`๐_{\lambda ,\mu }^k(S^1)`$ to the space of 1-forms $`_1`$, then one can integrate the result and obtain a non-local (i.e., non-differential) invariant linear map with values in $`_0`$. For instance, for $`๐_{0,1}^k(S^1)`$, the operator $`P_0`$ defined by (4.9) satisfies the required condition. One gets a $`\mathrm{Diff}(S^1)`$-invariant map:
$$A_{S^1}a_0(x)๐x,A๐_{0,1}^k(S^1).$$
(4.25)
It was proven in that there are no other $`\mathrm{Diff}(S^1)`$-invariant projections $`๐_{\lambda ,\mu }^k(S^1)_\nu `$ than the above ones and their compositions with $`C`$ and $`d`$. We use these operators to construct the generators of the symmetry algebras (but we do not use the classification result of in the proofs of our theorems).
## 5 Computing the algebras of symmetry
We will now investigate, case by case, the non-trivial algebras $`_{\lambda ,\mu }^k(S^1)`$. We will construct the generators of these algebras and calculate the multiplication tables. We then give an explicit identification of the algebras $`_{\lambda ,\mu }^k(S^1)`$ with the matrix algebras introduced in Section 3.1. The constructions of this section prove that the algebras $`_{\lambda ,\mu }^k(S^1)`$ are at least as big as stated in Section 3.
The proof of the second part of our classification theorems, namely that there are no other symmetries than we construct and study here, will be given in Section 6.
### 5.1 The algebra $`_{0,1}^k(S^1)`$
Let us start with the most particular algebra $`_{0,1}^k(S^1)`$ for all $`k`$. One has in this case
$$_{0,1}^k(S^1)=\mathrm{Span}(\mathrm{Id},C,P_0,P_0^{},P_1,L)$$
where the generators are defined by (4.8)-(4.12).
Let us now calculate the relations between these generators.
###### Proposition 5.1.
The multiplication table for the associative algebra $`_{0,1}^k(S^1)`$ is as follows:
$$\begin{array}{ccccccc}& \mathrm{Id}& P_0& C& P_0^{}& P_1& L\\ & & & & & & \\ & & & & & & \\ \mathrm{Id}& \mathrm{Id}& P_0& C& P_0^{}& P_1& L\\ & & & & & & \\ P_0& P_0& P_0& P_0^{}& P_0^{}& 0& 0\\ & & & & & & \\ C& C& P_0& \mathrm{Id}& P_0^{}& P_0^{}P_1P_0& L\\ & & & & & & \\ P_0^{}& P_0^{}& P_0& P_0& P_0^{}& P_0^{}P_0& 0\\ & & & & & & \\ P_1& P_1& 0& P_1& 0& P_1& L\\ & & & & & & \\ L& L& L& L& L& 0& 0\end{array}$$
(5.26)
###### Proof.
First, consider the product of $`P_1`$ and $`C`$. From the definition (4.8) one obtains
$$\begin{array}{ccc}\hfill (P_1C)(A)& =& P_1\left(\underset{i=0}{\overset{k}{}}(1)^i\left(\frac{d}{dx}\right)^ia_i(x)\right)\hfill \\ & =& P_1\left(\underset{i=0}{\overset{k}{}}\underset{j=0}{\overset{i}{}}(1)^i\left(\genfrac{}{}{0pt}{}{i}{j}\right)a_i^{(ij)}(x)\frac{d^j}{dx^j}\right)\hfill \\ & =& P_1\left(\underset{j=0}{\overset{k}{}}\underset{i=j}{\overset{k}{}}(1)^i\left(\genfrac{}{}{0pt}{}{i}{j}\right)a_i^{(ij)}(x)\frac{d^j}{dx^j}\right)\hfill \end{array}$$
and then from (4.12) it follows that
$$\begin{array}{ccc}& =& \left(\underset{j=1}{\overset{k}{}}\underset{i=j}{\overset{k}{}}(1)^{i+j1}\left(\genfrac{}{}{0pt}{}{i}{j}\right)a_i^{(i1)}(x)\right)d\hfill \\ & =& \left(\underset{i=1}{\overset{k}{}}\underset{j=1}{\overset{i}{}}(1)^{i+j1}\left(\genfrac{}{}{0pt}{}{i}{j}\right)a_i^{(i1)}(x)\right)d\hfill \\ & =& \left(\underset{i=1}{\overset{k}{}}(1)^ia_i^{(i1)}(x)\right)d\hfill \\ & =& P_1(A)\hfill \end{array}$$
as presented in table (5.26).
Now, consider the product $`CP_1`$. One then has from (4.8), (4.12) and (4.10):
$$\begin{array}{ccc}\hfill (CP_1)(A)& =& \underset{i=1}{\overset{k}{}}(1)^i\left(a_i^{(i1)}(x)\frac{d}{dx}+a_i^{(i)}(x)\right)\hfill \\ & =& P_0^{}P_1P_0.\hfill \end{array}$$
Furthermore, one obtains $`P_0^{}P_1=(P_0CP_1)(A)=P_0^{}P_0`$.
For other products of the the generators the results given in the table immediately follow from the definition. โ
Let us finally give an explicit isomorphism between the algebra $`_{0,1}^k(S^1)`$ and the matrix algebra $`๐^2`$ described in Section 3.1.
One checks using the multiplication table (5.26) that the formulรฆ
$$\begin{array}{ccc}\hfill \overline{a}& =& \frac{1}{2}(2P_1+P_0P_0^{}),\hfill \\ \hfill \overline{b}& =& \frac{1}{2}(P_0+P_0^{}),\hfill \\ \hfill \overline{c}& =& \frac{1}{2}(P_0P_0^{}),\hfill \\ \hfill \overline{d}& =& L,\hfill \end{array}$$
where $`\overline{a},\overline{b},\overline{c},\overline{d}`$ are the generators of $`๐`$, see Appendix 8, define an isomorphism of the associative algebras
$$\mathrm{Span}(P_0,P_0^{},P_1,L)๐.$$
The two more generators
$$\begin{array}{ccc}\hfill z_1& =& \mathrm{Id}+CP_0P_0^{}\hfill \\ \hfill z_2& =& \mathrm{Id}CP_0+P_0^{}2P_1\hfill \end{array}$$
are in the center and span the second summand $`^2`$.
The above generators are linearly independent if $`k3`$ and span the algebra
$$_{0,1}^k(S^1)๐^2,$$
in accordance with Theorem 3.1, 3, Theorem 3.2, 4 and Theorem 3.4, 7.
If $`k=2`$, then $`z_2=0`$ so that $`_{0,1}^2(S^1)๐`$, see Theorem 3.5, 5. Finally, if $`k=1`$, then $`z_1=z_2=0`$ and one has $`_{0,1}^1(S^1)๐`$ as stated Theorem 3.6, 2.
### 5.2 Algebras of symmetry in order $`k5`$
Assume that $`k5`$. We already investigated the algebra $`_{0,1}^k(S^1)`$. There are two more non-trivial algebras in this case, namely the algebra $`_{0,0}^k(S^1)`$ and the algebra $`_{1,1}^k(S^1)`$ which is isomorphic to $`_{0,0}^k(S^1)`$ by conjugation.
The algebra of symmetry $`_{0,0}^k(S^1)`$ is as follows
$$_{0,0}^k(S^1)=\mathrm{Span}(\mathrm{Id},P_0,S),$$
where $`P_0`$ and $`S`$ are as in (4.9) and (4.13), respectively. These operators are independent for $`k4`$. We have:
$$P_0S=SP_0=P_0,P_0^2=P_0,S^2=\mathrm{Id}.$$
Indeed, the first two relations are due to the fact that the scalar term of $`S(A)`$ is equal to $`a_0(x)`$, cf. eq. (4.13) and the explicit expression for $`S`$. Put
$$1=\mathrm{Id},\overline{a}_1=P_0\text{and}\overline{a}_2=\frac{1}{\sqrt{2}}(\mathrm{Id}S).$$
One obtains the generators of the algebra $`^3`$, cf. Appendix 8. Finally, one has
$$_{0,0}^k(S^1)^3,$$
as stated in Theorem 3.1, 2.
The algebras $`_{0,\mu }^k(S^1)`$ corresponding to the generic values of $`\mu `$ have only two generators: $`\mathrm{Id}`$ and $`P_0`$. These algebras are obviously isomorphic to $`^2`$.
### 5.3 Algebras of symmetry in order 4
Consider the modules of differential operators of order $`k=4`$:
$$A=a_4(x)\frac{d^4}{dx^4}+a_3(x)\frac{d^3}{dx^3}+a_2(x)\frac{d^2}{dx^2}+a_1(x)\frac{d}{dx}+a_0(x).$$
We will study all the exceptional modules systematically and investigate every non-trivial algebra of symmetry.
The generators of the algebras $`_{0,1}^4(S^1)`$ and $`_{0,0}^4(S^1)_{1,1}^4(S^1)`$ are the same as for $`k=5`$. Let us consider other interesting cases.
#### 5.3.1 The algebras $`_{0,\frac{5}{4}}^4(S^1)`$ and $`_{\frac{1}{4},1}^4(S^1)`$
We already constructed two generators of the algebra $`_{0,\frac{5}{4}}^4(S^1)`$, namely $`\mathrm{Id}`$ and $`P_0`$. One extra generator is obtained by the following procedure.
The values $`(\lambda ,\mu )=(0,\frac{5}{4})`$ satisfy the relation (4.21), so that the map
$$W:๐_{0,\frac{5}{4}}^4(S^1)_{\frac{3}{4}}$$
defined by (4.22) is $`\mathrm{Diff}(S^1)`$-invariant. There is a second-order bilinear differential operator
$$J:_{\frac{3}{4}}_0_{\frac{5}{4}}$$
defined by the second formula in (4.15). Applying the construction of Section 4.5 we consider the composition $`JW`$ defined as in Section 4.5 to obtain an element of algebra $`_{0,\frac{5}{4}}^4(S^1)`$.
Let us now compute the relations between the generators. The constructed map is given by the following explicit formula
$$\begin{array}{ccc}\hfill \left(JW\right)(A)& =& \left(\frac{16}{7}a_4^{\prime \prime }(x)\frac{12}{7}a_3^{}(x)+a_2(x)\right)\frac{d^2}{dx^2}\hfill \\ & & +\frac{4}{3}\left(\frac{12}{7}a_4^{\prime \prime \prime }(x)\frac{9}{7}a_3^{\prime \prime }(x)+\frac{3}{4}a_2^{}(x)\right)\frac{d}{dx}.\hfill \end{array}$$
(5.27)
Note that the right hand side is understood as an element of $`๐_{0,\frac{5}{4}}^4(S^1)`$. Since $`P_0(A)=a_0(x)`$, the product of $`JW`$ and $`P_0`$ vanishes:
$$\left(JW\right)P_0=P_0\left(JW\right)=0.$$
One also has the relations
$$P_0^2=P_0\text{and}\left(JW\right)^2=JW.$$
Finally, one gets the following answer:
$$_{0,\frac{5}{4}}^4(S^1)=\mathrm{Span}(\mathrm{Id},P_0,JW)^3,$$
as stated by Theorem 3.2, 2.
The conjugation establishes an isomorphism between the algebras $`_{\frac{1}{4},1}^4(S^1)`$ and $`_{0,\frac{5}{4}}^4(S^1)`$.
#### 5.3.2 The algebras $`_{0,3}^4(S^1)`$ and $`_{2,1}^4(S^1)`$
The algebra $`_{0,3}^4(S^1)`$ has the generators $`\mathrm{Id},P_0`$ and the following one constructed in Section 4.5.
Consider the projection $`V:๐_{0,3}^4(S^1)_0`$ defined by formula (4.19) and the third-order bilinear map $`J:_0_0_3`$, namely the first of the three operators (4.16). Their composition $`JV`$ is an element of $`_{0,3}^4(S^1)`$:
$$\left(JV\right)(A)=\left(6a_4^{\prime \prime }(x)a_3^{}(x)\right)\frac{d^2}{dx^2}\left(6a_4^{\prime \prime \prime }(x)a_3^{\prime \prime }(x)\right)\frac{d}{dx},$$
(5.28)
where the right hand side is understood as an element of $`๐_{0,3}^4(S^1)`$.
Finally, the algebra $`_{0,3}^4(S^1)`$ is of the form
$$_{0,3}^4(S^1)=\mathrm{Span}(\mathrm{Id},P_0,JV).$$
To obtain the isomorphism $`_{0,3}^4(S^1)๐`$ (see Theorem 3.2 part 3), one checks the following relations
$$P_0\left(JV\right)=\left(JV\right)P_0=(JV)^2=0.$$
Then the standard generators of $`๐`$ correspond to $`\{\mathrm{Id}P_0,JV,P_0\}`$.
The algebra $`_{2,1}^4(S^1)`$ is isomorphic to $`_{0,3}^4(S^1)`$ by conjugation.
#### 5.3.3 The algebra $`_{\frac{2}{3},\frac{5}{3}}^4(S^1)`$ and the Grozman operator
The conjugation map $`C`$ and $`\mathrm{Id}`$ are, of course, generators of symmetry of the module $`๐_{\frac{2}{3},\frac{5}{3}}^4(S^1)`$. One extra generator can be obtained as follows.
Consider the operator (4.19)
$$V:๐_{\frac{2}{3},\frac{5}{3}}^4(S^1)_{\frac{2}{3}}$$
and compose it with the Grozman operator $`G`$ given by (4.17); we obtain (up to a constant) the following operator:
$$\begin{array}{ccc}\hfill \left(GV\right)(A)& =& \left(a_3(x)2a_4^{}(x)\right)\frac{d^3}{dx^3}+\left(\frac{3}{2}a_3^{}(x)3a_4^{\prime \prime }(x)\right)\frac{d^2}{dx^2}\hfill \\ & & \left(\frac{3}{2}a_3^{\prime \prime }(x)3a_4^{\prime \prime \prime }(x)\right)\frac{d}{dx}\left(a_3^{\prime \prime \prime }(x)2a_4^{(IV)}(x)\right)\hfill \end{array}$$
(5.29)
which is a generator of $`_{\frac{2}{3},\frac{5}{3}}^4(S^1)`$.
The relations between the conjugation map and the above operator are:
$$\left(GV\right)C=C\left(GV\right)=GV.$$
One also has:
$$\left(GV\right)^2=GV.$$
One easily deduces from the above relations that the algebra
$$_{\frac{2}{3},\frac{5}{3}}^4(S^1)=\mathrm{Span}(\mathrm{Id},C,GV)$$
is, indeed, isomorphic to $`^3`$, cf. Theorem 3.2, 2.
### 5.4 Algebras of symmetry in order 3
Consider the differential operators of order $`k=3`$:
$$A=a_3(x)\frac{d^3}{dx^3}+a_2(x)\frac{d^2}{dx^2}+a_1(x)\frac{d}{dx}+a_0(x).$$
We will describe all the non-trivial algebras of symmetry.
#### 5.4.1 The algebra $`_{\frac{2}{3},\frac{5}{3}}^3(S^1)`$
The conjugation map $`C`$, as well as the identity $`\mathrm{Id}`$, are, of course, generators of the symmetry algebra $`_{\frac{2}{3},\frac{5}{3}}^3(S^1)`$. Let us construct one more generator. The principal symbol map is of the form:
$$\sigma :๐_{\frac{2}{3},\frac{5}{3}}^3(S^1)_{\frac{2}{3}}.$$
We compose it with the Grozman operator to obtain a new generator $`G\sigma `$. This operator is given by the same formula (5.29) as above, but with $`a_4(x)0`$. The relations between the generators are also the same as above, so that the symmetry algebra is $`_{\frac{2}{3},\frac{5}{3}}^3(S^1)=^3.`$
#### 5.4.2 The hyperbola $`(3\lambda +1)(3\mu 4)=1`$
Consider the class of modules $`๐_{\lambda ,\mu }^3(S^1)`$ with $`(\lambda ,\mu )`$ satisfying the quadratic relation
$$(3\lambda +1)(3\mu 4)=1,$$
(5.30)
see Theorem 3.4, 1.
First of all, we observe that this relation is precisely the relation (4.21) specified for $`k=3`$. The operator $`W:๐_{\lambda ,\mu }^3(S^1)_{\mu \lambda 1}`$ is then well-defined. Composing this operator with the Poisson bracket (4.14)
$$\{,\}:_{\mu \lambda 1}_\lambda _\mu ,$$
one obtains a generator of the algebra $`_{\lambda ,\mu }^3(S^1)`$. Let us denote this generator by $`๐ฒ`$:
$$\begin{array}{ccc}\hfill ๐ฒ(A)& =& (\mu \lambda 1)\left(\alpha _2a_3^{\prime \prime }(x)+\alpha _1a_2^{}(x)+\alpha _0a_1(x)\right)\frac{d}{dx}\hfill \\ & & \lambda \left(\alpha _2a_3^{\prime \prime \prime }(x)+\alpha _1a_2^{\prime \prime }(x)+\alpha _0a_1^{}(x)\right),\hfill \end{array}$$
(5.31)
where according to (4.23)
$$\begin{array}{ccc}\alpha _2\hfill & =\hfill & (3\lambda +1)^2,\hfill \\ \alpha _1\hfill & =\hfill & (3\lambda +1)(12\lambda ),\hfill \\ \alpha _0\hfill & =\hfill & 3\lambda ^2+3\lambda +1.\hfill \end{array}$$
In the generic case, $`(\lambda ,\mu )(0,1)`$ or $`(\frac{2}{3},\frac{5}{3})`$ the symmetry algebra has two generators: $`_{\lambda ,\mu }^3(S^1)=\mathrm{Span}(\mathrm{Id},S).`$ The generator $`๐ฒ`$ satisfies the relation
$$๐ฒ^2=\alpha _0(\mu \lambda 1)๐ฒ.$$
But $`\alpha _00`$ and if $`\mu \lambda 1=0`$, then (5.30) implies $`(\lambda ,\mu )=(0,1)`$. Finally, in the generic case, one obtains:
$$_{\lambda ,\mu }^3(S^1)^2.$$
###### Remark 5.2.
The module $`๐_{\frac{2}{3},\frac{5}{3}}^3(S^1)`$ belongs to the family (5.30). We have already considered this module separately, see Section 5.4.1. In this case, we have three generators: $`_{\frac{2}{3},\frac{5}{3}}^3(S^1)=\mathrm{Span}(\mathrm{Id},C,๐ฒ)`$ which are different from $`G\sigma `$. One checks, however, that the generator $`G\sigma `$ can be expressed in terms of the above ones:
$$G\sigma =\frac{1}{2}(\mathrm{Id}C)\frac{9}{4}๐ฒ.$$
#### 5.4.3 The line $`\mu \lambda =2`$
Consider the family of modules $`๐_{\lambda ,\mu }^3(S^1)`$ satisfying the property $`\mu \lambda =2`$ as in Theorem 3.4, 2.
The operator (4.19) is, in this case, $`V:๐_{\lambda ,\mu }^3(S^1)_0`$. Consider its composition with the second-order bilinear operator $`J:_0_\lambda _\mu `$ given by the first formula in (4.15). In the generic case, that is, where
$$(\lambda ,\mu )(0,2),(\frac{1}{2},\frac{3}{2}),(1,1),$$
one has $`_{\lambda ,\mu }^3(S^1)=\mathrm{Span}(\mathrm{Id},JV).`$
One obtains the following explicit formula for the constructed generator:
$$\left(JV\right)(A)=\left(3(\lambda +1)a_3^{\prime \prime }(x)a_2^{}(x)\right)\frac{d}{dx}\lambda \left(3(\lambda +1)a_3^{\prime \prime \prime }(x)a_2^{\prime \prime }(x)\right)$$
(5.32)
and immediately gets the following relation:
$$\left(JV\right)^2=0.$$
This implies $`_{\lambda ,\mu }^3(S^1)๐`$ in the generic case.
#### 5.4.4 The modules $`๐_{0,2}^3(S^1),๐_{1,1}^3(S^1)`$ and $`๐_{\frac{1}{2},\frac{3}{2}}^3(S^1)`$
Let us consider some exceptional modules still satisfying $`\mu \lambda =2`$.
In the case, $`(\lambda ,\mu )=(0,2)`$, one has an extra generator of symmetry, as compared with the preceding section. It is given by the operator $`P_0`$ as in (4.9). One then has from (5.32)
$$\left(JV\right)P_0=0,P_0\left(JV\right)=0.$$
The isomorphism
$$_{0,2}^3(S^1)=\mathrm{Span}(\mathrm{Id},JV,P_0)๐$$
is then obvious in accordance with Theorem 3.4, 4.
The conjugation map $`C`$ establishes an isomorphism $`_{1,1}^3(S^1)_{0,2}^3(S^1)`$, so that tha algebra $`_{1,1}^3(S^1)`$ is also isomorphic to $`๐`$.
In the interesting case $`๐_{\frac{1}{2},\frac{3}{2}}^3(S^1)`$, the extra generator is given by the conjugation map $`C`$, so that $`_{\frac{1}{2},\frac{3}{2}}^3(S^1)=\mathrm{Span}(\mathrm{Id},JV,C).`$ The relations between the generators are
$$\left(JV\right)C=JV,C\left(JV\right)=JV$$
as follows from (5.32) and (4.8). One obtains
$$_{\frac{1}{2},\frac{3}{2}}^3(S^1)๐ฑ_2,$$
as stated by Theorem 3.4, 5.
#### 5.4.5 The modules $`๐_{0,3}^3(S^1)`$ and $`๐_{2,1}^3(S^1)`$
The principal symbol map (4.18) is as follows $`\sigma :๐_{0,3}^3(S^1)_0`$. Compose this map with the third-order bilinear operator $`J:_0_0_3`$ defined by the first equation in (4.16). The explicit expression of the constructed generator is
$$\left(J\sigma \right)(A)=a_3^{}(x)\frac{d^2}{dx^2}a_3^{\prime \prime }(x)\frac{d}{dx}.$$
(5.33)
The symmetry algebra is then $`_{0,3}^3(S^1)=\mathrm{Span}(\mathrm{Id},J\sigma ,P_0).`$ One easily gets the relations
$$\left(J\sigma \right)P_0=P_0\left(J\sigma \right)=\left(J\sigma \right)^2=0$$
and, finally, $`_{0,3}^3(S^1)๐`$, see Theorem 3.4, 4. The algebra $`_{2,1}^3(S^1)`$ is isomorphic to the above one by conjugation.
### 5.5 Algebras of symmetry in orders 2
Consider now the modules of differential operators of order 2.
For all $`(\lambda ,\mu )`$, there is a generator of the algebra $`_{\lambda ,\mu }^2(S^1)`$ given by the composition of the projection (4.19)
$$V:๐_{\lambda ,\mu }^2(S^1)_{\mu \lambda 1}$$
and the Poisson bracket
$$\{,\}:_{\mu \lambda 1}_\lambda _\mu .$$
The explicit formula for this generator is as follows:
$$\begin{array}{ccc}\hfill ๐ฑ(A)& =& (\mu \lambda 1)\left((2\lambda +1)a_2^{}(x)+(\mu \lambda 2)a_1(x)\right)\frac{d}{dx}\hfill \\ & & \lambda \left((2\lambda +1)a_2^{\prime \prime }(x)+(\mu \lambda 2)a_1^{}(x)\right).\hfill \end{array}$$
(5.34)
This generator satisfies the following relation:
$$๐ฑ^2=(\mu \lambda 1)(\mu \lambda 2)๐ฑ.$$
In the generic case, the algebra of symmetry is $`_{\lambda ,\mu }^2(S^1)=\mathrm{Span}(\mathrm{Id},๐ฑ)`$. It is obviously isomorphic to $`^2`$.
If either $`\mu \lambda =1`$ or $`\mu \lambda =2`$, then $`๐ฑ^2=0`$. In these cases, $`_{\lambda ,\mu }^2(S^1)=\mathrm{Span}(\mathrm{Id},๐ฑ)๐`$, see Theorem 3.5, 1.
Of course, if $`\lambda =0`$, then there is one more generator, namely $`P_0`$. Then one has
$$P_0๐ฑ=๐ฑP_0=0$$
and so $`_{0,\mu }^2(S^1)=\mathrm{Span}(\mathrm{Id},๐ฑ,P_0)`$ which is isomorphic to $`^3`$ for generic $`\mu `$ while
$$_{0,2}^2(S^1)_{1,1}^2(S^1)๐$$
see Theorem 3.5, 2 and 3.5, 4.
In the case where $`\lambda +\mu =1`$, the conjugation map is well-defined. One checks that in this case the generator (5.34) is a linear combination of $`\mathrm{Id}`$ and $`C`$:
$$๐ฑ=\lambda (2\lambda +1)\left(C\mathrm{Id}\right)$$
so that $`_{\lambda ,\mu }^2(S^1)^2`$.
The exceptional case $`(\lambda ,\mu )=(\frac{1}{2},\frac{3}{2})`$ corresponds to (4.20). There are two invariant projections from $`๐_{\frac{1}{2},\frac{3}{2}}^2(S^1)`$ to $`_1`$ in this case. Composing one of them with the Poisson bracket
$$\{,\}:_1_{\frac{1}{2}}_{\frac{3}{2}}$$
one obtains a generator
$$a_2(x)\frac{d^2}{dx^2}+a_1(x)\frac{d}{dx}+a_0(x)a_2^{}(x)\frac{d}{dx}+\frac{1}{2}a_2^{\prime \prime }(x)$$
independent of $`\mathrm{Id}`$ and $`C`$. One easily gets an isomorphism $`_{\frac{1}{2},\frac{3}{2}}^2(S^1)๐ฑ_2`$.
## 6 Proof of the main theorems
In this section we prove that there are no other symmetries of the modules $`๐_{\lambda ,\mu }^k(S^1)`$ than those constructed above. In other words, we give here a complete classification of symmetries.
Let $`T`$ be a linear map (1.2) commuting with the $`\mathrm{Diff}(S^1)`$-action. There are two cases:
1. The map $`T`$ is local, that is, one has $`\mathrm{Supp}(T(A))\mathrm{Supp}(A)`$ for all $`A๐_{\lambda ,\mu }^k(S^1)`$. In this case, the famous Peetre theorem (see ) guarantees that $`T`$ is a differential operator in coefficients of $`A`$.
2. The map $`T`$ is non-local, that is, for some $`A๐_{\lambda ,\mu }^k(S^1)`$ vanishing in an open subset $`US^1`$, the operator $`T(A)`$ does not vanish on $`U`$.
These two cases are completely different and should be treated separately.
### 6.1 The identification
Let us fix a parameter $`x`$ on $`S^1`$ and the corresponding coordinate $`\xi `$ on the fibers of $`T^{}S^1`$.
For our computations, we will need to identify the spaces $`๐_{\lambda ,\mu }(S^1)`$ and $`๐ฎ_\delta (S^1)`$ using the map
$$\sigma _{\mathrm{tot}}:๐_{\lambda ,\mu }(S^1)๐ฎ_\delta (S^1)$$
(6.35)
assigns to an operator (2.4) the polynomial on $`T^{}S^1`$ given by (2.7).
The map $`\sigma _{\mathrm{tot}}`$ is an isomorphism of vector spaces but not an isomorphism of $`\mathrm{Diff}(S^1)`$-modules. It will, nevertheless, allow us to compare the $`\mathrm{Diff}(S^1)`$-action on both spaces.
### 6.2 The affine Lie algebra
We introduce our main tool that will allow un to use the results of the classic invariant theory.
Let $`x`$ be an affine parameter on $`S^1`$, the Lie algebra $`\mathrm{aff}`$ of affine transformations is the two-dimensional Lie algebra generated by the translations and linear vector fields:
$$\mathrm{aff}=\mathrm{Span}(\frac{d}{dx},x\frac{d}{dx}).$$
(6.36)
Since $`\mathrm{aff}`$ is a subalgebra of $`\mathrm{Vect}(S^1)`$, every $`\mathrm{Diff}(S^1)`$-invariant map has to commute with the $`\mathrm{aff}`$-action.
###### Proposition 6.1.
The $`\mathrm{aff}`$-action on $`๐_{\lambda ,\mu }^k(S^1)`$ depends only on the difference $`\mu \lambda `$ and coincides with the action on $`๐ฎ_\delta (S^1)`$ after identification $`(\text{6.35})`$.
###### Proof.
Straightforward. โ
The well-known result of invariant theory states that the associative algebra of differential operators on $`T^{}S^1`$ commuting with the $`\mathrm{aff}`$-action is generated by
$$E=\xi \frac{}{\xi }\text{and}D=\frac{}{x}\frac{}{\xi },$$
(6.37)
see . The operator $`E`$ is called the Euler field and the operator $`D`$ the divergence.
Every differential $`\mathrm{Diff}(M)`$-invariant operator (1.2) can therefore be expressed in terms of these operators. In local coordinates, any $`\mathrm{Diff}(M)`$-invariant map (1.2) is therefore of the form
$$T=T(E,D).$$
###### Example 6.2.
The expression
$$C=\mathrm{exp}(D)\mathrm{exp}(i\pi E)$$
for the conjugation map (4.8) is worth mentioning for the aesthetic reasons.
### 6.3 The local case: invariant differential operators
We consider the algebra $`_{\lambda ,\mu }^{k}{}_{}{}^{\mathrm{loc}}(S^1)`$ of local (and thus differential) $`\mathrm{Diff}(S^1)`$-invariant linear maps (1.2).
Let us restrict the map $`T`$ to the homogeneous component $`_{k\delta }`$ in (2.7). Since the Euler operator $`E`$ reduces to a constant, one has
$$T|_{_{\delta k}}=\underset{\mathrm{}=0}{\overset{k}{}}T_{k,\mathrm{}}D^{\mathrm{}}$$
(6.38)
where $`T_{k,\mathrm{}}`$ are some constants.
The operator $`T`$ has to commute with the $`\mathrm{Vect}(S^1)`$-action on $`๐_{\lambda ,\mu }^k(S^1)`$. Consider the Lie subalgebra $`\mathrm{sl}(2)\mathrm{Vect}(S^1)`$ generated by three vector fields:
$$\mathrm{sl}(2)=\mathrm{Span}(\frac{d}{dx},x\frac{d}{dx},x^2\frac{d}{dx},).$$
(6.39)
###### Remark 6.3.
Assume that $`x`$ is an affine parameter on $`S^1`$, that is, we identify $`S^1`$ with $`^1`$ with homogeneous coordinates $`(x_1:x_2)`$ and choose $`x=x_1/x_2`$. The vector fields (6.39) are then globally defined and correspond to the standard projective structure on $`^1`$, see, e.g., .
###### Proposition 6.4.
A linear map (1.2) written in the form (6.38) is $`\mathrm{sl}(2,)`$-invariant if and only if it satisfies the recurrence relation
$$\left(k+2\lambda 1\right)T_{k1,\mathrm{}1}\left(k+2\lambda \mathrm{}\right)T_{k,\mathrm{}1}\mathrm{}\left(2(\mu \lambda )2k+\mathrm{}1\right)T_{k,\mathrm{}}=0.$$
(6.40)
###### Proof.
The form (6.38) is already invariant with respect to the affine subalgebra (6.36) of $`\mathrm{sl}(2,)`$. It remains to impose the equivariance condition with respect to the vector field $`X=x^2\frac{d}{dx}`$. Let us compute the Lie derivative (2.5) along this vector field. Again, we use the identification (6.35) and express it in terms of symbols. One has
$$_X^{\lambda ,\mu }=L_X^{\mu \lambda }(2\lambda +E)\frac{}{\xi },$$
where
$$L_X^{\mu \lambda }=x^2\frac{}{x}2x\xi \frac{}{\xi }+2(\mu \lambda )x.$$
The equivariance condition $`[T,_X^{\lambda ,\mu }]=0`$ readily leads to the relation (6.40). โ
Let us now impose the equivariance condition with respect to the vector field
$$X=x^3\frac{d}{dx}.$$
This vector field is not globally defined on $`S^1`$ and has a singularity at $`x=\mathrm{}`$. Indeed, choose the coordinate $`z=\frac{1}{x}`$ in a vicinity of the point $`z=0`$, the vector field $`X`$ is written $`X=\frac{1}{z}\frac{d}{dz}`$, cf. Remark 6.3. However, every $`\mathrm{Diff}(S^1)`$-invariant differential operator $`T`$ has to commute with this vector field everywhere for $`x\mathrm{}`$.
###### Proposition 6.5.
A differential operator $`T`$ commutes with the action of $`X=x^3\frac{d}{dx}`$ if and only if $`T`$ satisfies the relation (6.40) together with the relation
$$\left(6\lambda +3k3\right)T_{k1,\mathrm{}1}+\mathrm{}\left(3(\mu \lambda )3k+\mathrm{}2\right)T_{k,\mathrm{}}=0.$$
(6.41)
Proof is similar to that of Proposition 6.4.
Let us now calculate the dimensions of the algebras $`_{\lambda ,\mu }^{k}{}_{}{}^{\mathrm{loc}}(S^1)`$.
###### Theorem 6.6.
The dimensions of the algebras $`_{\lambda ,\mu }^{k}{}_{}{}^{\mathrm{loc}}(S^1)`$ are given in the following table
$$\begin{array}{ccccccc}& & & & & & \\ k\hfill & 0& 1& 2& 3& 4& 5\\ & & & & & & \\ (\lambda ,\mu )\text{generic}\hfill & 1& 2& 2& 1& 1& 1\\ & & & & & & \\ \lambda =0,\text{or}\mu =1,\text{generic}\hfill & 1& 2& 3& 3& 2& 2\\ & & & & & & \\ \lambda +\mu =1,\text{generic}\hfill & 1& 2& 2& 2& 2& 2\\ & & & & & & \\ (3\lambda +1)(3\mu 4)=1,\text{or}\mu \lambda =2,\text{generic}\hfill & 1& 2& 2& 2& 1& 1\\ & & & & & & \\ (\lambda ,\mu )=(\frac{1}{4},1),(2,1),(0,\frac{5}{4}),(0,3)\hfill & 1& 2& 3& 3& 3& 2\\ & & & & & & \\ (\lambda ,\mu )=(0,0),(1,1)\hfill & 1& 2& 3& 3& 3& 3\\ & & & & & & \\ (\lambda ,\mu )=(\frac{2}{3},\frac{5}{3})\hfill & 1& 2& 2& 3& 3& 2\\ & & & & & & \\ (\lambda ,\mu )=(\frac{1}{2},\frac{3}{2})\hfill & 1& 2& 3& 3& 2& 2\\ & & & & & & \\ (\lambda ,\mu )=(0,1)\hfill & 1& 3& 4& 5& 5& 5\end{array}$$
###### Proof.
We will use the following result (which is similar to that of ).
###### Proposition 6.7.
Every $`\mathrm{Diff}(S^1)`$-invariant differential operator $`T`$ on $`๐_{\lambda ,\mu }^k(S^1)`$ with $`k3`$ is completely determined by its restriction to the subspace of third-order operators $`T|_{๐_{\lambda ,\mu }^3(S^1)}`$.
###### Proof.
Assume that $`k4`$ and $`T|_{๐_{\lambda ,\mu }^3(S^1)}=0`$ that is, $`T_{r,\mathrm{}}=0`$ for all $`\mathrm{}`$ and all $`r3`$. Then the system of equations (6.40,6.41) readily leads to $`T_{r,\mathrm{}}=0`$ for all $`(r,\mathrm{})`$. โ
The end of the proof of Theorem 6.6 is as follows.
We solve the system (6.40), (6.41) explicitly for $`k5`$ (we omit here the tedious computations) and obtain the result in this case.
It follows from Proposition 6.7 that the dimension of the algebras $`_{\lambda ,\mu }^{k}{}_{}{}^{\mathrm{loc}}(S^1)`$ of local symmetries with $`k3`$ can only decrease as $`k`$ becomes $`k+1`$. On the other hand, we have already constructed a set of generators of the algebra $`_{\lambda ,\mu }^k(S^1)`$ that gives a lower bound for the dimension. We thus conclude, by Proposition 6.7, that the constructed generators span the algebras $`_{\lambda ,\mu }^k(S^1)`$ for all $`k`$.
Theorem 6.6 is proved. โ
### 6.4 Non-local operators
Consider now the non-local case. We already constructed an example of a non-local linear map (1.2) commuting with the $`\mathrm{Vect}(S^1)`$-action, see formula (4.11). Let us show that there are no other such maps.
Let $`T`$ be a non-local linear map (1.2) commuting with the $`\mathrm{Vect}(S^1)`$-action. Assume that $`A๐_{\lambda ,\mu }^k(S^1)`$ vanishes on an open subset $`US^1`$, but $`T(A)๐_{\lambda ,\mu }^k(S^1)`$ does not vanish on $`U`$. Let $`X`$ be a vector field with support in $`U`$, then $`_X^{\lambda ,\mu }(A)=0`$. Since $`T`$ satisfies the relation of equivariance
$$_X^{\lambda ,\mu }(T(A))=T(_X^{\lambda ,\mu }(A)),$$
one gets
$$_X^{\lambda ,\mu }\left(T(A)\right)=0.$$
Consider the restriction $`T(A)|_U`$, this is an element of $`๐_{\lambda ,\mu }^k(U)`$. We just proved that $`T(A)|_U`$ is a $`\mathrm{Vect}(S^1)`$\- and thus a $`\mathrm{Diff}(S^1)`$-invariant differential operator $`_\lambda (U)_\mu (U)`$.
The classification of such invariant differential operators (on any manifold) is well known (see, e.g., and for proofs); the answer is as follows. There exists a unique non-trivial invariant differential operator, namely the de Rham differential. Therefore $`T(A)|_U`$ is proportional to one of the operators:
* $`\mathrm{Id}๐_{\lambda ,\lambda }^k(U)`$, so that $`\mu =\lambda `$, or
* $`d๐_{0,1}^k(U)`$, so that $`(\lambda ,\mu )=(0,1)`$.
In each case, one gets an invariant linear functional, namely, the operator $`T`$ is of the form
* $`T=t\mathrm{Id}`$, where $`t:๐_{\lambda ,\lambda }^k(U),`$ or
* $`T=td`$, where $`t:๐_{0,1}^k(U)`$,
respectively. The linear functional $`t`$ has to be $`\mathrm{Diff}(S^1)`$-invariant.
The space of symbols corresponding to the above modules are
$$_0\mathrm{}_k\text{and}_1\mathrm{}_{1k},$$
respectively. Projecting the functional $`t`$ to these modules of symbols one obtains a linear functional which is again $`\mathrm{Diff}(S^1)`$-invariant. Moreover, this functional is non-zero if and only if $`t`$ is non-zero.
###### Lemma 6.8.
There exists a unique (up to a multiplicative constant) non-trivial invariant functional $`_\lambda (S^1)`$ if and only if $`\lambda =1`$. This functional is the integral of 1-forms
$$:_1(S^1)$$
###### Proof.
The statement follows from the fact that any module of tensor fields on a manifold $`M`$ is irreducuble except the modules of differential forms $`\mathrm{\Omega }^k(M)`$, see . However, in the one-dimensional case, one can give simple direct arguments.
Let $`\tau :_\lambda (S^1)`$ be a $`\mathrm{Diff}(S^1)`$-invariant linear functional. Then $`\mathrm{ker}\tau _\lambda (S^1)`$ is a submodule. For every vector field $`X=X(x)\frac{d}{dx}`$ and every $`\lambda `$-density $`\phi =\varphi (x)(dx)^\lambda `$, one has
$$\tau (L_X\phi )=0.$$
Consider the function $`\psi (x)`$ defined by the expression
$$\varphi (x)=X(x)\varphi ^{}(x)+\lambda X^{}(x)\varphi (x).$$
It is easy to check that, in the case $`\lambda 1`$, for every $`\psi (x)`$ there exist $`X(x)`$ and $`\varphi (x)`$ such that the above equation is satisfied. If $`\lambda =1`$ then $`\psi (x)`$ has to satisfy the property $`\psi (x)๐x=0`$. โ
The lemma implies that the functional $`t`$ can be non-zero only in the second case and has to be proportional to (4.25). We conclude that every non-local $`\mathrm{Diff}(S^1)`$-invariant operator $`T`$ has to be proportional to the operator $`L`$ given by (4.11).
Theorems 3.1-3.6 are proved.
## 7 Conclusion and outlook
The classification theorems of Section 3 provide a number of particular examples of modules of differential operators. Some of these modules are known (however, the precise values of parameters $`(\lambda ,\mu )`$ are often implicit), other ones are new.
### 7.1 Known modules of differential operators
There are particular examples of modules of differential operators on $`S^1`$ that have been known for a long time. Some of these modules appear in our classification. Let us briefly mention here some interesting cases.
The family of modules $`๐_{\lambda ,\mu }^k(S^1)`$ with $`(\lambda ,\mu )`$ satisfying the condition $`\lambda +\mu =1`$ is, of course, the best known class. This is the only case when one can speak of conjugation and split a given differential operator into the sum of a symmetric and a skew-symmetric operators.
The module $`๐_{\frac{1}{2},\frac{3}{2}}^2(S^1)`$ is a well-known module of second-order differential operators. In contains a submodule of Sturm-Liouville operators
$$A=\frac{d^2}{dx^2}+a(x)$$
which is related to the Virasoro algebra (see ). This module also has a very interesting geometric meaning: it is isomorphic to the space of projective structures on $`S^1`$ (see and also ).
The value $`\lambda =\frac{1}{2}`$ is also related to the Lie superalgebras of Neveu-Schwarz and Ramond (see ). More precisely, the odd parts of these superalgebras consist of $`\frac{1}{2}`$-densities.
The above module is a particular case of the following series of modules.
The modules $`๐_{\frac{1k}{2},\frac{1+k}{2}}^k(S^1)`$, see formula (4.20), also have interesting geometric and algebraic meaning. These modules have already been studied in , it turns out that they are related to the space of curves in the projective space $`^{k1}`$ (see, e.g., ). These modules are also related to so-called Adler-Gelfand-Dickey bracket (see, e.g., and references therein).
### 7.2 New examples of modules of differential operators
Some of the particular modules $`๐_{\lambda ,\mu }^k(S^1)`$ provided by our classification theorems seem to be new.
The modules (4.21) can be characterized as the modules that are โabnormally closeโ to the corresponding $`\mathrm{Diff}(S^1)`$-modules of symbols. More precisely, in this case, the quotient-module
$$๐_{\lambda ,\mu }^k(S^1)/๐_{\lambda ,\mu }^{k3}(S^1)_{\delta k}_{\delta k+1}_{\delta k+2}$$
where $`\delta =\mu \lambda `$. The operator (4.22) has been known to the classics in some particular cases (cf. ).
An interesting particular case that belongs to the above family is the module $`๐_{\frac{2}{3},\frac{5}{3}}^3(S^1)`$. It can be characterized by three conditions: $`k=3`$ together with $`\lambda +\mu =1`$ and (4.21). This module is related to the Grozman operator (4.17).
The module $`๐_{\frac{2}{3},\frac{5}{3}}^4(S^1)`$ is also related to the Grozman operator. We strongly believe that this exceptional module has an interesting geometric and algebraic meaning.
Other examples of exceptional modules of Theorems 3.2 and 3.4, such as $`๐_{\frac{1}{4},1}^4(S^1)`$, $`๐_{2,1}^4(S^1)`$ or $`๐_{1,1}^3(S^1)`$, etc. are even more mysterious.
### 7.3 Results in the multi-dimensional case
For the sake of completeness, let us mention the results in the multi-dimensional case. Let $`M`$ be a smooth manifold, $`dimM2`$; the classification of invariant operators on $`๐_{\lambda ,\mu }^k(M)`$ was obtained in . The algebra of symmetries $`_{\lambda ,\mu }^k(M)`$ does not depend on the topology of $`M`$. This algebra is trivial for all $`(\lambda ,\mu )`$, except for the cases $`\lambda +\mu =1`$ and $`\lambda =0`$ or $`\mu =1`$. The generators are $`C,P_0`$ and $`\mathrm{Id}`$.
## 8 Appendix: generators of the algebras $`๐`$, $`๐`$ and $`^n`$
Let us introduce the generators of the matrix algebras $`๐`$, $`๐`$ and $`^n`$ that are useful in order to establish the isomorphisms with the algebras of symmetry $`_{\lambda ,\mu }^k(S^1)`$.
In the case of the algebra $`๐`$, we take the matrices
$$\stackrel{~}{a}=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right),\stackrel{~}{b}=\left(\begin{array}{cc}0& 0\\ 1& 0\end{array}\right).$$
The relations between the generators are:
$$\stackrel{~}{a}^2=\stackrel{~}{a},\stackrel{~}{a}\stackrel{~}{b}=\stackrel{~}{b}\stackrel{~}{a}=\stackrel{~}{b},\stackrel{~}{b}^2=0.$$
For the algebra $`๐`$, we consider
$$\overline{a}=\left(\begin{array}{cccc}1& 0& 0& 0\\ 0& 1& 0& 0\\ 0& 0& 0& 0\\ 0& 0& 0& 0\end{array}\right),\overline{b}=\left(\begin{array}{cccc}0& 0& 0& 0\\ 0& 0& 0& 0\\ 0& 0& 1& 0\\ 0& 0& 0& 1\end{array}\right),\overline{c}=\left(\begin{array}{cccc}0& 0& 0& 0\\ 0& 0& 0& 0\\ 0& 1& 0& 0\\ 0& 0& 0& 0\end{array}\right),\overline{d}=\left(\begin{array}{cccc}0& 0& 0& 1\\ 0& 0& 0& 0\\ 0& 0& 0& 0\\ 0& 0& 0& 0\end{array}\right).$$
The multiplication table for this algebra is
$$\begin{array}{ccccc}& \overline{a}& \overline{b}& \overline{c}& \overline{d}\\ & & & & \\ & & & & \\ \overline{a}& \overline{a}& 0& 0& \overline{d}\\ & & & & \\ \overline{b}& 0& \overline{b}& \overline{c}& 0\\ & & & & \\ \overline{c}& \overline{c}& 0& 0& 0\\ & & & & \\ \overline{d}& 0& \overline{d}& 0& 0\end{array}$$
In the case of the algebra $`^n`$, one can choose the generators $`1,\overline{a}_1,\mathrm{}\overline{a}_{n1}`$ with relations:
$$1\overline{a}_i=\overline{a}_i\mathrm{\hspace{0.17em}1}=\overline{a}_i,\overline{a}_i\overline{a}_j=\delta _{ij}.$$
Acknowledgments: We are pleased to thank D. Leites for his help. We are grateful to C. Duval, A. El Gradechi, P. Lecomte and S. Parmentier for enlightening discussions and their interest in this work. The third author is grateful to Swiss National Science Foundation for its support and to A. Alekseev for hospitality.
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# On differentiable compactifications of the hyperbolic plane and algebraic actions of "SL"โโข"("โขโโข")" on surfaces
## Introduction
An important role played by $`\text{SL}_2\text{(}\text{)}`$ is its isometric action on the hyperbolic plane $`^2`$, which can be described as the homogeneous space $`\text{SL}_2\text{(}\text{)}/\text{SO}_2\text{(}\text{)}`$, denoted by $``$. This action is real analytic and is, up to analytic change of coordinates, the only real analytic transitive action of $`\text{SL}_2\text{(}\text{)}`$ on the open disk.
The notion of asymptotic geodesics is a means of understanding the behaviour at infinity of this action, that is to say of giving a natural topological equivariant compactification of this action to an action on the closed disk.
One can ask whether there is a differentiable equivariant compactification of this action into the closed disk. The answer is positive, and there are two well known ways to achieve such a compactification.
The restriction to $`\text{SL}_2\text{(}\text{)}`$ of the natural action of $`\text{SL}_2\text{(}\text{)}`$ on the Riemann sphere $`\overline{}`$ has three orbits: two open hemispheres and between them a great circle. Considering the union of one open orbit and the circle, one gets an analytic equivariant compactification of $``$. We call it the conformal action. It corresponds to the continuous prolongation to the closed unit disk of the $`\text{SL}_2\text{(}\text{)}`$ action on Poincarรฉโs disk.
One can also realize the hyperbolic plane by taking a lorentzian scalar product $`Q`$ on $`^3`$: $`\text{SL}_2\text{(}\text{)}`$ acts isometrically on $`(^3,Q)`$, and when one projectivizes $`^3`$ it gives an analytic action of $`\text{SL}_2\text{(}\text{)}`$ on $`^2`$ with three orbits: an open disk (which is the hyperbolic plane), an open Moebius strip and between them a circle. By taking the action of $`\text{SL}_2\text{(}\text{)}`$ on the union of the disk and the circle we get another analytic equivariant compactification of $``$, called the projective action. It corresponds to the continuous prolongation to the closed unit disk of the $`\text{SL}_2\text{(}\text{)}`$ action on Kleinโs disk.
By uniqueness, we know that these two compactifications are topologically conjugate. However it is easy to check the following surely known but striking fact:
###### Proposition 0.1
The conformal and projective actions are not $`๐^1`$ conjugate, and in particular not $`๐^\omega `$ conjugate.
Proof: if we choose a point $`x`$ of the disk boundary and consider in Poincarรฉโs model the closure of the geodesics which have $`x`$ as an endpoint, we see that all of them are tangent, hence the differential in $`x`$ of the conformal action of the parabolic elements of $`\text{SL}_2\text{(}\text{)}`$ which fix $`x`$ have a common proper direction transversal to the boundary.
If we now consider the same geodesics in Kleinโs model, we see that no two of them are tangent and for each line of the tangent space in $`x`$, there is a closure of a geodesic tangent to it. Hence the differential in $`x`$ of the projective action of a parabolic element of $`\text{SL}_2\text{(}\text{)}`$ which fixes $`x`$ has no proper direction transversal to the boundary.
One can ask whether these two compactifications are the only ones. The answer, stated in a different way, was given by Schneider and Stowe : there exists a countable family of non-equivalent analytic compactifications of $``$, which can be described in terms of infinitesimal generators (see 3.2.1 page 3.2.1). These authors also describe all the analytic actions of $`\text{SL}_2\text{(}\text{)}`$ on compact surfaces with or without boundary and on $`^2`$.
However these new actions seem less natural than to the two compactifications we discussed before, which have well known explicit integral models. Both of these models come in a certain sense from the projectivization of a linear representation; they will be called algebraic in the following sense:
###### Definition 0.2
Let $`k`$ be a positive integer, possibly $`\mathrm{}`$ or $`\omega `$. An action $`\alpha `$ of a Lie group $`G`$ on a manifold possibly with boundary $`M`$ (where $`\alpha `$, $`G`$ and $`M`$ are assumed to be $`๐^k`$) is said to be $`๐^k`$-algebraic if there exists a continuous linear representation $`\stackrel{~}{\rho }`$ of $`G`$ on a real finite dimensional vectorial space $`V`$ and a $`๐^k`$ embedding $`\mathrm{\Phi }:M(V)`$ such that:
* $`\mathrm{\Phi }(M)`$ is a union of orbits for the action $`\rho `$ induced by $`\stackrel{~}{\rho }`$ on $`(V)`$,
* $`\alpha `$ coincides with $`\rho `$ via $`\mathrm{\Phi }`$, that is:
$$\mathrm{\Phi }\alpha (g)=\rho (g)\mathrm{\Phi }gG.$$
The pair $`(\stackrel{~}{\rho },\mathrm{\Phi })`$ is called a $`๐^k`$ algebraic realization of $`\alpha `$.
It is obvious that the projective action is algebraic. The Riemann sphere can be seen as a submanifold of the space of the 2-plans of $`^4`$ which, as a Grassmanian, can be embedded in a real projective space such that the conformal action of $`\text{SL}_2\text{(}\text{)}`$ extends to the projectivization of a linear representation. So the conformal action is algebraic too.
By studying the topology of all the algebraic continuous actions of $`\text{SL}_2\text{(}\text{)}`$ on surfaces and thus determining the regularity of the gluing of the orbits we prove (for a precise definition of โcompactificationโ see 3.2.1):
###### Theorem 0.3
The conformal and projective actions are the only $`๐^\omega `$ compactifications of $``$ which are algebraic.
With this material, we are also able to study all the analytic algebraic actions of $`\text{SL}_2\text{(}\text{)}`$ on surfaces and prove:
###### Theorem 0.4
The analytic algebraic actions of $`\text{SL}_2\text{(}\text{)}`$ on surfaces (with or without boundary) consist exactly of:
* the projective action (on $`^2`$),
* the conformal action (on $`๐^2`$),
* the standard product action on $`^1\times ^1`$,
* one action on the projective plane with an open dense orbit,
* a countable family of actions on the Klein bottle,
* a countable family of actions on the torus with two open cylindric orbits and two circular orbits,
* a countable family of actions on the torus with four open cylindric orbits and four circular orbits,
and of any subaction (i.e. union of orbits) of any one of these actions.
Remark: The realization of these actions as algebraic actions gives explicit global models for all of them.
## 1 The topology of low dimensional algebraic orbits
Our goal is in this section to describe the topology of all orbits of dimension less or equal to 2 which appear in the projectivization of a finite dimensional linear representation of $`\text{SL}_2\text{(}\text{)}`$.
### 1.1 Irreducible representations
All the irreducible representations of $`\text{SL}_2\text{(}\text{)}`$ are known; for a proof of the following theorem, see .
We define a family of linear representations of $`\text{SL}_2\text{(}\text{)}`$. For each non-negative integer $`n`$, $`\stackrel{~}{\rho }_n:\text{SL}_2\text{(}\text{)}_n[X,Y]`$, where $`_n[X,Y]`$ is the vector space of all homogenous polynomials of degree $`n`$ in $`X`$ and $`Y`$, is given by
$$\stackrel{~}{\rho }_n\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)P(X,Y)=P(aX+cY,bX+dY).$$
###### Theorem 1.1
The representation (of dimension $`n+1`$) $`\stackrel{~}{\rho }_n`$ is irreducible for any non-negative $`n`$ and any finite-dimensional irreducible representation of $`\text{SL}_2\text{(}\text{)}`$ is of this form.
### 1.2 Irreducible case
We start the study by the irreducible case.
The irreducible representation of dimension 1, $`\stackrel{~}{\rho }_0`$, is trivial: its associated projective action has one single (fixed !) point.
The irreducible representation of dimension 2, $`\stackrel{~}{\rho }_1`$, gives the obvious action of $`\text{SL}_2\text{(}\text{)}`$ on $`^1`$, which is transitive.
The irreducible representation of dimension 3, $`\stackrel{~}{\rho }_2`$, gives the projective action on $`^2`$, which has three orbits : one open disc, one circle and one Moebius strip. We can determine in which orbit lies the vector line given by a polynomial $`P=aX^2+bXY+cY^2`$ (we denote such a line by $`[aX^2+bXY+cY^2]`$) just by computing the discriminant $`\mathrm{\Delta }=b^24ac`$ (which plays the role of the Lorentzian scalar product in the description of the projective action given in the introduction). The open disk consists of the elements which are not factorizable over $``$ (i.e. of non-positive discriminant). The Moebius strip consists of those which are factorizable with two distinct factors (i.e. of non-negative discriminant). The circle consists of those which are squares (i.e. of zero discriminant).
We denote by $`^+`$ the upper half plane in $``$ and by $`^+`$ its boundary (in Riemannโs sphere $`\overline{}`$). We have a canonical identification between $`^+`$ and $`^1`$, which allows us to identify them.
It is important to notice that, since the map:
$`^+^+`$ $``$ $`(_3[X,Y])`$
$`z`$ $``$ $`[(zX+Y)(\overline{z}X+Y)]`$
is not differentiable on the boundary, it is not an analytic parametrization of the closed disk (union of the open disk orbit and of the circular orbit) and there is no reason to think that the conformal and projective actions on the closed disk are equal up to analytic coordinate change (we already saw that they are not).
Now we generalize this method for all irreducible representations. We shall fix a non-negative integer $`n`$. An element of $`(_n[X,Y])`$ factorizes into the following form:
$$\left[\underset{i=1}{\overset{k}{}}(t_iX+Y)^{\alpha _i}\underset{j=1}{\overset{l}{}}(z_jX+Y)^{\beta _j}(\overline{z_j}X+Y)^{\beta _j}\right]$$
(1)
where $`t_i`$โs are distinct elements of $`^+`$, $`z_j`$โs are disctinct elements of $`^+`$ and $`\alpha _i+2\beta _j=n`$.
Note that $`t_i`$โs are possibly infinite : for example $`[\mathrm{}X+Y]`$ denotes the projective element $`[X]`$.
The form (1) is efficient: since we have
$$\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)(zX+Y)=(cz+d)\left(\frac{az+b}{cz+d}X+Y\right)$$
where $`z\overline{}`$, the conformal action allows one to study all algebraic actions of $`\text{SL}_2\text{(}\text{)}`$topologically.
We shall first determine which orbits are of dimension 2 or less.
###### Lemma 1.2
The orbit of an element $`P`$ written under the form (1) is of dimension 2 or less if and only if: $`k+2l2`$.
Proof: We consider the different cases one by one. By โisometryโ we shall always mean โorientation-preserving isometryโ.
If $`l=1`$ and $`k=1`$, we can write $`P=[(tX+Y)^\alpha (zX+Y)^\beta (\overline{z}X+Y)^\beta ]`$ and the stabilizer of P is the set of the isometries of $`^+`$ (with the hyperbolic metric) which fix the point $`z`$ and the point of the boundary $`t`$, and hence consist only of the identity Id. Thus the orbit of $`P`$ is of the same dimension than $`\text{SL}_2\text{(}\text{)}`$, i.e. 3.
If $`l1`$ and $`k1`$, the same conclusion holds.
If $`l2`$, an element of the component of Id in the stabilizer of $`P`$ must fix at least two points of $`^+`$, hence it is discrete and the orbit of $`P`$ is of dimension 3.
If $`k3`$, an element of the component of Id of the stabilizer of $`P`$ must fix at least three points of the boundary $`^+`$, hence the same conclusion holds.
If $`l=0`$ and $`k=1`$, the stabilizer of $`P`$ is the set of the isometries of $`^+`$ which fix one given point (the only root of a representative polynomial for $`P`$) of the boundary, hence its dimension is 2. Thus the orbit of $`P`$ is one-dimensional.
If $`l=0`$ and $`k=2`$, the stabilizer of $`P`$ is the set of the isometries of $`^+`$ which fix two given points of the boundary, hence it is one-dimensional. Thus the dimension of the orbit of $`P`$ is 2.
If $`l=1`$ and $`k=0`$ the stabilizer of $`P`$ is the set of the isometries of $`^+`$ which fix one given point, hence it is one-dimensional. Thus the dimension of the orbit of $`P`$ is 2.
We have three cases of low dimensional orbits, namely the elliptic case ($`l=1`$ and $`k=0`$), the parabolic case ($`l=0`$ and $`k=1`$) and the hyperbolic case ($`l=0`$ and $`k=2`$).
###### Proposition 1.3
The topology of an orbit of dimension 2 or less of the action $`\rho _n`$ (obtained by projectivizing $`\stackrel{~}{\rho }_n`$) is given by the factorized form (1) of any one of its elements $`P`$ in the following way:
1. if $`l=0`$ and $`k=1`$: the orbit of $`P`$ is a circle
$$\left\{[(tX+Y)^n];t^+\right\}.$$
There is only one such orbit,
2. if $`l=0`$, $`k=2`$ and $`\alpha _1=\alpha _2`$: the orbit of $`P`$ is a Moebius strip
$$\left\{[(t_1X+Y)^\alpha (t_2X+Y)^\alpha ];t_1t_2^+\right\}$$
where $`t_1`$ and $`t_2`$ play the same role. There is one such orbit if $`n`$ is even, none if $`n`$ is odd,
3. if $`l=0`$, $`k=2`$ and $`\alpha _1\alpha _2`$: the orbit of $`P`$ is a cylinder
$$\left\{[(t_1X+Y)^{\alpha _1}(t_2X+Y)^{\alpha _2}];t_1t_2^+\right\}$$
where $`t_1`$ and $`t_2`$ play non-symmetric roles (inverting them maps an element of the orbit to another). There are $`\frac{n1}{2}`$ such orbits if $`n`$ is odd, $`\frac{n2}{2}`$ if $`n`$ is even,
4. if $`l=1`$ and $`k=0`$: the orbit of $`P`$ is a disc
$$\left\{[(zX+Y)^\beta (\overline{z}X+Y)^\beta ];z^+\right\}.$$
There is one such orbit if $`n`$ is even, none if $`n`$ is odd.
Proof: As $`\text{SL}_2\text{(}\text{)}`$ is transitive on $`^+`$ and doubly transitive on $`^+`$, each set described here is an orbit. Thanks to Lemma 1.2 there is no other case than the four mentionned. The computation of the number of orbits is easy with the condition $`\alpha _i+2\beta _j=n`$.
All we have to prove is that the topology of each of these sets is as claimed. The cases 1, 2, 4 can be deduced from the study of $`\rho _2`$ since the map
$`(_m[X,Y])`$ $``$ $`(_{\alpha m}[X,Y])`$
$`[P]`$ $``$ $`[P^\alpha ]`$
is a homeomorphism on its image.
The case 3 reduces to the elementary fact that
$$\left\{(x,y)๐^1\times ๐^1;xy\right\}$$
is a cylinder.
### 1.3 Notations for the reducible case
We shall now consider the reducible representations of $`\text{SL}_2\text{(}\text{)}`$. Since it is a semi-simple Lie group, its finite-dimensional representations are sums of irreducible representations. If we consider a representation $`\stackrel{~}{\rho }`$, we can write: $`\stackrel{~}{\rho }=\stackrel{~}{\rho }_{n_1}\stackrel{~}{\rho }_{n_2}\mathrm{}\stackrel{~}{\rho }_{n_p}`$ for some $`n_1,\mathrm{},n_p`$.
We denote by $`V=_{n_1}[X,Y]_{n_2}[X,Y]\mathrm{}_{n_p}[X,Y]`$ the vector space of $`\stackrel{~}{\rho }`$. Up to a permutation, we can assume that $`n_1n_2\mathrm{}n_p`$.
Moreover, as we want to consider together all the copies of a given irreducible representation which appears in $`\stackrel{~}{\rho }`$ we set $`I_1=i_1=1,i_21`$, $`I_2=i_2,i_31`$, $`\mathrm{}`$, $`I_r=i_r,i_{r+1}1=p`$ the integer intervals such that:
$$\underset{I_1}{\underset{}{n_1=\mathrm{}=n_{i_21}}}>\underset{I_2}{\underset{}{n_{i_2}=\mathrm{}=n_{i_31}}}>\mathrm{}>\underset{I_r}{\underset{}{n_{i_r}=\mathrm{}=n_p}}.$$
We say that $`I_s`$ is even, respectively odd if $`n_{i_s}`$ is even, respectively odd.
We write an element $`x`$ of $`(V)`$ under the factorized form:
$$x=\left[u_q\underset{i=1}{\overset{k_q}{}}(t_q^iX+Y)^{\alpha _q^i}\underset{j=1}{\overset{l_q}{}}(z_q^jX+Y)^{\beta _q^j}(\overline{z_q^j}X+Y)^{\beta _q^j}\right]_{1qp}$$
(2)
where the $`u_q`$โs are real numbers and for each $`q`$: $`\alpha _q^i+2\beta _q^j=n_q`$.
We call support of $`x`$ (or of the projective element $`[u_1,\mathrm{},u_p]`$) and denote by $`I(x)`$ the set of all the intervals $`I_s`$ such that there is at least one index $`iI_s`$, $`u_i0`$. We write $`qI(x)`$ instead of $`q_{II(x)}I`$.
We say that a support is even, respectively odd if all of its elements are even, respectively odd. We define an odd support the same way.
We denote by $`I_+(x)`$ the element of the support of $`x`$ which carries the greatest dimension (i.e. the lowest indices), $`I_{}(x)`$ the one which carries the lowest dimension. We denote by $`q_+(x)`$ (respectively $`q_{}(x)`$) the smallest (respectively the greatest) index $`q`$ such that $`u_q0`$. We have $`q_+(x)I_+(x)`$ and $`q_{}(x)I_{}(x)`$.
When there is no ambiguity, we write $`I_+`$, $`I_{}`$, $`q_+`$ and $`q_{}`$ instead of $`I_+(x)`$, $`I_{}(x)`$, $`q_+(x)`$ and $`q_{}(x)`$.
We denote by $`k(x)`$ (or $`k`$) the number of different $`t_q^i`$โs of $`^+`$ which arise in the factorized form (2) of $`x`$, and $`l(x)`$ (or $`l`$) the number of different $`z_q^j`$โs of $`^+`$.
With these notations we can now generalise the results of the previous section to reducible representations.
###### Lemma 1.4
Let $`x`$ be a element of the projective space $`(V)`$ whose orbit is of dimension 2 or less. Then $`k(x)+2l(x)2`$.
Proof: An element of the identity component of the stabilizer of $`x`$ is an isometry of $`^+`$ stabilizing $`l(x)`$ points and $`k(x)`$ points of the boundary, so we can conclude using the discussion in the proof of Lemma 1.2.
Until the end of the paper, we shall assume there is at least one index $`i`$ such that $`n_i>1`$ (otherwise the action of $`\text{SL}_2\text{(}\text{)}`$ is trivial).
### 1.4 Reducible elliptic case
We assume here that $`k=0`$ and $`l=1`$, that is to say we consider the orbit of an element
$$x=\left[u_q(zX+Y)^{\frac{n_q}{2}}(\overline{z}X+Y)^{\frac{n_q}{2}}\right]_{1qp}$$
which must be of even support.
###### Lemma 1.5
The orbit of an elliptic element is homeomorphic to a disk.
Proof: composing with an element of $`\text{SL}_2\text{(}\text{)}`$, we can assume $`z=ฤฑ`$. Thus the elements of the stabilizer of $`x`$ are exactly the matrices $`\left(\begin{array}{cc}a& b\\ b& a\end{array}\right)`$ where $`a^2+b^2=1`$.
Hence we can parametrize the orbit of $`x`$ by $`z^+`$.
### 1.5 Reducible parabolic case
Now we shall assume $`k=1`$ and $`l=0`$ and consider an element $`x=\left[u_qY^{n_q}\right]`$ (after possible composition with an element of $`\text{SL}_2\text{(}\text{)}`$).
###### Lemma 1.6
The orbit of a parabolic element with support reduced to a single element is homeomorphic to a circle.
The orbit of a parbolic element with support containing at least two elements is homeomorphic to a cylinder.
Proof: if $`A=`$$`\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)`$ stabilizes $`x`$, thus it stabilizes $`0`$ when acting projectively on $`^1`$ hence $`b=0`$ (and $`d=a^1`$).
Moreover we have
$$\left(\begin{array}{cc}a& 0\\ c& a^1\end{array}\right)x=\left[u_qa^{n_q}Y^{n_q}\right]_q.$$
If the support of $`x`$ consists of one single interval $`I_s`$ the condition $`b=0`$ is sufficient for $`A`$ to stabilize $`x`$. If $`d0`$,
$$\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)x=\left[u_q\left(\frac{b}{d}X+Y\right)^{n_q}\right]_{qI_s}$$
else
$$\left(\begin{array}{cc}a& b\\ c& 0\end{array}\right)x=\left[u_qX^{n_q}\right]_{qI_s}$$
Hence the orbit of $`x`$ is homeomorphic to $`^1`$.
If the support of $`x`$ consists of at least to intervals the stabilizer of $`x`$ consist of the matrices of the form $`A=`$$`\left(\begin{array}{cc}1& 0\\ c& 1\end{array}\right)`$ hence the orbit is of dimension 2.
If $`d0`$,
$$\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)x=\left[u_qd^{n_q}\left(\frac{b}{d}X+Y\right)^{n_q}\right]_q$$
else
$$\left(\begin{array}{cc}a& b\\ c& 0\end{array}\right)x=\left[u_qb^{n_q}X^{n_q}\right]_q$$
hence a point of the orbit of $`x`$ is determined by $`\frac{b}{d}^1`$ and a real non-zero parameter, $`b`$ or $`d`$. The case $`d0`$ gives a pair of disjoint copies of $`\times ^{}`$ which are glued along $`d=0`$ into a cylinder. If the support of $`x`$ is neither even nor odd this cylinder is naturally homeomorphic to the orbit of $`x`$, otherwise $`\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)`$$`x=`$ $`\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)`$ and it is naturally a 2-folded covering of the orbit of $`x`$ which is a cylinder too.
### 1.6 Reducible hyperbolic case
We shall assume $`k=2`$ and $`l=0`$ and consider an element
$$x=\left[u_qX^{\alpha _q}Y^{n_q\alpha _q}\right]_q$$
(note that we define $`\alpha _q`$ only when $`u_q0`$).
###### Lemma 1.7
With the notations of this section, a hyperbolic element has a 2 dimensional orbit if and only if $`2\alpha _qn_q`$ is constant, noted $`\delta `$. When this condition is satisfied, the orbit is a Moebius strip if $`\delta =0`$ and $`\alpha _{q_+}\alpha _q`$ is even for each $`q`$, a cylinder otherwise.
Proof: a stabilizing element of $`x`$ must stabilize $`0`$ and $`\mathrm{}`$ in $`\overline{}`$ hence can be written $`\left(\begin{array}{cc}a& 0\\ 0& a^1\end{array}\right)`$. As $`\left(\begin{array}{cc}a& 0\\ 0& a^1\end{array}\right)`$$`x=\left[u_qa^{2\alpha _qn_q}X^{\alpha _q}Y^{n_q\alpha _q}\right]`$ we see that if there are $`q_1`$, $`q_2`$ such that $`2\alpha _{q_1}n_{q_1}2\alpha _{q_1}n_{q_1}`$ thus the orbit of $`x`$ is 3-dimensional, and is 2-dimensional otherwise.
We shall assume we are in the latter case.
Thus the image of $`x`$ under the action of an element $`A\text{SL}_2\text{(}\text{)}`$ is given by the images $`t_1`$ and $`t_2`$ of $`0`$ and $`\mathrm{}`$ under the action of $`A`$ on $`^1`$. If $`\alpha _q=\frac{n_q}{2}`$ for all $`q`$ ($`x`$ is therefore of even support) and $`\alpha _{q_+}\alpha _q`$ is even for all $`q`$ thus exchanging $`t_1`$ and $`t_2`$ gives the same point of the orbit, else it does not.
## 2 Closure of low dimensional algebraic orbits
We shall now determine the closures of the orbits.
By the border of an orbit $`O`$ we mean the set $`\overline{O}O`$.
### 2.1 Elliptic case
We shall consider the orbit of the element $`x`$ which is elliptic, associated to $`ฤฑ`$ and $`[u_q]_q`$, that is : $`x=\left[u_q(ฤฑX+Y)^{\frac{n_q}{2}}(ฤฑX+Y)^{\frac{n_q}{2}}\right]_q.`$
###### Lemma 2.1
The border of the orbit of an elliptic element $`x`$ associated to a projective point $`[u_q]_q`$ is the circular parabolic orbit of $`[u_qY^{n_q}]_{qI_+(x)}`$. The union of these two orbits is a closed disk.
Proof: we have
$$\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)x=\left[u_q\left|cฤฑ+d\right|^{n_qn_{q_+}}\left(\frac{aฤฑ+b}{cฤฑ+d}X+Y\right)^{\frac{n_q}{2}}\left(\overline{\frac{aฤฑ+b}{cฤฑ+d}}X+Y\right)^{\frac{n_q}{2}}\right]_q$$
Since $`adbc=1`$ we can write:
$$\frac{aฤฑ+b}{cฤฑ+d}=\frac{ac+bd}{\left|cฤฑ+d\right|^2}+ฤฑ\frac{1}{\left|cฤฑ+d\right|^2}$$
thus $`\left|cฤฑ+d\right|^2=(\mathrm{Im}z)^1`$, and hence the orbit is the set of the elements
$$x(z)=\left[u_q(\mathrm{Im}z)^{\frac{n_{q_+}n_q}{2}}\left(zX+Y\right)^{\frac{n_q}{2}}\left(\overline{z}X+Y\right)^{\frac{n_q}{2}}\right]_q$$
where $`z^+`$.
If a sequence $`(x(z_i))_i`$ has a limit in $`(V)`$, necessarily $`(z_i)_i`$ has a limit in the closure of $`^+`$ in $`\overline{}`$. If this limit is in $`^+`$ we get a point of the orbit of $`x`$, otherwise it is a point $`t^+`$. In the latter case, if $`t`$ is finite, $`\mathrm{Im}z_i`$ has limit zero and $`(x(z_i))_i`$ has limit $`[u_q(tX+Y)^{n_q}]_{qI_+(x)}`$. If $`t=\mathrm{}`$, $`(x(z_i))_i`$ has limit $`[u_qX^{n_q}]_{qI_+(x)}`$, which we can write $`[u_q(\mathrm{}X+Y)^{n_q}]_{qI_+(x)}`$.
### 2.2 Parabolic case
The circular orbits are closed, so we consider only the two types of cylindric orbits; as the technic is the same than in the elliptic case, we shall not give much detail.
###### Lemma 2.2
Let $`x=[u_qY^{n_q}]_q`$ be of even non-reduced to a single element support. The border of the cylindric orbit of $`x`$ is the disjoint union of the orbits of $`[u_qY^{n_q}]_{qI_+(x)}`$ and $`[u_qY^{n_q}]_{qI_{}(x)}`$.
If the support of $`x`$ has a parity (i.e. is even or odd), the closure of the orbit of $`x`$ is a closed cylinder if $`n_q_{}>0`$ and a closed disk if $`n_q_{}=0`$.
If the support of $`x`$ is neither odd nor even, the closure of the orbit of $`x`$ is a Klein bottle if $`n_q_{}>0`$ and a projective plane if $`n_q_{}=0`$.
Proof: we shall consider the orbit of an element $`x=[u_qY^{n_q}]_q`$ whose support is even and has at least two elements. This orbit is described in Section 1.5, we can write it under the form:
$`\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)x`$ $`=`$ $`\left[u_qd^{n_qn_{q_\pm }}\left({\displaystyle \frac{b}{d}}X+Y\right)^{n_q}\right]_q\text{ if }d0,`$
$`\left(\begin{array}{cc}a& b\\ c& 0\end{array}\right)x`$ $`=`$ $`\left[u_qb^{n_qn_{q_\pm }}X^{n_q}\right]_q`$
where we choose $`\pm `$ to be $`+`$ (respectively $``$) if we want to study great (respectively small) values of the real parameter given for a choosen $`t=\frac{b}{d}^1`$ by $`d`$ (or $`b`$ if $`t=\mathrm{}`$).
For great values, we find a point of the circular orbit of $`[u_qY^{n_q}]_{qI_+(x)}`$, for small ones a point of the orbit of $`[u_qY^{n_q}]_{qI_{}(x)}`$ (which is a circle if $`n_{q_{}(x)}>0`$, a single point otherwise).
The way the cylindric orbit is glued on the circles of its border depends of the parity of the support of $`x`$: if it has a parity (i.e. is even or odd) the couples $`(b,d)`$ and $`(b,d)`$ of parameters give the same point, else they give two different points such that if one of them is close to a point of the border, the other is close to this point too: hence the cylinder will glue twice on each circle in its border.
### 2.3 Hyperbolic case
###### Lemma 2.3
The border of the orbit $`O`$ of an element $`x=\left[u_qX^{\alpha _q}Y^{n_q\alpha _q}\right]_q`$ (where $`2\alpha _qn_q`$ does not depend upon $`q`$) is the circular orbit of $`[u_qY^{n_q}]_{qI_+(x)}`$.
If $`O`$ is a Moebius strip, its closure is a closed Moebius strip.
If $`O`$ is a cylinder, its closure is a torus.
Proof: we can write this orbit as the set of all elements of the form
$$\left[u_q(t_1t_2)^{\alpha _{q_+}\alpha _q}\left(t_1X+Y\right)^{\alpha _q}\left(t_2X+Y\right)^{\beta _q}\right]_q$$
$$=\left[u_q\left(\frac{1}{t_2}\frac{1}{t_1}\right)^{\alpha _{q_+}\alpha _q}\left(X+\frac{1}{t_1}Y\right)^{\alpha _q}\left(X+\frac{1}{t_2}Y\right)^{\beta _q}\right]_q$$
with $`t_1,t_2^1`$. As before, this enables the description of the border of this orbit.
## 3 Classification of analytic algebraic action of $`\text{SL}_2\text{(}\text{)}`$ on surfaces
We shall now study the analyticity of the different topological surfaces obtained as a union of orbits and which are analytically conjugate (i.e. are equal up to an analytic change of coordinates).
### 3.1 Smoothness of polynomial-parametrized surfaces
We shall use many times the following result, which can be generalized (but we present here only the 2-dimensional version for simplicity).
###### Proposition 3.1
Let $`P:(x_1,x_2)(P_1(x_1,x_2),\mathrm{},P_n(x_1,x_2))`$ be a map defined on a neigborhood of 0 in $`^2`$ where the $`P_i`$โs are homogeneous non-constant polynomials. We assume $`P_1`$ to be of minimal degree and $`P_2[P_1]`$ of minimal degree among $`P_i`$โs with that property. If there exists some $`P_i[P_1,P_2]`$ then the image $`E`$ of $`P`$ is not a smooth 2-dimensional submanifold of $`^n`$ (more precisely, $`P`$ is singular at 0).
Proof: Assume that $`E`$ is a smooth 2-dimensional submanifold of $`^n`$. Thus there is a smooth implicit definition of $`E`$, that is to say a neighborhood $`U`$ of $`E`$ in $`^n`$ and a smooth map $`h:U^{n2}`$ of rank $`n2`$ everywhere such that $`E=\left\{xU;h(x)=0\right\}`$.
Moreover, assume there is a polynomial $`P_{i_0}[P_1,P_2]`$ (we choose it of minimal degree).
Let $`d`$ be the degree of $`P_1`$. We consider the taylor developpement of order 1 of $`h`$ in 0 and estimate it in $`(P_1(x_1,x_2),\mathrm{},P_n(x_1,x_2))`$. Noting $`h_j`$ the $`j^{th}`$ coordinate fonction of $`h`$ and $`_i`$ the derivation in the $`i^{th}`$ variable, we get for each $`j`$ :
$$0=\underset{i}{}_ih_j(0)P_i(x_1,x_2)+o(x_1,x_2^d)$$
where the sum is taken over the $`P_i`$โs of degree $`d`$, hence
$$\left(\begin{array}{ccc}_1h_1(0)& \mathrm{}& _nh_1(0)\\ \mathrm{}& & \mathrm{}\\ _1h_{n2}(0)& \mathrm{}& _nh_{n2}(0)\end{array}\right)\left(\begin{array}{c}P_1\\ P_2\text{ if it is of degree }d\text{, 0 otherwise}\\ \mathrm{}\\ P_i\text{ if it is of degree }d\text{, 0 otherwise}\\ \mathrm{}\\ P_n\text{ if it is of degree }d\text{, 0 otherwise}\end{array}\right)=0.$$
Each line in the second matrix is given by the coefficients of the polynomial.
First assume that $`P_2`$ and $`P_{i_0}`$ are both of degree $`d`$. Thus the family of $`P_i`$โs of degree $`d`$ is of rank at least 3, hence the jacobian matrix of $`h`$ at the point 0 is of rank at most $`n3`$ which prevent $`h`$ from being an implicit definition of $`E`$.
Next assume that $`P_2`$ is of degree $`d`$ and $`P_{i_0}`$ of degree $`d_0`$ greater than $`d`$. Thus $`h`$ is of corank at least 2 at the point 0: we have two independent linear combinations of the $`_ih(0)`$โs which must be zero and involve only the indices $`i`$ of degree $`d`$ polynomials. But we can now use the Taylor developpement of order $`d_0`$ to get for each $`j`$:
$$0=\underset{i}{}_ih_j(0)P_i(x_1,x_2)+Q_j(x_1,x_2)$$
where the sum is taken over all polynomials of degree $`d_0`$ which are not in $`[P_1,P_2]`$ and $`Q_j`$ is a polynomial of degree $`d_0`$ of $`[P_1,P_2]`$. Let $`S`$ be, in the vector space of all homogenous polynomials of degree $`d_0`$, a supplementary of the space $`^{d_0}[P_1,P_2]`$ of those of $`[P_1,P_2]`$. Let $`P_i^{}`$ be the projection of $`P_i`$ on $`S`$ along $`^{d_0}[P_1,P_2]`$. Thus we have for each $`j`$:
$$0=\underset{i}{}_ih_j(0)P_i^{}(x_1,x_2)$$
where the sum is taken over all polynomials of degree $`d_0`$ which are not in $`[P_1,P_2]`$. As before, it gives a linear combination of the $`_ih(0)`$โs which must be zero, and is independent of the two we get previously as $`P_{i_0}[P_1,P_2]`$. Hence $`h`$ is of corank at least 3 in O and the contradiction holds as before.
We can use the same proof for the case when $`P_2`$ is of degree greater than $`d`$.
### 3.2 Compactifications of the hyperbolic plane: the elliptic case
#### 3.2.1 Analytic non necessarily algebraic compactification
We shall start with a description of all analytic compactifications of $``$ into a closed disk, in the following sense:
###### Definition 3.2
A differentiable compactification of a differentiable action $`\alpha `$ of a Lie group $`G`$ on a manifold $`M`$ is a triple $`(N,\varphi ,\overline{\alpha })`$ where $`N`$ is a compact manifold with boundary, $`\varphi :MN`$ is an embedding and $`\overline{\alpha }`$ is a differentiable action of $`G`$ on $`N`$ such that $`\varphi (M)`$ is dense in $`N`$ and $`\overline{\alpha }`$ is a prolongation of the action induced by $`\alpha `$ on $`\varphi (M)`$.
The work of Schneider , Stowe exposed by Mitsumatsu gives immediately the classification of all such compactifications, which we recall in what follows.
We shall use the following basis for $`๐ฐ๐ฉ_2\text{(}\text{)}`$:
$$H=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right),K=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right),L=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right).$$
The infinitesimal generators for the projective compactification are given on $`\times _+`$ by
$`\overline{K}_{1+}`$ $`=`$ $`2{\displaystyle \frac{}{x}}`$
$`\overline{H}_{1+}`$ $`=`$ $`2\left((\mathrm{sin}x)(1+y){\displaystyle \frac{}{x}}+(\mathrm{cos}x)(2y+y^2){\displaystyle \frac{}{y}}\right)`$
$`\overline{L}_{1+}`$ $`=`$ $`2\left((\mathrm{cos}x)(1+y){\displaystyle \frac{}{x}}(\mathrm{sin}x)(2y+y^2){\displaystyle \frac{}{y}}\right).`$
and can be completed by adding a point at infinity.
###### Theorem 3.3 ()
By pulling back the restriction of the vector fields $`\overline{K}_{1+}`$, $`\overline{H}_{1+}`$, $`\overline{L}_{1+}`$ to $`\times _+^{}`$ by the map $`F_n(x,y)=(x,y^n)`$ where $`n`$ is a non-negative integer and by taking their continuous prolongations, we get analytic vector fields $`\overline{K}_{n+}`$, $`\overline{H}_{n+}`$, $`\overline{L}_{n+}`$ on $`\times _+`$. For any analytic compactifications of $``$ into a closed disc, there is an unique $`n`$ and a $`\times _+`$ chart in which these vector fields are the infinitesimal generators of the compactified action.
For example, $`\overline{K}_{2+}`$, $`\overline{H}_{2+}`$, $`\overline{L}_{2+}`$ are the infinitesimal generators for the conformal compactification.
#### 3.2.2 Analytic algebraic compactifications
We shall now study the algebraic analytic compactifications of $``$ into a closed disc, that is to say the elliptic orbits whose closure is an analytic submanifold with boundary in the projective space $`(V)`$.
We prove a more precise version of the theorem 0.3 exposed in the introduction:
###### Theorem 3.4
Let $`O`$ be the orbit of $`x=\left[u_q(ฤฑX+Y)^{\frac{n_q}{2}}(ฤฑX+Y)^{\frac{n_q}{2}}\right]_q.`$
If all the element of the family $`(\frac{n_{q_+}n_q}{2})_{qI(x)}`$ are even, thus $`\overline{O}`$ is an analytic submanifold with boundary and the action of $`\text{SL}_2\text{(}\text{)}`$ on this disk is conjugate to the projective action.
If there exists some $`q_{2+}`$ in $`I(x)`$ such that $`\frac{n_{q_+}n_{q_{2+}}}{2}=1`$, thus $`\overline{O}`$ is an analytic submanifold with boundary and the action of $`\text{SL}_2\text{(}\text{)}`$ on this disk is conjugate to the conformal action.
In all the other cases, $`\overline{O}`$ is not an analytic submanifold with boundary.
Proof: The methods used here will be useful through all the following sections.
We shall first consider the case when all the numbers $`\frac{n_{q_+}n_q}{2}`$, where $`q`$ is in $`I(x)`$, are even. A model for the projective compactification is given by the closure in $`(_2[X,Y])`$ of the orbit of $`[X^2+Y^2]`$, which is contained in the affine chart $`\left\{[aX^2+bXY+(1a)Y^2];a,b\right\}`$. The map
$`\phi :(_2[X,Y])`$ $``$ $`(V)`$
$`\left[aX^2+bXY+(1a)Y^2\right]`$ $``$ $`[u_q(a(1a){\displaystyle \frac{b^2}{4}})^{\frac{n_{q_+}n_q}{4}}`$
$`(aX^2+bXY+(1a)Y^2)^{\frac{n_q}{2}}]_q`$
is injective, analytic (thanks to the hypothesis) and realizes a conjugacy between the projective action and the dynamics on $`\overline{O}`$.
Moreover, it is an immersion since, noting $`s,t,u,v`$ the coefficients of the terms in $`X^{n_{q_+}},X^{n_{q_+}1}Y,Y^{n_{q_+}},XY^{n_{q_+}1}`$, we have $`\frac{s}{a}=\frac{n_{q_+}}{2}a^{\frac{n_{q_+}}{2}1},\frac{u}{a}=\frac{n_{q_+}}{2}(1a)^{\frac{n_{q_+}}{2}1}`$ and $`\frac{s}{b}=0,\frac{t}{b}=\frac{n_{q_+}}{2}a^{\frac{n_{q_+}}{2}1},\frac{u}{b}=0,\frac{v}{b}=\frac{n_{q_+}}{2}(1a)^{\frac{n_{q_+}}{2}1}.`$
Hence the differential of $`\phi `$ is of rank 2 everywhere.
This proves that $`\overline{O}`$ is an analytic submanifold with boundary and at the same time that the action of $`\text{SL}_2\text{(}\text{)}`$ on it is conjugate to the projective one.
Next we shall consider the case when there exists some $`q_{2+}`$ in $`I(x)`$ such that $`\frac{n_{q_+}n_{q_{2+}}}{2}=1`$. A model for the conformal action is given by the closure of $`^+`$ in the Riemmann sphere. We consider the map
$`\psi :\overline{^+}`$ $``$ $`(V)`$
$`a+ฤฑb`$ $``$ $`\left[u_qb^{\frac{n_{q_+}n_q}{2}}((a+ฤฑb)X+Y)^{\frac{n_q}{2}}((aฤฑb)X+Y)^{\frac{n_q}{2}}\right]_q`$
which is injective, analytic and realizes a conjugacy between the conformal action and the dynamics on $`\overline{O}`$. Notice that $`\psi (\mathrm{})=[u_qX^{n_q}]_{qI_+(x)}`$.
Moreover developping the expression of $`\psi (a+ฤฑb)`$, we see that a coefficient is $`n_{q_+}a`$ and another is $`u_{q_{2+}}b`$, so $`\psi `$ is everywhere of rank 2 and we can conclude as before.
For the last case, we use Proposition 3.1. We denote by $`\alpha `$ the smallest odd element of the family $`(\frac{n_{q_+}n_q}{2})_q`$, we denote by $`q_{2+}`$ an index realizing this minimum. By hypothesis $`\alpha >1`$. We can write an element of $`\overline{O}`$ under the form: $`\left[u_q(\mathrm{Im}z)^{\frac{n_{q_+}n_q}{2}}\left((\mathrm{Im}z^2+\mathrm{Re}z^2)X^2+2\mathrm{R}\mathrm{e}zXY+Y^2\right)^{\frac{n_q}{2}}\right]_q.`$ All coordinates are homogeneous polynomials in $`x=\mathrm{Re}z`$ and $`y=\mathrm{Im}z`$. Among them $`P_1=x`$ (we define it up to a multiplicative constant) is of minimal degree. Among those which are not in $`[P_1]`$, $`P_2=y^2`$ is of minimal degree. But $`P_3=y^\alpha [P_1,P_2]`$ hence $`\overline{O}`$ is not a smooth submanifold of $``$(V), therefore not an analytic one.
###### Remark 3.5
In this proof we can see more than stated: the embeddings $`\phi `$ and $`\psi `$ extend respectively to embeddings of a projective plane (union of the elliptic orbit of $`x`$, the hyperbolic orbit of $`\left[\left(\frac{1}{4}\right)^{\frac{n_{q_+}n_q}{4}}u_qX^{\frac{n_q}{2}}Y^{\frac{n_q}{2}}\right]_q`$ which is a Moebius strip and their common border, the circular orbit of $`[u_qY^{n_q}]_{qI_+(x)}`$) and a sphere (union of the elliptic orbits of $`x`$ and of $`\left[(1)^{\frac{n_{q_+}n_q}{2}}u_q(X^2+Y^2)^{\frac{n_q}{2}}\right]_q`$ and of their common border, the circular orbit of $`[u_qY^{n_q}]_{qI_+(x)}`$).
Moreover, we see that if we are in the third case, the map $`\phi `$ is not analytic but is a $`๐^{\frac{\alpha 1}{2}}`$ embedding of the projective action, so we can state the following fact concerning the differentiable case for elliptic orbits:
###### Theorem 3.6
The only algebraic differentiable compactifications of $``$ are equivalent to the projective or to the conformal ones. In the projective case there exist $`๐^k`$ non-analytic realizations for each finite $`k`$, but any $`๐^{\mathrm{}}`$ realization is in fact analytic. In the conformal case any $`๐^1`$ realization is in fact analytic.
### 3.3 Hyperbolic case
Here we shall consider the closure of a hyperbolic 2-dimensional orbit, which has the form
$$\overline{O}=\left\{\left[u_q(t_1t_2)^{\alpha _{q_+}\alpha _q}(t_1X+Y)^{\alpha _q}(t_2X+Y)^{n_q\alpha _q}\right]_q;t_1,t_2^1\right\}.$$
###### Theorem 3.7
If $`O`$ is a Moebius strip (i.e for each $`q`$, $`n_q`$ is even, $`\alpha _q=\frac{n_q}{2}`$ and $`\alpha _{q_+}\alpha _q`$ is even), $`\overline{O}`$ is an analytic submanifold; moreover its union with the elliptic orbit of $`\left[\left(\frac{1}{4}\right)^{\frac{n_{q_+}n_q}{4}}u_q(X^{n_q}+Y^{n_q})\right]`$ is still analytic and the dynamics is conjugate to the projective action of $`\text{SL}_2\text{(}\text{)}`$ on the projective plane.
If there is some $`q_{2+}`$ such that $`\alpha _{q_+}\alpha _{q_{2+}}=1`$, $`\overline{O}`$ is an analytic submanifold of $`(V)`$ and its dynamics is conjugate to the natural product action of $`\text{SL}_2\text{(}\text{)}`$ on $`^1\times ^1`$.
In all the other cases, $`\overline{O}`$ is not an analytic submanifold.
Proof: The first case is given by the map $`\phi `$ of the previous section (see Remark 3.5).
In the second case, we consider the map
$`\psi :^1\times ^1`$ $``$ $`(V)`$
$`(t_1,t_2)`$ $``$ $`\left[u_q(t_1t_2)^{\alpha _{q_+}\alpha _q}\left(t_1X+Y\right)^{\alpha _q}\left(t_2X+Y\right)^{n_q\alpha _q}\right]_q`$
which is analytic, injective as the orbit is by hypothesis a cylinder and is an immersion as the coefficient of the terms in $`XY^{n_{q_+}}`$ and $`Y^{n_{q_{2+}}}`$ of $`\psi (t_1,t_2)`$ are respectively $`\alpha _{q_+}t_1+(n_q\alpha _{q_+})t_2`$ and $`t_1t_2`$, which gives a partial jacobian matrix $`\left(\begin{array}{cc}\alpha _{q_+}& n_q\alpha _{q_+}\\ 1& 1\end{array}\right)`$ whose determinant is $`n_{q_+}0`$. Hence $`\overline{O}`$ is an analytic submanifold (without boundary) of $`(V)`$ and (see the topological study) its dynamics is conjugate to the product action of $`\text{SL}_2\text{(}\text{)}`$ on $`^1\times ^1`$.
For the last case we use Proposition 3.1. The only polynomial of degree 1 among the coordinates is $`P_1=\alpha t_1+\beta t_2`$ where we write $`\alpha `$ for $`\alpha _{q_+}`$ and $`\beta `$ for $`n_{q_+}\alpha _{q_+}`$. We can next choose $`P_2=\frac{\alpha (\alpha 1)}{2}t_1^2+\alpha \beta t_1t_2+\frac{\beta (\beta 1)}{2}t_2^2`$. Setting $`P_2^{}=(t_1t_2)^2`$, an easy computation gives $`[P_1,P_2]=[P_1,P_2^{}]`$.
If $`\alpha =\beta `$, as $`\overline{O}`$ is assumed to be a cylinder there must exist some index $`q_0`$ such that $`\alpha _{q_+}\alpha _{q_0}`$ is odd. Thus one of the coordinates has the form $`(t_1t_2)^{\alpha _{q_+}\alpha _{q_0}}`$ which is not in $`[P_1,P_2^{}]`$, hence from Proposition 3.1 we conclude that $`\overline{O}`$ is not an analytic submanifold of $`(V)`$.
If $`\alpha \beta `$, we see after an easy computation that the coordinate $`P_3=\frac{\alpha (\alpha 1)(\alpha 2)}{6}t_1^3+\frac{\alpha (\alpha 1)}{2}\beta t_1^2t_2+\alpha \frac{\beta (\beta 1)}{2}t_1t_2^2+\frac{\beta (\beta 1)(\beta 2)}{6}t_2^3`$ of the term $`X^3Y^{n_q3}`$ is not in $`[P_1,P_2]`$ and the conclusion still holds.
### 3.4 Parabolic case
We shall finally consider the closure of a parabolic orbit, which has the form $`\overline{O}=\left\{\left[u_qd^{n_qn_q_{}}(tX+Y)^{n_q}\right]_q;d\overline{}\text{ and }t^1\right\}`$ where $`d\overline{}`$ means $`d`$ is real or $`\pm \mathrm{}`$.
We shall prove some lemmas before stating the general result. Let $`q_2`$ (respectively $`q_{2+}`$) be an index such that $`n_{q_2}`$ (respectively $`q_{2+}`$) is minimal (respectively maximal) among $`n_q`$โs greater than $`n_q_{}`$ (respectively lesser than $`n_{q_+}`$).
###### Lemma 3.8
If $`n_q_{}=0`$ and $`\overline{O}`$ is a smooth submanifold of $`(V)`$, we must have $`n_{q_2}=1`$ and hence $`\overline{O}`$ is a projective plane.
Proof: We shall use Proposition 3.1 once again, around the point $`[u_q]_{qI_{}}`$ corresponding to $`d=0`$, $`t=0`$. The least-dimensional non-constant polynomial among the local coordinates is $`P_1=d^{n_{q_2}}`$. There is no other polynomial of the same degree, so we can choose $`P_2=tP_1[P_1]`$. If $`n_{q_2}>1`$, one of the coordinates can be written as $`t^2P_1[P_1,P_2]`$ and $`\overline{O}`$ can not be a smooth submanifold of $`(V)`$.
###### Lemma 3.9
If $`\overline{O}`$ is a smooth submanifold of $`(V)`$, we must have
* $`n_{q_+}n_{q_{2+}}=n_{q_2}n_q_{}`$,
* for each $`q`$, $`n_{q_+}n_{q_{2+}}`$ divides $`n_{q_+}n_q`$.
Proof: We use Proposition 3.1 twice.
We first look around the point $`[u_qY^{n_q}]_{qI_{}}`$ to prove that for each $`q`$, $`n_{q_2}n_q_{}`$ divides $`n_qn_q_{}`$. If $`n_q_{}=0`$, we have $`n_{q_2}=1`$ and the claim is obvious. If $`n_q_{}>0`$, we can choose $`P_1=t`$ and $`P_2=d^{n_{q_2}n_q_{}}`$. For each $`q`$ there is a coordinate which has the form $`d^{n_qn_q_{}}`$, hence by Proposition 3.1 $`n_{q_2}n_q_{}`$ must divide $`n_qn_q_{}`$.
In particular $`n_{q_2}n_q_{}`$ divides $`n_{q_+}n_{q_{2+}}`$.
We now look around the point $`[u_qY^{n_q}]_{qI_+}`$, where local coordinates are given by writting a point of $`\overline{O}`$ under the form $`\left[u_qe^{n_{q_+}n_q}(tX+Y)^{n_q}\right]_q`$ after a change of coordinates $`e=d^1`$. We can choose $`P_1=t`$ and $`P_2=e^{n_{q_+}n_{q_{2+}}}`$, thus as there is coordinates of the form $`e^{n_{q_+}n_q}`$, for all $`q`$, $`n_{q_+}n_{q_{2+}}`$ divides $`n_{q_+}n_q`$.
In particular $`n_{q_+}n_{q_{2+}}`$ divides $`n_{q_2}n_q_{}`$ and the conclusion holds.
It is easy to see that the necessary conditions given in the previous lemma are also sufficient if $`n_q_{}0`$ for $`\overline{O}`$ to be an analytic submanifold of $`(V)`$: around each point of $`\overline{O}`$ we can find local coordinates of the form $`P_{k,l}=d^{k(n_{q_+}n_{q_{2+}})}t^l`$ where $`k`$ and $`l`$ are integers and for some coordinates we have $`(k,l)=(0,1)`$ or $`(k,l)=(1,0)`$, hence writting $`P_{k,l}P_{1,0}^{}{}_{}{}^{k}P_{0,1}^{}{}_{}{}^{l}=0`$ we get an analytic implicit local definition of $`\overline{O}`$. If $`n_q_{}=0`$ the combination of the conditions of the two lemmas are also sufficient for $`\overline{O}`$ to be analytic since we can find local coordinates of the previous form or, around the points given by $`d=0`$, of the form $`P_{k,l}=d^kt^l`$ with $`k>0,kl`$; for some coordinates we have $`(k,l)=(1,1)`$ and $`(k,l)=(1,0)`$ hence we get an analytic implicit local definition of the form $`P_{k,l}P_{1,1}^{}{}_{}{}^{l}P_{1,0}^{}{}_{}{}^{kl}=0`$.
Moreover, if we map a point given by parameters $`d,t`$ from the closure of an analytic parabolic orbit to the point given by the same parameters on another such orbit closure of the same topology (projective plane, Klein bottle or cylinder) and with the same value for $`n_{q_2}n_q_{}`$ we build an analytic diffeomorphism between them:
$$\left[u_qd^{n_qn_q_{}}(tX+Y)^{n_q}\right]_q(d^{n_{q_2}n_q_{}},t)\left[u_q^{}d^{n_qn_q_{}^{}}(tX+Y)^{n_q}\right]_q.$$
Finally, if we consider the differential in the point $`x=[u_qY^{n_q_{}}]_{qI_{}}`$ of an element $`\left(\begin{array}{cc}a& 0\\ c& a^1\end{array}\right)`$ of the stabilizer of $`x`$ we find that its eigenvalues are $`a^2`$ and $`a^{(n_{q_2}n_q_{})}`$, so two closures of orbits with different values of $`n_{q_2}n_q_{}`$ can not be differentiably conjugate. Hence we can state:
###### Theorem 3.10
The conditions of Lemmas 3.8 and 3.9 are sufficient for $`\overline{O}`$ to be an analytic submanifold of $`(V)`$. Two analytic parabolic orbits are analyticaly conjugate if and only if they have the same topology and the same value for $`n_{q_2}n_q_{}`$ (and they are not even differentiably conjugate otherwise). In particular there is one parabolic algebraic action on the projective plane, a countable family of actions on the Klein bottle and a countable family of actions on the closed cylinder.
The last point we have to study in order to complete the proof of the results stated in the introduction is the way the cylindric orbits are glued together.
Let $`O`$ be a cylindric analytic orbit associated to a projective element $`[u_q]_q`$. Its boundary is the union of the two circular orbits associated to the projective elements $`[u_q]_{qI_+}`$ and $`[u_q]_{qI_{}}`$, which we call respectively the upper component and the lower component of the boundary.
An element of $`\overline{O}`$ can be writen $`\left[u_qd^{n_qn_q_{}}(tX+Y)^{n_q}\right]_q`$ around the lower component of the boundary. For each $`q`$ we denote by $`k_q`$ the integer $`\frac{n_qn_q_{}}{n_{q_2}n_q_{}}`$. The coordinates $`c_{q,l}=u_qd^{n_qn_q_{}}t^l`$ satisfy the implicit definition given previously:
$$\frac{1}{u_q}c_{q,l}\frac{1}{u_{q_2}}c_{q_2,0}^{}{}_{}{}^{k_q}\frac{1}{u_q_{}n_q_{}}c_{q_{},1}^{}{}_{}{}^{l}=0.$$
Let $`O^{}`$ be the cylindric analytic orbit associated with the projective element $`[u_q^{}]_q`$ where $`u_q^{}=(1)^{k_q}u_q`$. Thus the lower component of its boundary is the same than for $`O`$ and as around it the coordinates of $`O^{}`$ satisfy the same implicit parametrization, $`O`$ and $`O^{}`$ are analytically glued together around their lower component.
With the same method we see that $`O`$ and the orbit $`O^{\prime \prime }`$ associated with $`[u_q^{\prime \prime }]_q`$ where $`u_q^{\prime \prime }=(1)^{k_{q_+}k_q}u_q`$ are analytically glued around their common upper component.
If $`k_{q_+}`$ is even $`O^{}=O^{\prime \prime }`$ and $`O`$ together with $`O^{}`$ gives a torus with two open orbits, if $`k_{q_+}`$ is odd $`O^{}O^{\prime \prime }`$ but they are both glued analytically with $`O^{\prime \prime \prime }`$, the parabolic orbit associated with $`[(1)^{k_{q_+}}u_q]_q`$. Hence we have proven the last remaining result:
###### Theorem 3.11
Let $`O`$ be a parabolic, cylindric, analytic orbit associated to $`[u_q]_q`$.
If $`k_{q_+}=\frac{n_{q_+}n_q_{}}{n_{q_2}n_q_{}}`$ is even, the union of the two parabolic orbits associated to $`[u_q]_q`$ and $`[(1)^{k_q}u_q]_q`$ is a torus analytically embedded in $`(V)`$.
If $`k_{q_+}`$ is odd, the union of the four parabolic orbits associated to $`[u_q]_q`$, $`[(1)^{k_q}u_q]_q`$, $`[(1)^{k_{q_+}k_q}u_q]_q`$ and $`[(1)^{k_{q_+}}u_q]_q`$ is a torus analytically embedded in $`(V)`$.
UMPA, รNS Lyon
46, allรฉe dโItalie
69 364 Lyon cedex 07
France
bkloeckn@umpa.ens-lyon.fr
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# The gauge non-invariance of Classical Electromagnetism
## I Introduction
In Classical electromagnetism, the electric field $`๐`$ and the magnetic field $`๐`$ are related to the scalar $`V`$ and vector $`๐`$ potentials by the following definitions Jackson :
$$๐=\frac{๐}{t}Vand๐=\times ๐$$
(1)
One century ago, H.A. Lorentz noticed that the electromagnetic field remains invariant ($`๐^{}=๐`$ and $`๐^{}=๐`$) under the so-called gauge transformations Okun :
$$๐^{}=๐+fandV^{}=V\frac{f}{t}$$
(2)
where $`f(x,t)`$ is the gauge function.
Hence, this indeterminacy is believed to be an essential symmetry of Classical Electromagnetism Okun . Moreover, it is often related to the assertion that the potentials are not measurable quantities contrary to the fields. Hence, one must specify what is called a gauge condition, that is a supplementary equation which is injected in the Maxwell equations expressed in function of the electromagnetic potentials in order to supress this indeterminacy. It is common to say that these gauge conditions are mathematical conveniences that lead to the same determination of the electromagnetic field. In this context, the choice of a specific gauge condition is motivated from the easiness in calculations compared to another one. In a certain manner, although their mathematical expressions are different, it is supposed that they are equivalent as the fields are invariant with respect to the gauge transformations. Furthermore, no physical meaning is ascribed to the gauge conditions as the potentials are assumed not to have oneโฆ
Despite these assertions which are shared by a large majority of physicists, a definition for the potentials dating back to Maxwell was recalled recently and which resolves, according to our point of view, the question of indeterminacy by giving them a physical interpretation. Moreover, we showed that the Coulomb and Lorenz gauge conditions were, in fact, not equivalent because they must be interpreted as physical constraints that is electromagnetic continuity equations Guyon ; Rousseaux . In addition, we were able to demonstrate that the Coulomb gauge condition is the galilean approximation of the Lorenz gauge condition within the magnetic limit of Lรฉvy-Leblond & Le Bellac Levy ; ARQS ; RL . So, to โmake a gauge choiceโ that is choosing a gauge condition is, as a consequence of our findings, not related to the fact of fixing a special couple of potentials. Gauge conditions are completely uncorrelated to the supposed indeterminacy of the potentials. Hence, we proposed to rename โgauge conditionโ by โconstraintโ Guyon ; Rousseaux ; RL .
In this article, we would like to reexamine the common belief concerning the assumed indeterminacy of the potentials with the assumption that the โconstraintโ do not fix the value of the potentials. Indeed, we will show that gauge transformations introduce paradoxes which imply their rejection. This point of view was expressed already by the school of De Broglie : by the master himself Broglie or by his followers like Costa De Beauregard Costa
## II The case of a stationnary electric field
Imagine a one dimensional stationary problem defined by the following potentials :
$$๐=\mathrm{๐}andV(x)=Ex$$
(3)
One finds easily :
$$๐=\mathrm{๐}and๐=E๐_๐ฑ$$
(4)
The electric and magnetic fields are constants in time.
Now, we can perform a gauge transformation with this particular gauge function :
$$f(x,t)=Ext$$
(5)
The new potentials are :
$$๐^{}=๐tandV^{}=V\frac{f}{t}=0$$
(6)
Of course, the electric field is unchanged but is the underlying physics expressed by the potentials the same ? We believe that an electric field can be created by two very different physical processes that is time variation of a vector potential (like in induction phenomena) or space variation of a scalar potential (like in the electron gun). We are in front of the first paradox : how can a physical quantity (here, the vector potential) be a function of time in a stationary problem ? In the case of a capacitor for example, the static electric field is created by static electric charges on the plates of the capacitor. If one admits that one can describe this static electric field by a time-dependent vector potential, one must admit that the sources (charges or currents) of this electric field are time-dependent which is not experimentally the case.
The unconvinced reader could argue that we proposed a gauge function which depends explicitly on time in a stationary problem. As a matter of fact, if $`f(x,t)`$ does not depend on time, the scalar potential is invariant $`V^{}=V`$ and so is not indetermined with respect to the gauge transformations.
In a stationary problem, the electric field is expressed as $`๐=V`$ whereas the magnetic field is still defined as previously. In this case, the vector potential could still be indetermined. So, are two vector potentials differing from a gradient physically equivalent ?
## III The case of a vector potential equal to a gradient
Now, one will apply the so-called Stokes-Helmholtz-Hodge decomposition to the vector potential :
$$๐=๐_{\mathrm{๐ฅ๐จ๐ง๐ ๐ข๐ญ๐ฎ๐๐ข๐ง๐๐ฅ}}+๐_{\mathrm{๐ญ๐ซ๐๐ง๐ฌ๐ฏ๐๐ซ๐ฌ๐}}$$
(7)
with :
$$๐=g+\times ๐$$
(8)
where $`g`$ is a scalar and $`๐`$ a vector. The decomposition is unique up to the additive gradient of a harmonic function with the following properties Helmholtz :
$$.๐_{\mathrm{๐ก๐๐ซ๐ฆ๐จ๐ง๐ข๐}}=\mathrm{๐}\times ๐_{\mathrm{๐ก๐๐ซ๐ฆ๐จ๐ง๐ข๐}}=\mathrm{๐}$$
(9)
If we use gauge transformations, we can notice that only the longitudinal (and/or harmonic) part of the vector potential and the scalar potential are affected by these transformations. The transverse part remains unchanged. Moreover, the magnetic field depends only on the transverse part. So, if there is indeterminacy, it must imply indeterminacy of the longitudinal (and/or harmonic) part. As a consequence, the longitudinal (and/or harmonic) part cannot have a physical meaning if it is indetermined with respect to the gauge transformations.
Usually, the vector potential is equal to its transverse part in most of the problems of Classical Electromagnetism. For example, the vector potential for a magnet is expressed by :
$$๐=๐_{\mathrm{๐ญ๐ซ๐๐ง๐ฌ๐ฏ๐๐ซ๐ฌ๐}}=\frac{\mu _0}{4\pi }\times (\frac{๐ฆ}{r})$$
(10)
where $`๐ฆ`$ is the strengh of the poles (the so-called magnetic mass or moment).
In this case, we observe a magnetic field by definition. And, if the vector potential varies in time, it creates an electric field again by definition : the time integral of the electric field could be considered as a direct measure of the vector potential without the presence of static charges that is of a scalar potential. If not, one can use the superposition theorem to evaluate first the part of the electric field associated to the static charge (that is the scalar potential) and then the part associated to the current (that is the vector potential) if and only if the separation is possibleโฆ
At this stage, the question is to know whether a vector potential only equal to a gradient can have a physical effect when the electromagnetic field is null.
Outside a solenoid, the vector potential is precisely equal to a gradient as expressed by the following formula in cylindrical coordinates $`(r,\theta )`$ Jackson :
$$๐=\frac{\mathrm{\Phi }}{2\pi r}๐_\theta =\frac{\mathrm{\Phi }\theta }{2\pi }$$
(11)
where $`\mathrm{\Phi }`$ is the flux of magnetic field inside the solenoid or the circulation of the vector potential outside the solenoid. More precisely, its curl and divergence are null so the vector potential outside a solenoid is of a harmonic-type according to the indeterminacy of the Stokes-Helmholtz-Hodge decomposition Helmholtz . Moreover, we point out forcefully that the mathematical indeterminacy due to the gauge transformations is discarded by the boundary conditions which give a physical determination to the vector potential outside a solenoid : the vector potential vanishes far from its current sources.
The external region of a solenoid is one of the numerous experimental configurations to observe the well-known Aharonov-Bohm effect Feynman . So, we are in front of the second paradox : the Aharonov-Bohm effect contradicts the fact that a vector potential equal only to a gradient could not have a physical effect.
However, one usually argues that the Aharanov-Bohm effect is a quantum effect and that the potentials can have a meaning in quantum physics and not in classical physics. Moreover, the prediction of the Aharonov-Bohm effect pointed out an unanticipated way in which the vector potential can affect a measurement on a charged particle in a region of zero field, but the predicted phase is the result of a path integral which insures that the result is gauge independent : it is not the result of a local (gauge-dependent) value of the vector potential.
Why add in the path integral a (mathematical) gradient to the vector potential which is already a (physical) gradient ?
A vector potential equal only to a gradient should not have a physical effect as imply by gauge transformations because we could cancel the longitudinal (and/or harmonic) part by adding the gradient of the appropriate gauge function.
In the solenoid example, we can take as a gauge function $`f(r,\theta ,t)=\mathrm{\Phi }\theta /2\pi `$ and still there is experimentally an effect despite the fact that all the potentials and so the fields cancel outside the solenoid.
However, one can found in the litterature some theoretical arguments against this gauge transformation according to the fact that the existence of the solenoid implies that the the space is not simply-connected. Yet, it is true that in a multiply connected region, the function $`\alpha `$ which characterises the longitudinal (and/or harmonic) part of the vector potential in the general case becomes multivalued but the longitudinal (and/or harmonic) part of the vector potential (the only one which is different from zero outside the solenoid) is not multivalued as one take the gradient of $`\alpha `$.
Another argument is to remember that the gauge transformations were introduced by Lorentz without any constraint on the connectness of the space.
Having the Aharonov-Bohm effect in mind, we can recall now a very simple experiment which cannot be explained with Maxwell equations expressed in function of the electromagnetic field only and which shows the physical character of a longitudinal (and/or harmonic) vector potential in classical physics.
Letโs take again the geometry of the solenoid. If the current varies with time the magnetic field is still null outside the solenoid but because the vector potential is not null outside the solenoid and varies with time, it creates an electric field outside the solenoid. If we denote the flux of the magnetic field inside the solenoid (or the circulation of the vector potential outside the solenoid) $`\mathrm{\Phi }=LI`$ where $`L`$ is the inductance of the solenoid and $`I`$ the current intensity, the electric field is expressed by :
$$๐=\frac{๐}{t}=\frac{L}{2\pi r}\frac{dI}{dt}๐_\theta $$
(12)
If we apply Maxwell equations expressed in function of the fields with the prescription that the magnetic field is null outside the solenoid, we only find that the electric field is lamellar outside the solenoid which is supposed to be infinite ($`\times ๐=0`$ because $`๐/t=\mathrm{๐}`$ even in this time-dependent problem because $`๐=\mathrm{๐}`$ outside the solenoid)โฆ
This experiment is carried out very easily. It demonstrates that a vector potential only equal to a gradient can have a physical effect in Classical Electromagnetism when it varies in time and thus creates by definition an electric field. Of course, for a finite solenoid, the leaking magnetic field is not null outside the solenoid. However, it creates a leaking electric field which is negligeable and opposite with respect to the electric field created by the contribution of the vector potential due to the ideal solenoidโฆ
Another example of a physical vector potential which is equal to a gradient appears in the well-known Meissner effect in supraconductivity and it was discussed nicely by Tonomura Tonomura .
## IV The case of a uniform magnetic field
Another drawback of the gauge transformations can be illustrated by the following example : one often finds in textbooks that we can describe a uniform magnetic field $`๐=B๐_z`$ by either the so-called symmetric โgaugeโ $`๐_๐ฌ=\mathrm{๐}/\mathrm{๐}๐\times ๐ซ`$ or by the so-called Landau โgaugeโ Okun . This two โgaugesโ are related by a gauge transformation :
$$๐_1=\frac{1}{2}๐\times ๐ซ=\frac{1}{2}[By,Bx\mathrm{,0}]$$
(13)
becomes either :
$$๐_2=[0,Bx\mathrm{,0}]or๐_3=[By\mathrm{,0,0}]$$
(14)
with the gauge functions $`\pm f=\pm xy/2`$.
However, there is no discussion in the litterature of the following issue. As a matter of fact, if we consider a solenoid with a current along $`๐_\theta `$, the magnetic field is uniform (along $`๐_๐ณ`$) and could be described by the symmetric โgaugeโ or the Landau โgaugeโ. Yet, the vector potential in the Landau โgaugeโ $`๐_2`$ is along $`๐_๐ฒ`$ whereas the vector potential in the symmetric gauge is along $`๐_\theta `$. We advocate that only the symmetric โgaugeโ is valid in this case because it does respect the symmetry of the currents ($`๐=J๐_\theta `$) whereas the Landau โgaugeโ does not. Moreover, the symmetric โgaugeโ (or the Landau โgaugeโ) is not, in fact, a gauge condition but a solution describing a uniform magnetic field under the Coulomb constraint ($`.๐_\mathrm{๐}=\mathrm{๐}`$). In order to understand this last point, one can picture an analogy between Fluid Mechanics and Classical Electromagnetism. Indeed, the solenoid is analogous to a cylindrical vortex core with vorticity $`๐ฐ`$ and we know that the velocity inside the core is given by $`๐ฎ=\mathrm{๐}/\mathrm{๐}๐ฐ\times ๐ซ`$ which is analogous to the symmetric gauge for an incompressible flow ($`.๐ฎ=\mathrm{๐}`$). Outside the vortex core, the velocity is given by GHP :
$$๐ฎ=\frac{\mathrm{\Gamma }\theta }{2\pi }=\frac{\mathrm{\Gamma }}{2\pi r}๐_\theta $$
(15)
where $`\mathrm{\Gamma }`$ is the flux of vorticity inside the vortex or the circulation of the velocity outside the vortex. One recovers the analogue formula for the vector potential outside a solenoidโฆ
Of course, if the problem we are considering does not feature the cylindrical geometry (two horizontal plates with opposite surface currents for example, analogous to a plane Couette flow GHP ), one of the Landau gauges $`๐_2`$ or $`๐_3`$ must be used instead of the symmetric gauge $`๐_1`$ according to the necessity of respecting the underlying distribution/symmetry of the currents which is at the origin of both the vector potential and the magnetic field. To give a magnetic vector field without specifying its current source is an ill-posed problem which was interpreted so far by attributing an indeterminacy to the vector potential which is wrong.
Now, how can we test experimentally this argument based on symmetry ? If the current of the solenoid varies with time, it will create an electric field which is along $`๐_\theta `$ as the vector potential because the electric field is minus the time derivative of the vector potential. If the currents in the horizontal plates change with time, a horizontal electric field will appear for the same reason. The author rejects all the arguments based on the fact that one can define through a gauge transformations a time dependent scalar potential which would explain the oberved electric field. Indeed, a scalar potential is physically defined with respect to charge distributions and not current distributions.
Contrary to the common belief, it is possible to discriminate experimentally between two vector potentials (related by a gauge transformation) creating a uniform magnetic field. We have shown that the symmetry of the current source implies a certain distribution of the vector potential which is at the origin of an electric field when the intensity is time-dependent. Its orientation is dictated by the vector potential alone which, eventhough is not observable by itself, has observable consequences.
## V Conclusions
What is the meaning of gauge transformations ? We believe that it is only a structural feature (that is linearity) of the definitions of the potentials from the fields. The potentials of Classical Electromagnetism do have a physical meaning as recalled recently Guyon ; Rousseaux ; RL ; Mead ; K and should be considered as the starting point of Classical Electromagnetism RL ; Mead . If we defined the fields from the potentials and not the contrary, the gauge transformations loose their sense because they imply the paradoxes raised in this article. As a conclusion, we must reject gauge transformations. Gauge invariance is preserved but in a weaker sense : the potentials are defined up to a constant. This constant is equal to zero when the sources are confined to a certain region of space : one assumes that the potentials vanish at infinity far from their sources. If the domain is bounded like in a Faraday cage, the surface potentials are given by the contribution of all the sources outside the region of interrest Mourier : one makes the assumption that their knowledge is not important as one measures only differences of potentials according to their definition in function of a constant of reference RL . We recall for example that the vector potential is an electromagnetic impulsion that is a difference of electromagnetic momentum. In mechanics, a momentum is indetermined as it is defined with respect to a reference but the impulsion constructed from this momentum has a definite value so is not indeterminedโฆ
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# SEMILEPTONIC AND ELECTROWEAK PENGUIN RESULTS FROM BABAR
## 1 Introduction
An important goal of the study of semileptonic charmless $`B`$-meson decays is the measurement of $`|V_{ub}|`$, which essentially measures one side of the Unitarity Triangle of the CKM matrix. Both inclusive and exclusive analyses have been used to measure $`|V_{ub}|`$, here we report on recent results using exclusive decays. The measurement of exclusive branching fractions is also useful for distinguishing among theoretical calculations of the form factors, as we shall see.
The the radiative decay $`BX_s\gamma `$ is studied as a probe of New Physics. Since this decay occurs at the one-loop level, the branching fraction is sensitive to models with additional heavy particles that can participate in the loop. In contrast, the shape of the photon energy spectrum is quite insensitive to contributions from New Physics, but it is rather sensitive to two important parameters of Heavy Quark (HQ) theory: the $`b`$-quark mass $`m_b`$ and the quantity $`\mu _\pi ^2`$, which is related to the Fermi motion of the $`b`$-quark inside the hadron.
In this report, we present recent results from the BaBar experiment $`^\mathrm{?}`$ on semileptonic charmless $`B`$-meson decays to exclusive states and on the $`BX_s\gamma `$ process. Beyond providing information on the CKM matrix and probing the possibility of New Physics, these analyses also provide insight into $`B`$-meson decay dynamics and QCD. All results presented herein are preliminary.
## 2 Exclusive semileptonic charmless $`B`$ decays
Semileptonic charmless $`B`$-meson decays to exclusive final states can be used to measure $`|V_{ub}|`$ by exploiting the dependence of the branching fraction on the CKM matrix element. In the case of $`B^0\pi ^{}\mathrm{}^+\nu `$, we have:
$$\frac{d\mathrm{\Gamma }(B^0\pi ^{}\mathrm{}^+\nu )}{dq^2}=\frac{G_F^2}{24\pi ^3}|V_{ub}|^2p_\pi ^3|f_+(q^2)|^2$$
(1)
Here $`G_F`$ is the Fermi coupling constant, $`p_\pi `$ is the pion momentum in the center-of-mass frame, $`q`$ is the invariant mass of the lepton-neutrino pair and $`f_+(q^2)`$ is the form factor, which is calculated theoretically. The goal is to measure the branching fraction in bins of $`q^2`$, which allows one to distinguish among form factor calculations as well as extract the value of $`|V_{ub}|`$.
Experimentally, the branching fractions of exclusive $`bu`$ decays are small and backgrounds from $`bc`$ transitions are substantial. BaBar has used two different methods for overcoming the experimental difficulties: 1) an โuntaggedโ analysis, based on 83 million $`B\overline{B}`$ pairs, where a premium is placed on high quality neutrino reconstruction using the missing momentum in the event; and 2) a โtaggedโ analysis (232 million $`B\overline{B}`$ pairs for the $`B^0\pi ^{}\mathrm{}^+\nu `$ state, 88 million $`B\overline{B}`$ pairs for $`B^+\pi ^0\mathrm{}^+\nu `$ state), where backgrounds are reduced by requiring the other $`B`$-meson in the event be โtaggedโ via a $`D^{()}\mathrm{}\nu `$ decay.
The untagged analysis relies on good neutrino reconstruction to perform its measurement of the branching fractions of $`B\pi \mathrm{}\nu `$ and $`B\rho \mathrm{}\nu `$. The neutrino momentum is inferred from the event missing momentum and strict requirements are placed to ensure good neutrino reconstruction. For example, the event missing mass is required to be compatible with zero: since its resolution broadens linearly with missing energy, we require $`|m_{\mathrm{miss}}^2/2E\mathrm{miss}|<0.4`$ GeV. The variables used to distinguish signal from background are $`m_{\mathrm{ES}}=\sqrt{s/4|\stackrel{}{p}_B^{}|^2}`$ and $`\mathrm{\Delta }E=E_B^{}\sqrt{s}/2`$, where $`\sqrt{s}`$ is the total energy in the $`\mathrm{{\rm Y}}(4S)`$ center-of-mass frame. Figure 1 shows the distribution of these two variables for the $`B\pi \mathrm{}\nu `$ modes in five bins of $`q^2`$.
Branching fraction measurements from both the tagged and untagged analyses are reported in Table 1.
The statistics of the untagged sample permits the study of the $`q^2`$ dependence of the branching fraction and an investigation of several form factor calculations. Figure 1 shows the differential decay rates along with the predictions of four theoretical calculations: LCSR1 $`^\mathrm{?}`$, LQCD1 $`^\mathrm{?}`$, LQCD2 $`^\mathrm{?}`$ and ISGW II $`^\mathrm{?}`$. The $`\chi ^2`$ probabilities are good ($`50`$%) for the first three calculations, while it is marginal (3%) for the ISGW II prediction. We extract the value of $`|V_{ub}|`$ using the $`B\pi \mathrm{}\nu `$ data and the LQCD2 calculation over the full $`q^2`$ range $`025`$ GeV<sup>2</sup>. The BK parametrization $`^\mathrm{?}`$ is used to extrapolate the LQCD2 form factor calculation to low $`q^2`$. We obtain $`|V_{ub}|=(3.82\pm 0.14\pm 0.24\pm 0.11_{0.52}^{+0.88})\times 10^3`$, where the uncertainties are due to statistics, systematics, form factor shape and form factor normalization, respectively.
## 3 $`BX_s\gamma `$
BaBar has performed two analyses of the $`BX_s\gamma `$ channel: a fully inclusive measurement, where no requirements are made on the hadronic state ($`X_s`$) and a semi-inclusive analysis, which aims to reconstruct a large part of the total $`BX_s\gamma `$ rate by summing many exclusively reconstructed modes. The two approaches are complementary: the fully inclusive method requires a lepton tag to reduce continuum background, but nevertheless suffers from significant backgrounds from $`B\overline{B}`$ events. The semi-inclusive analysis, which sums 38 exclusive decay modes, has the advantage of reduced backgrounds due to the kinematic handles provided by fully reconstructed $`B`$ candidates. This analysis however, has a significant systematic uncertainty due to the missing fraction, the part of the $`BX_s\gamma `$ rate that it does not reconstruct. Both of these analysis are based on approximately 89 million $`B\overline{B}`$ pairs.
Figure 3 shows the resulting photon energy spectra for the two analyses. The semi-inclusive analysis has better photon energy resolution for two reasons: 1) the energy is measured in the $`B`$-meson rest frame and 2) the photon energy is actually inferred from the hadronic invariant mass, which has quite good resolution.
We present in Table 2 the energy moments of the photon spectrum calculated above a certain energy threshold, measured in the $`B`$-meson rest frame. A correction is applied to the fully-inclusive values to bring them into this frame. These moments may be directly compared to theoretical calculations to give information on HQ parameters. A fit to the semi-inclusive spectrum was performed to extract the HQ parameters $`m_b`$ and $`\mu _\pi ^2`$. Two theoretical schemes were used to perform the fits: the kinetic scheme $`^\mathrm{?}`$, which gives:
$$m_b=4.69_{0.04}^{+0.05}\mathrm{GeV}\mathrm{and}\mu _\pi ^2=0.30_{0.05}^{+0.07}\mathrm{GeV}^2;$$
(2)
and the shape function scheme $`^\mathrm{?}`$, which yields:
$$m_b=4.65\pm 0.04\mathrm{GeV}\mathrm{and}\mu _\pi ^2=0.19_{0.05}^{+0.06}\mathrm{GeV}^2;$$
(3)
where the errors are the sum of statistical and systematic, but do not include theoretical uncertainties. We note that the parameter $`\mu _\pi ^2`$ is not defined the same way in the two schemes. The spectrum fit also yields the total inclusive branching fraction down to $`E_\gamma >1.6`$ GeV. Averaging the results from the two theoretical schemes gives: $`B(bs\gamma ,E_\gamma >1.6\mathrm{GeV})=(3.38\pm 0.19_{0.410.08}^{+0.64+0.07})\times 10^4`$. We note that the branching fraction result is compatible with the Standard Model calculation $`^\mathrm{?}`$ and with the experimental world average $`^\mathrm{?}`$.
## Acknowledgments
A heartfelt thanks the organizers for a stimulating and enjoyable conference. Thanks also go to my BaBar colleagues for their assistance in preparing these results.
## References
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# Reducible boundary conditions in coupled channels
## 1. Introduction
Quantum-mechanical Hamiltonians in coupled channels are a natural generalization of the Schrรถdinger operators with zero-range interactions . They provide a simple example of matrix-valued Fermi pseudopotentials and can be described using the tools of the extension theory . At the same time, such systems can be viewed as special quantum graphs , so that one can use the general technique in order to describe all possible interactions in channels, their spectral characteristics etc. In the present paper, we develop the formalism of coupled channels using the quantum graph approach.
Such an approach gives a possibility to use the self-adjoint extension theory and the symplectic technique in order to describe all possible boundary conditions. This can be done in many ways including the transfer matrix formalism, which is widely used in scalar one-dimensional point interactions . Using the representations obtained we show that a wide class of the Hamiltonians in question admits decoupling, i.e. by a certain unitary transformation one can reduce them to the direct sum of well-studied Schrรถdinger operators with point interactions; we call such Hamiltonians as well as the corresponding boundary conditions *reducible*. Such conditions can be formulated in various terms, including continuity properties of functions from the domain of the Hamiltonian. An essential feature of the matrix Hamiltonians considered as quantum graphs is the presence of isometric parts (channels). In general, we show that the possibility of the reduction to scalar problems is always connected with certain invariance properties of the boundary conditions with respect to channel permutations. More precisely, it is proved that the matrix Hamiltonian is reducible iff the boundary conditions are invariant under the cyclic coordinate shift in a certain orthonormal basis; similar correspondence was found recently in connection with the inverse scattering problem on graphs . Although such a decoupling is not a generic property, many โstandardโ boundary conditions appear to be reducible, in particular, the so-called $`\delta `$, $`\delta _p`$, $`\delta ^{}`$, and $`\delta _s^{}`$ couplings as well as Kirchhoffโs boundary conditions are reducible (in our opinion, this may illustrate the difference between the general quantum graphs and the coupled channels: the model interactions on graphs appear to be trivial in channels, although one can use the same technique for their study). The reduction permits us to describe the spectrum of the simplest matrix Kronig-Penney Hamiltonians (periodically coupled channels) and to show how their parameters influence various spectral effects like embedded eigenvalues or the number of gaps.
## 2. Parameterization of boundary conditions in coupled channels
Consider a free particle on a graph, with the Hamiltonian acting on each edge as $`\psi _j\psi _j^{\prime \prime }`$, where $`j`$ is the edge index. Assume that the graph has the simplest structure, i.e. consists of $`n`$ half-lines $`[0,\mathrm{})`$ coupled at the origin. The boundary conditions take the form $`A\psi (0)=B\psi ^{}(0)`$, where $`A`$ and $`B`$ are $`n\times n`$ matrices satisfying the following two conditions :
$$AB^{}=BA^{},$$
(1a)
the block matrix $`(AB)`$ has maximal rank. (1b)
The condition (1b) may be rewritten in an equivalent form $`det(AA^{}+BB^{})0`$ or $`det(B\pm iA)0`$ .
Now let us consider the free motion on $`n`$ lines coupled at some point $`q`$. The Hamiltonian of the problem is the operator $`H=d^2/dx^2`$ acting in the space $`L^2(,^n)`$, and the coupling is described by some boundary conditions at $`q`$. Formally one can consider the $`n`$ coupled lines as $`2n`$ coupled half-lines, so that all possible boundary conditions take the form
$$A\left(\begin{array}{c}f(q)\\ f(q+)\end{array}\right)=B\left(\begin{array}{c}f^{}(q)\\ f^{}(q+)\end{array}\right),$$
(2)
where $`A`$ and $`B`$ are $`2n\times 2n`$ matrices satisfying (1). From the other point of view, the nature of coupled channels requests other types of parameterization , namely, the transfer matrix formalism,
$$\left(\begin{array}{c}f^{}(q+)\\ f(q+)\end{array}\right)=\left(\begin{array}{cc}C_{11}& C_{12}\\ C_{21}& C_{22}\end{array}\right)\left(\begin{array}{c}f^{}(q)\\ f(q)\end{array}\right).$$
(3)
Below we will use mainly boundary conditions of the form (2). Nevertheless, in many situations it is useful to know the connection between these two types of parameterization. The following proposition generalizes a construction of .
###### Proposition 1.
The boundary conditions (3) define a self-adjoint operator in $`L^2(,^n)`$ iff the matrices $`C_{jk}`$, $`j,k=1,2`$, obey
$$\begin{array}{c}C_{12}C_{11}^{}C_{11}C_{12}^{}=0,C_{21}C_{22}^{}C_{22}C_{21}^{}=0,\\ C_{11}C_{22}^{}C_{12}C_{21}^{}=E_n.\end{array}$$
(4)
The conditions (4) are equivalent to
$$\begin{array}{c}C_{11}^{}C_{21}C_{21}^{}C_{11}=0,C_{12}^{}C_{22}C_{22}^{}C_{12}=0,\\ C_{11}^{}C_{22}C_{21}^{}C_{12}=E_n\end{array}$$
(5)
###### Proof.
Substituting the equalities
$`\left(\begin{array}{c}f^{}(q+)\\ f(q+)\end{array}\right)`$ $`=\left(\begin{array}{cc}0& 0\\ 0& E_n\end{array}\right)\left(\begin{array}{c}f(q)\\ f(q+)\end{array}\right)+\left(\begin{array}{cc}0& E_n\\ 0& 0\end{array}\right)\left(\begin{array}{c}f^{}(q)\\ f^{}(q+)\end{array}\right),`$
$`\left(\begin{array}{c}f^{}(q)\\ f(q)\end{array}\right)`$ $`=\left(\begin{array}{cc}0& 0\\ E_n& 0\end{array}\right)\left(\begin{array}{c}f(q)\\ f(q+)\end{array}\right)\left(\begin{array}{cc}E_n& 0\\ 0& 0\end{array}\right)\left(\begin{array}{c}f^{}(q)\\ f^{}(q+)\end{array}\right),`$
into (3) we obtain
$$\left(\begin{array}{cc}C_{12}& 0\\ C_{22}& E_n\end{array}\right)\left(\begin{array}{c}f(q)\\ f(q+)\end{array}\right)=\left(\begin{array}{cc}C_{11}& E_n\\ C_{21}& 0\end{array}\right)\left(\begin{array}{c}f^{}(q)\\ f^{}(q+)\end{array}\right).$$
(6)
These boundary conditions define a self-adjoint operator iff the conditions (1) are satisfied. Eq. (1b) holds due to the presence of the blocks with $`E_n`$, and Eq. (1a) takes the form
$$\left(\begin{array}{cc}C_{12}C_{11}^{}& C_{12}C_{21}^{}\\ C_{22}C_{11}^{}1& C_{22}C_{21}^{}\end{array}\right)=\left(\begin{array}{cc}C_{11}C_{12}^{}& C_{11}C_{22}^{}1\\ C_{21}C_{12}^{}& C_{21}C_{22}^{}\end{array}\right),$$
which is exactly (4). These conditions means the equality
$$\left(\begin{array}{cc}C_{11}& C_{12}\\ C_{21}& C_{22}\end{array}\right)\left(\begin{array}{cc}C_{22}^{}& C_{12}^{}\\ C_{21}^{}& C_{11}^{}\end{array}\right)=E_{2n},$$
which is equivalent to
$$\left(\begin{array}{cc}C_{22}^{}& C_{12}^{}\\ C_{21}^{}& C_{11}^{}\end{array}\right)\left(\begin{array}{cc}C_{11}& C_{12}\\ C_{21}& C_{22}\end{array}\right)=E_{2n}$$
and results in (5). โ
If $`n=1`$ (i.e. we have just one channel), the conditions (4) are well known . The blocks $`C_{jk}`$ are just complex numbers, and the conditions $`C_{12}\overline{C_{11}}=C_{11}\overline{C_{12}}`$ and $`C_{21}\overline{C_{22}}=C_{22}\overline{C_{21}}`$ mean that $`\mathrm{arg}C_{11}=\mathrm{arg}C_{12}=:\theta _1`$ and $`\mathrm{arg}C_{21}=\mathrm{arg}C_{22}=:\theta _2`$. Put $`C_{11}=ae^{i\theta _1}`$, $`C_{12}=be^{i\theta _1}`$, $`C_{21}=ce^{i\theta _2}`$, and $`C_{22}=de^{i\theta _2}`$, where $`a,b,c,d`$ and $`\theta _1,\theta _2[0,2\pi )`$. The third condition in (4) reads as $`(adbc)e^{i(\theta _2\theta _1)}=1`$, which means that $`\theta _1=\theta _2=:\theta `$, and the boundary conditions take the form
$$\left(\begin{array}{c}f^{}(q+)\\ f(q+)\end{array}\right)=e^{i\theta }\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)\left(\begin{array}{c}f^{}(q)\\ f(q)\end{array}\right),\theta [0,2\pi ),a,b,c,d,adbc=1.$$
It is reasonable to call boundary conditions admitting the representation (3) *connecting*. Clearly, some boundary conditions are not connecting, for example, the direct sum of the Dirichlet at $`q`$ and $`q+`$. Below we discuss some less obvious examples.
Let us return to the parameterization (2). The use of the values
$$\mathrm{\Gamma }_1f=(f(q),f(q+)),\mathrm{\Gamma }_2f=(f^{}(q),f^{}(q+))$$
(7)
has its origin in the theory of self-adjoint extensions of symmetric operators . We recall briefly some notions from the theory of abstract boundary values . Let $`S`$ be a symmetric operator in a certain Hilbert space with the domain $`domS`$, $`S^{}`$ be its adjoint with the domain $`domS^{}`$. Let $`V`$ be some auxiliary Hilbert space, and $`\mathrm{\Gamma }_1`$, $`\mathrm{\Gamma }_2`$ be linear maps from $`domS^{}`$ to $`V`$ such that
$$f,S^{}gS^{}f,g=\mathrm{\Gamma }_1f,\mathrm{\Gamma }_2g\mathrm{\Gamma }_2f,\mathrm{\Gamma }_1g,\text{for any}f,gdomS^{},$$
(8)
and for any $`(v_1,v_2)V\times V`$ there exists $`fdomS^{}`$ with $`\mathrm{\Gamma }_1f=v_1`$, $`\mathrm{\Gamma }_2f=v_2`$. If $`fdomS^{}`$, the values $`\mathrm{\Gamma }_1f`$ and $`\mathrm{\Gamma }_2f`$ are called *boundary values* of $`f`$, and the triple $`(V,\mathrm{\Gamma }_1,\mathrm{\Gamma }_2)`$ is called a *boundary triple* of the operator $`S`$. Boundary triple exists iff $`S`$ has equal deficiency indices (i.e. has self-adjoint extensions), and in this case the dimension of $`V`$ coincides with this deficiency index, see Chapter 3 in for detailed discussion. If the space $`V`$ is finite-dimensional (i.e. if the deficiency indices of $`S`$ are finite), then all self-adjoint extensions of $`S`$ are restrictions of $`S^{}`$ on the elements $`fdomS^{}`$ satisfying abstract boundary conditions $`A\mathrm{\Gamma }_1f=B\mathrm{\Gamma }_2f`$, where $`A`$ and $`B`$ are matrices satisfying (1). To obtain a one-to-one parameterization of the self-adjoint extensions one can normalize $`A`$ and $`B`$ by choosing unitary matrix $`U`$ with
$$A=1U,B=i(1+U).$$
(9)
Unitary $`2n\times 2n`$ matrices form a $`4n^2`$-dimensional real manifold, which is exactly the number of real parameters in the problem.
Let $`q`$ be fixed. Denote by $`S`$ the operator acting in $`L^2(,^n)`$ as $`d^2/dx^2`$ on functions from $`C_0^{\mathrm{}}(\{q\},^n)`$; this operator is symmetric and has deficiency indices $`(2n,2n)`$. The adjoint operator $`S^{}`$ acts outside $`q`$ in the same way on functions from $`W^{2,2}(\{q\},^n)`$, so that the usual integration by parts in (8) leads to $`V=^{2n}`$ and $`\mathrm{\Gamma }_1`$, $`\mathrm{\Gamma }_2`$ in the form (7), see . The unitary matrix $`U`$ in (9) is particularly useful in approximation problems , and we will actively use the representation (9) for the boundary conditions (2). The choice of a boundary triple is not unique: the dimension of $`V`$ is invariant, so one can always assume $`V=^{2n}`$, and starting with given boundary operators $`\mathrm{\Gamma }_1`$, $`\mathrm{\Gamma }_2`$ one can describe all possible boundary triples by means of suitable linear transformations . From the point of view of spectral problems it may be reasonable to take as a boundary triple for $`S`$ the set $`(V,\stackrel{~}{\mathrm{\Gamma }}_1,\stackrel{~}{\mathrm{\Gamma }}_2)`$ with
$$\stackrel{~}{\mathrm{\Gamma }}_1f=\left(\begin{array}{c}\stackrel{~}{\mathrm{\Gamma }}_{11}f\\ \stackrel{~}{\mathrm{\Gamma }}_{12}f\end{array}\right)=\left(\begin{array}{c}f^{}(q)f^{}(q+)\\ f(q+)f(q)\end{array}\right),\stackrel{~}{\mathrm{\Gamma }}_2f=\left(\begin{array}{c}\stackrel{~}{\mathrm{\Gamma }}_{21}f\\ \stackrel{~}{\mathrm{\Gamma }}_{22}f\end{array}\right)=\left(\begin{array}{c}\frac{f(q)+f(q+)}{2}\\ \frac{f^{}(q)+f^{}(q+)}{2}\end{array}\right),$$
(10)
so that all possible self-adjoint boundary conditions at $`q`$ take the form
$$L\stackrel{~}{\mathrm{\Gamma }}_1f=M\stackrel{~}{\mathrm{\Gamma }}_2f,$$
(11)
where $`L`$, $`M`$ are matrices satisfying the same conditions as $`A`$ and $`B`$ in (1), respectively. Denote the corresponding Hamiltonian by $`H^{L,M}`$. If the boundary conditions (2) and (11) are equivalent, then one can choose the corresponding pairs of matrices $`(A,B)`$ and $`(L,M)`$ in such a way that they satisfy
$$\{\begin{array}{cc}\hfill L& =\frac{1}{2}(AD_1+BD_2),\hfill \\ \hfill M& =BD_1AD_2,\hfill \\ \hfill D_1& =\left(\begin{array}{cc}0& E_n\\ 0& E_n\end{array}\right),\hfill \\ \hfill D_2& =\left(\begin{array}{cc}E_n& 0\\ E_n& 0\end{array}\right),\hfill \end{array}\text{and}\{\begin{array}{cc}\hfill A& =LK_2\frac{1}{2}MK_1,\hfill \\ \hfill B& =\frac{1}{2}MK_2+LK_1,\hfill \\ \hfill K_1& =\left(\begin{array}{cc}E_n& E_n\\ 0& 0\end{array}\right),\hfill \\ \hfill K_1& =\left(\begin{array}{cc}0& 0\\ E_n& E_n\end{array}\right);\hfill \end{array}$$
(12)
We emphasize that due to the non-uniqueness of the parameterization this correspondence is not unique.
For the sake of completeness we describe also the resolvents of Hamiltonians in coupled channels, which is useful in spectral problems. Let $`q_1,\mathrm{},q_m`$, $`m`$, be points of $``$, $`q_1<\mathrm{}<q_m`$. Consider the operator in $`L^2(,^n)`$ acting as $`d^2/dx^2`$ on functions $`fW^{2,2}(\{q_1,\mathrm{},q_m\},^m)`$ satisfying
$$L^{(s)}\stackrel{~}{\mathrm{\Gamma }}_1^sf=M^{(s)}\stackrel{~}{\mathrm{\Gamma }}_2^{(s)}f,$$
$$\stackrel{~}{\mathrm{\Gamma }}_1^{(s)}f=\left(\begin{array}{c}f^{}(q_s)f^{}(q_s+)\\ f(q_s+)f(q_s)\end{array}\right),\stackrel{~}{\mathrm{\Gamma }}_2^{(s)}f=\frac{1}{2}\left(\begin{array}{c}f(q_s)+f(q_s+)\\ f^{}(q_s)+f^{}(q_s+)\end{array}\right),s=1,\mathrm{},m,$$
where for each $`s`$ the $`2n\times 2n`$ matrices $`L^{(s)}`$ and $`M^{(s)}`$ satisfy the same conditions as $`A`$ and $`B`$ before. Denote by $`L`$ and $`M`$ the $`2mn\times 2mn`$ block matrices $`diag(L^{(1)},\mathrm{},L^{(m)})`$ and $`diag(M^{(1)},\mathrm{},M^{(m)})`$, respectively. The Hamiltonian described will be denoted by $`H^{L,M}`$. Note that the operator $`H^0H^{E_n,0}`$ is just the free Laplacian. The following proposition is a variant of the Krein resolvent formula expressed in terms of boundary conditions .
###### Proposition 2.
For $`\zeta _+`$ denote by $`Q(\zeta )`$ the $`2mn\times 2mn`$ matrix consisting of the $`2n\times 2n`$ blocks $`Q^{(l,s)}(\zeta )`$,
$$Q^{(l,s)}(\zeta )=\frac{e^{\sqrt{\zeta |q_lq_s|}}}{2}\left(\begin{array}{cc}\frac{1}{\sqrt{\zeta }}E_n& sign(q_lq_s)E_n\\ sign(q_lq_s)E_n& \sqrt{\zeta }E_n\end{array}\right),l,s=1,\mathrm{},m;$$
here and below we assume that $`sign0=0`$ and that the square root branch is chosen by the condition $`\mathrm{}\sqrt{\zeta }>0`$ for $`\zeta (0,+\mathrm{})`$. If such $`\zeta `$ is a regular value of $`H^{L,M}`$, then the matrix $`MQ(\zeta )L`$ is invertible and for any $`f=(f_1,\mathrm{},f_n)L^2(,^n)`$ the following relation holds:
$$\begin{array}{c}(H^{L,M}\zeta )^1f=(H^0\zeta )^1f\hfill \\ \hfill \underset{s,l=1}{\overset{m}{}}(\underset{j,k=1}{\overset{n}{}}\alpha _{2n(s1)+j,2n(l1)+k}(\zeta )g_{\overline{\zeta }}^{(l)},f_kg_\zeta ^{(s)}e_j\\ \hfill \underset{j,k=1}{\overset{n}{}}\alpha _{2n(s1)+j,2n(l1)+n+k}(\zeta )h_{\overline{\zeta }}^{(l)},f_kg_\zeta ^{(s)}e_j\\ \hfill \underset{j,k=1}{\overset{n}{}}\alpha _{2n(s1)+n+j,2n(l1)+k}(\zeta )g_{\overline{\zeta }}^{(l)},f_kh_\zeta ^{(s)}e_j\\ \hfill \underset{j,k=1}{\overset{n}{}}\alpha _{2n(s1)+n+j,2n(l1)+n+k}(\zeta )h_{\overline{\zeta }}^{(l)},f_kh_\zeta ^{(s)}e_j),\end{array}$$
(13)
where the numbers $`\alpha _{jk}(\zeta )`$, $`j,k=1,\mathrm{},2mn`$, are the entries of $`\left(MQ(\zeta )L\right)^1M`$,
$$g_\zeta ^{(s)}(x)=\frac{1}{2\sqrt{\zeta }}e^{\sqrt{\zeta }|xq_s|},$$
$$h_\zeta ^{(s)}(x)=\frac{sign(xq_s)}{2}e^{\sqrt{\zeta }|xq_s|},s=1,\mathrm{},m,$$
and $`e_j`$, $`j=1,\mathrm{},n`$, is the standard basis of $`^n`$.
Using (12) and (6) one can easily express the resolvent in terms of the parameters $`A`$ and $`B`$ in (2) or $`C_{jk}`$, $`j,k=1,2`$, in (3).
## 3. Decoupling of the single-vertex graph
To emphasize the specifics of the problems with coupled channels let us return to the case of $`n`$ half-lines coupled at the origin. The Hamiltonian of the problem is $`d^2/dx^2`$ acting in $`L^2((0,\mathrm{}),^n)_{j=1}^nL^2(0,+\mathrm{})`$ on functions $`fW^{2,2}((0,\mathrm{}),^n)`$ satisfying the boundary conditions $`Af(0)=Bf^{}(0)`$ with suitable $`A`$ and $`B`$ from (1). We normalize $`A`$ and $`B`$ by choosing them in the form (9) with suitable $`U๐ฐ(n)`$; denote the corresponding Hamiltonian by $`H_U`$. Choose $`\mathrm{\Theta }๐ฐ(n)`$ such that the matrix $`\mathrm{\Theta }^1U\mathrm{\Theta }`$ is diagonal. Denote by the same letter $`\mathrm{\Theta }`$ the associated unitary transformation of $`L^2((0,\mathrm{}),^n)`$, $`(\mathrm{\Theta }f)(x)=\mathrm{\Theta }f(x)`$. For $`x(0,\mathrm{})`$ there holds $`(\mathrm{\Theta }f)^{\prime \prime }(x)=\mathrm{\Theta }f^{\prime \prime }(x)`$. This means that $`\mathrm{\Theta }`$ reduces the boundary conditions to a direct sum,
$$diag(1e^{i\theta _1},\mathrm{},1e^{i\theta _n})g(0)=idiag(1+e^{i\theta _1},\mathrm{},1+e^{i\theta _n})g^{}(0),g=\mathrm{\Theta }f,$$
where $`e^{i\theta _j}`$, $`j=1,\mathrm{},n`$, are the eigenvalues of $`U`$. In other words, the operator $`H_U`$ appears to be unitarily equivalent to the direct sum $`_{j=1}^nH_j`$, where each $`H_j`$ is a self-adjoint operator in $`L^2(0,\mathrm{})`$ acting as $`d^2/dx^2`$ on functions $`g_jW^{2,2}(0,\mathrm{})`$ satisfying the boundary conditions
$$(1e^{i\theta _j})g_j(0)=i(1+e^{i\theta _j})g_j^{}(0),j=1,\mathrm{},n.$$
Therefore, the spectral properties of $`H_U`$ are determined by the eigenvalues of $`U`$ only, i.e. by $`n`$ parameters from $`๐^1`$; moreover, two such $`n`$-tuples differing only by the order of terms are equivalent. This means, in particular, that the inverse scattering problem on the single vertex graph has multiple solution; note that matrix $`U`$ can be still uniquely recovered from the scattering data .
Although the above schema gives a complete result, its applicability is rather restricted if the graph contains more than one vertex. One can find some generalizations for โstar-shapedโ graphs, i.e. if instead of the Hilbert space $`L^2(0,\mathrm{})`$ one deals with the space $`L^2(G)`$, where $`G`$ is some graph; this models $`n`$ identical graphs $`G`$ coupled at a certain point. But even in this case the transformation $`\mathrm{\Theta }`$ mentioned above is non-local and leads in general to non-local boundary conditions at other vertices of the partial graph $`G`$. To obtain a reasonable gain from such a procedure one should consider only diagonalizing transformations preserving the structure of $`G`$. We illustrate such a possibility by problems with coupled channels.
## 4. Decoupling of channels
Denote $`=_{j=1}^nL^2()L^2(,^n)`$. Let $`Q`$ be a uniformly discrete subset of $``$, i.e.
$$\underset{\begin{array}{c}p,qQ,\\ pq\end{array}}{inf}|pq|=d>0.$$
On the domain $`domS=C_0^{\mathrm{}}(Q,^n)`$ consider the operator $`S=d^2/dx^2`$; the adjoint operator $`S^{}`$ acts in the same way on the domain $`domS^{}=W^{2,2}(Q,^n)`$. To obtain self-adjoint operators one should introduce boundary conditions at all points of $`Q`$ as described in Section 2; the uniform discreteness of $`Q`$ guarantees that the operator obtained is self-adjoint . Such an operator can be interpreted as the Hamiltonian of a free particle in $`n`$ channels coupled at the points of $`Q`$. We say that such a Hamiltonian $`H`$ is *reducible* iff there exists $`\mathrm{\Theta }๐ฐ(n)`$ such that the unitary transformation $`f\mathrm{\Theta }f`$ reduces $`H`$ to a direct sum of $`n`$ one-dimensional point interaction Hamiltonians.
All possible boundary conditions at $`qQ`$ have the form
$$\begin{array}{c}\left(1U(q)\right)\mathrm{\Gamma }_1f=i\left(1+U(q)\right)\mathrm{\Gamma }_2f(\mathrm{\Gamma }_1i\mathrm{\Gamma }_2)f=U(q)(\mathrm{\Gamma }_1+i\mathrm{\Gamma }_2)f,\\ \mathrm{\Gamma }_1f=(f(q),f(q+)),\mathrm{\Gamma }_2f=(f^{}(q),f(q+)),U(q)๐ฐ(2n),\end{array}$$
(14)
Denote the Hamiltonian corresponding to these boundary conditions by $`H_{Q,U}`$. Clearly, in the case of finite $`Q`$ all possible Hamiltonians with point interactions are parameterized by $`4n^2|Q|`$ real parameters.
Representing $`U(q)`$ in the block form,
$$U(q)=\left(\begin{array}{cc}U_{11}(q)& U_{12}(q)\\ U_{21}(q)& U_{22}(q)\end{array}\right),$$
we conclude that $`H_{Q,U}`$ is reducible if and only if the $`n\times n`$ blocks $`U_{jk}(q)`$, $`j,k=1,2`$, $`qQ`$, can be diagonalized simultaneously in some orthogonal basis, i.e. if there exists $`\mathrm{\Theta }๐ฐ(n)`$ with
$$\begin{array}{c}\left(\begin{array}{cc}\mathrm{\Theta }^1& 0\\ 0& \mathrm{\Theta }^1\end{array}\right)\left(\begin{array}{cc}U_{11}(q)& U_{12}(q)\\ U_{21}(q)& U_{22}(q)\end{array}\right)\left(\begin{array}{cc}\mathrm{\Theta }& 0\\ 0& \mathrm{\Theta }\end{array}\right)=\left(\begin{array}{cc}\mathrm{\Lambda }_{11}(q)& \mathrm{\Lambda }_{12}(q)\\ \mathrm{\Lambda }_{21}(q)& \mathrm{\Lambda }_{22}(q)\end{array}\right),\\ \mathrm{\Lambda }_{jk}(q)=diag\left(\lambda _{jk}(q,s)\right),\\ \lambda _{jk}(q,s)\text{ are the eigenvalues of }U_{jk}(q),s=1,\mathrm{},n,j,k=1,2,qQ.\end{array}$$
(15)
The unitary transformation $`f\mathrm{\Theta }f`$ reduces $`H_{Q,U}`$ to the direct sum of one-dimensional Hamiltonians, as due to (15) and to the equalities $`(\mathrm{\Theta }f)(q\pm )=\mathrm{\Theta }\left(f(q\pm )\right)`$ and $`(\mathrm{\Theta }f)^{}(q\pm )=\mathrm{\Theta }\left(f^{}(q\pm )\right)`$ the boundary conditions (14) for $`g=\mathrm{\Theta }f`$ take the form
$$\left(\begin{array}{c}g_s(q)+ig_s^{}(q)\\ g_s(q+)ig_s^{}(q+)\end{array}\right)=\left(\begin{array}{cc}\lambda _{11}(q,s)& \lambda _{12}(q,s)\\ \lambda _{21}(q,s)& \lambda _{22}(q,s)\end{array}\right)\left(\begin{array}{c}g_s(q)ig_s^{}(q)\\ g_s(q+)+ig_s^{}(q+)\end{array}\right),s=1,\mathrm{},n.$$
(16)
For the generic interaction, the condition for the reducibility of the boundary conditions can be formulated as follows:
###### Proposition 3.
The boundary conditions (14) are reducible if and only if all the blocks $`U_{jk}(q)`$, $`j,k=1,2`$, $`qQ`$, are normal and commute with each other.
This means that the reducible boundary conditions are parameterized, roughly speaking, by $`n|Q|`$ unitary $`2\times 2`$ matrices $`\left(\lambda _{jk}(q,s)\right)_{j,k=1,2}`$, $`qQ`$, $`s=1,\mathrm{},n`$ (up to permutations), and a unitary $`n\times n`$ matrix $`\mathrm{\Theta }`$ which diagonalizes the boundary conditions. This means that the most general reducible boundary conditions involve $`n|Q|dim_{}๐ฐ(2)+dim_{}๐ฐ(n)=4n|Q|+n^2`$ real parameters.
An analogue of Proposition 3 can be given in terms of the transfer matrix (3).
###### Proposition 4.
The Hamiltonian $`H_{Q,U}`$ given by the boundary conditions (3) is reducible iff the blocks $`C_{jk}`$, $`j,k=1,2`$, are normal and commute for all $`qQ`$. In particular, if $`Q`$ consists of a single point $`q`$ and $`C_{jk}`$ are self-adjoint, then $`H_{Q,U}`$ is reducible.
###### Proof.
The first part is obvious. Assuming that that the blocks $`C_{jk}`$, $`j,k=1,2`$, are self-adjoint (like it was done in ), one concludes from (4) and (5) that they all commute with each other and, therefore, can be diagonalized simultaneously. But this means that the corresponding Hamiltonian is reducible. โ
It is useful also to have โquantitativeโ reducibility criteria in terms of boundary conditions. The corresponding matrix $`U`$ may be difficult to find, but the reducibility can be found by other means. To illustrate this, we consider the Hamiltonian $`\stackrel{~}{H}`$ given by its quadratic form $`\stackrel{~}{Q}(f,f)=f^{},f^{}+f(q),Af(q)`$, $`dom\stackrel{~}{Q}=W^{1,2}(,^n)`$, where $`A`$ is a $`n\times n`$ self-adjoint matrix. This Hamiltonian may be viewed as the so-called matrix $`\delta `$-potential and corresponds to the boundary conditions $`f(q)=f(q+)=:f(q)`$, $`f^{}(q+)f^{}(q)=Af(q)`$, cf. . Clearly, an orthogonal transformation which diagonalizes $`A`$ will reduce the boundary conditions to a direct sum; the Hamiltonian $`\stackrel{~}{H}`$ is unitarily equivalent to the operator $`d^2/dx^2`$ with the boundary conditions $`g_j(q)=g_j(q+)=:g_j(q)`$, $`g_j^{}(q+)g_j^{}(q)=\alpha _jg_j(q)`$, $`j=1,\mathrm{},n`$, respectively, where $`\alpha _j`$ are the eigenvalues of $`A`$. In other words, $`\stackrel{~}{H}`$ is isomorph to the direct sum of the usual one-dimensional $`\delta `$-perturbations. Let us try to generalize this example.
###### Proposition 5.
Let $`Q`$ consist of a single point $`q`$. If there exist $`\alpha ,\beta `$, $`|\alpha |+|\beta |>0`$, and $`c,c^{}\{1,1\}`$ such that
$$\alpha \left(f^{}(q+)+c^{}f^{}(q)\right)=\beta \left(f(q+)+cf(q)\right)$$
(17)
for all $`fdomH_U`$, then $`H_UH_{Q,U}`$ is reducible.
###### Proof.
Consider first the case $`c=c^{}=1`$. Assume first $`\alpha =0`$. Put $`D=\{fdomS^{}:f(q)=f(q+)=:f(q)\}`$. Clearly, $`domH_UD`$, and for arbitrary $`f,gD`$ there holds
$$\begin{array}{c}\mathrm{\Gamma }_1f,\mathrm{\Gamma }_2g\mathrm{\Gamma }_2f,\mathrm{\Gamma }_1gf(q),g^{}(q)+f(q+),g^{}(q+)\hfill \\ \hfill +f^{}(q),g(q)f^{}(q+),g(q+)=\mathrm{\Gamma }_1^{}f,\mathrm{\Gamma }_2^{}g\mathrm{\Gamma }_2^{}f,\mathrm{\Gamma }_1^{}g,\end{array}$$
where $`\mathrm{\Gamma }_1^{}f=f(q)`$, $`\mathrm{\Gamma }_2^{}f=f^{}(q+)f^{}(q)`$. Denote by $`S_0`$ the restriction of $`S^{}`$ to the set $`domS_0=\{fD:\mathrm{\Gamma }_1^{}f=\mathrm{\Gamma }_2^{}f=0\}\{fdomS^{}:f(q)=f(q+)=0,f^{}(q+)=f^{}(q)\}`$. Clearly, this is a symmetric operator, and the set $`D`$ is the domain of its adjoint $`S_0^{}`$. Therefore, $`(^n,\mathrm{\Gamma }_1^{}`$, $`\mathrm{\Gamma }_2^{})`$ is a boundary triple for this new operator $`S_0`$. As $`H_U`$ is a self-adjoint extension of $`S_0`$, there exists $`V๐ฐ(n)`$ so that $`H_U`$ is determined by the boundary conditions $`(\mathrm{\Gamma }_1^{}i\mathrm{\Gamma }_2^{})f=V(\mathrm{\Gamma }_1^{}+i\mathrm{\Gamma }_2^{})f`$, $`fdomS_0^{}`$. Let $`\mathrm{\Theta }`$ be a unitary transformation which diagonalizes $`V`$. Clearly, $`\mathrm{\Theta }`$ induces a unitary transformation of $``$, and the components of the function $`g=\mathrm{\Theta }f`$, $`fdomH_U`$, satisfy
$$g_j(q)=g_j(q+)=:g(q),(1e^{i\theta _j})g_j(q)=i(1+e^{i\theta _j})(g_j^{}(q+)g_j^{}(q)),$$
(18)
$$e^{i\theta _j}\text{ are eigenvalues of }V,j=1,\mathrm{},n.$$
(19)
Therefore, $`H_U`$ is reducible.
Consider now the case $`\alpha 0`$. Put $`\gamma =\beta /\alpha `$. We use the boundary triple (10). Denote by $`D`$ the set $`\{fdomS^{}:\stackrel{~}{\mathrm{\Gamma }}_{11}f\gamma \stackrel{~}{\mathrm{\Gamma }}_{12}f=0\}`$. The condition (17) means the inclusion $`domH_UD`$. Let $`f,gD`$, then $`\stackrel{~}{\mathrm{\Gamma }}_1f,\stackrel{~}{\mathrm{\Gamma }}_2g\stackrel{~}{\mathrm{\Gamma }}_2f,\stackrel{~}{\mathrm{\Gamma }}_1g=\mathrm{\Gamma }_1^{}f,\mathrm{\Gamma }_2^{}g\mathrm{\Gamma }_2^{}f,\mathrm{\Gamma }_1^{}g`$ with $`\mathrm{\Gamma }_1^{}f=\stackrel{~}{\mathrm{\Gamma }}_{12}f`$, $`\mathrm{\Gamma }_2^{}f=\gamma \stackrel{~}{\mathrm{\Gamma }}_{21}f+\stackrel{~}{\mathrm{\Gamma }}_{22}f`$. Denote by $`S_0`$ the restriction of $`S^{}`$ to the set $`domS_0=\{fD:\mathrm{\Gamma }_1^{}f=\mathrm{\Gamma }_2^{}f=0\}`$; this is a symmetric operator, $`D=domS_0^{}`$, and $`(^n,\mathrm{\Gamma }_1^{},\mathrm{\Gamma }_2^{})`$ is a boundary triple for $`S_0`$. As $`H_U`$ is a self-adjoint extension of $`S_0`$, there exists $`V๐ฐ(n)`$ such that $`H_U`$ is determined by the boundary conditions $`(\mathrm{\Gamma }_1^{}i\mathrm{\Gamma }_2^{})f=V(\mathrm{\Gamma }_1^{}+i\mathrm{\Gamma }_2^{})f`$, $`fD`$. Let $`\mathrm{\Theta }๐ฐ(n)`$ such that $`\mathrm{\Theta }^1V\mathrm{\Theta }`$ is diagonal. Noting that $`\mathrm{\Theta }`$ commutes with all the operators $`\stackrel{~}{\mathrm{\Gamma }}_{jk}`$, $`\mathrm{\Gamma }_j^{}`$, $`j,k=1,2`$, we reduce the boundary conditions to a direct sum for $`g=\mathrm{\Theta }f`$.
Now let $`c^{}=1`$, $`c=1`$. Denote by $`D`$ the set of functions $`fdomS^{}`$ satisfying (17) and use again the boundary triple (10), then for any $`f,gD`$ there holds $`\stackrel{~}{\mathrm{\Gamma }}_1f,\stackrel{~}{\mathrm{\Gamma }}_2g\stackrel{~}{\mathrm{\Gamma }}_2f,\stackrel{~}{\mathrm{\Gamma }}_1g=\mathrm{\Gamma }_1^{}f,\mathrm{\Gamma }_2^{}g\mathrm{\Gamma }_2^{}f,\mathrm{\Gamma }_1^{}g`$ with $`\mathrm{\Gamma }_1^{}f=\stackrel{~}{\mathrm{\Gamma }}_{12}f`$, $`\mathrm{\Gamma }_2^{}f=\stackrel{~}{\mathrm{\Gamma }}_{22}f`$. Denote by $`S_0`$ the symmetric operator which is the restriction of $`S^{}`$ to the domain $`domS_0=\{fD:\mathrm{\Gamma }_1^{}f=\mathrm{\Gamma }_2^{}f=0\}`$, then $`D=domS_0^{}`$. Taking into account the fact that $`H_U`$ is a self-adjoint extension of $`S_0`$ we proceed with the proof as in the previous case. The rest combinations of $`c`$ and $`c^{}`$ can be considered in the same way. โ
To formulate an important corollary we recall that a function $`f:`$ is called anticontinuous at $`q`$ if there exist the limits $`f(q\pm )`$ and $`f(q+)+f(q)=0`$.
###### Corollary 6.
Let the set $`Q`$ consist of a single point $`q`$. If one of the following conditions is satisfied:
* all functions from $`domH_{Q,U}`$ are continuous,
* all functions from $`domH_{Q,U}`$ are anticontinuous,
* derivatives of all functions from $`domH_{Q,U}`$ are continuous,
* derivatives of all functions from $`domH_{Q,U}`$ are anticontinuous,
then $`H_{Q,U}`$ is reducible.
We emphasize again that the last proposition and the corollary apply to channels coupled at one point only. Of course, this works also for channels which are identically coupled at several points.
## 5. Permutation-invariant boundary conditions
Let us return to the general Hamiltonian $`H_{Q,U}`$ with boundary conditions at point of a discrete set $`Q`$ (see the beginning of the previous section). The aim of this section is to discuss a correspondence between the reducibility and invariance under channel permutations.
The most general version of this correspondence can be formulated as follows:
###### Proposition 7.
A matrix-valued point interaction Hamiltonian $`H_{Q,U}`$ with a point interaction supported by a uniformly discrete set $`Q`$ is reducible if and only if there exists a unitary $`n\times n`$ matrix $`\mathrm{\Theta }`$ with non-degenerate eigenvalues such that the boundary conditions at all points of $`Q`$ are invariant under the transformation $`f\mathrm{\Theta }f`$.
###### Proof.
At each point $`qQ`$ there exists $`U=U(q)๐ฐ(2n)`$ such that all the functions from the domain of $`H_{Q,U}`$ are characterized by the condition (14).
If the Hamiltonian is reducible, all the blocks $`U_{jk}(q)`$, $`j,k=1,2`$, $`qQ`$, are diagonal in some orthogonal basis and. therefore, commute with any matrix which is diagonal in this basis. Taking an arbitrary diagonal unitary matrix with non-degenerate eigenvalues we show that the condition formulated in the proposition is necessary. Let us show that this condition is also sufficient.
The invariance of boundary conditions under $`\mathrm{\Theta }`$ means that all the blocks $`U_{jk}(q)`$, $`j,k=1,\mathrm{},n`$, commute with $`\mathrm{\Theta }`$. This means that the invariant subspaces of $`\mathrm{\Theta }`$ are such for *all* blocks at *all* points of $`Q`$. As these subspaces are one-dimensional and orthogonal to each other, all the blocks are diagonal in the eigenbasis of $`\mathrm{\Theta }`$. โ
This proposition shows that the reducibility is an effect which is closely connected with non-uniqueness in the inverse scattering or spectral problems on graphs : If the Hamiltonian is invariant under a unitary transformation with certain properties, then there exists another graph (in our case, the union of real lines with marked points) having the same spectrum.
An important example is provided by Hamiltonians which are invariant under certain channel permutations.
###### Corollary 8.
Let $`\sigma `$ be a permutation of order $`n`$ (i.e. $`\sigma ^n=\text{id}`$ and $`\sigma ^k\text{id}`$ for all $`k\{1,\mathrm{},n1\}`$) and $`H_{Q,U}`$ be invariant under the transformation $`f_jf_{\sigma (j)}`$, $`j=1,\mathrm{},n`$, then $`H_{Q,U}`$ is reducible.
###### Proof.
Indeed, the minimal polynom of the transformation is $`\lambda ^n1`$, which means that all eigenvalues are simple. โ
Actually, this situation is in some sense generic, as the following proposition shows.
###### Proposition 9.
The Hamiltonian $`H_{Q,U}`$ is reducible iff there exists an orthonormal basis $`(h_1,\mathrm{},h_n)`$ in $`^n`$, so that all boundary conditions are invariant under the transformation $`h_jh_{(j1)\mathrm{mod}n}`$.
###### Proof.
The minimal polynom of the transformation described is again $`\lambda ^n1`$, which means that all eigenvalues are simple. This shows that the existence of such transformation is sufficient for the boundary conditions to be reducible. Let us show that this condition is also necessary.
Let $`H_{Q,U}`$ be reducible, then there exists an orthonormal basis $`G=(g_j,j=1,\mathrm{},n)`$ in $`^n`$ in which all blocks $`U_{jk}(q)`$ are diagonal. In this basis, define a linear transformation $`\mathrm{\Xi }`$ by its matrix $`diag(\lambda _1,\mathrm{},\lambda _n)`$, $`\lambda _j=\mathrm{exp}(2\pi ij/n)`$, $`j=1,\mathrm{},n`$. Clearly, the blocks $`U_{jk}(q)`$ commute with $`\mathrm{\Xi }`$ (as all these matrices are diagonal). From the other side, in the basis
$$h_j=\frac{1}{\sqrt{n}}\underset{k=1}{\overset{n}{}}\overline{\lambda }_j^kg_k,j=1,\mathrm{},n,$$
the transformation $`\mathrm{\Xi }`$ has the matrix
$$\left(\begin{array}{ccccc}0& 1& 0& \mathrm{}& 0\\ 0& 0& 1& \mathrm{}& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& 0& \mathrm{}& 1\\ 1& 0& 0& \mathrm{}& 0\end{array}\right),$$
which is exactly the matrix of the transformation
$$\underset{j=1}{\overset{n}{}}h_j,fh_j\underset{j=2}{\overset{n}{}}h_j,fh_{j1}+h_1,fh_n.$$
Therefore, the boundary conditions at all points are invariant under the cyclic shift of coordinates with respect to the basis $`(h_1,\mathrm{},h_n)`$. โ
Of certain interest are Hamiltonians (and the corresponding boundary conditions) which are invariant under *all* channel permutations. Clearly, this means that the blocks $`U_{jk}(q)`$ are of the form $`U_{jk}(q)=a_{jk}(q)E_n+b_{jk}(q)J_n`$, where $`J_n`$ is the $`n\times n`$ matrix whose all entries are equal to $`1`$, and the complex numbers $`a_{jk}(q)`$ and $`b_{jk}(q)`$ obey the condition $`\left(a_{jk}(q)\right)_{j,k=1,2},\left(a_{jk}(q)+nb_{jk}(q)\right)_{j,k=1,2}๐ฐ(2)`$. (Clearly, the spectrum of $`J_n`$ consists of a simple eigenvalue $`n`$ and a $`(n1)`$-fold degenerate eigenvalue $`0`$.) This class includes the frequently used $`\delta `$, $`\delta _s^{}`$, $`\delta _p`$, and $`\delta ^{}`$ couplings, which we consider in greater detail (some different notation is used, see ). For more detailed discussion of the origin of these coupling types we refer to the works and references therein. The corresponding boundary conditions for a function $`fW^{2,2}(\{q\},^n)`$ are as follows:
$`\delta (q,\alpha ):`$ $`\{\begin{array}{cc}\hfill f_j(q)=f_k(q+)& =:f(q),j,k=1,\mathrm{},n,\hfill \\ \hfill {\displaystyle \underset{j=1}{\overset{n}{}}}\left(f_j^{}(q+)f_j^{}(q)\right)& =\alpha f(q),\hfill \end{array}`$ (20a)
$`\delta _s^{}(q,\beta ):`$ $`\{\begin{array}{cc}\hfill f_j^{}(q)=f_k^{}(q+)& =:f^{}(q),j,k=1,\mathrm{},n,\hfill \\ \hfill {\displaystyle \underset{j=1}{\overset{n}{}}}\left(f_j(q+)+f_j(q)\right)& =\beta f^{}(q),\hfill \end{array}`$ (20b)
$`\delta _p(q,\alpha ):`$ $`\{\begin{array}{cc}\hfill \pm f_j^{}(q\pm )f_k^{}(q\pm )& ={\displaystyle \frac{\alpha }{2n}}\left(f_j(q\pm )f_k(q\pm )\right),j,k=1,\mathrm{},n,\hfill \\ \hfill {\displaystyle \underset{j=1}{\overset{n}{}}}\left(f_j(q)+f_j(q+)\right)& =0,\hfill \end{array}`$ (20c)
$`\delta ^{}(q,\beta ):`$ $`\{\begin{array}{cc}\hfill f_j(q\pm )f_k(q\pm )& ={\displaystyle \frac{\beta }{2n}}\left(\pm f_j^{}(q\pm )f_k^{}(q\pm )\right),j,k=1,\mathrm{},n,\hfill \\ \hfill {\displaystyle \underset{j=1}{\overset{n}{}}}\left(f_j^{}(q+)f_j^{}(q)\right)& =0,\hfill \end{array}`$ (20d)
where $`\alpha `$ and $`\beta `$ are real parameters. The $`\delta (q,0)`$-coupling corresponds to the so-called *Kirchhoff boundary conditions* at $`q`$; they appear, for example, if one considers the coupled channels as a limit of shrinking manifolds . For the sake of brevity we denote the introducing of boundary conditions as a formal sum, for example, under the operator
$$H=\frac{d^2}{dx^2}+\delta _s^{}(q_1,\beta )+\delta (q_2,\alpha )$$
(21)
we mean the operator which acts as $`ff^{\prime \prime }`$ on functions $`fW^{2,2}(\{q_1,q_2\},^n)`$ satisfying the boundary condition (20b) for $`q=q_1`$ and (20a) for $`q=q_2`$. In one-dimensional case we use a more traditional way of writing, for example,
$$H=\frac{d^2}{dx^2}+\beta \delta _s^{}(xq_1)+\alpha \delta (\alpha q_2)$$
(22)
will denote the same operator as in (21) *assuming that $`n=1`$*. In fact, one can consider the expression (22) as a self-adjoint operator if one uses the theory of distributions with discontinuous test functions , see also .
###### Proposition 10.
Let $`Q`$, $`Q_s^{}`$, $`Q_p`$, $`Q^{}`$ be non-intersecting discrete subsets of $``$, and their union $`P:=QQ_s^{}Q_pQ^{}`$ be uniformly discrete. Denote by $`H`$ the self-adjoint operator in $`L^2(,^n)`$, $`n>1`$, of the form
$$\frac{d^2}{dx^2}+\underset{qQ}{}\delta (q,\alpha _q)+\underset{qQ_s^{}}{}\delta _s^{}(q,\beta _q)+\underset{qQ_p}{}\delta _p(q,\alpha _q)+\underset{qQ^{}}{}\delta ^{}(q,\beta _q),$$
where $`\alpha _q`$, $`\beta _q`$ are real parameters. Then $`H`$ is unitarily equivalent to the direct sum $`_{k=1}^nH_k`$, where $`H_k`$ are self-adjoint operators in $`L^2()`$, namely,
$$H_1=\frac{d^2}{dx^2}+\underset{qQ}{}\frac{\alpha _q}{n}\delta (xq)+\underset{qQ_s^{}}{}\frac{\beta _q}{n}\delta _s^{}(xq)+\underset{qQ_p}{}\frac{\alpha _q}{n}\delta _p(xq)+\underset{qQ^{}}{}\frac{\beta _q}{n}\delta ^{}(xq),$$
(23)
i.e. the operator $`d^2/dx^2`$ acting on functions $`fW^{2,2}(P)`$ satisfying
$`f(q)=f(q+)`$ $`=:f(q),`$ $`f^{}(q+)f^{}(q)`$ $`={\displaystyle \frac{\alpha _q}{n}}f(q),`$ $`qQ,`$
$`f^{}(q)+f^{}(q+)`$ $`=0,`$ $`f(q)+f(q+)`$ $`={\displaystyle \frac{\beta _q}{n}}f(q+),`$ $`qQ_s^{},`$
$`f(q)+f(q+)`$ $`=0,`$ $`f^{}(q+)f^{}(q)`$ $`={\displaystyle \frac{\alpha _q}{n}}f^{}(q+),`$ $`qQ_p,`$
$`f^{}(q)=f^{}(q+)`$ $`=:f^{}(q),`$ $`f(q+)f(q)`$ $`={\displaystyle \frac{\beta _q}{n}}f^{}(q),`$ $`qQ^{}.`$
and the operators $`H_2,\mathrm{},H_n`$ are equal to each other and act as $`g(x)g^{\prime \prime }(x)`$, $`xP`$, on functions $`gW^{2,2}(P)`$ satisfying the following boundary conditions:
$$\begin{array}{c}g(q)=g(q+)=0,qQ,g^{}(q)=g^{}(q+)=0,qQ_s^{},\\ \alpha _qg(q)+2ng^{}(q)=\alpha _qg(q+)2ng^{}(q+)=0,qQ_p,\\ 2ng(q)+\beta _qg^{}(q)=2ng(q+)\beta _qg^{}(q+)=0,qQ^{}.\end{array}$$
(24)
###### Proof.
We recall that the uniform discreteness of $`P`$ guarantees the self-adjointness of $`H`$ .
As it was shown in , the boundary conditions (20) can be written as (14) with $`U(q)=a_qE_{2n}+b_qJ_{2n}`$, where
$$a_q=\{\begin{array}{cc}1\hfill & \text{for }\delta (q,\alpha _q),\hfill \\ 1\hfill & \text{for }\delta _s^{}(q,\beta _q),\hfill \\ \frac{2ni\alpha _q}{2n+i\alpha _q}\hfill & \text{for }\delta _p(q,\alpha _q),\hfill \\ \frac{2n+i\beta _q}{2ni\beta _q}\hfill & \text{for }\delta ^{}(q,\beta _q),\hfill \end{array}b_q=\{\begin{array}{cc}\frac{2}{2n+i\alpha _q}\hfill & \text{for }\delta (q,\alpha _q),\hfill \\ \frac{2}{2ni\beta _q}\hfill & \text{for }\delta _s^{}(q,\beta _q),\hfill \\ \frac{2}{2n+i\alpha _q}\hfill & \text{for }\delta _p(q,\alpha _q),\hfill \\ \frac{2}{2ni\beta _q}J_{2n}\hfill & \text{for }\delta ^{}(q,\beta _q).\hfill \end{array}$$
(25)
The $`n\times n`$ blocks of $`U`$ are of a rather simple form, namely, $`U_{11}(q)=U_{22}(q)=a_qE_n+b_qJ_n`$, $`U_{12}(q)=U_{21}(q)=b_qJ_n`$. Let $`\mathrm{\Xi }`$ be a linear transformation which diagonalizes $`J_n`$, then at each point $`qP`$ the components of the functions $`g:=\mathrm{\Xi }f`$, $`fdomH`$, satisfy
$$\left(\begin{array}{c}g_k(q)+ig_k^{}(q)\\ g_k(q+)ig_k^{}(q+)\end{array}\right)=V_k(q)\left(\begin{array}{c}g_k(q)ig_k^{}(q)\\ g_k(g+)+ig_k^{}(q+)\end{array}\right),k=1,\mathrm{},n$$
(26)
with
$$V_1(q)=\left(\begin{array}{cc}a_q+nb_q& nb_q\\ nb_q& a_q+nb_q\end{array}\right),V_k=\left(\begin{array}{cc}a_q& 0\\ 0& a_q\end{array}\right),k=2,\mathrm{},n,$$
which is exactly (23) and (24). โ
All the boundary conditions (24) are obviously non-connecting; this means that none of the couplings (20) admits the representation (3). This is connected with the fact that these couplings are actually invariant also under half-channel permutation.
## 6. Periodically coupled channels
Let us illustrate the separability effects by periodic problems with point interactions. Periodically coupled channels provide simple examples of periodic quantum graphs, so that the general powerful technique for their analysis is available . The previous discussion gives a possibility to describe the spectrum of some periodic Hamiltonians by other means: one can easily reduce the spectral problem for periodically coupled channels to the spectral problem for periodic scalar Hamiltonians with point interactions, i.e. to the well-studied generalized Kronig-Penney models . We restrict ourselves by considering some examples.
###### Example 11 (Permutation-invariant delta-potential).
In $`L^2(,^2)`$ consider the periodic delta-potential invariant under channel permutation; this corresponds to the boundary conditions
$$f(q)=f(q+)=:f(q),f^{}(q+)f^{}(q)=(\alpha E_2+\beta J_2)f(q),\alpha ,\beta ,q\pi .$$
Elementary considerations show that this Hamiltonian $`H`$ is unitarily equivalent to the direct sum $`H_1H_2`$,
$$H_1=\frac{d^2}{dx^2}+\alpha \underset{n}{}\delta (x\pi n),H_2=\frac{d^2}{dx^2}+(\alpha +2\beta )\underset{n}{}\delta (x\pi n),$$
so that the spectrum of $`H`$ is the union of the spectra of $`H_1`$ and $`H_2`$. If both $`\alpha `$ and $`\alpha +2\beta `$ have the same sign, then the spectrum of $`H`$ has an infinite number of gaps. For example, for $`\alpha ,\alpha +2\beta >0`$ the spectra of $`H_1`$ and $`H_2`$ consist of the bands $`(a_m,m^2)`$ and $`(b_m,m^2)`$, $`m=1,2,\mathrm{}`$, respectively, where $`a_m,b_m>(m1)^2`$, see Theorem III.2.2.3 in . The spectrum of $`H`$ consists then of the bands $`(\mathrm{min}(a_m,b_m),m^2)`$, $`m=1,2,\mathrm{}`$.
Let us show that $`H`$ has only a finite number of gaps if $`\alpha (\alpha +2\beta )<0`$. To be definite, assume that $`\alpha >0`$ and $`\alpha +2\beta <0`$ (the second case can be considered in the same way). The spectrum of $`H_1`$ consists of the bands $`(a_m,m^2)`$, $`m=1,2,\mathrm{}`$, where $`a_m=m^22m2\alpha /\pi 1+O(1/m)`$, $`m\mathrm{}`$, and the spectrum of $`H_2`$ consists of the bands $`(A_m,B_m)`$, $`m=1,2,\mathrm{}`$, where $`A_1<B_1<0`$, $`B_m>A_m=(m1)^2`$, $`m=2,3,\mathrm{}`$, $`B_m=(m1)^2+2m+2(\alpha +2\beta )1+O(1/m)`$, $`m\mathrm{}`$, see Theorem III.2.2.3 in . Obviously, for large $`m`$ there holds $`B_m>a_m`$, which means that the gaps are overlapped by the large bands.
###### Example 12 (Periodic $`\delta `$-coupling).
For any real $`\alpha `$ the operator
$$H=\frac{d^2}{dx^2}+\underset{l}{}\delta (\pi l,\alpha )$$
acting in $`L^2(,^n)`$ is unitarily equivalent to the direct sum $`H_\alpha (_{j=1}^{n1}H_D)`$, where
$$H_\alpha =\frac{d^2}{dx^2}+\frac{\alpha }{n}\underset{l}{}\delta (x\pi l)$$
and $`H_D`$ is the Laplace operator in $`L^2()`$ acting on functions satisfying the Dirichlet boundary conditions at the points $`\pi l`$, $`l`$. Therefore, the spectrum of $`H`$ consists of the spectrum of $`H_\alpha `$ and of the infinitely degenerate eigenvalues $`m^2`$, $`m`$.
For $`\alpha 0`$, the spectrum of $`H_\alpha `$ consists of values $`k^2`$ satisfying the Kronig-Penney equation $`\left|\mathrm{cos}\pi k+\alpha /(2nk)\mathrm{sin}\pi k\right|1`$, $`\mathrm{}k0`$, and the band edges are given by the values $`k^2`$ with $`\mathrm{cos}\pi k+\alpha /(2nk)\mathrm{sin}\pi k=\pm 1`$, see \[1, Theorem III.2.3.1\]. In particular, the Dirichlet eigenvalues are situated on the band edges.
###### Example 13 (Periodic Kirchhoff coupling).
Let us emphasize a particular case of the previous example. If $`\alpha =0`$ (Kirchhoff couplings), then $`H_\alpha `$ is just the free Laplacian. The spectrum of the initial operator $`H`$, i.e. of channels periodically coupled by the Kirchhoff boundary conditions, consists of the semiaxis $`[0,+\mathrm{})`$ and embedded Dirichlet eigenvalues $`m^2`$, $`m=1,2,\mathrm{}`$.
The existence of eigenvalues in the spectrum of a periodic problem on the graph in our toy situation is connected closely with the reducibility of the boundary conditions. Nevertheless, such effects appear in much more general structures . It is known that a periodic graph can have eigenvalues only in the case of compactly supported solutions . The existence of such solutions is possible only in the case of the so-called analytically disjoint couplings, which can produce even stronger spectral effects .
## Acknowledgments
The author thanks Vladimir Geyler, Volodymyr Mikhailets, Olaf Post, and Nader Yeganefar for stimulating discussions and valuable remarks. The paper was considerably improved following the comments of one of the anonymous referees. The author is indebted him very much. The work was partially supported by the Sonderforschungsbereich โRaum $``$ Zeit $``$ Materieโ(SFB 647, Berlin), INTAS, and the program of cooperation between the Deutsche Forschungsgemeinschaft and the Russian Academy of Sciences.
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# 1 Introduction
## 1 Introduction
Supersymmetric field theories in four space-time dimensions are leading candidates to provide the next stage of unification of fundamental interactions beyond the description offered by the Standard Model. They are natural low-energy descendants of higher-dimensional fundamental theories such as superstrings and ten- or eleven-dimensional supergravities. However, supersymmetry must be broken at low energies, and understanding the pattern of masses and couplings which describes this breaking is one of the central issues of the theory of fundamental interactions. Moreover, the forthcoming experiments will be able to test TeV-scale supersymmetry breakdown, hence exploration of theoretical possibilities leading to realistic predictions becomes more and more relevant. One of the obvious questions which arises in this context is how large a hierarchy between the scales parametrizing supersymmetry breakdown is tolerable. This issue has been raised recently in a series of papers on split supersymmetry, , where it has been found that among fermionic superpartners just light gauginos and higgsinos are sufficient to keep the model within experimental limits and to retain a number of interesting predictions.
However, the question about a natural mechanism generating a hierarchy between supersymmetry breaking terms remains open. In general, within the framework of N=1 supergravity in 4d generating a significant hierarchy between supersymmetry breaking terms is problematic. This is more or less expected in the scenarios of gravity mediated supersymmetry breakdown, since the hidden sector breakdown is characterized by a single scale, like condensation scale of strong gauge dynamics, and mediation is modulated only by expectation values of moduli fields, which cannot differ to much as they are determined by the same potential which switches-on the supersymmetry breakdown. Another argument based on the particular structure of N=1 supergravity is the observation, that Majorana gaugini masses are forbidden by R-symmetry, and this symmetry is broken when the gravitino becomes massive. This is because the N=1 gravitino mass term arises by means of the super-higgs effect. In fact, typically the fermions which supply the helicity 1/2 components to the gravitino come from the chiral multiplets (we neglect general D-type breaking as it needs a non-trivial F-component). As a consequence the gravitino mass term depends on the nonzero expectation value of the superpotential, which always breaks R-symmetry.
An interesting proposal which avoids this problem within the framework of higher-dimensional locally supersymmetric theories has been put forward in . There gravitini can obtain Dirac-type mass through mixing with additional degrees of freedom from the gravitational multiplet, which are there due to the N=2 superesymmetry in the bulk. The interesting feature of this mechanism is that at tree level it decouples gravitino mass from the scale of the supersymmetry breaking in the gauge sector, which gives a hope for creating hierarchy between supersymmetry breaking masses. In addition, one can break the R-symmetry continously in the gravity sector, using the brane terms, which is equivalent to adjusting continously boundary conditions. This breaking is communicated in loops to the gauge sector living on the branes.
This mechanism has been analysed in detail at the level of string theory construction, . In this paper we give a detailed description at the level of five-dimensional supergravity. It is interesting to note, that the case with broken supersymmetry but exact R-symmetry corresponds precisely to flipped supergravity of .
To illustrate the difference between four-dimensional and higher-dimensional superhiggs effects, at the end of the paper we describe the superhiggs effect in N=2 supergravity with flipped boundary conditions.
## 2 Rโsymmetry breakdown
The model we discuss here is the simple $`N=2`$ five-dimensional supergravity with branes, basic features of which are summarized in the appendix. In this paper we consider for simplicity flat geometry, hence we do not put in any cosmological term in the bulk, and there are no explicit brane tensions at the fixed points. However, one can freely enhance the model by non-zero gravitini masses localized on the branes, which fully respect fiveโdimensional supersymmetry (for details see ). Thus, the relevant boundary Lagrangian can be written as
$$e_4_{brane}=\underset{i}{}e_4\delta (yy_i)\overline{\mathrm{\Psi }}_\mu ^A\gamma ^{\mu \nu }(M_i+\gamma _5\overline{M}_i)_A^B\mathrm{\Psi }_{\nu B},$$
(1)
where $`(M_i)_{AB}`$ and $`(\overline{M}_i)_{AB}`$ are symmetric matrices which denote gravitini masses on the branes. In addition, one needs to modify the supersymmetry transformation of the fifth component of the gravitino
$$\delta \mathrm{\Psi }_5^A\delta \mathrm{\Psi }_5^A+2\delta _\alpha ^{\mathrm{\hspace{0.33em}5}}\underset{i}{}\delta (yy_i)(M_i+\gamma _5\overline{M}_i)_B^A\gamma _5\eta ^B.$$
(2)
The N=2 five-dimensional supergravity is invariant under the $`SU(2)_R`$ symmetry. Metric tensor and graviphoton form singlets, while gravitini and parameters of supersymmetry transformations form doublets with respect to this symmetry. One can check that the orbifold projections at the given brane, generated by the $`๐_\mathrm{๐}`$ operator and the gravitini masses on the brane, break $`SU(2)_R`$ symmetry to a $`U(1)`$ subgroup. If the projections breaks Rโsymmetry to the same subgroup at each brane, the $`U(1)_R`$ symmetry remains unbroken in the effective theory. In the other case all generators of the $`SU(2)_R`$ symmetry are broken. To be more specific, let us find explicitly the unbroken generator for a given projection.
Let us assume equal $`๐_\mathrm{๐}`$ operators on both branes: $`Q_0=Q_\pi =\sigma _3`$. In addition, let us allow (locally on each brane) only the even components of gravitini to have localized mass terms, i.e.
$`(M_0)_A^B={\displaystyle \frac{1}{2}}\alpha _0(\sigma _1)_A^B,(M_\pi )_A^B={\displaystyle \frac{1}{2}}\alpha _\pi (\sigma _1)_A^B,`$
$`(\overline{M}_0)_A^B={\displaystyle \frac{1}{2}}\mathrm{i}\alpha _0(\sigma _2)_A^B,(\overline{M}_\pi )_A^B={\displaystyle \frac{1}{2}}\mathrm{i}\alpha _\pi (\sigma _2)_A^B,`$ (3)
where $`\alpha _{0,\pi }`$ are real coefficients<sup>1</sup><sup>1</sup>1 In general, such terms can also be generated by the condensition of the superpotentials localized on the branes, $`\alpha _i=W_i`$..
Then the boundary conditions for the gravitini read
$`ฯต^1(y)\delta (y)\gamma _5(\mathrm{\Psi }_{})_\mu ^A=\delta (y)\alpha _0\sigma _1(\mathrm{\Psi }_+)_\mu ^A,`$
$`ฯต^1(y)\delta (y\pi r_c)\gamma _5(\mathrm{\Psi }_{})_\mu ^A=\delta (y\pi r_c)\alpha _\pi \sigma _1(\mathrm{\Psi }_+)_\mu ^A,`$ (4)
where we have decomposed the gravitini into the $`๐_\mathrm{๐}`$-even $`(+)`$ and $`๐_\mathrm{๐}`$-odd $`()`$ components
$$(\mathrm{\Psi }_\pm )_\alpha ^A=\frac{1}{2}(\delta \pm \gamma _5\sigma _3)_B^A\mathrm{\Psi }_\alpha ^B.$$
(5)
One can check that the unique $`U(1)`$ subgroup that leaves the boundary conditions invariant is generated by
$$P_B^A=\left(\frac{2\alpha _0}{\alpha _0^2+1}\sigma _1+\frac{\alpha _0^21}{\alpha _0^2+1}\sigma _3\right)_B^A,$$
(6)
for the projection acting on the $`y=0`$ brane, and by
$$P_B^A=\left(\frac{2\alpha _\pi }{\alpha _\pi ^2+1}\sigma _1+\frac{\alpha _\pi ^21}{\alpha _\pi ^2+1}\sigma _3\right)_B^A,$$
(7)
for the brane at $`y=\pi r_c`$. It is important to note that the brane action localized at each brane is not invariant on its own under the respective unbroken $`U(1)_R`$ symmetry. To see the invariance of the full brane plus bulk action, one needs to include the relevant contributions from the five-dimensional bulk action (see ).
## 3 General solution of the gravitini equation of motion and compactification
To compactify the model to four dimensions one needs to solve 5d equations of motion for gravitini. They take the following form in the bulk
$`\gamma ^{\mu \nu \rho }_\nu \mathrm{\Psi }_\rho ^A\gamma ^5\gamma ^{\mu \nu }_5\mathrm{\Psi }_\nu ^A=0,`$
$`\gamma ^5\gamma ^{\mu \nu }_\mu \mathrm{\Psi }_\nu ^A=0,`$ (8)
where we have chosen the gauge $`\mathrm{\Psi }_5=0`$. The boundary conditions are given by (2). The solution can be expressed as a linear combination of the sine and cosine functions
$`(\mathrm{\Psi }_+)_\mu ^A`$ $`={\displaystyle \underset{n}{}}A^{(n)}\mathrm{cos}(m_n|y|)\left(\begin{array}{c}\psi _{\mu R}^{(n)}\\ \chi _{\mu L}^{(n)}\end{array}\right)^A+{\displaystyle \underset{n}{}}B^{(n)}\mathrm{sin}(m_n|y|)\left(\begin{array}{c}\widehat{\psi }_{\mu R}^{(n)}\\ \widehat{\chi }_{\mu L}^{(n)}\end{array}\right)^A`$ (13)
$`(\mathrm{\Psi }_{})_\mu ^A`$ $`=ฯต(y){\displaystyle \underset{n}{}}A^{(n)}\mathrm{sin}(m_n|y|)\left(\begin{array}{c}\psi _{\mu L}^{(n)}\\ \chi _{\mu R}^{(n)}\end{array}\right)^A+ฯต(y){\displaystyle \underset{n}{}}B^{(n)}\mathrm{cos}(m_n|y|)\left(\begin{array}{c}\widehat{\psi }_{\mu L}^{(n)}\\ \widehat{\chi }_{\mu R}^{(n)}\end{array}\right)^A`$ (18)
where $`\psi _\mu ^{(n)}`$, $`\widehat{\psi }_\mu ^{(n)}`$, $`\chi _\mu ^{(n)}`$, $`\widehat{\chi }_\mu ^{(n)}`$ denote 4d gravitini in the flat space, which satisfy
$`\gamma ^{\mu \rho \nu }_\rho \psi _\nu ^{(n)}m_n\gamma ^{\mu \nu }\psi _\nu ^{(n)}=0`$
$`\gamma ^{\mu \rho \nu }_\rho \chi _\nu ^{(n)}m_n\gamma ^{\mu \nu }\chi _\nu ^{(n)}=0,`$ (19)
with additional conditions $`\gamma ^{\mu \nu }_\mu \psi _\nu ^{(n)}=0`$ and $`\gamma ^{\mu \nu }_\mu \chi _\nu ^{(n)}=0`$. The symplectic Majorana condition implies $`\overline{\psi }_\mu ^{(n)}=(\chi _\mu ^{(n)})^TC`$ (hatted spinors share the same properties).
The boundary condition (2) at the point $`y=0`$ implies $`B^{(n)}\widehat{\psi }_\mu ^{(n)}=\alpha _0A^{(n)}\chi _\mu ^{(n)}`$ and $`B^{(n)}\widehat{\chi }_\mu ^{(n)}=\alpha _0A^{(n)}\psi _\mu ^{(n)}`$, hence one can write
$`(\mathrm{\Psi }_+)_\mu ^A`$ $`={\displaystyle \underset{n}{}}A^{(n)}\left(\mathrm{cos}(m_n|y|)\left(\begin{array}{c}\psi _{\mu R}^{(n)}\\ \chi _{\mu L}^{(n)}\end{array}\right)^A+\alpha _0\mathrm{sin}(m_n|y|)\left(\begin{array}{c}\chi _{\mu R}^{(n)}\\ \psi _{\mu L}^{(n)}\end{array}\right)^A\right)`$ (24)
$`(\mathrm{\Psi }_{})_\mu ^A`$ $`=ฯต(y){\displaystyle \underset{n}{}}A^{(n)}\left(\mathrm{sin}(m_n|y|)\left(\begin{array}{c}\psi _{\mu L}^{(n)}\\ \chi _{\mu R}^{(n)}\end{array}\right)^A+\alpha _0\mathrm{cos}(m_n|y|)\left(\begin{array}{c}\chi _{\mu L}^{(n)}\\ \psi _{\mu R}^{(n)}\end{array}\right)^A\right).`$ (29)
The boundary condition at $`y=\pi r_c`$ implies in turn
$`(1\alpha _0\alpha _\pi )\mathrm{sin}(m_n\pi r_c)\psi _{\mu L}^{(n)}=(\alpha _0+\alpha _\pi )\mathrm{cos}(m_n\pi r_c)\chi _{\mu L}^{(n)}`$
$`(1\alpha _0\alpha _\pi )\mathrm{sin}(m_n\pi r_c)\chi _{\mu R}^{(n)}=(\alpha _0+\alpha _\pi )\mathrm{cos}(m_n\pi r_c)\psi _{\mu R}^{(n)}.`$ (30)
We shall solve these equations considering separately various cases for the gravitini masses.
* Let us start with $`\alpha _0=\alpha _\pi =0`$
The condition (3) gives the following quantization of the masses:
$$\mathrm{sin}(m_n\pi r_c)=0m_n=\frac{\mathrm{n}}{r_c},$$
(31)
where $`\mathrm{n}=0,1,2,\mathrm{}`$ The zero mode does exist and supersymmetry remains unbroken. The solution (24) takes the form
$`(\mathrm{\Psi }_+)_\mu ^A`$ $`={\displaystyle \underset{n}{}}A^{(n)}\mathrm{cos}(m_n|y|)\left(\begin{array}{c}\psi _{\mu R}^{(n)}\\ \chi _{\mu L}^{(n)}\end{array}\right)^A`$ (34)
$`(\mathrm{\Psi }_{})_\mu ^A`$ $`=ฯต(y){\displaystyle \underset{n}{}}A^{(n)}\mathrm{sin}(m_n|y|)\left(\begin{array}{c}\psi _{\mu L}^{(n)}\\ \chi _{\mu R}^{(n)}\end{array}\right)^A`$ (37)
and it is invariant under the symmetry $`U(1)SU(2)_R`$ generated by $`(\sigma _3)_B^A`$. The gravitini $`\psi _\mu ^{(n)}`$ have a negative charge, say $`1`$, with respect to this symmetry while the ones denoted by $`\chi _\mu ^{(n)}`$ have a positive charge $`+1`$. In fact, ona can check that the original $`SU(2)_R`$ symmetry is broken down to this $`U(1)`$ subgroup by the boundary conditions imposed on the brane. We have obtained the Dirac masses in the effective theory, hence the effective four-dimensional action is invariant under the $`U(1)_R`$ symmetry related to the unbroken $`N=1`$ supersymmetry.
* In the second step let us discuss the case $`\alpha _0=\alpha _\pi `$.
Again, the boundary conditions imply
$$\mathrm{sin}(m_n\pi r_c)=0m_n=\frac{\mathrm{n}}{r_c},$$
(38)
and the solution
$`(\mathrm{\Psi }_+)_\mu ^A`$ $`={\displaystyle \underset{n}{}}A^{(n)}\left(\mathrm{cos}(m_n|y|)\left(\begin{array}{c}\psi _{\mu R}^{(n)}\\ \chi _{\mu L}^{(n)}\end{array}\right)^A+\alpha _0\mathrm{sin}(m_n|y|)\left(\begin{array}{c}\chi _{\mu R}^{(n)}\\ \psi _{\mu L}^{(n)}\end{array}\right)^A\right)`$ (43)
$`(\mathrm{\Psi }_{})_\mu ^A`$ $`=ฯต(y){\displaystyle \underset{n}{}}A^{(n)}(\mathrm{sin}(m_n|y|)\left(\begin{array}{c}\psi _{\mu L}^{(n)}\\ \chi _{\mu R}^{(n)}\end{array}\right)^A`$ (49)
$`+\alpha _0\mathrm{cos}(m_n|y|)\left(\begin{array}{c}\chi _{\mu L}^{(n)}\\ \psi _{\mu R}^{(n)}\end{array}\right)^A)`$
preserves $`N=1`$ supersymmetry. Again, the $`U(1)SU(2)_R`$ survives compactification. The unbroken generator is given by (6). Gravitini $`\psi _\mu ^{(n)}`$ are negatively charged, while $`\chi _\mu ^{(n)}`$ have are positively charged with respect to this symmetry, and we have obtained the Dirac mass terms in the effective theory.
* The choice $`\alpha _0=1/\alpha _\pi `$ corresponds to the flipped supergravity.
The boundary conditions imply
$$\mathrm{cos}(m_n\pi r_c)=0m_n=\frac{\mathrm{n}+\frac{1}{2}}{r_c},$$
(50)
and the solution of the equations of motion is
$`(\mathrm{\Psi }_+)_\mu ^A`$ $`={\displaystyle \underset{n}{}}A^{(n)}\left(\mathrm{cos}(m_n|y|)\left(\begin{array}{c}\psi _{\mu R}^{(n)}\\ \chi _{\mu L}^{(n)}\end{array}\right)^A+\alpha _0\mathrm{sin}(m_n|y|)\left(\begin{array}{c}\chi _{\mu R}^{(n)}\\ \psi _{\mu L}^{(n)}\end{array}\right)^A\right)`$ (55)
$`(\mathrm{\Psi }_{})_\mu ^A`$ $`=ฯต(y){\displaystyle \underset{n}{}}A^{(n)}(\mathrm{sin}(m_n|y|)\left(\begin{array}{c}\psi _{\mu L}^{(n)}\\ \chi _{\mu R}^{(n)}\end{array}\right)^A`$ (61)
$`+\alpha _0\mathrm{cos}(m_n|y|)\left(\begin{array}{c}\chi _{\mu L}^{(n)}\\ \psi _{\mu R}^{(n)}\end{array}\right)^A).`$
In this case supersymmetry is broken by the boundary conditions, nevertheless the $`U(1)SU(2)_R`$ symmetry remains unbroken and the unbroken generator is given by (6). Again, the gravitini $`\psi _\mu ^{(n)}`$ have the negative charge, while the $`\chi _\mu ^{(n)}`$ have the positive charge, and we have obtained the Dirac mass terms in the effective theory.
* Finally, we shall treat the remaining cases.
To solve the boundary conditions (3), one needs to change the basis of the four-dimensional gravitini to:
$$\stackrel{~}{\psi }_\mu ^{(n)}=\frac{1}{\sqrt{2}}\left(\psi _\mu ^{(n)}+\chi _\mu ^{(n)}\right),\stackrel{~}{\chi }_\mu ^{(n)}=\frac{1}{\sqrt{2}}\left(\psi _\mu ^{(n)}\chi _\mu ^{(n)}\right).$$
(62)
Then, the equation (3) reads
$`(1\alpha _0\alpha _\pi )\mathrm{sin}(m_n\pi r_c)\stackrel{~}{\psi }_\mu ^{(n)}=(\alpha _0+\alpha _\pi )\mathrm{cos}(m_n\pi r_c)\stackrel{~}{\psi }_\mu ^{(n)}`$
$`(1\alpha _0\alpha _\pi )\mathrm{sin}(m_n\pi r_c)\stackrel{~}{\chi }_\mu ^{(n)}=(\alpha _0+\alpha _\pi )\mathrm{cos}(m_n\pi r_c)\stackrel{~}{\chi }_\mu ^{(n)},`$ (63)
which eventually leads to the following quantization of the KK masses:
$`m_{\stackrel{~}{\psi }}={\displaystyle \frac{1}{r_c}}\left(\mathrm{n}+{\displaystyle \frac{1}{\pi }}\mathrm{arctan}\left({\displaystyle \frac{\alpha _0+\alpha _\pi }{1\alpha _0\alpha _\pi }}\right)\right),\mathrm{for}{\displaystyle \frac{\alpha _0+\alpha _\pi }{1\alpha _0\alpha _\pi }}0,`$
$`m_{\stackrel{~}{\chi }}={\displaystyle \frac{1}{r_c}}\left(\mathrm{n}+1{\displaystyle \frac{1}{\pi }}\mathrm{arctan}\left({\displaystyle \frac{\alpha _0+\alpha _\pi }{1\alpha _0\alpha _\pi }}\right)\right),\mathrm{for}{\displaystyle \frac{\alpha _0+\alpha _\pi }{1\alpha _0\alpha _\pi }}0,`$
$`m_{\stackrel{~}{\psi }}={\displaystyle \frac{1}{r_c}}\left(\mathrm{n}+1+{\displaystyle \frac{1}{\pi }}\mathrm{arctan}\left({\displaystyle \frac{\alpha _0+\alpha _\pi }{1\alpha _0\alpha _\pi }}\right)\right),\mathrm{for}{\displaystyle \frac{\alpha _0+\alpha _\pi }{1\alpha _0\alpha _\pi }}<0,`$
$`m_{\stackrel{~}{\chi }}={\displaystyle \frac{1}{r_c}}\left(\mathrm{n}{\displaystyle \frac{1}{\pi }}\mathrm{arctan}\left({\displaystyle \frac{\alpha _0+\alpha _\pi }{1\alpha _0\alpha _\pi }}\right)\right),\mathrm{for}{\displaystyle \frac{\alpha _0+\alpha _\pi }{1\alpha _0\alpha _\pi }}<0.`$ (64)
The solution takes the form
$`(\mathrm{\Psi }_+)_\mu ^A`$ $`={\displaystyle \underset{n}{}}A_{\stackrel{~}{\psi }}^{(n)}\left(\mathrm{cos}(m_{\stackrel{~}{\psi }}|y|)+\alpha _0\mathrm{sin}(m_{\stackrel{~}{\psi }}|y|)\right)\left(\begin{array}{c}\stackrel{~}{\psi }_{\mu R}^{(n)}\\ \stackrel{~}{\psi }_{\mu L}^{(n)}\end{array}\right)^A`$ (70)
$`+{\displaystyle \underset{n}{}}A_{\stackrel{~}{\chi }}^{(n)}\left(\mathrm{cos}(m_{\stackrel{~}{\chi }}|y|)\alpha _0\mathrm{sin}(m_{\stackrel{~}{\chi }}|y|)\right)\left(\begin{array}{c}\stackrel{~}{\chi }_{\mu R}^{(n)}\\ \stackrel{~}{\chi }_{\mu L}^{(n)}\end{array}\right)^A,`$
$`(\mathrm{\Psi }_{})_\mu ^A`$ $`=ฯต(y){\displaystyle \underset{n}{}}A_{\stackrel{~}{\psi }}^{(n)}\left(\alpha _0\mathrm{cos}(m_{\stackrel{~}{\psi }}|y|)\mathrm{sin}(m_{\stackrel{~}{\psi }}|y|)\right)\left(\begin{array}{c}\stackrel{~}{\psi }_{\mu L}^{(n)}\\ \stackrel{~}{\psi }_{\mu R}^{(n)}\end{array}\right)^A`$ (76)
$`ฯต(y){\displaystyle \underset{n}{}}A_{\stackrel{~}{\chi }}^{(n)}\left(\alpha _0\mathrm{cos}(m_{\stackrel{~}{\chi }}|y|)+\mathrm{sin}(m_{\stackrel{~}{\chi }}|y|)\right)\left(\begin{array}{c}\stackrel{~}{\chi }_{\mu L}^{(n)}\\ \stackrel{~}{\chi }_{\mu R}^{(n)}\end{array}\right)^A.`$
In this case supersymmetry is broken and the orbifold projections break down $`SU(2)_R`$ symmetry to different subgroups at different branes, hence no $`U(1)`$ invariance survives in the effective theory. In particular, non-vanishing Majorana mass terms for gravitini are generated.
## 4 Limitations of four-dimensional description
Let us recall that in the effective theory, at energies below the compactification scale, one observes the zero modes of the particles that form N=1 massless supergravity multiplet and N=1 chiral supermultiplet. The effective N=1 supersymmetric action is determined by a Kรคhler potential $`K`$ and a superpotential $`W`$. Reduction of the five-dimensional bosonic action in the flat case leads to the following form of the Kรคhler function
$$K=3\mathrm{log}\left(T+\overline{T}\right),$$
(77)
where
$$T=r+\mathrm{i}A,$$
(78)
$`r`$ denotes the proper radius of the fifth dimension in the original 5d coordinates and $`A`$ denotes the axion. The only form of the superpotential that leads to the vanishing scalar potential is a constant $`W=\omega `$. In the flat compactification performed in the previous section the effective scalar potential vanishes. As a consequence, the proper radius of the fifth dimension, hence, the vacuum expectation value of the $`T`$ field is undetermined. In the previous section we have denoted the proper radius by $`r_c`$, assuming that there exists some mechanism (in fact unknown) which determines this value. Then we performed rescaling of the fifth coordinate that the expectation value of the $`e_5^{\widehat{5}}`$ is 1. In such a case the curvature scalar in five and four dimensions are equal and one do not need the Weyl rescaling, turning from five-dimensional to four-dimensional theory.
Now, we would like to keep the freedom of choosing the vacuum expectation value of the $`T`$ field. Hence, we assume that the value of $`e_5^{\widehat{5}}`$ is undetermined. Then the proper radius is $`r_ce_5^{\widehat{5}}`$. It is more convenient to put the value of $`r_c`$ equal to 1 (now, it is only a free parameter) and identify the proper radius ($`T`$ field) with $`e_5^{\widehat{5}}`$. To obtain the canonical curvature scalar in four dimensions, the following Weyl rescaling is needed:
$$g_{\mu \nu }r^1g_{\mu \nu }.$$
(79)
Then the mass of the lowest Kaluza-Klein mode of the gravitino changes to<sup>2</sup><sup>2</sup>2In fact, one can argue that $`r^{\frac{3}{2}}`$ gives the effective physical radius of the fifth dimension, see .
$$m_{3/2}=\frac{1}{2r_c}r^{\frac{3}{2}}=\frac{1}{2}r^{\frac{3}{2}}.$$
(80)
In this paper we mostly use, for convenience, mass terms corresponding to the 5d canonical normalization of the gravitational action, however, the need for the final Weyl rescaling is always understood.
In the four-dimensional supergravity gravitino mass is proportional to the vacuum expectation value of the 4d superpotential
$$m_{3/2}=e^{\frac{K}{2}}W,$$
(81)
The calculations made in our effective four-dimensional model lead to
$$m_{3/2}=\omega T+\overline{T}^{\frac{3}{2}}=\omega 2r^{\frac{3}{2}}$$
(82)
and agree with the five-dimensional gravitini mass of the lowest Kaluza-Klein mode (80) for $`\omega =\sqrt{2}`$.
One can calculate the vacuum expectation value of the superpotential which leads to spontaneously broken supersymmetry in the effective four-dimensional supergravity, for a given set of boundary conditions in five-dimensional models:
$$W=\frac{2\sqrt{2}}{\pi }\mathrm{arctan}\left(\frac{\alpha _0+\alpha _\pi }{1\alpha _0\alpha _\pi }\right).$$
(83)
Note that the four-dimensional supergravity presented above describes effective theory at energies below the compactification scale, where one observes the lightest modes of the particles. In fact, this formalism can be valid for the scale of supersymmetry breaking much smaller than the comactification scale. In the other case a gap between the masses of the first and second Kaluza-Klein states of the gravitini is relatively small and it is difficult to find a proper scale below which one observes only the lightest gravitino. In the limiting case (flipped supergravity) these masses are equal, hence above four-dimensional description totally breaks down. We have obtained a novel and unique four-dimensional theory that consists of one massless graviton and two massive gravitini. In addition, the mass terms for the gravitini in the effective Lagrangian are of the Dirac type, hence they are invariant under the $`U(1)_R`$ symmetry.
## 5 Coupling to the matter localized on the branes
In the N=1 four-dimensional supergravity left and right components of the gravitino have opposite charges with respect to the $`U(1)_R`$ symmetry. The complete theory including gauge fields and chiral matter can be arranged to be invariant under the $`R`$-symmetry by the apropriate choice of the superpotential, and the gravitational sector is invariant under this symmetry because gravitino mass terms, which in principle could break it, are absent. However, gravitino couples to the superpotential:
$$_4W(\mathrm{\Phi },\overline{\mathrm{\Phi }})\overline{\psi _\mu }\gamma ^{\mu \nu }\psi _\nu ,$$
(84)
and the nonzero vacuum expectation value of the superpotential spontaneously breaks supersymmetry as well as Rโsymmetry. Effectively, one obtains the Majorana masses for the chiral gravitini $`m_{eff}W`$. In the matter and gauge sectors, supersymmetry breaking manifests itself through masses of scalars and masses of gaugini. The first ones arise at tree level from the explicit coupling to the Fโterms $`_4|F|^2\mathrm{\Phi }^2`$, where
$$F^i=K^{i\overline{j}}D_{\overline{j}}\overline{W}e^{K/2},$$
(85)
and are of the same order as the gravitino mass. The masses of gaugini are generated by loop corrections. To be more specific let us consider the relevant coupling in the fourโdimensional supergravity
$$_4\frac{1}{4}\overline{\psi }_\mu \gamma ^{\nu \rho }\gamma ^\mu \lambda \overline{\psi }_\nu \gamma ^\rho \lambda .$$
(86)
The Fierz rearrangement leads to the following form useful for loop calculations
$`_4{\displaystyle \frac{1}{16}}\overline{\psi }_\mu \gamma ^{\mu \nu }(1+\gamma _5)\psi _\nu \overline{\lambda }(1\gamma _5)\lambda +{\displaystyle \frac{1}{16}}\overline{\psi }_\mu \gamma ^{\mu \nu }(1\gamma _5)\psi _\nu \overline{\lambda }(1+\gamma _5)\lambda `$
$`{\displaystyle \frac{3}{16}}\overline{\psi }_\mu (1+\gamma _5)\psi ^\mu \overline{\lambda }(1\gamma _5)\lambda {\displaystyle \frac{3}{16}}\overline{\psi }_\mu (1\gamma _5)\psi ^\mu \overline{\lambda }(1+\gamma _5)\lambda `$
$`+{\displaystyle \frac{1}{8}}\overline{\psi }_\mu \gamma ^\mu \gamma _5\psi _\nu \overline{\lambda }\gamma ^\nu \gamma _5\lambda +{\displaystyle \frac{1}{16}}\overline{\psi }_\mu \gamma ^\rho \gamma _5\psi ^\mu \overline{\lambda }\gamma _\rho \gamma _5\lambda {\displaystyle \frac{\mathrm{i}}{16}}\overline{\psi }_\mu ฯต^{\mu \rho \nu \sigma }\gamma _\sigma \psi _\nu \overline{\lambda }\gamma _\rho \gamma _5\lambda .`$ (87)
One can check that only terms in the two first lines in (5) can contribute to the effective mass terms for gaugini.
Let us turn to the five dimensional case. In a most general situation in the presence of arbitrary boundary terms one does not know the exact structure of coupling of the five dimensional supergravity to branes<sup>3</sup><sup>3</sup>3See for a discusion of brane-bulk couplings.. However, one can expect that the effective theory should reconstruct the four-dimensional structure described above, with the modification, that the fields that enter (5) are the fermionic modes which have a nonzero amplitude on the brane with the gauge sector in question. Taking the general solution for the gravitini (70), one can check that only one half of the fermionic degrees of freedom couples to the specific brane. For example, the combination that couples to the brane at the point $`y=0`$ is given by
$$\psi _{\mu R}^{0(n)}=\frac{1}{\sqrt{2}}\left(\stackrel{~}{\psi }_{\mu R}^{(n)}+\stackrel{~}{\chi }_{\mu R}^{(n)}\right),\psi _{\mu L}^{0(n)}=\frac{1}{\sqrt{2}}\left(\stackrel{~}{\psi }_{\mu L}^{(n)}\stackrel{~}{\chi }_{\mu L}^{(n)}\right).$$
(88)
The orthogonal combination
$$\psi _{\mu R}^{\pi (n)}=\frac{1}{\sqrt{2}}\left(\stackrel{~}{\psi }_{\mu R}^{(n)}\stackrel{~}{\chi }_{\mu R}^{(n)}\right),\psi _{\mu L}^{\pi (n)}=\frac{1}{\sqrt{2}}\left(\stackrel{~}{\psi }_{\mu L}^{(n)}+\stackrel{~}{\chi }_{\mu L}^{(n)}\right),$$
(89)
decouples from the brane. Notice that the gravitini in the new basis form Majorana spinors such that right and left handed combinations have opposite charges under the $`U(1)_R`$ symmetry preserved by boundary condition given by (6). Of course, spinors in the new basis are not eigenstates of the mass matrix and the mass terms in the Lagrangian take the form
$`_{mass}={\displaystyle \frac{1}{2}}{\displaystyle \underset{n}{}}(\overline{M}\overline{\psi }_\mu ^{0(n)}\gamma ^{\mu \nu }\psi _\nu ^{0(n)}+\overline{M}\overline{\psi }_\mu ^{\pi (n)}\gamma ^{\mu \nu }\psi _\nu ^{\pi (n)}`$
$`\overline{m}_n\overline{\psi }_\mu ^{0(n)}\gamma ^{\mu \nu }\psi _\nu ^{\pi (n)}\overline{m}_n\overline{\psi }_\mu ^{\pi (n)}\gamma ^{\mu \nu }\psi _\nu ^{0(n)}).`$ (90)
The masses
$$\overline{m}_n=\frac{\mathrm{n}+\frac{1}{2}}{r_c},$$
(91)
mix $`\psi _\mu ^{0(n)}`$ and $`\psi _\mu ^{\pi (n)}`$ states and do not violate the $`U(1)_R`$ symmetry, since the left/right handed component of $`\psi _\mu ^{0(n)}`$ has the same charge as the right/left handed component of $`\psi _\mu ^{\pi (n)}`$. The terms which depends on $`\alpha _{0/\pi }`$ form Majorana mass terms that have the same form at each Kaluza-Klein level:
$$\overline{M}=\frac{1}{r_c}\left(\frac{1}{2}\frac{1}{\pi }\mathrm{arctan}\left(\frac{\alpha _0+\alpha _\pi }{1\alpha _0\alpha _\pi }\right)\right),$$
(92)
and break the $`U(1)_R`$ symmetry. One should note that in the flipped limit $`\overline{M}=0`$, which agrees with the fact that the $`U(1)_R`$ symmetry remains unbroken in that case.
Let us consider a vector supermultiplet localized on the brane at $`y=0`$. In the effective four-dimensional theory only one half of the gravitini degrees of freedom couples to this supermultiplet, precisely the same modes which couple to the brane ($`\psi _\mu ^{0(n)}`$). To be able to close the diagrams that produce one-loop effective masses for the gaugini, see Figure 1, one needs a nonzero $`\overline{M}`$.
## 6 Super-higgs effect in the presence of flipped boundary conditions
In this section we shall present in some detail the super-higgs mechanism arising in supergravity spontaneously broken by non-trivial boundary conditions (the Scherk-Schwarz mechanism). We shall explicitly show that the longitudinal degrees of freedom for massive gravitini come from the super-higgs mechanism that occurs at each level of the Kaluza-Klein tower. The fifth component of the five-dimensional gravitini is absorbed by the four-dimensional gravitini. We shall avoid artificial diagonalization of infinitely dimensional matrices known from the earlier work. Our final results agree for instance with those of when they overlap. To start with, let us concentrate on the gravitini equation of motion in the bulk:
$`\gamma ^{\mu \nu \rho }_\nu \mathrm{\Psi }_\rho ^A+\gamma ^5\gamma ^{\mu \nu }_\nu \mathrm{\Psi }_5^A\gamma ^5\gamma ^{\mu \nu }_5\mathrm{\Psi }_\nu ^A=0,`$
$`\gamma ^5\gamma ^{\mu \nu }_\mu \mathrm{\Psi }_\nu ^A=0.`$ (93)
We performe the calculation for the flipped supergravity ($`\alpha _0=1/\alpha _\pi `$), hence the boundary conditions take the form
$`ฯต^1(y)\delta (y)\gamma _5(\mathrm{\Psi }_{})_\mu ^A=\delta (y)\alpha _0\sigma _1(\mathrm{\Psi }_+)_\mu ^A,`$
$`ฯต^1(y)\delta (y\pi r_c)\gamma _5(\mathrm{\Psi }_{})_\mu ^A=\delta (y\pi r_c)(1/\alpha _0)\sigma _1(\mathrm{\Psi }_+)_\mu ^A.`$ (94)
One can easily find solutions:
$`(\mathrm{\Psi }_+)_\mu ^A`$ $`={\displaystyle \underset{n}{}}A^{(n)}\left(\mathrm{cos}(m_n|y|)\left(\begin{array}{c}\psi _{\mu R}^{(n)}\\ \chi _{\mu L}^{(n)}\end{array}\right)^A+\alpha _0\mathrm{sin}(m_n|y|)\left(\begin{array}{c}\chi _{\mu R}^{(n)}\\ \psi _{\mu L}^{(n)}\end{array}\right)^A\right)`$ (99)
$`(\mathrm{\Psi }_{})_\mu ^A`$ $`=ฯต(y){\displaystyle \underset{n}{}}A^{(n)}\left(\mathrm{sin}(m_n|y|)\left(\begin{array}{c}\psi _{\mu L}^{(n)}\\ \chi _{\mu R}^{(n)}\end{array}\right)^A+\alpha _0\mathrm{cos}(m_n|y|)\left(\begin{array}{c}\chi _{\mu L}^{(n)}\\ \psi _{\mu R}^{(n)}\end{array}\right)^A\right)`$ (104)
$`(\mathrm{\Psi }_+)_5^A`$ $`=ฯต(y){\displaystyle \underset{n}{}}A^{(n)}\left(\mathrm{sin}(m_n|y|)\left(\begin{array}{c}\psi _R^{(n)}\\ \chi _L^{(n)}\end{array}\right)^A\alpha _0\mathrm{cos}(m_n|y|)\left(\begin{array}{c}\chi _R^{(n)}\\ \psi _L^{(n)}\end{array}\right)^A\right)`$ (109)
$`(\mathrm{\Psi }_{})_5^A`$ $`={\displaystyle \underset{n}{}}A^{(n)}\left(\mathrm{cos}(m_n|y|)\left(\begin{array}{c}\psi _L^{(n)}\\ \chi _R^{(n)}\end{array}\right)^A+\alpha _0\mathrm{sin}(m_n|y|)\left(\begin{array}{c}\chi _L^{(n)}\\ \psi _R^{(n)}\end{array}\right)^A\right),`$ (114)
where $`\psi _\mu ^{(n)}`$, $`\chi _\mu ^{(n)}`$ and $`\psi ^{(n)}`$, $`\chi ^{(n)}`$ denote 4d gravitini and fermions in the flat space, which satisfy
$`\gamma ^{\mu \rho \nu }_\rho \psi _\nu ^{(n)}\gamma ^{\mu \rho }_\rho \psi ^{(n)}m_n\gamma ^{\mu \nu }\psi _\nu ^{(n)}=0,`$
$`\gamma ^{\mu \rho \nu }_\rho \chi _\nu ^{(n)}\gamma ^{\mu \rho }_\rho \chi ^{(n)}m_n\gamma ^{\mu \nu }\chi _\nu ^{(n)}=0,`$ (115)
with the additional conditions $`\gamma ^{\mu \nu }_\mu \psi _\nu ^{(n)}=\gamma ^{\mu \nu }_\mu \chi _\nu ^{(n)}=0`$. One can easily find the normalization constant: $`A^{(n)}=1/\sqrt{\pi r_c(1+\alpha _0^2)}`$.
The boundary conditions (6) imply the quantization of the masses:
$$m_n=\frac{1}{r_c}\left(\mathrm{n}+\frac{1}{2}\right),\mathrm{for}\mathrm{n}๐.$$
(116)
Let us investigate the effective four-dimensional theory. Putting the solutions (99) into the supergravity action (A.1) leads to the following four-dimensional Lagrangian describing gravitini
$`_{3/2}=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _{\pi r_c}^{\pi r_c}}\overline{\mathrm{\Psi }}_\alpha ^A\gamma ^{\alpha \beta \gamma }_\beta \mathrm{\Psi }_{\gamma A}{\displaystyle \frac{1}{2}}e_5^1e_4\alpha _0\overline{\mathrm{\Psi }}_\mu ^A\gamma ^{\mu \nu }(\sigma _1+\mathrm{i}\gamma _5\sigma _2)_A^B\mathrm{\Psi }_{\nu B}|_{y=0}`$
$`{\displaystyle \frac{1}{2\alpha _0}}e_5^1e_4\overline{\mathrm{\Psi }}_\mu ^A\gamma ^{\mu \nu }(\sigma _1+\mathrm{i}\gamma _5\sigma _2)_A^B\mathrm{\Psi }_{\nu B}|_{y=\pi r_c}`$
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _{\pi r_c}^{\pi r_c}}\left(\overline{\mathrm{\Psi }}_\mu ^A\gamma ^{\mu \nu \rho }_\nu \mathrm{\Psi }_{\rho A}+\overline{\mathrm{\Psi }}_5^A\gamma ^5\gamma ^{\mu \nu }_\mu \mathrm{\Psi }_{\nu A}+\overline{\mathrm{\Psi }}_\mu ^A\gamma ^5\gamma ^{\mu \nu }_\nu \mathrm{\Psi }_{5A}\overline{\mathrm{\Psi }}_\mu ^A\gamma ^5\gamma ^{\mu \nu }_5\mathrm{\Psi }_{\nu A}\right)`$
$`{\displaystyle \frac{1}{2}}e_5^1e_4\alpha _0\overline{\mathrm{\Psi }}_\mu ^A\gamma ^{\mu \nu }(\sigma _1+\mathrm{i}\gamma _5\sigma _2)_A^B\mathrm{\Psi }_{\nu B}|_{y=0}{\displaystyle \frac{1}{2\alpha _0}}e_5^1e_4\overline{\mathrm{\Psi }}_\mu ^A\gamma ^{\mu \nu }(\sigma _1+\mathrm{i}\gamma _5\sigma _2)_A^B\mathrm{\Psi }_{\nu B}|_{y=\pi r_c}`$
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{n}{}}\left(\overline{\psi }_\mu ^{(n)}\gamma ^{\mu \nu \rho }_\nu \psi _\rho ^{(n)}+\overline{\chi }_\mu ^{(n)}\gamma ^{\mu \nu \rho }_\nu \chi _\rho ^{(n)}m_n\overline{\psi }_\mu ^{(n)}\gamma ^{\mu \nu }\psi _\nu ^{(n)}m_n\overline{\chi }_\mu ^{(n)}\gamma ^{\mu \nu }\chi _\nu ^{(n)}\right)`$ (117)
$`{\displaystyle \frac{1}{2}}{\displaystyle \underset{n}{}}\left(\overline{\psi }^{(n)}\gamma ^{\mu \nu }_\mu \psi _\nu ^{(n)}+\overline{\chi }^{(n)}\gamma ^{\mu \nu }_\mu \chi _\nu ^{(n)}\overline{\psi }_\mu ^{(n)}\gamma ^{\mu \nu }_\nu \psi ^{(n)}\overline{\chi }_\mu ^{(n)}\gamma ^{\mu \nu }_\nu \chi ^{(n)}\right).`$
The variational principle leads to the four-dimensional equation of motion (6). One can remove from the Lagrangian $`\psi ^{(n)}`$ and $`\chi ^{(n)}`$ fields by the following redefinition:
$`\psi _\mu ^{(n)}\psi _\mu ^{(n)}{\displaystyle \frac{1}{m_n}}_\mu \psi ^{(n)}`$
$`\chi _\mu ^{(n)}\chi _\mu ^{(n)}{\displaystyle \frac{1}{m_n}}_\mu \chi ^{(n)}.`$ (118)
Also the equation (6) reduces to the standard Rarita-Schwinger equation
$`\gamma ^{\mu \rho \nu }_\rho \psi _\nu ^{(n)}m_n\gamma ^{\mu \nu }\psi _\nu ^{(n)}=0`$
$`\gamma ^{\mu \rho \nu }_\rho \chi _\nu ^{(n)}m_n\gamma ^{\mu \nu }\chi _\nu ^{(n)}=0.`$ (119)
One should note that the transformations (6) are the part of the supersymmetry transformations of the gravitini with the parameters $`\psi ^{(n)}`$ and $`\chi ^{(n)}`$. In the more general case, when one considers the reduction of the full five-dimensional action including the interaction term between gravitini and graviphoton, the redefinitions which remove the fermions $`\psi ^{(n)}`$ and $`\chi ^{(n)}`$ from the four-dimensional Lagrangian should also include terms with graviphoton in the same manner as they appear in the full supersymmetry transformation of the four-dimensional gravitino.
## 7 Summary
The scenarios of split supersymmetry have demonstrated that the current phwenomenological constraints can safely be satisfied in models with a large hierarchy between supersymmetry breaking terms. Using simple locally supersymmetric five-dimensional models we have demonstrated at field theoretical level how the scenario proposed by Antoniadis and Dimopoulos in realizes such a hierarchy, in fact - an arbitrary hierarchy, between gravitini mass terms and masses of gaugini. Indeed, it turns out that for the special choice of boundary conditions realized by a set of brane sources there appears an unbroken R-symmetry (with supersymmetry broken at the same time) which forbids gaugino masses while gravitini masses are non-vanishing. Departure from this symmetric set of boundary conditions breaks R-symmetry, and gaugino masses can be generated at one-loop order, however the magnitude of the resulting soft masses is proportional to the R-symmetry breaking Majorana-type gravitini mass, which is continously deformable to zero (at the R-symmetric point). In contrast to N=1 supergravity all gravitini are massive but R-symmetry can stay unbroken, since in the limiting case with boundary sources of supersymmetry breaking absent, the superhiggs effect is contained within the gravitational sector. While construction of the working field theoretical extension of the Standard Model along the lines discussed here may be a formidable task, the scenario is certainly interesting, as it allows one to avoid constraints imposed by the tight framework of N=1 4d supergravity.
Acknowledgements
This work was partially supported by the EC 6th Framework Programme MRTN-CT-2004-503369, by the Polish State Committee for Scientific Research grant KBN 1 P03D 014 26 and by POLONIUM 2005. R.M. gratefully acknowledges financial support from the European Network for Theoretical Astroparticle Physics (ENTApP), member of ILIAS, EC contract number RII-CT-2004-506222.
Appendix A: Supergravity on $`๐^\mathrm{๐}/๐_\mathrm{๐}`$
Let us define five-dimensional, N=2 supergravity on $`_4\times ๐^\mathrm{๐}/๐_\mathrm{๐}`$, where $`_4`$ denotes four-dimensional Minkowski space-time. Simple supergravity multiplet contains: metric tensor (represented by the vielbein $`e_\alpha ^m`$), two gravitini $`\mathrm{\Psi }_\alpha ^A`$ and vector field $`A_\alpha `$ โ the graviphoton. The pair of gravitini satisfies symplectic Majorana condition $`\overline{\mathrm{\Psi }}^A\mathrm{\Psi }_A^{}\gamma _0=(ฯต^{AB}\mathrm{\Psi }_B)^TC`$. Five-dimensional Lagrangian reads
$`_{grav}=`$ $`{\displaystyle \frac{1}{2}}R{\displaystyle \frac{3}{4}}_{\alpha \beta }^{\alpha \beta }{\displaystyle \frac{1}{2\sqrt{2}}}A_\alpha _{\beta \gamma }_{\delta ฯต}ฯต^{\alpha \beta \gamma \delta ฯต}`$ (A.1)
$`{\displaystyle \frac{1}{2}}\overline{\mathrm{\Psi }}_\alpha ^A\gamma ^{\alpha \beta \gamma }_\beta \mathrm{\Psi }_{\gamma A}`$
$`+{\displaystyle \frac{3\mathrm{i}}{8\sqrt{2}}}\left(\overline{\mathrm{\Psi }}_\gamma ^A\gamma ^{\alpha \beta \gamma \delta }\mathrm{\Psi }_{\delta A}+2\overline{\mathrm{\Psi }}^{\alpha A}\mathrm{\Psi }_A^\beta \right)_{\alpha \beta },`$
with supersymmetry transformations
$`\delta e_\alpha ^m={\displaystyle \frac{1}{2}}\overline{\eta }^A\gamma ^m\mathrm{\Psi }_{\alpha A},\delta A_\alpha ={\displaystyle \frac{\mathrm{i}}{2\sqrt{2}}}\overline{\mathrm{\Psi }}_\alpha ^A\eta _A,`$
$`\delta \mathrm{\Psi }_\alpha ^A=_\alpha \eta ^A{\displaystyle \frac{\mathrm{i}}{4\sqrt{2}}}\left(\gamma _\alpha ^{\beta \gamma }4\delta _\alpha ^\beta \gamma ^\gamma \right)_{\beta \gamma }\eta ^A.`$ (A.2)
One should note at this point that the above Lagrangian is invariant under the $`SU(2)_R`$ symmetry, that acts on the symplectic indices. The graviton and the graviphoton form singlets with respect to this symmetry, while the pair of gravititni and the parameters of the supersymmetry transformations $`\eta ^A`$ transform as doublets.
We pass on to the orbifold $`๐^\mathrm{๐}/๐_\mathrm{๐}`$ by identifying $`(x_\mu ,y)`$ with $`(x_\mu ,y)`$ and choosing the action of $`๐_\mathrm{๐}`$ on the fields. In the bosonic sector we have chosen even parity for $`e_\mu ^a`$, $`e_5^5`$, $`A_5`$ and odd parity for $`e_\mu ^5`$, $`e_5^a`$, $`A_\mu `$. In the fermionic sector $`๐_\mathrm{๐}`$ operators $`Q_0`$ and $`Q_\pi `$ acts on the fields as follows
$`\mathrm{\Psi }_\mu ^A(y)=\gamma _5(Q_0)_B^A\mathrm{\Psi }_\mu ^B(y),\mathrm{\Psi }_\mu ^A(\pi r_cy)=\gamma _5(Q_\pi )_B^A\mathrm{\Psi }_\mu ^B(\pi r_c+y),`$
$`\mathrm{\Psi }_5^A(y)=\gamma _5(Q_0)_B^A\mathrm{\Psi }_5^B(y),\mathrm{\Psi }_5^A(\pi r_cy)=\gamma _5(Q_\pi )_B^A\mathrm{\Psi }_5^B(\pi r_c+y),`$ (A.3)
$`\eta ^A(y)=\gamma _5(Q_0)_B^A\eta ^B(y),\eta ^A(\pi r_cy)=\gamma _5(Q_\pi )_B^A\eta ^B(\pi r_c+y).`$
The symplectic Majorana condition and the normalization $`(Q_{0,\pi })^2=1`$ imply that $`๐_\mathrm{๐}`$ operators can be written as the following linear combinations of the Pauli matrices: $`Q_{0,\pi }=(q_{0,\pi })_i\sigma ^i`$, where $`(q_{0,\pi })_i`$ form real unit vector. In general, one can choose different $`Q_i`$ operators at each orbifold fixed point ($`y=0`$ or $`y=\pi r_c`$).
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# Molecular Hydrogen in the Damped Ly๐ผ Absorber of Q1331+170
## 1 Introduction
Damped Ly$`\alpha `$ absorbers (hereafter DLAs) are defined as quasar absorption systems with neutral hydrogen column density, $`N(\mathrm{HI})>2\times 10^{20}\mathrm{cm}^2`$. They are generally considered to be progenitors of present-day galaxies, in the form of fast rotating disks (Prochaska & Wolfe 1997), or merging protogalactic clumps (Haehnelt et al. 1998). They serve as an important gas reservoir for star formation at high redshifts (e.g. Storrie-Lombardi & Wolfe 2000).
Molecular gas, especially molecular hydrogen, is an important ingredient of star formation, since stars form by cooling via molecular lines. The search for $`\mathrm{H}_2`$ in DLAs can be carried out by observing the $`\mathrm{H}_2`$ absorption lines in the Lyman X$`{}_{}{}^{1}\mathrm{\Sigma }_{g}^{+}`$B$`{}_{}{}^{1}\mathrm{\Sigma }_{u}^{+}`$ and Werner X$`{}_{}{}^{1}\mathrm{\Sigma }_{g}^{+}`$C$`{}_{}{}^{1}\mathrm{\Sigma }_{u}^{\pm }`$ bands, with rest wavelengths from $``$ 1220ร
to the Lyman limit (Abgrall et al. 1993a, 1993b). Ge & Bechtold (1999) conducted a survey for $`\mathrm{H}_2`$ absorption in 13 DLAs, using the MMT Spectrograph. Two have been confirmed with significant $`\mathrm{H}_2`$ fraction; upper limits in the range of $`10^6`$ to $`10^4`$ were put on the $`\mathrm{H}_2`$ fraction in other systems (Ge & Bechtold 1997, 1999, Ge et al. 2001). More recently, based on the high resolution spectra obtained with UVES at the ESO VLT, Ledoux et al. (2003) searched for $`\mathrm{H}_2`$ in a sample of 33 DLAs, with firm detections of $`\mathrm{H}_2`$ absorption in 8 of them. By contrast, almost every line of sight through Galactic or Magellanic gas clouds with comparable H I column density has strong $`\mathrm{H}_2`$ absorption (e.g. Shull et al. 2000).
In the interstellar medium, $`\mathrm{H}_2`$ molecules are expected to form on the surface of dust grains by physical or chemical adsorption (Cazaux & Tielens 2002). Therefore one contributing factor to the lack of $`\mathrm{H}_2`$ in DLAs is the low dust-to-gas ratios in these systems (e.g. Pei et al. 1991, Pettini et al. 1994, Vladilo 1998, Prochaska & Wolfe 2002), perhaps the result of the fact that quasar surveys tend to select objects which are bright and blue (e.g. Gregg et al. 2002, Khare et al. 2004). A second contributing factor to the lack of $`\mathrm{H}_2`$ may be the high gas temperature in DLAs which makes $`\mathrm{H}_2`$ formation onto dust grains inefficient (Petitjean et al. 2000, Liszt 2002). This is consistent with the suggestion that unlike the local ISM, the gas in DLAs might be predominantly warm (Lu et al. 1996, Prochaska & Wolfe 1999). Finally, $`\mathrm{H}_2`$ molecules are easily dissociated by UV photons in the energy range of 11.3 โ 13.6 eV. Therefore the low $`\mathrm{H}_2`$ fraction in DLAs may be the result of a strong UV radiation field (e.g. Black et al. 1987), which could in turn be a signature of high star formation rate (Wolfe et al. 2003a, 2003b, 2004). Recently, Hirashita et al. (2003) showed that the deficiency of $`\mathrm{H}_2`$ in DLAs may also be a result of the expected low volume filling factor of $`\mathrm{H}_2`$ in environments with a strong UV radiation field and low dust-to-gas ratio, rather than a real absence of $`\mathrm{H}_2`$.
Q1331+170 (at $`z_{\mathrm{em}}=2.084`$) is one of the first high redshift quasars discovered (Baldwin et al. 1973), and has a damped Ly$`\alpha `$ absorber at $`z=1.7765`$ (Carswell et al. 1975) with a neutral hydrogen column density of $`N_{\mathrm{HI}}=1.5\times 10^{21}\mathrm{cm}^2`$ (Chaffee et al. 1988, Prochaska & Wolfe 1999). The 21-cm absorption line was also detected (Wolfe & Davis 1979, Chengalur & Kanekar 2000), with a spin temperature of $`T_\mathrm{s}770\mathrm{K}`$. This implies the presence of a cold neutral medium (CNM) phase conducive to the efficient formation of $`\mathrm{H}_2`$ onto dust grains (Petitjean et al. 2000). One of the most interesting features in the absorption spectrum of Q1331+170 is the presence of strong C I and C I\* multiple lines near $`\lambda 1656\mathrm{\AA }`$ (Meyer et al. 1986, Songaila et al. 1994). C I absorption has been found to be rare in quasar absorbers as compared to Galactic diffuse clouds as a result of higher UV radiation fields, or lower dust abundances (Ge et al. 1997, 2001). Since these conditions would also be inhospitable to an appreciable formation of $`\mathrm{H}_2`$, the damped Ly$`\alpha `$ absorber toward Q1331+170 is an excellent candidate for a search for $`\mathrm{H}_2`$ absorption.
Detecting C I is also of cosmological interest, since it is insensitive to the local physical conditions, but depends strongly on the cosmic background level (Liszt 2002). Therefore the fine structure transitions of C I can be used to measure the CMB temperature as a function of redshift (e.g. Bahcall et al. 1973). In the case of the DLA toward Q1331+170, the excitation temperature of neutral carbon determined by Songaila et al. (1994) is $`7.4\pm 0.8\mathrm{K}`$ for one of two velocity components at $`z=1.77654`$, in agreement with the predicted value of $`T_{\mathrm{CMB}}=7.566\mathrm{K}`$ at that redshift. This result implies that local excitation of C I is negligible. Since the population of $`\mathrm{H}_2`$ at different $`J`$ states can also be used to determine physical parameters, such as kinetic temperature, neutral hydrogen number density, and UV radiation field, observations of molecular hydrogen and its rotational excitation provide an independent verification of the observed C I excitation.
In this paper, we present our detection of redshifted $`\mathrm{H}_2`$ in the high resolution absorption spectrum of Q1331+170 taken with the Space Telescope Imaging Spectrograph (STIS) aboard the Hubble Space Telescope (HST). Sec. 2 describes the observations and basic data reductions. In Sec. 3, we describe the Voigt profile fitting of $`\mathrm{H}_2`$ lines, and present the best-fit physical parameters. Discussions are presented in Sec. 4, and conclusions in Sec. 5.
## 2 Observations and continuum fitting
In HST GO programs 7271 and 9172, we used HST/STIS to obtain the near UV spectrum of Q1331+170 in the wavelength range of 2500ร
to 3120ร
. The spectrum was obtained with the E230M grating with $`0.2^{\mathrm{}}\times 0.2^{\mathrm{}}`$ aperture, and the NUV/MAMA detector in the ACCUM operating mode, resulting in spectral resolution of $`R=30,000`$ (or $`10\mathrm{km}\mathrm{s}^1`$) at $`\lambda `$ 2700ร
. The total exposure time was 52,298 seconds, yielding spectra with S/N$``$7 above the Lyman limit of the damped Ly$`\alpha `$ absorber ($``$2530ร
). The data sets are listed in Table 1.
The data were pre-processed with the IRAF calibration routine CALSTIS (v2.7) for nonlinearity correction, dark subtraction, flat-fielding, wavelength calibration, and extraction of the 1-D spectrum. Rebinning was performed on the raw images for Nyquist sampling (i.e. 2 pixels per resolution element). Instead of using automatic wavelength calibration files taken for each Mode Selection Mechanism (MSM) setting, we obtained additional exposures of comparison lamps during each orbit. Typical observing conditions may cause thermal drifts of $`0.1\mathrm{pixel}`$ per hour (see Sec. 11.2 of Space Telescope Imaging Spectrograph Instrument Handbook for Cycle 13 for more details). For our longest single exposure, this corresponds to a wavelength uncertainty of $`0.007\mathrm{\AA }`$ at $`\lambda 2700\mathrm{\AA }`$ introduced by wavelength calibration, or a velocity uncertainty of $`0.8\mathrm{km}\mathrm{s}^1`$.
The entire flux-calibrated spectrum of Q1331+170 is shown in Fig. 1, smoothed by 5 pixels. The continuum was determined by a 4th order polynomial fit using the data longward of the Lyman limit of the DLA ($``$ 2530ร
). Regions with strong absorption features identified by eye were excluded. The dotted line in Fig. 1 shows the $`1\sigma `$ uncertainties of the observed spectrum. The continuum fit suffers from considerable uncertainties, due to the numerous absorption features in the Ly$`\alpha `$ forest, as well as $`\mathrm{H}_2`$ absorption lines in the Lyman and Werner bands. For this reason, we checked our continuum fit according to the following prescription. We separate the whole spectrum into several continuous wavelength regions, and for each region, we construct the histogram of all the data points with normalized flux greater than 1. Ideally, the histogram can be well represented by a (half-)Gaussian distribution with a standard deviation corresponding to the normalized $`1\sigma `$ error of our spectrum. Next, we adjust the best-fit continuum in the region under consideration until the calculated Gaussian standard deviation is approximately equal to the average normalized $`1\sigma `$ error in the same region. The above procedure was performed in an iterative manner, since an adjustment of the continuum also changes the (normalized) $`1\sigma `$ error used for comparison. This procedure can be used to determine both the overall amplitude and shape of the continuum.
We estimate the accuracy of continuum fitting by calculating the reduced $`\chi ^2`$-deviation between the probability distribution function (PDF) of normalized flux adopting the best-fit continuum and that with continuum level scaled by a certain factor (ranging from 50% lower to 50% higher). Here only pixels with normalized flux greater than 1 are considered. This calculation shows that the best-fit continuum is accurate within 15% at 83% confidence level, and accurate within 25% at 98% confidence level, based on the $`\chi ^2`$ distribution.
The two dashed lines in Fig. 1 correspond to the best-fit continuum scaled 15% lower and 15% higher. This represents a very conservative estimate of the uncertainty in the continuum fit; clearly the dashed lines are inconsistent with the expected continuum level in portions of the spectrum. The effects of the possible uncertainties in continuum fitting on our final results are discussed in Sec. 3.3.
## 3 Analysis of molecular hydrogen lines
### 3.1 Detection of molecular hydrogen
Previous spectroscopic studies show that there are at least five redshift systems with detectable metal absorption lines along the line of sight toward Q1331+170. They include the damped Ly$`\alpha `$ system at $`z=1.7765`$, as well as absorbers at $`z=0.7441`$ (Sargent et al. 1988), $`z=1.3284`$ (Steidel & Sargent 1992), $`z=1.4458`$ (Sargent et al. 1988) and $`z=1.7864`$ (Lanzetta et al. 1987).
We searched for molecular hydrogen absorption associated with the $`z=1.7765`$ DLA and two other absorbers at $`z=1.4458`$ and $`z=1.7864`$, since their redshifts allow considerable overlap between the wavelength coverage of our spectrum and the redshifted Lyman and Werner bands of $`\mathrm{H}_2`$. By comparing with the synthetic $`\mathrm{H}_2`$ spectrum, we detected strong molecular hydrogen absorption associated with the $`z=1.7765`$ damped Ly$`\alpha `$ system, whereas the $`\mathrm{H}_2`$ line pattern does not match the observed spectrum for the other two absorbers.
For the damped Ly$`\alpha `$ system at $`z=1.7765`$, we calculated the significance level of $`\mathrm{H}_2`$ lines for each $`J`$ state. We define the significance level, $`\mathrm{SL}`$ associated with state $`J`$ as
$$\mathrm{SL}=\frac{_i[1f_i(\lambda )]๐\lambda }{[_i\sigma _i^2(\lambda )๐\lambda ]^{1/2}},$$
(1)
where $`f_i(\lambda )`$ and $`\sigma _i(\lambda )`$ are the normalized flux and $`1\sigma `$ error within an individual $`\mathrm{H}_2`$ line, and the summation is over all the selected $`\mathrm{H}_2`$ lines for the same $`J`$ state. The results show that based on our spectrum, $`J=0,1,2,3,4`$ lines are detectable at $`4.4\sigma `$, $`5.1\sigma `$, $`9.0\sigma `$, $`8.9\sigma `$ and $`3.7\sigma `$ significance level, respectively. The $`J5`$ lines were not detected at $`3\sigma `$ significance level. The number of unblended $`\mathrm{H}_2`$ lines included in the above calculation is two for $`J=0`$, three for $`J=1`$, six for $`J=2`$, ten for $`J=3`$, two for $`J=4`$, and five for $`J=5`$, respectively (see also Table 2). Therefore, we conclude that we detect molecular hydrogen absorption from the $`J=0`$ to $`J=4`$ states, associated with the damped Ly$`\alpha `$ system at $`z=1.7765`$.
In Fig. 2, we mark the expected positions of metal lines from all the known absorbers, including the Milky Way interstellar medium, on the observed spectrum (solid line) with $`1\sigma `$ errors (dotted line). The atomic line list was taken from Table 2 of Prochaska et al. (2001). Also marked on Fig. 2 are all the $`\mathrm{H}_2`$ lines from $`J=0`$ to $`4`$ in the Lyman and Werner bands for the $`z=1.7765`$ absorber.
### 3.2 Voigt profile fitting of molecular hydrogen lines
We used VPFIT (version 5) to simultaneously fit Voigt profiles for all the $`\mathrm{H}_2`$ lines. All lines with $`J=0,\mathrm{\hspace{0.25em}1},\mathrm{\hspace{0.25em}2},\mathrm{\hspace{0.25em}3}`$ and $`4`$ were fit simultaneously, assuming one $`b`$-parameter and redshift for all lines and different column densities for each $`J`$ state. These free parameters were adjusted independently to minimize the $`\chi ^2`$ value of the overall fitting. We adopted the STIS instrument profile for the echelle E230M and $`0.2^{\mathrm{}}\times 0.2^{\mathrm{}}`$ aperture. Those $`H_2`$ lines with large reduced $`\chi ^2`$ deviations ($`\chi _\nu ^2>1.2`$) were excluded, to avoid possible cases in which $`\mathrm{H}_2`$ absorption may be contaminated by Ly$`\alpha `$ forest lines. This procedure was iterated until no more lines were rejected. Our final model includes 20 different regions with 26 molecular hydrogen lines for simultaneous fitting (two $`J=0`$ lines, three $`J=1`$ lines, seven $`J=2`$ lines, twelve $`J=3`$ lines and two $`J=4`$ lines). We also note that the physical parameters determined from Voigt profile fitting (see below) are insensitive to our adopted maximum reduced $`\chi ^2`$ deviation of $`1.2`$.
Table 2 provides the results for all the $`\mathrm{H}_2`$ lines included in the final fitting. Rest-frame wavelengths, $`f`$-values and damping parameters $`\mathrm{\Gamma }`$ are taken from Abgrall et al. (1993a, 1993b). The sixth column gives the reduced $`\chi ^2`$ value for all the 20 fitted regions. We show the best-fit Voigt profiles for different $`\mathrm{H}_2`$ lines in Fig. 3 (solid lines), superposed on the observed spectrum with $`1\sigma `$ errors. The Voigt profiles for all the $`\mathrm{H}_2`$ lines with $`J=0`$ to $`4`$ in the Lyman and Werner bands are also shown by the dashed line in Fig. 2.
The best-fit results of $`\mathrm{H}_2`$ column densities, $`b`$-parameter and redshift are listed in Table 3. The $`1\sigma `$ errors were derived from $`\chi ^2`$ fitting by VPFIT, and are mainly dependent on photon statistics. The total molecular hydrogen column density is $`N_{\mathrm{H}_2}=(4.45\pm 0.36)\times 10^{19}\mathrm{cm}^2`$. We can place an upper limit on the $`J=5`$ column density as $`N_{J=5}6.3\times 10^{13}\mathrm{cm}^2`$, at $`2\sigma `$ significance level. The $`1\sigma `$ error in redshift given in Table 3 corresponds to the uncertainty in line centroiding, which is equivalent with $`0.9\mathrm{km}\mathrm{s}^1`$ in velocity space. We mentioned in Sec. 2 that the velocity uncertainty caused by wavelength calibration is $`0.8\mathrm{km}\mathrm{s}^1`$. Therefore the total $`1\sigma `$ error is $`1.2\mathrm{km}\mathrm{s}^1`$, which is about 10% of the spectral resolution.
The molecular hydrogen fraction is defined as
$$f_{\mathrm{H}_2}=\frac{2N_{\mathrm{H}_2}}{N_{\mathrm{HI}}+2N_{\mathrm{H}_2}}.$$
(2)
Adopting a neutral hydrogen column density of $`N_{\mathrm{HI}}=(1.50\pm 0.14)\times 10^{21}\mathrm{cm}^2`$ (Prochaska & Wolfe 1999), we derived a molecular hydrogen fraction of $`f_{\mathrm{H}_2}=(5.6\pm 0.7)\%`$. This is the largest value reported so far in any redshifted damped Ly$`\alpha `$ system. The profiles of the Lyman series lines are shown in Fig. 1, assuming $`N_{\mathrm{HI}}=1.50\times 10^{21}\mathrm{cm}^2`$ and $`b=19.6\mathrm{km}\mathrm{s}^1`$. The $`b`$-parameter is directly scaled from the $`b`$-parameter determined from Voigt profile fitting of molecular hydrogen lines, taking into account the difference in particle mass between the two species. The synthetic line profiles are in broad agreement with the observations. The fact that H I lines are heavily blended with other $`\mathrm{H}_2`$ lines makes Voigt profile fitting to derive $`N_{\mathrm{HI}}`$ difficult. Therefore we simply adopt the neutral hydrogen column density from previous work (Prochaska & Wolfe 1999).
### 3.3 Uncertainties due to curve-of-growth and continuum fitting
The uncertainties in the measured $`\mathrm{H}_2`$ column densities at different $`J`$ states depend on the location of the corresponding lines on the curve of growth. We show in Fig. 4 the curve of growth for the 26 $`\mathrm{H}_2`$ lines used in our simultaneous Voigt profile fitting. Note that the equivalent widths are calculated with the best-fit parameters shown in Table 3, since some of the lines are blended and therefore not appropriate for direct measurements of equivalent widths from the observed spectrum. Lines with different $`J`$ values are marked with different symbols. The background dots in Fig. 4 bracket the full range of the curve of growth for $`\mathrm{H}_2`$ lines, taking into account the variations in the atomic data (Abgrall et al. 1993a, 1993b). The figure shows that the $`J=0`$, $`1`$ lines lie on the damped part of the curve of growth, while $`J=2`$, $`3`$ lines on the flat part, and $`J=4`$ lines on the linear part of the curve of growth. Since the determination of column densities is very sensitive to the choice of $`b`$-parameter on the flat part of the curve of growth, the uncertainties in $`N_{J=2}`$ and $`N_{J=3}`$ are relatively larger than the column densities at lower $`J`$ states. Although $`J=4`$ lines are on the linear part of the curve of growth, the $`1\sigma `$ error in $`N_{J=4}`$ is also relatively large due to the low signal-to-noise ratio for these weaker lines.
Uncertainties in column density are also likely introduced by uncertainties in continuum fitting. To investigate this effect, we scaled our continuum by a factor ranging from 25% lower to 25% higher than the best-fit continuum. Then we did simultaneous Voigt profile fitting in the same way as described in Sec. 3.2. The results are shown in Fig. 5, for the measured column densities at different $`J`$ states. We find that the errors in column density at all $`J`$ states caused by $`1\sigma `$ uncertainties in continuum fit are considerably smaller than the $`1\sigma `$ error in column density derived from Voigt profile fitting. In the extreme case that if our best-fit continuum is under-estimated by a factor of 25% (which can be excluded by 98% confidence level, see Sec. 3.2), the corresponding uncertainties in $`\mathrm{H}_2`$ column densities at $`J=2`$ and $`3`$ can be relatively large. Fortunately, the physical state in the $`\mathrm{H}_2`$ absorber is not sensitive to $`\mathrm{N}_{J=2,\mathrm{\hspace{0.25em}3}}`$ (see Sec. 4.1 for more details). Therefore we conclude that our final results are not critically dependent on the uncertainties in continuum fitting.
### 3.4 Excitation temperature
Once the $`\mathrm{H}_2`$ column densities at different $`J`$ states are determined, we can calculate the excitation temperature for each $`J1`$ state by the Boltzmann law
$$\frac{N_J}{N_0}=\frac{g_J}{g_0}\mathrm{exp}[\frac{ฯต_Jฯต_0}{kT_{\mathrm{ex}}}],$$
(3)
where $`ฯต_Jฯต_0`$ is the excitation energy of level $`J`$ relative to the ground state, $`k`$ is the Boltzmann constant, the degeneracy is $`g_J=(2I+1)\times (2J+1)`$, and $`I`$ is the nuclear spin (0 for even $`J`$ and 1 for odd $`J`$). Fig. 6 shows the plot of $`\mathrm{log}\frac{N_J}{g_J}`$ for each $`J`$ state as a function of its excitation potential, with $`1\sigma `$ errors. The populations of the five detected $`J`$ states can be well fitted by a single temperature Boltzmann distribution, with $`\chi _\nu ^2=0.20`$. The one-temperature model is shown by the solid line in Fig. 6, with a best-fit excitation temperature of $`T_{\mathrm{ex}}=(152\pm 10)\mathrm{K}`$. The upper limit put on $`N_{J=5}`$ is also consistent with this model. Increasing or decreasing the best-fit continuum by 15% results in a similar excitation model with the best-fit excitation temperature of $`T_{\mathrm{ex}}=(153\pm 10)\mathrm{K}`$ or $`T_{\mathrm{ex}}=(150\pm 10)\mathrm{K}`$, respectively. The insensitivity of $`T_{\mathrm{ex}}`$ to the continuum fitting partly results from the logarithmic dependence of $`T_{\mathrm{ex}}`$ on the ratio of column densities.
A unique temperature for rotational excitation of $`\mathrm{H}_2`$ is not characteristic of Galactic $`\mathrm{H}_2`$ clouds, and the populations of rotational states of $`\mathrm{H}_2`$ often require at least two different excitation temperatures (e.g. Spitzer & Cochran 1973). $`T_{\mathrm{ex}}`$ for the low $`J`$ states approximates the kinetic temperature of the gas, while $`T_{\mathrm{ex}}`$ for high $`J`$ states results from UV pumping (Spitzer et al. 1974). For the damped Ly$`\alpha `$ absorber of Q1331+170, a unique excitation temperature for all states from $`J=0`$ to $`4`$ indicates that the rotational populations of $`\mathrm{H}_2`$ are probably dominated by collisional excitation in a gas with kinetic temperature around $`150\mathrm{K}`$. Moreover, the local UV radiation field does not significantly populate the high $`J`$ states of $`\mathrm{H}_2`$ molecules. We quantify the UV radiation field below (see Sec. 4.1).
## 4 Discussion
### 4.1 UV radiation field
A simple analysis on the rate equation of the $`J=4`$ state can yield an estimate of the UV photo-absorption rate of $`\mathrm{H}_2`$. Collisions do not provide an important population source for $`J4`$ states (e.g. Browning et al. 2003), unless the gas phase temperature is comparable with or higher than $`T\frac{ฯต_{J4}ฯต_0}{k}1700\mathrm{K}`$. However, such a high temperature can be excluded for the damped Ly$`\alpha `$ absorber toward Q1331+170, either by the distribution of $`\mathrm{H}_2`$ at different $`J`$ states (see Fig. 6), or by the observation of 21 cm absorption line (Wolfe & Davis 1979, Chengalur & Kanekar 2000).
Previous work shows that the $`J=4`$ state is mainly populated by direct formation into this level and pumping from $`J=0`$, while the dominant depopulation mechanism is spontaneous decay (e.g., Dalgarno & Wright 1972, Jura 1975). Assuming a steady state, these effects can be combined into the following equation
$$p_{4,0}R_{\mathrm{abs}}n(\mathrm{H}_2,J=0)+0.19R_{\mathrm{dust}}n(\mathrm{HI})n(\mathrm{H})=A_{42}n(\mathrm{H}_2,J=4),$$
(4)
where $`R_{\mathrm{abs}}`$ is the photo-absorption rate, $`n(\mathrm{HI})`$ is the neutral hydrogen number density, $`n(\mathrm{H})n(\mathrm{HI})+2n(\mathrm{H}_2)`$ is the total hydrogen number density, $`R_{\mathrm{dust}}`$ is the $`\mathrm{H}_2`$ formation rate on the surface of dust grains, $`p_{4,0}=0.26`$ is the UV pumping efficiency from $`J=0`$ to $`J=4`$ (Jura 1975), and $`A_{42}=2.75\times 10^9\mathrm{s}^1`$ is the spontaneous transition probability from $`J=4`$ to $`J=2`$ (Wolniewicz et al. 1998). For simplicity, we assume 19% of all $`\mathrm{H}_2`$ formation results in the population at $`J=4`$ (Jura 1975), although more recent work suggests that this fraction may vary for different types of dust grains (Takahashi & Uehara 2001).
The equilibrium between the formation and photo-dissociation of $`\mathrm{H}_2`$ can be written as
$$R_{\mathrm{dust}}n(\mathrm{HI})n(\mathrm{H})0.11R_{\mathrm{abs}}n(\mathrm{H}_2),$$
(5)
where we assume that approximately 11% of photo-absorption leads to photo-dissociation (Jura 1974). Combining Eqn. 4 and 5, we get $`R_{\mathrm{abs}}=(7.6\pm 2.4)\times 10^{13}\mathrm{s}^1`$. We note that the photo-absorption rate derived this way is independent of the neutral hydrogen density, since the relevant term $`Rn(\mathrm{HI})n(\mathrm{H})`$ cancels from Eqn. 4 and 5. Therefore $`R_{\mathrm{abs}}`$ is independent of the adopted H I column density (see Sec. 3.5).
The photo-absorption rate, $`R_{\mathrm{abs}}`$ derived above is based on the $`\mathrm{H}_2`$ population at the $`J=4`$ state. Since $`J=4`$ lines are optically thin (see Sec. 3.3), the photo-absorption rate within the $`\mathrm{H}_2`$ cloud should be approximately the same as that outside the cloud. Therefore $`R_{\mathrm{abs}}`$ is also characteristic of the UV background intensity in the absorberโs environment. With this in mind, we use $`R_{\mathrm{abs}}`$ derived above to estimate the ambient UV radiation field at $`1000\mathrm{\AA }`$ from the expression,
$$R_{\mathrm{abs}}=4.0\times 10^{10}S_{\mathrm{shield}}\frac{J_{1000\mathrm{\AA }}}{J_{1000\mathrm{\AA },}}\mathrm{s}^1,$$
(6)
where $`J_{1000\mathrm{\AA },}3.2\times 10^{20}\mathrm{ergs}\mathrm{s}^1\mathrm{cm}^2\mathrm{Hz}^1\mathrm{Sr}^1`$ is the UV intensity at $`1000\mathrm{\AA }`$ in the solar vicinity (Hirashita & Ferrara 2005), and $`S_{\mathrm{shield}}`$ corrects for the the $`\mathrm{H}_2`$ self-shielding and/or dust extinction. Since the $`J=4`$ lines are optically thin, the self-shielding of $`\mathrm{H}_2`$ can be reasonably ignored, and we estimate the importance of dust extinction following Hirashita & Ferrara (2005), i.e.
$$S_{\mathrm{shield}}\mathrm{exp}[0.879\kappa (\frac{N_{\mathrm{HI}}}{10^{21}\mathrm{cm}^2})],$$
(7)
where $`\kappa `$ is the dust-to-gas ratio, defined as the ratio of the extinction optical depth to neutral hydrogen column density in units of $`10^{21}\mathrm{cm}^2`$ (Pei et al. 1991). Vladilo (1998) showed that $`\kappa `$ can be estimated by
$$\kappa =1.7\times 10^{[\mathrm{Zn}/\mathrm{H}]}(110^{[\mathrm{Fe}/\mathrm{Zn}]}).$$
(8)
Assuming that ZnII and FeII are the dominant ionization states for these two species, we adopt $`[\mathrm{Fe}/\mathrm{H}]=1.22`$ and $`[\mathrm{Fe}/\mathrm{Zn}]=0.87`$ for the Q1331+170 DLA from Prochaska & Wolfe (1999). This gives $`\kappa 0.088`$, and therefore $`S_{\mathrm{shield}}0.89`$, with our adopted H I column density (see Sec. 3.2). Inserting the above values into Eqn. 6, we get
$$\frac{J_{1000\mathrm{\AA }}}{J_{1000\mathrm{\AA },}}(2.1\pm 0.7)\times 10^3,$$
(9)
i.e. the ambient UV radiation field of the $`\mathrm{H}_2`$ cloud associated with DLA 1331+170 is about three orders of magnitude weaker than that in the Solar vicinity. Assuming the spectrum of the UV radiation field follows $`J_\nu \nu ^{0.5}`$, we estimate the UV intensity at the Lyman limit as
$$J_{912\mathrm{\AA }}(7.1\pm 2.3)\times 10^{23}\mathrm{ergs}\mathrm{s}^1\mathrm{cm}^2\mathrm{Hz}^1\mathrm{Sr}^1.$$
(10)
This value is consistent with the mean metagalactic UV background intensity of $`J_{912\mathrm{\AA }}7.6_{3.0}^{+9.4}\times 10^{23}\mathrm{ergs}\mathrm{s}^1\mathrm{cm}^2\mathrm{Hz}^1\mathrm{Sr}^1`$ at similar redshifts determined from proximity effect in the Ly$`\alpha `$ forest (Scott et al. 2002), implying that the ambient UV radiation of the $`\mathrm{H}_2`$ absorber may originate externally from the cloud. If instead the UV field is stellar in origin, $`J_{912\mathrm{\AA }}`$ would be lower than the metagalactic value, so our assumption of $`J_\nu \nu ^{0.5}`$ gives a reasonable upper limit to the ionizing radiation field. Based on the above calculations, the existence of any significant internal source of UV emission is not allowed, suggesting negligible star formation activity associated with this $`\mathrm{H}_2`$ absorber. Finally, we emphasize that our estimate of $`J_{912\mathrm{\AA }}`$ is based on the assumption of QSO-only spectral shape for the UV background intensity. This is a reasonable assumption if star-forming galaxies do not make significant contribution to the metagalactic UV radiation field (e.g. Leithere et al. 1995, Hurwitz et al. 1997, Deharveng et al. 2001, Giallongo et al. 2002, Fernรกndez-Soto et al. 2003).
The ambient UV intensity for DLA 1331+170 is much smaller than that in other DLAs, which is usually close to the typical Galactic value (e.g. Ge & Bechtold 1997, Levshakov et al. 2002, Ledoux et al. 2002). Hirashita & Ferrara (2005) estimated that the probable range of radiation field is $`0.5\mathrm{log}(\frac{J_{1000\mathrm{\AA }}}{J_{1000\mathrm{\AA },}})1.5`$, for a sample of DLAs which have been previously searched for $`\mathrm{H}_2`$ absorption (Ledoux et al. 2003). The associated surface star formation rate (SFR) is $`5\times 10^3\mathrm{M}_{}\mathrm{yr}^1\mathrm{kpc}^2\mathrm{\Sigma }_{\mathrm{SFR}}5\times 10^2\mathrm{M}_{}\mathrm{yr}^1\mathrm{kpc}^2`$ (Hirashita & Ferrara 2005, see also Wolfe et al. 2003a, 2003b). Also, we note that Reimers et al. (2003) reported the detection of $`\mathrm{H}_2`$ associated with a sub-DLA system, with local UV radiation field inferred to be $`300`$ times stronger than the mean Galactic value. Although most DLAs may be associated with considerable on-going star formation, the DLA toward Q1331+170 at least represents a special system, which is in a stage either before the onset of considerable star formation, or after star forming activities have ceased.
For comparison, we note that the C II\* 1336$`\mathrm{\AA }`$ absorption feature is useful for probing the ambient UV intensity near $`1500\mathrm{\AA }`$ (Wolfe et al. 2003a, 2003b). Thus it is interesting to compare the UV radiation field inferred by C II\* with that determined from $`\mathrm{H}_2`$. Using the C II\* absorption line as a diagnostic of the UV radiation field follows the idea that $`[\mathrm{CII}]\mathrm{\hspace{0.25em}158}\mu m`$ emission is the dominant coolant in DLAs, and is balanced by the grain photoelectric heating in a steady state. Since the grain photoelectric heating rate is proportional to the ambient UV intensity, a measurement of the strength of C II\* 1336$`\mathrm{\AA }`$ absorption provides an estimate of the UV radiation field at $`1500\mathrm{\AA }`$. Following Wolfe et al. (2003a), we express the heating rate per H atom as
$$\mathrm{\Gamma }=10^{24}\kappa ฯต(\frac{J_{1500\mathrm{\AA }}}{10^{19}\mathrm{ergs}\mathrm{s}^1\mathrm{cm}^2\mathrm{Hz}^1\mathrm{Sr}^1})\mathrm{ergs}\mathrm{s}^1\mathrm{H}^1,$$
(11)
where $`\kappa =0.088`$ is the dust-to-gas ratio determined from Eqn. 8, $`ฯต0.049`$ is the heating efficiency (Bakes & Tielens 1994, Wolfire et al. 1995), and $`J_{1500\mathrm{\AA }}`$ is the UV intensity at $`1500\mathrm{\AA }`$. For the DLA towards Q1331+170, the C II\* 1336$`\mathrm{\AA }`$ absorption line strength implies a cooling rate of $`l_\mathrm{c}2.24\times 10^{27}\mathrm{ergs}\mathrm{s}^1\mathrm{H}^1`$ (Wolfe et al. 2003a). Equating $`\mathrm{\Gamma }`$ and $`l_\mathrm{c}`$ implies $`J_{1500\mathrm{\AA }}5.4\times 10^{20}\mathrm{ergs}\mathrm{cm}^2\mathrm{s}^1\mathrm{Hz}^1\mathrm{Sr}^1`$, not too far from the mean Galactic value and much higher than the UV radiation field inferred from $`\mathrm{H}_2`$ absorption. This discrepancy is not expected, since the large difference in the UV intensity determined from $`\mathrm{H}_2`$ and C II\* is difficult to be interpreted by any local sources of absorption. Note that $`\mathrm{H}_2`$ and C II\* absorption probe UV intensity near $`1000\mathrm{\AA }`$ and $`1500\mathrm{\AA }`$, respectively, which are both above the Lyman limit. This indicates that photoionization of H I is not relevant. We also mentioned above that $`J_{1000\mathrm{\AA }}`$ determined from $`\mathrm{H}_2`$ absorption is characteristic of the UV radiation field outside the $`\mathrm{H}_2`$ cloud, therefore molecular hydrogen absorption in the Lyman and Werner bands does not account for the discrepancy either. However, the very low ambient UV intensity determined for the $`\mathrm{H}_2`$ absorber is reasonable, otherwise the residual excitation temperature inferred from C I/C I\* absorption would be far below the cosmic microwave background (CMB) temperature predicted by the standard cosmology (see Sec. 4.7).
### 4.2 The structure of the $`\mathrm{H}_2`$ absorber
To study the structure of the $`\mathrm{H}_2`$ cloud associated with the DLA of Q1331+170, we use CLOUDY (version c9400) to construct a grid of models with different total hydrogen number density $`n(\mathrm{H})n(\mathrm{HI})+2n(\mathrm{H}_2)`$, and examine which model provides the best match with the observed $`\mathrm{H}_2`$ column density. We adopt a metallicity of $`[\mathrm{Zn}/\mathrm{H}]=1.22`$ (Prochaska & Wolfe 1999) for all the elements in the absorber. To take into account differential dust depletion, we determine the dust-to-gas ratio $`\kappa `$ according to Eqn. 8. The ambient UV intensity is characterized by the metagalactic UV radiation field, assuming a QSO-like spectral shape, as implied by results in Sec. 4.1. For simplicity, we assume plane parallel geometry with a uniform total hydrogen number density. The cloud is illuminated from both sides in our model, and we run the calculation until the neutral hydrogen column density reaches the observed value of $`N(\mathrm{HI})=1.5\times 10^{21}\mathrm{cm}^2`$.
The model with a total hydrogen number density, $`n(\mathrm{H})n(\mathrm{HI})+2n(\mathrm{H}_2)0.22`$, gives the best match to the observed $`\mathrm{H}_2`$ fraction for the $`z=1.7765`$ DLA system. We show in Fig. 7.1 the thermal structure of the cloud, with the electron temperature determined from the balance between cooling and heating. In Fig. 7.2, we give the structure of the cloud traced by H I, H II, $`\mathrm{H}_2`$ and $`e`$. We find that the neutral hydrogen number density remains almost constant from the boundary deep into the central region. The molecular hydrogen fraction is very low in an outer shell with a thickness of $`46\mathrm{pc}`$, and increases sharply from $`<10^5`$ in the shell to $`10^2`$ in the center. We note that the outer shell is very thin, compared with the depth of the center region ($`1.2\mathrm{kpc}`$). The $`\frac{n(\mathrm{p})}{n(\mathrm{HI})}`$ ratio decreases monotonically with the depth into the cloud, and the metagalactic UV radiation field does produce some fraction of ionization in the outer shell of the cloud ($`\frac{n(\mathrm{p})}{n(\mathrm{H})}20\%`$). Fig. 7.1 shows that the electron temperature in the center $`\mathrm{H}_2`$ core is much lower than that in the outer shell, and has a mean value of $`T_\mathrm{e}140\mathrm{K}`$, in broad agreement with the excitation temperature determined in Sec. 3.4.
Fig. 7.3 shows the ionization structure of carbon for our model. The dominant ionization state is C I and C II. Although the ambient UV radiation results in some fraction of C III before itโs shielded by molecular hydrogen, the model calculation shows that the column density of highly ionized carbon species (e.g. C IV) is completely negligible.
Compared with observations, these results suggest that the Q1331+170 absorber has multiple velocity components with different carbon ionization. First, the observed value of $`N_{\mathrm{CII}}/N_{\mathrm{CI}}2647`$ for Q1331+170 (Liszt 2002) is much larger than the model prediction of $`3.5`$ for the $`\mathrm{H}_2`$ cloud. Second, there is strong CIV absorption associated with this damped Ly$`\alpha `$ system, with $`N(\mathrm{CIV})>10^{15}\mathrm{cm}^2`$ (Prochaska & Wolfe 1999). Therefore the bulk of the observed C II and C IV are associated with other components, either in a warm/hot phase, or embedded in a UV radiation field much stronger than that derived for the $`\mathrm{H}_2`$ absorber (see also Sec. 4.1).
Dessauges-Zavadsky et al. (2004) reported multiple velocity components associated with the DLA of Q1331+170, and they did find non-negligible differences between the component structure of low-ion transitions and that of intermediate-ion transitions. We note that the $`\mathrm{H}_2`$ cloud identified in our spectrum is different from the five intermediate-ion velocity components by $`56\mathrm{km}\mathrm{s}^1`$, $`28\mathrm{km}\mathrm{s}^1`$, $`15\mathrm{km}\mathrm{s}^1`$, $`15\mathrm{km}\mathrm{s}^1`$, and $`34\mathrm{km}\mathrm{s}^1`$ in the rest frame, respectively. On the other hand, it agrees with one low-ion component with detected C I absorption, with a velocity difference of only $`1.6\mathrm{km}\mathrm{s}^1`$ in the rest frame. Songaila et al. (1994) reported two C I absorption features associated with the same system, which are different from the $`\mathrm{H}_2`$ cloud by $`1.4\mathrm{km}\mathrm{s}^1`$ and $`19\mathrm{km}\mathrm{s}^1`$ in the rest frame. The $`1\sigma `$ uncertainty is $`1.2\mathrm{km}\mathrm{s}^1`$, caused by wavelength calibration and line centroiding (see Sec. 3.2), and therefore we conclude that the molecular hydrogen absorption identified in our STIS spectrum is consistent with one low-ion component reported by Dessauges-Zavadsky et al. (2004), as well as one C I component reported by Songaila et al. (1994). This is as expected since $`\mathrm{H}_2`$ tends to be co-spatial with neutral carbon and/or other low ion species.
### 4.3 Mass and dynamical state of the $`\mathrm{H}_2`$ cloud
Based on the model shown in Fig. 7.2, we estimate the mass of the entire cloud as $`M1.36\times \frac{4}{3}\pi R^3[n(\mathrm{HI})m_\mathrm{H}+n(\mathrm{H}_2)m_{\mathrm{H}_2}]6.5\times 10^7\mathrm{M}_{}`$, where we take $`R1.2\mathrm{kpc}`$, $`n(\mathrm{HI})0.17\mathrm{cm}^3`$, $`n(\mathrm{H}_2)0.05\mathrm{cm}^3`$, and the factor of $`1.36`$ corrects for the abundance of helium (Dickman 1978). We refer to this mass as $`M_{\mathrm{LTE}}`$, since our model is constructed based on the assumption of local thermodynamic equilibrium. Here we have ignored the small mass fraction contributed by H II. We also assume spherical geometry for the absorbing region.
The masses of giant molecular clouds (GMCs) have been measured for the Milky Way (e.g. Solomon et al. 1987, Digel et al. 1996, Heyer et al. 2001), and nearby galaxies such as LMC (Mizuno et al. 2001a), SMC (Mizuno et al. 2001b), and M33 (Engargiola et al. 2003). These masses are either determined based on the assumption of self-gravitational equilibrium, i.e. virial mass, $`M_{\mathrm{vir}}`$, or determined by adopting a constant CO-to-$`\mathrm{H}_2`$ conversion factor and then integrating the $`\mathrm{H}_2`$ column density over the projected area of the cloud, i.e. CO luminosity mass, $`M_{\mathrm{CO}}`$. In all cases, the mass spectrum of GMCs has been found to follow approximately a power law, with the power index ranging from $`1.8`$ for the Milky Way (Heyer et al. 2001) to $`2.6`$ for M33 (Engargiola et al. 2003). The mass of the $`\mathrm{H}_2`$ cloud associated with the DLA towards Q1331+170 is much larger than the values determined for most GMCs in either the Milky Way or nearby galaxies. For example, it is almost two orders of magnitude greater than the value of $`7\times 10^5\mathrm{M}_{}`$ for the most massive GMC in M33 (Engargiola et al. 2003). Similarly, if we adopt the GMC mass spectrum of $`\frac{dM}{dN}M^{1.8}`$ in the Milky Way, and take the limiting CO luminosity of $`138\mathrm{K}\mathrm{km}\mathrm{s}^1\mathrm{pc}^2`$, corresponding to $`566\mathrm{M}_{}`$ (Heyer et al. 2001), we find that the fraction of GMCs with masses greater than $`6.5\times 10^7\mathrm{M}_{}`$ (the value derived for the $`\mathrm{H}_2`$ cloud associated with DLA 1331+170) is only $`9\times 10^5`$ in the Milky Way. This may simply result from our assumption of spherical symmetry, or from the fact that the spectral resolution of our STIS spectrum is not high enough to resolve individual molecular clouds in velocity space along the line-of-sight. Note that the typical velocity dispersions measured for Galactic GMCs peak at $`\sigma _v1\mathrm{km}\mathrm{s}^1`$ (Heyer et al. 2001), much smaller than the spectral resolution of $`10\mathrm{km}\mathrm{s}^1`$ for our spectrum and the $`b`$-parameter of $`20\mathrm{km}\mathrm{s}^1`$ for the H I gas in DLA 1331+170. The total $`\mathrm{H}_2`$ mass inside a $`2\mathrm{kpc}`$ galactocentric radius has been determined for a sample of 17 nearby galaxies (Paglione et al. 2001), based on the observed CO emission and a standard CO-to-$`\mathrm{H}_2`$ conversion factor of $`1.6\times 10^{20}\mathrm{cm}^2\mathrm{K}^1\mathrm{km}^1\mathrm{s}`$. These masses range from $`1.6\times 10^8\mathrm{M}_{}`$ to $`2.8\times 10^9\mathrm{M}_{}`$, much larger than our value of $`6.5\times 10^7\mathrm{M}_{}`$. This is consistent with the general deficiency of molecular hydrogen in damped Ly$`\alpha `$ systems, especially when considering that the DLA towards Q1331+170 presents the largest $`\mathrm{H}_2`$ fraction in all high-z DLAs.
It is interesting to compare the mass derived above, i.e. $`M_{\mathrm{LTE}}`$, with the mass predicted by self-gravitational equilibrium, i.e. $`M_{\mathrm{vir}}`$. We estimate the virial mass by
$$M_{\mathrm{vir}}K\frac{\sigma _v^2R}{G}=9.6\times 10^8\mathrm{M}_{}.$$
(12)
Here $`\sigma _v`$ is the one-dimensional velocity dispersion which we approximate as the $`b`$-parameter determined in Sec. 3.2, $`R1.2\mathrm{kpc}`$ is the size of the cloud (see Sec. 4.2), $`G`$ is the gravitational constant, and $`K`$ is a dimensionless factor of order unity that depends on the geometry as well as the density profile of the cloud. We adopt $`K8.7`$ from Solomon et al. (1987). Correspondingly, the gravitational parameter, $`\alpha _\mathrm{G}=\frac{M_{\mathrm{vir}}}{M_{\mathrm{LTE}}}14`$, indicating that the $`\mathrm{H}_2`$ cloud can not be bound by self-gravity. This is contrary to the situation for most Galactic GMCs, of which the self-gravitational equilibrium state has long been established (e.g. Solomon et al. 1987). More recently, Heyer et al. (2001) found that the dynamical state of Galactic GMCs varies with the cloud mass, in that $`M_{\mathrm{CO}}>10^4\mathrm{M}_{}`$ clouds are bound by self-gravity. One possible interpretation to the large $`\alpha _\mathrm{G}`$ value is that the $`\mathrm{H}_2`$ cloud towards Q1331+170 is in hydrostatic equilibrium, i.e. bound by the pressure of an external medium. Here we emphasize that the above result is subject to uncertainties due to the unknown geometry and density profile of the cloud.
### 4.4 Molecular hydrogen fraction
In Sec. 3.2, we derived a molecular hydrogen fraction of $`f_{\mathrm{H}_2}=(5.6\pm 0.7)\%`$, for the $`z=1.7765`$ DLA toward Q1331+170, which is the largest value reported so far in any redshifted damped Ly$`\alpha `$ absorber. The large value of $`f_{\mathrm{H}_2}`$ is consistent with the relatively high dust abundance for this system. In Sec. 4.2, we derived the dust-to-gas ratio to be $`\stackrel{~}{k}0.088`$. Ledoux et al. (2003) showed that there exists a threshold at $`\stackrel{~}{k}0.03`$ above which the molecular hydrogen fraction tends to be large ($`f_{\mathrm{H}_2}>10^4`$). Moreover, the model described in Sec. 4.2 suggests a low gas phase temperature, which is also supported by the distribution of $`\mathrm{H}_2`$ at different $`J`$ states. The low kinetic temperature is representative of a cold neutral medium (CNM) phase, in which efficient formation of $`\mathrm{H}_2`$ molecules onto dust grains is possible. Third, the ambient UV radiation field of the $`\mathrm{H}_2`$ cloud is extremely low, implying that destruction of $`\mathrm{H}_2`$ molecules by photo-dissociation is ineffective.
It is interesting to compare the damped Ly$`\alpha `$ absorber toward Q1331+170 with other high-redshift damped systems. In Table 4, we compiled data for all the 42 DLAs which have been searched for molecular hydrogen absorption, including Q1331+170. In the upper panel of Fig. 8, we show the distribution of damped Ly$`\alpha `$ absorbers with respect to neutral hydrogen column density and molecular hydrogen fraction. The figure clearly shows the bi-model distribution of DLAs, i.e. for a certain value of $`N_{\mathrm{HI}}`$, the molecular hydrogen fraction of the absorber can be either as large as $`10^110^3`$, or smaller than the detection limit ($`10^510^7`$). This can be interpreted by the self-shielding effect of molecular hydrogen: once molecular hydrogen begins to form, it shields itself from subsequent photodissociation by UV photons, and thus the molecular hydrogen fraction increases significantly (see also Ge & Bechtold 1999). The lower panel of Fig. 8 shows the relation between the molecular hydrogen fraction and the dust abundance in the same system. The dust abundance is characterized by the column density ratio of undepleted element to depleted element. For element which is not depleted, we take Zn as a representative, whereas when Zn is not available in the literature, we adopt S instead. The element which is depleted to the dust grain is represented by Fe, and when Fe abundance has not been measured, we take Cr instead. The lower panel of Fig. 8 shows a strong correlation between the dust abundance and the molecular hydrogen fraction, similar to previous results (e.g. Ledoux et al. 2003). Kendallโs $`\tau `$ test shows the correlation is present at $`>99.9\%`$ confidence level. In Fig. 8, the damped Ly$`\alpha `$ absorber toward Q1331+170 is indicated by a solid circle.
### 4.5 Dust extinction
The reddening of QSOs with intervening DLAs compared with non-DLA QSOs has been detected and quoted as evidence for the presence of dust in DLAs (e.g. Fall & Pei 1989, Pei et al. 1991). More recently, Murphy & Liske (2004) showed that there is no evidence for the dust reddening of QSOs with DLAs, based on a much larger sample selected from the Sloan Digital Sky Survey Data Release 2 (SDSS DR2) catalog. However, since the SDSS main QSO sample is largely color-selected (Richards et al. 2002), the non-detection by Murphy & Liske (2004) may simply be a selection effect against heavily reddened objects. In this section, we estimate the dust reddening of Q1331+170, due to the damped Ly$`\alpha `$ absorber at $`z1.7765`$, and we also compare the reddening of QSOs for which the intervening DLAs have different molecular hydrogen fractions.
Similar to Pei (1992), we calculate the color excess, $`E_{\mathrm{B}\mathrm{V}}`$, of Q1331+170 by
$$E_{\mathrm{B}\mathrm{V}}=1.086\frac{\kappa }{1+R_\mathrm{V}}(\frac{N_{\mathrm{HI}}}{10^{21}\mathrm{cm}^2}),$$
(13)
where $`\kappa 0.088`$ for DLA 1331+170 is the dust-to-gas ratio determined from Eqn. 8, and $`R_\mathrm{V}`$ is the ratio of total-to-selective extinction, which depends on the exact extinction curve adopted in our calculation. If we assume the SMC extinction curve, then $`R_\mathrm{V}=2.93`$ (Pei 1992) and $`E_{\mathrm{B}\mathrm{V}}0.037\pm 0.005`$. We also calculate $`E_{\mathrm{B}\mathrm{V}}`$ by adopting the mean $`N_{\mathrm{HI}}/E_{\mathrm{B}\mathrm{V}}`$ values of $`36.1\pm 3.3`$ based on the new determination of the SMC extinction curve (for the SMC bar sample) by Gordon et al. (2003), where $`N_{\mathrm{HI}}`$ is in units of $`10^{21}\mathrm{cm}^2`$. This gives the same value of $`E_{\mathrm{B}\mathrm{V}}`$ within $`1\sigma `$ errors. There are two reasons for adopting the SMC extinction curve. First, the DLAs are believed to be in an early stage of chemical evolution, which is well characterized by the SMC extinction curve (e.g. Pei 1992). Second, DLAs show no evidence for the $`2175\mathrm{\AA }`$ dust feature observed in either Galactic or LMC extinction curve (e.g. Fall & Pei 1989, but see Malhotra 1997 which reported the detection of this feature based on a sample of Mg II absorbers).
In column 7 of Table 4, we give $`E_{\mathrm{B}\mathrm{V}}`$ for each QSO of which the intervening DLA has previous measurements of both depleted and undepleted metal absorption. These values are calculated similar to the Q1331+170 case. However, we use a different normalization when estimating the dust-to-gas ratio by Eqn. 8, depending on the types of depleted/undepleted elements adopted in the calculation (see Table 1 of Vladilo 1998). For several QSOs with more than one intervening DLAs (Q0013-004, Q0405-443, Q0841+129 and Q2059-360), we only list in Table 4 $`E_{\mathrm{B}\mathrm{V}}`$ contributed by individual absorbers, which does not represent the total extinction towards the background QSO. Fig. 9 presents the distribution of $`E_{\mathrm{B}\mathrm{V}}`$ for DLAs with firm $`\mathrm{H}_2`$ detections (solid line) as well as the distribution for DLAs for which only upper limits of $`f_{\mathrm{H}_2}`$ have been put (dashed line). The extinction for Q1331+170 is marked with an arrow. The two distributions in Fig. 9 are different at $`99.8\%`$ confidence level, based on Kolmogorov-Smirnov test. Specifically, the sub-sample of DLAs with firm $`\mathrm{H}_2`$ detections shows more extinction compared with other DLAs. The mean color excess for absorbers with firm $`\mathrm{H}_2`$ detections is $`E_{\mathrm{B}\mathrm{V}}=0.038\pm 0.001`$, significantly greater than the value of $`E_{\mathrm{B}\mathrm{V}}=0.014\pm 0.001`$ for cases in which only upper limits of $`f_{\mathrm{H}_2}`$ have been reported. Whereas only $`12.5\%`$ of DLAs with firm $`\mathrm{H}_2`$ detections have $`E_{\mathrm{B}\mathrm{V}}`$ smaller than $`0.01`$, this fraction is as large as $`75\%`$ for other DLAs. We have also examined the relation between $`E_{\mathrm{B}\mathrm{V}}`$ and $`f_{\mathrm{H}_2}`$, but no significant correlation has been found.
The intrinsic QSO continuum in the UV/optical band is reddened by the dust extinction due to an intervening DLA. Assuming the QSO continuum follows $`F_\nu \nu ^\alpha `$ or equivalently $`F_\lambda \lambda ^{\alpha 2}`$, we get $`\alpha 4.1`$ by fitting a power law to the dereddened continuum (between $`2600\AA `$ and $`3000\AA `$) of the QSO, Q1331+170, while $`\alpha =4.2`$ for the observed, reddened continuum adopted in Sec. 2. For dereddening of the Galactic extinction, we use the extinction map of Schlegel et al. (1998). For reddening of the intervening damped Ly$`\alpha `$ absorber at $`z1.7765`$, we adopt the FM parametrization (Fitzpatrick & Massa 1990) determined for the SMC bar sample (Gordon et al. 2003). The observed spectrum has been shifted to the rest frame of the absorber before dereddening. The value of $`4.1`$ for the dereddened spectral slope indicates an extremely red intrinsic spectrum below the rest frame wavelength of $`970\AA `$ for Q1331+170.
### 4.6 CO-to-$`\mathrm{H}_2`$ ratio
According to our model calculations in Sec. 4.2, the predicted CO column density is $`6.7\times 10^{10}\mathrm{cm}^2`$. This is consistent with the observed $`2\sigma `$ upper limit of $`N_{\mathrm{CO}}<1.1\times 10^{13}\mathrm{cm}^2`$ (Levshakov et al. 1988), determined from UV absorption lines of CO ($`A^1\mathrm{\Pi }X^1\mathrm{\Sigma }^1`$). Adopting the total molecular hydrogen column density in Table 3, we obtained the observed CO-to-$`\mathrm{H}_2`$ column density ratio of $`\frac{\mathrm{N}_{\mathrm{CO}}}{\mathrm{N}_{\mathrm{H}_2}}<2.5\times 10^7`$.
The upper limit of the CO-to-$`\mathrm{H}_2`$ column density ratio for DLA 1331+170 is similar to the typical value of CO-to-$`\mathrm{H}_2`$ measured from UV absorption of CO in Galactic diffuse clouds, where the effects of CO self-shielding are not important (e.g. Crenny & Federman 2004). For dense molecular clouds in the Milky Way, the CO/$`\mathrm{H}_2`$ conversion factor is given in terms of the quantity $`\frac{I_{\mathrm{CO}}}{N_{\mathrm{H}_2}}`$, where $`I_{\mathrm{CO}}`$ is related to the antenna temperature $`T_\mathrm{A}`$ directly measurable from millimeter CO emission lines (e.g. Strong et al. 1988). $`\frac{I_{\mathrm{CO}}}{N_{\mathrm{H}_2}}`$ usually does not provide information on the corresponding column density ratio, since the profiles of the optically thick CO lines give the velocity width of molecular clouds rather than the CO column densities. The ratio of CO to $`\mathrm{H}_2`$ column densities in typical dense molecular clouds in the Milky Way is $`10^4`$, three orders of magnitude higher than the upper limit determined for the DLA toward Q1331+170 (e.g. Frerking et al. 1992, Lacy et al. 1994). This could be associated with the relatively low number density of the $`\mathrm{H}_2`$ cloud (see Sec. 4.2), since the CO abundance is controlled by collisional reactions involved with various species, including O, C, $`\mathrm{H}_2`$, CH and CH<sub>2</sub> (e.g. Tielens & Hollenbach 1985).
### 4.7 Measurement of CMB Temperature at $`z=1.77654`$
Songaila et al. (1994) reported the detection of C I/C I\* absorption features associated with a velocity component at $`z=1.77654`$. As described in Sec. 4.2, the redshift of this component agrees with the redshift of the $`\mathrm{H}_2`$ absorber identified in our spectrum at $`1.2\sigma `$. Therefore, the C I and $`\mathrm{H}_2`$ probably arises in the same gas cloud. Theoretically, C I and $`\mathrm{H}_2`$ tend to be co-spatial as well (Ge et al. 1997, 2001). With the physical parameters determined in Sec. 4.1 and 4.2, we can estimate the local contribution to the C I excitation in this system and then put constraints on the CMB temperature at the absorberโs redshift. The equilibrium between the population and de-population of the C I $`{}_{}{}^{3}P_{0}^{}`$ and $`{}_{}{}^{3}P_{1}^{}`$ fine structure can be written as
$$N_0(B_{01}I_\nu +\mathrm{\Gamma }_{01}+\underset{j}{}R_{01}^jn_j)=N_1(A_{10}+B_{10}I_\nu +\mathrm{\Gamma }_{10}+\underset{j}{}R_{10}^jn_j),$$
(14)
where $`R_{01}`$ and $`R_{10}`$ are collisional excitation and de-excitation rates, with $`j`$ representing different collision partners (H, He, $`e`$, $`p`$, or $`H_2`$), $`\mathrm{\Gamma }_{01}`$ and $`\mathrm{\Gamma }_{10}`$ are UV pumping rates, $`A_{10}=7.93\times 10^8\mathrm{s}^1`$ is the probability of the C I fine structure transition $`{}_{}{}^{3}P_{1}^{}^3P_0`$ (Bahcall & Wolf 1968). $`B_{01}I_\nu `$ and $`B_{10}I_\nu `$ represent excitation or de-excitation due to the absorption of ambient microwave photons, which can be expressed as
$$B_{01}I_\nu =3B_{10}I_\nu =\frac{2.38\times 10^7}{\mathrm{exp}(23.6\mathrm{K}/T_{\mathrm{ex}})1}\mathrm{s}^1.$$
(15)
The collisional rates, $`R_{01}^j`$ and $`R_{10}^j`$ are taken from Launay & Roueff (1977), Johnson et al. (1987), Roueff & Le Bourlot (1990), Staemmler & Flower (1991), and Schrรถder et al. (1991), and are related through the Milne relation,
$$R_{10}^j=\frac{1}{3}R_{01}^j\mathrm{exp}(23.6\mathrm{K}/T_\mathrm{K}).$$
(16)
We take number densities of different species from the model described in Sec. 4.2. The various collision terms are calculated to be $`R_{01}^{\mathrm{HI}}n(\mathrm{HI})=8.6\times 10^{10}\mathrm{s}^1`$, $`R_{01}^{\mathrm{He}}n(\mathrm{He})=3.3\times 10^{13}\mathrm{s}^1`$, $`R_{01}^{\mathrm{H}_2}n(\mathrm{H}_2)=4.2\times 10^{13}\mathrm{s}^1`$, $`R_{01}^\mathrm{p}n(\mathrm{p})=2.3\times 10^{12}\mathrm{s}^1`$, and $`R_{01}^\mathrm{e}n(\mathrm{e})=1.4\times 10^{12}\mathrm{s}^1`$, assuming a kinetic temperature of $`152\mathrm{K}`$. We determine the UV pumping rates for C I by scaling from the mean Galactic values (e.g. Ge et al. 1997), and get $`\mathrm{\Gamma }_{01}=8.7\times 10^{13}\mathrm{s}^1`$ and $`\mathrm{\Gamma }_{10}=2.8\times 10^{13}\mathrm{s}^1`$. Inserting these values into Eqn. 14 and 15, and taking $`\frac{N(^3P_1)}{N(^3P_0)}=0.125\pm 0.042`$ from Songaila et al. (1994), we get $`T_{\mathrm{ex}}=(7.2\pm 0.8)\mathrm{K}`$ at $`z=1.77654`$. Note that the uncertainty in the calculated $`T_{\mathrm{ex}}`$ is associated with the errors in the measured C I/C I\* column densities. The excitation temperature, $`T_{\mathrm{ex}}`$, calculated above is consistent with the predicted CMB temperature of $`2.725\times (1+z)=7.566\mathrm{K}`$ at $`z=1.77654`$, within $`1\sigma `$.
## 5 Conclusions
The main results of this paper are summarized as follows:
1. Using the NUV spectrum with spectral resolution of $`10\mathrm{km}\mathrm{s}^1`$ obtained with the E230M grating of HST/STIS, we detected strong molecular hydrogen absorption associated with the $`z=1.7765`$ damped Ly$`\alpha `$ system toward Q1331+170. The total $`\mathrm{H}_2`$ column density is $`N_{\mathrm{H}_2}=(4.45\pm 0.36)\times 10^{19}\mathrm{cm}^2`$, determined from simultaneous Voigt profile fitting of 26 $`\mathrm{H}_2`$ lines.
2. The molecular hydrogen fraction was determined to be $`f_{\mathrm{H}_2}=\frac{2N_{\mathrm{H}_2}}{N_{\mathrm{HI}}+2N_{\mathrm{H}_2}}=(5.6\pm 0.7)\%`$, which is the largest value reported so far in any redshifted damped Ly$`\alpha `$ system. This is a combined effect of a relatively high dust-to-gas ratio, a low gas temperature, and an extremely low ambient UV radiation field.
3. We detect rotationally excited transitions of $`\mathrm{H}_2`$, with $`J=04`$. The relative column densities of the $`J=0,1,2,3,4`$ states can be fit with a single temperature of $`T_{\mathrm{ex}}=152\pm 10\mathrm{K}`$. This suggests that the derived excitation temperature represents the kinetic temperature of the gas, and that photo-excitation is negligible.
4. The photo-absorption rate was estimated to be $`R_{\mathrm{abs}}=(7.6\pm 2.4)\times 10^{13}\mathrm{s}^1`$, corresponding to an ambient UV radiation field of $`J_{912\mathrm{\AA }}7\times 10^{23}\mathrm{ergs}\mathrm{s}^1\mathrm{cm}^2\mathrm{Hz}^1\mathrm{Sr}^1`$. This is comparable with the metagalactic UV background intensity at the same redshift, and implies an extremely low star-formation rate in the DLA environment.
5. We constructed a simple model to describe the structure of the $`\mathrm{H}_2`$ absorber, with a best-fit total hydrogen number density of $`n_\mathrm{H}0.2\mathrm{cm}^3`$ and an electron temperature of $`T_\mathrm{e}140\mathrm{K}`$. The total mass of the model cloud is $`6.5\times 10^7\mathrm{M}_{}`$, greater than masses of most GMCs in the Milky Way and nearby galaxies. We also find that this mass is considerably smaller than the virial mass of the cloud, indicating that the $`\mathrm{H}_2`$ cloud is not in a state of self-gravitational equilibrium. Our model is consistent with the population of $`\mathrm{H}_2`$ at different $`J`$ states, as well as the observational upper limit put on the CO column density. The observed CO-to-$`\mathrm{H}_2`$ column density ratio is $`\frac{\mathrm{N}_{\mathrm{CO}}}{\mathrm{N}_{\mathrm{H}_2}}<2.5\times 10^7`$, characteristic of the typical value measured for diffuse molecular clouds in the Galactic ISM.
6. Adopting the SMC-like extinction curve, we calculate the extinction of Q1331+170 to be $`E_{\mathrm{B}\mathrm{V}}=0.037\pm 0.005`$, due to the intervening damped Ly$`\alpha `$ absorber at $`z1.7765`$. We also find that the extinction of QSOs by foreground DLAs with firm $`\mathrm{H}_2`$ detections is considerably greater than those for which only upper limits of $`f_{\mathrm{H}_2}`$ have been put.
7. The redshift of the $`\mathrm{H}_2`$ absorber identified in our spectrum is consistent with a C I/C I\* absorber reported by Songaila et al. (1994). We calculated the local contribution to C I excitation, including UV pumping and collisions. The residual excitation temperature was determined to be $`T_{\mathrm{ex}}=(7.2\pm 0.8)\mathrm{K}`$, consistent with the predicted $`T_{\mathrm{CMB}}`$ of $`7.566\mathrm{K}`$ at $`z=1.77654`$.
We are grateful to R.I. Thompson, D.J. Eisenstein, X. Fan, D. Welty, J. Black, R. Davรฉ, J. Bieging and the anonymous referee for their suggestions, which have greatly improved this work. We also thank R.F. Carswell and G. Ferland for making the Voigt profile fitting program VPFIT and the photoionization code CLOUDY publicly available. This work is based on observations made with the NASA/ESA Hubble Space Telescope, obtained at the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS 5-26555. These observations were supported by NASA through Grants HST-GO-09172.01 and HST-GO-07271.01 from the Space Telescope Science Institute.
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# Nonextensive aspects of self-organized scale-free gas-like networks
## I Introduction
A common feature in several recent approaches to complex networks with statistical mechanical methods is the definition of network Hamiltonians newman\_sm ; vicsek ; berg04 . Such Hamiltonians currently depend on the number of links either on a global level, or on the degree of individual nodes. This definition allows to borrow powerful concepts from statistical physics such as the maximum entropy principle newman\_sm , which may provide the most probable ensembles of networks. Further, Hamiltonians allow to define both thermodynamical ensembles (microcanonical, canonical, grand canonical) vicsek ; dormen03 and a partition function, which opens the way to compute degree correlation functions in a formalism most familiar to physicists newman\_sm ; berg04 .
However, these approaches do not yet aim to explain the structure of degree distributions, and mainly address random networks. A conceptually different approach has been taken recently in tsallis\_sm , where networks are embedded in some metric space and the definition of entropy in networks is broadened. In this work it was noted that the characteristic distribution of the relevant degree of freedom โ the degree of nodes โ appears to coincide to distribution functions known for nonextensive systems tsallis88 ; gellman . More precisely, in tsallis\_sm it was found that, for some preferential attachment network growth models, the resulting degree distributions are of the $`q`$-exponential type (defined later on). In the usual preferential attachment model, the probability of a new node (i.e., being added to a network) to attach its link to a pre-existing node $`i`$ is proportional to this nodeโs number of links, or degree $`k_i`$, i.e., $`p_Ak_i`$. This is also true for networks embedded in $`R^n`$, where the linking probabilities are made dependent on the relative distance of node $`i`$ to the new node, i.e., $`p_Ak_i/d_{ij}^\alpha `$. Here $`\alpha `$ is a free parameter that defines the connecting horizon of a new node to the system. For large $`\alpha `$ the node will link with high probability to a nearby node, whereas distance becomes irrelevant for $`\alpha 0`$.
A problem which has not yet been explored in the literature is that of a definition of an interaction between nodes, for example in the way one would think of an interaction of gas molecules. In classical statistical mechanics interactions/collisions between gas particles result in a transfer of momentum from one particle to another, under the constraint that momentum is conserved. In elastic interactions this results in a change of direction and speed of particles after a collision, in inelastic ones also in a change of masses of the particles. In this paper we find that the class of self-organizing networks as introduced in sneppen opens the possibility to define an โinteractionโ between nodes of a network. In analogy to the momentum transfer in the classical situation, in the network case the interactions are defined by a transfer of links. This enables one to think of a network as some sort of gas, which turns out to be describable by distribution functions characteristic of nonextensive statistical mechanics.
## II Model
Let us consider the following gas-like system. Particles have links among them. The total number of links of a given node is a characteristic quantity of the node, such as the momentum of an ordinary particle. The particles of this โgasโ have no momentum but only their degree. Neither do the particles have an absolute position in space. They possess only a relative distance $`d_{ij}`$ to each other which is given by the shortest number of links between them (sometimes called chemical distance). Particles can interact โnon-elasticallyโ. Upon an interaction one particle ceases to exist and transfers all its links to the other. If the interacting particles $`i`$ and $`j`$ have both had links to a third particle $`k`$ before the interaction, the remaining particle $`i`$ will keep its link to $`k`$, while the links of the disappeared particle $`j`$ to $`k`$ will be removed, meaning that links are only counted once (and are not weighted). Consider these interactions taking place in a chemostat, such that the number of particles in a closed system is constant. This means that, for every merging interaction, a new particle will be introduced to the system. The interaction is characterized by the probability that two particles meet and transfer links. Given a โmetricโ (relative network distance) this probability can be made distance-dependent as in tsallis\_sm . To do this we introduce a power like potential.
In what follows we numerically study the distribution of the characteristic degree of freedom (the degree of nodes) as a function of the range of the interaction. As in tsallis\_sm , we find that the distribution of this (nonextensive) โgasโ is described by a $`q`$-exponential.
### II.1 Network dynamics
There has been a recent suggestion to model scale-free networks of constant size, the so-called self-organized networks sneppen . The idea is that at any given timestep one single pair of nodes, say $`i`$ and $`j`$, merge together to become one single node. This node keeps the name of one of the original nodes, say for example $`i`$. This node now gains all the links of the other node $`j`$, resulting in a change of degree for node $`i`$ according to
$`k_i`$ $``$ $`k_i+k_jN_{\mathrm{common}},\mathrm{if}(i,j)\mathrm{not}\mathrm{first}\mathrm{neighbors}`$
$`k_i`$ $``$ $`k_i+k_jN_{\mathrm{common}}2,\mathrm{if}(i,j)\mathrm{first}\mathrm{neighbors}`$
where $`N_{\mathrm{common}}`$ is the number of nodes, which shared links to both of $`i`$ and $`j`$ before the merger. In the case that $`i`$ and $`j`$ were first neighbors before the merger, i.e., they had been previously linked, the removal of this link will be taken care of by the term $`2`$ in Eq. (LABEL:update). Next, to keep the system at constant size, a new node is created, and is linked to $`r`$ randomly chosen nodes from the existing network. Let us note here that the smallest degree found in a network can only be larger or equal to $`r`$. This will have consequences for the normalization of distribution functions as will be discussed below. In the following (except for Fig. 2) we will restrict ourselves to $`r=2`$, for simplicity. Thus the smallest degree will always be 2. (This is not at all an important restriction; as an alternative the actual number of links can also be a random number picked, e.g. from a uniform distribution with an average of $`r`$, as in sneppen ). After that we address the next timestep. Nodes in the network start with a small number of links, and gain links through merger-interactions. Links to the $`N_{\mathrm{common}}`$ common neighbors of two merging nodes are lost as mentioned above, which reduces the number of links. Gains and losses eventually lead to an effective balance over time as shown for instance in Fig. 1 a. The number of links of the best connected node in the system is followed over time. After about 1000 timesteps a stationary state is reached. The situation for an individual link is shown in Fig. 1 b. A node starts with $`r=2`$ links when it is introduced to the system. It gains links through mergers over time. When the node is taken up in a merger it disappears from the system and, as said before, a new one starts with $`r=2`$ links again. Networks with these rewiring scheme lead to scale-free degree distributions sneppen , i.e. the power exponent of the cumulative degree distribution tail behaves as, $`P(k)k^\gamma `$. In sneppen two schemes were discussed: The case where only nodes being first-neighbors can merge, and the case where any two nodes โ directly connected or not โ can merge with the same probability for each possible pair $`(i,j)`$. The neighbor scheme leads to an exponent $`\gamma 1.3`$, the random scheme to $`\gamma 0.5`$. The cumulative degree distribution, for the neighbor scheme is shown in Fig. 2.
These distributions can be fit by q-exponentials,
$$P(k)=e_{q_c}^{(k2)/\kappa }(k=2,3,4,\mathrm{}),$$
(2)
where the $`q`$-exponential function is defined, for $`1+(1q_c)x0`$, as
$$e_{q_c}^x[1+(1q_c)x]^{1/(1q_c)}$$
(3)
with $`\kappa >0`$ some characteristic number of links, and $`\gamma 1/(q_c1)`$ being the tail exponent of the (asymptotic) power-law distribution. Whenever we talk about q-values corresponding to a cumulative distribution, we use the notation $`q_c`$, where $`c`$ stands for cumulative. We have normalized with the value corresponding to the smallest possible degree (which in our case equals 2) in order to have $`P(2)=1`$.
A convenient procedure to perform a two-parameter fit of this kind is to take the $`q`$-logarithm of the distribution $`P`$, defined by $`Z_q\mathrm{ln}_qP(k)\frac{[P(k)]^{1q_c}1}{1q_c}`$. This is done for a series of different values of $`q_c`$. The function $`Z_q`$ which is best fit with a straight line determines the value of $`q_c`$, the slope being $`\kappa `$. The situation for the $`N=2^{14}`$ data of Fig. 2 is shown in Fig. 3 for $`q_c`$ running between 1 and 2.6.
We numerically verify with good precision that the Ansatz in Eq. (2) for the cumulative degree distribution is a satisfactory one (it could even be the exact one for the present problem). This reveals a connection tsallis\_sm ; doye of scale-free network dynamics to nonextensive statistical mechanics tsallis88 ; gellman . Let us be more specific. Consider the entropy
$`S_q`$ $``$ $`{\displaystyle \frac{1_2^{\mathrm{}}๐k[p(k)]^q}{q1}}`$
$`[S_1`$ $`=`$ $`S_{BG}{\displaystyle _2^{\mathrm{}}}dkp(k)\mathrm{ln}p(k)],`$ (4)
where we assume $`k`$ as a continuous variable for simplicity, and $`BG`$ stands for Boltzmann-Gibbs. If we extremize $`S_q`$ with the constraints TsallisMendesPlastino
$$_2^{\mathrm{}}๐kp(k)=1$$
(5)
and
$$\frac{_2^{\mathrm{}}๐kk[p(k)]^q}{_2^{\mathrm{}}๐k[p(k)]^q}=K,$$
(6)
we obtain
$$p(k)=\frac{e_q^{\beta (k2)}}{_2^{\mathrm{}}๐k^{}e_q^{\beta (k^{}2)}}=\beta (2q)e_q^{\beta (k2)}(k2),$$
(7)
where $`\beta `$ is determined through Eq. (6). Both positivity of $`p(k)`$ and normalization constraint (5) impose $`q<2`$. The corresponding cumulative distribution $`P(>k)`$ is then given by
$$P(>k)1_2^k๐k^{}p(k^{})=[1(1q)\beta (k2)]^{\frac{2q}{1q}}.$$
(8)
This expression can be rewritten precisely as the Ansatz (2) with
$$q_c\frac{1}{2q};\kappa \frac{1}{(2q)\beta }.$$
(9)
### II.2 Network distance and distance-dependent re-linking potential
Unlike the two schemes in the original presentation of self-organizing networks, the neighbor and the random scheme, we would like to define a distance-dependent merging probability. This needs a definition of distance on the graph. For simplicity we define the distance $`d_{ij}`$, between two nodes $`i`$ and $`j`$ on an undirected graph as the shortest distance, given all links are of unit length. This distance is obtained, for instance, from the Dijkstra algorithm dijk . We randomly choose a node (denoted by $`i`$) and then choose the second node (denoted by $`j`$) of the merger with probability
$$p_{ij}=๐ฉd_{ij}^\alpha (\alpha 0),$$
(10)
where $`๐ฉ=1/_jd_{ij}^\alpha `$ is a normalization that makes $`p_{ij}`$ a probability, and $`d_{ij}`$ is the shortest distance (path) on the network connecting nodes $`i`$ and $`j`$; $`\alpha `$ is a real number. Obviously, tuning $`\alpha `$ from $`0`$ toward large values switches the model from the random to the neighbor scheme in sneppen .
## III Results
Realizing this distance-dependent potential in a numerical simulation we find the degree distributions given in Fig. 4. All following data was obtained from averages over 1000 identical realizations of degree distributions of networks with a number of nodes, $`N=2^9`$. Networks have been recorded after 5 network updates. A network update is performed when $`N`$ mergers have taken place. This corresponds to the $`5\times 2^9`$ timesteps shown in Fig. 1.
From these degree distributions we obtain the index $`q_c`$ and the characteristic degree $`\kappa `$. Their dependence on $`\alpha `$ is given in Fig. 5. The $`\alpha `$-dependence of $`q_c`$ shows the expected tendency. Our result in the limit $`\alpha 0`$, $`q_c(0)=2.8`$ corresponds to an exponent $`\gamma =0.55`$, which is about 10 percent lower than the reported value in sneppen . The reason for this small discrepancy seems to be a finite-size effect.
To demonstrate that this might indeed be so, in Fig. 5 b we plot the $`q_c(\alpha )`$-dependence for $`N=2^7`$, $`N=2^8`$ and $`2^9`$ networks for comparison. As network size increases the value of $`q_c`$ approaches the expected value of $`3`$ sneppen in the small $`\alpha `$ limit. For the $`\alpha \mathrm{}`$ limit, the expected value is recovered for the $`N=2^9`$ network, for smaller networks, there is still a visible size effect. This size effect is related to the problem that the finite size cutoff plays a relatively large role, and interferes considerably in the fits in small networks. Simulations on larger networks are certainly desirable, but inaccessible to our present computational power.
## IV Discussion
To summarize, we have explored the possibility to making some connection between a nonextensive gas and a self-organized scale-free network. We have shown that the characteristic degree distributions are well described by q-exponentials whose parameters vary with the interaction range, i.e. $`\alpha `$. The limiting cases $`\alpha 0`$ and $`\alpha \mathrm{}`$ reproduce the situations given in sneppen , namely the neighbor merging and the random merging.
In the present work we have used the networkโs intrinsic metric space, i.e., its adjacency matrix, to measure distances ($`d_{ij}`$) between nodes to be merged. This is in variance with what is done in tsallis\_sm , where the network is embedded in a โgeographicalโ metric space (with distances $`r_{ij}`$). Both models can of course be unified by introducing both merging (with probability $`1/d_{ij}^\alpha `$) and distance-dependent linking (with probability $`1/r_{ij}^{\alpha _A}`$, where $`A`$ stands for attachment). The degree distribution of such a unified model could still be of the $`q`$-exponential form with $`q_c(\alpha ,\alpha _A)`$. Of course, $`q_c(\alpha ,0)=q_c(\alpha )`$ as given in the present paper. It would not be surprising if $`q_c(\alpha ,\alpha _A)`$ was a monotonically decreasing function of both variables $`\alpha `$ and $`\alpha _A`$, with the maximal value being $`q_c(0,0)`$, and with say $`q_c(\alpha ,\mathrm{})=1,\alpha >0`$. In such a case, the interval spanned by $`q_c`$ would clearly be wider than that of the present model.
S.T. would like to thank the SFI and in particular J.D. Farmer for their great hospitality and support during Sept-Oct of 2004, when this work was initiated. Support from SI International and AFRL/USA is acknowledged as well.
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# Quantum-information theoretic properties of nuclei and trapped Bose gases
## 1 Introduction
Information-theoretic methods are used in recent years for the study of quantum mechanical systems. <sup>-</sup> The quantity of interest is Shannonโs information entropy for a probability distribution $`p(x)`$
$$S=p(x)\mathrm{ln}p(x)๐x$$
(1)
where $`p(x)๐x=1`$.
An important step is the discovery of an entropic uncertainty relation (EUR), which for a three-dimensional system has the form
$$S=S_r+S_k3(1+\mathrm{ln}\pi )6.434$$
(2)
where $`S_r`$ is the information entropy in position-space of the density distribution $`\rho (\text{r})`$ of a quantum system
$$S_r=\rho (\text{r})\mathrm{ln}\rho (\text{r})๐\text{r}$$
(3)
and $`S_k`$ is the information entropy in momentum-space of the corresponding momentum distribution $`n(\text{k})`$
$$S_k=n(\text{k})\mathrm{ln}n(\text{k})๐\text{k}$$
(4)
The density distributions $`\rho (\text{r})`$ and $`n(\text{k})`$ are normalized to one. Inequality (2), for the information entropy sum in conjugate spaces, is a joint measure of uncertainty of a quantum mechanical distribution, since a highly localized $`\rho (\text{r})`$ is associated with a diffuse $`n(\text{k})`$, leading to low $`S_r`$ and high $`S_k`$ and vice-versa. Expression (2) is an information-theoretical relation stronger than Heisenbergโs. $`S`$ is measured in bits if the base of the logarithm is 2 and nats (natural units of information) if the logarithm is natural.
In previous work we proposed a universal property of $`S`$ for the density distributions of nuclei, electrons in atoms and valence electrons in atomic clusters. This property has the form
$$S=a+b\mathrm{ln}N$$
(5)
where $`N`$ is the number of particles of the system and the parameters $`a,b`$ depend on the system under consideration. It is noted that recently we have obtained the same form for systems of correlated bosons in a trap. This concept was also found to be useful in a different context. Using the formalism in phase-space of Ghosh, Berkowitz and Parr, we found that the larger the information entropy the better the quality of the nuclear density distribution.
In previous work we employed one-body density distributions in the definition of $`S`$. In the present paper we introduce two-body density distributions $`\rho (\text{r}_1,\text{r}_2)`$ and the corresponding two-body momentum distributions $`n(\text{k}_1,\text{k}_2)`$. Our aim is to investigate the properties of $`S`$ at the two-body level for correlated densities. The correlated nucleon systems or the trapped Bose gas, in a good approximation, are studied using the lowest order approximation. Short-range correlations (SRC) are taken into account employing the Jastrow correlation function. Thus it is of interest to examine how $`S_2`$ is affected qualitatively and quantitatively by the same form of correlations in comparison with $`S_1`$, in view of the fact that the quantities $`\rho (\text{r}_1,\text{r}_2)`$ and $`n(\text{k}_1,\text{k}_2)`$ carry more direct information for correlations than the quantities $`\rho (\text{r})`$ and $`n(\text{k})`$ which are only indirectly affected by correlations. The above procedure is repeated for an alternative measure of information i.e. Onicescuโs information energy $`E`$. So far, only the mathematical aspects of this concept have been developed, while the physical aspects have been neglected.
A well known measure of distance of two discrete probability distributions $`p_i^{(1)},p_i^{(2)}`$ is the Kullback-Leibler relative entropy
$$K(p_i^{(1)},p_i^{(2)})=\underset{i}{}p_i^{(1)}\mathrm{ln}\frac{p_i^{(1)}}{p_i^{(2)}}$$
(6)
which for continuous probability distributions $`\rho ^{(1)},\rho ^{(2)}`$ is defined as
$$K=\rho ^{(1)}(x)\mathrm{ln}\frac{\rho ^{(1)}(x)}{\rho ^{(2)}(x)}dx$$
(7)
which can be easily extended for 3-dimensional systems.
Our aim is to calculate the relative entropy (distance) between $`p^{(1)}`$ (correlated) and $`p^{(2)}`$ (uncorrelated) densities both at the one- and the two-body levels in order to assess the influence of SRC (through the correlation parameter $`y`$) on the distance $`K`$. It is noted that this is done for both systems under consideration: nuclei and trapped Bose gases. An alternative definition of distance of two probability distributions was introduced by Rao and Lin, i.e. a symmetrized version of $`K`$, the Jensen-Shannon divergence $`J`$
$$J(p^{(1)},p^{(2)})=H\left(\frac{p^{(1)}+p^{(2)}}{2}\right)\frac{1}{2}H\left(p^{(1)}\right)\frac{1}{2}H\left(p^{(2)}\right)$$
(8)
where $`H(p)=_ip_i\mathrm{ln}p_i`$ stands for Shannonโs entropy. We expect for strong SRC the amount of distinguishability of the correlated from the uncorrelated distributions is larger than the corresponding one with small SRC. We may also see the effect of SRC on the number of trials $`L`$ needed to distinguish $`p^{(1)}`$ and $`p^{(2)}`$ (in the sense described in ).
In addition to the above considerations, we connect $`S_r`$ and $`S_k`$ with fundamental quantities i.e. the root mean square radius and kinetic energy respectively. We also argue on the effect of SRC on EUR and we propose a universal relation for $`S`$, by extending our formalism from the one- and two-body level to the $`N`$-body level, which holds exactly for uncorrelated densities in trapped Bose gas, almost exactly for uncorrelated densities in nuclei (due to the additional exchange term compared to Bose gas) and it is conjectured to hold approximately for correlated densities both in nuclei and Bose gases.
The plan of the present paper is the following. In Sec. 2 we review the formulas of Kullback-Leibler relative entropy entropy $`K`$ and Jensen-Shannon divergence $`J`$, while in Sec. 3 Onicescuโs information energy $`E`$ is described. In Sec. 4 we present the formalism of density distributions used in present work and their applications to Shannonโs and Onicescuโs entropies. In Sec. 5 we introduce SRC in nuclei. In Sec. 6 we apply the formulas of $`K`$ and $`J`$ in correlated distributions. In Sec. 7 we present our numerical results and discussion. Finally, Sec. 8 contains our main conclusions.
## 2 Kullback-Leibler relative entropy and Jensen-Shannon divergence
The Kullback-Leibler relative information entropy $`K`$ for continuous distributions $`\rho _i^{(1)}`$ and $`\rho _i^{(2)}`$ is defined by relation (7). It measures the difference of $`\rho _i^{(1)}`$ form the reference (or apriori) distribution $`\rho _i^{(2)}`$. It satisfies: $`K0`$ for any distributions $`\rho _i^{(1)}`$ and $`\rho _i^{(2)}`$. It is a measure which quantifies the distinguishability (or distance) of $`\rho _i^{(1)}`$ from $`\rho _i^{(2)}`$, employing a well-known concept in standard information theory. In other words it describes how close $`\rho _i^{(1)}`$ is to $`\rho _i^{(2)}`$ by carrying out observations or coin tossing, namely $`L`$ trials (in the sense described in ). We expect for strong SRC the amount of distinguishability of the correlated $`\rho _i^{(1)}`$ and the uncorrelated distributions $`\rho _i^{(2)}`$ is larger than the corresponding one with small SRC.
However, the distance $`K`$ does not satisfy the triangle inequality and in addition is i) not symmetric ii) unbounded and iii) not always well defined. To avoid these difficulties Rao and Lin introduced a symmetrized version of $`K`$ (recently discused in ), the Jensen-Shannon divergence $`J`$ defined by relation (8). $`J`$ is minimum for $`\rho ^{(1)}=\rho ^{(2)}`$ and maximum when $`\rho ^{(1)}`$ and $`\rho ^{(2)}`$ are two distinct distributions, when $`J=\mathrm{ln}2`$. In our case $`J`$ can be easily generalized for continuous density distributions. For $`J`$ minimum the two states represented by $`\rho ^{(1)}`$ and $`\rho ^{(2)}`$ are completely indistinguishable, while for $`J`$ maximum they are completely distinguishable. It is expected that for strong SRC the amount of distinguishability can be further examined by using Wooterโs criterion. Two probability distributions $`\rho ^{(1)}`$ and $`\rho ^{(2)}`$ are distinguishable after $`L`$ trials $`(L\mathrm{})`$ if and only if $`\left(J(\rho ^{(1)},\rho ^{(2)})\right)^{\frac{1}{2}}>\frac{1}{\sqrt{2L}}`$.
The present work is a first step to examine the problem of comparison of probability distributions (for nuclei and bosonic systems) which is an area well developed in statistics, known as information geometry.
## 3 Onicescuโs information energy
Onicescu tried to define a finer measure of dispersion distributions than that of Shannonโs information entropy. Thus, he introduced the concept of information energy $`E`$. For a discrete probability distribution $`(p_1,p_2,\mathrm{},p_k)`$ the information energy $`E`$ is defined by
$$E=\underset{i}{\overset{k}{}}p_i^2$$
(9)
which is extended for a continuous density distribution $`\rho (x)`$ as
$$E=\rho ^2(x)๐x$$
(10)
The meaning of (10) can be seen by the following simple argument: For a Gaussian distribution of mean value $`\mu `$, standard deviation $`\sigma `$ and normalized density
$$\rho (x)=\frac{1}{\sqrt{2\pi }\sigma }\text{exp}\left[\frac{(x\mu )^2}{2\sigma ^2}\right]$$
(11)
relation (10) gives
$$E=\frac{1}{2\pi \sigma ^2}_{\mathrm{}}^{\mathrm{}}\text{exp}\left[\frac{(x\mu )^2}{\sigma ^2}\right]๐x=\frac{1}{2\sigma \sqrt{\pi }}$$
(12)
$`E`$ is maximum if one of the $`p_i`$โs equals 1 and all the others are equal to zero i.e. $`E_{max}=1`$, while $`E`$ is minimum when $`p_1=p_2=\mathrm{}=p_k=\frac{1}{k}`$, hence $`E_{min}=\frac{1}{k}`$ (total disorder). The fact that $`E`$ becomes minimum for equal probabilities (total disorder), by analogy with thermodynamics, it has been called information energy, although it does not have the dimension of energy.
It is seen from (12) that the greater the information energy, the more concentrated is the probability distribution, while the information content decreases. $`E`$ and information content are reciprocal, hence one can define the quantity
$$O=\frac{1}{E}$$
(13)
as a measure of the information content of a quantum system corresponding to Onicescuโs information energy.
Relation (10) is extended for a 3-dimensional spherically symmetric density distribution $`\rho (\text{r})`$
$`E_r={\displaystyle \rho ^2(\text{r})๐\text{r}}`$
$`E_k={\displaystyle n^2(\text{k})๐\text{k}}`$ (14)
in position and momentum space respectively, where $`n(\text{k})`$ is the corresponding density distribution in momentum space.
$`E_r`$ has dimension of inverse volume, while $`E_k`$ of volume. Thus the product $`E_rE_k`$ is dimensionless and can serve as a measure of concentration (or information content) of a quantum system. It is also seen from (12),(13) that $`E`$ increases as $`\sigma `$ decreases (or concentration increases) and the information (or uncertainty) decreases. Thus $`O`$ and $`E`$ are reciprocal. In order to be able to compare $`O`$ with Shannonโs entropy $`S`$, we redifine $`O`$ as
$$O=\frac{1}{E_rE_k}$$
(15)
as a measure of the information content of a quantum system in both position and momentum spaces, inspired by Onicescuโs definition.
## 4 Density Matrices and Information entropies
Let $`\mathrm{\Psi }(๐ซ_1,๐ซ_2,\mathrm{},๐ซ_A)`$ be the wave function that describes the nuclei or the trapped Bose gases and depends on 3A coordinates as well as on spin and isospin (in nuclei). The one-body density matrix is defined in
$$\rho (๐ซ_1,๐ซ_1^{})=\mathrm{\Psi }^{}(๐ซ_1,๐ซ_2,\mathrm{},๐ซ_A)\mathrm{\Psi }(๐ซ_1^{},๐ซ_2,\mathrm{},๐ซ_A)๐๐ซ_2\mathrm{}๐๐ซ_A$$
(16)
while the two-body density matrix by
$$\rho (๐ซ_1,๐ซ_2;๐ซ_1^{},๐ซ_2^{})=\mathrm{\Psi }^{}(๐ซ_1,๐ซ_2,\mathrm{},๐ซ_A)\mathrm{\Psi }(๐ซ_1^{},๐ซ_2^{},\mathrm{},๐ซ_A)๐๐ซ_3\mathrm{}๐๐ซ_A$$
(17)
The above density matrices are related by
$$\rho (๐ซ_1,๐ซ_1^{})=\frac{1}{A1}\rho (๐ซ_1,๐ซ_2;๐ซ_1^{},๐ซ_2)๐๐ซ_2$$
(18)
where the integration is carried out over the radius vectors $`๐ซ_2,\mathrm{},๐ซ_A`$ and summation over spin (or isospin) variables is implied. The corresponding definitions in momentum space are similar. The two-body density distribution $`\rho (๐ซ_1,๐ซ_2)`$ which is a key quantity in the present work, is defined as the diagonal part of the two-body density matrix
$$\rho (๐ซ_1,๐ซ_2)=\rho (๐ซ_1,๐ซ_2;๐ซ_1^{},๐ซ_2^{})_{๐ซ_1^{}=๐ซ_1,๐ซ_2^{}=๐ซ_2}$$
(19)
and expresses the joint probability of finding two nucleons or two atoms at the positions $`๐ซ_1`$ and $`๐ซ_2`$, respectively. The density distribution is given by the diagonal part of the one-body density matrix, that is
$$\rho (๐ซ_1)=\rho (๐ซ_1,๐ซ_1^{})|_{๐ซ_1=๐ซ_1^{}}$$
(20)
or by the equivalent integral
$$\rho (๐ซ_1)=\frac{1}{A1}\rho (๐ซ_1,๐ซ_2)๐๐ซ_2$$
(21)
The two-body momentum distribution $`n(๐ค_1,๐ค_2)`$ is given by a particular Fourier transform of the $`\rho (๐ซ_1,๐ซ_2;๐ซ_1^{},๐ซ_2^{})`$, that is
$$n(๐ค_1,๐ค_2)=\frac{1}{(2\pi )^6}\rho (๐ซ_1,๐ซ_2;๐ซ_1^{},๐ซ_2^{})\mathrm{exp}[i๐ค_1(๐ซ_1๐ซ_1^{})]\mathrm{exp}[i๐ค_2(๐ซ_2๐ซ_2^{})]๐๐ซ_1๐๐ซ_1^{}๐๐ซ_2๐๐ซ_2^{}$$
(22)
In the independent particle model, where the nucleons are considered to move independently in nuclei, the $`\mathrm{\Psi }(๐ซ_1,๐ซ_2,\mathrm{},๐ซ_A)`$ is a Slater determinant. In this case it is easy to show that the two-body density matrix is given by the relation
$`\rho _{SD}(๐ซ_1,๐ซ_2;๐ซ_1^{},๐ซ_2^{})`$ $`=`$ $`{\displaystyle \underset{i,j}{}}\varphi _i(๐ซ_1)\varphi _i(๐ซ_1^{})\varphi _j(๐ซ_2)\varphi _j(๐ซ_2^{}){\displaystyle \underset{i,j}{}}\varphi _i(๐ซ_1)\varphi _j(๐ซ_1^{})\varphi _j(๐ซ_2)\varphi _i(๐ซ_2^{})`$ (23)
$`=`$ $`\rho _{SD}(๐ซ_1,๐ซ_1^{})\rho _{SD}(๐ซ_2,๐ซ_2^{})\rho _{SD}(๐ซ_1,๐ซ_2^{})\rho _{SD}(๐ซ_2,๐ซ_1^{})`$
where $`\varphi _i(๐ซ)`$ is the single-particle wave function normalized to one and
$$\rho _{SD}(๐ซ_1,๐ซ_1^{})=\underset{i}{}\varphi _i(๐ซ_1)\varphi _i(๐ซ_1^{})$$
In Bose gases the many-body ground-state wave function $`\mathrm{\Psi }(๐ซ_1,๐ซ_2,\mathrm{},๐ซ_A)`$ is a product of $`A`$ identical single-particle ground-state wave functions i.e.
$$\mathrm{\Psi }(๐ซ_1,๐ซ_2,\mathrm{},๐ซ_A)=\varphi _0(๐ซ_1)\varphi _0(๐ซ_2)\mathrm{}\varphi _0(๐ซ_A)$$
(24)
where $`\varphi _0(๐ซ_1)`$ is the normalized to one ground-state single-particle wave function describing bosonic atoms. The two-body density matrix in a Bose gas, is given by the relation
$$\rho _0(๐ซ_1,๐ซ_2;๐ซ_1^{},๐ซ_2^{})=\rho _0(๐ซ_1,๐ซ_1^{})\rho _0(๐ซ_2,๐ซ_2^{})$$
(25)
where
$$\rho _0(๐ซ_1,๐ซ_1^{})=\varphi _0(๐ซ_1)\varphi _0(๐ซ_1^{})$$
(26)
We consider that the atoms of the Bose gases are confined in an isotropic HO well, where $`\varphi _0(๐ซ)=(1/(\pi b^2))^{3/4}\mathrm{exp}[r^2/(2b^2)]`$.
As the mean field approach fails to incorporate the interparticle correlation which is necessary for the description of the correlated nuclei or trapped Bose gases, we introduce the repulsive interactions through the Jastrow correlation function $`f(๐ซ_1๐ซ_2)`$ . The correlated nucleon systems or the Bose gases, in a good approximation, can be studied using the lowest order approximation, where the correlated two-body density matrices in nuclei and Bose gases have the following forms respectively
$$\rho (๐ซ_1,๐ซ_2;๐ซ_1^{},๐ซ_2^{})=N\rho _{SD}(๐ซ_1,๐ซ_2;๐ซ_1^{},๐ซ_2^{})f(๐ซ_1๐ซ_2)f(๐ซ_1^{}๐ซ_2^{})$$
(27)
$$\rho (๐ซ_1,๐ซ_2;๐ซ_1^{},๐ซ_2^{})=N\rho _0(๐ซ_1,๐ซ_2;๐ซ_1^{},๐ซ_2^{})f(๐ซ_1๐ซ_2)f(๐ซ_1^{}๐ซ_2^{})$$
(28)
In the present work, in the case of nuclei and trapped Bose gas, the normalization factor $`N`$, is calculated by the normalization condition
$$\rho (\text{r}_1,\text{r}_2)๐\text{r}_1๐\text{r}_2=1$$
(29)
The same holds for $`n(\text{k}_1,\text{k}_2)`$
$$n(\text{k}_1,\text{k}_2)๐\text{k}_1๐\text{k}_2=1$$
(30)
The Jastrow correlation function $`f(\text{r}_1\text{r}_2)`$ both in the case of nuclei and trapped Bose gas is taken to be of the form
$$f(\text{r}_1\text{r}_2)=1\text{exp}[y\frac{(\text{r}_1\text{r}_2)^2}{b^2}]$$
(31)
The uncorrelated case corresponds to $`y\mathrm{}`$, while SRC increase as $`y`$ decreases. The above ansatz has the advantage that it leads to analytical forms for the $`\rho (\text{r}_1,\text{r}_2)`$, $`n(\text{k}_1,\text{k}_2)`$, $`\rho (\text{r})`$ and $`n(\text{k})`$.
The one-body Shannon information entropy both in position- and momentum- space are defined in (3) and (4), where the total sum is
$$S_1=S_{1r}+S_{1k}$$
(32)
The two-body Shannon information entropy both in position- and momentum- space and in total are defined respectively
$$S_{2r}=\rho (๐ซ_1,๐ซ_2)\mathrm{ln}\rho (๐ซ_1,๐ซ_2)๐๐ซ_1๐๐ซ_2$$
(33)
$$S_{2k}=n(๐ค_1,๐ค_2)\mathrm{ln}n(๐ค_1,๐ค_2)๐๐ค_1๐๐ค_2$$
(34)
$$S_2=S_{2r}+S_{2k}$$
(35)
The one-body Onicescu information entropy is already defined in (3) and (15), where the generalization to the two-body information entropy is straightforward and is given by
$$O_2=\frac{1}{E_{2r}E_{2k}}$$
(36)
where
$`E_{2r}`$ $`=`$ $`{\displaystyle \rho ^2(๐ซ_1,๐ซ_2)๐๐ซ_1๐๐ซ_2}`$
$`E_{2k}`$ $`=`$ $`{\displaystyle n^2(๐ค_1,๐ค_2)๐๐ค_1๐๐ค_2}`$ (37)
It is easy to prove that in the case of the uncorrelated trapped Bose gas
$$S_2=2S_1$$
(38)
and
$$O_2=O_1^2$$
(39)
It is worth noting that the above relations hold only approximately in finite nuclei (see Table 1), due to the additional exchange term, originating from the antisymmetry of the nuclear wave function. There is an exception in the case of <sup>4</sup>He, where it holds exactly due to the absence of the exchange term.
## 5 Introduction of SRC in nuclei
We consider that the single particle wave functions, which describe the nucleons is harmonic oscillator type. In order to incorporate the nucleon-nucleon (or atom-atom) correlations, as we mention in the previous section, we apply the lowest order approximation. In this case the two-body density distribution, for <sup>4</sup>He, takes the following form
$$\rho ^{{}_{}{}^{4}He}(๐ซ_1,๐ซ_2)=\rho _{SD}^{{}_{}{}^{4}He}(๐ซ_1,๐ซ_2)+\rho _{cor}^{{}_{}{}^{4}He}(๐ซ_1,๐ซ_2)$$
(40)
The first term of the right-hand side of Eq. (40) which represents the uncorrelated part of the two-body sensity distribution, has the form
$$\rho _{SD}^{{}_{}{}^{4}He}(๐ซ_1,๐ซ_2)=\frac{1}{\pi ^3b^6}\mathrm{exp}[r_{1b}^2]\mathrm{exp}[r_{2b}^2]$$
(41)
and the second term which represents the correlated part of the two-body density distribution, is written
$`\rho _{cor}^{{}_{}{}^{4}He}(๐ซ_1,๐ซ_2)`$ $`=`$ $`{\displaystyle \frac{1}{\pi ^3b^6}}\mathrm{exp}[r_{1b}^2]\mathrm{exp}[r_{2b}^2]`$ (42)
$`\times \left(N\left(1\mathrm{exp}[y(๐ซ_{1b}๐ซ_{2b})^2]\right)^21\right)`$
where $`\text{r}_b=๐ซ/b`$.
In the above expression $`b`$ is the width of the HO potential and $`N`$ is the normalization constant which ensures that $`\rho _{cor}^{{}_{}{}^{4}He}(๐ซ_1,๐ซ_2)๐๐ซ_1๐๐ซ_2=1`$ and has the form
$$N=\left(1\frac{2}{(1+2y)^{3/2}}+\frac{1}{(1+4y)^{3/2}}\right)^1$$
(43)
The density distribution can be written also in the form
$$\rho ^{{}_{}{}^{4}He}(r)=\rho _{SD}^{{}_{}{}^{4}He}(r)+\rho _{cor}^{{}_{}{}^{4}He}(r)$$
(44)
The two-body momentum distribution is given also by the formula
$$n^{{}_{}{}^{4}He}(๐ค_1,๐ค_2)=n_{SD}^{{}_{}{}^{4}He}(๐ค_1,๐ค_2)+n_{cor}^{{}_{}{}^{4}He}(๐ค_1,๐ค_2)$$
(45)
where, as in the case of two-body density distribution, the uncorrelated part has the form
$$n_{SD}^{{}_{}{}^{4}He}(๐ค_1,๐ค_2)=\frac{b^6}{\pi ^3}\mathrm{exp}[k_{1b}^2]\mathrm{exp}[k_{2b}^2]$$
(46)
and the correlated part is written as
$`n_{cor}^{{}_{}{}^{4}He}(๐ค_1,๐ค_2)`$ $`=`$ $`{\displaystyle \frac{b^6}{\pi ^3}}\mathrm{exp}[k_{1b}^2]\mathrm{exp}[k_{2b}^2]`$
$`\times \left(N(1{\displaystyle \frac{1}{(1+4y)^{3/2}}}\mathrm{exp}[{\displaystyle \frac{y}{1+4y}}(๐ค_{1b}๐ค_{2b})^2])^21\right)`$
where $`\text{k}_b=\text{k}b`$.
The momentum distribution is given also by the relation
$$n^{{}_{}{}^{4}He}(k)=n_{SD}^{{}_{}{}^{4}He}(k)+n_{cor}^{{}_{}{}^{4}He}(k)$$
(48)
In the present work, we extend our calculations in nuclei heavier than $`{}_{}{}^{4}\text{He}`$ (<sup>12</sup>C, <sup>16</sup>O and <sup>40</sup>Ca) based on the fact that the high-momentum tails of $`n(k)`$ are almost the same for all nuclei with $`A4`$. Inspired by previous work we suggest a practical method to calculate the one- and two-body density and momentum distributions for nuclei heavier than <sup>4</sup>He. The theoretical scheme of the method combines the mean-field predictions of the two-body density distributions and two-body momentum distributions of various nuclei with their correlated part of <sup>4</sup>He. Specifically, in our treatment we consider the following forms
$$\rho ^A(๐ซ_1,๐ซ_2)=\rho _{SD}^A(๐ซ_1,๐ซ_2)+\rho _{cor}^{{}_{}{}^{4}\text{He}}(๐ซ_1,๐ซ_2)$$
(49)
$$n^A(๐ค_1,๐ค_2)=n_{SD}^A(๐ค_1,๐ค_2)+n_{cor}^{{}_{}{}^{4}\text{He}}(๐ค_1,๐ค_2)$$
(50)
From the above expressions it is obvious that the uncorrelated part of the $`\rho (\text{r}_1,\text{r}_2)`$ and $`n(\text{k}_1,\text{k}_2)`$ originate from the independent particle model for every nucleus separately, where the correlated part in each nucleus is that coming from the nucleus <sup>4</sup>He. The $`\rho (\text{r})`$ and $`n(\text{k})`$ have a similar form.
It should be emphasized that in the uncorrelated case the additional information which is contained in $`\rho (\text{r}_1,\text{r}_2)`$ and $`n(\text{k}_1,\text{k}_2)`$ in nuclei, compared to the trapped Bose gas is the statistical correlations which come from the antisymmetry character of the many-body wave function of nuclei. Moreover, in the correlated case the $`\rho (\text{r}_1,\text{r}_2)`$ and $`n(\text{k}_1,\text{k}_2)`$ contain additional information which originate from the character of the nuleon-nucleon interaction, making our model more realistic and the description more complete. It is of interest to study how the correlations (both statistical and dynamical) affect quantitatively and qualitatively the various kinds of information entropy.
## 6 Application of the Formalism of Relative Entropy and Jensen-Shannon divergence for Correlated Densities
The relative entropy is a measure of distinguishability or distance of two states. It is defined, generalizing (7), by
$$K=\psi ^2(\text{r})\mathrm{ln}\frac{\psi ^2(\text{r})}{\varphi ^2(\text{r})}d\text{r}$$
(51)
In our case $`\psi (\text{r})`$ is the correlated case and $`\varphi (\text{r})`$ the uncorrelated one. Thus
$$K_{1r}=\rho (\text{r})\mathrm{ln}\frac{\rho (\text{r})}{\rho ^{}(\text{r})}d\text{r}$$
(52)
where $`\rho (\text{r})`$ is the correlated one-body density and $`\rho ^{}(\text{r})`$ is the uncorrelated one-body density.
A corresponding formula holds in momentum-space
$$K_{1k}=n(\text{k})\mathrm{ln}\frac{n(\text{k})}{n^{}(\text{k})}d\text{k}$$
(53)
where $`n(\text{k})`$ is the correlated one-body density and $`n^{}(\text{k})`$ is the uncorrelated one.
For the two-body case we have
$$K_{2r}=\rho (\text{r}_1,\text{r}_2)\mathrm{ln}\frac{\rho (\text{r}_1,\text{r}_2)}{\rho ^{}(\text{r}_1,\text{r}_2)}d\text{r}_1d\text{r}_2$$
(54)
where $`\rho (\text{r}_1,\text{r}_2)`$ is the correlated two-body density in position-space and $`\rho ^{}(\text{r}_1,\text{r}_2)`$ is the uncorrelated one.
The generalization to momentum- space is straightforward
$$K_{2k}=n(\text{k}_1,\text{k}_2)\mathrm{ln}\frac{n(\text{k}_1,\text{k}_2)}{n^{}(\text{k}_1,\text{k}_2)}d\text{k}_1d\text{k}_2$$
(55)
where $`n(\text{k}_1,\text{k}_2)`$ is the correlated two-body density in momentum-space and $`n^{}(\text{k}_1,\text{k}_2)`$ is the uncorrelated one.
For the Jensen-Shannon divergence $`J`$ we may write formulas for $`J_1`$ (one-body) and $`J_2`$ (two-body), employing definition (8) and putting the corresponding correlated $`\rho ^{(1)}`$ and uncorrelated $`\rho ^{(2)}`$ distributions in position- and momentum- spaces. We calculate $`K`$ and $`J`$ in position- and momentum- spaces, for nuclei and bosons.
## 7 Numerical results and discussion
For the sake of symmetry and simplicity we put the width of the HO potential $`b=1`$. Actually for $`b=1`$ in the case of uncorrelated case it is easy to see that $`S_{1r}=S_{1k}`$ and also $`S_{2r}=S_{2k}`$ (the same holds for Onicescu entropy), while when $`b1`$ there is a shift of the values of $`S_{1r}`$ and $`S_{1k}`$ by an additive factor $`\mathrm{ln}b^3`$. However, the value of $`b`$ does not affect directly the total information entropy $`S`$ (and also $`O`$). $`S`$ and $`O`$ are just functions of the correlation parameter $`y`$.
In Fig. 1 we present the Shannon information entropy $`S_1`$ using relation (32) and $`S_2`$ using relation (35) in nuclei $`(^{12}\text{C})`$ and trapped Bose gas as functions of the correlation parameter $`\mathrm{ln}(\frac{1}{y})`$. It is seen that $`S_1`$ and $`S_2`$ increase almost linearly with the strength of SRC i.e. $`\mathrm{ln}(\frac{1}{y})`$ in both systems. The relations $`S_2=2S_1`$ and $`O_2=O_1^2`$ hold exactly for the uncorrelated densities in trapped Bose gas, while the above relations are almost exact for the uncorrelated densities in nuclei and in the case of correlated densities both in nuclei and trapped Bose gas. A similar behavior is seen for all nuclei considered in the present work (<sup>4</sup>He, <sup>16</sup>O, <sup>40</sup>Ca).
Values of $`S_1`$, $`S_2`$, $`O_1`$, $`O_2`$ for various nuclei in the uncorrelated case, are shown in Table 1. The relations (38) and (39) are satisfied exactly only in the case of <sup>4</sup>He. However, for the other nuclei, due to the additional exchange term in the nuclear wave function, the relations (38) and (39) hold only approximately (the differences are of order $`0.03\%0.09\%`$ for $`S`$ and $`0.14\%0.96\%`$ for $`O`$).
In Fig. 2 we present the decomposition of $`S`$ in coordinate and momentum spaces, for the sake of comparison i.e. $`S_{1r}`$, $`S_{1k}`$, $`S_{2r}`$, $`S_{2k}`$ for $`{}_{}{}^{16}\text{O}`$ and trapped Bose gas employing (3), (4), (33), (34). The most striking feature concluded from the above Figures is the similar behavior between $`S_{1r}`$ and $`S_{2r}`$ and also $`S_{1k}`$ and $`S_{2k}`$ respectively.
In Fig. 3 we plot the Onicescu information entropy both one-body $`(O_1)`$ and two-body $`(O_2)`$ for nuclei $`(^{12}\text{C},^{40}\text{Ca})`$ and trapped Bose gas (relations (15), (36)). We conclude by noting once again the strong similarities of the behavior between one- and two-body Onicescu entropy.
It is interesting to observe the correlation of the rms radii $`\sqrt{r^2}`$ with $`S_r`$ as well as the corresponding behavior of the mean kinetic energy $`T`$ with $`S_k`$, as functions of the strength of SRC $`\mathrm{ln}(\frac{1}{y})`$ for the $`{}_{}{}^{16}\text{O}`$ nucleus and trapped Bose gas. This is done in Fig. 4 for $`\sqrt{r^2}`$ and Fig. 5 for $`T`$ after apllying the suitable rescaling. The corresponding curves are similar for nuclei and trapped Bose gas.
A well-known concept in information theory is the distance between the probability distributions $`\rho _i^{(1)}`$ and $`\rho ^{(2)}`$, in our case the correlated and the uncorrelated distributions respectively. A measure of distance is the Kullback-Leibler relative entropy $`K`$ defined previously. The correlated and uncorrelated cases are compared for the one-body case $`(K_1)`$ in Fig. 6 and the the two-body case $`(K_2)`$ in Fig. 7 for nuclei $`(^4\text{He},^{16}\text{O},^{40}\text{Ca})`$ and trapped Bose gas, decomposing in position- and momentum-spaces according to (52)-(55). It is seen that $`K_{1r}`$, $`K_{2r}`$ increase as the strength of SRC increases, while $`K_{1k}`$, $`K_{2k}`$ have a maximum at a certain value of $`\mathrm{ln}(\frac{1}{y})`$ depending on the system under consideration.
Calculations are also carried out for the Jensen-Shannon divergence for one-body density distribution ($`J_1`$ entropy) as function of $`\mathrm{ln}(\frac{1}{y})`$ for nuclei and trapped Bose gas, decomposed in position- and momentum- spaces (Fig. 8). We observe again that $`J_1`$ increases with the strength of SRC in position-space, while in most cases in momentum-space there is a maximum for a certain value of $`\mathrm{ln}(\frac{1}{y})`$. It is verified that $`0<J<\mathrm{ln}2`$ as expected theoretically.
It is noted that the dependence of the various kinds of information entropy on the correlation parameter $`\mathrm{ln}(\frac{1}{y})`$ is studied up to the value $`\mathrm{ln}(\frac{1}{y})=0`$ $`(y=1)`$, which is already unrealistic corresponding to strong SRC. In addition, lowest order approximation does not work well beyond that value. In this case three-body terms should be included but this prospect is out of the scope of the present work.
For very strong SRC the momentum distribution $`n(k)`$ exhibits a similar behavior with the mean field $`(y\mathrm{})`$. This is illustrated in Fig. 9, where we present $`n(k)`$ for various values of $`\mathrm{ln}(\frac{1}{y})`$. It is seen that for small and large SRC the tail of $`n(k)`$ disappears. That is why for small and large SRC the relative entropy ($`K_{1k}`$ and $`J_{1k}`$) is small, while in between shows a maximum (Fig. 6, 8). A similar trend of $`n(\text{k}_1,\text{k}_2)`$ for large SRC explains also the maximum of the relative entropy $`K_{2k}`$ in Fig. 7.
## 8 Conclusions and final comment
Our main conclusions are the following
* Increasing the SRC (i.e. the parameter $`\mathrm{ln}(\frac{1}{y})`$) the information entropies $`S`$, $`O`$, $`K`$ and $`J`$ increase. A comparison leads to the conclusion that the correlated systems have larger values of entropies than the uncorrelated ones.
* There is a similar behavior of the entropies as functions of correlations for both systems (nuclei and trapped Bose gas) although they obey different statistics (fermions and bosons).
* There is a correlation of $`\sqrt{r^2}`$ with $`S_r`$ and $`T`$ with $`S_k`$ in the sense that they have the same behavior as a function of the correlation parameter $`\mathrm{ln}(\frac{1}{y})`$. These results can lead us to relate the theoretical quantities $`S_r`$ and $`S_k`$ with experimental ones like charge form factor, charge density distribution, and momentum distribution, radii, etc. A recent paper addressed in that problem.
* The relations $`S_2=2S_1`$ and $`O_2=O_1^2`$ hold exactly for the uncorrelated densities in trapped Bose gas while the above relations are almost exact for the uncorrelated densities and in the case of correlated densities both in nuclei and trapped Bose gas. In previous work we proposed the universal relation $`S_1=S_r+S_k=a+b\mathrm{ln}N`$ where $`N`$ is the number of particles of the system either fermionic (nucleus, atom, atomic cluster) or bosonic (correlated atoms in a trap). Thus in our case
$$S_2=2(a+b\mathrm{ln}N)$$
For 3-body distributions $`\rho (\text{r}_1,\text{r}_2,\text{r}_3)`$ and $`n(\text{k}_1,\text{k}_2,\text{k}_3)`$
$$S_3=3(a+b\mathrm{ln}N)$$
and generalizing for the $`N`$-body distributions $`\rho (\text{r}_1,\text{r}_2,\mathrm{},\text{r}_N)`$ and $`n(\text{k}_1,\text{k}_2,\mathrm{},\text{k}_N)`$
$$S_N=N(a+b\mathrm{ln}N)$$
This is exact for the uncorrelated trapped Bose gas, almost exact in correlated nuclei $`(N=1,2)`$ and it is conjectured that it holds approximately for correlated systems (which has still to be proved for $`N3`$).
* The entropic uncertainty relation (EUR) is
$$S=S_r+S_k6.434$$
It is well-known that the lower bound is attained for a Gaussian distribution (i.e. the case of $`{}_{}{}^{4}\text{He}`$ uncorrelated). In all cases studied in the present work EUR is verified.
A final comment seems appropriate. In general, the calculation of $`\rho (๐ซ_1,๐ซ_2)`$ and $`n(๐ค_1,๐ค_2)`$ is a problem very hard to be solved, especially in the case of nuclei, in the framework of short range correlations. Just a few works are addressed in that problem. In the present work we tried to treat the problem in an approximate but self-consistent way in the sense that the calculations of $`\rho (๐ซ_1,๐ซ_2)`$ and $`n(๐ค_1,๐ค_2)`$ are based in the same $`\rho (๐ซ_1,๐ซ_2;๐ซ_1^{},๐ซ_2^{})`$, which is the generating function of the above quantities. As a consequence the information entropy $`S_2=S_{2r}+S_{2k}`$ is derived also in a self-consistent way and there is a direct link between $`S_{2r}`$ and $`S_{2k}`$, as well as the other kinds of information entropies which are studied in the present work.
## Acknowledgments
The work of Ch. C. Moustakidis was supported by the Greek State Grants Foundation (IKY) under contract (515/2005) while the work of K. Ch. Chatzisavvas by Herakleitos Research Scolarships (21866). One of the authors (Ch. C. M.) would like to thank Prof. Vergados for his hospitality in the University of Ioannina where the earlier part of this work was performed.
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# Optimal slit orientation for long multi-object spectroscopic exposures
## 1 Introduction
The introduction of multi-object spectrographs (MOS) in optical spectroscopy is quite often thought of as just a set of traditional long-slit spectrographs (LS) and some consequences are overlooked. The main difference between MOS and LS spectroscopy is the constraints imposed by the geometry of the instrument. While in LS spectroscopy, the slit orientation can be chosen arbitrarily, this is no longer the case in MOS spectroscopy: slit position angles can not be changed during a set of observations.
As observers tend to concentrate on very faint objects (typical R-band magnitude of 23-24 with 8-10 meter class telescopes), the integration times are getting longer and longer. Sometimes a whole night (8-10 hours or more) is spent on a single mask. These observations require a rethinking of the optimal observation strategy.
Since many MOS instruments lack an atmospheric dispersion corrector (e..g. VIMOS on VLT, DEIMOS on Keck), atmospheric dispersion is a serious problem. Long observations (many hours) span a large range of zenith distance. Thus, a differential refraction of a few arc seconds is quite common.
As most observations require high signal to noise or good spectral resolution, using sufficiently wide slits to compensate for this effect is not acceptable. Very wide slits have a devastating effect on background limited exposures (since the sky background grows linearly with slit width, while the object signal grows much slower) and a wide slit also blurs the spectra.
For short exposures, one can observe โclose to the parallactic angleโ, i.e. align the slit with the atmospheric dispersion direction. This way the photons from the object enter the slit, the dispersion only introduces an additional tilt in the resulting spectra, which is easy to correct for for most applications. If the goal of the observation is to extract spatial information, too, then extra care is required to correct for this effect.
In the case of longer observations, the direction of dispersion projected on the sky, i.e. the parallactic angle, varies in time. For single object observations (i.e. long slit spectroscopy), one can compensate by periodically realigning the slit. For MOS exposures, this is not possible. For masks, a single slit orientation *must* be chosen for the whole exposure.
In this paper, I examine how the effective slit loss can be estimated and I demonstrate the effect using our sample observation of the Lockman Hole using the Keck telescope. I also describe how to use our Web-based service to find the best strategy for a particular observation.
I will start with the current best determination of the atmospheric refraction. I also work out a simplified formula that is sufficient in many applications. I simulate different observational strategies and show the effect of atmospheric refraction on the efficiency of observations.
## 2 Atmospheric dispersion and large telescopes
The importance of atmospheric refraction in spectroscopic observations was emphasized by Filippenko (filippenko1982 (1982)). The paper discusses the optimal strategy (and the effect of non optimal strategies) for *short* long-slit spectroscopic observations. Even though the paper uses a formula to calculate the index of refraction that became obsolete, it is still the strategy to be followed for *short* integrations. The refined formula to calculate the refraction introduces only negligible changes. On the other hand, the paper does not discuss the optimal strategy for *long* integrations.
Cohen and Cromer (cohen1988 (1988)) calculates the magnitude of differential refraction for the Keck and the Norris spectrographs for realistic observing scenarios. The paper determines the limits beyond which the atmospheric dispersion degrades the data, but does not discuss how to optimize observations that go beyond these limits.
Donnelly et. al (donnelly1989 (1989)) discusses optimal observational strategies for *fiber* spectrographs for *long* exposures. Unfortunately these results can not be directly applied to slit spectrographs: In many (but not all) projects, the slit orientation can be chosen arbitrarily, thus, there is an extra degree of freedom to minimize the effect of atmospheric refraction. This is not possible for fiber spectrographs, thus, this is not discussed in this paper.
In Table 1 I review all current and known future 8m telescopes and optical MOS spectrographs. There are only 2 instruments without an ADC in operation or in planning:
DEIMOS on Keck is heavily red optimized, thus, atmospheric dispersion is not a significant issue for many projects. The DEIMOS Slitmask design page (http://www.ucolick.org/$``$phillips/deimos\_ref/masks.html) provides preliminary tools to evaluate the effect of atmospheric dispersion on slit loss for short integrations. No guidelines are provided for long integrations.
The VIMOS manual (vimos2005 (2005)) discusses the effect of atmospheric dispersion in MOS mode. They arrive at the conclusion that the only generic way to minimize slit losses is to orient the slits North-South and observe within $`\pm `$2 hours of the meridian. Even though the detailed study of the effect of atmospheric dispersion on VIMOS (Cuby et al. cuby1996 (1998)) makes no explicit statement about deviating from these constraints, a casual reading of the manual by an unexperienced observer may leave the impression that this slit orientation is the *only* valid strategy.
This conclusion is clearly valid for the sample observation used in the manual: observing between airmasses of 1.7 and 1.4 in the UV/blue. On the other hand, for some observing projects, e.g. limiting UV/blue spectroscopy to low airmass, the advantage of N-S orientation diminishes. Thus, an additional freedom is available in *some* cases to maximize scientific return of the observations.
## 3 Sample observation
Throughout this paper, I will use a sample observation of the Lockman Hole ($`\delta `$=+57:35:25.0, J2000) using the Keck telescope (latitude of +19:46:36). I assume multislit spectroscopy with 1.0โ slit width and 1.2โ seeing. I assume an ambient temperature of 2.5 degrees C, ambient pressure of 61.5 kPa (615 mbar) and a relative humidity of 40%.
The observations consist of 1 night long integration on a mask, which for this field implies an hour angle range between 18 hours and 3 hours (the asymmetry is due to the mechanical constraints of the telescope), which covers an airmass range of 1.27 (hour angle of 0) to 3.5 (hour angle of 18 hours). I will concentrate on the DEIMOS multiobject spectrograph, which is red-optimized. Thus, I will concentrate on the 4500โฆ9500 ร
wavelength range.
In the calculation I assume that the seeing does *not* depend on the wavelength and I also assume that the alignment and guiding is done in the R-band (approximately 7000 ร
).
## 4 Atmospheric dispersion
The most up-to-date atmospheric dispersion determination, the Ciddor formula (ciddor1996 (1996)) is reviewed in the Appendix. The most important formulas are:
The *differential* refraction (as a function of wavelength, relative to the alignment/guiding effective wavelength, $`\lambda _0`$) in radians is:
$$\mathrm{\Delta }R(\lambda )R(\lambda )R(\lambda _0)=(n(\lambda )n_0)\mathrm{tan}z_a$$
(1)
In figure 1 I plot $`\mathrm{\Delta }R`$ as a function of wavelength at different wavelengths for different airmasses.
The index of refraction, $`n_{as}`$ of standard air is
$$10^8(n_{as}1)=\frac{5792105\mu m^2}{238.0185\mu m^2\sigma ^2}+\frac{167917\mu m^2}{57.362\mu m^2\sigma ^2}$$
(2)
where $`\sigma `$ is the wave number (reciprocal of the *vacuum* wavelength) in inverse micrometers.
If we are only interested in differential refraction, we can write a simpler formula that is sufficiently accurate for many applications:
$$\mathrm{\Delta }R(\lambda )\frac{pT_0}{p_0T}(n_{as}(\lambda )n_{as}(\lambda _0))\mathrm{tan}z_a$$
(3)
The error introduced by this approximation (as well as the error introduced by the old Edlรฉn formula) is shown is figure 2. As one can see, even the *absolute* refraction is well reproduced by this simpler formula, while the old Edlรฉn formula is significantly different. For the calculation of the *differential* refraction, there is no practical difference between the formulas in the wavelength range considered.
## 5 Slit loss
Now that we have the atmospheric dispersion, we can also calculate the slit loss. I assume a point source that has a surface brightness profile of
$$\mu (r)=\frac{1}{2\pi \sigma ^2}e^{r^2/2\sigma ^2}$$
(4)
If we assume a slit of width $`2a`$, that is sufficiently long, the fraction of light entering the slit from an object that is displaced by $`x_0`$ *perpendicular* to the slit (a displacement parallel to the slit does not affect the amount of light entering the slit) is
$$I(x_0)=\frac{1}{\sigma \sqrt{2\pi }}\underset{a+x_0}{\overset{a+x_0}{}}e^{x^2/2\sigma ^2}๐x$$
(5)
The perpendicular displacement, $`x_0`$, depends on the differential refraction, $`\mathrm{\Delta }R`$, and the angle between the slit and and parallactic angle. In figure 3 I consider two configurations, an East-West oriented slit (i.e. the slits are perpendicular to the parallactic angle at the meridian) and a North-South orientation (slits are parallel to the parallactic angle at the meridian).
As we can see, the N-S slit orientation results in a low slit loss at hour angle of 0 (minimal airmass) that does not depend on wavelength (as we are observing close to the parallactic angle). On the other hand, as we are moving away from the optimal configuration, the situation deteriorates rapidly. This is due to the fact that the atmospheric dispersion increases as the airmass increases *and* the slit orientation is moving away from the ideal, parallactic angle โ both effects increase the slit loss.
In the alternative configuration, i.e. East-West slit orientation, the slit loss is never optimal. Even at low airmass, a significant fraction of the light is lost (e.g. slit loss is 42% at 5000 ร
, instead of 33%), *but* the slit loss does not deteriorate so quickly. This is due to the fact that as the airmass increases, the dispersion increases, but the slit is getting closer to the parallactic angle, thus the *projected* dispersion is not increasing so rapidly. In fact, in our particular configuration (Lockman Hole and Keck), the slit loss is actually smaller at high airmass (42%, 38% and 39% at 5000 ร
at the airmass of 1.27, 1.56 and 3.50, respectively).
To evaluate the overall effect of the slit orientation on the signal level achievable, we can also calculate the โaverageโ slit loss of a long exposure. This is shown in figure 4.
In figure 5 I show the effect of sky position angle on average slit loss at 4500ร
, using my example observations and a southern field (Chandra Deep Field South). As we can see, for short exposures around meridian passage, the N-S orientation is optimal. As we go to longer and longer integration times, the effect of slit orientation becomes smaller and smaller, while at extremely long integration times, the E-W orientation becomes ideal for the Lockman Hole field. For asymmetric cases (for example our actual observations between hour angles 18 and 3) the optimal slit orientation is neither E-W nor N-S.
In contrast, for the southern field, the optimal slit orientation remains North-South. This is due to the fact that for the CDFS field the airmass is never below 1.48 from Keck, thus, the differential refraction is comparable to the slit width even at meridian passing.
## 6 Field differential refraction
All results presented so far only hold for objects in the center of the field of view of the instrument. For wide-field spectrographs, an additional effect is important: the zenith distance is not constant across the field of view, thus the differential refraction is not constant, either. This is an achromatic effect, which we can calculate using equation 10:
$$\delta R(n1)\frac{d\mathrm{tan}z_a}{dz_a}\delta z_a=(n1)\mathrm{sec}^2z_a\delta z_a$$
(6)
where $`\delta R`$ is the variation in the differential refraction across the field and $`\delta z_a`$ is the field of view. As we have seen, $`(n1)`$ is on the order of $`3\times 10^4`$, $`\mathrm{sec}z_a`$ is on the order of 1, thus a 1000 arc second field of view (nearly 17 arc minutes) introduces a *field* differential refraction on the order of 0.3 arc second. It is important to point out that this effect makes it imperative to realign the masks periodically unless guiding is near the center of the field or the telescope compensates for off-axis guiding.
## 7 Web interface
I make our code used in our calculations available to the community in both source code and web application form at http://www.xray.mpe.mpg.de/$``$szgyula/slitloss/. The program uses the full Ciddor formula to evaluate the slit loss, but for comparison all three formulas are available to calculate the differential refraction.
## 8 Conclusions
I demonstrated that choosing the optimal slit orientation for multiobject spectroscopy, using long exposure times, requires care and should be evaluated individually for each project. For every field and expected duration, one has to find a balance.
I reviewed the most recent determination of atmospheric dispersion. For typical cases, I found that a simplified version of the most up-to-date Ciddor formula can be used, due to the fact that alignment/guiding removes the effect of dispersion in zeroth order. The simplified formula only depends on pressure and temperature. The effect of relative humidity and CO<sub>2</sub> concentration is very small. Furthermore, the differential refraction follows very simple scaling rules, i.e. it scales linearly with pressure and the inverse of temperature (in Kelvins).
For short exposures, the optimal strategy is, as expected, still to orient the slits with the parallactic angle. On the other hand, for longer exposures, this is not always the right strategy. There are two effects to consider, the increasing differential refraction and the changing angle between the slits and the parallactic angle. Depending on the configuration, these effects can work against each other, thus resulting in a long, relatively stable observation that is never optimal or these can amplify each other, thus resulting in an optimal short observation that deteriorate very fast.
It is also important to point out that alignment/guiding is crucial. One has to select the effective wavelength of these to maximize the science output. It is absolutely worthwhile to spend a few extra minutes every few hours using some standard filters instead of using no filters at all, especially with alignment stars with unknown spectral types. This latter approach runs the risk of using very blue stars for alignment, thus the alignment will only be optimal in the blue, where one may not be observing.
Naturally, in the long run, the use of atmospheric dispersion correctors should be considered. As I have shown, these can improve the throughput by as much as a factor of two for instruments operating in the blue. The cost of these from an observational point of view is small (few photons are lost in the extra optical elements) as only two very weak prisms are sufficient in most cases to produce a โtunableโ prism to compensate for differential refraction in first order. This can solve the problem of differential refraction, but does not eliminate the slit loss completely. Finite slit widths will always โcutโ the object signal. As the seeing can be wavelength dependent (this effect was completely ignored in this paper), so can the slit loss be wavelength dependent. Thus, accurate spectrophotometry still requires very wide slits, and consequently, very low spectral resolution and a significantly degraded signal to noise ratio.
Finally, I provide a web based service to the community and we also release the software developed.
###### Acknowledgements.
Part of this work was supported by the German *Deutsche Forschungsgemeinschaft, DFG* project number Ts 17/2โ1.
## Appendix A Atmospheric dispersion
Atmospheric dispersion (i.e. the apparent displacement of object), $`R`$, is defined as
$$R=z_tz_a$$
(7)
where $`z_t`$ and $`z_a`$ are the true and apparent zenith distances, respectively.
Assuming that the index of refraction depends only on height, using Snellโs law, we can write that
$$n(h)\mathrm{sin}(z(h))const.$$
(8)
where $`h`$ is the height, $`z(h)`$ is the apparent zenith distance at height $`h`$.
$$n\mathrm{sin}z_a=\mathrm{sin}z_t$$
(9)
Thus, the apparent zenith distance at the telescope only depends on the index of refraction at the observatory. Assuming that $`R`$ is small, i.e. $`\mathrm{sin}RR`$ and $`\mathrm{cos}R1`$, we can write
$$R(n1)\mathrm{tan}z_a$$
(10)
As no telescope points accurately enough (i.e. with less than 0.1 arc second accuracy required by slit based spectroscopy), all observations start with an โalignmentโ. This step guarantees that all objects are centered on the slit *at a particular wavelength*, $`\lambda _0`$ (determined by the filter used for the alignment). During the exposure, the guiding subsystem and periodic realignments maintain this condition. As the guider typically does not operate at the same wavelength, special care is required to compensate for atmospheric refraction in the guider system. Even with an ideal guider, the mask gets misaligned due to other effects, e.g. the open loop instrument rotator (especially on modern, altitude-azimuth mounted telescopes). Thus, a periodic realignment is mandatory. For simplicity, I will assume that the telescope is equiped with an ideal alignment and guiding system, thus the objects are always centered on the slit at wavelength $`\lambda _0`$. In section 6 I discuss why this can not hold for wide field of view instruments.
As a consequence, the *absolute* magnitude of atmospheric refraction is irrelevant, since alignment/guiding automatically compensates for it. The relevant quantity, the *differential* refraction (as a function of wavelength, relative to the alignment/guiding effective wavelength, $`\lambda _0`$) in radians is:
$$\mathrm{\Delta }R(\lambda )R(\lambda )R(\lambda _0)=(n(\lambda )n_0)\mathrm{tan}z_a$$
(11)
The calculation of the refractive index is a crucial part of our calculation. Unfortunately, there are still old formulas in use, most notably the Cauchy formula and the old and new Edlรฉn formulas that are at least 50 years old. These formulas are known to be inaccurate, but they still crop up in the literature and astronomical applications. The current best formula is the Ciddor (ciddor1996 (1996)) formula, presented below.
The index of refraction, $`n_{as}`$ of standard air, i.e. dry air at 15C temperature, using the International Temperature Scale of 1990 (Saunders, saunders1990 (1990)) , 101325 Pa pressure and 450 ppm (part per million) CO<sub>2</sub> concentration, is
$$10^8(n_{as}1)=\frac{5792105\mu m^2}{238.0185\mu m^2\sigma ^2}+\frac{167917\mu m^2}{57.362\mu m^2\sigma ^2}$$
(12)
where $`\sigma `$ is the wave number (reciprocal of the *vacuum* wavelength) in inverse micrometers. In the range of 3500ร
โฆ24000ร
, $`10^8(n_{as}1)`$ is in the range of 28612โฆ27289.
If the CO<sub>2</sub> concentration is $`x_c`$ ppm instead of 450 ppm, the index of refraction, $`n_{axs}`$, is
$$n_{axs}1=(n_{as}1)\left(1+0.534\times 10^6\left(x_c450\right)\right).$$
(13)
This formula is accurate to $`10^8`$ for the refractive index up to 600 ppm CO<sub>2</sub> concentrations in the range of 360-2500 nm. In this range, the effect of CO<sub>2</sub> variation on $`n_{axs}`$ is on the order of $`10^7`$.
For water vapor at the โstandard conditionsโ, i.e. at 20 C and 1333 Pa, the index of refraction, $`n_{ws}`$, is
$`10^8(n_{ws}1)=1.022(295.235\mu m^2+2.6422\mu m^2\sigma ^2`$
$`0.032380\mu m^4\sigma ^4+0.004028\mu m^6\sigma ^6)`$ (14)
The formula is accurate to $`2\times 10^7`$ in the range of 350-1200 nm. In the optical/near-IR range $`\sigma `$ is between 0.5 and 3. Thus, the value of $`10^8(n_{ws}1)`$ is 324โฆ302 in this wavelength range.
The saturation vapor pressure of water vapor, $`p_{vs}`$, at temperature $`T`$ (in Kelvins), over liquid water is
$`p_{vs}=\mathrm{exp}(1.2378847\times 10^5K^2T^2`$
$`1.9121316\times 10^2K^1T+33.937110476343.1645K/T)`$ (15)
Considering a temperature range of -20 Cโฆ40 C, the saturated vapor pressure is 0.1โฆ7.4 kPa.
The enhancement factor of water vapor in air is
$$f=1.00062+3.14\times 10^8Pa^1p+5.6\times 10^7^{}C^2t^2$$
(16)
where $`p`$ is the pressure and $`t=T273.15K`$. The deviation of $`f`$ from 1 is at most $`4\times 10^4`$.
The molar fraction of water vapor in moist air is
$$x_w=\frac{fhp_{vs}}{p}$$
(17)
where $`h`$ is the fractional humidity (between 0 and 1). The range of $`x_w`$ is 0โฆ0.25, but the high value (0.25) is assuming unrealistic conditions, i.e. 40 C temperature, 100% humidity and Mt. Everest type ambient pressure (30 kPa). In most cases, accepting a range of 0โฆ0.05 (20 C, 60 kPa) is more realistic.
The compressibility of the moist air, $`Z`$ is
$`Z=`$ (18)
$`1{\displaystyle \frac{p}{T}}[1.58123\times 10^6KPa^12.9331\times 10^8Pa^1t+`$
$`1.1043\times 10^{10}K^1Pa^1t^2+`$
$`\left(5.707\times 10^6KPa^12.051\times 10^8Pa^1t\right)x_w+`$
$`(1.9898\times 10^4KPa^12.376\times 10^6Pa^1t)x_w^2]+`$
$`\left({\displaystyle \frac{p}{t}}\right)^2\left(1.83\times 10^{11}K^2Pa^20.765\times 10^8K^2Pa^2x_w^2\right)`$
where $`p`$ is pressure (in Pascals), T is temperature (in Kelvins) and $`t=T273.15K`$. Considering realistic pressures (30โฆ100 kPa), temperatures (-20โฆ40 C) and water vapor molar fractions ($`x_w<0.1`$), the compressibility is very close to one: $`|Z1|<2\times 10^3`$ (using a very conservative upper limit estimate).
At standard conditions ($`p=101325Pa`$, t=15 C, dry air), the standard compressibility is
$$Z_00.9995922115$$
(19)
At the saturated water vapor conditions, i.e. $`p=1333Pa`$ and $`t=20^{}C`$, the compressibility is
$$Z_10.9999952769$$
(20)
The molar mass of air is
$$M_a=10^3\left(28.9635+12.011\times 10^6\left(x_c400\right)\right)$$
(21)
in kg/mol units ($`x_c`$ is the CO<sub>2</sub> concentration in ppm, as used above).
The density of air, $`\rho `$ (in kg/m<sup>3</sup> units) is
$$\rho =\frac{pM_a}{ZRT}\left(1x_w\left(1\frac{M_w}{M_a}\right)\right)$$
(22)
where $`R=8.314510Jmol^1K^1`$, the gas constant, $`M_w=0.018015kg/mol`$, the molar mass of water vapor.
At standard conditions ($`p_0=101325Pa`$ and $`t_0=15^{}`$C) the density of dry air ($`x_w=0`$) only depends on the CO<sub>2</sub> concentration
$$\rho _{axs}=\frac{p_0M_a}{Z_0RT_0}=\frac{28.9635\times 10^33p_0}{Z_0RT_0}\left(1+\frac{12.011}{28.9635}\times \frac{x_c400}{10^6}\right)$$
(23)
At water vapor standard conditions ($`p_1=`$1333 Pa, $`t_1=`$20 C, $`w_s=1`$), the saturated water vapor density is
$$\rho _{ws}=\frac{p_1M_w}{Z_1RT_1}0.00985235$$
(24)
For actual conditions, the air density of the air component is
$$\rho _a=\frac{pM_a}{ZRT}\left(1x_w\right)$$
(25)
and
$$\frac{\rho _a}{\rho _{axs}}=\frac{Z_0}{Z}\frac{p}{p_0}\frac{T_0}{T}(1x_w)$$
(26)
The water vapor component is
$$\rho _v=\frac{pM_wx_w}{ZRT}$$
(27)
and
$$\frac{\rho _v}{\rho _{ws}}=\frac{p}{p_1}\frac{T_1}{T}\frac{Z_1}{Z}x_w=\frac{T_1}{T}\frac{Z_1}{Z}\frac{p_{vs}}{p_1}fh$$
(28)
Considering that $`Z`$ and $`f`$ are very close to one, assuming realistic ranges for the parameters, we can place a very conservative upper limit on the vapor density to standard saturated water vapor density ratio: $`\rho _v/\rho _{ws}<`$6.5.
Finally, the refractive index is
$$n1=\frac{\rho _a}{\rho _{axs}}\left(n_{axs}1\right)+\frac{\rho _v}{\rho _{ws}}\left(n_{ws}1\right)$$
(29)
As I have shown above, $`n_{ws}`$ varies by at most $`2\times 10^7`$ as a function of wavelength in the optical/near-IR range and $`\rho _v/\rho _{ws}<`$6.5. Thus, the second, water related term can change by at most $`1.3\times 10^6`$, thus the differential refraction can change by at most this much. This limits the effect of water in the atmosphere to 0.3 arc seconds in the most extreme case: alignment/guiding in the K-band, observing in the UV, close to 100% humidity. In a realistic case, the effect is much smaller so the second term can be ignored most of the time. The first term in the equation *does depend* on humidity through the air density, but in realistic cases the only water related term, $`1x_w`$ varies by a few percent, thus the effect is very small.
As I have shown, the variation in $`n_{axs}`$ introduced by CO<sub>2</sub> concentration variation is less than $`10^7`$. Thus, the differential refraction variation introduced is 0.03 arc seconds or less. For most applications, this is negligible.
If we are only interested in differential refraction and are considering these approximations, we can write a much simpler formula that is sufficiently accurate for many applications:
$$\mathrm{\Delta }R(\lambda )\frac{pT_0}{p_0T}(n_{as}(\lambda )n_{as}(\lambda _0))\mathrm{tan}z_a$$
(30)
In Table 2 the importance of atmospheric condition variations is shown. The effect of relative humidity and CO<sub>2</sub> concentration is negligible in most observations.
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# The noncommutative degenerate electron gas
## References
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# Critical regularity for elliptic equations from Littlewood-Paley theory
## 1 Introduction
In this paper we prove the following local regularity result for (complex pseudodifferential) elliptic equations. Let $`B_R`$ denote a ball of radius $`R`$ and center $`0`$ in $`๐^n`$, $`n=2`$, $`3`$, $`\mathrm{}`$.
###### Theorem 1.1
Let $`uH^s(B_1)`$ solve
(1.1)
$$\mathrm{\Delta }^\alpha u+Vu=0\mathrm{๐๐}B_1$$
with $`VL^{n/2\alpha }(B_1)`$, $`0<2\alpha <n`$, and $`0<2s<n`$. Assume also that
(1.2)
$$2\alpha \frac{n}{2}<s<2\alpha .$$
Then
(1.3)
$$uH^{s+\epsilon }(B_{1/2})$$
for some $`\epsilon >0`$, $`\epsilon =\epsilon (n,\alpha ,s)`$.
By $`H^s`$ we denote the Sobolev Hilbert space of order $`s๐^1`$. We write $`uH^s(B_R)`$ if $`u`$ is a distribution on $`B_R`$ such that
$$u=\stackrel{~}{u}|_{B_R}\mathrm{for}\mathrm{some}\stackrel{~}{u}H^s(๐^n).$$
For any $`\alpha >0`$ the operator $`\mathrm{\Delta }^\alpha `$ is a nonlocal pseudodifferential operator defined on $`H^s(๐^n)`$,
$$\mathrm{\Delta }^\alpha :H^s(๐^n)H^{s2\alpha }(๐^n),$$
for any $`s๐^n`$.
Thus the theorem states that, if a distribution $`uH^s(B_1)`$ has an extension $`\stackrel{~}{u}`$ with $`\mathrm{\Delta }^\alpha \stackrel{~}{u}H^{s2\alpha }(๐^n)`$ satisfying (1.1) in $`B_1`$, then we can find a distribution from $`H^{s+\epsilon }(๐^n)`$, $`\epsilon >0`$, coinciding with $`u`$ in $`B_{1/2}`$. The particular case of integer $`s=\alpha `$ arises in the calculus of variations. In this case our theorem says that equation (1.1) improves the regularity of $`H^\alpha `$-solutions if $`0<2\alpha <n`$.
We will prove the theorem using only two simple facts from Littlewood-Paley theory, namely, the Plansherel isometry and the Bernstein inequality. Regularity for (1.1) for the end-point relations between parameters $`n`$, $`\alpha `$, and $`s`$ cannot apparently be established with such simple tools. The only non-obvious assumption on the parameters in Theorem 1.1 is the lower bound for $`s`$ in (1.2). Together with the Sobolev embedding it garantees that the product $`Vu`$ is defined as a distribution.
Theorem 1.1 can be derived from the results of Y.Y.Li , see Theorem 1.3 there. Equation (1.1) is treated in as an integral equation in the physical space and the frequency space is not used there at all. Technique developed in the present paper does not depend on the structure of the fundamenatal solution of $`\mathrm{\Delta }^\alpha `$. In particular, it allows to establish the local regularity for more general pseudodifferential equations on smooth manifolds .
The function $`V`$ is integrable with the critical power in the theorem meaning the following: if $`VL^p(B_1)`$ with $`p<n/2\alpha `$ then in general (1.3) does not hold for any $`\epsilon >0`$ as the family of examples below shows. If $`VL^p(B_1)`$ with $`p>n/2\alpha `$ then the improved regularity (1.3) is easy to prove. Indeed, in this case Sobolev and Holder inequalities imply at once that
(1.4)
$$\mathrm{\Delta }^\alpha u=f\mathrm{in}B_1\mathrm{with}fL^p(B_1),\frac{1}{p}<\frac{2\alpha }{n}+\frac{1}{2}\frac{s}{n}.$$
Now (1.3) is the straightforward consequence of the Calderon-Zygmund estimate and Sobolev inequality.
The main purpose and application of Theorem 1.1 is deriving the full regularity for quasilinear (complex pseudodifferential) elliptic equations with the critical growth nonlinearity. For such application any $`\epsilon >0`$ in (1.3) works equally well. It is for this reason that we do not care about the sharp value of $`\epsilon (n,\alpha ,s)`$ in Theorem 1.1. For example, consider a weak solution $`uH^\alpha `$, $`0<2\alpha <n`$, of the equation
$$\mathrm{\Delta }^\alpha u+g(x,u)=0\mathrm{in}B_1.$$
Assume that $`g`$ is a smooth, possibly complex valued function of the critical growth:
$$|g(x,t)|C\left(1+|t|^{(n+2\alpha )/(n2\alpha )}\right)\mathrm{for}\mathrm{all}xB_1,t๐^1.$$
We can write
$`g(x,u)`$ $`=`$ $`{\displaystyle \frac{g(x,u)}{1+|u|}}+\left({\displaystyle \frac{g(x,u)}{1+|u|}}{\displaystyle \frac{|u|}{u}}\right)u`$
$`=`$ $`f+Vu`$
with $`f`$ as in (1.4) and $`VL^{n/2\alpha }`$. Now the application of Theorem 1.1 combined with Calderon-Zygmund and Sobolev inequalities improve the integrability of $`u`$. Then Schauder estimates imply that $`uC^{\mathrm{}}(B_{1/2})`$.
This way of proving the regularity for the critical semilinear equations was suggested for $`\alpha =1`$ by Brezis and Kato , see also Appendix B in . These authors improved integrability of $`u`$ using Moserโs iteration technique. Earlier Trudinger (also in the case $`\alpha =1`$) had already used Moserโs iterations to prove the full regularity for the nonlinear problem directly. The case of an integer $`\alpha >1`$ has attracted recent attention in , due to its applications in conformal geometry. In a related paper Y. Y. Li proved the full regularity for the equation
(1.5)
$$\mathrm{\Delta }^\alpha u+u^{(n+2\alpha )/(n2\alpha )}=0,u>0,0<2\alpha <n.$$
The main goal in was to establish Liouville-type theorems for (1.5) in $`๐^n`$ using the moving spheres method. Earlier Liouville theorems were proved for (1.5) in , , using the moving plane method. The Littlewood-Paley approach was used in and in to give proofs of regularity of Holder-continuous harmonic maps and harmonic maps from surfaces into spheres respectively.
Elliptic equations with supercritical nonlinearity do not improve the regularity of solutions. For example, for $`\alpha =1`$ and any $`p>(n+2)/(n2)`$, $`n3`$, the function
$$u(x)=\frac{A}{|x|^a},a=2/(p1),A=(a(n+a2))^{1/(p1)},$$
satisfies
$$\mathrm{\Delta }u+u^p=0\mathrm{in}B_1,uH^1(B_1).$$
However, $`u`$ is not smooth in $`B_{1/2}`$. Pohozaev in investigated local regularity for supercritical semilinear problems, and established some sharp low regularity results.
Acknowledgments. This work was done when the author was visiting the Australian National University in 2003 by the invitation of Neil Trudinger and Xu-Jia Wang. The author also wishes to thank O. V. Besov, M. L. Goldman, S. I. Pohozaev, and other participants of the Fall, 2004 seminar on function spaces at the Steklov Institute for their comments.
## 2 Proof of Theorem 1.1
Let $`\{\widehat{\phi }_j\}_{j=\mathrm{}}^+\mathrm{}`$ be the standard smooth partition of unity in the Littlewood-Paley theory , , , . Thus $`\widehat{\phi }_j=\widehat{\phi }(/2^j)`$ is supported in, say, the ring
$$\{\xi :2^j3/5|\xi |2^j5/3\}\left(B_{2^{j+1}}B_{2^{j1}}\right).$$
Let $`P_j`$ denote the Littlewood-Paley projection,
$$(P_jf)^{}=\widehat{\phi }_j\widehat{f},f๐ฎ^{}.$$
We also set
$$P_{a<<b}=\underset{j=a+1}{\overset{b1}{}}P_j.$$
Distributions with the localised Fourier transform enjoy the important Bernstein inequality: for $`f๐ฎ^{}`$ and $`1pq\mathrm{}`$
$$f_q2^{nj((1/p)(1/q))}f_p\mathrm{provided}\mathrm{supp}\widehat{f}B_{2^j}.$$
For $`s๐^1`$ the Sobolev space $`H^s(๐^n)`$ consists of distributions with the finite norm
$$f_{H^s}=P_0f_2+\left(\underset{j=1}{\overset{\mathrm{}}{}}2^{2js}P_jf_2^2\right)^{1/2}.$$
The Plansherel isometry implies that for $`s=1`$, $`2`$, $`\mathrm{}`$ the space $`H^s(๐^n)`$ consists of distributions with all derivatives up to the order $`s`$ lying in $`L^2(๐^n)`$.
Proof. (of Theorem 1.1) 1. First, we localise the problem. Take a cutoff function $`\eta _\rho `$,
$$\eta _\rho =1\mathrm{in}B_\rho ,\eta _\rho =0\mathrm{outside}B_{2\rho }.$$
The commutator of the multiplication by $`\eta _\rho `$ and $`\mathrm{\Delta }^\alpha `$ is a pseudodifferential operator of order $`2\alpha 1`$. For interger $`\alpha `$ this is just the Leibnitz formula for the derivative of the product. Hence, for some $`FH^{s2\alpha +1}(๐^n)`$ we obtain
(2.1) $`\mathrm{\Delta }^\alpha (\eta _\rho u)`$ $`=`$ $`\mathrm{\Delta }^\alpha (\eta _\rho \stackrel{~}{u})`$
$`=`$ $`\eta _\rho \mathrm{\Delta }^\alpha \stackrel{~}{u}+F`$
$`=`$ $`(\eta _{2\rho }V)(\eta _\rho u)+F\mathrm{in}๐^n.`$
To economize on notations denote $`u\eta _\rho `$ by $`u`$ and $`V\eta _{2\rho }`$ by $`V`$. Then in (2.1) we have $`uH^s(๐^n)`$, $`\mathrm{supp}(u)B_{2\rho }`$. Moreover, the $`L^{n/2\alpha }`$-norm of $`V`$ is small when $`\rho `$ is small. In the proof we, by making this norm small enough, will establlish that
$$u(=\eta _\rho u)H^{s+\epsilon }(๐^n).$$
Statement (1.3) then follows by covering $`B_{1/2}`$ with small balls. Therefore the goal is to choose a suitable $`\rho `$ so that for some constant $`C>0`$, $`C=C(u,V,\rho ,n,\alpha ,s)`$,
(2.2)
$$P_ku_2\frac{C}{2^{(s+\epsilon )k}}\mathrm{for}\mathrm{all}k1.$$
Clearly it is enough to prove (2.2) only for large $`k`$.
2. The product in the right hand side of (2.1) is an integrable function as a result of (1.2). Hence, applying the Littlewood-Paley projection, we derive that
(2.3) $`2^{2\alpha k}P_ku_2`$ $``$ $`P_k(Vu)_2+P_kF_2`$
$``$ $`P_k(Vu)_2+C_F2^{(2\alpha s1)k}.`$
Thus to prove (2.2) we need to estimate $`P_k(Vu)`$. We take into account the localisation of the Littlewood-Paley projections in the frequency space. It implies that for $`f,g๐ฎ^{}`$ the distribution $`P_k(P_ifP_jg)`$ vanishes identically if
$$\left(B_{2^{i+1}}B_{2^{i1}}+B_{2^{j+1}}B_{2^{j1}}\right)\left(B_{2^{k+1}}B_{2^{k1}}\right)=\mathrm{}.$$
Consequently for a fixed $`k๐`$
(2.4) $`P_k(Vu)`$ $`=`$ $`{\displaystyle \underset{i,j๐}{}}P_k(P_iVP_ju)`$
$`=`$ $`\left\{{\displaystyle \underset{i,jLL}{}}+{\displaystyle \underset{i,jLH}{}}+{\displaystyle \underset{i,jHL}{}}+{\displaystyle \underset{i,jHH}{}}\right\}P_k(P_iVP_ju)`$
$`=`$ $`I+II+III+IV,`$
where $`LL`$, $`LH`$, $`HL`$, and $`HH`$ are the low-low, low-high, high-low, and high-high frequencies interaction zones on the integer lattice:
$`LL`$ $`=`$ $`\{i,j๐:k5i,jk+7,\mathrm{min}\{i,j\}k+5\},`$
$`LH`$ $`=`$ $`\{i,j๐:i<k5,k3jk+3\},`$
$`HL`$ $`=`$ $`\{i,j๐:k3ik+3,j<k5\},`$
$`HH`$ $`=`$ $`\{i,j๐:i,j>k+5,|ij|3\}.`$
We are going to estimate the four terms in (2.4) separately. For brevity set
$$\delta =V_{n/2\alpha }.$$
As mentioned above, we can make $`\delta `$ as small as we wish by choosing a small enough $`\rho `$ in (2.1). We will always assume that $`k`$ is big enough, say $`k10`$.
3. By properties of $`P_k`$ and the Bernstein inequality
$`I_2`$ $``$ $`{\displaystyle \underset{i,jLL}{}}P_iVP_ju_2`$
$``$ $`{\displaystyle \underset{i,jLL}{}}P_iV_{\mathrm{}}P_ju_2`$
$``$ $`2^{nk(2\alpha /n)}\delta {\displaystyle \underset{j=k5}{\overset{k+7}{}}}P_ju_2.`$
Term $`II`$ is estimated exactly the same way. It is convinient to record the final estimate in the following form
(2.5)
$$I_2+II_2\delta 2^{(2\alpha s)k}\underset{j=k5}{\overset{k+7}{}}2^{sj}P_ju_2$$
4. To estimate $`III`$ we distinguish two cases. First, assume that
(2.6)
$$n4\alpha ,$$
and hence $`n/2\alpha 2`$. Apply the Holder inequality to derive
$`III_2`$ $``$ $`P_{k3k+3}VP_0u_2+{\displaystyle \underset{j=1}{\overset{k5}{}}}P_{k3k+3}VP_ju_2`$
$``$ $`P_{k3k+3}V_2P_0u_{\mathrm{}}`$
$`+{\displaystyle \underset{j=1}{\overset{k5}{}}}P_{k3k+3}V_2P_ju_{\mathrm{}}`$
$`=`$ $`X+Y.`$
From the Bernstein inequalities we deduce that
$`X`$ $``$ $`2^{nk((2\alpha /n)(1/2))}V_{n/2\alpha }P_0u_2`$
$``$ $`2^{2\alpha k(n/2)k}\delta P_0u_2,`$
and similarly
$`Y`$ $``$ $`{\displaystyle \underset{j=1}{\overset{k5}{}}}2^{nk((2\alpha /n)(1/2))}\delta \mathrm{\hspace{0.17em}2}^{nj/2}P_ju_2`$
$``$ $`\delta {\displaystyle \underset{j=1}{\overset{k5}{}}}2^{2\alpha k(n/2)k}\mathrm{\hspace{0.17em}2}^{(n/2)jsj}\mathrm{\hspace{0.17em}2}^{sj}P_ju_2.`$
Consequently, in the case of (2.6), we can write the final estimate for $`III`$ as
(2.7) $`III_2`$ $``$ $`\delta 2^{(2\alpha (n/2))k}P_0u_2`$
$`+\delta 2^{(2\alpha s)k}{\displaystyle \underset{j=1}{\overset{k5}{}}}\left(2^{sj}P_ju_2\right)2^{((n/2)s)(jk)}.`$
Next assume that
(2.8)
$$n>4\alpha .$$
Hence
$$\frac{2\alpha }{n}+\frac{n4\alpha }{2n}=\frac{1}{2},\mathrm{and}\frac{n}{2\alpha },\frac{2n}{n4\alpha }>2.$$
By the Holder inequality
$`III_2`$ $``$ $`P_{k3k+3}V_{n/2\alpha }P_0u_{2n/(n4\alpha )}`$
$`+{\displaystyle \underset{j=1}{\overset{k5}{}}}P_{k3k+3}V_{n/2\alpha }P_ju_{2n/(n4\alpha )}`$
$`=`$ $`Z+W.`$
The Bernstein inequalities imply that
$$Z\delta P_0u_2,$$
and
$`W`$ $``$ $`\delta {\displaystyle \underset{j=1}{\overset{k5}{}}}2^{nj((1/2)(1/2)+(2\alpha /n))}P_ju_2`$
$``$ $`\delta {\displaystyle \underset{j=1}{\overset{k5}{}}}2^{(2\alpha s)j}\mathrm{\hspace{0.17em}2}^{sj}P_ju_2.`$
Consequently, in the case of (2.8), the final estimate for $`III`$ can be written as
(2.9) $`III_2`$ $``$ $`\delta P_0u_2`$
$`+\delta 2^{(2\alpha s)k}{\displaystyle \underset{j=1}{\overset{k5}{}}}2^{(2\alpha s)(jk)}\left(2^{sj}P_ju_2\right).`$
5. To estimate $`IV`$ we also need to consider two cases. First assume that (2.6) holds. By the Holder inequality
$`P_k(P_iVP_ju)_2`$ $``$ $`2^{nk/2}P_k(P_iVP_ju)_1`$
$``$ $`2^{nk/2}P_iV_{n/2\alpha }P_ju_{n/(n2\alpha )}.`$
According to (2.6) we have
$$\frac{n}{n2\alpha }2.$$
Therefore we can continue with the help of Bernstein inequality and derive that
$$P_k(P_iVP_ju)_22^{nk/2}\delta 2^{nj((1/2)1+(2\alpha /n))}P_ju_2.$$
After the summation over $`i`$ and $`j`$ lying in the $`HH`$ zone we discover that
(2.10)
$$IV_2\delta 2^{(2\alpha s)k}\underset{j=k}{\overset{\mathrm{}}{}}2^{((n/2)2\alpha +s)(kj)}\left(2^{sj}P_ju_2\right)$$
provided (2.6) holds.
Next assume that (2.8) holds. Then define $`q`$, $`1q2`$ by writing
$$\frac{1}{q}=\frac{1}{2}+\frac{2\alpha }{n}.$$
Bernstein and Holder inequalities imply that
$`P_k(P_iVP_ju)_2`$ $``$ $`2^{nk(2\alpha /n)}P_k(P_iVP_ju)_q`$
$``$ $`2^{nk(2\alpha /n)}P_iVP_ju_q`$
$``$ $`2^{2\alpha k}\delta P_ju_2.`$
Summing this estimate over $`i`$ and $`j`$ in the $`HH`$ region, we conclude that in the case of (2.8)
(2.11)
$$IV_2\delta 2^{(2\alpha s)k}\underset{j=k}{\overset{\mathrm{}}{}}2^{s(kj)}\left(2^{sj}P_ju_2\right).$$
6. Now we can prove the desired estimate (2.2). If (2.6) holds, then substitute (2.5), (2.7), and (2.10) into (2.4). If (2.8) holds then use (2.5), (2.9), and (2.11). To express the result define
$$\theta =\{\begin{array}{ccc}\mathrm{min}\{1,(n/2)s,(n/2)+s2\alpha \}\hfill & \mathrm{for}& n4\alpha \hfill \\ \mathrm{min}\{1,s,2\alpha s\}\hfill & \mathrm{for}& 4\alpha <n.\hfill \end{array}$$
According to assumptions of the theorem, $`\theta >0`$. Then we derive from (2.3) that for $`k10`$
(2.12) $`2^{sk}P_ku_2`$ $``$ $`C_1(u,\rho )2^{\theta k}`$
$`+C_2(n,\alpha ,s)\delta {\displaystyle \underset{j=0}{\overset{\mathrm{}}{}}}\left(2^{sj}P_ju_2\right)2^{\theta |jk|}.`$
For convenience set
$$a_k=2^{sk}P_ku_2,k=0,1,\mathrm{}.$$
We intend to use elemetary iteration Lemma 2.1 below to bound the sequence $`\{a_k\}`$. First take $`\epsilon =\theta /2`$. Next, find $`\rho >0`$ such that in (2.12) we have
$$\stackrel{~}{\delta }\stackrel{def}{=}C_2\delta <(12^{\epsilon /100})/2.$$
Then utilising (2.12) we can choose $`J=J(u,\rho )`$ such that
$$a_k\frac{1}{2^{\epsilon k}}+\stackrel{~}{\delta }\underset{j=0}{\overset{\mathrm{}}{}}\frac{a_j}{2^{2\epsilon |kj|}}\mathrm{for}kJ$$
with $`\stackrel{~}{\delta }`$ satisfying (2.13). Now, utilising the definition of $`H^s`$-norm find $`K=K(u,\rho )`$ such that
$$a_k1\mathrm{for}kK.$$
Finally set
$$S=J+K.$$
All assumptions of Lemma 2.1 now hold and we derive (2.2).
The following lemma is a statement about number sequences. The proof of the lemma is a careful but straightforward iteration of its assumptions. Actually we establish a stronger statement: the proof shows that (2.15) holds even if in (2.1) we replace $`2\epsilon `$ by any $`\epsilon ^{}>\epsilon `$.
###### Lemma 2.1
Let $`\epsilon >0`$, let $`\delta `$ satisfy
(2.13)
$$0<\delta <(12^\epsilon )/2,$$
and let the sequence $`\{a_k\}`$ satisfy
$`0a_k1\mathrm{for}kS,`$
(2.14) $`a_k{\displaystyle \frac{1}{2^{\epsilon k}}}+\delta {\displaystyle \underset{j0}{}}{\displaystyle \frac{a_j}{2^{2\epsilon |kj|}}}\mathrm{for}kS,`$
with some $`S0`$. Then
(2.15)
$$a_k\frac{M}{2^{\epsilon k}},k=0,1,\mathrm{},$$
with a constant $`M0`$, $`M=M(\epsilon ,\delta ,S,\{a_k\}_l^{\mathrm{}})`$.
Proof. 1. From the bounds on $`a_k`$ we derive at once that
(2.16)
$$a_k\frac{A}{2^{\epsilon k}}+\delta \underset{jS}{}\frac{a_j}{2^{2\epsilon |kj|}}\mathrm{for}\mathrm{all}kS$$
with a constant $`A>0`$, $`A=A(\epsilon ,\delta ,S,\{a_k\}_l^{\mathrm{}})`$. Define
$$C_\epsilon =2/\left(12^\epsilon \right).$$
Then, replacing $`a_j`$ in (2.16) by $`1`$, we also have
(2.17) $`a_k`$ $``$ $`{\displaystyle \frac{A}{2^{\epsilon k}}}+\delta {\displaystyle \frac{2}{12^{2\epsilon }}}`$
$``$ $`{\displaystyle \frac{A}{2^{\epsilon k}}}+\delta C_\epsilon \mathrm{for}\mathrm{all}kS.`$
2. We claim that for any $`kS`$ and any $`N0`$ the estimate
(2.18)
$$a_k\frac{A}{2^{\epsilon k}}\left(1+\delta C_\epsilon +\mathrm{}+(\delta C_\epsilon )^N\right)+(\delta C_\epsilon )^{N+1}$$
holds. Indeed, for $`N=0`$ and all $`kS`$ this is just (2.17). Assume now that (2.18) holds for some $`N`$ and all $`kS`$. Then substitute (2.18) into (2.16) to discover that for any $`kS`$
$`a_k`$ $``$ $`{\displaystyle \frac{A}{2^{\epsilon k}}}+\delta A\left(1+\delta C_\epsilon +\mathrm{}+(\delta C_\epsilon )^N\right){\displaystyle \underset{jS}{}}{\displaystyle \frac{1}{2^{\epsilon j}2^{2\epsilon |jk|}}}`$
$`+\delta (\delta C_\epsilon )^{N+1}{\displaystyle \underset{jS}{}}{\displaystyle \frac{1}{2^{2\epsilon |jk|}}}`$
$``$ $`{\displaystyle \frac{A}{2^{\epsilon k}}}\left(1+\delta C_\epsilon +\mathrm{}+(\delta C_\epsilon )^{N+1}\right)+(\delta C_\epsilon )^{N+2},`$
because for $`kS`$
$`{\displaystyle \underset{jS}{}}{\displaystyle \frac{1}{2^{\epsilon j}2^{2\epsilon |jk|}}}`$ $`=`$ $`{\displaystyle \underset{j=S}{\overset{k}{}}}{\displaystyle \frac{2^{\epsilon j}}{2^{2\epsilon k}}}+{\displaystyle \underset{j=k+1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{2^{2\epsilon k}}{2^{3\epsilon j}}}`$
$``$ $`{\displaystyle \frac{1}{2^{\epsilon k}}}\left({\displaystyle \frac{1}{12^\epsilon }}\right)+{\displaystyle \frac{1}{2^{\epsilon k}}}\left({\displaystyle \frac{1}{12^{2\epsilon }}}\right)`$
$``$ $`{\displaystyle \frac{C_\epsilon }{2^{\epsilon k}}}.`$
Hence (2.18) is proved.
3. Finally, sending $`N`$ to infinity in (2.18), we deduce according to (2.13) that $`\delta C_\epsilon <1`$ and
$$a_k\left(\frac{A}{1\delta C_\epsilon }\right)\frac{1}{2^{\epsilon k}}\mathrm{for}\mathrm{all}kS.$$
Thus (2.15) holds.
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# Study of the process ๐โบโข๐โปโ๐โบโข๐โป in the energy region 400<โ๐ <1000 MeV.
## I Introduction
The cross section of the $`e^+e^{}\pi ^+\pi ^{}`$ process in the energy region $`\sqrt{s}<1000`$ MeV can be described within the vector meson dominance model (VDM) framework and is determined by the transitions $`V\pi ^+\pi ^{}`$ of the light vector mesons ($`V=\rho ,\omega ,\rho ^{},\rho ^{\prime \prime }`$) into the final state. The main contribution in this energy region comes from the $`\rho \pi ^+\pi ^{}`$ and from the G-parity violating $`\omega \pi ^+\pi ^{}`$ transitions. Studies of the $`e^+e^{}\pi ^+\pi ^{}`$ reaction allow us to determine the $`\rho `$ and $`\omega `$ meson parameters and provide information on the $`G`$-parity violation mechanism.
At low energies the $`e^+e^{}\pi ^+\pi ^{}`$ cross section gives the dominant contribution to the celebrated ratio $`R(s)=\sigma (e^+e^{}\text{hadrons})/\sigma (e^+e^{}\mu ^+\mu ^{})`$, which is used for calculation of the dispersion integrals. For example, for evaluation of the electromagnetic running coupling constant at the $`Z`$-boson mass $`\alpha _{em}(s=m_Z^2)`$, or for determination of the hadronic contribution $`a_\mu ^{hadr}`$ to the anomalous magnetic moment of the muon, which nowadays is measured with very high accuracy $`5\times 10^6`$ bnl1 ; bnl2 .
Assuming conservation of the vector current (CVC) in the isospin symmetry limit, the spectral function of the $`\tau ^\pm \pi ^\pm \pi ^0\nu _\tau `$ decay can be related to the isovector part of the $`e^+e^{}\pi ^+\pi ^{}`$ cross section. The spectral function was determined with high precision in Ref.aleph ; opal ; cleo2 . The comparison of the $`e^+e^{}\pi ^+\pi ^{}`$ cross section with what follows from the spectral function provides an accurate test of the CVC hypothesis.
The process $`e^+e^{}\pi ^+\pi ^{}`$ in the energy region $`\sqrt{s}<1000`$ MeV was studied in several experiments augu ; ausl ; bena ; quen ; vas1 ; buki ; vas2 ; vas3 ; kur1 ; kur2 ; spec ; olya ; kmd2 ; kloe during more than 30 years. In present work the results of the $`e^+e^{}\pi ^+\pi ^{}`$ cross section measurement with SND detector at $`390\sqrt{s}980`$ MeV are reported.
## II Experiment
The SND detector sndnim operated from 1995 to 2000 at the VEPP-2M vepp2 collider in the energy range $`\sqrt{s}`$ from 360 to 1400 MeV. The detector contains several subsystems. The tracking system includes two cylindrical drift chambers. The three-layer spherical electromagnetic calorimeter is based on NaI(Tl) crystals. The muon/veto system consists of plastic scintillation counters and two layers of streamer tubes. The calorimeter energy and angular resolutions depend on the photon energy as $`\sigma _E/E(\%)=4.2\%/\sqrt[4]{E(\text{GeV})}`$ and $`\sigma _{\varphi ,\theta }=0.82^{}/\sqrt{E(\mathrm{GeV})}0.63^{}`$. The tracking system angular resolution is about $`0.5^{}`$ and $`2^{}`$ for azimuthal and polar angles respectively.
In 1996 โ 2000 the SND detector collected data in the energy region $`\sqrt{s}<980`$ MeV with integrated luminosity about $`10.0\text{pb}^1`$. The beam energy was calculated from the magnetic field value in the bending magnets of the collider. The accuracy of the energy setting is about 0.1 MeV. The beam energy spread varies in the range from 0.06 MeV at $`\sqrt{s}=360`$ MeV to 0.35 MeV at $`\sqrt{s}=970`$ MeV.
## III Data Analysis
The cross section of the $`e^+e^{}\pi ^+\pi ^{}`$ process was measured in the following way.
1. The collinear events $`e^+e^{}e^+e^{},\pi ^+\pi ^{},\mu ^+\mu ^{}`$ were selected;
2. The selected events were sorted into the two classes: $`e^+e^{}`$ and $`\pi ^+\pi ^{},\mu ^+\mu ^{}`$ using the energy deposition in the calorimeter layers;
3. The $`e^+e^{}e^+e^{}`$ events were used for integrated luminosity determination. The events of the $`e^+e^{}\mu ^+\mu ^{}`$ process were subtracted according to the theoretical cross section, integrated luminosity and detection efficiency;
4. In order to determine the cross section of the $`e^+e^{}\pi ^+\pi ^{}`$ process, the number of $`e^+e^{}\pi ^+\pi ^{}`$ events in each energy point were normalized on the integrated luminosity and divided by the detection efficiency and radiative correction.
The detection efficiency was obtained from Monte Carlo (MC) simulation sndnim . The MC simulation of SND is based on UNIMOD unimod package. The SND geometrical model description comprises about 10000 distinct volumes and includes details of the SND design. The primary generated particles are tracked through the detector media taking into account the following effects: ionization losses, multiple scattering, bremsstrahlung of electrons and positrons, Compton effect and Rayleigh scattering, $`e^+e^{}`$ pair production by photons, photo-effect, unstable particles decays, interaction of stopped particles, nuclear interaction of hadrons union ; umnuc1 ; umnuc2 . After that the signals produced in each detector element are simulated. The electronics noise, signals pile up, the actual time and amplitude resolutions of the electronics channels and broken channels were taken into account during processing the Monte Carlo events to provide the adaptable account of variable experimental conditions.
The MC simulation of the processes $`e^+e^{}e^+e^{},\mu ^+\mu ^{},\pi ^+\pi ^{}`$ was based on the formula obtained in the Ref.berkl ; arbuzqed ; arbuzhad . The simulation of the process $`e^+e^{}e^+e^{}`$ was performed with the cut $`30^{}<\theta _{e^\pm }<150^{}`$ on the polar angles of the final electron and positron.
The $`e^+e^{}e^+e^{}`$, $`\mu ^+\mu ^{}`$ and $`\pi ^+\pi ^{}`$ events are differed by energy deposition in the calorimeter. In $`e^+e^{}e^+e^{}`$ events the electrons produce the electromagnetic shower with the most probable energy losses about 0.92 of the initial particle energy. The distributions of the energy deposition of the electrons with the different energies are shown in Fig.2. The experimental and simulated spectra are in good agreement. Muons lose their energy by ionization of the calorimeter material through which they pass and their energy deposition spectra are well modeled in simulation (Fig.2). The similar ionization losses are experienced by charged pions and this part of the charged pion energy deposition is well described by simulation (Fig.4). But pions lose their energy also due to nuclear interactions which is not so accurately reproduced in simulation. This leads to some difference in energy deposition spectra in experiment and simulation for charged pions (Fig.4).
The discrimination between electrons and pions in the SND detector is based on difference in longitudinal energy deposition profiles (deposition in calorimeter layers) for these particles. To use in the most complete way the correlations between energy depositions in the calorimeter layers, the corresponding separation parameter was based on the neural network approach neural1 . For each energy point the neural network โ multilayer perceptron was constructed. The network had input layer consisting of 7 neurons, two hidden layers with 20 neurons each and the output layer with one neuron. As the input data the network used the energy depositions of the particles in calorimeter layers and the polar angle of one of the particles. The output signal $`R_{e/\pi }`$ is a number in the interval from -0.5 to 1.5. The network was trained by using simulated $`e^+e^{}\pi ^+\pi ^{}`$ and $`e^+e^{}e^+e^{}`$ events. The distribution of the discrimination parameter $`R_{e/\pi }`$ is shown in Fig.5. The $`e^+e^{}e^+e^{}`$ events are located in the region $`R_{e/\pi }>0.5`$, while $`e^+e^{}\pi ^+\pi ^{},\mu ^+\mu ^{}`$ events at $`R_{e/\pi }<0.5`$.
### III.1 Selection criteria
During the experimental runs, the first-level trigger sndnim selects events with one or more tracks in tracking system and with two clusters in calorimeter with the spatial angle between the clusters more than $`100^{}`$. The threshold on energy deposition in cluster was equal to 25 MeV. The threshold on the total energy deposition in the calorimeter was set equal to 140 MeV in the energy region $`\sqrt{s}850`$ MeV, and to 100 MeV, or was absent at all, below 850 MeV. During processing of the experimental data the event reconstruction is performed sndnim ; phi98 . For further analysis, events containing two charged particles with $`|z|<10`$ cm and $`r<1`$ cm were selected. Here $`z`$ is the coordinate of the charged particle production point along the beam axis (the longitudinal size of the interaction region depends on beam energy and varies from 1.5 to 2.5 cm); $`r`$ is the distance between the charged particle track and the beam axis in the $`r\varphi `$ plane. The polar angles of the charged particles were bounded by the criterion: $`55^{}<\theta <125^{}`$ and the energy deposition of each of them was required to be greater than 50 MeV. The following cuts on the acollinearity angles in the azimuthal and polar planes were applied: $`|\mathrm{\Delta }\varphi |<10^{}`$ and $`|\mathrm{\Delta }\theta |<10^{}`$. In the event sample selected under these conditions one has the $`e^+e^{}e^+e^{}`$, $`\pi ^+\pi ^{}`$, $`\mu ^+\mu ^{}`$ events, cosmic muons background and a small contribution from the $`e^+e^{}\pi ^+\pi ^{}\pi ^0`$ reaction at $`\sqrt{s}m_\omega `$. The muon system $`veto`$ was used for suppression of the cosmic muon background ($`veto=0`$).
### III.2 The background from the cosmic muons and from the $`e^+e^{}\pi ^+\pi ^{}\pi ^0`$ process.
The number of background events from the $`e^+e^{}\pi ^+\pi ^{}\pi ^0`$ process was estimated in the following way:
$`N_{3\pi }(s)=\sigma _{3\pi }(s)ฯต_{3\pi }(s)IL(s),`$ (1)
where $`\sigma _{3\pi }(s)`$ is the cross section of the $`e^+e^{}\pi ^+\pi ^{}\pi ^0`$ process with the radiative corrections taken into account, $`IL(s)`$ is the integrated luminosity, $`ฯต_{3\pi }(s)`$ is the detection probability for the background process obtained from the simulation under the selection criteria described above. The values of $`\sigma _{3\pi }(s)`$ were taken from the SND measurements pi3omeg . Although $`\sigma _{3\pi }(m_\omega )1300`$ nb, the $`e^+e^{}3\pi `$ process contribution to the total number of the collinear events at the $`\omega `$ resonance peak is less than 0.3 %. The leading role in the suppression of this background was played by the cuts on the acollinearity angles $`\mathrm{\Delta }\theta `$ and $`\mathrm{\Delta }\varphi `$. In order to check the estimation (1), the events containing two and more photons with energy depositions more than 200 MeV were considered. The constraint on the photons energy deposition greatly suppresses not the $`e^+e^{}3\pi `$ events, as a result of the fact that our selection criteria select the $`e^+e^{}3\pi `$ events with collinear charged pions and therefore the neutral pion in this events has relatively low energy. In order to obtain $`e^+e^{}3\pi `$ events number $`n_{3\pi }`$, the invariant mass spectrum $`m_{\gamma \gamma }`$ (Fig.7) was fitted by the sum of Gaussian and the second order polynomial: $`G(m_{\gamma \gamma })\times n_{3\pi }+P_2(m_{\gamma \gamma })\times (nn_{3\pi })`$. The value of $`n_{3\pi }`$ agrees with events number calculated according to (1).
The cosmic muon background was suppressed by the muon/veto system. The $`z`$ coordinate distribution for the charged particle production point along the beam axis is shown in Fig.7 for collinear events. The $`e^+e^{}`$ annihilation events have the Gaussian distribution peaked at $`z=0`$, while the cosmic background distribution is nearly uniform and clearly extends outside the peak. As the Fig.7 shows, the muon system $`veto`$ ($`veto=1`$) separates cosmic muons from the $`e^+e^{}`$ annihilation events. The residual events number of the cosmic muon background was estimated from the following formula:
$`N_\mu =\nu _\mu \times T.`$ (2)
Here $`\nu _\mu 1.3\times 10^3`$ Hz is the frequency of cosmic background registration under the applied selection criteria, $`T`$ is the time of data taking. The value of $`\nu _\mu `$ was obtained by using data collected in special runs without beams in collider. The first-level trigger counting rate in these runs was 2 Hz. The contribution of the cosmic background to the total number of selected collinear events depends on energy $`\sqrt{s}`$ and varies from 0.1 % to 1 %.
The $`e^+e^{}\pi ^+\pi ^{}\pi ^0`$ events are concentrated in the $`R_{e/\pi }`$ discrimination parameter region $`R_{e/\pi }<0.5`$. The cosmic background events at the energies $`\sqrt{s}>600`$ also fall in the area $`R_{e/\pi }<0.5`$, because the energy deposition of the cosmic muons is much lower than the energy deposition in the $`e^+e^{}e^+e^{}`$ events. For the lower center of mass energies the cosmic background moves to the area $`R_{e/\pi }>0.5`$, because in this case the energy depositions are close.
### III.3 Detection efficiency
The $`\mathrm{\Delta }\varphi `$ and $`\mathrm{\Delta }\theta `$ distributions of the $`e^+e^{}e^+e^{}`$ and $`e^+e^{}\pi ^+\pi ^{}`$ events are shown in Fig.9,9, 11 and 11. Experiment and simulation agree rather well. As a measure of systematic uncertainty due to $`\mathrm{\Delta }\theta `$ cut the following value was used:
$`\delta _{\mathrm{\Delta }\theta }={\displaystyle \frac{\delta _{\mathrm{\Delta }\theta }^{\pi \pi }}{\delta _{\mathrm{\Delta }\theta }^{ee}}},`$ (3)
where
$$\delta _{\mathrm{\Delta }\theta }^x=\frac{n_x(|\mathrm{\Delta }\theta |<10^{})}{N_x(|\mathrm{\Delta }\theta |<20^{})}\text{ }/\text{ }\frac{m_x(|\mathrm{\Delta }\theta |<10^{})}{M_x(|\mathrm{\Delta }\theta |<20^{})},\text{ }x=\pi \pi (ee).$$
Here $`n_x(|\mathrm{\Delta }\theta |<10^{})`$ and $`m_x(|\mathrm{\Delta }\theta |<10^{})`$ are the numbers of experimental and simulated events, selected under the condition $`|\mathrm{\Delta }\theta |<10^{}`$, while $`N_x(|\mathrm{\Delta }\theta |<20^{})`$ and $`M_x(|\mathrm{\Delta }\theta |<20^{})`$ are the numbers of experimental and simulated events with $`|\mathrm{\Delta }\theta |<20^{}`$. The $`\delta _{\mathrm{\Delta }\theta }`$ does not depend on energy, its average value is equal to 0.999 and it has systematic spread of 0.4 %. This systematic spread was added to the error of the cross section measurement in each energy point. Systematic error due to the $`\mathrm{\Delta }\varphi `$ cut is significantly lower and was neglected.
The polar angle distributions for the $`e^+e^{}e^+e^{}`$ and $`e^+e^{}\pi ^+\pi ^{}`$ processes are shown in Fig.13 and 13. The ratio of these $`\theta `$ distributions is shown in Fig.15. The experimental and simulated distributions are in agreement. In order to estimate the systematic inaccuracy due to the $`\theta `$ angle selection cut the following ratio was used:
$`\delta _\theta ={\displaystyle \frac{\delta (\theta _x)}{\delta (55^{})}},`$ (4)
where
$$\delta (\theta _x)=\frac{N_{\pi \pi }(\theta _x<\theta <180^{}\theta _x)}{N_{ee}(\theta _x<\theta <180^{}\theta _x)}/\frac{M_{\pi \pi }(\theta _x<\theta <180^{}\theta _x)}{M_{ee}(\theta _x<\theta <180^{}\theta _x)},\text{ }50^{}<\theta _x<90^{}.$$
Here $`N_{\pi \pi }(\theta _x<\theta <180^{}\theta _x)`$, $`N_{ee}(\theta _x<\theta <180^{}\theta _x)`$, $`M_{\pi \pi }(\theta _x<\theta <180^{}\theta _x)`$, $`M_{ee}(\theta _x<\theta <180^{}\theta _x)`$ are the experimental and simulated $`e^+e^{}\pi ^+\pi ^{}`$ and $`e^+e^{}e^+e^{}`$ event numbers in the angular range $`\theta _x<\theta <180^{}\theta _x`$. The maximal difference of $`\delta _\theta `$ from unity was found to be 0.8%. This value was taken as a systematic error $`\sigma _\theta =0.8`$% associated with the angular selection cut.
In the tracking system the particle track can be lost due to reconstruction inefficiency. The probabilities to find the track was determined by using experimental data themselves. It was found to be $`\epsilon _e0.996`$ for electrons and $`\epsilon _\pi 0.995`$ for pions. In simulation these values actually do not differ from unity, while in reality the track finding probability for electrons is slightly greater then for pions. So the detection efficiency was multiply by the correction coefficient:
$`\delta _{rec}=\left[{\displaystyle \frac{\epsilon _\pi }{\epsilon _e}}\right]^2=0,997`$ (5)
Pions can be lost due to the nuclear interaction in the detector material before the tracking system, for example, via the reaction $`\pi ^\pm N\pi ^\pm N`$ with the final pion scattered at the large angle or via charge exchange reaction $`\pi ^\pm N\pi ^0N`$. As a measure of systematic inaccuracy associated to this effect the difference from unity of the following quantity was used:
$`\delta _{nucl}=\left[\left(1{\displaystyle \frac{n}{3N}}\right)/\left(1{\displaystyle \frac{m}{3M}}\right)\right]^2,`$ (6)
where $`N`$ and $`M`$ is the pions numbers in experiment and simulation; $`n`$ and $`m`$ is the pions numbers in experiment and simulation which had a track in the drift chamber nearest to the beam-pipe, but the corresponding track in the second drift chamber and associated cluster in the calorimeter were not found. The particle loss probability was divided by 3 โ the ratio of amounts of the matter between the drift chambers and before the tracking system. The deviation of $`\delta _{nucl}`$ from 1 was taken as a systematic error $`\sigma _{nucl}=0.2`$ %.
Uncertainties in simulation of pions nuclear interactions imply that the cut on the particles energy deposition leads to an inaccuracy in detection efficiency of the $`e^+e^{}\pi ^+\pi ^{}`$ process. In order to take into account this inaccuracy, the detection efficiency was multiplied by the correction coefficients. The correction coefficients was obtained by using events of the $`e^+e^{}\pi ^+\pi ^{}\pi ^0`$ reaction phi98 ; pi3omeg ; dplphi98 . Pions energies in the $`e^+e^{}\pi ^+\pi ^{}\pi ^0`$ events were determined via the kinematic fit. The pion energies were divided into the 10 MeV wide bins . For each bin the correction coefficient (Fig.15) was obtained:
$`\delta _{E>50}=\left[{\displaystyle \frac{n_i/N_i}{m_i/M_i}}\right]^2,`$ (7)
where $`i`$ is the bin number, $`N_i`$ and $`M_i`$ are the pions numbers in experiment and simulation selected in the $`i`$th bin by the kinematic fit without any cut on the energy deposition in the calorimeter ; $`n_i`$ and $`m_i`$ are the pions numbers in experiment and simulation under the condition that the pion energy deposition is greater than 50 MeV. To estimate systematic errors in determination of these correction coefficients, the ratio of the probability that both pions in simulated $`e^+e^{}\pi ^+\pi ^{}`$ events have energy deposition more than 50 MeV to the quantity $`(m_i/M_i)^2`$ was consider. This ratio is 0.994 at $`\sqrt{s}>420`$ MeV and about 0.97 at $`\sqrt{s}<420`$ MeV. The difference of this ratio from unity was taken as a systematic error $`\sigma _{E>50}`$ of the $`\delta _{E>50}`$ correction coefficient determination: $`\sigma _{E>50}=0.6`$ % at $`\sqrt{s}>420`$ MeV and $`\sigma _{E>50}=3`$ % at $`\sqrt{s}<420`$ MeV.
In the energy region $`\sqrt{s}=840`$$`970`$ MeV the probability to hit the muon/veto system for muons and pions varies from 1% upto 93%, and from 0.5% to 3% respectively. The usage of the muon system $`veto`$ for events selection ($`veto=0`$) leads to inaccuracy in the measured cross section determination due to the uncertainty in the simulation of the muons and pions traversing through the detector at $`\sqrt{s}>840`$ MeV. In order to obtain the necessary corrections, the events close to the median plane $`\varphi <10^{}`$, $`170^{}\varphi <190^{}`$, $`\varphi >350^{}`$, where the cosmic background is minimal, were used. The $`e^+e^{}\pi ^+\pi ^{}`$ cross section was measured with ($`veto=0`$) and without ($`veto0`$) using the muon system, and the following correction coefficient was obtained for each energy point:
$`\delta _{veto}={\displaystyle \frac{\sigma (e^+e^{}\pi ^+\pi ^{};veto0)}{\sigma (e^+e^{}\pi ^+\pi ^{};veto=0)}}`$ (8)
It was found that $`\delta _{veto}=0.95`$ at $`\sqrt{s}=970`$ MeV and quickly rises up to 1 for lower energies.
The detection efficiencies of the processes $`e^+e^{}\pi ^+\pi ^{}`$, $`\mu ^+\mu ^{}`$ and $`e^+e^{}`$ after all applied corrections are shown in Fig.16. The detection efficiency of the $`e^+e^{}e^+e^{}`$ reaction does not depend on energy, while for $`e^+e^{}\mu ^+\mu ^{}`$ and $`\pi ^+\pi ^{}`$ processes it does. The decrease of the $`e^+e^{}\mu ^+\mu ^{}`$ process detection efficiency at $`\sqrt{s}>800`$ MeV is caused by the fact that the probability for muons to hit the muon system rises with energy. The detection efficiency of the $`e^+e^{}\pi ^+\pi ^{}`$ process at $`\sqrt{s}>500`$ MeV is determined mainly by the cuts on the pions angles. Below 500 MeV the detection efficiency decreases due to the cut on the pions energy deposition in the calorimeter. The statistical error $`1\%`$ of the detection efficiency determination was added to the cross section measurement error in each energy point. The total systematic error of the detection efficiency determination $`\sigma _{eff}=\sigma _{E>50}\sigma _{nucl}\sigma _\theta `$ is $`\sigma _{eff}=1`$ % at $`\sqrt{s}420`$ MeV and $`\sigma _{eff}=3.1`$ % at $`\sqrt{s}<420`$ MeV.
### III.4 Measurement of the $`e^+e^{}\pi ^+\pi ^{}`$ cross section.
The number of selected events in the regions $`R_{e/\pi }<0.5`$ and $`R_{e/\pi }>0.5`$ are:
$`N=N_{\pi \pi }+N_{ee}+N_{\mu \mu }+N_\mu +N_{3\pi },`$ (9)
$`M=M_{\pi \pi }+M_{ee}+M_{\mu \mu }+M_\mu +M_{3\pi }.`$ (10)
Here $`N`$ and $`M`$ are the events numbers in the regions $`R_{e/\pi }<0.5`$ and $`R_{e/\pi }>0.5`$ respectively. $`N_\mu `$, $`M_\mu `$ and $`N_{3\pi }`$, $`M_{3\pi }`$ are the number of background events due to cosmic muons and the $`e^+e^{}\pi ^+\pi ^{}\pi ^0`$ process, calculated as was described above. The $`e^+e^{}\mu ^+\mu ^{}`$ process events number can be written as:
$`N_{\mu \mu }=\sigma _{\mu \mu }\times \epsilon _{\mu \mu }\times (1ฯต_{\mu \mu })\times IL,`$ (11)
$`M_{\mu \mu }=\sigma _{\mu \mu }\times \epsilon _{\mu \mu }\times ฯต_{\mu \mu }\times IL,`$ (12)
where $`\sigma _{\mu \mu }`$ is the $`e^+e^{}\mu ^+\mu ^{}`$ process cross section obtained according to Ref.arbuzqed , $`\epsilon _{\mu \mu }`$ is the process detection efficiency, $`ฯต_{\mu \mu }`$ is the probability for the $`e^+e^{}\mu ^+\mu ^{}`$ process events to have $`R_{e/\pi }>0.5`$. $`IL`$ is the integrated luminosity:
$`IL={\displaystyle \frac{M_{ee}}{\sigma _{ee}\times \epsilon _{ee}\times ฯต_{ee}}},`$ (13)
where $`\epsilon _{ee}`$ and $`ฯต_{ee}`$ are the detection efficiency and the probability to have $`R_{e/\pi }>0.5`$ for the process $`e^+e^{}e^+e^{}`$, $`\sigma _{ee}`$ is the process cross section with the $`30^{}<\theta <150^{}`$ angular cut for the electron and positron in the final state. The cross section $`\sigma _{ee}`$ was calculated by using BHWIDE 1.04 bhwide code with accuracy 0.5 %. The $`e^+e^{}\pi ^+\pi ^{}`$ process events number with $`R_{e/\pi }>0.5`$ and the $`e^+e^{}e^+e^{}`$ process events number with $`R_{e/\pi }<0.5`$ can be written in the following way:
$$N_{ee}=\frac{1ฯต_{ee}}{ฯต_{ee}}\times M_{ee}=\lambda _{ee}\times M_{ee},\text{ }M_{\pi \pi }=\frac{1ฯต_{ee}}{ฯต_{ee}}\times N_{\pi \pi }=\lambda _{\pi \pi }\times N_{\pi \pi }.$$
The $`e^+e^{}e^+e^{}`$ process events number with $`R_{e/\pi }>0.5`$ and the $`e^+e^{}\pi ^+\pi ^{}`$ process events number with $`R_{e/\pi }<0.5`$ are equal to:
$`M_{ee}={\displaystyle \frac{MM_\mu \lambda _{\pi \pi }\times (NN_\mu )}{\kappa \mathrm{\Delta }\times \lambda _{\pi \pi }}},`$ (14)
$`N_{\pi \pi }=NN_\mu M_{ee}\times \mathrm{\Delta }.`$ (15)
Here
$$\mathrm{\Delta }=\lambda _{ee}+\frac{\sigma _{\mu \mu }\times \epsilon _{\mu \mu }\times (1ฯต_{\mu \mu })+N_{3\pi }/IL}{\sigma _{ee}\times \epsilon _{ee}\times ฯต_{ee}},$$
$$\kappa =1+\frac{\sigma _{\mu \mu }\times \epsilon _{\mu \mu }\times ฯต_{\mu \mu }+M_{3\pi }/IL}{\sigma _{ee}\times \epsilon _{ee}\times ฯต_{ee}}.$$
The percentage of each process in the selected events in dependence on energy $`\sqrt{s}`$ is shown in Fig.18. The experimental angular distributions agree with the sum of distributions for each process weighted according to its contribution (Fig.18).
The $`e^+e^{}\pi ^+\pi ^{}`$ process cross section is calculated from the following formula:
$`\sigma _{\pi \pi }={\displaystyle \frac{N_{\pi \pi }}{IL\times \epsilon _{\pi \pi }\times (1ฯต_{\pi \pi })}}={\displaystyle \frac{\sigma _{ee}\times \epsilon _{ee}\times ฯต_{ee}}{\epsilon _{\pi \pi }\times (1ฯต_{\pi \pi })}}\times \left[{\displaystyle \frac{\kappa \mathrm{\Delta }\times \lambda _{\pi \pi }}{\frac{MM_\mu }{NN_\mu }\lambda _{\pi \pi }}}\mathrm{\Delta }\right].`$ (16)
In order to estimate the systematic uncertainty due to $`e\pi `$ discrimination, the pseudo $`\pi \pi `$ and pseudo $`ee`$ events in the experiment and simulation were formed. The pseudo $`\pi \pi `$ events were constructed by using pions from the $`e^+e^{}\pi ^+\pi ^{}\pi ^0`$ reaction. In order to construct the pseudo $`\pi \pi `$ event with the pions having energy $`E_0`$, two charged pions with energies $`E_\pi `$ such that $`|E_0E_\pi |<10`$ MeV were used from two separate $`e^+e^{}\pi ^+\pi ^{}\pi ^0`$ events. Of course, such pseudo $`\pi \pi `$ events are in general not collinear but this is irrelevant for our purposes here. The pseudo $`ee`$ event was constructed analogously from the particles of two separate collinear events such that their partners in these events have energy depositions in the calorimeter layers typical for electrons. Fig.20 and 20 show probabilities for the discrimination parameter to have values less than some magnitude in experiment and simulation for such pseudo events. Using these distributions, the corrections to the probabilities for the separation parameter $`R_{e/\pi }`$ to be greater or less than 0.5 was obtained. The difference between cross sections measured with and without these corrections was taken as a systematic error and its value does not exceed 0.5 % for different energy points.
The obtained cross sections together with the radiative corrections $`\delta _{rad}`$, including the initial and final state radiation, are presented in Table 1. The $`\delta _{rad}`$ radiative correction was calculated according to Ref.arbuzhad . The accuracy of its determination is 0.2 %. Having at hand the radiative corrections the Born cross section for the $`e^+e^{}\pi ^+\pi ^{}`$ process can be extracted as follows
$`\sigma _0(s)={\displaystyle \frac{\sigma _{\pi \pi }(s)}{\delta _{rad}(s)}}`$ (17)
The value of $`\delta _{rad}(s)`$ depends on the cross section at lower energies, so it was calculated iteratively. The iteration stops then its value changes by not more than 0.1 % in consecutive iterations. The form factor values
$$|F_\pi (s)|^2=\frac{3s}{\pi \alpha ^2\beta ^3}\sigma _{\pi \pi }(s),\text{ }\beta =\sqrt{14m_\pi ^2/s}$$
are also listed in Table 1. To evaluate the value of $`R(s)=\sigma (e^+e^{}\text{hadrons})/\sigma (e^+e^{}\mu ^+\mu ^{})`$, which is used in dispersion integrals calculation, the bare cross section $`e^+e^{}\pi ^+\pi ^{}`$ is used (the cross section without vacuum polarization contribution but taking into account the final state radiation):
$`\sigma _{\pi \pi }^{pol}(s)=\sigma _0(s)\times |1\mathrm{\Pi }(s)|^2\times \left(1+{\displaystyle \frac{\alpha }{\pi }}a(s)\right),`$ (18)
where $`\mathrm{\Pi }(s)`$ is the polarization operator calculated according to the Ref.arbuzqed from the known $`e^+e^{}\text{hadrons}`$ cross section fedor . The last factor takes into account the final state radiation, and $`a(s)`$ has the form shw
$$a(s)=\frac{1+\beta ^2}{\beta }\left[4Li_2\left(\frac{1\beta }{1+\beta }\right)+2Li_2\left(\frac{1\beta }{1+\beta }\right)3\mathrm{ln}\frac{2}{1+\beta }\mathrm{ln}\frac{1+\beta }{1\beta }2\mathrm{ln}\beta \mathrm{ln}\frac{1+\beta }{1\beta }\right]$$
$$3\mathrm{ln}\frac{4}{1\beta ^2}4\mathrm{ln}\beta +\frac{1}{\beta ^3}\left[\frac{5}{4}(1+\beta ^2)^22\right]\times \mathrm{ln}\frac{1+\beta }{1\beta }+\frac{3}{2}\frac{1+\beta ^2}{\beta ^2}.$$
Here
$$Li_2(x)=\underset{0}{\overset{x}{}}๐t\mathrm{ln}(1t)/t.$$
The values of $`\sigma _{\pi \pi }^{pol}(s)`$ are listed in Table 1.
The total systematic error of the cross section determination is:
$$\sigma _{sys}=\sigma _{eff}\sigma _{sep}\sigma _{IL}\sigma _{rad}.$$
Here $`\sigma _{eff}`$ is the systematic error of the detection efficiency determination, $`\sigma _{sep}`$ is the systematic error associated with the $`e\pi `$ separation, $`\sigma _{IL}`$ is the systematic error of the integrated luminosity determination, and $`\sigma _{rad}`$ is the uncertainty of the radiative correction calculation. The magnitudes of various contributions to the total systematic error are shown in Table 2. The total systematic error of the cross section determinations is $`\sigma _{sys}=1.3`$ % at $`\sqrt{s}420`$ MeV and $`\sigma _{sys}=3.2`$ % at $`\sqrt{s}<420`$ MeV.
## IV The $`e^+e^{}\pi ^+\pi ^{}`$ cross section analysis
### IV.1 Theoretical framework
In the framework of the vector meson dominance model, the cross section of the $`e^+e^{}\pi ^+\pi ^{}`$ process is
$`\sigma _{\pi \pi }(s)={\displaystyle \frac{4\pi \alpha ^2}{s^{3/2}}}P_{\pi \pi }(s)|A_{\pi \pi }(s)|^2`$ (19)
Here $`P_{\pi \pi }(s)`$ is the phase space factor:
$$P_{\pi \pi }(s)=q_\pi ^3(s),\text{ }q_\pi (s)=\frac{1}{2}\sqrt{s4m_\pi ^2}.$$
Amplitudes of the $`\gamma ^{}\pi ^+\pi ^{}`$ transition have the form:
$`|A_{\pi \pi }(s)|^2=\left|\sqrt{{\displaystyle \frac{3}{2}}}{\displaystyle \frac{1}{\alpha }}{\displaystyle \underset{V=\rho ,\omega ,\rho ^{},\rho ^{\prime \prime }}{}}{\displaystyle \frac{\mathrm{\Gamma }_Vm_V^3\text{ }\sqrt{m_V\sigma (V\pi ^+\pi ^{})}}{D_V(s)}}{\displaystyle \frac{e^{i\varphi _{\rho V}}}{\sqrt{q_\pi ^3(m_V)}}}\right|^2,`$ (20)
where
$$D_V(s)=m_V^2si\text{ }\sqrt{s}\mathrm{\Gamma }_V(s),\text{ }\mathrm{\Gamma }_V(s)=\underset{f}{}\mathrm{\Gamma }(Vf,s).$$
Here $`f`$ denotes the final state of the $`V`$ vector meson decay, $`m_V`$ is the vector meson mass, $`\mathrm{\Gamma }_V=\mathrm{\Gamma }_V(m_V)`$. The following forms of the energy dependence of the vector mesons total widths were used:
$$\mathrm{\Gamma }_\omega (s)=\frac{m_\omega ^2}{s}\frac{q_\pi ^3(s)}{q_\pi ^3(m_\omega )}\mathrm{\Gamma }_\omega B(\omega \pi ^+\pi ^{})+\frac{q_{\pi \gamma }^3(s)}{q_{\pi \gamma }^3(m_\omega )}\mathrm{\Gamma }_\omega B(\omega \pi ^0\gamma )+\frac{W_{\rho \pi }(s)}{W_{\rho \pi }(m_\omega )}\mathrm{\Gamma }_\omega B(\omega 3\pi ),$$
$$\mathrm{\Gamma }_V(s)=\frac{m_V^2}{s}\frac{q_\pi ^3(s)}{q_\pi ^3(m_V)}\mathrm{\Gamma }_V\text{ }(V=\rho ,\rho ^{},\rho ^{\prime \prime })$$
Here $`q_{\pi \gamma }=(sm_\pi ^2)/2\sqrt{s}`$, $`W_{\rho \pi }(s)`$ is the phase space factor for the $`\rho \pi \pi ^+\pi ^{}\pi ^0`$ final state phi98 ; pi3omeg ; dplphi98 . In the energy dependence of the $`\rho ,\rho ^{},\rho ^{\prime \prime }`$ mesons widths only the $`V\pi ^+\pi ^{}`$ decays were taken into account. Such approach is justified in the energy region $`\sqrt{s}<1000`$ MeV. Nowadays the $`\rho ^{},\rho ^{\prime \prime }`$ decays are rather poorly known and therefore the same approximation was used also for the fitting of the data above 1000 MeV. The $`\omega `$-meson mass and width were taken from the SND measurements: $`m_\omega =782.79`$ MeV, $`\mathrm{\Gamma }_\omega =8.68`$ MeV pi3omeg .
The relative decay probabilities were calculated as follows
$$B(VX)=\frac{\sigma (VX)}{\sigma (V)},\text{ }\sigma (V)=\underset{X}{}\sigma (VX),\text{ }\sigma (VX)=\frac{12\pi B(Ve^+e^{})B(VX)}{m_V^2}.$$
In the analysis presented here we have used $`\sigma (\omega \pi ^0\gamma )=155.8`$ nb, $`\sigma (\omega 3\pi )=1615`$ nb obtained in the SND experiments pi3omeg ; pi0gam .
The parameter $`\varphi _{\rho V}`$ is the relative interference phase between the vector mesons $`V`$ and $`\rho `$, so $`\varphi _{\rho \rho }=0`$. The phases $`\varphi _{\rho V}`$ can deviate from $`180^{}`$ or $`0^{}`$, and their values can be energy dependent due to mixing between vector mesons. The phases $`\varphi _{\rho \rho ^{}}`$ and $`\varphi _{\rho \rho ^{\prime \prime }}`$ were fixed at $`180^{}`$ and $`0^{}`$, because these values are consistent with the existing experimental data for the $`e^+e^{}\pi ^+\pi ^{}`$ reaction.
Taking into account the $`\rho \omega `$ mixing, the $`\omega \pi ^+\pi ^{}`$ and $`\rho \pi ^+\pi ^{}`$ transition amplitudes can be written in the following way thrhoom ; akozi
$`A_{\omega \pi ^+\pi ^{}}+A_{\rho \pi ^+\pi ^{}}={\displaystyle \frac{g_{\gamma \rho }^{(0)}g_{\rho \pi \pi }^{(0)}}{D_\rho (s)}}\left[1{\displaystyle \frac{g_{\gamma \omega }^{(0)}}{g_{\gamma \rho }^{(0)}}}\epsilon (s)\right]+{\displaystyle \frac{g_{\gamma \omega }^{(0)}g_{\rho \pi \pi }^{(0)}}{D_\omega (s)}}\left[\epsilon (s)+{\displaystyle \frac{g_{\omega \pi \pi }^{(0)}}{g_{\rho \pi \pi }^{(0)}}}\right],`$ (21)
where
$$\epsilon (s)=\frac{\mathrm{\Pi }_{\rho \omega }}{D_\omega (s)D_\rho (s)},\text{ }|g_{V\gamma }|=\left[\frac{3m_V^3\mathrm{\Gamma }_VB(Ve^+e^{})}{4\pi \alpha }\right]^{1/2},\text{ }|g_{V\pi \pi }|=\left[\frac{6\pi m_V^2\mathrm{\Gamma }_VB(V\pi ^+\pi ^{})}{q_\pi ^3(m_V)}\right]^{1/2}.$$
The superscript $`(0)`$ denotes the coupling constants of the bare, unmixed state. $`\mathrm{\Pi }_{\rho \omega }`$ is the polarization operator of the $`\rho \omega `$ mixing:
$`\mathrm{\Pi }_{\rho \omega }(s)=\text{Re}(\mathrm{\Pi }_{\rho \omega }(s))+i\text{ }\text{Im}(\mathrm{\Pi }_{\rho \omega }(s)).`$ (22)
The $`\text{Im}(\mathrm{\Pi }_{\rho \omega }(s))`$ can be written as
$`\text{Im}(\mathrm{\Pi }_{\rho \omega }(s))=\sqrt{s}\{{\displaystyle \frac{g_{\rho \pi \pi }^{(0)}g_{\omega \pi \pi }^{(0)}q_\pi ^3(s)}{6\pi s}}+{\displaystyle \frac{g_{\rho \pi \gamma }^{(0)}g_{\omega \pi \gamma }^{(0)}q_{\pi \gamma }^3(s)+g_{\rho \eta \gamma }^{(0)}g_{\omega \eta \gamma }^{(0)}q_{\eta \gamma }^3(s)}{3}},\}`$ (23)
where
$$g_{VP\gamma }=\left[\frac{3\mathrm{\Gamma }_VB(VP\gamma )}{q_{P\gamma }^3(m_V)}\right]^{1/2}.$$
We neglected the contributions to $`\text{Im}(\mathrm{\Pi }_{\rho \omega }(s))`$ due to $`VP`$ intermediate state ($`V=\omega ,\rho `$,$`P=\pi ,\eta `$). The $`\text{Re}(\mathrm{\Pi }_{\rho \omega }(s))`$ can be represented as
$`\text{Re}(\mathrm{\Pi }_{\rho \omega }(s))=\text{Re}(\mathrm{\Pi }_{\rho \omega }^\gamma (s))+\text{Re}(\mathrm{\Pi }_{\rho \omega }^{}(s)),`$ (24)
where
$`\text{Re}(\mathrm{\Pi }_{\rho \omega }^\gamma (s))={\displaystyle \frac{4\pi g_{\rho \gamma }^{(0)}g_{\omega \gamma }^{(0)}}{s}}`$ (25)
represents the one-photon contribution to the $`\text{Re}(\mathrm{\Pi }_{\rho \omega }(s))`$. Let us assume that the energy dependence of the $`\text{Re}(\mathrm{\Pi }_{\rho \omega }^{}(s))`$ is negligible, then it can be expressed by using the measured branching ratio
$`B(\omega \pi ^+\pi ^{})={\displaystyle \frac{\mathrm{\Gamma }_\rho (m_\omega )}{\mathrm{\Gamma }_\omega }}\left|\epsilon (m_\omega )+{\displaystyle \frac{g_{\omega \pi \pi }^{(0)}}{g_{\rho \pi \pi }^{(0)}}}\right|^2`$ (26)
as follows
$`\text{Re}(\mathrm{\Pi }_{\rho \omega }^{})={\displaystyle \frac{4\pi g_{\rho \gamma }^{(0)}g_{\omega \gamma }^{(0)}}{m_\omega ^2}}+{\displaystyle \frac{g_{\omega \pi \pi }^{(0)}}{g_{\rho \pi \pi }^{(0)}}}(m_\omega ^2m_\rho ^2)+`$
$`+\sqrt{{\displaystyle \frac{\mathrm{\Gamma }_\omega B(\omega \pi ^+\pi ^{})}{\mathrm{\Gamma }_\rho (m_\omega )}}\left|D_\omega (m_\omega )D_\rho (m_\omega )\right|^2\left[{\displaystyle \frac{g_{\rho \pi \gamma }^{(0)}g_{\omega \pi \gamma }^{(0)}q_{\pi \gamma }^3(m_\omega )+g_{\rho \eta \gamma }^{(0)}g_{\omega \eta \gamma }^{(0)}q_{\eta \gamma }^3(m_\omega )}{3}}+{\displaystyle \frac{g_{\omega \pi \pi }^{(0)}}{g_{\rho \pi \pi }^{(0)}}}m_\omega \mathrm{\Gamma }_\omega \right]^2}`$ (27)
Equation (21) can be rewritten as follows
$`A_{\omega \pi ^+\pi ^{}}+A_{\rho \pi ^+\pi ^{}}=\sqrt{{\displaystyle \frac{3}{2}}}{\displaystyle \frac{1}{\alpha }}{\displaystyle \underset{V=\omega ,\rho }{}}{\displaystyle \frac{\mathrm{\Gamma }_Vm_V^3\text{ }\sqrt{m_V\sigma (V\pi ^+\pi ^{})}}{D_V(s)}}{\displaystyle \frac{f_{V\pi \pi }(s)}{\sqrt{q_\pi (m_V)}}},`$ (28)
where
$$f_{V\pi \pi }(s)=\frac{r_{V\pi \pi }(s)}{r_{V\pi \pi }(m_V)},\text{ }$$
and
$$r_{\rho \pi \pi }(s)=1\frac{g_{\gamma \omega }^{(0)}}{g_{\gamma \rho }^{(0)}}\epsilon (s),\text{ }r_{\omega \pi \pi }(s)=\epsilon (s)+\frac{g_{\omega \pi \pi }^{(0)}}{g_{\rho \pi \pi }^{(0)}}$$
The theoretical value of the phase $`\varphi _{\rho \omega }`$ can be calculated from the above given expressions: $`\varphi _{\rho \omega }=\mathrm{arg}(f_{\omega \pi \pi }(m_\omega ))\mathrm{arg}(f_{\rho \pi \pi }(m_\rho ))101^{}`$. The phase $`\varphi _{\rho \omega }`$ almost does not depend on energy. In this calculation we assumed that the $`\omega \pi ^+\pi ^{}`$ transition proceeds only via the $`\rho \omega `$ mixing, that is $`g_{\omega \pi \pi }^{(0)}=0`$. In order to determine the $`g_{\rho \pi \pi }^{(0)}`$, $`g_{\gamma V}^{(0)}`$ and $`g_{VP\gamma }^{(0)}`$ coupling constants, the corresponding measured decay widths were used.
### IV.2 Fit to the experimental data
The $`\rho ^{}`$ and $`\rho ^{\prime \prime }`$ parameters were determined from the fit to the $`e^+e^{}\pi ^+\pi ^{}`$ cross section measured at the energy region $`\sqrt{s}<2400`$ MeV by OLYA and DM2 detectors olya ; dm2 , together with the isovector part of the $`e^+e^{}\pi ^+\pi ^{}`$ cross section calculated by assuming the CVC hypothesis from the spectral function of the $`\tau ^{}\pi ^{}\pi ^0\nu _\tau `$ decay measured by CLEO II cleo2 :
$`\sigma _{\pi \pi }(m_i)={\displaystyle \frac{4(\pi \alpha )^2}{m_i}}{\displaystyle \frac{B(\tau \pi \pi ^0\nu _\tau )}{B(\tau e\overline{\nu }_e\nu _\tau )}}{\displaystyle \frac{m_\tau ^8}{12\pi |V_{ud}|^2}}{\displaystyle \frac{1}{S_{EW}}}{\displaystyle \frac{1}{m_i(m_\tau ^2m_i^2)^2(m_\tau ^2+2m_i^2)}}{\displaystyle \frac{1}{N}}{\displaystyle \frac{N_i}{\mathrm{\Delta }m_i}},`$ (29)
where $`m_i`$ is the central value of the $`\pi ^{}\pi ^0`$ pair invariant mass for the $`i`$-th bin, $`\mathrm{\Delta }m_i`$ is the bin width, $`N_i`$ is the number of entries in the $`i`$-th bin, $`N`$ is the total number of entries, $`|V_{ud}|`$ is the CKM matrix element, $`S_{EW}=1.0194`$ is the radiative correction aleph ; cleo2 ; radtau .
The obtained $`\rho ^{}`$ and $`\rho ^{\prime \prime }`$ parameters were used in the fitting to the SND data (Table 3, Fig.21). The free parameters of the fit were $`m_\rho `$, $`\mathrm{\Gamma }_\rho `$, $`\sigma (\rho \pi ^+\pi ^{})`$, $`\sigma (\omega \pi ^+\pi ^{})`$, $`\varphi _{\rho \omega }`$ and $`\sigma (\rho ^{}\pi ^+\pi ^{})`$. The first fit was performed with $`\sigma (\rho ^{\prime \prime }\pi ^+\pi ^{})`$, $`\rho ^{}`$ and $`\rho ^{\prime \prime }`$ masses and widths fixed at the values obtained from the fit to the CLEO II and DM2 data. The second and third fits were done without the $`\rho ^{\prime \prime }`$ meson. The $`\rho ^{}`$ mass and width were fixed by using results of the fit to the CLEO II and DM2 data (the second variant in the Table 3) and to the OLYA data (the third variant in the Table 3). The values of the $`\rho `$ and $`\omega `$ parameters exhibit a rather weak model dependence.
## V Discussion.
The comparison of the $`e^+e^{}\pi ^+\pi ^{}`$ cross section obtained in SND experiment with other results bena ; quen ; olya ; kmd2 ; kloe is shown in Fig.23,23,25 and 25. In the energy region $`\sqrt{s}<600`$ MeV all experimental data are in agreement (Fig.23). Above 600 MeV the OSPK(ORSAY-ACO)bena and DM1 quen points lay about 10 % lower than the SND ones (Fig.23). The SND cross section exceeds the OLYA and CMD measurements olya by $`6\pm 1`$ % in this energy region (Fig.25). The systematic error of OLYA measurement is 4 % and the OLYA data agree with the SND result. The systematic uncertainty of CMD result is 2 %, so the difference between the SND and CMD results is about 2.5 of joint systematic error. At the same time the SND and CMD data below 600 MeV agree well (Fig.23). The average deviation between CMD2 kmd2 and SND data is $`1.4\pm 0.5`$ %, the systematic inaccuracies of these measurements are 0.6 % and 1.3 % respectively. In the KLOE experiment at $`\varphi `$-factory DAF$`\mathrm{\Phi }`$NE the form factor $`|F_\pi (s)|^2`$ was measured by using โradiative returnโ method with systematic error of 0.9 % kloe . In Ref.kloe the bare form factor is listed. So in order to compare the KLOE result with the SND one, the form factor was appropriately dressed by us. The results of this comparison are shown in Fig.25. The KLOE measurement is in conflict with the SND result as well as with the CMD2 one.
The $`\rho `$-meson parameters $`m_\rho `$, $`\mathrm{\Gamma }_\rho `$, $`\sigma (\rho \pi ^+\pi ^{})`$ were determined from study of the $`e^+e^{}\pi ^+\pi ^{}`$ cross section. The $`\rho `$ meson mass and width were found to be
$$m_\rho =774.9\pm 0.04\pm 0.05\text{ MeV},$$
$$\mathrm{\Gamma }_\rho =146.5\pm 0.8\pm 1.5\text{ MeV}.$$
The systematic errors is related to the accuracy of the collider energy determination, to the model uncertainty and to the error of the cross section determination. The $`\rho `$-meson parameters were studied in other $`e^+e^{}`$ experiments by using the processes $`e^+e^{}\pi ^+\pi ^{}`$ kmd2 ; olya , $`e^+e^{}\rho \pi \pi ^+\pi ^{}\pi ^0`$ kloe3pi ; dplphi98 and the $`\tau ^{}\pi ^{}\pi ^0\nu _\tau `$ decay cleo2 ; aleph . The SND results are in agreement with these measurements as is shown in Fig.27 and 27.
The parameter $`\sigma (\rho \pi ^+\pi ^{})`$ was found to be
$$\sigma (\rho \pi ^+\pi ^{})=1220\pm 7\pm 16\text{ nb},$$
which corresponds to
$$B(\rho e^+e^{})\times B(\rho \pi ^+\pi ^{})=(4.991\pm 0.028\pm 0.066)\times 10^5,$$
$$\mathrm{\Gamma }(\rho e^+e^{})=7.31\pm 0.021\pm 0.11\text{ keV}.$$
The systematic error includes the systematic uncertainties in the cross section measurement and the model dependence. A comparison of the $`\mathrm{\Gamma }(\rho e^+e^{})`$ obtained in this work with other experimental results kmd2 ; olya ; bena and with the PDG world average pdg is shown in Fig.29. The SND result exceeds all previous measurements. It differs by about 1.5 standard deviations from the CMD2 measurement kmd2 and by 2 standard deviations from the PDG world average pdg . The difference of the $`\rho `$-meson leptonic widths obtained by SND and CMD2 should be attributed mainly to the difference in the total widths of the $`\rho `$-meson rather then to the difference in the cross section values. The value $`\sigma (\rho \pi ^+\pi ^{})=1198`$ nb, which can be obtained by using CMD2 cross section data reported in Ref.kmd2 , agrees with the SND result within the measurements errors.
The parameter $`\sigma (\omega \pi ^+\pi ^{})`$ was found to be
$$\sigma (\omega \pi ^+\pi ^{})=29.9\pm 1.2\pm 1.0\text{ nb},$$
which corresponds to
$$B(\omega e^+e^{})\times B(\omega \pi ^+\pi ^{})=(1.247\pm 0.062\pm 0.042)\times 10^6.$$
The systematic error is related to the model dependence, to the error of the cross section determination and to the accuracy of the collider energy determination. In the previous studies of the $`e^+e^{}\pi ^+\pi ^{}`$ reaction the relative probability of the $`\omega \pi ^+\pi ^{}`$ decay was also reported. The comparison of $`B(\omega \pi ^+\pi ^{})=0.0175\pm 0.0011`$ obtained by using the SND data and the PDG value of the $`\omega e^+e^{}`$ decay width pdg with the results of other experiments is shown in Fig.29. The SND result is the most precise.
The phase $`\varphi _{\rho \omega }`$ was found to be
$$\varphi _{\rho \omega }=113.5\pm 1.3\pm 1.7\text{ degree}.$$
This value differs by 6 standard deviations from $`101^{}`$ expected under assumption that the $`\omega \pi ^+\pi ^{}`$ transition proceeds through the $`\rho \omega `$ mixing mechanism. If instead of the phase $`\varphi _{\rho \omega }`$, the ratio $`g_{\omega \pi \pi }^{(0)}/g_{\rho \pi \pi }^{(0)}`$ is the free parameter of the fit it follows that
$$\frac{g_{\omega \pi \pi }^{(0)}}{g_{\rho \pi \pi }^{(0)}}=0.11\pm 0.01.$$
This ratio corresponds to the too large direct transition width $`\mathrm{\Gamma }^{(0)}(\omega \pi ^+\pi ^{})=1.82\pm 0.33`$ MeV, while the natural expectation is $`\mathrm{\Gamma }^{(0)}(\omega \pi ^+\pi ^{})\alpha ^2\mathrm{\Gamma }_\rho 8`$ keV. Let us note, that the analysis of the OLYA and CMD2 data kmd2 ; olya give the similar values of the $`\varphi _{\rho \omega }`$ phase. This result can point out that the considerable direct transition $`\omega \pi ^+\pi ^{}`$ exists. On the other hand this discrepancy can be attributed also to inadequacies of the applied theoretical model.
The comparison of the phase $`\mathrm{arg}(A_{\rho \pi ^+\pi ^{}}+A_{\rho ^{}\pi ^+\pi ^{}})`$ with the $`\pi \pi `$ scattering phase in the P-wave pwa1 ; pwa2 is shown in Fig.31. These phases must be equal in the purely elastic scattering region. The agreement is satisfactory, in any case in the energy region $`\sqrt{s}m_\rho `$ no essential difference is observed.
The comparison of the $`e^+e^{}\pi ^+\pi ^{}`$ cross section, obtained under the CVC hypothesis from the $`\tau `$ spectral function of the $`\tau ^{}\pi ^{}\pi ^0\nu _\tau `$ decay cleo2 ; aleph with isovector part of the cross section measured in this work is shown in Fig.31. The cross section obtained by SND was undressed from the vacuum polarization and the contribution from the $`\omega \pi ^+\pi ^{}`$ decay was excluded. The cross section calculated from the $`\tau `$ spectral function was multiplied by the coefficient which takes into account the difference of the $`\pi ^\pm `$ and $`\pi ^0`$ masses:
$$\delta =\left(\frac{q_\pi (s)}{q_{\pi ^\pm }(s)}\right)^3\frac{|A_{\pi ^+\pi ^{}}(s)|^2}{|A_{\pi ^0\pi ^\pm }(s)|^2},\text{ }q_{\pi ^\pm }(s)=\frac{1}{2\sqrt{s}}\left[(s(m_{\pi ^0}+m_{\pi ^\pm })^2)(s(m_{\pi ^0}m_{\pi ^\pm })^2)\right]^{1/2}.$$
The average deviation of the SND and $`\tau `$ data is about 1.5 %. For almost all energy points this deviation is within the joint systematic error $`1.6\%`$. The 10% difference between $`e^+e^{}`$ and $`\tau `$ data at $`\sqrt{s}>800`$ MeV, which was claimed in Ref.eetau , is absent.
Using the $`\sigma _{\pi \pi }^{pol}(s)`$ cross section (Table 1), the contribution to the anomalous magnetic moment of the muon, due to the $`\pi ^+\pi ^{}(\gamma )`$ intermediate state in the vacuum polarization, was calculated via the dispersion integral:
$$a_\mu (\pi \pi ,390\text{MeV}\sqrt{s}970\text{MeV})=\left(\frac{\alpha m_\mu }{3\pi }\right)^2_{s_{min}}^{s_{max}}\frac{R(s)K(s)}{s^2}๐s,$$
where $`s_{max}=970`$ MeV, $`s_{min}=390MeV`$, $`K(s)`$ is the known kernel and
$$R(s)=\frac{\sigma _{\pi \pi }^{pol}}{\sigma (e^+e^{}\mu ^+\mu ^{})},\text{ }\sigma (e^+e^{}\mu ^+\mu ^{})=\frac{4\pi \alpha ^2}{3s}.$$
The integral was evaluated by using the trapezoidal rule. To take into account the numerical integration errors, the correction method suggested in Ref.aki was applied. As a result we obtained
$$a_\mu (\pi \pi ,390\text{MeV}\sqrt{s}970\text{MeV})=(488.7\pm 2.6\pm 6.6)\times 10^{10}.$$
This is about 70 % of the total hadronic contribution to the anomalous magnetic moment of the muon $`(g2)/2`$.
If the integration is performed for the energy region corresponding to the CMD2 measurements kmd2 , then the result is $`a_\mu (\pi \pi )=(385.6\pm 5.2)\times 10^{10}`$, which is 1.8 % (1 standard deviation) higher than the CMD2 result: $`a_\mu (\pi \pi )=(378,6\pm 3.5)\times 10^{10}`$. So no considerable difference between the SND and CMD2 results is observed.
## VI Conclusion
The cross section of the process $`e^+e^{}\pi ^+\pi ^{}`$ was measured in the SND experiment at the VEPP-2M collider in the energy region $`390<\sqrt{s}<980`$ MeV with accuracy 1.3 % at $`\sqrt{s}420`$ MeV and 3.4 % at $`\sqrt{s}<420`$ MeV. The measured cross section was analyzed in the framework of the generalized vector meson dominance model. The following $`\rho `$-meson parameters were obtained: $`m_\rho =774.9\pm 0.4\pm 0.5`$ MeV, $`\mathrm{\Gamma }_\rho =146.5\pm 0.8\pm 1.5`$ MeV and $`\sigma (\rho \pi ^+\pi ^{})=1220\pm 7\pm 16`$ nb. The parameters of the $`G`$-parity suppressed process $`e^+e^{}\omega \pi ^+\pi ^{}`$ were measured with high precision. The measured value $`\sigma (\omega \pi ^+\pi ^{})=29.9\pm 1.4\pm 1.0`$ nb corresponds to the relative probability $`B(\omega \pi ^+\pi ^{})=1.75\pm 0.11\%`$. The relative interference phase between the $`\rho `$ and $`\omega `$ mesons was found to be $`\varphi _{\rho \omega }=113.5\pm 1.3\pm 1.7`$ degree. This result is in conflict with the naive expectation from the $`\rho \omega `$ mixing $`\varphi _{\rho \omega }=101^{}`$. The SND result agrees with the cross section calculated from the $`\tau `$ spectral function data within the accuracy of the measurements. Using measured cross section, the contribution to the anomalous magnetic moment of the muon due to the $`\pi ^+\pi ^{}(\gamma )`$ intermediate state in the vacuum polarization was calculated: $`a_\mu (\pi \pi ,390\text{MeV}\sqrt{s}970\text{MeV})=(488.7\pm 2.6\pm 6.6)\times 10^{10}.`$
###### Acknowledgements.
The authors are grateful to N.N.Achasov for useful discussions. The work is supported in part by grants Sci.School-1335.2003.2, RFBR 04-02-16181-a, 04-02-16184-a, 05-02-16250-a.
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# Sur le spectre du laplacien des fibrรฉs en tores qui sโeffondrent
## Introduction
Lโune des questions centrales de lโanalyse sur les variรฉtรฉs est lโestimation du spectre du laplacien agissant sur les fonctions dโune variรฉtรฉ compacte en fonction dโinvariants gรฉomรฉtriques, et fait lโobjet dโun vaste littรฉrature. Des rรฉsultats importants ont en particulier รฉtรฉ obtenus concernant la minoration de la premiรจre valeur propre du laplacien. J. Cheeger en a donnรฉ une estimation en fonction dโune constante isopรฉrimรฉtrique:
###### Dรฉfinition 1
Soit $`(M,g)`$ une variรฉtรฉ riemannienne compacte connexe de dimension $`n`$. On dรฉfinit la constante de Cheeger $`h`$ de la variรฉtรฉ $`(M,g)`$ par
$$h(M,g)=\underset{S}{inf}\frac{VolS}{\mathrm{min}(VolM_1,VolM_2)},$$
$`S`$ parcourt lโensemble des sous-variรฉtรฉs fermรฉes de dimension $`n1`$ de $`M`$ qui partitionnent $`M`$ en deux sous-variรฉtรฉs $`M_1`$ et $`M_2`$ dont le bord est $`S`$.
###### Thรฉorรจme 2 (\[Ch70\])
Si $`(M,g)`$ une variรฉtรฉ compacte connexe, alors $`\lambda _1(M,g)\frac{h(M,g)^2}{4}`$, oรน $`\lambda _1(M,g)`$ est la premiรจre valeur propre du laplacien agissant sur les fonctions de $`M`$.
Si on ajoute une hypothรจse sur la courbure de Ricci, on peut encadrer la premiรจre valeur propre en fonction de la constante de Cheeger :
###### Thรฉorรจme 3 (\[Bu82\])
Soit $`(M,g)`$ une variรฉtรฉ compacte connexe de dimension $`n`$. Si $`Ric(M,g)>(n1)ag`$ alors $`\lambda _1(M,g)\tau (n)(\sqrt{a}h(M,g)+h(M,g)^2)`$, oรน $`a`$ est une constante positive et $`\tau (n)`$ une constante ne dรฉpendant que de la dimension.
Une autre approche est de minorer cette valeur propre en fonction de bornes sur la courbure et le diamรจtre de la variรฉtรฉ. Un premier rรฉsultat a รฉtรฉ obtenu par M. Gromov:
###### Thรฉorรจme 4 (\[Gr80\])
Soit $`a`$ et $`d`$ deux rรฉels strictement positifs et $`n`$ un entier. Il existe une constante $`c(n,a,d)>0`$ telle que si $`(M,g)`$ est une variรฉtรฉ riemannienne de dimension $`n`$ dont le diamรจtre et la courbure de Ricci vรฉrifient $`diam(M,g)d`$ et $`Ric(M,g)ag`$, alors
$$\lambda _1(M,g)c.$$
Il existe dโautres contributions ร ce type de minoration (\[LY80\], \[BBG85\]), et lโestimation du thรฉorรจme 4 a รฉtรฉ affinรฉe. En particulier, Cheng et Zhou donnent dans \[CZ95\] :
$$c=\frac{\pi ^2}{d^2}e^{\frac{1}{2}c_n\sqrt{ad^2}},\text{ avec }c_n=\mathrm{max}\{\sqrt{n1},\sqrt{2}\}.$$
Le spectre du laplacien de Hodge-de Rham agissant sur les formes diffรฉrentielles a รฉtรฉ moins รฉtudiรฉ, et la question se pose naturellement de savoir si les thรฉorรจmes 2 et 4 se gรฉnรฉralisent aux formes. Rappelons que si $`(M,g)`$ est une variรฉtรฉ connexe compacte orientable, le laplacien est dรฉfini sur lโensemble $`\mathrm{\Omega }^p(M)`$ des formes diffรฉrentielles de degrรฉ $`p`$ de $`M`$, $`0pn`$ par
$$\mathrm{\Delta }^p=\mathrm{d}\delta +\delta \mathrm{d},$$
$`\mathrm{d}`$ est la diffรฉrentielle extรฉrieure et $`\delta `$ son adjoint formel pour la norme $`L^2`$ sur $`\mathrm{\Omega }^p(M)`$. La dimension du noyau de $`\mathrm{\Delta }^p`$ est รฉgale au $`p`$-iรจme nombre de Betti de $`M`$, et son spectre est discret. On notera
$$0=\lambda _{p\mathrm{,0}}(M,g)<\lambda _{p\mathrm{,1}}(M,g)\lambda _{p\mathrm{,2}}(M,g)\mathrm{}\lambda _{p,k}(M,g)\mathrm{}$$
les valeurs propres de $`\mathrm{\Delta }^p`$, en rรฉpรฉtant les valeurs propres non nulles sโil y a multiplicitรฉ. Si $`p=0`$, on retrouve le laplacien agissant sur les fonctions.
Lโรฉtude du laplacien de Hodge-de Rham montre que le comportement du spectre pour $`1pn1`$ est diffรฉrent du spectre du laplacien agissant sur les fonctions. En particulier, le thรฉorรจme 4 ne se gรฉnรฉralise pas aux formes diffรฉrentielles, mรชme avec lโhypothรจse plus forte de courbure sectionnelle bornรฉe. En effet, C. Colbois et G. Courtois ont donnรฉ dans \[CC90\] des exemples de variรฉtรฉs admettant des petites valeurs propres, cโest-ร -dire quโon peut faire varier leur mรฉtrique de sorte que la ou les premiรจres valeurs propres non nulles tendent vers zรฉro, le diamรจtre et la courbure sectionnelles restant bornรฉs.
Le thรฉorรจme 4 prรฉcรฉdent soulรจve la question de savoir ร quelles conditions une variรฉtรฉ admet une petite valeur propre ร diamรจtre et courbure bornรฉs. Un dรฉbut de rรฉponse a รฉtรฉ donnรฉ par B. Colbois et G. Courtois dans \[CC90\] en montrant quโon peut obtenir pour les formes diffรฉrentielles un rรฉsultat semblable au thรฉorรจme 4 si on se donne une hypothรจse supplรฉmentaire sur le rayon dโinjectivitรฉ de la variรฉtรฉ:
###### Thรฉorรจme 5
Pour tous rรฉels $`a`$, $`d`$, et $`r`$ strictement positifs et tout entier $`n`$, il existe une constante $`c^{}(n,a,d,r)>0`$ telle que si $`(M,g)`$ est une variรฉtรฉ riemannienne compacte connexe de dimension $`n`$ vรฉrifiant $`diam(M,g)d`$, $`|K(M,g)|a`$ et $`inj(M,g)r`$, oรน $`K(M,g)`$ et $`inj(M,g)`$ dรฉsignent respectivement la courbure sectionnelle et le rayon dโinjectivitรฉ de $`M`$, alors
$$\lambda _{p\mathrm{,1}}(M,g)c^{},$$
pour tout $`p`$.
Ce rรฉsultat a รฉtรฉ amรฉliorรฉ par S. Chanillo et F. Trรจves, qui donnent une valeur explicite de la constante $`c^{}`$ du thรฉorรจme 5, et en mettant en รฉvidence le rรดle du rayon dโinjectivitรฉ dans cette constante (\[CT97\]):
$$\lambda _{p\mathrm{,1}}(M,g)c^{\prime \prime }(n,a,d)r^{4n^2+4n2}.$$
(6)
Une consรฉquence immรฉdiate est que si $`(M_i^n,g_i)`$ est une suite de variรฉtรฉs riemanniennes de dimension $`n`$ de diamรจtres et courbures uniformรฉment bornรฉs telle que $`\lambda _{p\mathrm{,1}}(M_i,g_i)`$ tende vers zรฉro pour un certain $`p`$, $`1pn`$, alors $`inj(M_i,g_i)`$ tend aussi vers zรฉro (remarque: si le diamรจtre et la courbure sectionnelle sont bornรฉs, minorer le rayon dโinjectivitรฉ est รฉquivalent ร minorer le volume). Plus simplement, si on se donne une variรฉtรฉ compacte $`M`$ et quโon fait varier sa mรฉtrique ร diamรจtre et courbure bornรฉe de sorte que $`\lambda _{p\mathrm{,1}}`$ tende vers zรฉro, alors son rayon dโinjectivitรฉ tend vers zรฉro, cโest-ร -dire quโelle sโeffondre. Cette condition a des consรฉquences topologiques importantes. En particulier, si $`(M,g_i)`$ tend pour la distance de Gromov-Haussdorf vers une variรฉtรฉ de dimension infรฉrieure, K. Fukaya a montrรฉ (entre autres choses) dans \[Fu87b\] et \[Fu89\] que $`M`$ possรจde une structure de fibrรฉ dont la base est la variรฉtรฉ limite:
###### Thรฉorรจme 7
Soit $`(M_i,g_i)`$ une suite de variรฉtรฉs compactes de dimension $`n`$ et $`(N,h)`$ une variรฉtรฉ riemannienne compacte de dimension $`m<n`$. Si pour tout $`i`$ la courbure sectionnelle de $`M_i`$ vรฉrifie $`|K(M,g_i)|1`$, et si $`(M,g_i)`$ converge vers $`(N,h)`$ pour la distance de Gromov-Hausdorff, alors pour tout $`i`$ suffisamment grand il existe une fibration $`\pi _i:M_iN`$ dont la fibre est une infranilvariรฉtรฉ.
Cependant, il nโest pas suffisant que le rayon dโinjectivitรฉ tende vers zรฉro pour pour que la premiรจre valeur propre non nulle du laplacien tende vers zรฉro. On peut par exemple considรฉrer le produit riemannien dโune variรฉtรฉ $`(N,h)`$ quelconque munie dโune mรฉtrique fixรฉe avec un tore plat dont on fait tendre le diamรจtre vers zรฉro: le rayon dโinjectivitรฉ du produit tend vers zรฉro, mais son spectre est la rรฉunion des spectres de $`(N,h)`$ et du tore, et donc ne contient pas de petites valeurs propres. Ce la motive la
###### Question 8
ร quelles conditions une variรฉtรฉ qui sโeffondre admet-elle une โ ou plusieurs โ petites valeurs propres?
Par ailleurs, la minoration (6) nโest *a priori* pas optimale, en particulier en ce qui concerne lโexposant du rayon dโinjectivitรฉ. Il se pose donc aussi la
###### Question 9
Peut-on amรฉliorer la minoration de la premiรจre valeur propre non nulle du laplacien de Hodge-de Rham?
Des rรฉponses prรฉcises aux questions 8 et 9 ont รฉtรฉ apportรฉes dans des situations topologiques ou gรฉomรฉtriques particuliรจres. Cโest le cas des limites adiabatiques : si $`(M,g)`$ est une variรฉtรฉ riemannienne compacte, $`A`$ une distribution de sous-espaces de $`TM`$ et $`B`$ la distribution orthogonale ร $`A`$, on peut รฉcrire la mรฉtrique $`g`$ sous la forme $`g=g_Ag_B`$$`g_A`$ et $`g_B`$ sont des mรฉtriques sur $`A`$ et $`B`$ respectivement, et dรฉfinir sur $`M`$ la famille de mรฉtriques
$$g_t=t^2g_Ag_B.$$
La limite de $`(M,g_t)`$ quand $`t0`$ est appelรฉe \<\< limite adiabatique \>\>. Dans le cas oรน $`M`$ est un fibrรฉ riemannien, la distribution $`A`$ est la distribution verticale $`T^VM`$, $`B`$ la distribution horizontale $`T^HM`$, et on fait tendre le diamรจtre de la fibre vers zรฉro par des homothรฉties de rapport $`t`$. Pour un tel effondrement, Mazzeo et Melrose ont montrรฉ (\[MM90\]) que le nombre de valeurs propres de lโordre de $`t^2`$ quand $`t`$ tend vers zรฉro peut se calculer ร lโaide de la suite spectrale de Leray du fibrรฉ (notons que Forman, รlvarez et Kordyukov ont gรฉnรฉralisรฉ ce rรฉsultat aux feuilletages riemanniens dans \[Fo95\] et \[ALK00\]). Cependant, dans une situation de limite adiabatique, la courbure nโest en gรฉnรฉral pas bornรฉe, en particulier si la fibre nโest pas un tore. Dโautre part, on peut en gรฉnรฉral effondrer le fibrรฉ sur sa base autrement que par homothรฉtie de la fibre.
Par ailleurs, B. Colbois et G. Courtois ont รฉtudiรฉ dans \[CC00\] le cas oรน $`M`$ est une variรฉtรฉ de dimension $`n+1`$ qui tend pour la distance de Gromov-Hausdorff vers une variรฉtรฉ $`(N,h)`$ de dimension $`n`$ et donnent des estimations plus prรฉcises de la ou les premiรจres valeurs propres non nulles en fonction de la topologie de $`M`$ et de la gรฉomรฉtrie de lโeffondrement. Selon le thรฉorรจme 7, on sait que $`M`$ est nรฉcessairement un fibrรฉ en cercle sur $`N`$, et lโhypothรจse dโorientabilitรฉ sur $`M`$ assure que ce fibrรฉ est principal. La topologie de ce fibrรฉ est caractรฉrisรฉe par sa classe dโEuler $`[e]H^2(N,)`$ (cf. \[BT82\], p. 72). En notant $`e(M)`$ le reprรฉsentant harmonique de la classe de cohomologie $`[e]`$ pour la mรฉtrique $`h`$ sur $`N`$, on a alors:
###### Thรฉorรจme 10
Soit $`a`$ et $`d`$ deux rรฉels strictement positifs et $`(N,h)`$ une variรฉtรฉ riemannienne de dimension $`n`$. Il existe des constantes $`\epsilon _0(n,a,d,(N,h))>0`$ et $`C_i(n,a,d,(N,h))>0`$ pour $`i=\mathrm{1,2,3}`$ telles que si $`(M,g)`$ est une variรฉtรฉ riemanienne de dimension $`n+1`$ vรฉrifiant $`diam(M,g)d`$, $`|K(M,g)|a`$ et $`\epsilon =d_{GH}((M,g),(N,h))\epsilon _0`$, alors on a, pour $`1pn`$,
1. $`\lambda _{p,m_p+1}(M,g)C_1`$;
2. Si $`e0`$, alors $`C_2e(M)_2^2\epsilon ^2\lambda _{\mathrm{1,1}}(M,g)C_3e(M)_2^2\epsilon ^2`$;
3. Si $`dimH^2(N,)=1`$, alors
$$C_2e(M)_2^2\epsilon ^2\lambda _{p,k}(M,g)C_3e(M)_2^2\epsilon ^2\mathrm{pour}1km_p;$$
avec $`m_p=b_p(N)+b_{p1}(N)b_p(M)`$.
Remarque 11. Si $`e=0`$, par exemple si le fibrรฉ est trivial, on a $`m_p=0`$ pour tout $`p`$. Le thรฉorรจme se rรฉduit alors ร 10.1. Dโautre part, si $`e0`$ alors $`m_1=1`$. $`\lambda _{\mathrm{1,1}}(M,g)`$ est donc la seule petite valeur propre pour les $`1`$-formes.
Remarque 12. B. Colbois et G. Courtois montrent aussi dans \[CC00\] que si $`dimH^2(N,)1`$, on ne peut pas trouver de constante $`C_2`$ vรฉrifiant 10.3 qui soit indรฉpendante de la topologie de $`M`$.
Remarque 13. Le thรฉorรจme isole les rรดles de la topologie (par lโintermรฉdiaire de $`e(M)`$) et du rayon dโinjectivitรฉ (qui est de lโordre de $`\epsilon `$ quand $`\epsilon `$ tend vers zรฉro) dans la minoration de la premiรจre valeur propre non nulle du laplacien. On voit donc que dans la situation dโun fibrรฉ en cercle qui sโeffondre, si le laplacien admet une petite valeur propre, elle se comporte asymptotiquement comme le carrรฉ du rayon dโinjectivitรฉ.
Dans des situations topologiques et gรฉomรฉtriques plus gรฉnรฉrales, on ne connait pas de rรฉsultat semblable. Cependant, en ce qui concerne la question 8, J. Lott a apportรฉ une contribution importante dans \[Lo02b\] et \[Lo02a\] en dรฉfinissant un opรฉrateur limite pour le laplacien, cโest-ร -dire un opรฉrateur dont le spectre est la limite du spectre de laplacien de Hodge-de Rham quand la variรฉtรฉ sโeffondre, le nombre de valeurs propres petites ou nulles รฉtant alors donnรฉ par la multiplicitรฉ de la valeur propre nulle dans le spectre de lโopรฉrateur limite. Dans le chapitre 1, nous exposerons la construction de cet opรฉrateur limite, ainsi que les conditions nรฉcessaires ร lโexistence de petites valeurs propres que Lott en dรฉduit.
Les chapitres 2 et 3 seront consacrรฉs ร lโรฉtude, pour des fibrรฉs munis dโune structure homogรจne โ construits plus prรฉcisรฉment comme quotients dโun groupe de Lie rรฉsoluble $`G`$ par un rรฉseau cocompact $`\mathrm{\Gamma }`$ โ, du spectre du laplacien $`\mathrm{\Delta }_{inv}^p`$ restreint ร lโespace de dimension finie $`\mathrm{\Omega }^p(M)^G`$ des formes diffรฉrentielles invariantes lors dโeffondrements par des mรฉtriques homogรจnes, afin de mettre en รฉvidence comment varie prรฉcisรฉment le nombre de petites valeurs propres en fonction de la topologie du fibrรฉ et de la gรฉomรฉtrie de lโeffondrement. En effet, les petites valeurs propres de $`\mathrm{\Delta }_{inv}^p`$ sont aussi petites valeurs propres de $`\mathrm{\Delta }^p`$, et dans le cas oรน le groupe $`G`$ est nilpotent, J. Lott a montrรฉ rรฉciproquement que la recherche de petites valeurs propres de $`\mathrm{\Delta }`$ se ramรจne ร lโรฉtude de $`\mathrm{\Delta }_{inv}^p`$:
###### Proposition 14 (\[Lo02b\])
Il existe des constantes $`a(n)`$, $`a^{}(n)`$ et $`c(n)`$ strictement positives telles que si $`M`$ est une infranilvariรฉtรฉ de dimension $`n`$ munie dโune mรฉtrique homogรจne pour laquelle $`R_{\mathrm{}}diam(M)^2a^{}`$, oรน $`R_{\mathrm{}}`$ est la norme du tenseur de courbure, et si $`\alpha `$ est une forme propre du laplacien sur $`M`$ dont la valeur propre $`\lambda `$ vรฉrifie $`\lambda <adiam(M)^2cR_{\mathrm{}}`$, alors $`\alpha `$ est invariante.
Cependant, ce rรฉsultat ne se gรฉnรฉralise pas ร tous les groupes rรฉsolubles. On en verra un exemple au paragraphe 2.7.2.
Dans le chapitre 2, nous รฉtudierons des fibrรฉs de fibre $`T^n`$ sur la base $`S^1`$. Le fait quโune variรฉtรฉ qui sโeffondre sur un cercle admette une structure de solvariรฉtรฉ est dรฉjร connu (\[Pe89\], \[Tu97\]). Nous en ferons une construction explicite dans un cas simple: les fibrรฉs considรฉrรฉs seront dรฉfinis comme suspension dโun diffรฉomorphisme linรฉaire de la fibre $`T^n`$ reprรฉsentรฉ par un รฉlรฉment $`ASL_n()`$. De plus, nous ferons lโhypothรจse simplificatrice que $`A`$ admet un logarithme rรฉel. Les principales propriรฉtรฉs du spectre que nous allons mettre en รฉvidence sont donnรฉes par le
###### Thรฉorรจme 15
Soit $`n2`$, $`ASL_n()`$ et $`B\mathrm{M}_n()`$ tels que $`A=\mathrm{exp}(B)`$, $`d`$ la dimension du sous-espace caractรฉristique associรฉ ร la valeur propre $`0`$ de $`B`$ et $`d^{}`$ la dimension de son noyau. Il existe un groupe $`G=G(B)GL_{n+2}()`$ et un rรฉseau $`\mathrm{\Gamma }G`$ tel que $`\mathrm{\Gamma }\backslash G`$ soit homรฉomorphe au fibrรฉ $`M`$ de fibre $`T^n`$ de fibration $`p:MS^1`$ construit par suspension du diffรฉomorphisme $`A`$. Si de plus on suppose que les mรฉtriques sur $`M`$ sont homogรจnes et telles que $`p`$ soit une submersion riemannienne pour une mรฉtrique de volume $`1`$ sur $`S^1`$, alors:
1. $`dimKer\mathrm{\Delta }_{inv}^1=d^{}+1`$ et $`\mathrm{\Delta }_{inv}^1`$ admet $`nd^{}`$ valeurs propres non nulles distinctes ou non.
2. Pour tout $`a>0`$, il existe une constante $`c(B,a)>0`$ telle que pour toute mรฉtrique invariante $`g`$ sur $`M`$ dont la courbure sectionnelle vรฉrifie $`|K(M,g)|<a`$, on a $`\lambda _{1,i}^{inv}(M,g)<c`$, pour tout $`i`$.
3. Si $`dn`$, alors pour tout $`a>0`$, il existe une constante $`c^{}(B,a)>0`$ telle que pour toute mรฉtrique $`g`$ sur $`M`$ la courbure sectionnelle vรฉrifie $`|K(M,g)|<a`$, on a $`\lambda _{1,dd^{}+1}^{inv}(M,g)>c^{}`$.
Si $`d=n`$ et $`d^{}n`$, alors $`G`$ est nilpotent, et il existe une suite de mรฉtriques $`g_\epsilon `$ sur $`M`$ telle que la courbure soit uniformรฉment bornรฉe et $`\lambda _{1,i}^{inv}(M,g_\epsilon )0`$ quand $`\epsilon 0`$, pour tout $`0<ind^{}`$.
Si $`d=d^{}=n`$, alors $`G=^n`$, $`M`$ est un tore et les formes harmoniques sont exactement les formes invariantes.
4. Pour tout $`kdd^{}`$, il existe une famille de mรฉtriques $`g_\epsilon ^k`$ sur $`M`$ de courbure et diamรจtre uniformรฉment bornรฉs par rapport ร $`\epsilon `$ et une constante $`c^{\prime \prime }(B,k)>0`$ telle que $`\lambda _{1,i}^{inv}(M,g_\epsilon ^k)0`$ pour $`ik`$ quand $`\epsilon 0`$, et $`\lambda _{1,k+1}^{inv}(M,g_\epsilon ^k)>c^{\prime \prime }`$ si $`k<n`$.
Remarque 16. Le point 15.4 montre que le nombre de petites valeurs propres peut, quand la topologie le permet, fortement varier avec la gรฉomรฉtrie de lโeffondrement. Cette situation diffรจre donc du cas oรน la fibre est de dimension $`1`$, รฉtudiรฉ par B. Colbois et G. Courtois dans \[CC00\].
Remarque 17. La dรฉmonstration de 15.4 met en รฉvidence le fait que dans le cas dโun effondrement par homothรฉties de la fibre (situation de limite adiabatique), il nโy a pas de petites valeurs propres.
Remarque 18. Une consรฉquence immรฉdiate du rรฉsultat 15.3 est que si tout le spectre de $`\mathrm{\Delta }_{inv}^1`$ tend vers zรฉro, alors le groupe $`G`$ est nilpotent.
Dโautre part, ce thรฉorรจme permet de donner une condition nรฉcessaire et suffisante sur $`B`$ pour lโexistence de petites valeurs propres pour les $`1`$-formes:
###### Corollaire 19
Sous les hypothรจses du thรฉorรจme 15, il existe une suite de mรฉtriques homogรจnes $`g_\epsilon `$ sur $`M`$ telle que $`\lambda _{\mathrm{1,1}}^{inv}(M,g_\epsilon )0`$ quand $`\epsilon 0`$ si et seulement si $`dd^{}`$ (*i.e.* si la rรฉduite de Jordan de $`B`$ a un bloc nilpotent non nul).
Remarque 20. Les rรฉsultats 15 et 19 montrent que le comportement asymptotique du spectre de $`\mathrm{\Delta }_{inv}^1`$ est essentiellement liรฉ ร la nature, de la partie nilpotente de la rรฉduite de Jordan de $`B`$. En effet, le bloc nilpotent est caractรฉrisรฉ par les invariants $`d`$ et $`d^{}`$ de $`B`$.
Dans le cas des $`p`$-formes, $`p2`$, la situation est plus complexe, et en particulier le corollaire 19 nโest pas vrai pour $`n`$ et $`p`$ quelconque (on verra au paragraphe 2.7.1 quโil peut exister une petite valeur propre pour les $`2`$-formes alors que $`B`$ nโa pas de bloc nilpotent). On peut cependant donner une condition nรฉcessaire ร lโexistence de petites valeurs propres, ร savoir que $`B`$ nโest pas semi-simple. En effet, on a le:
###### Thรฉorรจme 21
Sous les hypothรจse du thรฉorรจme 15, si $`B`$ est semi-simple, alors il existe $`c^{\prime \prime }(B,a)>0`$ tel que pour toute mรฉtrique homogรจne $`g`$ sur $`M`$ dont la courbure sectionnelle vรฉrifie $`|K(M,g)|<a`$, on a $`\lambda _{p\mathrm{,1}}^{inv}(M,g)>c^{\prime \prime }`$.
Dans le chapitre 3, nous allons nous intรฉresser ร des fibrรฉs en tore $`T^k`$ sur le tore $`T^2`$ muni dโune structure de fibrรฉ principal. Leur topologie est assez simple. On sait en effet quโun fibrรฉ principal en tore sur le tore peut รชtre muni dโune structure de nilvariรฉtรฉ (\[Pe89\]). On va ici montrer un rรฉsultat plus prรฉcis en utilisant le fait que la base est de dimension 2 : ce fibrรฉ est diffรฉomorphe au produit dโun tore et dโune nilvariรฉtรฉ de dimension $`3`$. Puis on calculera le spectre du laplacien restreint aux formes diffรฉrentielles invariantes en supposant le fibrรฉ muni dโune mรฉtrique homogรจne.
On obtient le rรฉsultat suivant:
###### Thรฉorรจme 22
Soit $`M`$ un fibrรฉ principal non trivial en tore $`T^n`$ sur le tore $`T^2`$ . Alors
1. $`M`$ est une nilvariรฉtรฉ et, si $`n2`$, $`M`$ est homรฉomorphe ร $`N\times T^{n1}`$, oรน $`N`$ est une nilvariรฉtรฉ de dimension $`3`$.
2. Il existe un vecteur $`V`$ vertical (tangent ร la fibre $`T^n`$) tel que si $`M`$ est muni dโune mรฉtrique homogรจne, alors pour tout $`p[1,n+1]`$, $`\mathrm{\Delta }_{inv}^p`$ admet une unique valeur propre non nulle $`\lambda `$. Sa multiplicitรฉ est $`C_n^{p1}`$, et $`\lambda =Vol(B)^2|V|^2`$, oรน $`Vol(B)`$ est le volume de la base du fibrรฉ pour la mรฉtrique induite.
Remarque 23. Le produit du 22.1 nโest pas nรฉcessairement riemannien pour les mรฉtriques considรฉrรฉes. Le spectre ne peut donc pas se dรฉduire directement de la formule de Kรผnneth.
Remarque 24. On voit que contrairement ร la situation du thรฉorรจme 15, un effondrement par homothรฉties de la fibre produit une petite valeur propre, et que $`\lambda `$ est alors proportionnelle au carrรฉ du diamรจtre de la fibre, ร topologie fixรฉe.
Remarque 25. Dans le cas oรน $`n=1`$, la remarque prรฉcรฉdente rejoint les rรฉsultats de B. Colbois et G. Courtois qui รฉtudient dans \[CC00\] le spectre des fibrรฉs en cercles sur des bases quelconques et sans restrictions sur la mรฉtrique. Mais si la dimension de la fibre est plus grande ($`n2`$), un phรฉnomรจne nouveau apparaรฎt: il existe dans ce cas des effondrements du fibrรฉ tels que $`\lambda `$ ne tende pas vers zรฉro. Nous en donnerons des exemples au paragraphe 3.2. Si $`n2`$, le nombre de petites valeurs propres ne dรฉpend donc pas uniquement de la topologie. Cependant, on a pas de libertรฉ sur ce nombre comme en 15.4.
Dans une deuxiรจme partie, du chapitre 4 au chapitre 7, nous allons nous intรฉresser ร la question 9. Le thรฉorรจme 10 ne considรจre que des fibrรฉs dont la fibre est de dimension fixรฉe รฉgale ร $`1`$ ; on peut se demander dans quelle mesure il se gรฉnรฉralise aux fibres de dimension plus grande. Nous allons ici nous intรฉresser plus prรฉcisรฉment aux situation de fibrations principales sโeffondrant sur leur base, la fibre รฉtant alors un tore (ce sont les seules infranilvariรฉtรฉs possรฉdant une structure de groupe). Notre but dans cette seconde partie est dโarriver au rรฉsultat suivant:
###### Thรฉorรจme 26
Soit $`a`$ et $`d`$ deux rรฉels strictement positifs, $`n`$ un entier et $`(N,h)`$ une variรฉtรฉ riemannienne de dimension strictement infรฉrieure ร $`n`$. Il existe des constantes $`\epsilon _0(n,a,d,(N,h))>0`$ et $`C(n,a,d,(N,h))>0`$ telles que si $`(M,g)`$ est une variรฉtรฉ riemannienne de dimension $`n`$ vรฉrifiant $`diam(M,g)d`$, $`|K(M,g)|a`$ et si $`\pi :(M,g)(N,h)`$ est une fibration principale de fibre $`T^k`$ qui soit une $`\epsilon `$-approximation de Hausdorff avec $`\epsilon <\epsilon _0`$, alors
$$\lambda _{\mathrm{1,1}}(M,g)CVol^2(M,g).$$
Remarque 27. On obtient une minoration en fonction du volume de la variรฉtรฉ et pas en fonction du rayon dโinjectivitรฉ, mais avec un exposant รฉgal ร 2, indรฉpendamment de la dimension de la variรฉtรฉ.
Ce rรฉsultat soulรจve deux questions qui restent ouvertes:
###### Question 28
Peut-on obtenir une minoration en fonction du rayon dโinjectivitรฉ avec un exposant indรฉpendant de la dimension?
###### Question 29
Peut-on gรฉnรฉraliser ce rรฉsultat aux $`p`$-formes diffรฉrentielles, pour tout $`p`$ ?
Nous commencerons dans le chapitre 4 par discuter de la topologie des fibrรฉs principaux en tore, dans le but de dรฉterminer un invariant topologique gรฉnรฉralisant la classe dโEuler des fibrรฉs en cercle et qui pourra รชtre utilisรฉ pour contrรดler le spectre.
Dans le chapitre 5, nous รฉtudierons comment, dans le cas dโun fibrรฉ principal en tore $`T^k`$ muni dโune mรฉtrique invariante, on peut se ramener ร lโรฉtude des petites valeurs propres du laplacien ร celle du spectre du laplacien resteint aux formes diffรฉrentielles invariantes par lโaction de $`T^k`$. Dans le cas des formes diffรฉrentielles dโun fibrรฉ en cercle, Colbois et Courtois donnent, avec certaines hypothรจses sur la mรฉtrique, une minoration du spectre sur lโorthogonal des formes invariantes :
###### Thรฉorรจme 30 (\[CC00\])
Soit $`(N^n,h)`$ une variรฉtรฉ compacte, $`S^1M^{n+1}N`$ un $`S^1`$-fibrรฉ principal, et $`g_\epsilon `$ une mรฉtrique $`S^1`$-invariante sur $`M`$ telle que $`|K(M,g_\epsilon )|<a`$, $`diam(M,g_\epsilon )<d`$ et que la longueur des fibres soit infรฉrieure ร $`\epsilon `$. Il existe des constantes $`C=C(n,a,d,(N,h))`$ et $`\rho =\rho (n)`$ telles que toute forme propre de valeur propre $`\lambda \frac{C}{\epsilon ^{\frac{1}{2\rho }}}`$ est $`S^1`$-invariante.
Nous allons dรฉmontrer un rรฉsultat semblable ร celui de Colbois et Courtois, mais sans utiliser dโhypothรจse sur le diamรจtre et la courbure du fibrรฉ, et en donnant une minoration plus prรฉcise du spectre du laplacien restreint ร lโorthogonal des formes invariantes:
###### Thรฉorรจme 31
Soit $`(M,g)`$ un fibrรฉ principal en cercle muni dโune mรฉtrique $`S^1`$-invariante. On note $`l_0`$ le maximum des longueurs des fibres.
Soit $`\lambda `$ une valeur propre de $`\mathrm{\Delta }^p`$. Si $`\lambda <\left({\displaystyle \frac{2\pi }{l_0}}\right)^2`$, alors les formes propres associรฉes sont $`S^1`$-invariantes.
Remarque 32. La constante $`\left({\displaystyle \frac{2\pi }{l_0}}\right)^2`$ du thรฉorรจme est optimale, en ce sens quโil existe des fibrรฉs pour lequelles elle est รฉgale ร une valeur propre associรฉe ร une forme propre non invariante. En effet, $`\left(\frac{2\pi }{l}\right)^2`$ est la premiรจre valeur propre du cercle de longueur $`l`$. Dans le cas dโun fibrรฉ trivial $`M=N\times S^1`$ muni dโune mรฉtrique produit, les formes diffรฉrentielles de la formes $`\alpha f`$$`\alpha `$ est une $`p`$-forme harmonique de $`N`$ et $`f`$ une fonction propre de $`S^1`$ seront des formes propres de $`\mathrm{\Delta }^p`$, de mรชme valeur propre que $`f`$. Comme la fonction $`f`$ nโest pas constante, les formes propres associรฉs ร la valeur propre $`\left(\frac{2\pi }{l}\right)^2`$ ne sont pas toutes invariantes. Comme $`\left(\frac{2\pi }{l}\right)^2`$ est le $`\lambda _{\mathrm{0,1}}`$ du cercle de longueur $`l`$, la constante $`\left(\frac{2\pi }{l_0}\right)^2`$ du thรฉorรจme 31 peut sโinterprรฉter comme la borne infรฉrieure sur lโensemble des fibres de la premiรจre valeur propre du laplacien agissant sur les fonctions de la fibre. Dans le cas des fibrรฉs en tore, on peut montrer un rรฉsultat semblable, faisant intervenir la premiรจre valeur propre du laplacien $`\mathrm{\Delta }^0`$ restreint au tore :
###### Thรฉorรจme 33
Soit $`k^{}`$, $`T^kM\stackrel{\pi }{}N`$ un fibrรฉ en tore $`T^k`$, $`\overline{g}`$ une mรฉtrique invariante sur $`T^k`$ et $`f`$ une fonction sur $`N`$ strictement positive. Supposons que $`M`$ est muni dโune mรฉtrique $`T^k`$-invariante $`g`$ telle que pour tout $`xN`$, la restriction $`\overline{g}_x`$ de $`g`$ ร la fibre $`\pi ^1(x)`$ vรฉrifie $`\overline{g}_xf(x)\overline{g}`$.
Soit $`\lambda `$ une valeur propre du laplacien agissant sur les formes diffรฉrentielles de $`M`$. Si $`\lambda <(\underset{xN}{sup}f(x))^1\lambda _{\mathrm{0,1}}(T^k,\overline{g})`$, alors les formes propres associรฉes sont $`T^k`$-invariantes.
Remarque 34. Comme pour le thรฉorรจme 31, la constante est optimale dans le cas dโun fibrรฉ trivial muni dโune mรฉtrique produit.
On peut dรฉduire du thรฉorรจme 33 une inรฉgalitรฉ en fonction du maximum des diamรจtres des fibres. Cependant, on doit ajouter une hypothรจse sur la mรฉtrique $`g`$. On verra en effet que la donnรฉe dโune borne sur le diamรจtre des fibres ne permet pas de majorer la fonction $`f`$ du thรฉorรจme 33.
###### Corollaire 35
Soit $`k^{}`$, $`T^kM\stackrel{\pi }{}N`$ un fibrรฉ en tore $`T^k`$, et $`\overline{g}`$ une mรฉtrique invariante sur $`T^k`$. Supposons que $`M`$ est muni dโune mรฉtrique $`T^k`$-invariante telle que pour tout $`xN`$, la restriction de la mรฉtrique ร la fibre $`\pi ^1(x)`$ soit proportionnelle ร $`\overline{g}`$.
Soit $`\lambda `$ une valeur propre du laplacien. Si $`\lambda <\left({\displaystyle \frac{\pi }{d_0}}\right)^2`$, oรน $`d_0`$ est le maximum des diamรจtres des fibres pour la mรฉtrique induite par $`g`$, alors les formes propres associรฉes sont $`T^k`$-invariantes.
Remarque 36. La dรฉmonstration des deux thรฉorรจmes met en รฉvidence le fait que si la multiplicitรฉ dโune valeur propre est impaire, alors le sous-espace propre associรฉ contient des formes invariantes.
Le chapitre 6 aura pour but de montrer que pour obtenir le thรฉorรจme 26, on peut se ramener ร une situation gรฉomรฉtrique plus simple. En particulier, on montrera quโune mรฉtrique de courbure et diamรจtre bornรฉs sur le fibrรฉ est proche dโune mรฉtrique invariante pour laquelle les fibres sont totalement gรฉodรฉsiques :
###### Thรฉorรจme 37
Soient $`a`$ et $`d`$ deux rรฉels strictement positifs, $`n`$ un entier et $`(N,h)`$ une variรฉtรฉ riemannienne de dimension strictement infรฉrieure ร $`n`$. Il existe des constantes $`\epsilon _0(n,a,d,(N,h))>0`$, $`\tau (n,a,d,(N,h))>0`$ et $`\tau ^{}(n,a,d,(N,h))>0`$ telles que si $`(M,g)`$ est une variรฉtรฉ riemannienne de dimension $`n`$ vรฉrifiant $`|K(N,h)|a`$, $`|K(M,g)|a`$, $`diam(M,g)d`$ et si $`\pi :(M,g)(N,h)`$ une fibration principale de fibre $`T^k`$ qui soit une $`\epsilon `$-approximation de Hausdorff avec $`\epsilon <\epsilon _0`$, alors il existe des mรฉtriques $`\stackrel{~}{g}`$ et $`\stackrel{~}{h}`$ sur $`M`$ et $`N`$ respectivement et une fibration principale $`\pi ^{}:(M,\stackrel{~}{g})(N,\stackrel{~}{h})`$ telles que
1. Lโaction de $`T^k`$ sur $`(M,\stackrel{~}{g})`$ est isomรฉtrique;
2. Les fibres de la fibration $`\pi ^{}`$ sont totalement gรฉodรฉsiques;
3. $`{\displaystyle \frac{1}{\tau }}g\stackrel{~}{g}\tau g`$ et $`{\displaystyle \frac{1}{\tau }}h\stackrel{~}{h}\tau h`$;
4. La restriction de $`\stackrel{~}{g}`$ ร la fibre est telle que $`diam(\pi ^1(x))\tau ^{}\epsilon `$, pour tout $`xN`$.
Remarque 38. On verra aussi que si lโon suppose que la mรฉtrique $`g`$ sur $`M`$ est $`T^k`$-invariante, alors on peut remplacer dans le thรฉorรจme 37 lโhypothรจse sur la courbure de $`(M,g)`$ par lโhypothรจse plus faible $`K(M,g)a`$.
Enfin, dans le chapitre 7, nous dรฉmontrerons le thรฉorรจme 26 en utilisant les rรฉsultats des chapitres 4 ร 6, et nous discuterons de la possibilitรฉ dโexprimer la constante $`C`$ du thรฉorรจme 26 en fonction dโinvariants gรฉomรฉtriques de $`(N,h)`$.
Une grande partie des chapitres 2 et 3 a รฉtรฉ publiรฉe sous une forme lรฉgรจrement diffรฉrente dans \[Ja03\], et lโarticle \[Ja04\] contient โ entre autres choses โ la dรฉmonstration du thรฉorรจme 26.
## Chapitre 1 Existence de petites valeurs propres du laplacien
### 1.1 Un exemple : la nilvariรฉtรฉ dโHeisenberg de dimension 3
Nous allons commencer par prรฉsenter un exemple de variรฉtรฉ compacte admettant une petite valeur propre. Cet exemple nous sera utile par la suite car il illustre un certain nombre de phรฉnomรจnes.
On considรจre le groupe dโHeisenberg de dimension 3
$$G=\left\{\left(\begin{array}{ccc}1& x& z\\ 0& 1& y\\ 0& 0& 1\end{array}\right),x,y,z\right\}.$$
(1.1)
Cโest un groupe nilpotent, diffรฉomorphe ร $`^3`$, et son centre vรฉrifie $`Z(G)=[G,G]=\{\left(\begin{array}{c}10z\\ 010\\ 001\end{array}\right),z\}`$. On note $`\mathrm{\Gamma }`$ le sous-groupe de $`G`$ formรฉ des matrices ร coefficients entiers. Cโest un sous-groupe discret cocompact de $`G`$, et on dรฉfinit la nilvariรฉtรฉ dโHeisenberg de dimension 3 par $`N=\mathrm{\Gamma }\backslash G`$. Sa topologie peut รชtre dรฉcrite de trois maniรจres :
* par construction cโest une nilvariรฉtรฉ, cโest-ร -dire le quotient dโun groupe nilpotent par un sous-groupe cocompact;
* cโest aussi un fibrรฉ en tore sur le cercle dont les fibres sont, dans la paramรฉtrisation donnรฉe par (1.1), les sous-variรฉtรฉs dโรฉquation $`x=c^{te}`$. Ces sous-variรฉtรฉs de $`G`$ sont diffรฉomorphes ร $`^2`$, et leurs quotients dans $`N`$ sont des tores;
* enfin, cโest un fibrรฉ principal en cercle sur le tore, les fibres รฉtant dรฉfinies comme les orbites de lโaction du centre $`Z(G)`$ sur $`N`$.
Soit $`X`$, $`Y`$ et $`Z`$ les champs de vecteurs invariants ร gauche sur $`N`$ engendrรฉs en $`(\mathrm{0,0,0})`$ par $`/x`$, $`/y`$ et $`/z`$. Ces champs passent au quotient sur $`N`$, et vรฉrifient $`[X,Y]=Z`$ et $`[X,Z]=[Y,Z]=0`$. On va, ร lโaide de ces champs, construire une famille de mรฉtrique pour laquelle le diamรจtre et la courbure seront uniformรฉment bornรฉs: soit $`\alpha `$, $`\beta `$ et $`\gamma `$ trois rรฉels positifs fixรฉs. Pour tout $`\epsilon ]\mathrm{0,1}]`$, on dรฉfinit sur $`N`$ la mรฉtrique $`g_\epsilon `$ comme รฉtant la mรฉtrique invariante ร gauche telle quโen tout point, la base $`(X_\epsilon ,Y_\epsilon ,Z_\epsilon )`$ dรฉfinie par $`X_\epsilon =\epsilon ^\alpha X`$, $`Y_\epsilon =\epsilon ^\beta Y`$ et $`Z_\epsilon =\epsilon ^\gamma Z`$ soit orthonormรฉe. Les crochets de Lie entre les champs de vecteurs de cette base sont
$$[X_\epsilon ,Y_\epsilon ]=Z_\epsilon ^{\gamma \alpha \beta }\text{ et }[X_\epsilon ,Z_\epsilon ]=[Y_\epsilon ,Z_\epsilon ]=0.$$
(1.2)
Comme les paramรจtres $`\alpha `$, $`\beta `$ et $`\gamma `$ sont positifs, les normes de $`X`$, $`Y`$ et $`Z`$ pour la mรฉtrique $`g_\epsilon `$ sont infรฉrieures ร 1, donc le diamรจtre de $`N`$ reste bornรฉ quand $`\epsilon `$ varie. Dโautre part, comme les champs $`X_\epsilon `$, $`Y_\epsilon `$ et $`Z_\epsilon `$ sont invariants, le tenseur de courbure peut sโรฉcrire en fonction des coefficients des crochets de Lie entre ces champs. On impose donc ร $`\alpha `$, $`\beta `$ et $`\gamma `$ de vรฉrifier $`\tau =\gamma \alpha \beta 0`$, de sorte que la courbure reste elle aussi bornรฉe.
Comme la mรฉtrique sur $`N`$ est invariante ร gauche, lโespace des $`1`$-formes diffรฉrentielles invariantes est stable par le laplacien. On va calculer son spectre en restriction ร cet espace. On peut dรฉduire de (1.2) que
$$\mathrm{d}X_\epsilon ^{\mathrm{}}=\mathrm{d}Y_\epsilon ^{\mathrm{}}=0\text{ et }\mathrm{d}Z_\epsilon ^{\mathrm{}}=\epsilon ^\tau X_\epsilon ^{\mathrm{}}Y_\epsilon ^{\mathrm{}},$$
(1.3)
$`U^{\mathrm{}}`$ dรฉsigne la $`1`$-forme duale de $`U`$ pour la mรฉtrique $`g_\epsilon `$, et donc que
$$\delta (X_\epsilon ^{\mathrm{}}Y_\epsilon ^{\mathrm{}})=\epsilon ^\tau Z_\epsilon ^{\mathrm{}}\text{ et }\delta (X_\epsilon ^{\mathrm{}}Z_\epsilon ^{\mathrm{}})=\delta (Y_\epsilon ^{\mathrm{}}Z_\epsilon ^{\mathrm{}})=0.$$
(1.4)
En restriction aux $`1`$-formes invariantes, lโopรฉrateur $`\mathrm{d}\delta `$ est nul. On a donc finalement
$$\mathrm{\Delta }X_\epsilon ^{\mathrm{}}=\mathrm{\Delta }Y_\epsilon ^{\mathrm{}}=0\text{ et }\mathrm{\Delta }Z_\epsilon ^{\mathrm{}}=\epsilon ^{2\tau }Z_\epsilon ^{\mathrm{}}.$$
(1.5)
La forme diffรฉrentielle $`Z_\epsilon ^{\mathrm{}}`$ est donc une forme propre de valeur propre $`\lambda =\epsilon ^{2\tau }`$. On voit que si $`\tau >0`$, par exemple si $`\alpha =1`$, $`\beta =1`$ et $`\gamma =3`$, cette valeur propre tend vers 0. La variรฉtรฉ $`N`$ admet donc une petite valeur propre. On peut noter quโen revanche, si $`\tau =0`$, par exemple si $`\alpha =\beta =1`$ et $`\gamma =2`$, la valeur propre $`\lambda `$ ne tend pas vers zรฉro. La question 8 doit donc รชtre formulรฉe de maniรจre plus prรฉcise:
###### Question 1.6
Comment varie le nombre de petites valeurs propres avec la gรฉomรฉtrie de lโeffondrement ?
### 1.2 Opรฉrateur limite
#### 1.2.1 Le cas des fonctions
Avant dโaborder la construction faite par J. Lott dโun opรฉrateur limite pour le laplacien de Hodge-de Rham, rappelons le rรฉsultat obtenu par K. Fukaya dans le cas des fonctions. Dans \[Fu87a\], Fukaya montre lโexistence โ sous certaines conditions โ dโun opรฉrateur limite pour le laplacien agissant sur les fonctions.
Soit $`(n,d)`$ lโensemble des variรฉtรฉs riemanniennes $`(M^n,g)`$ de dimension $`n`$ telles que leur courbure sectionnelle et leur diamรจtre vรฉrifient respectivement $`|K(M,g)|1`$ et $`diam(M,g)d`$, et $`\lambda _k(M,g)`$ la $`k`$-iรจme valeur propre du laplacien agissant sur les fonctions de $`M`$, les valeurs propres รฉtant rรฉpรฉtรฉes sโil y a multiplicitรฉ. Le problรจme est dโรฉtendre continument pour tout $`k`$ la fonction $`(M,g)\lambda _k(M,g)`$ ร lโadhรฉrence de $`(n,d)`$ dans lโensemble des espaces mรฉtriques compacts. Cโest impossible si on munit cet ensemble de la topologie de Gromov-Hausdorff, mais Fukaya considรจre lโensemble des espaces mรฉtriques mesurรฉs โ *i.e.* muni dโune mesure de probabilitรฉ โ muni de la topologie de Hausdorff mesurรฉe dรฉfinie comme suit: soit $`(X_i,\mu _i)`$ une suite dโespaces mรฉtriques mesurรฉs. Elle converge vers $`(X,\mu )`$ sโil existe des applications mesurables $`\mathrm{\Psi }_i:X_iX`$ et une suite strictement positive $`\epsilon _i`$ vรฉrifiant
1. $`\underset{i\mathrm{}}{lim}\epsilon _i=0`$ ;
2. Lโ$`\epsilon _i`$-voisinage de $`\mathrm{\Psi }_i(X_i)`$ est $`X`$ ;
3. Pour tout $`p,q`$ dans $`X_i`$, on a
$$|d(\mathrm{\Psi }_i(p),\mathrm{\Psi }_i(q))d(p,q)|<\epsilon _i;$$
4. Pour toute fonction $`f`$ sur $`X`$, on a
$$\underset{i\mathrm{}}{lim}f\mathrm{\Psi }_id\mu _i=fd\mu ,$$
cโest-ร -dire que la mesure $`(\mathrm{\Psi }_i)_{}(\mu _i)`$ converge vers $`\mu `$ pour la topologie faible-$``$.
Il obtient alors:
###### Thรฉorรจme 1.7
Soit $`\overline{(n,d)}`$ lโadhรฉrence de $`(n,d)`$ dans lโensemble des espaces mรฉtriques mesurรฉs, en munissant chaque variรฉtรฉ $`(M^n,g)`$ de sa mesure riemannienne normalisรฉe.
1. La fonction $`\lambda _k`$ sโรฉtend continuement ร lโensemble $`\overline{(n,d)}\backslash \{(point\mathrm{,1})\}`$.
2. Pour tout $`(X,\mu )\overline{(n,d)}`$, il existe un opรฉrateur auto-adjoint $`P_{(X,\mu )}`$ agissant sur $`L^2(X,\mu )`$ tel que $`\lambda _k(X,\mu )`$ soit รฉgal ร la $`k`$-iรจme valeur propre de $`P_{(X,\mu )}`$.
3. Supposons que $`lim_i\mathrm{}(M_i,g_i)=(X,\mu )`$. Soit $`\phi _{k,i}`$ une fonction propre normalisรฉe associรฉe ร la valeur propre $`\lambda _k(M_i,g_i)`$ et $`\mathrm{\Lambda }_{k,i}=\{\phi \mathrm{\Psi }_i,\phi L^2(X,\mu ),P_{(X,\mu )}\phi =\lambda _k(X,\mu )\phi \}`$. Alors $`lim_i\mathrm{}d(\mathrm{\Lambda }_{k,i},\phi _{k,i})=0`$.
Remarque 1.8. Dans le cas oรน lโespace limite $`(X,\mu )`$ est une variรฉtรฉ riemannienne, la mesure limite $`\mu `$ et lโopรฉrateur limite $`P_{(X,\mu )}`$ ne sont pas nรฉcessairement la mesure riemanienne ni le laplacien de la variรฉtรฉ, comme le montre lโexemple suivant :
Exemple 1.9. On considรจre le tore $`T^2=\{(s,t),s,tS^1\}`$, $`c:S^1`$ une fonction positive $`C^{\mathrm{}}`$, et la famille de mรฉtriques $`g_\epsilon (c)`$ dรฉfinie sur $`T^2`$ par
$$g_\epsilon (c)=\mathrm{d}s^2\epsilon ^2c(s)^2\mathrm{d}t^2.$$
Si $`fC^{\mathrm{}}(T^2)`$, on a alors
$$\mathrm{\Delta }_{(T^2,g_\epsilon (c))}f(s,t)=c(s)^1\frac{}{s}\left(c(s)\frac{}{s}f(s,t)\right)\epsilon ^2c(s)^2c(s)^2\frac{^2}{t^2}f(s,t).$$
Si on note $`\lambda _k(c)`$ la $`k`$-iรจme valeur propre de lโopรฉrateur $`P_c`$ dรฉfini sur $`S^1`$ par $`P_c(f)(s)=c(s)^1\frac{\mathrm{d}}{\mathrm{d}s}\left(c(s)\frac{\mathrm{d}}{\mathrm{d}s}f(s)\right)`$, alors
$$\underset{\epsilon 0}{lim}\lambda _k(T^2,g_\epsilon (c))=\lambda _k(c),$$
et
$$\underset{\epsilon 0}{lim}(T^2,g_\epsilon (c))=(S^1,\mu )$$
avec $`\mu =c\mathrm{d}s`$.
Remarque 1.10. Les fonctions contenues dans $`\mathrm{\Lambda }_{k,i}`$ ont la propriรฉtรฉ dโรชtre constantes sur chacun des $`\mathrm{\Psi }_i^1(x)`$, $`xX`$. Cela signifie, dans le cas oรน les $`\mathrm{\Psi }_i`$ dรฉfinissent des fibrations, que les fonctions propres $`\phi _{k,i}`$ sont approximรฉes par des fonctions constantes sur les fibres.
#### 1.2.2 Construction de lโopรฉrateur
Rรฉcemment, J. Lott a gรฉnรฉralisรฉ le rรฉsultat de Fukaya aux formes diffรฉrentielles (\[Lo02b\], \[Lo02a\]). Nous allons ici prรฉsenter la construction de lโopรฉrateur limite en nous restreignant par soucis de clartรฉ au cas dโun fibrรฉ $`F(M,g)(N,h)`$ sur une variรฉtรฉ riemannienne dont la fibre est une nilvariรฉtรฉ $`F=\mathrm{\Gamma }\backslash G`$, et qui tend pour la distance de Gromov-Hausdorff vers sa base. Nous noterons $`^{aff}`$ la connexion sur $`F`$ telle que les champs invariants ร gauche soient parallรจles.
La premiรจre difficultรฉ est de dรฉterminer un espace sur lequel va agir lโopรฉrateur limite. Dans la construction de Fukaya, lโopรฉrateur $`P_{(X,\mu )}`$ du thรฉorรจme 1.7 agit sur $`L^2(N)`$, cโest-ร -dire sur un espace de fonctions ร valeurs rรฉelles sur la base. Dans le cas des formes diffรฉrentielles, on utilise un espace de formes diffรฉrentielle sur la base, mais ร valeur dans un espace plus grand que $``$. Plus prรฉcisement, on considรจre un fibrรฉ vectoriel graduรฉ $`E=_{j=0}^mE^j`$ sur la base dont chaque fibre est munie dโun produit scalaire graduรฉ โ *i.e.* tel que les $`E^i`$ soient orthogonaux entre eux โ notรฉ $`h_E`$, et on munit ce fibrรฉ dโune superconnexion $`A^{}`$ de degrรฉ 1, cโest-ร -dire dโun opรฉrateur de la forme
$$A^{}=A_{[0]}^{}+A_{[1]}^{}+A_{[2]}^{}$$
* $`A_{[0]}^{}C^{\mathrm{}}(N,\mathrm{Hom}(E^{},E^{+1}))`$ ;
* $`A_{[1]}^{}`$ est une connexion sur $`E`$ qui prรฉserve la graduation ;
* $`A_{[2]}^{}\mathrm{\Omega }^2(N,\mathrm{Hom}(E^{},E^1))`$.
On peut รฉtendre cette superconnexion par la rรจgle de Leibniz ร un opรฉrateur sur lโespace $`\mathrm{\Omega }(N,E)`$ des formes diffรฉrentielles sur $`N`$ ร valeur dans $`E`$. La mรฉtrique $`h_N`$ et le produit scalaire $`h_E`$ permettent de construire un produit scalaire sur $`\mathrm{\Omega }(N,E)`$, et donc de dรฉfinir un opรฉrateur adjoint ร $`A^{}`$ notรฉ $`(A^{})^{}`$, et un laplacien sur $`\mathrm{\Omega }(N,E)`$ par $`\mathrm{\Delta }_E=A^{}(A^{})^{}+(A^{})^{}A^{}`$. On notera $`\mathrm{\Delta }_E^p`$ la restriction de $`\mathrm{\Delta }_E`$ ร $`_{a+b=p}\mathrm{\Omega }^a(N,E^b)`$.
Pour construire le fibrรฉ sur lequel agit lโopรฉrateur limite, J. Lott se ramรจne dโabord en utilisant \[CFG92\] au cas oรน $`(M,g)`$ est un fibrรฉ affine riemannien:
###### Dรฉfinition 1.11
Un fibrรฉ riemannien $`F(M,g)(N,h)`$ est une fibrรฉ affine riemannien si:
* le groupe de structure du fibrรฉ est contenu dans le groupe $`\mathrm{Aff}(F)`$ des diffรฉomorphismes de $`F`$ qui prรฉservent $`^{aff}`$ ;
* $`M`$ est muni dโune distribution horizontale $`T^HM`$ dont lโholonomie est dans $`\mathrm{Aff}(F)`$ ;
* chaque fibre est munie dโune mรฉtrique $`g_{F_b}`$ qui est parallรจle par rapport ร la connexion affine $`^{aff}`$ sur la fibre ;
* $`N`$ est muni dโune mรฉtrique $`h_N`$ ;
* la mรฉtrique $`g`$ sur $`M`$ sโรฉcrit $`g=h_Ng_{F_b}`$ relativement ร la distribution $`T^HM`$.
En dรฉcomposant les formes diffรฉrentielles de $`M`$ en leurs parties verticale et horizontale, on peut รฉcrire $`\mathrm{\Omega }^{}(M)\mathrm{\Omega }^{}(N,W)`$, oรน $`W`$ est un fibrรฉ sur $`N`$ dont la fibre est isomorphe ร $`\mathrm{\Omega }^{}(F)`$. Dans le cas dโun fibrรฉ affine, le groupe de structure du fibrรฉ $`M`$ prรฉserve $`^{aff}`$, et par consรฉquent lโaction de ce groupe sur $`\mathrm{\Omega }^{}(F)`$ prรฉserve le sous-espace des formes invariantes. On peut donc dรฉfinir le sous-fibrรฉ $`E`$ de $`W`$ de fibre $`\mathrm{\Lambda }^{}(๐ซ^{})`$, ainsi que lโespace $`\mathrm{\Omega }^{}(N,E)`$ sur lequel va agir lโopรฉrateur limite. On a de plus un plongement $`\mathrm{\Omega }^{}(N,E)\mathrm{\Omega }^{}(M)`$ qui permet dโidentifier chaque รฉlรฉment de $`\mathrm{\Omega }^{}(N,E)`$ ร une forme diffรฉrentielle sur $`M`$ parallรจle pour la connexion $`^{aff}`$.
Selon les rรฉsultats donnรฉs dans \[BL95\] sur les superconnexions, et en utilisant le fait que sur un tel fibrรฉ, lโespace des formes parallรจles le long de la fibre est stable par lโaction de la diffรฉrentielle extรฉrieure, cette diffรฉrentielle induit par lโintermรฉdiaire du plongement $`\mathrm{\Omega }(N,E)\mathrm{\Omega }(M)`$ une superconnexion sur $`E`$ telle que $`A_{[0]}^{}`$ soit la diffรฉrentielle sur $`\mathrm{\Lambda }^{}(๐ซ^{})`$, $`A_{[1]}^{}`$ soit la connexion sur le fibrรฉ $`E`$ induite par $`T^HM`$ et $`A_{[2]}^{}`$ soit le produit intรฉrieur $`i_T`$ par la forme de courbure $`T`$ de la distribution $`T^HM`$. Dโautre part, la mรฉtrique $`g_{F_b}`$ induit un produit scalaire $`h_E`$ sur les fibres du fibrรฉ vectoriel $`E`$.
La structure de fibrรฉ affine riemannien induit donc ร la fois une superconnexion et un produit scalaire sur $`E`$, et permet par consรฉquent de dรฉfinir un laplacien $`\mathrm{\Delta }_E`$ sur $`\mathrm{\Omega }(N,E)`$. En notant $`\sigma (\mathrm{\Delta }_E^p)`$ le spectre du laplacien $`\mathrm{\Delta }_E^p`$, J. Lott montre que dans cette situation, les petites valeurs propres du laplacien $`\mathrm{\Delta }_M`$ sur $`M`$ sont celle de $`\mathrm{\Delta }_E^p`$ (\[Lo02b\], thรฉorรจme 1):
###### Thรฉorรจme 1.12
Si $`M`$ est un fibrรฉ affine riemannien sur $`N`$ et si on note $`R_M`$ et $`R_F`$ les tenseurs de courbures respectivement sur $`M`$ et sur $`F_b`$ pour la mรฉtrique $`g_{F_b}`$, $`diam(F)`$ la borne supรฉrieure des diamรจtres des fibres et $`\mathrm{\Pi }`$ la seconde forme fondamentale des fibres, alors il existe des constantes $`a`$, $`a^{}`$ et $`c`$ qui ne dรฉpendent que de $`dim(M)`$ telles que si $`R_F_{\mathrm{}}diam(F)^2a^{}`$ alors pour tout $`pdim(M)`$,
$`\sigma (\mathrm{\Delta }_M^p)[0,adiam(F)^2c(R_M_{\mathrm{}}+\mathrm{\Pi }_{\mathrm{}}^2+T_{\mathrm{}}^2)[`$ $`=`$
$`\sigma (\mathrm{\Delta }_E^p)[0,adiam(F)^2c(R_M_{\mathrm{}}+\mathrm{\Pi }_{\mathrm{}}^2+T_{\mathrm{}}^2)[.`$
On est donc ramenรฉ ร la recherche dโun opรฉrateur limite sur $`\mathrm{\Omega }(N,E)`$. On ne peut considรฉrer sรฉparรฉment la limite de la superconnexion $`A^{}`$ et de la mรฉtrique $`h_E`$. En effet, $`A^{}`$ ne dรฉpend pas de la mรฉtrique sur $`M`$ et donc elle est constante au cours de lโeffondrement, et $`h_E`$ dรฉgรฉnรจre donc sa limite ne permet pas de dรฉfinir un opรฉrateur adjoint comme $`(A^{})^{}`$. Lโidรฉe de Lott est de considรฉrer lโensemble $`(๐ฎ_E\times _E)`$, oรน $`๐ฎ_E`$ est lโespace des superconnexions de degrรฉ 1 sur $`E`$ et $`_E`$ lโespace des produits scalaires euclidiens sur $`E`$, et de quotienter cet ensemble par le groupe $`๐ข_E`$ des $`GL(E)`$-transformations de jauge sur $`E`$ qui prรฉservent la graduation. Il obtient alors le rรฉsultat de compacitรฉ suivant (\[Lo02b\], thรฉorรจme 3):
###### Thรฉorรจme 1.13
Soit $`(N,h_N)`$ une variรฉtรฉ riemannienne fixรฉ, $`n>dim(N)`$, $`a>0`$ et $`\epsilon >0`$. Il existe une partie compacte $`K(n,a,\epsilon )(๐ฎ_E\times _E)/๐ข_E`$ telle que si $`(M^n,g)`$ est une variรฉtรฉ riemannienne de dimension $`n`$ telle que $`R_M_{\mathrm{}}a`$ et $`d_{GH}(M,N)\epsilon `$, alors la classe dโรฉquivalence du couple $`(A^{},h_E)`$ induit par $`g`$ appartient ร $`K`$.
รtant donnรฉe une suite de mรฉtriques $`(g_i)_i`$ qui effondre $`M`$ sur sa base, on peut extraire de la suite $`[(A_i^{},h_{E,i})]_i`$ dโรฉlรฉments de $`(๐ฎ_E\times _E)/๐ข_E`$ une sous-suite qui converge vers une superconnexion limite qui permet de construire un laplacien limite $`\mathrm{\Delta }_{\mathrm{}}`$, dont le spectre est la limite du spectre du laplacien sur $`M`$.
Remarque 1.14. La suite $`[(A_i^{},h_{E,i})]`$ ne converge pas nรฉcessairement. On peut par exemple au paragraphe 1.1 prendre $`\alpha =0`$, $`\beta =1`$ et choisir pour $`\gamma `$ une fonction qui oscille entre $`1`$ et $`2`$. La premiรจre valeur propre va osciller entre $`1`$ et $`t`$ sans converger. Lโopรฉrateur $`\mathrm{\Delta }_E`$ ne converge donc pas, ni la classe de $`(A_t^{},h_{E,t})`$.
Remarque 1.15. Pour lโรฉtude du spectre sur les fonctions, on peut se resteindre au fibrรฉ $`E^0`$, qui est un fibrรฉ trivial en droite rรฉelle sur $`N`$. Cependant, le produit scalaire $`h_{E^0}`$ nโest pas trivial. Il correspond ร la mesure sur lโespace limite dans le travail de Fukaya.
#### 1.2.3 Petites valeurs propres
Pour dรฉterminer la dimension du noyau de $`\mathrm{\Delta }_{\mathrm{}}`$, on peut calculer la cohomologie $`H^{}(A^{})`$ pour lโaction sur $`\mathrm{\Omega }(N,E)`$ de la superconnexion limite $`A^{}`$. J. Lott calcule une majoration de cette dimension en utilisant la thรฉorie des suites spectrales, et en remarquant le fait suivant: le terme $`A_{[0]}^{}`$ de la superconnexion vรฉrifie $`(A_{[0]}^{})^2=0`$, et dรฉfinit donc un complexe diffรฉrentiel sur les fibre de $`E`$, dont la cohomologie $`H^{}(A_{[0]}^{})`$ est un fibrรฉ vectoriel graduรฉ sur $`N`$. De plus, la connexion $`A_{[1]}^{}`$ sur le fibrรฉ $`E`$ passe au quotient sur $`H^{}(A_{[0]}^{})`$ en une connexion plate, cโest-ร -dire telle que $`(A_{[1]}^{})^2=0`$, et dรฉfinit donc aussi un complexe diffรฉrentiel dont la cohomologie est $`H^{}(N,H^{}(A_{[0]}^{}))`$. Les premiers termes de la suite de Leray $`(_r^,,d_r)`$ sont
$`_0^,`$ $`=`$ $`\mathrm{\Omega }^{}(N,E^{}),`$
$`_1^,`$ $`=`$ $`\mathrm{\Omega }^{}(N,H^{}(A_{[0]}^{}))\text{ et}`$
$`_2^,`$ $`=`$ $`H^{}(N,H^{}(A_{[0]}^{})),`$
avec $`d_0=A_{[0]}^{}`$ et $`d_1=A_{[1]}^{}`$. Lott en dรฉduit:
###### Proposition 1.16
$$dimKer\mathrm{\Delta }_{\mathrm{}}^p\underset{a+b=p}{}dim\left(H^a(N,H^b(A_{[0]}^{}))\right).$$
Cette formule peut se simplifier dans certains cas, en particulier pour les $`1`$-formes:
###### Corollaire 1.17
$$dimKer\mathrm{\Delta }_{\mathrm{}}^1b_1(N)+dim(F).$$
On peut noter que cette majoration est trรจs gรฉnรฉrale, en ce sens quโon ne fait pas dโhypothรจse sur la structure du fibrรฉ, ni sur la gรฉomรฉtrie de lโeffondrement. Cependant, le thรฉorรจme รฉnoncรฉ en 15 dans lโintroduction montre que la topologie impose des restrictions sur le nombre de petites valeurs propres possible. En particulier, dans la situation du thรฉorรจme 15, le cas dโรฉgalitรฉ de la majoration donnรฉe par le corollaire 1.17 nโest atteint que si $`G`$ est nilpotent (cf. remarque Introduction).
Dans le cas dโun fibrรฉ en cercle, on obtient aussi une expression simple:
###### Corollaire 1.18
Si $`M`$ est un fibrรฉ en cercle sur $`N`$, alors
$$dim\mathrm{\Delta }_{\mathrm{}}^pb_p(N)+b_{p1}(N)$$
.
Cependant, on sait dรฉjร (\[CC00\]) quโil y a nรฉcessairement รฉgalitรฉ dans lโinรฉgalitรฉ ci-dessus.
Dโautre part, Lott montre quโune petite valeur propre ne peut รชtre obtenue que selon trois mรฉcanismes (\[Lo02b\], Th. 5):
###### Thรฉorรจme 1.19
Soit $`g_i`$ une suite de mรฉtriques qui effondre $`M`$ sur $`N`$. Supposons que $`lim_i\mathrm{}\lambda _{1,p}(M,g_i)=0`$. Alors au moins lโune des trois conditions suivantes est vรฉrifiรฉe:
1. Il existe $`q[0,p]`$ tel que $`b_q(F)<dim\mathrm{\Lambda }^q(๐ซ^{}))`$ ;
2. Il existe $`q[0,p]`$ tel que lโholonomie du fibrรฉ de base $`N`$ et de fibre $`H^q(F)`$ nโest pas semi-simple ;
3. La suite spectrale de Leray qui calcule la cohomologie $`H^{}(M,)`$ du fibrรฉ $`M`$ ne dรฉgรฉnรจre pas au rang $`2`$.
Les trois situations qui interviennent dans ce thรฉorรจme sont illustrรฉes par les trois descriptions topologiques de la nilvariรฉtรฉ dโHeisenberg quโon a donnรฉes en 1.1
Dans le premier cas, le terme $`A_{[0],i}^{}`$ de la superconnexion dรฉgรฉnรจre quand $`i`$ tend vers lโinfini โ cโest-ร -dire que si on note $`(A_{[0],\mathrm{}}^{},h_E)`$ la limite de $`(A_{[0],i}^{},h_{E,i})`$ dans $`(๐ฎ_E\times _E)/๐ข_E`$, la dimension du noyau de $`A_{[0],\mathrm{}}^{}`$ est plus grande que celle du noyau de $`A_{[0],i}^{}`$ โ donc le laplacien restreint ร la fibre admet une petite valeur propre. Gรฉomรฉtriquement, cela signifie que la nilvariรฉtรฉ $`F`$ nโest pas un tore. Un exemple simple est donnรฉ par $`M=N\times F`$ muni dโune mรฉtrique produit, oรน $`F`$ est la nilvariรฉtรฉ dโHeisenberg de dimension $`3`$, et de considรฉrer sur $`F`$ la suite de mรฉtriques dรฉfinie en 1.1 avec $`(\alpha ,\beta ,\gamma )=(\mathrm{1,1,3})`$.
La deuxiรจme condition signifie que le terme $`A_{[1],i}^{}`$ de la superconnexion dรฉgรฉnรจre quand $`i`$ tend vers lโinfini. Lโexemple le plus simple est donnรฉ par la variรฉtรฉ dโHeisenberg de dimension $`3`$ vue comme fibrรฉ en tore sur le cercle : comme la fibre est plate, $`A_{[0],i}^{}=0`$, et comme la base est de dimension $`1`$, les $`2`$-formes sur la base, et donc $`A_{[2],i}^{}`$, sont nulles. Cependant, si on prend $`\alpha =0`$, $`\beta =1`$ et $`\gamma =2`$ en 1.1, on a bien une petite valeur propre.
La troisiรจme condition est illustrรฉe par les situations de fibrรฉ principal. Si on considรจre la variรฉtรฉ dโHeisenberg comme un fibrรฉ principal en cercle sur le tore $`T^2`$, on a $`A_{[0],i}^{}=0`$ (la fibre est plate) et lโholonomie de fibrรฉ $`E`$ est nรฉcessairement semi-simple car les fibres des fibrรฉs $`(E_i)_{i=\mathrm{0,1}}`$ sont de dimension $`1`$. On verra dโautres exemples de fibrรฉs principaux dans le chapitre 3.
Le thรฉorรจme 1.19 donne une condition nรฉcessaire ร lโexistence de petites valeurs propres, ce qui rรฉpond ร la question 8, mais pas ร la question 1.6.
Dans le cas particulier dโune variรฉtรฉ $`M`$, sโeffondrant sur un cercle, Lott donne le corollaire suivant (\[Lo02b\], Cor. 4):
###### Corollaire 1.20
Soit $`g_i`$ une suite de mรฉtriques qui effondre $`M`$ sur $`S^1`$. Supposons que $`lim_i\mathrm{}\lambda _{1,p}(M,g_i)=0`$. Alors au moins lโune des deux conditions suivantes est vรฉrifiรฉe:
1. Il existe $`q[0,p]`$ tel que $`b_q(F)<dim\mathrm{\Lambda }^q(๐ซ^{}))`$ ;
2. Il existe $`q[0,p]`$ tel que si on note $`\mathrm{\Phi }^{}Aut(H^{}(Z))`$ lโaction de lโholonomie sur le fibrรฉ $`H^{}(Z)`$, alors la rรฉduite de Jordan de $`\mathrm{\Phi }^q`$ ou $`\mathrm{\Phi }^{q1}`$ contient un bloc unipotent non trivial.
Le chapitre suivant sera consacrรฉ ร une รฉtude des situations de fibrรฉs en tore sur le cercle, qui illustrent le point 2 du thรฉorรจme 1.19.
## Chapitre 2 Effondrements homogรจnes de fibrรฉs en tores sur le cercle
### 2.1 Structure homogรจne
Nous commenรงons par dรฉmontrer le dรฉbut du thรฉorรจme 15 en construisant le groupe $`G`$ et le rรฉseau $`\mathrm{\Gamma }`$ qui nous intรฉressent. Considรฉrons un fibrรฉ $`M`$ en tore $`T^n`$ sur le cercle qui est la suspension dโun diffรฉomorphisme linรฉaire $`\phi `$ reprรฉsentรฉ par la matrice $`ASL_n()`$. Un tel fibrรฉ sera homรฉomorphe ร
$$M:=T^n\times [\mathrm{0,1}]_{/(x\mathrm{,0})(\phi (x)\mathrm{,1})},$$
(2.1)
Pour construire $`G`$, on va munir $`^{n+1}`$ dโune structure de groupe telle que $`^{n+1}\backslash ^{n+1}=M`$. Si on note $`(x_1,\mathrm{},x_n,y)`$ les รฉlรฉments de $`^{n+1}`$, une telle structure devra vรฉrifier
$$(k_1,\mathrm{},k_n\mathrm{,0})(x_1,\mathrm{},x_n,y)=(x_1+k_1,\mathrm{},x_n+k_n,y)$$
(2.2)
de sorte que les sous-espaces de $`^{n+1}`$ dโรฉquation $`y=c^{te}`$ passent au quotient comme des tores $`T^n`$, et
$$(0,\mathrm{}\mathrm{,0},l)(x_1,\mathrm{},x_n,y)=(A^l\left(\begin{array}{c}x_1\\ \text{.}\text{.}\text{.}\\ x_n\end{array}\right),y+l)$$
(2.3)
de sorte que la structure de fibrรฉ soit bien celle dรฉfinie par (2.1). Cette structure est effectivement rรฉalisรฉe en dรฉfinissant $`G`$ comme lโimage du plongement
$$(x_1,\mathrm{},x_n,y)\left(\begin{array}{ccc}A^y& \begin{array}{c}0\\ \text{.}\text{.}\text{.}\\ 0\end{array}& \begin{array}{c}x_1\\ \text{.}\text{.}\text{.}\\ x_n\end{array}\\ & & \\ 0& \begin{array}{c}\\ 1\\ 0\end{array}& \begin{array}{c}\\ y\\ 1\end{array}\end{array}\right).$$
(2.4)
Comme on se restreint aux matrices $`A`$ qui admettent un logarithme $`B`$, lโexpression $`A^y`$ est bien dรฉfinie en posant $`A^y=\mathrm{exp}(yB)`$. On peut facilement vรฉrifier que cette application est injective, que son image $`G`$ est bien un sous-groupe de $`GL_{n+2}()`$ et que sa structure est bien celle dรฉfinie par (2.2) et (2.3). Enfin, lโimage de $`^{n+1}`$ par cette application est bien un sous-groupe discret de $`G`$, quโon notera $`\mathrm{\Gamma }`$. La variรฉtรฉ $`M`$ est donc homรฉomorphe au quotient $`\mathrm{\Gamma }\backslash G`$.
Remarque: on peut vรฉrifier que si $`A=\left(\begin{array}{cc}1& 1\\ 0& 1\end{array}\right)`$, le groupe $`G`$ obtenu est isomorphe au groupe dโHeisengerg de dimension $`3`$ tel quโil est prรฉsentรฉ dans lโexemple du paragraphe 1.1.
### 2.2 Laplacien
Soit $`X_i`$ et $`Y`$ les champs invariants ร gauche engendrรฉs en $`I_{n+2}`$ respectivement par
$$\frac{}{x_i}=\left(\begin{array}{cc}0& \begin{array}{cc}\begin{array}{c}0\\ \text{.}\text{.}\text{.}\\ \text{.}\end{array}& \begin{array}{c}0\\ \text{.}\text{.}\text{.}\\ 0\end{array}\\ \text{.}& 1\\ \begin{array}{c}\text{.}\\ \text{.}\text{.}\text{.}\\ 0\end{array}& \begin{array}{c}0\\ \text{.}\text{.}\text{.}\\ 0\end{array}\end{array}\\ & \\ 0& 0\end{array}\right)\text{ et }\frac{}{y}=\left(\begin{array}{cc}B& \begin{array}{cc}\begin{array}{c}0\\ \text{.}\text{.}\text{.}\\ 0\end{array}& \begin{array}{c}0\\ \text{.}\text{.}\text{.}\\ 0\end{array}\end{array}\\ & \\ 0& \begin{array}{cc}0& \mathrm{\hspace{0.25em}\hspace{0.25em}1}\\ 0& \mathrm{\hspace{0.25em}\hspace{0.25em}0}\end{array}\end{array}\right).$$
(2.5)
Ces champs vรฉrifient $`[X_i,X_j]=0`$ et $`[Y,X_i]=_{j=1}^nb_{ji}X_j`$. On peut remarquer que lโapplication $`X[Y,X]`$ est un endomorphisme de lโespace $`\mathrm{\Gamma }(T_VM)^^G`$ des champs de vecteurs invariants verticaux, cโest-ร -dire lโespace engendrรฉ par les $`X_i`$, et dont la matrice est $`B`$. On notera $`f`$ cet endomorphisme.
On fixe une mรฉtrique homogรจne $`g`$ sur $`M`$ en se donnant une base $`(V_i)_{i[1,n]}`$ de lโespace $`\mathrm{\Gamma }(T_VM)^^G`$, cette mรฉtrique รฉtant telle que $`(V_1,\mathrm{},V_n,Y)`$ soit orthonormรฉe en tout point. On notera $`(V_1^{\mathrm{}},\mathrm{},V_n^{\mathrm{}},Y^{\mathrm{}})`$ sa base duale, et $`C`$ la matrice de $`f`$ dans la base $`(V_1,\mathrm{},V_n)`$. On va dรฉterminer le spectre du laplacien $`\mathrm{\Delta }_{inv}^1`$ restreint ร lโensemble $`\mathrm{\Omega }^1(M)^G`$ des $`1`$-formes invariantes ร gauche en fonction des coefficients de $`C`$. Plus prรฉcisรฉment, on a le
###### Lemme 2.6
La matrice du laplacien $`\mathrm{\Delta }_{inv}^1`$ dans la base $`(V_1^{\mathrm{}},\mathrm{},V_n^{\mathrm{}},Y^{\mathrm{}})`$ est
$$\mathrm{\Delta }_{inv}^1:\left(\begin{array}{cc}C^tC& \begin{array}{c}0\\ \text{.}\text{.}\text{.}\\ 0\end{array}\\ 0\mathrm{}0& 0\end{array}\right).$$
Dรฉmonstration: Les crochets de Lie entre les vecteurs de la nouvelle base sont
$$[V_i,V_j]=0\text{ et }[Y,V_i]=\underset{j=1}{\overset{n}{}}c_{ji}V_j.$$
(2.7)
Soit $`\alpha `$ une $`1`$-forme diffรฉrentielle invariante. Sa diffรฉrentielle extรฉrieure est dรฉterminรฉe par la relation $`\mathrm{d}\alpha (U_1,U_2)=U_1\alpha (U_2)U_2\alpha (U_1)\alpha ([U_1,U_2])`$, oรน $`U_1`$ et $`U_2`$ sont des champs de vecteur. Si ces champs sont invariants ร gauche, cette relation devient: $`\mathrm{d}\alpha (U_1,U_2)=\alpha ([U_1,U_2])`$. On en dรฉduit:
$$\mathrm{d}Y^{\mathrm{}}=0\text{ et }\mathrm{d}V_i^{\mathrm{}}=\underset{j=1}{\overset{n}{}}c_{ij}Y^{\mathrm{}}V_j^{\mathrm{}}.$$
(2.8)
La matrice de la diffรฉrentielle extรฉrieure $`\mathrm{d}:\mathrm{\Omega }^1(M)^G\mathrm{\Omega }^2(M)^G`$ sera, dans les bases $`(V_1^{\mathrm{}},\mathrm{},V_n^{\mathrm{}},Y^{\mathrm{}})`$ et $`(Y^{\mathrm{}}V_i^{\mathrm{}},V_i^{\mathrm{}}V_j^{\mathrm{}})`$,
$$\mathrm{d}:\left(\begin{array}{cc}^tC& \begin{array}{c}0\\ \text{.}\text{.}\text{.}\\ 0\end{array}\\ 0& \begin{array}{c}0\\ \text{.}\text{.}\text{.}\\ 0\end{array}\end{array}\right).$$
(2.9)
Les deux bases sont orthonormรฉes, donc la matrice dans ces bases de la divergence $`\delta :\mathrm{\Omega }^2(M)^G\mathrm{\Omega }^1(M)^G`$ sera donc la transposรฉe de la matrice ci-dessus.
Comme la diffรฉrentielle restreinte ร $`\mathrm{\Omega }^0(M)^G`$ est nulle, le laplacien $`\mathrm{\Delta }=\delta \mathrm{d}+\mathrm{d}\delta `$ se rรฉduit sur $`\mathrm{\Omega }^1(M)^G`$ ร lโopรฉrateur $`\delta \mathrm{d}`$. On en dรฉduit la matrice du laplacien $`\mathrm{\Delta }^{inv}`$ restreint ร $`\mathrm{\Omega }^1(M)^G`$ est, dans la base $`(V_1^{\mathrm{}},\mathrm{},V_n^{\mathrm{}},Y^{\mathrm{}})`$,
$$\left(\begin{array}{cc}C& 0\\ 0\mathrm{}0& 0\mathrm{}0\end{array}\right)\left(\begin{array}{cc}^tC& \begin{array}{c}0\\ \text{.}\text{.}\text{.}\\ 0\end{array}\\ 0& \begin{array}{c}0\\ \text{.}\text{.}\text{.}\\ 0\end{array}\end{array}\right)=\left(\begin{array}{cc}C^tC& \begin{array}{c}0\\ \text{.}\text{.}\text{.}\\ 0\end{array}\\ 0\mathrm{}0& 0\end{array}\right).$$
(2.10)
Remarque: On a fait ici le calcul pour un $`Y`$ fixรฉ, cโest-ร -dire pour un certain choix de connexion du fibrรฉ. Mais si on choisit $`Y^{}`$ tel que $`Y^{}Y\mathrm{\Gamma }(T_VM)^^G`$ et une mรฉtrique telle que $`(V_1,\mathrm{},V_n,Y^{})`$ soit orthonormรฉe, le rรฉsultat sera le mรชme car on aura toujours $`[Y^{},V_i]=[Y,V_i]=_{j=1}^nb_{ji}V_j`$.
On peut noter que la mรฉtrique intervient par la rรฉรฉcriture de la matrice $`B`$ dans une base orthonormรฉe. Ce travail de renormalisation correspond dans le travail de J. Lott au passage au quotient de $`(๐ฎ_E\times _E)`$ par le groupe de transformation de jauge $`๐ข_E`$. Si deux mรฉtriques donnent la mรชme matrice $`C`$, cela signifie que les deux รฉlรฉments correspondants dans $`(๐ฎ_E\times _E)`$ appartiennent ร la mรชme classe dans $`(๐ฎ_E\times _E)/๐ข_E`$.
### 2.3 Courbure
Nous allons dรฉmontrer dans cette partie un lemme qui nous servira ร faire le lien entre le contrรดle de la courbure et lโexistence de petites valeurs propres.
###### Lemme 2.11
Soit $`a`$ la borne supรฉrieure de la valeur absolue de la courbure sectionnelle de $`(M,g)`$. Il existe des constantes $`\tau (n)>0`$ et $`\kappa (B)`$ telle que
$$\tau ^1a<Tr(C^tC)<\tau a+\kappa .$$
Dรฉmonstration: Rappelons tout dโabord lโexpression suivante (dont le lecteur pourra trouver la dรฉmonstration dans \[CE75\]) de la courbure sectionnelle $`K(U,V)`$, oรน $`U`$ et $`V`$ sont deux champs invariants ร gauche dโun groupe de Lie quelconque:
$`K(U,V)`$ $`=`$ $`{\displaystyle \frac{1}{4}}ad_U^{}V+ad_V^{}U^2ad_U^{}U,ad_V^{}V`$
$`{\displaystyle \frac{3}{4}}[U,V]^2{\displaystyle \frac{1}{2}}[[U,V],V],U{\displaystyle \frac{1}{2}}[[V,U],U],V.`$
Nous allons appliquer cette relation aux champs de la base $`(V_i,Y)`$. Pour cela, remarquons dโabord que les matrices de $`ad_Y`$ et $`ad_{V_i}`$ sont, dans cette base
$$ad_Y:\left(\begin{array}{cc}C& \begin{array}{c}0\\ \text{.}\text{.}\text{.}\\ 0\end{array}\\ 0\mathrm{}0& 0\end{array}\right)\text{ et }ad_{V_i}:\left(\begin{array}{cc}0& \begin{array}{c}c_{1i}\\ \text{.}\text{.}\text{.}\\ c_{ni}\end{array}\\ 0\mathrm{}0& 0\end{array}\right)$$
(2.13)
On en dรฉduit $`ad_{V_i}^{}Y=0`$, $`ad_Y^{}Y=0`$, $`ad_Y^{}V_i=_jc_{ij}V_j`$ et $`ad_{V_i}^{}V_j=c_{ji}Y`$, et donc que
$`K(Y,V_i)`$ $`=`$ $`{\displaystyle \frac{1}{4}}ad_Y^{}V_i^2{\displaystyle \frac{3}{4}}[Y,V_i]^2{\displaystyle \frac{1}{2}}[[V_i,Y],Y],V_i`$ (2.14)
$`=`$ $`{\displaystyle \frac{1}{4}}{\displaystyle \underset{j}{}}\left(c_{ij}^23c_{ji}^22c_{ij}c_{ji}\right)`$
$`=`$ $`{\displaystyle \underset{j}{}}c_{ji}^2+{\displaystyle \frac{1}{4}}{\displaystyle \underset{j}{}}(c_{ij}c_{ji})^2`$
et
$`K(V_i,V_j)`$ $`=`$ $`{\displaystyle \frac{1}{4}}ad_{V_i}^{}V_j+ad_{V_j}^{}V_i^2ad_{V_i}^{}V_i,ad_{V_J}^{}V_j`$ (2.15)
$`=`$ $`{\displaystyle \frac{1}{4}}(c_{ij}+c_{ji})^2c_{ii}c_{jj}.`$
Dโautre part, comme $`C`$ est la matrice de $`f`$, le terme de degrรฉ $`n2`$ du polynรดme caractรฉristique est indรฉpendant de la mรฉtrique choisie. Le calcul montre que son coefficient est $`\kappa =_{ij}(c_{ii}c_{jj}c_{ij}c_{ji})`$. On peut en dรฉduire que
$$\underset{i,j=1}{\overset{n}{}}K(V_i,V_j)+\kappa =\underset{i,j=1}{\overset{n}{}}\left(\frac{1}{4}(c_{ij}+c_{ji})^2c_{ij}c_{ji}\right)=\frac{1}{4}\underset{i,j=1}{\overset{n}{}}(c_{ij}c_{ji})^2,$$
(2.16)
et donc que
$$\underset{i,j=1}{\overset{n}{}}c_{ji}^2=\underset{i,j=1}{\overset{n}{}}K(V_i,V_j)\underset{i=1}{\overset{n}{}}K(Y,V_i)+\kappa (n^2+n)a+\kappa ,$$
(2.17)
ce qui montre lโune des deux inรฉgalitรฉs du lemme. La seconde dรฉcoule du fait que la courbure sectionnelle sโรฉcrit comme un polynรดme homogรจne de degrรฉ deux relativement aux $`c_{ij}`$.
### 2.4 Petites valeurs propres
Nous allons maintenant dรฉmontrer les rรฉsultats concernant le spectre de $`\mathrm{\Delta }^{inv}`$.
Dรฉmonstration de 15.1 : Si $`U`$ est un vecteur colonne tel que $`C^tCU=0`$, alors $`{}_{}{}^{t}UC^tCU=0`$, et donc $`^tCU=0`$. Par consรฉquent, $`dimKerC^tC=dimKer^tC=dimKerC=d^{}`$. Comme $`dimKer\mathrm{\Delta }_{inv}^1=1+dimKerC^tC`$, on a bien $`dimKer\mathrm{\Delta }_{inv}^1=d^{}+1`$.
Dรฉmonstration de 15.2 : Cโest une consรฉquence directe du lemme 2.11: comme la trace de $`\mathrm{\Delta }_{inv}^1`$ est celle de $`C^tC`$, cette trace est majorรฉe en fonction de $`a`$ et $`B`$. Comme les valeurs propres de $`\mathrm{\Delta }_{inv}^1`$ sont positives, chacune est majorรฉe.
Dรฉmonstration de 15.3 :
Supposons que $`dn`$. Soit $`E_0`$ le sous-espace caractรฉristique de $`f`$ associรฉ ร la valeur propre $`0`$. On notera $`E_0^{}`$ son orthogonal pour la dualitรฉ dans lโespace des $`1`$-formes invariantes verticales. Comme $`dn`$, lโespace $`E_0^{}`$ est de dimension non nulle. On va montrer que le quotient de Rayleigh est uniformรฉment minorรฉ sur $`E_0^{}`$, pour ensuite appliquer le principe du minimax.
Remarques: comme les formes et les mรฉtriques considรฉrรฉes sont invariantes, la norme ponctuelle dโune forme ne dรฉpendra pas du point oรน on la calcule, ce qui permet dโรฉcrire que $`R(\alpha )=\frac{d\alpha ^2}{\alpha ^2}=\frac{|d\alpha |^2}{|\alpha |^2}`$. Dโautre part, il faut noter que la notion dโorthogonalitรฉ pour la dualitรฉ est indรฉpendante de la mรฉtrique. En particulier, comme $`E_0`$ est dรฉfini indรฉpendamment de la mรฉtrique, $`E_0^{}`$ le sera aussi.
Soit $`V^{\mathrm{}}E_0^{}`$ et $`(V_i)`$ une base orthonormรฉe de $`\mathrm{\Gamma }(T_VM)^^G`$ telle que $`(V_1,\mathrm{},V_d)`$ soit une base orthonormรฉe de $`E_0`$ (si $`d=0`$ et donc $`E_0=0`$, on choisit alors $`(V_i)`$ orthonormรฉe quelconque, la suite de la dรฉmonstration restant valide). Lโespace $`E_0`$ est stable par $`f`$, donc $`E_0^{}`$ est stable par $`{}_{}{}^{t}f`$, et la matrice de $`(^tf)_{|E_0^{}}`$ dans la base $`(V_{d+1}^{\mathrm{}},\mathrm{},V_n^{\mathrm{}})`$ est $`{}_{}{}^{t}D`$, oรน $`D`$ est une sous-matrice de $`C`$. Comme la relation (2.8) peut sโรฉcrire $`\mathrm{d}V^{\mathrm{}}=Y^{\mathrm{}}(^tf)(V^{\mathrm{}})`$, on a
$$|\mathrm{d}V^{\mathrm{}}|^2=|(^tf)(V^{\mathrm{}})|^2\lambda |V^{\mathrm{}}|^2,$$
(2.18)
$`\lambda `$ est la plus petite valeur propre de $`D^tD`$. Dโune part, le dรฉterminant de cette matrice vรฉrifie
$$DetD^tD=(Det^tD)^2=(Det(^tf)_{|E_0^{}})^2,$$
(2.19)
et donc $`DetD^tD`$ est indรฉpendant du choix de la base $`(V_i)`$. Dโautre part, $`Det^t(f_{|E_0^{}})`$ est non nul. En effet, si $`{}_{}{}^{t}f_{|E_0^{}}^{}(\alpha )=0`$, alors $`\alpha f=0`$, donc $`\alpha `$ est orthogonal ร lโimage de $`f`$, qui contient les sous-espaces caractรฉristiques de $`f`$ autres que $`E_0`$, et par consรฉquent $`\alpha `$ est nul. On en dรฉduit que $`\lambda `$ est uniformรฉment minorรฉe: sโil existe une suite de mรฉtriques telle que $`\lambda 0`$, alors la plus grande valeur propre de $`D^tD`$ tend vers lโinfini (car $`DetD^tD`$ est constant), ce qui est impossible puisque la courbure est bornรฉe et que $`TrC^tCTrD^tD`$ (car $`D`$ est une sous-matrice de $`C`$), et donc que la somme des valeurs propres de $`D^tD`$ est bornรฉe.
On a montrรฉ que le quotient de Rayleigh de $`\alpha E_0^{}`$ est minorรฉ par une constante $`c(f,a)`$ indรฉpendante de la mรฉtrique et du choix de $`\alpha `$. Comme $`dimE_0^{}=nd`$, le principe du minimax nous dit donc que les $`n+1d`$ plus grandes valeurs propres de $`\mathrm{\Delta }_{inv}^1`$ sont minorรฉes par $`c`$. Comme $`dimKer\mathrm{\Delta }_{inv}^1=d^{}+1`$ et que $`dim\mathrm{\Omega }^1(M)^G=n+1`$, on en dรฉduit que $`\lambda _{dd^{}+\mathrm{1,1}}^{inv}>c`$.
Si $`d=n`$, alors il existe $`PGL_n()`$ tel que $`P^1BP`$ soit triangulaire supรฉrieure avec des $`0`$ sur la diagonale, et comme $`P^1AP=P^1\mathrm{exp}(B)P=\mathrm{exp}(P^1BP)`$, la matrice $`P^1AP`$ sera triangulaire supรฉrieure avec des $`1`$ sur la diagonale. On en dรฉduit, en posant $`P^{}=\left(\begin{array}{cc}P& 0\\ 0& I\end{array}\right)GL_{n+2}()`$, que le groupe $`P_{}^{}{}_{}{}^{1}GP^{}`$, oรน $`G`$ est le groupe construit au paragraphe 2.1, est constituรฉ de matrices triangulaires supรฉrieures avec des 1 sur la diagonale. Cโest donc un groupe nilpotent.
Lโexistence dโun effondrement tel que toutes les valeurs propres de $`\mathrm{\Delta }_{inv}^1`$ tendent vers zรฉro dรฉcoulera du 15.4
Dรฉmonstration de 15.4 : On vient de dรฉmontrer que si $`d`$ (et donc $`d^{}`$) est nul, il nโy a pas de petites valeurs propres.
Supposons que $`d>0`$. Pour simplifier, nous allons montrer le rรฉsultat dans le cas oรน la partie nilpotente de la rรฉduite de Jordan de $`B`$ ne comporte quโun seul bloc de Jordan, la construction de $`g_\epsilon ^k`$ รฉtant semblable dans le cas gรฉnรฉral.
On construit une base $`(V_1,\mathrm{},V_n)`$ de $`\mathrm{\Gamma }(T_VM)^^G`$ en choisissant une base de Jordan $`(V_1,\mathrm{},V_d)`$ de $`E_0`$ (en particulier, $`(V_1,\mathrm{},V_d^{})`$ sera une base de $`Kerf`$) que lโon on la complรจte de maniรจre quelconque en une base $`(V_1,\mathrm{},V_d^{})`$ de $`\mathrm{\Gamma }(T_VM)^^G`$. On notera $`C`$ la matrice de $`f`$ dans cette base. La matrice $`C`$ nโest pas de Jordan, mais sa restriction ร $`E_0`$, cโest-ร -dire le bloc carrรฉ supรฉrieur droit de taille $`d`$, lโest. Elle est de la forme
$$C=\left(\begin{array}{ccc}\begin{array}{cccc}0& \mathrm{}& \mathrm{}& 0\\ \mathrm{}& \mathrm{}& & \mathrm{}\\ & & \mathrm{}& \mathrm{}\\ & & & 0\\ & & & \mathrm{}\\ & & & \\ \mathrm{}& & & \mathrm{}\\ 0& \mathrm{}& \mathrm{}& 0\end{array}& \begin{array}{cccc}0& \mathrm{}& \mathrm{}& 0\\ \mathrm{}& & & \mathrm{}\\ 0& & & \\ 1& \mathrm{}& & \mathrm{}\\ 0& \mathrm{}& 0& 0\\ \mathrm{}& \mathrm{}& 1& 0\\ \mathrm{}& & 0& 1\\ 0& \mathrm{}& 0& 0\end{array}& C_1\\ & & \\ \begin{array}{cccc}\mathrm{}& & & \mathrm{}\\ 0& \mathrm{}& \mathrm{}& 0\end{array}& \begin{array}{cccc}\mathrm{}& & & \mathrm{}\\ 0& \mathrm{}& \mathrm{}& 0\end{array}& C_2\end{array}\right)$$
(2.20)
$$\stackrel{\underset{}{}}{d^{}\mathrm{colonnes}}\stackrel{\underset{}{}}{dd^{}\mathrm{colonnes}}$$
$`C_2`$ est un bloc carrรฉ de taille $`nd`$ et de dรฉterminant non nul.
Soit $`kdd^{}`$. On pose $`V_i^\epsilon =\nu _i(\epsilon )V_i`$, avec $`\nu _i(\epsilon )=\epsilon ^1`$ pour $`id^{}+k`$, et $`\nu _i(\epsilon )=\epsilon ^{(1+d^{}+ki)}V_i`$ pour $`i<d^{}+k`$. La matrice $`C_\epsilon `$ de $`f`$ dans cette base vรฉrifiera
$$c_{ij}^\epsilon =\frac{\nu _j(\epsilon )}{\nu _i(\epsilon )}c_{ij},$$
(2.21)
donc $`c_{ij}^\epsilon =c_{ij}`$ pour $`id^{}+k`$ (en tenant compte du fait que $`c_{ij}=0`$ pour $`id`$ et $`jd`$), et $`c_{ij}^\epsilon 0`$ quand $`\epsilon 0`$, pour $`i<d^{}+k`$. La matrice $`C_\epsilon `$ tend donc vers une matrice $`C_0`$ de la forme
$$C_0=\left(\begin{array}{ccc}\begin{array}{cccc}0& \mathrm{}& \mathrm{}& 0\\ \mathrm{}& \mathrm{}& & \mathrm{}\\ & & \mathrm{}& \mathrm{}\end{array}& \begin{array}{cccc}0& \mathrm{}& \mathrm{}& 0\\ \mathrm{}& & & \mathrm{}\\ 0& & & \end{array}& 0\\ & & \\ \begin{array}{cccc}& & & 0\\ & & & \mathrm{}\\ & & & \\ \mathrm{}& & & \mathrm{}\\ 0& \mathrm{}& \mathrm{}& 0\end{array}& \begin{array}{cccc}1& \mathrm{}& & \mathrm{}\\ 0& \mathrm{}& 0& 0\\ \mathrm{}& \mathrm{}& 1& 0\\ \mathrm{}& & 0& 1\\ 0& \mathrm{}& \mathrm{}& 0\end{array}& C_1^{}\\ & & \\ \begin{array}{cccc}\mathrm{}& & & \mathrm{}\\ 0& \mathrm{}& \mathrm{}& 0\end{array}& \begin{array}{cccc}\mathrm{}& & & \mathrm{}\\ 0& \mathrm{}& \mathrm{}& 0\end{array}& C_2^{}\end{array}\right)$$
(2.22)
$$\stackrel{\underset{}{}}{d^{}+k\mathrm{colonnes}}\stackrel{\underset{}{}}{dd^{}k\mathrm{colonnes}}$$
Comme les $`\lambda _{i\mathrm{,1}}^{inv}`$ sont ces fonctions continues de $`C`$, il suffit de calculer la dimension du noyau de $`C_0^tC_0`$, qui est รฉgale ร celle de $`KerC_0`$. Dโune part, les $`d^{}+k`$ premiรจres colonnes de $`C_0`$ sont nulles, donc $`dimKerC_0d^{}+k`$, et dโautre part, comme $`detC_20`$, la famille formรฉe par les lignes $`d^{}+k`$ ร $`d1`$ et les $`nd`$ derniรจres lignes de $`C_0`$ est libre, donc $`dimKerC_0n(nd)((d1)(d^{}+k1))=d^{}+k`$. De mรชme, $`dimKerC=d^{}`$, donc on a bien exactement $`k`$ petites valeurs propres.
Si la partie nilpotente de la rรฉduite de Jordan de $`B`$ contient plusieurs blocs de Jordan, on obtient le rรฉsultat en procรฉdant de la mรชme maniรจre pour annuler le nombre souhaitรฉ de lignes dans $`C`$.
Remarques: la famille de matrice $`C_\epsilon `$ est uniformement bornรฉe par rapport ร $`\epsilon `$, et le lemme 2.11 donne la majoration $`|K(M,g_\epsilon ^k)|\tau Tr(C_\epsilon ^tC_\epsilon )`$, pour tout $`\epsilon `$. La courbure sectionnelle du fibrรฉ est donc bien uniformรฉment bornรฉe au cours de lโeffondrement. Dโautre part, on voit que si on effondre le fibrรฉ par homothรฉtie de la fibre, par exemple en posant $`\nu _i(\epsilon )=\epsilon ^1`$ pour tout $`i`$, la matrice $`C_\epsilon `$ est indรฉpendante de $`\epsilon `$, et donc il nโy a pas de petite valeur propre.
Dรฉmonstration du corollaire 19 : Si $`d=d^{}`$ et $`dn`$, alors $`\lambda _{\mathrm{1,1}}^{inv}`$ est uniformรฉment minorรฉ dโaprรจs 15.3. Si $`d=d^{}`$ et $`d=n`$, alors $`B=0`$ et toutes les valeurs propres de $`\mathrm{\Delta }_{inv}^1`$ sont nulles.
Si $`dd^{}`$, alors 15.4 garantit lโexistence dโune petite valeur propre.
Dรฉmonstration du thรฉorรจme 21 : Comme $`B`$ est semi-simple, son orbite par conjugaison est fermรฉe (\[CM93\], p. 28). Comme la courbure est bornรฉe, la norme de $`C`$ reste bornรฉe quand la base $`(V_1^{\mathrm{}},\mathrm{},V_n^{\mathrm{}},Y^{\mathrm{}})`$ โ et donc la mรฉtrique โ varie. La matrice $`C`$ est par construction dans lโorbite par conjugaison de $`B`$, donc elle prend finalement ses valeurs au cours de lโeffondrement dans une partie compacte $`K`$ de cette orbite.
La base orthonormรฉe $`(V_1^{\mathrm{}},\mathrm{},V_n^{\mathrm{}},Y^{\mathrm{}})`$ de $`\mathrm{\Omega }^1(M)^G`$ engendre, par produit extรฉrieur, une base orthonormรฉe de $`\mathrm{\Omega }^{}(M)^G`$. Les coefficients de la matrice de la diffรฉrentielle extรฉrieure $`\mathrm{d}`$ dans cette base, et donc ceux de la matrice de $`\mathrm{\Delta }=\mathrm{d}\delta +\delta \mathrm{d}`$, sont des fonctions continues de $`CK`$. Par consรฉquent, quand la mรฉtrique varie, la matrice de $`\mathrm{\Delta }`$ prend ses valeurs dans un compact image de $`K`$. Sโil existe une famille de mรฉtriques telle que $`\lambda _{1,p}^{inv}`$ tende vers zรฉro pour un $`p[1,n]`$, alors la matrice de $`\mathrm{\Delta }`$ tend vers une matrice de rang strictement infรฉrieur, ce qui est impossible puisque, par compacitรฉ, la matrice limite sera dans lโimage de $`K`$, donc de mรชme rang que $`\mathrm{\Delta }`$.
Par consรฉquent, lโopรฉrateur $`\mathrm{\Delta }`$ restreint ร $`\mathrm{\Omega }^{}(M)^G`$ nโadmet pas de petite valeur propre.
### 2.5 Variรฉtรฉs de petites dimensions
En petite dimension, on peut รชtre plus prรฉcis que les thรฉorรจmes 15 et 21, et mettre en รฉvidence un lien simple entre lโexistence de petites valeurs propres et la structure du groupe $`G`$:
###### Corollaire 2.23
Supposons que $`n=2\text{ ou }3`$. Sโil existe $`p[1,n]`$ et une suite de mรฉtriques homogรจnes sur $`M`$ telle que la courbure sectionnelle associรฉe soit uniformรฉment bornรฉe et que $`\lambda _{1,p}^{inv}`$ tende vers $`0`$, alors $`G`$ est nilpotent.
Remarque 2.24. Cโest par exemple la situation exposรฉe en 1.1, oรน on a $`p=1`$ et $`n=2`$.
Dรฉmonstration du corollaire 2.23 : Montrons dโabord que sโil existe $`p`$ tel que $`\lambda _{p\mathrm{,1}}^{inv}0`$, alors $`dd^{}`$.
Si $`p=1`$, cela dรฉcoule du corollaire 19. Si $`p=n`$, on est ramenรฉ ร la situation $`p=1`$ par dualitรฉ de Hodge.
Reste les cas $`p=2`$ et $`n=3`$. On a dรฉjร calculรฉ les matrices de $`\delta :\mathrm{\Omega }^2(M)^G\mathrm{\Omega }^1(M)^G`$ et $`\mathrm{d}:\mathrm{\Omega }^1(M)^G\mathrm{\Omega }^2(M)^G`$. On en dรฉduit que la matrice de $`\delta \mathrm{d}`$, en restriction ร $`\mathrm{\Omega }^2(M)^G`$ est de la forme, dans les bases introduites au paragraphe 2.2,
$$\mathrm{d}\delta :\left(\begin{array}{cc}{}_{}{}^{t}CC& 0\\ 0& 0\end{array}\right).$$
(2.25)
Comme la variรฉtรฉ est de dimension $`4`$, lโopรฉrateur de Hodge $``$ est une isomรฉtrie de $`\mathrm{\Omega }^2(M)^G`$. En restriction ร $`\mathrm{\Omega }^2(M)^G`$, on aura $`\delta \mathrm{d}=\mathrm{d}\delta `$, et donc $`\delta \mathrm{d}`$ et $`\mathrm{d}\delta `$ ont mรชme spectre. Dโautre part, dโaprรจs la thรฉorie de Hodge, le spectre du laplacien est la rรฉunion des spectres de $`\delta \mathrm{d}`$ et $`\mathrm{d}\delta `$, on dรฉduit de ce qui prรฉcรจde quโune petite valeur propre non nulle de $`\mathrm{\Delta }_{inv}^2`$ sera une petite valeur propre non nulle de $`\mathrm{d}\delta |\mathrm{\Omega }^2(M)^G`$, et donc une petite valeur propre de $`{}_{}{}^{t}CC`$. Le raisonnement appliquรฉ ร $`C^tC`$ dans la dรฉmonstration de 15.4 reste valable pour $`{}_{}{}^{t}CC`$. On peut donc conclure que si $`\lambda _{\mathrm{2,1}}^{inv}`$ tend vers zรฉro, alors $`dd^{}`$.
Supposons que $`dd^{}`$. Alors le noyau de $`B`$ est non trivial, par consรฉquent $`d>d^{}>0`$ et la multiplicitรฉ de la valeur propre $`0`$ de $`B`$ est au moins รฉgale ร deux. Si $`n=3`$ la troisiรจme valeur propre est รฉgale ร la trace de $`B`$ qui est nulle puisquโelle est rรฉelle et que $`\mathrm{exp}(TrB)=det(\mathrm{exp}B)=detA=1`$. Donc $`d=n`$, et $`G`$ est nilpotent, dโaprรจs 15.3.
### 2.6 Homologie du fibrรฉ
Nous allons ici montrer que lโon peut calculer le premier nombre de Betti du fibrรฉ $`M`$ indรฉpendamment de la cohomologie, en utilisant le fait que le rรฉseau $`\mathrm{\Gamma }`$ est isomorphe au groupe fondamental de $`M`$.
###### Thรฉorรจme 2.26
Soit $`M=\mathrm{\Gamma }\backslash G`$ un fibrรฉ en tore $`T^n`$ sur le cercle construit selon 15, dรฉfinit par une matrice $`ASL_n()`$. Alors le premier nombre de Betti de $`M`$ est $`b_1(M)=1+dimKer(AI)`$.
On en dรฉduit:
###### Corollaire 2.27
Si $`1`$ nโest pas valeur propre de $`A`$, alors les $`1`$-formes harmoniques de $`M`$ sont $`G`$-invariantes.
On verra au paragraphe 2.7.2 un exemple qui montre โ entre autres choses โ quโon peut effectivement, dans certains cas, avoir des formes harmoniques qui ne sont pas invariantes.
Dรฉmonstration du thรฉorรจme 2.26 : Comme $`G`$ est simplement connexe, le rรฉseau $`\mathrm{\Gamma }`$ est isomorphe au groupe fondamental du quotient $`M=\mathrm{\Gamma }\backslash G`$. Pour dรฉterminer le premier nombre de Betti de $`M`$, on va calculer lโabรฉlianisรฉ de son groupe fondamental.
Rappelons que le groupe $`\mathrm{\Gamma }`$ est de la forme:
$$(x_1,\mathrm{},x_n,y)\left(\begin{array}{ccc}A^y& \begin{array}{c}0\\ \text{.}\text{.}\text{.}\\ 0\end{array}& \begin{array}{c}x_1\\ \text{.}\text{.}\text{.}\\ x_n\end{array}\\ & & \\ 0& \begin{array}{c}\\ 1\\ 0\end{array}& \begin{array}{c}\\ y\\ 1\end{array}\end{array}\right),$$
(2.28)
avec $`(x_1,\mathrm{},x_n,y)^{n+1}`$. Dans la suite de la dรฉmonstration, nous noterons les รฉlรฉments de $`\mathrm{\Gamma }`$ indiffรฉremment sous la forme de vecteurs lignes ou de vecteurs colonnes.
Soit $`g=(x_1,\mathrm{},x_n,y)`$ et $`g^{}=(x_1^{},\mathrm{},x_n^{},y^{})`$ deux รฉlรฉments de $`\mathrm{\Gamma }`$. Leur inverse respectives sont:
$$g^1=\left(\begin{array}{c}A^y\left(\begin{array}{c}x_1\\ \text{.}\text{.}\text{.}\\ x_n\end{array}\right)\\ y\end{array}\right)\text{ et }g^1=\left(\begin{array}{c}A^y^{}\left(\begin{array}{c}x_1^{}\\ \text{.}\text{.}\text{.}\\ x_n^{}\end{array}\right)\\ y^{}\end{array}\right).$$
(2.29)
Le calcul de leur commutateur $`[g,g^{}]=gg^{}g^1g^1`$ donne:
$$[g,g^{}]=\left(\begin{array}{c}\left(\begin{array}{c}x_1\\ \text{.}\text{.}\text{.}\\ x_n\end{array}\right)+A^y\left(\begin{array}{c}x_1^{}\\ \text{.}\text{.}\text{.}\\ x_n^{}\end{array}\right)A^y^{}\left(\begin{array}{c}x_1\\ \text{.}\text{.}\text{.}\\ x_n\end{array}\right)\left(\begin{array}{c}x_1^{}\\ \text{.}\text{.}\text{.}\\ x_n^{}\end{array}\right)\\ 0\end{array}\right),$$
(2.30)
soit $`[g,g^{}]=((A^yI)\left(\begin{array}{c}x_1^{}\\ \text{.}\text{.}\text{.}\\ x_n^{}\end{array}\right)(A^y^{}I)\left(\begin{array}{c}x_1\\ \text{.}\text{.}\text{.}\\ x_n\end{array}\right)\mathrm{,0})`$.
On voit que $`[\mathrm{\Gamma },\mathrm{\Gamma }]((AI)^n\mathrm{,0})`$. Rรฉciproquement, si on fixe $`g=(0,\mathrm{}\mathrm{0,1})`$ et quโon fait varier $`(x_1^{},\mathrm{},x_n^{})`$, on obtient que $`[\mathrm{\Gamma },\mathrm{\Gamma }]((AI)^n\mathrm{,0})`$, et donc $`[\mathrm{\Gamma },\mathrm{\Gamma }]`$ est exactement $`((AI)^n\mathrm{,0})`$.
Lโimage de $`^n`$ par $`(AI)`$ est un sous-rรฉseau dโindice fini du rรฉseau des entiers de $`Im(AI)`$, donc le quotient de $`(^n\mathrm{,0})`$ par $`((AI)^n\mathrm{,0})`$ est de la forme $`Z^k\times H`$, oรน $`k=codimIm(AI)`$ et $`H`$ est un groupe fini โ รฉventuellement trivial โ, et finalement lโabรฉlianisรฉ $`\mathrm{\Gamma }^{}=\mathrm{\Gamma }/[\mathrm{\Gamma },\mathrm{\Gamma }]`$ du groupe $`\mathrm{\Gamma }`$ est donc de la forme $`Z^{k+1}\times H`$. Le premier nombre de Betti de $`M`$ est donc
$$b_1(M)=1+dimKer(AI).$$
(2.31)
Dรฉmonstration du corollaire 2.27 : Si $`1`$ nโest pas valeur propre de $`A`$, alors $`0`$ nโest pas valeur propre de $`B`$, et par consรฉquent, en vertu du point 15.1, $`b_1(M)=dimKer\mathrm{\Delta }_{inv}^1=1`$. Toutes les $`1`$-formes harmoniques sont donc dans le noyau de $`\mathrm{\Delta }_{inv}^1`$, et en particulier sont $`G`$-invariantes.
### 2.7 Exemples
#### 2.7.1 Petites valeurs propres pour les $`2`$-formes diffรฉrentielles
Nous allons donner ici un exemple de fibrรฉ en tore sur le cercle pour lequel $`d=d^{}=0`$ et $`\mathrm{\Delta }_{inv}^2`$ admet une petite valeur propre. Cet exemple montre que le corollaire 19 ne se gรฉnรฉralise pas ร $`n`$ et $`p`$ quelconque.
On dรฉfinit le fibrรฉ considรฉrรฉ par la matrice
$$A=\left(\begin{array}{cc}A^{}& A^{\prime \prime }\\ & \\ 0& A^{}\end{array}\right),$$
(2.32)
avec
$$A^{}=\left(\begin{array}{cc}2& 1\\ 1& 1\end{array}\right)\text{ et }A^{\prime \prime }=\left(\begin{array}{cc}0& 1\\ 0& 0\end{array}\right).$$
(2.33)
###### Fait 2.34
La matrice $`A`$ est semblable ร une matrice de la forme
$$\left(\begin{array}{cccc}e^\lambda & 1& 0& 0\\ 0& e^\lambda & 0& 0\\ & & & \\ 0& 0& e^\lambda & 1\\ 0& 0& 0& e^\lambda \end{array}\right),$$
(2.35)
$`\lambda `$ est un rรฉel non nul.
Dรฉmonstration : La matrice $`A^{}`$ admet deux valeurs propres rรฉelles positives, qui sont inverses lโune de lโautre car $`DetA^{}=1`$. On notera $`\lambda `$ le rรฉel positif tel que ces deux valeurs propres soient $`e^\lambda `$ et $`e^\lambda `$. Elles sont aussi valeurs propres de $`A`$ avec la multiplicitรฉ deux. On peut vรฉrifier que le polynรดme caractรฉristique de $`A`$ est son polynรดme minimal. Les sous-espaces propres de $`A`$ sont donc tous les deux de dimension 1, et par consรฉquent, les deux blocs de sa rรฉduite de Jordan sont $`\left(\begin{array}{cc}e^\lambda & 1\\ 0& e^\lambda \end{array}\right)`$ et $`\left(\begin{array}{cc}e^\lambda & 1\\ 0& e^\lambda \end{array}\right)`$.
###### Fait 2.36
Il existe une suite de mรฉtrique $`g_\epsilon `$ sur $`M=\mathrm{\Gamma }\backslash G(B)`$ et une suite de matrices $`C_\epsilon `$ associรฉes telles que
$$C_\epsilon =\left(\begin{array}{cccc}\lambda & ฯต& 0& 0\\ 0& \lambda & 0& 0\\ & & & \\ 0& 0& \lambda & ฯต\\ 0& 0& 0& \lambda \end{array}\right).$$
(2.37)
Dรฉmonstration : Comme on a
$$\mathrm{exp}\left(\begin{array}{cc}\lambda & e^\lambda \\ 0& \lambda \end{array}\right)=\left(\begin{array}{cc}e^\lambda & 1\\ 0& e^\lambda \end{array}\right)\text{ et }\mathrm{exp}\left(\begin{array}{cc}\lambda & e^\lambda \\ 0& \lambda \end{array}\right)=\left(\begin{array}{cc}e^\lambda & 1\\ 0& e^\lambda \end{array}\right),$$
La matrice $`A`$ admet un logarithme semblable ร
$$C=\left(\begin{array}{cccc}\lambda & e^\lambda & 0& 0\\ 0& \lambda & 0& 0\\ & & & \\ 0& 0& \lambda & e^\lambda \\ 0& 0& 0& \lambda \end{array}\right).$$
(2.38)
Soit $`(V_1,V_2,V_3,V_4)`$ la base dans laquelle la matrice de lโendomorphisme $`f`$ est รฉgal ร $`C`$. Si on pose $`V_i^\epsilon =\epsilon ^\alpha V_i`$ pour $`i=\mathrm{1,3}`$, $`V_2^\epsilon =\epsilon ^{\alpha +1}e^\lambda V_2`$ et $`V_4^\epsilon =\epsilon ^{\alpha +1}e^\lambda V_4`$$`\alpha `$ est un rรฉel strictement positif, la matrice de $`f`$ dans cette base sera
$$C_\epsilon =\left(\begin{array}{cccc}\lambda & ฯต& 0& 0\\ 0& \lambda & 0& 0\\ & & & \\ 0& 0& \lambda & ฯต\\ 0& 0& 0& \lambda \end{array}\right).$$
(2.39)
Il suffit donc de dรฉfinir $`g_\epsilon `$ en posant que la base $`(V_1^\epsilon ,V_2^\epsilon ,V_3^\epsilon ,V_4^\epsilon ,Y)`$ est orthonormรฉe. Le fait que la courbure reste bornรฉe quand $`\epsilon 0`$ dรฉcoule du lemme 2.11
###### Fait 2.40
La valeur propre $`\lambda _{\mathrm{2,1}}^{inv}(M,g_\epsilon )`$ tend vers zรฉro quand $`\epsilon 0`$.
Dรฉmonstration :
On va calculer la matrice de $`\mathrm{d}:\mathrm{\Omega }^2(M)^G\mathrm{\Omega }^3(M)^G`$ dans des bases de la forme $`(V_i^{\mathrm{}}V_j^{\mathrm{}},V_i^{\mathrm{}}Y^{\mathrm{}})`$ et $`(V_i^{\mathrm{}}V_j^{\mathrm{}}Y^{\mathrm{}},V_i^{\mathrm{}}V_j^{\mathrm{}}V_k^{\mathrm{}})`$.
En utilisant (2.8), on obtient que $`\mathrm{d}V_i^{\mathrm{}}Y^{\mathrm{}}=0`$ pour tout $`i`$, et que
$`\mathrm{d}(V_1^{\mathrm{}}V_2^{\mathrm{}})`$ $`=`$ $`(\lambda V_1^{\mathrm{}}+\epsilon V_2^{\mathrm{}})Y^{\mathrm{}}V_2^{\mathrm{}}V_1^{\mathrm{}}\lambda V_2^{\mathrm{}}Y^{\mathrm{}}`$ (2.41)
$`=`$ $`2\lambda V_1^{\mathrm{}}V_2^{\mathrm{}}Y^{\mathrm{}},`$
$`\mathrm{d}(V_1^{\mathrm{}}V_3^{\mathrm{}})`$ $`=`$ $`(\lambda V_1^{\mathrm{}}+\epsilon V_2^{\mathrm{}})Y^{\mathrm{}}V_3^{\mathrm{}}V_1^{\mathrm{}}(\lambda V_3^{\mathrm{}}+\epsilon V_4^{\mathrm{}})Y^{\mathrm{}}`$ (2.42)
$`=`$ $`\epsilon V_2^{\mathrm{}}V_3^{\mathrm{}}Y^{\mathrm{}}\epsilon V_1^{\mathrm{}}V_4^{\mathrm{}}Y^{\mathrm{}},`$
$`\mathrm{d}(V_1^{\mathrm{}}V_4^{\mathrm{}})`$ $`=`$ $`(\lambda V_1^{\mathrm{}}+\epsilon V_2^{\mathrm{}})Y^{\mathrm{}}V_4^{\mathrm{}}V_1^{\mathrm{}}(\lambda V_4^{\mathrm{}})Y^{\mathrm{}}`$ (2.43)
$`=`$ $`\epsilon V_2^{\mathrm{}}V_4^{\mathrm{}}Y^{\mathrm{}},`$
$`\mathrm{d}(V_2^{\mathrm{}}V_3^{\mathrm{}})`$ $`=`$ $`\lambda V_2^{\mathrm{}}Y^{\mathrm{}}V_3^{\mathrm{}}V_2^{\mathrm{}}(\lambda V_3^{\mathrm{}}+\epsilon V_4^{\mathrm{}})Y^{\mathrm{}}`$ (2.44)
$`=`$ $`\epsilon V_2^{\mathrm{}}V_4^{\mathrm{}}Y^{\mathrm{}},`$
$`\mathrm{d}(V_2^{\mathrm{}}V_4^{\mathrm{}})`$ $`=`$ $`\lambda V_2^{\mathrm{}}Y^{\mathrm{}}V_4^{\mathrm{}}V_2^{\mathrm{}}(\lambda V_4^{\mathrm{}})Y^{\mathrm{}}=0,`$ (2.45)
$`\mathrm{d}(V_3^{\mathrm{}}V_4^{\mathrm{}})`$ $`=`$ $`(\lambda V_3^{\mathrm{}}+\epsilon V_4^{\mathrm{}})Y^{\mathrm{}}V_4^{\mathrm{}}V_3^{\mathrm{}}(\lambda V_4^{\mathrm{}})Y^{\mathrm{}}`$ (2.46)
$`=`$ $`2\lambda V_3^{\mathrm{}}V_4^{\mathrm{}}Y^{\mathrm{}}.`$
La matrice de $`\mathrm{d}`$ dans les bases
$$(V_1^{\mathrm{}}V_2^{\mathrm{}},V_1^{\mathrm{}}V_3^{\mathrm{}},V_1^{\mathrm{}}V_4^{\mathrm{}},V_2^{\mathrm{}}V_3^{\mathrm{}},V_2^{\mathrm{}}V_4^{\mathrm{}},V_3^{\mathrm{}}V_4^{\mathrm{}},V_i^{\mathrm{}}Y^{\mathrm{}})$$
et
$$\begin{array}{c}(V_1^{\mathrm{}}V_2^{\mathrm{}}Y^{\mathrm{}},V_1^{\mathrm{}}V_3^{\mathrm{}}Y^{\mathrm{}},V_1^{\mathrm{}}V_4^{\mathrm{}}Y^{\mathrm{}},V_2^{\mathrm{}}V_3^{\mathrm{}}Y^{\mathrm{}},\hfill \\ V_2^{\mathrm{}}V_4^{\mathrm{}}Y^{\mathrm{}},V_3^{\mathrm{}}V_4^{\mathrm{}}Y^{\mathrm{}},V_i^{\mathrm{}}V_j^{\mathrm{}}V_k^{\mathrm{}})\hfill \end{array}$$
est de la forme
$$\left(\begin{array}{cc}\begin{array}{cccccc}\hfill 2\lambda & \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0\\ \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0\\ \hfill 0& \hfill \epsilon & \hfill 0& \hfill 0& \hfill 0& \hfill 0\\ \hfill 0& \hfill \epsilon & \hfill 0& \hfill 0& \hfill 0& \hfill 0\\ \hfill 0& \hfill 0& \hfill \epsilon & \hfill \epsilon & \hfill 0& \hfill 0\\ \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 2\lambda \end{array}& 0\\ & \\ 0& 0\end{array}\right).$$
(2.47)
On voit que quand $`\epsilon `$ tend vers zรฉro, cette matrice tend vers une matrice de rang strictement infรฉrieur. On peut en dรฉduire comme au paragraphe 2.4 que lโopรฉrateur $`\delta \mathrm{d}`$ admet une petite valeur propre, qui sera aussi petite valeur propre de $`\mathrm{\Delta }`$.
#### 2.7.2 Structure homogรจne non abรฉlienne sur le tore
Nous allons ici รฉtudier plus en dรฉtail un exemple particulier de fibrรฉ en tore $`T^2`$ sur le cercle, pour mettre en รฉvidence plusieurs de ses propriรฉtรฉs.
Ce fibrรฉ est construit par le thรฉorรจme 15, avec la donnรฉe de
$$A=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)\text{ et }B=\left(\begin{array}{cc}0& 2\pi \\ 2\pi & 0\end{array}\right).$$
(2.48)
La matrice $`\mathrm{exp}(xB)`$ est de la forme
$$\left(\begin{array}{cc}\mathrm{cos}2\pi x& \mathrm{sin}2\pi x\\ \mathrm{sin}2\pi x& \mathrm{cos}2\pi x\end{array}\right).$$
(2.49)
Cโest donc la matrice dโune rotation dโangle $`2\pi x`$ ; nous la noterons $`R(2\pi x)`$.
Le groupe $`G`$ sโรฉcrit:
$$\left(\begin{array}{cc}R(2\pi x)& \begin{array}{cc}0& y\\ 0& z\end{array}\\ & \\ 0& \begin{array}{cc}1& x\\ 0& 1\end{array}\end{array}\right),x,y,z.$$
(2.50)
Un premiรจre remarque est de constater que la variรฉtรฉ $`M=\mathrm{\Gamma }\backslash G`$ est un tore:
###### Fait 2.51
$`\mathrm{\Gamma }`$ est isomorphe ร $`^3`$ et $`\mathrm{\Gamma }\backslash G`$ est diffรฉomorphe ร $`T^3`$.
En effet, le rรฉseau $`\mathrm{\Gamma }`$ sโรฉcrit:
$$\left(\begin{array}{cc}\begin{array}{cc}1& 0\\ 0& 1\end{array}& \begin{array}{cc}0& y\\ 0& z\end{array}\\ & \\ 0& \begin{array}{cc}1& x\\ 0& 1\end{array}\end{array}\right),x,y,z.$$
(2.52)
On peut vรฉrifier que $`\mathrm{\Gamma }`$ est bien abรฉlien. Par ailleurs, comme la topologie du fibrรฉ est entiรจrement dรฉterminรฉe par la matrice $`A`$, le fibrรฉ est bien trivial. On peut remarquer le fait que $`\mathrm{\Gamma }`$ soit abรฉlien est cohรฉrent avec le fait que ce soit le groupe fondamental dโun tore.
On a donc construit un groupe de Lie rรฉsoluble simplement connexe dont un sous-groupe cocompact est commutatif. En comparaison, on a pour les groupes nilpotents le rรฉsultat suivant (\[Ra72\]):
###### Thรฉorรจme 2.53
Soit $`N_1`$ et $`N_2`$ deux groupes de Lie nilpotents simplement connexes, et $`\mathrm{\Gamma }_1`$, $`\mathrm{\Gamma }_2`$ deux sous-groupes cocompacts de $`N_1`$ et $`N_2`$ respectivement. Alors tout isomorphisme entre $`\mathrm{\Gamma }_1`$ et $`\mathrm{\Gamma }_2`$ sโรฉtend en un isomorphisme entre $`N_1`$ et $`N_2`$.
En particulier, si un groupe nilpotent simplement connexe contient un sous-groupe cocompact isomorphe ร $`^n`$, alors il est abรฉlien. Le groupe $`G`$ illustre donc le fait que ce thรฉorรจme ne se gรฉnรฉralise pas aux groupes rรฉsolubles.
Dโautre part, comme $`M=\mathrm{\Gamma }\backslash G`$ et un tore, son premier nombre de Betti est รฉgal ร sa dimension, donc $`b_1(M)=3`$. Or, selon le thรฉorรจme 15, si on munit $`G`$ dโune mรฉtrique invariante le noyau de $`\mathrm{\Delta }_{inv}^p`$ est de dimension $`1`$. On en conclut :
###### Fait 2.54
Soit $`g`$ une mรฉtrique $`G`$-invariante ร gauche sur $`M`$. Alors il existe sur $`M`$ des $`1`$-formes harmoniques qui ne sont pas $`G`$-invariantes.
On voit donc que la proposition 14 et le corollaire 2.27 ne se gรฉnรฉralisent pas ร toutes les solvariรฉtรฉs. Rรฉciproquement, le groupe $`G`$ illustre la situation oรน la multiplicitรฉ de la valeur propre $`1`$ dans $`A`$ est strictement supรฉrieure ร la multiplicitรฉ de $`0`$ dans $`B`$.
Enfin, on peut remarquer que dโaprรจs les formules (2.14) et (2.15), pour la mรฉtrique invariante sur $`M`$ telle que la base $`(/x,/y,/z)`$ soit orthonormรฉe ร lโorigine, la courbure sur $`M`$ est nulle (on peut vรฉrifier que cโest en gรฉnรฉral faux pour une mรฉtrique invariante quelconque). On va reformuler ce rรฉsultat et en donner une dรฉmonstration trรจs simple qui ne fait pas appel aux formules du paragraphe 2.3:
###### Fait 2.55
Il existe sur $`^3`$ une mรฉtrique invariante pour la structure canonique de groupe abรฉlien, et invariante pour lโaction ร gauche du groupe $`G`$.
Dรฉmonstration : On considรจre sur $`^3`$ la mรฉtrique euclidienne canonique, et on note $`x`$, $`y`$ et $`z`$ les cordonnรฉes canoniques. Si $`a`$, $`b`$ et $`c`$ sont des rรฉels fixรฉs, la paramรฉtrisation de $`G`$ donnรฉe par (2.50) dรฉfinit lโaction ร gauche de $`(a,b,c)`$ comme รฉtant
$$\left(\begin{array}{c}a\\ b\\ c\end{array}\right)\left(\begin{array}{c}x\\ y\\ z\end{array}\right)=\left(\begin{array}{c}a+x\\ \left(\begin{array}{c}b\\ c\end{array}\right)+R(2\pi a)\left(\begin{array}{c}y\\ z\end{array}\right)\end{array}\right).$$
(2.56)
On voit que lโaction de $`(a,b,c)`$ est la composรฉe dโune rotation et dโune translation. Cโest donc une isomรฉtrie pour la norme euclidienne canonique. Par consรฉquent, cette norme est invariante ร gauche pour lโaction de $`G`$.
Remarquons pour finir quโon peut facilement gรฉnรฉraliser cet exemple en dimension supรฉrieure en construisant une matrice $`A`$ contenant un bloc de la forme $`\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)`$, par exemple :
$$A=\left(\begin{array}{cccc}1& 0& 0& 0\\ 0& 1& 0& 0\\ 0& 0& 1& 1\\ 0& 0& 0& 1\end{array}\right)\text{ et }B=\left(\begin{array}{cccc}0& 2\pi & 0& 0\\ 2\pi & 0& 0& 0\\ 0& 0& 0& 1\\ 0& 0& 0& 0\end{array}\right).$$
(2.57)
Le fibrรฉ obtenu dans ce cas sera une solvariรฉtรฉ ayant une topologie de nilvariรฉtรฉ, mais dont les formes harmoniques ne sont pas toutes invariantes.
## Chapitre 3 Effondrements homogรจnes de fibrรฉs principaux en tores sur le tore
### 3.1 Topologie et spectre du fibrรฉ
Nous allons ici dรฉmontrer le thรฉorรจme 22.
Dรฉmonstration de 22.1 : Soit $`M`$ un fibrรฉ principal de base $`T^2`$ et de fibre $`F`$. La base du fibrรฉ peut sโรฉcrire
$$[\mathrm{0,1}]\times [\mathrm{0,1}]_/,$$
$``$ est la relation dโรฉquivalence engendrรฉe par $`(x\mathrm{,0})(x\mathrm{,1})`$ et $`(0,y)(1,y)`$. Le fibrรฉ $`M`$ peut alors se dรฉfinir par la donnรฉe, pour tout point $`p`$ du bord $`K`$ de $`K=[\mathrm{0,1}]\times [\mathrm{0,1}]`$ dโun diffรฉomorphisme $`\phi _p`$ de la fibre, et en posant
$$M=K\times F_{/(p,x)(q,\phi _q^1\phi _p(x)),xF,p,qK,pq}.$$
(3.1)
Lโhypothรจse de principalitรฉ se traduit ici par le fait que pour tout $`p,qK`$ tels que $`pq`$ et pour tout $`g,xF`$, on a
$$(p,gx)(q,g\phi _q^1\phi _p(x)),$$
(3.2)
ce qui impose aux $`\phi _q^1\phi _p`$ dโรชtre des translations ร droite sur la fibre. On peut donc, sans perte de gรฉnรฉralitรฉ se restreindre, pour le choix des $`\phi _p`$, au groupe des translations de la fibre, qui est isomorphe ร $`F`$. Le fibrรฉ est donc dรฉterminรฉ par la donnรฉe dโune application de $`K`$ dans $`F`$. Comme sa topologie ne dรฉpend pas de la classe dโhomotopie de cette application, lโensemble des fibrรฉs principaux de fibre $`F`$ sur le tore $`T^2`$ est parametrรฉ par le groupe fondamental de $`F`$. Il sโagit en fait dโun exemple de classe dโobstruction (\[St51\], ยง 35) qui est, dans le cas gรฉnรฉral dโun $`F`$-fibrรฉ principal sur une variรฉtรฉ compacte $`N`$ un รฉlรฉment de $`H_2(N,\pi _1(F))`$ et qui mesure lโobstruction du fibrรฉ ร รชtre trivial.
Considรฉrons maintenant un fibrรฉ principal $`M`$ de fibre $`T^n`$, et $`(a_1,\mathrm{},a_n)\pi _1(T^n)=^n`$ sa classe dโobstruction. Nous allons munir $`^{n+2}`$ dโune structure de groupe telle que la topologie du quotient ร gauche par $`^{n+2}`$ soit celle du fibrรฉ. Pour ce faire, nous choisirons le reprรฉsentant $`\gamma :KT^n`$ de la classe dโobstruction de la maniรจre suivante:
$$\begin{array}{c}\gamma _{|\{0\}\times [\mathrm{0,1}]}=\gamma _{|[\mathrm{0,1}]\times \{0\}}=\gamma _{|[\mathrm{0,1}]\times \{1\}}=0,\\ \\ \gamma (1,t)=(ta_1,\mathrm{},ta_n),t[\mathrm{0,1}],\end{array}$$
(3.3)
de sorte quโun รฉlรฉment $`(x_1,\mathrm{},x_n)T^n`$ de la fibre au dessus de $`(0,t)K`$ sera identifiรฉ ร lโรฉlรฉment $`(x_1+ta_1,\mathrm{},x_n+ta_n)`$ au dessus de $`(1,t)`$. Si on note $`(x_1,\mathrm{},x_n,y_1,y_2)`$ les รฉlรฉments de $`^{n+2}`$, on veut donc dรฉfinir sur cet ensemble un produit tel que
$$(k_1,\mathrm{},k_n\mathrm{,0,0})(x_1,\mathrm{},x_n,y_1,y_2)=(x_1+k_1,\mathrm{},x_n+k_n,y_1,y_2)$$
(3.4)
de sorte que dโune part les sous-espaces de $`^{n+2}`$ dโรฉquation $`(y_1,y_2)=c^{te}`$ passent au quotient comme des tores, et tel que
$$\begin{array}{c}(0,\mathrm{}\mathrm{,0},l_1,l_2)(x_1,\mathrm{},x_n,y_1,y_2)=\hfill \\ \hfill (x_1+y_2a_1l_1,\mathrm{},x_n+y_2a_nl_1,y_1+l_1,y_2+l_2),\end{array}$$
(3.5)
de sorte que la structure de fibrรฉ en tore sera bien celle dรฉfinie par (3.3).
On peut effectivement construire une telle structure de groupe en plongeant $`^{n+2}`$ dans $`M_{n+3}()`$ par lโapplication suivante:
$$(x_1,\mathrm{},x_n,y_1,y_2)\left(\begin{array}{cccc}I_n& \begin{array}{c}0\\ \text{.}\text{.}\text{.}\\ 0\end{array}& \begin{array}{c}a_1y_1\\ \text{.}\text{.}\text{.}\\ a_ny_1\end{array}& \begin{array}{c}x_1\\ \text{.}\text{.}\text{.}\\ x_n\end{array}\\ & & & \\ 0& \begin{array}{c}\\ 1\\ 0\\ 0\end{array}& \begin{array}{c}\\ 0\\ 1\\ 0\end{array}& \begin{array}{c}\\ y_1\\ y_2\\ 1\end{array}\end{array}\right).$$
(3.6)
Notons $`G`$ lโimage de cette application. Cโest un sous-groupe de $`M_{n+3}()`$, et le quotient $`\mathrm{\Gamma }\backslash G`$$`\mathrm{\Gamma }`$ est le rรฉseau des entiers de $`G`$ est diffรฉomorphe ร la variรฉtรฉ $`M`$, qui est donc une nilvariรฉtรฉ.
Supposons maitenant que $`n2`$. On pose $`d=\mathrm{pgcd}(a_1,\mathrm{},a_n)`$ et $`a_i^{}=a_i/d`$. Soit $`P=(p_{ij})M(n,)`$ une matrice telle que $`p_{i1}=a_i^{}`$ et que ses vecteurs colonnes forment une base du rรฉseau $`^n`$. On a alors
$$P^1\left(\begin{array}{c}a_1\\ \text{.}\text{.}\text{.}\\ a_n\end{array}\right)=\left(\begin{array}{c}d\\ 0\\ \text{.}\text{.}\text{.}\\ 0\end{array}\right)$$
(3.7)
et
$$\left(\begin{array}{cc}P^1& 0\\ & \\ 0& I\end{array}\right)\left(\begin{array}{cccc}I_n& \begin{array}{c}0\\ \text{.}\text{.}\text{.}\\ 0\end{array}& \begin{array}{c}a_1y_1\\ \text{.}\text{.}\text{.}\\ a_ny_1\end{array}& \begin{array}{c}x_1\\ \text{.}\text{.}\text{.}\\ x_n\end{array}\\ & & & \\ 0& \begin{array}{c}\\ 1\\ 0\\ 0\end{array}& \begin{array}{c}\\ 0\\ 1\\ 0\end{array}& \begin{array}{c}\\ y_1\\ y_2\\ 1\end{array}\end{array}\right)\left(\begin{array}{cc}P& 0\\ & \\ 0& I\end{array}\right)=\left(\begin{array}{cccc}I_n& \begin{array}{c}0\\ \text{.}\text{.}\text{.}\\ 0\end{array}& \begin{array}{c}dy_1\\ 0\\ \text{.}\text{.}\text{.}\\ 0\end{array}& \begin{array}{c}x_1^{}\\ \text{.}\text{.}\text{.}\\ x_n^{}\end{array}\\ & & & \\ 0& \begin{array}{c}\\ 1\\ 0\\ 0\end{array}& \begin{array}{c}\\ 0\\ 1\\ 0\end{array}& \begin{array}{c}\\ y_1\\ y_2\\ 1\end{array}\end{array}\right)$$
(3.8)
avec $`\left(\begin{array}{c}x_1^{}\\ \text{.}\text{.}\text{.}\\ x_n^{}\end{array}\right)=P^1\left(\begin{array}{c}x_1\\ \text{.}\text{.}\text{.}\\ x_n\end{array}\right)`$. On peut voir que le groupe $`P_{}^{}{}_{}{}^{1}GP^{}`$, avec $`P^{}=\left(\begin{array}{cc}P& 0\\ 0& I\end{array}\right)`$, est isomorphe ร $`^{n1}\times G^{}`$, oรน $`G^{}`$ est le groupe
$$\left\{\left(\begin{array}{cccc}1& 0& dy_1& x\\ 0& 1& 0& y_1\\ 0& 0& 1& y_2\\ 0& 0& 0& 1\end{array}\right),x,y_1,y_2\right\}.$$
(3.9)
Dโautre part, comme $`detP^{}=1`$ et que $`P^{}`$ est ร coefficients entiers, le rรฉseau des matrices ร coefficients entiers de $`P_{}^{}{}_{}{}^{1}GP^{}`$ est exactement $`P_{}^{}{}_{}{}^{1}\mathrm{\Gamma }P^{}`$, oรน $`\mathrm{\Gamma }`$ est le rรฉseau des entiers de $`G`$. La variรฉtรฉ $`M=\mathrm{\Gamma }\backslash G`$, qui est diffรฉomorphe ร $`P_{}^{}{}_{}{}^{1}\mathrm{\Gamma }P^{}\backslash P_{}^{}{}_{}{}^{1}GP^{}`$ peut donc sโรฉcrire
$$M(^{n1}\times \mathrm{\Gamma }^{})\backslash (^{n1}\times G^{})T^{n1}\times N,$$
(3.10)
$`N=\mathrm{\Gamma }^{}\backslash G^{}`$, en notant $`\mathrm{\Gamma }^{}=`$ le rรฉseau des entiers de $`G^{}`$.
Ce calcul montre quโon peut se ramener au cas oรน les $`a_i`$, $`i2`$ sont nuls. On supposera dans la suite que cโest le cas, et on posera $`a_1=a`$.
Dรฉmonstration de 22.2 : Soient $`X_i`$, $`Y_1`$ et $`Y_2`$ les champs de vecteurs invariants ร gauche engendrรฉs en $`I_{n+3}`$ respectivement par
$$\frac{}{x_i}=\left(\begin{array}{cc}0& \begin{array}{ccc}\begin{array}{c}0\\ \text{.}\text{.}\text{.}\\ \text{.}\end{array}& \begin{array}{c}0\\ \text{.}\text{.}\text{.}\\ \text{.}\end{array}& \begin{array}{c}0\\ \text{.}\text{.}\text{.}\\ 0\end{array}\\ \text{.}& \text{.}& 1\\ \begin{array}{c}\text{.}\\ \text{.}\text{.}\text{.}\\ 0\end{array}& \begin{array}{c}\text{.}\\ \text{.}\text{.}\text{.}\\ 0\end{array}& \begin{array}{c}0\\ \text{.}\text{.}\text{.}\\ 0\end{array}\end{array}\\ & \\ 0& 0\end{array}\right),\frac{}{y_1}=\left(\begin{array}{cc}0& \begin{array}{ccc}\begin{array}{c}0\\ \\ \text{.}\text{.}\text{.}\\ 0\end{array}& \begin{array}{c}a\\ 0\\ \text{.}\text{.}\text{.}\\ 0\end{array}& \begin{array}{c}0\\ \\ \text{.}\text{.}\text{.}\\ 0\end{array}\end{array}\\ & \\ 0& \begin{array}{ccc}0& 0& \mathrm{\hspace{0.33em}\hspace{0.25em}1}\\ 0& 0& \mathrm{\hspace{0.33em}\hspace{0.25em}0}\\ 0& 0& \mathrm{\hspace{0.33em}\hspace{0.25em}0}\end{array}\end{array}\right)\text{ et }\frac{}{y_2}=\left(\begin{array}{cc}0& 0\\ & \\ 0& \begin{array}{ccc}0& 0& 0\\ 0& 0& 1\\ 0& 0& 0\end{array}\end{array}\right).$$
Ces champs vรฉrifient $`[X_i,X_j]=0,[X_i,Y_j]=0\text{ et }[Y_1,Y_2]=aX_1`$. On notera $`V`$ le vecteur $`aX_1`$, dont on peut remarquer quโil est non nul (si $`a`$ est nul, le fibrรฉ est trivial).
Soit $`g`$ une mรฉtrique homogรจne sur $`M`$, et $`(V_i)`$ une base de $`\mathrm{\Gamma }(T_VM)^^G`$ telle que $`(V_1,\mathrm{},V_n,Y_1,Y_2)`$ soit orthonormรฉe en tout point et que $`V_1`$ soit colinรฉaire ร $`V`$. Les crochets de Lie entre les vecteurs de cette base sont:
$$[V_i,V_j]=0,[V_i,Y_j]=0,\text{ et }[Y_1,Y_2]=\eta V_1,\eta ^{}.$$
(3.11)
On en dรฉduit:
$$\mathrm{d}V_1^{\mathrm{}}=\eta Y_1^{\mathrm{}}Y_2^{\mathrm{}}\text{ et }\mathrm{d}V_i^{\mathrm{}}=\mathrm{d}Y_j^{\mathrm{}}=0,i>1,$$
(3.12)
$`(V_1^{\mathrm{}},\mathrm{},V_n^{\mathrm{}},Y_1^{\mathrm{}},Y_2^{\mathrm{}})`$ est la base duale de $`(V_1,\mathrm{},V_n,Y_1,Y_2)`$. Les formes de cette base de $`\mathrm{\Omega }^1(M)^G`$ engendrent, par produit extรฉrieur, une base de $`\mathrm{\Omega }^{}(M)^G`$ composรฉe de formes propres de $`\mathrm{\Delta }_{inv}`$. En effet, il dรฉcoule de (3.12) quโelles sont toutes fermรฉes sauf celles de la forme $`V_1V_{i_1}\mathrm{}V_{i_k}`$ ($`i_j1`$), dont la diffรฉrentielle vaut:
$$\mathrm{d}(V_1V_{i_1}\mathrm{}V_{i_k})=\eta Y_1Y_2V_{i_1}\mathrm{}V_{i_k},$$
(3.13)
et, puisque $`\delta =(1)^{n(p+1)+1}\mathrm{d}`$, elles sont toutes cofermรฉes sauf celles de la forme $`Y_1Y_2V_1V_{i_1}\mathrm{}V_{i_k}`$ dont la codiffรฉrentielle vaut:
$$\delta (Y_1Y_2V_1V_{i_1}\mathrm{}V_{i_k})=\eta V_1V_{i_1}\mathrm{}V_{i_k}.$$
(3.14)
En restriction ร $`\mathrm{\Omega }^p(M)^G`$, les formes de la base sont donc harmoniques, sauf $`C_{n1}^{p1}`$ formes cofermรฉes et $`C_{n1}^{p2}`$ formes fermรฉes qui sont des formes propres de valeur propre รฉgale ร $`\eta ^2`$. Lโopรฉrateur $`\mathrm{\Delta }_{inv}^p`$ admet donc une unique valeur propre non nulle, de multiplicitรฉ $`C_{n1}^{p1}+C_{n1}^{p2}=C_n^{p1}`$ et รฉgale ร $`\eta ^2=|V|^2`$.
Si on choisit une base orthonormรฉe de la forme $`(V_1,\mathrm{},V_n,Y_1^{},Y_2^{})`$, avec $`Y_i^{}=Y_i+_{k=1}^n\xi _kV_k`$, on aura toujours $`[Y_1^{},Y_2^{}]=[Y_1,Y_2]`$. Le rรฉsultat ne dรฉpend donc pas du choix de la connexion sur le fibrรฉ. Remarquons enfin que si lโon choisit une autre mรฉtrique sur la base (en se donnant deux champs horizontaux quelconques $`Y_1^{}`$ et $`Y_2^{}`$ et en les supposant orthogonaux) on obtiendra le mรชme rรฉsultat en remplaรงant $`V`$ par $`V^{}=[Y_1^{},Y_2^{}]`$, avec comme valeur propre $`\eta ^2=|V^{}|^2=Vol(B)^2|V|^2`$.
### 3.2 Effondrements des fibrรฉs principaux sur le tore $`T^2`$
Dans cette partie, nous allons montrer que les fibrรฉs construits dans le thรฉorรจme 22 peuvent admettre, si $`n2`$, un effondrement ร diamรจtre et courbure bornรฉs pour lequel $`\lambda `$ ne tend pas vers zรฉro.
Soit $`M`$ un tel fibrรฉ. Nous allons dโabord montrer le lemme suivant qui nous permettra de contrรดler la courbure:
###### Lemme 3.15
Pour toute mรฉtrique homogรจne $`g`$ sur $`M`$, la courbure sectionnelle de $`M`$ vรฉrifie $`|K(M,g)|\frac{3}{4}|V|^2`$.
Dรฉmonstation : On se place dans la mรชme base $`(V_1,\mathrm{},V_n,Y_1,Y_2)`$ que celle utilisรฉe dans la dรฉmonstration de 22.2. De (3.11), on peut rapidement dรฉduire que $`ad_{V_i}=0`$, et que $`ad_{Y_i}^{}Y_j=0`$ car pour tout vecteur $`U`$, $`ad_{Y_i}^{}Y_j,U=Y_j,ad_{Y_i}U=Y_j,[Y_i,U]=0`$. De plus, comme $`ad_{Y_i}^{}V_j,U=V_j,[Y_i,U]`$, on aura $`ad_{Y_i}^{}V_j=0`$ pour $`j1`$, $`ad_{Y_1}^{}V_1=\mu Y_2`$ et $`ad_{Y_2}^{}V_1=\mu Y_1`$.
La formule (2.3) donne donc:
$$K(V_i,V_j)=0,K(Y_1,Y_2)=\frac{3}{4}[Y_1,Y_2]^2,$$
(3.16)
$$K(V_1,Y_i)=\frac{\mu ^2}{4},\text{ et }K(V_i,Y_j)=0\text{ pour }i1.$$
(3.17)
Comme dโune part $`\mu ^2=V^2`$, et dโautre part $`V=[Y_1,Y_2]`$, la majoration du lemme en dรฉcoule immรฉdiatement.
Un effondrement ร base fixe du fibrรฉ $`M`$ par des mรฉtriques homogรจnes est dรฉterminรฉ par une famille de bases $`(V_1^\epsilon ,\mathrm{},V_n^\epsilon )`$ de $`\mathrm{\Gamma }(T_VM)^^G`$ (remarque: on ne suppose plus ici que $`V_1^\epsilon `$ est colinรฉraire ร $`V`$). Nous allons prรฉsenter ici des exemples dโeffondrements associรฉs ร des familles de bases de la forme $`(V_1^\epsilon ,\mathrm{},V_n^\epsilon )=(\epsilon ^{\alpha _1}V_1,\mathrm{},\epsilon ^{\alpha _n}V_n)`$, oรน $`(V_1,\mathrm{},V_n)`$ est une base fixรฉe. Si $`b_i`$ et $`b_i^\epsilon `$ sont les coefficients de $`V`$ dans les bases respectives $`(V_1,\mathrm{},V_n)`$ et $`(V_1^\epsilon ,\mathrm{},V_n^\epsilon )`$, on aura
$$b_i^\epsilon =\epsilon ^{\alpha _i}b_i.$$
(3.18)
Exemple 3.19. Si $`\alpha _i>0`$, pour tout $`i`$, alors le diamรจtre de la fibre tend vers $`0`$, ainsi que les $`b_i^\epsilon `$. On a donc un effondrement ร courbure bornรฉe du fibrรฉ sur la base $`T^2`$, et la valeur propre $`\lambda =V^2=_{i=1}^n(b_i^\epsilon )^2`$ tend vers zรฉro.
On peut cependant construire des effondrements pour lesquels le comportement du spectre est diffรฉrent, et en particulier tels quโil nโy ait pas de petites valeurs propres :
Exemple 3.20. Supposons que $`\alpha _i=0`$, pour tout $`i>1`$, $`\alpha _1>0`$, et que les composantes de $`V_1`$ dans la base $`(X_1,\mathrm{},X_n)`$ soient irrationnelles entre elles. Une droite de la fibre de direction $`V_1`$ sera donc dense dans la fibre, et par consรฉquent, il suffit que seul $`\alpha _1`$ soit non nul pour que la fibre sโeffondre sur un point. On aura alors $`b_i^\epsilon =b_i`$ pour $`i>1`$, et $`b_1^\epsilon 0`$. La courbure reste donc bornรฉe et $`\lambda _{i>1}b_i^20`$.
Ce dernier exemple justifie remarque Introduction faite dans lโintroduction.
### 3.3 Exemples de fibrรฉs principaux sur des bases de dimension strictement supรฉrieures ร 2
Les exemples du thรฉorรจme 22 sont relativement simples, en ce sens que lโessentiel de la topologie est contenu dans la structure de fibrรฉ en cercle de la nilvariรฉtรฉ $`N`$. On peut cependant facilement contruire des fibrรฉs principaux en tore dont le comportement du spectre est un peu plus riche:
Exemple 3.21. Considรฉrons deux fibrรฉs principaux $`M_1`$ et $`M_2`$, de fibre respective $`T^{k_1}`$ et $`T^{k_2}`$ et de base $`T^2`$, muni dโune structure homogรจne et dโune mรฉtrique invariante. Soient $`\lambda _1`$ et $`\lambda _2`$ leur valeur propre associรฉe dรฉfinie par le thรฉorรจme 22. La variรฉtรฉ $`M`$ dรฉfinie par le produit riemannien $`M=M_1\times M_2`$ est un fibrรฉ principal de fibre $`T^{k_1}\times T^{k_2}=T^{k_1+k_2}`$ et de base $`T^2\times T^2=T^4`$. Dโaprรจs la formule de Kรผnneth, il admet $`\lambda _1`$ et $`\lambda _2`$ comme valeurs propres. En choissant sur $`M_1`$ et $`M_2`$ des suites de mรฉtriques qui effondrent ces variรฉtรฉs sur $`T^2`$ la suite des mรฉtriques produits effondre $`M`$ sur $`T^4`$. On voit que sur cet exemple, on peut avoir deux valeurs propres non nulles distinctes en restriction aux formes invariantes. De plus, on peut choisir des effondrements sur $`M_1`$ et $`M_2`$ tels que ces deux valeurs propres tendent vers zรฉro ร des vitesses diffรฉrentes, ou mรชme, en choisissant pour $`M_1`$ et $`M_2`$ les effondrements dรฉcrits dans les exemples 3.2 et 3.2 respectivement, tels que seule lโune de ces valeurs propres tende vers zรฉro. Enfin, on peut remarquer que considรฉrer un effondrement de $`M`$ par homothรฉtie de la fibre revient ร considรฉrer des homothรฉtie des fibres de $`M_1`$ et $`M_2`$, et quโon retrouve le fait remarquรฉ en Introduction que cet effondrement produit des petites valeurs propres.
## Chapitre 4 Topologie des fibrรฉs principaux en tores
Nous allons dans cette partie nous attacher ร dรฉcrire la topologie des fibrรฉs principaux en tore, et en particulier ร construire un invariant diffรฉrentiel qui permettra, comme la classe dโEuler dans la cas des fibrรฉs en cercles, dโรฉtudier le comportement du spectre du laplacien lors dโun effondrement.
Soit $`M`$ un fibrรฉ principal en tore $`T^k`$ sur une base $`N`$. Le tore $`T^k=^k/^k`$ peut sโรฉcrire comme le produit de $`k`$ cercles : $`T^k=_{i=1}^kS_{(i)}^1`$. Lโaction de $`T^k`$ sur $`M`$ induit une action de chacun des $`S_{(i)}^1`$. On peut donc dรฉfinir les variรฉtรฉs
$$M_i=M/\underset{ji}{}S_{(j)}^1,$$
(4.1)
chaque $`M_i`$ รฉtant un fibrรฉ en cercle de base $`N`$ sur lequel agit $`S_{(i)}^1`$. Rรฉciproquement, la donnรฉe des $`k`$ fibrรฉs en cercles $`(M_i)_{ik}`$ sur $`N`$ permet de construire un fibrรฉ en tore $`T^k`$ en prenant la somme de Whitney $`_{i=1}^kM_i`$ de ces fibrรฉs en cercles, ce fibrรฉ en tore รฉtant diffรฉomorphe au fibrรฉ $`M`$. Comme la structure dโun fibrรฉ en cercle est dรฉterminรฉ par sa classe dโEuler, la topologie de $`M`$ est dรฉterminรฉe par la donnรฉe dโun $`k`$-uplet $`(e_1,\mathrm{},e_k)(H^2(N,))^k`$ de classes dโEuler. Cependant, la dรฉcomposition de $`T^k`$ en produit de cercles nโest pas unique. En effet, pour chaque base $`(a_1,\mathrm{},a_k)`$ du rรฉseau $`^k`$, on peut รฉcrire $`T^k`$ comme le produit de la famille de cercles $`(a_i/a_i)_i`$, auquel correspond en gรฉnรฉral un $`k`$-uplet diffรฉrent de classes dโEuler.
En homologie simpliciale, on peut dรฉfinir la classe dโobstruction $`[c]`$ dโun fibrรฉ $`FMN`$, oรน $`[c]`$ est un รฉlรฉment de $`H_2(N,\pi _1(F))`$ qui est une mesure de lโobstruction du fibrรฉ ร admettre une section (voir \[St51\], ยง35). Si la fibre $`F`$ est un tore, les groupes dโhomotopies $`\pi _n(F)`$ sont triviaux pour $`n2`$, et par consรฉquent cette classe dโobstruction est nulle si et seulement si le fibrรฉ admet une section (\[St51\], ยง29 et ยง35) ce qui, dans le cas dโun fibrรฉ principal est รฉquivalent ร รชtre trivial. Dans le cas dโune fibre $`T^k`$, cette classe dโobstruction est un รฉlรฉment de $`H_2(N,^k)`$. On veut dรฉfinir un objet semblable pour la cohomologie de Rham.
Dans le cas dโun fibrรฉ en cercle de classe dโEuler $`[e]`$, on a la propriรฉtรฉ suivante (\[BT82\], p.72) : si $`\omega `$ est une $`1`$-forme verticale invariante dont lโintรฉgrale sur chaque fibre vaut $`1`$, alors $`\mathrm{d}\omega `$ est une $`2`$-forme horizontale qui dรฉpend du choix de la connexion sur le fibrรฉ, mais qui est, au signe prรจs, le relevรฉ dโune รฉlรฉment de $`[e]`$. Dans le cas dโun fibrรฉ en tore, on va construire une invariant qui gรฉnรฉralise cette propriรฉtรฉ.
Rappelons tout dโabord que si on note $`๐ข`$ lโalgรจbre de Lie de $`T^k`$, lโaction de $`T^k`$ sur le fibrรฉ $`M`$ induit un plongement de $`๐ข`$ dans lโespace des champs de vecteurs verticaux invariants de $`M`$. Ce plongement sโรฉtend naturellement aux tenseurs sur $`๐ข`$ en une application
$$(\stackrel{p}{}๐ข)(\stackrel{q}{}๐ข^{})\mathrm{\Gamma }\left((\stackrel{p}{}T^VM)(\stackrel{q}{}T_{}^{V}{}_{}{}^{}M)\right)=\mathrm{\Gamma }(T_{}^{V}{}_{q}{}^{p}M).$$
Le tenseur obtenu ne dรฉfinit pas de maniรจre canonique un รฉlรฉment de $`\mathrm{\Gamma }(T_q^pM)`$ si $`q0`$, mais si on se donne une connexion sur le fibrรฉ $`M`$, la partie covariante du tenseur est bien dรฉfinie sur $`TM`$ en imposant ร sa partie horizontale dโรชtre nulle. Par abus de langage, si on se donne par exemple un รฉlรฉment de $`๐ข^{}`$, on dira quโil \<\< induit \>\> une $`1`$-forme verticale sur $`M`$, en prรฉcisant la connexion utilisรฉe sโil y a ambiguรฏtรฉ.
On va montrer le rรฉsultat suivant, qui permet de dรฉfinir une gรฉnรฉralisation de la classe dโEuler aux fibrรฉ principaux en tore:
###### Proposition 4.2
Soit $`\overline{\omega }๐ข^{}`$, $`\omega `$ la $`1`$-forme diffรฉrentielle sur $`M`$ induite par $`\overline{\omega }`$ et $`\alpha _\omega `$ la $`2`$-forme diffรฉrentielle sur $`N`$ telle que $`\mathrm{d}\omega =\pi ^{}(\alpha _\omega )`$. Alors lโapplication $`e:๐ข^{}H^2(N,)`$ donnรฉe par $`\overline{\omega }[\alpha _\omega ]`$ est bien dรฉfinie (c.-ร -d. que la classe de cohomologie de $`\alpha _\omega `$ ne dรฉpend pas du choix de la connexion) et linรฉaire.
Dรฉmonstration: On pose $`T^k=^k/^k`$ et on note $`(\overline{\omega }_i)_{i[1,k]}`$ les formes coordonnรฉes de $`^k`$ passรฉes au quotient sur $`T^k`$. En utilisant la dรฉcomposition $`T^k=_{i=1}^kS_{(i)}^1`$, on dรฉfinit la famille de fibrรฉs en cercle $`(M_i)`$ comme en (4.1). On a alors des projections $`M\stackrel{\pi _i}{}M_i\stackrel{\pi _i^{}}{}N`$ qui vรฉrifient $`\pi _i^{}\pi _i=\pi _j^{}\pi _j=\pi `$. Chaque forme $`\omega _i`$ induite sur $`M`$ par $`\overline{\omega }_i`$ est le relevรฉ $`\pi _i^{}(\omega _i^{})`$ de la forme de connexion du fibrรฉ en cercle $`M_i`$. On peut รฉcrire $`\mathrm{d}\omega _i=\pi _i^{}(\mathrm{d}\omega _i^{})=\pi ^{}(e_i)`$$`e_i`$ est une $`2`$-forme sur $`N`$. Or, on sait que $`e_i`$ appartient ร la classe dโEuler du fibrรฉ $`M_i`$, indรฉpendamment du choix de la connexion sur $`M_i`$ donc du choix de la connexion sur $`M`$. Si on dรฉfinit $`e:๐ข^{}H^2(N,)`$ par $`e(\overline{\omega }_i)=e_i`$ en lโรฉtendant par linรฉaritรฉ, on aura bien, par linรฉaritรฉ de la diffรฉrentielle extรฉrieure, $`\mathrm{d}\omega =\pi ^{}(e(\overline{\omega }))`$ pour tout $`\overline{\omega }๐ข^{}`$, en notant $`\omega `$ la forme induite sur $`M`$.
Remarque 4.3. Si $`k=1`$ et si $`\overline{\omega }`$ est la forme volume du cercle de longueur $`1`$, alors $`e(\overline{\omega })`$ est la classe dโEuler du fibrรฉ.
Remarque 4.4. La dรฉmonstration 4.2 met en รฉvidence le lien entre lโinvariant $`e`$ et la famille de classe dโEuler associรฉe ร une dรฉcomposition particuliรจre en somme de Whitney de $`M`$: ร chaque dรฉcomposition possible est associรฉe une base de $`๐ข^{}`$, et la famille de classe dโEuler est lโimage de cette base par $`e`$.
Exemple 4.5. Au chapitre 3, on a considรฉrรฉ des fibrรฉs principaux en tore $`T^k`$ non triviaux dont la base est un tore $`T^2`$. Comme $`H^2(T^2)`$ est de dimension $`1`$, le noyau de $`e`$ est de dimension $`k1`$, ce qui signifie quโon peut dรฉcomposer le fibrรฉ en une somme de Whitney de $`k`$ fibrรฉs en cercles dont $`k1`$ sont triviaux. On retrouve donc le fait que le fibrรฉ peut sโรฉcrire comme le produit dโun fibrรฉ en cercle et dโun tore de dimension $`k1`$.
Dans la suite, si $`\omega `$ est une forme induite par un รฉlรฉment $`\overline{\omega }`$ de $`๐ข^{}`$, on รฉcrira parfois par abus de langage \<\< $`e(\omega )`$ \>\> au lieu de \<\< $`e(\overline{\omega })`$ \>\>. De plus, verra parfois $`e`$ comme une application de $`๐ข^{}`$ dans lโespace $`^2(N,h)`$ des $`2`$-formes harmoniques de $`N`$, en utilisant le fait que $`^2(N,h)`$ est canoniquement isomorphe ร $`H^2(N,)`$.
## Chapitre 5 Formes invariantes et petites valeurs propres
Nous allons ici dรฉmontrer les rรฉsultats 31, 33 et 35 รฉnoncรฉs dans lโintroduction. Ceux-ci sโappuient essentiellement sur le
###### Lemme 5.1
Soit $`(M,g)`$ une variรฉtรฉ riemannienne, $`\phi _t`$ un flot agissant par isomรฉtrie sur $`M`$ et $`X`$ le champ de vecteur associรฉ. On suppose de plus que $`X`$ nโest pas uniformรฉment nul.
Soit $`E`$ un sous-espace de $`\mathrm{\Omega }^{}(M)`$ stable par le laplacien et par $`(\phi _t^{})_t`$, $`\lambda `$ une valeur propre du laplacien restreint ร $`E`$, et $`E_\lambda `$ lโespace propre associรฉ.
Sโil existe $`T`$ tel que $`\phi _{t+T}^{}=\phi _t^{}`$ pour tout $`t`$ et que $`\lambda <\left({\displaystyle \frac{2\pi }{TX_{\mathrm{}}}}\right)^2`$, alors $`(\phi _t^{})_t`$ agit trivialement sur $`E_\lambda `$.
Dรฉmonstration du lemme 5.1 : Remarquons tout dโabord quโon peut supposer que $`T=2\pi `$, le rรฉsultat gรฉnรฉral se dรฉduisant par simple changement de variable. Dโautre part, par thรฉorie de Hodge, on peut se restreindre ร lโรฉtude des formes cofermรฉes. On supposera donc que $`\lambda `$ est une valeur propre de $`\delta d_{|Ker\delta }`$ et $`E_\lambda `$ dรฉsignera le sous-espace propre associรฉ dans $`Ker\delta `$.
Soit $`\omega E_\lambda `$. On sait que $`_X\omega =i_X\mathrm{d}\omega +\mathrm{d}i_X\omega `$.
Dโune part, on a $`_X\omega =\underset{t0}{lim}{\displaystyle \frac{\phi _t^{}\omega \omega }{t}}`$. Puisque $`\phi _t`$ est une isomรฉtrie, les formes $`\phi _t^{}\omega `$ et $`\frac{\phi _t^{}\omega \omega }{t}`$ sont dans $`Ker\delta `$, et donc $`_X\omega `$ aussi car $`Ker\delta `$ est fermรฉ. Dโautre part, $`\mathrm{d}i_X\omega `$ est dans $`Im\mathrm{d}`$.
Comme $`Ker\delta `$ et $`Im\mathrm{d}`$ sont orthogonaux, les formes $`_X\omega `$ et $`\mathrm{d}i_X\omega `$ sont orthogonales. Le thรฉorรจme de Pythagore donne donc :
$$_X\omega _2^2+\mathrm{d}i_X\omega _2^2=i_X\mathrm{d}\omega _2^2.$$
(5.2)
On en dรฉduit
$$_X\omega _2^2i_X\mathrm{d}\omega _2^2i_X^2\mathrm{d}\omega _2^2i_X^2\lambda \omega _2^2.$$
(5.3)
On va maintenant majorer la norme de $`i_X`$ dโune part, et รฉvaluer celle de $`_X\omega `$ dโautre part.
##### majoration de $`i_X`$
Soit $`\alpha \mathrm{\Omega }^p(M)`$.
Soit $`mM`$, et $`(X_1,\mathrm{},X_n)`$ une base orthonormรฉe de $`T_mM`$ telle que $`X=\mu X_1`$. On pose $`\alpha _m={\displaystyle \underset{i_1<\mathrm{}<i_p}{}}\alpha _{i_1,\mathrm{},i_p}X_{i_1}^{\mathrm{}}\mathrm{}X_{i_p}^{\mathrm{}}`$. On a alors
$$i_X(\alpha _m)=\underset{i_1<\mathrm{}<i_p}{}\alpha _{i_1,\mathrm{},i_p}i_X(X_{i_1}^{\mathrm{}}\mathrm{}X_{i_p}^{\mathrm{}}).$$
(5.4)
Si $`i_1=1`$, alors $`i_X(X_{i_1}^{\mathrm{}}\mathrm{}X_{i_p}^{\mathrm{}})=\mu X_{i_2}^{\mathrm{}}\mathrm{}X_{i_p}^{\mathrm{}}`$.
Si $`i_11`$, alors $`i_X(X_{i_1}^{\mathrm{}}\mathrm{}X_{i_p}^{\mathrm{}})=0`$.
Donc $`i_X(\alpha _m)=\mu _{i_2<\mathrm{}<i_p}\alpha _{1,i_2,\mathrm{},i_p}X_{i_2}^{\mathrm{}}\mathrm{}X_{i_p}^{\mathrm{}}`$, et
$$|i_X(\alpha _m)|^2=\mu ^2\underset{i_2<\mathrm{}<i_p}{}\alpha _{1,i_2,\mathrm{},i_p}^2\mu ^2\underset{i_1<\mathrm{}<i_p}{}\alpha _{i_1,\mathrm{},i_p}^2=|X|^2|\alpha _m|^2.$$
On en dรฉduit
$`i_X(\alpha )_2^2`$ $`=`$ $`{\displaystyle _M}|i_X(\alpha _m)|^2dv{\displaystyle _M}|X|^2|\alpha _m|^2dvX_{\mathrm{}}{\displaystyle _M}|\alpha _m|^2dv`$
$``$ $`X_{\mathrm{}}\alpha _2^2.`$
et donc
$$i_XX_{\mathrm{}}$$
(5.5)
##### calcul de $`_x\omega `$
Comme $`\phi _t`$ est une isomรฉtrie, $`\phi _t^{}`$ agit par isomรฉtrie sur $`E_\lambda `$. Si on suppose que $`t\phi _t`$ est $`2\pi `$-pรฉriodique, $`\phi _t`$ induit donc un morphisme $`S^1SO(E_\lambda )`$ quโon peut dรฉcomposer en somme de reprรฉsentation irrรฉductibles.
Les reprรฉsentations irrรฉductibles de $`S^1`$ sont :
* la reprรฉsentation triviale $`S^1SO(1),tId`$;
* les rotations du plan $`S^1SO(2),tR(kt)=\left(\begin{array}{cc}\mathrm{cos}kt& \mathrm{sin}kt\\ \mathrm{sin}kt& \mathrm{cos}kt\end{array}\right),kZ^{}`$.
Supposons maintenant que la dรฉcomposition fasse apparaรฎtre au moins une rotation. On peut choisir une forme $`\omega 0`$ situรฉe dans le sous-espace stable associรฉ. On a alors
$$_X\omega =\underset{t0}{lim}\frac{\phi _t^{}\omega \omega }{t}=|k|\omega \omega .$$
(5.6)
Avec (5.3) et (5.5), on en dรฉduit :
$$\omega ^2\lambda X_{\mathrm{}}^2\omega ^2,$$
(5.7)
et finalement
$$\lambda \frac{1}{X_{\mathrm{}}^2}.$$
(5.8)
La conclusion du lemme en dรฉcoule immรฉdiatement.
Dรฉmonstration du thรฉorรจme 31 : Cโest une application directe du lemme 5.1 avec $`E=\mathrm{\Omega }^p(M)`$, le flot $`\phi _t`$ รฉtant induit par lโaction de $`S^1`$. Le champ $`X`$ est alors un champ vertical $`S^1`$-invariant.
Si on paramรจtre $`\phi _t`$ de maniรจre ร รชtre $`2\pi `$-pรฉriodique, la longueur $`l`$ dโune fibre sera รฉgale ร $`2\pi |X|`$, la norme de $`X`$ ne dรฉpendant pas du point choisi sur la fibre. On a donc
$$|X|=\frac{l}{2\pi }$$
(5.9)
et par consรฉquent
$$X_{\mathrm{}}=\frac{l_0}{2\pi }.$$
(5.10)
Si $`\lambda <\left(\frac{2\pi }{l_0}\right)^2`$, alors $`\lambda <\frac{1}{X_{\mathrm{}}^2}`$ et donc les formes propres de $`E_\lambda `$ sont $`S^1`$-invariantes.
Dรฉmonstration du thรฉorรจme 33 : Remarquons tout dโabord que dans la dรฉmonstration du thรฉorรจme 31, la dรฉcomposition de $`\mathrm{\Omega }^p(M)`$ en รฉlรฉments irrรฉductibles est une dรฉcomposition en sรฉrie de Fourier par rapport ร la fibre $`S^1`$, la reprรฉsentation triviale et les reprรฉsentations $`\theta R(k\theta )`$ correspondant respectivement aux fonctions constantes et aux fonctions $`\frac{2\pi }{k}`$-pรฉriodiques du cercle. On va reprendre cette idรฉe et lโappliquer au tore $`T^k`$ pour dรฉcomposer $`\mathrm{\Omega }^p(M)`$ en somme en sous-espaces de formes invariantes dans une direction et pรฉriodiques dans une autre pour ensuite appliquer le lemme 5.1 ร ces sous-espaces.
Plus prรฉcisรฉment, si on pose $`T^k=^k/^k`$, $`^k`$ รฉtant muni dโune mรฉtrique euclidienne (pas nรฉcessairement la mรฉtrique euclidienne canonique), et si $`\mathrm{\Gamma }`$ est le rรฉseau dual de $`^k`$, $`\mathrm{\Gamma }=\{\gamma ^k,\gamma ,\gamma ^{},\gamma ^{}^k\}`$, une base des fonctions propres de $`T^k`$ est donnรฉe par $`f_\gamma =\mathrm{cos}(2\pi \gamma ,x)`$ et $`g_\gamma =\mathrm{sin}(2\pi \gamma ,x)`$, $`\gamma \mathrm{\Gamma }`$ (voir par exemple \[GHL87\], p. 200). Si on note $`\gamma ^{}`$ lโorthogonal de $`\gamma `$ dans $`^n`$, les fonctions $`f_\gamma `$ et $`g_\gamma `$ sont invariantes sous lโaction de $`\gamma ^{}`$. On peut remarquer que, si $`\mathrm{\Gamma }`$ dรฉpend de la mรฉtrique choisie sur $`^k`$, ce nโest pas le cas de lโensemble des $`\gamma _{\gamma 0}^{}`$. En effet, cโest lโensemble des hyperplans vectoriels de $`^k`$ engendrรฉs par des รฉlรฉments de $`^k`$. Notons $`A`$ cet ensemble. En regroupant les fonctions propres en fonction de leur direction invariante, on obtient la dรฉcomposition
$$C^{\mathrm{}}(T^k)=\overline{\underset{VA}{}C^{\mathrm{}}(T^k)^V}$$
(5.11)
$``$ reprรฉsente les fonctions constantes, et oรน $`C^{\mathrm{}}(T^k)^V`$ est lโespace des fonctions $`C^{\mathrm{}}`$ dโintรฉgrale nulle (cโest-ร -dire orthogonales aux fonctions constantes) et invariantes dans la direction $`V`$. Par construction, chaque $`C^{\mathrm{}}(T^k)^V`$ est invariant par $`\mathrm{\Delta }`$ et par lโaction de $`T^k`$.
De la mรชme faรงon, $`\mathrm{\Omega }^p(M)`$ peut sโรฉcrire
$$\mathrm{\Omega }^p(M)=\overline{\underset{VA}{}\mathrm{\Omega }^p(M)^V\mathrm{\Omega }^p(M)^{T^k}},$$
(5.12)
$`\mathrm{\Omega }^p(M)^{T^k}`$ est lโespace des $`p`$-formes $`T^k`$-invariantes et $`\mathrm{\Omega }^p(M)^V`$ lโespace des $`p`$-formes diffรฉrentielles invariantes par lโaction de $`V`$ et orthogonales ร $`\mathrm{\Omega }^p(M)^{T^k}`$. Comme $`V`$ ne dรฉpend pas de la mรฉtrique sur chaque fibre, chacun des $`\mathrm{\Omega }^p(M)^V`$ est bien dรฉfini, et sera de plus invariant par $`\mathrm{\Delta }`$ et par lโaction de $`T^k`$. On a ainsi partiellement dรฉcomposรฉ les formes diffรฉrentielles de $`M`$ en sรฉrie de Fourier par rapport ร la fibre.
Soit $`VA`$, $`\lambda _M^V`$ une valeur propre du laplacien restreint ร $`\mathrm{\Omega }^p(M)^V`$ et $`\lambda _{T^k}^V`$ la premiรจre valeur propre de $`\mathrm{\Omega }^p(T^k)^V`$ ( remarque : elle est non nulle car les formes harmoniques du tore plat sont les formes invariantes ). On choisit sur $`T^k`$ un champ invariant $`X_{T^k}`$ orthogonal ร $`V`$ et on note $`X_M`$ le champ vertical sur $`M`$ induit par $`X_{T^k}`$. Lโaction sur $`\mathrm{\Omega }^p(M)^V`$ du flot $`\phi _t`$ associรฉ ร $`X_M`$ est pรฉriodique, et si on note $`T`$ sa pรฉriode, $`\lambda _{T^k}^V`$ est par construction exactement $`\left({\displaystyle \frac{2\pi }{TX_{T^k}_{\mathrm{}}}}\right)^2`$. Dโautre part, comme $`\overline{g}_xf(x)\overline{g}`$, les normes de $`X_{T^k}_{\mathrm{}}`$ et $`X_M_{\mathrm{}}`$ sont liรฉs :
$$X_M_{\mathrm{}}(\underset{B}{sup}f)^{1/2}X_{T^k}_{\mathrm{}},$$
(5.13)
donc si $`\lambda _M^V<(sup_Bf)^1\lambda _{T^k}^V`$, alors
$$\lambda _M^V<(\underset{B}{sup}f)^1\left(\frac{2\pi }{TX_{T^k}_{\mathrm{}}}\right)^2\left(\frac{2\pi }{TX_M_{\mathrm{}}}\right)^2,$$
(5.14)
et le lemme 5.1 sโapplique avec $`E=\mathrm{\Omega }^p(M)^V`$ et les formes propres associรฉes ร $`\lambda _M^V`$ sont $`\phi _t^{}`$-invariantes, donc $`T^k`$-invariantes.
Si $`\lambda _M`$ est une valeur propre du laplacien agissant sur $`\mathrm{\Omega }^p(M)`$ et que $`\lambda `$ est strictement infรฉrieure ร la premiรจre valeur propre de $`T^k`$, alors elle sera *a fortiori* infรฉrieure ร tous les $`\lambda _{T^k}^V`$ et donc les formes propres de $`\lambda _M^V`$ dans $`\mathrm{\Omega }^p(M)^V`$ sont $`T^k`$-invariantes. Comme les $`\mathrm{\Omega }^p(M)^V`$ sont stables par le laplacien, lโespace propre de $`\lambda _M`$ est la somme des espaces propres restreints aux $`\mathrm{\Omega }^p(M)^V`$. Par consรฉquent, toutes les formes propres associรฉes ร $`\lambda _M`$ sont $`T^k`$-invariantes.
Dรฉmonstration du corollaire 35 : Lโhypothรจse sur la mรฉtrique peut sโรฉcrire $`\overline{g}_x=f(x)\overline{g}`$, oรน $`f`$ est une fonction positive sur $`N`$. On peut alors appliquer le thรฉorรจme 33. Il reste ร montrer que si $`\lambda <\left({\displaystyle \frac{\pi }{d_0}}\right)^2`$ alors $`\lambda <(\underset{xB}{sup}f(x))^1\lambda _{\mathrm{0,1}}(T^k,\overline{g})`$.
Comme la mรฉtrique resteinte ร la fibre $`\pi ^1(x)`$ est $`f(x)\overline{g}`$, la premiรจre valeur propre de laplacien restreint ร cette fibre est $`\frac{\lambda _{\mathrm{0,1}}(T^k,\overline{g})}{f(x)}`$. De plus, la premiรจre valeur propre dโun tore plat de diamรจtre $`d`$ est minorรฉe par $`(\frac{\pi }{d})^2`$, par consรฉquent
$$\frac{\lambda _{\mathrm{0,1}}(T^k,\overline{g})}{f(x)}=\lambda _{\mathrm{0,1}}(T^k,f(x)\overline{g})\left(\frac{\pi }{d_x}\right)^2,$$
(5.15)
$`d_x`$ est le diamรจtre de la fibre $`\pi ^1(x)`$ pour la distance intrinsรจque, et donc
$$\frac{\lambda _{\mathrm{0,1}}(T^k,\overline{g})}{sup_Bf}\left(\frac{\pi }{d_0}\right)^2,$$
(5.16)
ce qui achรจve la dรฉmonstration.
Lโexemple suivant montre que si on ne suppose pas que les fibres sont homothรฉtiques entre elles, une majoration du diamรจtre des fibres ne permet pas de majorer la fonction $`f`$ du thรฉorรจme 33.
Exemple 5.17. On considรจre sur $`^2`$ muni de son systรจme de coordonnรฉs canonique la famille de mรฉtrique $`g_t=(\mathrm{d}x+t\mathrm{d}y)^2+\mathrm{d}y^2`$. Ces mรฉtriques passent au quotient sur le tore $`T^2=^2/^2`$. Quel que soit $`t_0`$, il nโexiste pas de constante $`c>0`$ telle que $`g_tcg_{t_0}`$ pour tout $`t`$. Cependant, le diamรจtre de $`(T^2,g_t)`$ reste bornรฉ quand $`t`$ varie. En effet, le diffรฉomorphisme linรฉaire $`\left(\begin{array}{cc}1& 1\\ 0& 1\end{array}\right)`$ est une isomรฉtrie de $`(T^2,g_t)`$ dans $`(T^2,g_{t+1})`$, et par consรฉquent le diamรจtre de $`(T^2,g_t)`$ est une fonction pรฉriodique de $`t`$.
Dรฉmonstration de la remarque Introduction : Il suffit de remarquer que dans le lemme 5.1, si la dimension du sous-espace propre $`E_\lambda `$ est impaire, la dรฉcomposition de cet espace en espace de reprรฉsentations irrรฉductibles contient nรฉcessairement une reprรฉsentation triviale, et donc $`E_\lambda `$ contient des formes invariantes par $`\phi _t`$. Dans le cas du thรฉorรจme 31, les espaces propres de dimension impaire du laplacien agissant sur $`\mathrm{\Omega }^{}(M)`$ contiennent donc des formes $`S^1`$-invariantes.
Dans le cas du thรฉorรจme 33, si un espace propre du laplacien est de dimension impaire, lโun des รฉlรฉments de la dรฉcomposition de Fourier de cet espace sera aussi de dimension impaire, et la remarque prรฉcรฉdente sโapplique.
## Chapitre 6 Gรฉomรฉtrie des fibrรฉs principaux en tores
### 6.1 Mรฉtriques adaptรฉes
Nous allons ici montrer quโon peut, dans le but dโobtenir le thรฉorรจme 26, se ramener ร une situation gรฉomรฉtrique pour laquelle lโรฉtude du spectre dโun fibrรฉ en tore est plus simple. Cette situation est une gรฉnรฉralisation de la notion de mรฉtrique adaptรฉe dรฉfinie dans le cas des fibrรฉs en cercle par B. Colbois et G. Courtois (\[CC00\]) :
###### Dรฉfinition 6.1
On dit que le couple de mรฉtriques $`(g,h)`$ dรฉfinies sur $`M`$ et $`N`$ respectivement est adaptรฉ ร la fibration principale $`T^kM^n\stackrel{\pi }{}N`$ si:
1. $`\pi :(M,g)(N,h)`$ est une submersion riemannienne ;
2. Lโaction de $`T^k`$ sur $`M`$ est isomรฉtrique ;
3. Les fibres sont totalement gรฉodรฉsiques ;
4. Toute $`1`$-forme verticale $`\omega `$ induite par un รฉlรฉment de $`๐ข^{}`$ vรฉrifie $`\mathrm{d}\omega =\pi ^{}(e(\omega ))`$.
On veut montrer quโune mรฉtrique de courbure bornรฉe sur un fibrรฉ principal en tore est proche dโune mรฉtrique adaptรฉe:
###### Thรฉorรจme 6.2
Soient $`a`$ et $`d`$ deux rรฉels strictement positifs, et $`T^k(M^n,g)\stackrel{\pi }{}(N,h)`$ un fibrรฉ principal en tore. Il existe des constantes $`\epsilon _0(n,a,d,(N,h))>0`$, $`\tau (n,a,d,(N,h))>0`$, $`\tau ^{}(n,a,d,(N,h))>0`$ et $`c(n,a,d,(N,h))>0`$ telles que si $`|K(N,h)|a`$, $`|K(M,g)|a`$, $`diam(M,g)d`$ et si $`\pi `$ est une $`\epsilon `$-approximation de Hausdorff avec $`\epsilon <\epsilon _0`$, alors il existe des mรฉtriques $`\stackrel{~}{g}`$ et $`\stackrel{~}{h}`$ sur $`M`$ et $`N`$ respectivement et une fibration $`\pi ^{}:(M,\stackrel{~}{g})(N,\stackrel{~}{h})`$ telles que
1. Le couple $`(\stackrel{~}{g},\stackrel{~}{h})`$ est adaptรฉ ร la fibration $`\pi ^{}`$;
2. $`{\displaystyle \frac{1}{\tau }}g\stackrel{~}{g}\tau g`$ et $`{\displaystyle \frac{1}{\tau }}h\stackrel{~}{h}\tau h`$;
3. La restriction de $`\stackrel{~}{g}`$ ร la fibre est telle que $`diam(\pi ^1(x))\tau ^{}\epsilon `$, pour tout $`xN`$;
4. La courbure sectionnelle de $`(M,\stackrel{~}{g})`$ vรฉrifie $`|K(X,Y)|c`$, pour toute paire de vecteurs horizontaux orthonormรฉs $`(X,Y)`$.
On pourra alors appliquer le rรฉsultat de J. Dodziuk selon lequel si deux mรฉtriques sont proches, alors les spectres du laplacien pour ces deux mรฉtriques sont proches aussi:
###### Thรฉorรจme 6.3 (\[Do82\])
Soit $`g`$ et $`\stackrel{~}{g}`$ deux mรฉtriques riemanniennes sur une variรฉtรฉ $`M`$ de dimension $`n`$, et $`\tau `$ une constante positive. Si les deux mรฉtriques vรฉrifient $`\frac{1}{\tau }g\stackrel{~}{g}\tau g`$, alors
$$\frac{1}{\tau ^{3n1}}\lambda _{p,k}(M,g)\lambda _{p,k}(M,\stackrel{~}{g})\tau ^{3n1}\lambda _{p,k}(M,g),$$
pour tout entiers $`k0`$ et $`p[0,n]`$.
Remarque 6.4. Le thรฉorรจme 6.2 implique en particulier le thรฉorรจme 37. La restriction sur la gรฉomรฉtrie imposรฉe par le point 4 de la conclusion du thรฉorรจme 6.2 permet de mieux contrรดler le spectre du laplacien.
Remarque 6.5. En vertu dโun thรฉorรจme de Hermann (\[He60\], \[Be87\] p. 249), le fait que les fibres soient totalement gรฉodรฉsiques implique quโelles sont isomรฉtriques entre elles. On va voir dans la dรฉmonstration du thรฉorรจme 6.2 que rรฉciproquement, sur les fibrรฉs considรฉrรฉs, si la mรฉtrique est invariante et que les fibres sont isomรฉtriques alors elles sont totalement gรฉodรฉsiques.
### 6.2 Situation de mรฉtrique invariante
Nous allons dans un premier temps montrer que si on suppose quโon a sur $`M`$ une mรฉtrique invariante, elle est proche dโune mรฉtrique qui vรฉrifie les points (i) ร (iii) de la dรฉfinition 6.1. Plus prรฉcisรฉment :
###### Proposition 6.6
Soit $`T^k(M^n,g)\stackrel{\pi }{}(N,h)`$ un fibrรฉ principal en tore muni dโune mรฉtrique invariante $`g`$ tel que $`\pi `$ soit une submersion riemannienne. Pour tout $`a>0`$ et $`d>0`$, il existe des constantes $`\tau (n,a,d)>0`$ et $`c(n,a)>0`$ telles que si $`|K(N,h)|a`$, $`K(M,g)a`$ et $`diam(M,g)d`$, alors il existe une mรฉtrique invariante $`\stackrel{~}{g}`$ sur $`M`$ telle que la fibration $`\pi :(M,\stackrel{~}{g})(N,h)`$ soit une submersion riemannienne ร fibres totalement gรฉodรฉsiques et
$$\frac{1}{\tau }g\stackrel{~}{g}\tau g.$$
Remarque 6.7. On peut noter quโon utilise non pas une hypothรจse de courbure bornรฉe sur $`M`$ mais seulement que la courbure sectionnelle est minorรฉe.
Pour montrer la proposition 6.6, on utilisera les deux lemmes suivants. Le premier est une application directe de la formule de OโNeill:
###### Lemme 6.8
Soit $`a>0`$ et $`T^k(M^n,g)\stackrel{\pi }{}(N,h)`$ un fibrรฉ principal en tore muni dโune mรฉtrique invariante $`g`$ tel que $`p`$ soit une submersion riemannienne, $`|K(N,h)|a`$, et $`K_{(M,g)}(X,Y)a`$ pour tout couple $`(X,Y)`$ de vecteurs horizontaux orthonormรฉs. Alors, pour toute $`1`$-forme diffรฉrentielle verticale $`\omega `$ induite par lโaction de $`T^k`$, on a :
1. $`|\mathrm{d}\omega (X,Y)|_x^2{\displaystyle \frac{8a}{3}}|\omega |_x^2`$, pour tout $`xM`$ et tout couple de vecteurs horizontaux orthonormรฉs $`X`$ et $`Y`$ ;
2. $`\mathrm{d}\omega _{\mathrm{}}{\displaystyle \frac{4an(n1)}{3}}\omega _{\mathrm{}}`$.
Dรฉmonstration: Soit $`xM`$, $`y=\pi (x)`$, $`\stackrel{~}{X}`$ et $`\stackrel{~}{Y}`$ deux champs de $`N`$ orthonormรฉs en $`y`$, et $`X`$ et $`Y`$ les relevรฉs de $`\stackrel{~}{X}`$ et $`\stackrel{~}{Y}`$ ร $`M`$. La formule de OโNeill (\[GHL87\] p. 127, \[Be87\] p. 241) donne
$$K_N(\stackrel{~}{X},\stackrel{~}{Y})=K_M(X,Y)+\frac{3}{4}\left|[X,Y]^V\right|^2,$$
(6.9)
$`[X,Y]^V`$ dรฉsigne la composante verticale de $`[X,Y]`$. Dโautre part on a, en utilisant le fait que $`\omega `$ est verticale,
$`\mathrm{d}\omega (X,Y)`$ $`=`$ $`X\omega (Y)Y\omega (X)\omega ([X,Y])`$ (6.10)
$`=`$ $`\omega ([X,Y]).`$
On en dรฉduit :
$`|\mathrm{d}\omega (X,Y)|_x^2`$ $`=`$ $`|\omega ([X,Y])|_x^2|\omega |_x^2|[X,Y]^V|_x^2`$ (6.11)
$``$ $`{\displaystyle \frac{4}{3}}|\omega |_x^2(K_y(\stackrel{~}{X},\stackrel{~}{Y})K_x(X,Y)).`$
Comme chacun des couples $`(\stackrel{~}{X},\stackrel{~}{Y})`$ et $`(X,Y)`$ est orthonormรฉ en $`x`$ et $`y`$, on a les majorations $`|K_y(\stackrel{~}{X},\stackrel{~}{Y})|a`$ et $`K_x(X,Y)a`$, et donc
$$|\mathrm{d}\omega (X,Y)|_x^2\frac{8a}{3}|\omega |_x^2.$$
(6.12)
Et comme lโinรฉgalitรฉ prรฉcรฉdente est vraie quel que soit le choix de $`(\stackrel{~}{X},\stackrel{~}{Y})`$, il en dรฉcoule finalement
$$|\mathrm{d}\omega |_x^2\frac{4an(n1)}{3}|\omega |_x^2\frac{4an(n1)}{3}\omega ^2,$$
(6.13)
ce qui achรจve la dรฉmonstration.
Le second lemme montre, dans le cas dโun fibrรฉ en cercle, quโร courbure bornรฉe, la longueur des fibres varie peu dโune fibre ร lโautre.
###### Lemme 6.14
Soit $`S^1(M^n,g)\stackrel{\pi }{}(N,h)`$ un fibrรฉ principal en cercle sur $`N`$, tel que $`g`$ soit invariante et $`\pi `$ soit une submersion riemannienne. Pour tout $`a>0`$ et $`d>0`$, il existe $`\tau (n,a,d)>0`$ tel que si $`|K(N,h)|a`$, $`K(M,g)a`$ et $`diam(M,g)d`$, alors pour tout $`x,yN`$, on a
$$\frac{1}{\tau }l_yl_x\tau l_y,$$
$`l_x`$ et $`l_y`$ dรฉsignent les longueurs des fibres au dessus de $`\pi ^1(x)`$ et $`\pi ^1(y)`$ respectivement.
Dรฉmonstration: On choisit sur le fibrรฉ $`M`$ une $`1`$-forme verticale $`\omega `$ dont lโintรฉgrale sur chaque fibre est รฉgale ร $`1`$. Soit $`U`$ le champ vertical induit par lโaction de $`S^1`$ qui vรฉrifie $`\omega (U)=1`$. La norme $`|U|`$ de ce champ est constante sur chaque fibre, et sโรฉcrit $`|U|=\pi ^{}f`$, oรน $`f`$ est une fonction sur $`N`$. De plus, en tout point $`x`$ de $`N`$, la norme $`f(x)`$ de $`U`$ est รฉgale ร la longueur de la fibre au dessus de $`x`$. On va montrer que $`f`$ est bornรฉe en fonction de $`a`$ et $`d`$. Remarque : $`\omega `$ nโest pas la forme duale de $`U`$ pour la mรฉtrique. Sa norme ponctuelle sur la fibre $`\pi ^1(x)`$ est $`|\omega |=\frac{1}{f}`$, et on a $`U^{\mathrm{}}=f^2\omega `$.
Soit $`xN`$, et $`\stackrel{~}{X}`$ un vecteur unitaire tangent ร $`N`$ en $`x`$. Soit $`\stackrel{~}{X}_i`$ une base orthonormรฉe de champs de vecteurs au voisinage de $`x`$, telle que $`D_{\stackrel{~}{X}_1}\stackrel{~}{X}_1=0`$ sur ce voisinage, et $`\stackrel{~}{X}_{1|x}=\stackrel{~}{X}`$. On relรจve cette base ร $`T^HM`$ en notant $`X_i=\pi ^{}(\stackrel{~}{X}_i)`$ et $`X=X_1`$. Ces champs vรฉrifient
$$[X_i,U]=0.$$
(6.15)
En effet, ces crochets de Lie sont dรฉterminรฉs par $`[X_i,U]=\frac{\mathrm{d}}{\mathrm{d}t}(\mathrm{\Phi }_t)_{}X_i`$, oรน $`\mathrm{\Phi }_t`$ est le flot induit par le champ $`U`$. Par dรฉfinition de $`U`$, ce flot correspond ร lโaction de $`S^1`$ sur $`M`$. Les crochets de Lie sont donc nuls, car les champs $`X_i`$ sont $`S^1`$-invariants. On notera par ailleurs $`U^{}`$ le champ de norme $`1`$ dรฉfini par $`U^{}=U/|U|`$.
On va calculer la courbure sectionnelle $`K(X,U)`$ en fonction de $`f`$ et de ses variations (Remarque: le champ $`U`$ nโest pas normรฉ, mais ce calcul est plus simple que si on utilise le champ $`U^{}`$). Cette courbure sโรฉcrit
$`K(X,U)`$ $`=`$ $`R(X,U)X,U`$ (6.16)
$`=`$ $`D_UD_XXD_XD_UXD_{[X,U]}X,U,`$
$`,`$ dรฉsigne la mรฉtrique. Le vecteur $`D_XX`$ est horizontal et vaut $`\pi ^{}(D_{\stackrel{~}{X}_1}\stackrel{~}{X}_1)`$ (voir \[Be87\] p. 239), et par consรฉquent $`D_XX=0`$ au voisinage de $`x`$. Comme dโautre part $`[X,U]=0`$, on est ramenรฉ ร calculer
$$K(X,U)=D_XD_UX,U.$$
(6.17)
Pour ce faire, on utilisera la formule suivante, qui caractรฉrise la connexion de Levi-Civita:
$`2D_{Z_1}Z_2,Z_3`$ $`=`$ $`Z_1Z_2,Z_3+Z_2Z_3,Z_1Z_3Z_1,Z_2`$ (6.18)
$`+[Z_1,Z_2],Z_3[Z_1,Z_3],Z_2[Z_2,Z_3],Z_1.`$
En utilisant lโorthogonalitรฉ de $`(X_1,\mathrm{},X_n,U)`$ et le fait que $`[X_i,U]`$=0, on obtient
$$2D_UX,U=XU,U=Xf^2=2f\mathrm{d}f(X)$$
et
$$2D_UX,X_i=[X,X_i],U,$$
et donc
$$D_UX=\mathrm{d}f(X)U^{}\frac{1}{2}\underset{i=1}{\overset{n}{}}[X,X_i],UX_i.$$
(6.19)
Le premier terme peut sโรฉcrire $`\mathrm{d}f(X)U^{}=\frac{\mathrm{d}f(X)}{f}U`$, par consรฉquent
$`D_X(\mathrm{d}f(X)U^{})`$ $`=`$ $`\left(X{\displaystyle \frac{\mathrm{d}f(X)}{f}}\right)U+{\displaystyle \frac{\mathrm{d}f(X)}{f}}D_XU`$
$`=`$ $`\left({\displaystyle \frac{(D_X\mathrm{d}f)(X)}{f}}+{\displaystyle \frac{\mathrm{d}f(D_XX)}{f}}{\displaystyle \frac{(\mathrm{d}f(X))^2}{f^2}}\right)U`$
$`+{\displaystyle \frac{\mathrm{d}f(X)}{f}}D_XU`$
$`=`$ $`{\displaystyle \frac{Hessf(X,X)}{f}}U{\displaystyle \frac{(\mathrm{d}f(X))^2}{f^2}}U+{\displaystyle \frac{\mathrm{d}f(X)}{f}}D_XU,`$
en utilisant le fait que $`D_XX=0`$. De plus, la relation (6.18) donne
$$2D_XU,U=X|U|^2=2f\mathrm{d}f(X),$$
et donc
$$D_X(\mathrm{d}f(X)U^{}),U=fHessf(X,X).$$
(6.21)
La dรฉrivation des termes suivants de (6.19) donne
$$D_X([X,X_i],UX_i)=(X[X,X_i],U)X_i+[X,X_i],UD_XX_i.$$
Quand on calcule le produit scalaire de cette expression avec $`U`$, le premier terme sโannule, et comme la relation (6.18) donne $`D_XX_i,U=[X,X_i],U`$, il reste
$$D_X([X,X_i],UX_i),U=[X,X_i],U^2.$$
(6.22)
La somme des รฉquations (6.21) et (6.22) donne
$$K(X,U)=fHessf(X,X)+\frac{1}{2}\underset{i=1}{\overset{n}{}}[X,X_i],U^2.$$
(6.23)
En normalisant le vecteur $`U`$, on obtient
$$K(X,U^{})=\frac{Hessf(X,X)}{f}+\frac{1}{2f^2}\underset{i=1}{\overset{n}{}}[X,X_i],U^2.$$
(6.24)
Les derniers termes peuvent รชtre majorรฉs en fonction de la courbure. En effet, on a
$$[X,X_i],U=U^{\mathrm{}}([X,X_i])=f^2\omega ([X,X_i])=f^2\mathrm{d}\omega (X,X_i),$$
(6.25)
et donc, en vertu du lemme 6.8,
$$[X,X_i],U^2f^4\frac{8a}{3}|\omega |^2=\frac{8a}{3}f^2.$$
(6.26)
Par hypothรจse, la courbure sectionnelle de $`M`$ est minorรฉe par $`a`$. On a donc finalement :
$$\frac{Hessf(X,X)}{f}\left(\frac{4n}{3}+1\right)a.$$
(6.27)
Soient $`x`$ et $`y`$ deux points de $`N`$, et $`\gamma `$ une gรฉodรฉsique minimisante joignant ces deux points. Notons $`\mu `$ la fonction dรฉfinie par
$$\mu (t)=\mathrm{ln}f\gamma (t).$$
(6.28)
En dรฉrivant $`\mu `$ par rapport ร $`t`$, on obtient :
$$\mu ^{}(t)=\frac{\mathrm{d}f(\gamma ^{}(t))}{f\gamma (t)}$$
(6.29)
et
$`\mu ^{\prime \prime }(t)`$ $`=`$ $`\left({\displaystyle \frac{\mathrm{d}f(\gamma ^{}(t))}{f\gamma (t)}}\right)^2+{\displaystyle \frac{\mathrm{d}f(D_{\gamma ^{}(t)}\gamma ^{}(t))+D\mathrm{d}f(\gamma ^{}(t),\gamma ^{}(t))}{f\gamma (t)}}`$ (6.30)
$`=`$ $`\mu ^{}(t)^2+{\displaystyle \frac{Hessf(\gamma ^{}(t),\gamma ^{}(t))}{f\gamma (t)}}.`$
On a donc, en vertu de la majoration (6.27) :
$$\mu ^{\prime \prime }(t)\frac{Hessf(\gamma ^{}(t),\gamma ^{}(t))}{f\gamma (t)}\left(\frac{4n}{3}+1\right)a.$$
(6.31)
Supposons que $`x`$ est un point oรน $`f`$, et donc $`\mu `$, atteint son minimum. On a alors $`\mu ^{}(0)=0`$ et donc, en remarquant que $`d`$ majore le diamรจtre de $`N`$,
$$\mu ^{}(t)\left(\frac{4n}{3}+1\right)at,$$
(6.32)
et
$$\mu (t)\frac{4n+3}{6}at^2\frac{4n+3}{6}ad^2.$$
(6.33)
Le rapport $`\frac{f(y)}{f(x)}`$ est donc majorรฉ par une constante $`\tau =\mathrm{exp}(\frac{4n+3}{6}ad^2)`$. Comme on a montrรฉ cette majoration en prenant pour $`x`$ un point oรน $`f`$ atteint son minimum, elle sera vraie *a fortiori* pour un $`x`$ quelconque.
Remarque : si les fibre sont isomรฉtriques, alors le champ $`U`$ est de norme constante. Il est aisรฉ de vรฉrifier ร lโaide de (6.18) que $`D_UU`$ est alors nul, cโest-ร -dire que les fibres sont totalement gรฉodรฉsiques.
Dรฉmonstration de la proposition 6.6: Le but est en fait de gรฉnรฉraliser le lemme 6.14 aux fibrรฉs en tore pour montrer que $`g`$ est proche dโune mรฉtrique pour laquelle toutes les fibres sont isomรฉtriques.
Soit $`\overline{U}๐ข`$ non nul et $`U`$ le champ vertical induit par $`\overline{U}`$ sur $`M`$. Soit $`x_0N`$. On choisit $`\overline{U}`$ de sorte que $`|U|=1`$ au dessus de $`x_0`$. De plus, on impose ร $`\overline{U}`$ dโavoir un coefficient directeur rationnel, cโest-ร -dire que $`\overline{U}`$ est colinรฉaire ร un vecteur de $`^k๐ข`$. Lโaction du flot associรฉ ร $`U`$ induit alors une fibration $`S^1M\stackrel{\pi }{}(M^{},g^{})`$. On peut contrรดler la courbure de ce fibrรฉ. En effet, utilisant la formule de OโNeill, on peut รฉcrire
$`K_M^{}(\stackrel{~}{X},\stackrel{~}{Y})`$ $`=`$ $`K_M(X,Y){\displaystyle \frac{3}{4}}|[X,Y]^V|^2`$ (6.34)
$`=`$ $`K_M(X,Y){\displaystyle \frac{3}{4}}{\displaystyle \frac{|\mathrm{d}\omega (X,Y)|^2}{|\omega |^2}},`$
$`\stackrel{~}{X}`$ et $`\stackrel{~}{Y}`$ sont deux vecteurs orthonormรฉs de $`M^{}`$, $`X`$ et $`Y`$ leurs relevรฉs respectifs sur $`M`$, et $`\omega `$ la $`1`$-forme induite par lโaction de $`T^k`$ telle que $`\omega (U)=1`$. Le lemme 6.8 permet de contrรดler le dernier terme en fonction de la courbure de $`M`$, et donc
$$K_M^{}(\stackrel{~}{X},\stackrel{~}{Y})a2a=3a.$$
(6.35)
Le lemme 6.14 assure alors quโil existe une constante $`\tau (n,a,d)`$ telle que
$$\frac{1}{\tau }|U|\tau .$$
(6.36)
Notons $`\stackrel{~}{g}`$ la mรฉtrique invariante sur $`M`$ obtenue en modifiant $`g`$ dans la direction verticale de sorte que les fibres soient isomรฉtriques ร $`\pi ^1(x_0)`$, et en conservant la distribution horizontale et la mรฉtrique horizontale associรฉes ร $`g`$. Pour cette nouvelle mรฉtrique, la norme de $`U`$ est uniformรฉment รฉgale ร 1. La relation (6.36) peut sโรฉcrire
$$\frac{1}{\tau ^2}\stackrel{~}{g}(U,U)g(U,U)\tau ^2\stackrel{~}{g}(U,U),$$
(6.37)
Par continuitรฉ, (6.37) sโรฉtend ร nโimporte quel vecteur vertical. Comme $`g`$ et $`\stackrel{~}{g}`$ sont identiques sur la direction horizontale, on aura finalement
$$\frac{1}{\tau ^2}\stackrel{~}{g}g\tau ^2\stackrel{~}{g}.$$
(6.38)
Pour conclure, remarquons que si la mรฉtrique sur $`M`$ est telle que les fibres soient isomรฉtriques, la fibration $`S^1M\stackrel{\pi }{}(M^{},g^{})`$ induite par le champ $`U`$ est ร fibre totalement gรฉodรฉsique. Par continuitรฉ, les fibres du fibrรฉ $`T^kMN`$ sont elles aussi totalement gรฉodรฉsiques.
### 6.3 Cas gรฉnรฉral
On va maintenant dรฉmontrer le thรฉorรจme 6.2. Pour ce faire, on va sโinspirer dโune dรฉmonstration dโun thรฉorรจme de Lott (\[Lo02b\], thรฉorรจme 2), qui utilise les rรฉsultats de \[CFG92\].
Soit $`g`$ une mรฉtrique sur $`M`$ vรฉrifiant les hypothรจse du thรฉorรจme 6.2. Tout dโabord, en utilisant un rรฉsultat de rรฉgularisation dโAbresch (\[CFG92\], thรฉorรจme 1.12), on construit une mรฉtrique $`g_1`$ sur $`M`$ telle que $`\frac{1}{\tau _1}gg_1\tau _1g`$, $`|K(M,g_1)|a`$ et $`D^iRA_i(n,a,\tau _1)`$, oรน $`\tau _1>1`$ est un rรฉel fixรฉ, et $`D`$ et $`R`$ dรฉsignent respectivement la dรฉrivรฉe covariante et le tenseur de courbure pour la mรฉtrique $`g_1`$.
On applique ensuite le thรฉorรจme 2.6 de \[CFG92\], qui assure lโexistence de constantes $`ฯต_0(n,(N,h))`$, $`\kappa (n,A)`$, $`\kappa ^{}(n,A,(N,h))`$ et $`\kappa _i(n,A,(N,h))`$ et dโune fibration $`\pi ^{}:(M,g_1)(N,h)`$ tels que si $`\pi `$ est une $`\epsilon `$-approximation de Hausdorff avec $`\epsilon <ฯต_0(n,(N,h))`$, alors :
* pour tout $`xN`$, le diamรจtre de $`\pi ^1(x)`$ pour la mรฉtrique $`g^{}`$ est infรฉrieur ร $`\kappa \epsilon `$;
* la seconde forme fondamentale de la fibre vรฉrifie $`II_{\pi ^1(x)}_{\mathrm{}}\kappa ^{}`$ pour tout $`xN`$;
* la submersion $`\pi ^{}`$ est $`\kappa _i`$-rรฉguliรจre, cโest-ร -dire que $`D^i\pi ^{}_{\mathrm{}}\kappa _i`$, pour tout $`i`$.
Enfin, pour une telle fibration $`\pi ^{}`$, les parties 3 et 4 de \[CFG92\] donnent la construction dโun mรฉtrique $`g_2`$ sur $`M`$ qui est $`T^k`$-invariante et telle que $`|D^i(g_2g_1)|c(n,A,(N,h),i)`$, pour tout $`i`$. Cette derniรจre รฉgalitรฉ assure lโexistence dโune constante $`\tau _2(n,A,d,(N,h))`$ telle que $`\frac{1}{\tau _2}g_1g_2\tau _2g_1`$, et permet aussi de contrรดler la courbure pour la mรฉtrique $`g_2`$.
On peut alors appliquer la proposition 6.6, qui nous donne une mรฉtrique $`g_3`$ qui vรฉrifie les points (i) ร (iii) de la dรฉfinition 6.1.
Pour obtenir la mรฉtrique $`\stackrel{~}{g}`$ du thรฉorรจme 6.2, il reste ร modifier la distribution horizontale de maniรจre ร ce que $`\stackrel{~}{g}`$ vรฉrifie le point (iv) de la dรฉfinition 6.1. Remarquons tout dโabord que comme lโapplication $`e:๐ข^{}^2(N)`$ est linรฉaire, il suffit de montrer lโรฉgalitรฉ $`\mathrm{d}\omega =\pi ^{}(e(\omega ))`$ pour une base de $`๐ข^{}`$. Soit $`(\omega _i)`$ une base de $`๐ข^{}`$ orthonormรฉe pour la mรฉtrique $`g_3`$. Pour chaque $`i`$, $`\mathrm{d}\omega _i`$ sโรฉcrit
$$\mathrm{d}\omega _i=\pi ^{}(\alpha _i+\mathrm{d}\beta _i),$$
(6.39)
$`\alpha _i`$ est une forme harmonique et $`\beta _i`$ une forme cofermรฉe. On dรฉfinit une nouvelle forme verticale $`\omega _i^{}=\omega _i\pi ^{}(\beta _i)`$. Cette forme vรฉrifie
$$\mathrm{d}\omega _i^{}=\mathrm{d}\omega _i\pi ^{}(\mathrm{d}\beta _i)=\pi ^{}(\alpha _i)^2(N).$$
(6.40)
Lโintersection des noyaux des formes $`\omega _i^{}`$ dรฉfinit une nouvelle distribution horizontale. On dรฉfinit $`\stackrel{~}{g}`$ comme รฉtant la mรฉtrique sur $`M`$ telle que $`\pi ^{}:(M,\stackrel{~}{g})(N,h)`$ soit une submersion riemannienne et $`\stackrel{~}{g}=g_3`$ sur lโespace vertical. Cette mรฉtrique vรฉrifie le point (iv) de la dรฉfinition 6.1 du fait de (6.40).
On doit encore vรฉrifier $`\stackrel{~}{g}`$ est proche de $`g_3`$. Remarquons que
$$\stackrel{~}{g}g_3=\underset{i}{}(\omega ^2\omega ^2)=\underset{i}{}(2\pi ^{}(\beta _i)\omega _i+\pi ^{}(\beta _i)^2).$$
(6.41)
Or, B. Colbois et G. Courtois ont montrรฉ dans \[CC00\] (lemme A.32) quโil existe une constante $`\kappa (n,a,d,(N,h))>0`$ telle que les formes $`\beta _i`$ telles quโon les a dรฉfinies vรฉrifient $`\beta _i_{\mathrm{}}\kappa `$, ce qui permet de conclure. Remarque: le lemme A.32 de \[CC00\] utilise le fait que pour la mรฉtrique $`g_3`$, la norme de la seconde forme fondamentale est contrรดlรฉ et que la submersion $`\pi ^{}`$ est $`\kappa _i`$-rรฉguliรจre. Il nโest donc pas รฉvident quโon puisse obtenir le thรฉorรจme 6.2 en supposant que la mรฉtrique initiale $`g`$ est invariante et en se passant des rรฉsultats de \[CFG92\].
Enfin, il reste ร montrer que la courbure de $`(M,\stackrel{~}{g})`$ reste bornรฉe dans la direction horizontale. Soit $`xN`$, $`\stackrel{~}{X}`$ et $`\stackrel{~}{Y}`$ deux vecteurs orthonormรฉs tangents ร $`N`$ en $`x`$, $`\overline{\omega }`$ une $`1`$-forme invariante de $`T^k`$, $`\omega `$ la $`1`$-forme induite sur $`M`$ pour la distribution horizontale associรฉe ร $`g`$ et $`\omega ^{}`$ la $`1`$-forme induite pour la distribution associรฉe ร $`\stackrel{~}{g}`$. Ces deux formes vรฉrifient $`\mathrm{d}\omega ^{}=\pi ^{}(\alpha )`$ et $`\mathrm{d}\omega =\pi ^{}(\alpha +\mathrm{d}\beta )`$, oรน $`\alpha `$ est une $`2`$-forme harmonique de $`N`$ et $`\beta `$ une $`1`$-forme de $`N`$.
Dโaprรจs la formule de OโNeill, il suffit pour contrรดler la courbure sectionnelle $`K_{(M,\stackrel{~}{g})}(\pi ^{}(X),\pi ^{}(Y))`$ de majorer la norme de $`[\pi ^{}(X),\pi ^{}(Y)]^V`$. Or, on peut รฉcrire dโune part,
$$\omega ^{}([\pi ^{}(X),\pi ^{}(Y)]^V)=\mathrm{d}\omega ^{}(\pi ^{}(X),\pi ^{}(Y))=\alpha (X,Y).$$
(6.42)
Dโautre part, on a
$$\alpha _{\mathrm{}}\tau ^{}(n,a,d)\alpha _2,$$
(6.43)
dโaprรจs \[Li80\], car $`\alpha `$ est harmonique, et
$$\alpha _2\alpha +\mathrm{d}\beta _2=\mathrm{d}\omega _2\mathrm{d}\omega _{\mathrm{}},$$
(6.44)
en utilisant le fait quโune forme harmonique est le plus petit รฉlรฉment de sa classe de cohomologie pour la norme $`L^2`$. Enfin, le lemme 6.8 permet de contrรดler la norme de $`\mathrm{d}\omega `$ en fonction de $`a`$ et $`\omega _{\mathrm{}}`$, et la norme de $`\omega `$ est contrรดlรฉ en fonction de $`\omega ^{}_{\mathrm{}}`$. Comme la majoration de $`\omega ^{}([\pi ^{}(X),\pi ^{}(Y)]^V)`$ obtenue est indรฉpendante du choix de $`\overline{\omega }`$, on a bien une majoration de $`|[\pi ^{}(X),\pi ^{}(Y)]^V|`$ en fonction de $`n`$, $`a`$ et $`d`$.
## Chapitre 7 Petites valeurs propres des fibrรฉs principaux en tores
### 7.1 Minoration du spectre des $`1`$-formes par le volume du fibrรฉ
Les rรฉsultats des chapitres prรฉcรฉdents nous permettent maintenant de dรฉmontrer le thรฉorรจme 26. On a vu quโon pouvait se ramener au cas dโun fibrรฉ muni dโune mรฉtrique adaptรฉe. On va donc montrer le rรฉsultat du thรฉorรจme 26 pour un fibrรฉ vรฉrifiant les conclusions du thรฉorรจme 6.2:
###### Thรฉorรจme 7.1
Soient $`a>0`$, $`d>0`$ deux rรฉels, $`n`$, $`k`$ et $`m`$ trois entiers tels que $`n=k+m`$, et $`(N^m,h)`$ une variรฉtรฉ riemannienne. Il existe des constantes $`c(n,a,d,(N,h))`$ et $`\epsilon (n,a,d,(N,h))`$ strictement positives telles que si $`\overline{g}`$ est une mรฉtrique sur le tore $`T^k`$ telle que $`diam(T^k)<\epsilon `$ et si $`T^kM^nN`$ est un fibrรฉ principal muni dโun couple de mรฉtriques $`(g,h)`$ adaptรฉ au fibrรฉ et tel que $`g=\overline{g}`$ en restriction ร la fibre, $`diam(M,g)<d`$ et $`|K_M(X,Y)|a`$ pour toute paire $`(X,Y)`$ de vecteurs horizontaux orthonormรฉs, alors on a
$$\lambda _{\mathrm{1,1}}(M,g)cVol^2(T^k).$$
On sโest ici donnรฉ comme hypothรจse que la mรฉtrique sur $`N`$ est fixรฉe. En effet, une hypothรจse sur la courbure ne nous sera pas suffisante. On verra au paragraphe suivant dans quelle mesure on peut espรฉrer obtenir le mรชme rรฉsultat avec des hypothรจses plus faibles.
Dรฉmonstration :
Dans un premier temps, nous allons dรฉmontrer le thรฉorรจme dans le cas oรน le fibrรฉ $`M`$ ne contient pas de sous-fibrรฉ trivial, cโest-ร -dire quand lโapplication $`e:๐ข^{}^2(N,h)`$ est injective. Nous gรฉnรฉraliserons ensuite le rรฉsultat ร un fibrรฉ principal quelconque. Dโautre part, on se restreindra aux formes $`T^k`$-invariantes, en vertu des rรฉsultats du chapitre 5 (corollaire 35). En effet, le spectre des formes orthogonales aux formes invariantes sera minorรฉ en fonction de la constante $`\epsilon `$ du thรฉorรจme, et on pourra toujours choisir cette constante suffisamment petite de sorte que le spectre des formes orthogonales aux formes invariantes soit plus grand que le terme $`cVol^2(T^k)`$.
Supposons donc $`e`$ injective. La dรฉmonstration se dรฉroule en deux รฉtapes. Dโabord, on se ramรจne ร lโรฉtude des valeurs propres de lโopรฉrateur $`e^{}e`$, lโadjoint รฉtant dรฉfini en munissant $`^2(N)`$ de sa norme $`L^2`$ :
###### Fait 7.2
Il existe $`\epsilon (n,a,\lambda _{\mathrm{0,1}}(N,h),\lambda _{\mathrm{1,1}}(N,h))>0`$ et $`c(n,a,\lambda _{\mathrm{0,1}}(N,h))>0`$ tel que pour toute $`1`$-forme $`\phi `$ sur $`M`$ $`T^k`$-invariante et orthogonale ร $`Ker\mathrm{\Delta }^1(M,g)`$, si le quotient de Rayleigh de $`\phi `$ vรฉrifie $`R(\phi )<\epsilon `$, alors il existe une forme $`\omega `$ induite par un รฉlรฉment de $`๐ข^{}`$ telle que $`e(\omega )^2c\epsilon \omega ^2`$.
Dรฉmonstration : Soit $`\phi `$ une $`1`$-forme diffรฉrentielle $`T^k`$-invariante de $`M`$. On peut รฉcrire
$$\phi =\pi ^{}(\alpha )+\underset{i=1}{\overset{k}{}}\pi ^{}(a_i)\omega _i,$$
(7.3)
$`\alpha `$ est une $`1`$-forme de $`N`$, $`a_i`$ des fonctions de $`N`$ et $`\omega _i`$ les $`1`$-formes verticales induites par une base orthonormรฉe de $`๐ข^{}`$. On a alors :
$$\mathrm{d}\phi =\pi ^{}(\mathrm{d}\alpha +\underset{i=1}{\overset{k}{}}a_ie_i)+\underset{i=1}{\overset{k}{}}\mathrm{d}\pi ^{}(a_i)\omega _i,$$
(7.4)
$`e_i`$ dรฉsigne lโimage de $`\omega _i`$ par lโapplication $`e:๐ข^{}^2(N)`$. De plus, pour tout $`i`$ on a
$$\delta (\pi ^{}(a_i)\omega _i)^2=(\delta (\pi ^{}(a_i)\omega _i),\delta (\pi ^{}(a_i)\omega _i))=(\pi ^{}(a_i)\omega _i,\mathrm{d}\delta (\pi ^{}(a_i)\omega _i)),$$
$`(,)`$ dรฉsigne le produit scalaire $`L^2`$. Comme $`\pi ^{}(a_i)\omega _i`$ est une forme $`T^k`$-invariante, $`\delta (\pi ^{}(a_i)\omega _i)`$ est une fonction invariante, cโest-ร -dire que cโest le relevรฉ dโune fonction sur $`N`$. Par consรฉquent, $`\mathrm{d}\delta (\pi ^{}(a_i)\omega _i)`$ est le relevรฉ dโune $`1`$-forme sur $`N`$, et est donc orthogonale ร $`\pi ^{}(a_i)\omega _i`$. Finalement, on a:
$$\delta \phi =\pi ^{}(\delta \alpha ).$$
(7.5)
On peut calculer prรฉcisรฉment quelles sont les $`1`$-formes harmoniques de $`M`$. On sait dรฉjร que les formes harmoniques sont invariantes, donc de la forme donnรฉe en (7.3). Si $`\phi `$ est harmonique, on a de plus $`\mathrm{d}\phi =0`$ et $`\delta \phi =0`$, donc
$$\delta \alpha =0,\mathrm{d}a_i=0\text{ pour tout }i\text{, et }\mathrm{d}\alpha +\underset{i=1}{\overset{k}{}}a_ie_i=0.$$
(7.6)
Comme les fonctions $`a_i`$ sont constantes, $`_{i=1}^ka_ie_i`$ est une $`2`$-forme harmonique de $`N`$, donc orthogonale ร la forme exacte $`\mathrm{d}\alpha `$. On a donc $`\mathrm{\Delta }\alpha =0`$ et $`_{i=1}^ka_ie_i=0`$. Comme $`e`$ est injective, les $`e_i`$ forment une famille libre, et donc $`a_i=0`$ pour tout $`i`$. On obtient finalement que les formes harmoniques de $`M`$ sont les relevรฉs des formes harmoniques de $`N`$.
Supposons que $`\phi `$ est de norme $`1`$, cโest-ร -dire que $`\alpha ^2+_{i=1}^ka_i^2=1`$. Le quotient de Rayleigh de $`\phi `$ sโรฉcrit alors
$$R(\phi )=\delta \alpha ^2+\mathrm{d}\alpha +\underset{i=1}{\overset{k}{}}a_ie_i^2+\underset{i=1}{\overset{k}{}}\mathrm{d}\pi ^{}(a_i)^2.$$
(7.7)
Supposons que $`R(\phi )<\epsilon `$ pour un $`\epsilon >0`$ donnรฉ. Comme $`\alpha ^2+_{i=1}^ka_i^2`$ est รฉgal ร $`1`$, lโun des termes de la somme est plus grande que $`\frac{1}{1+k}`$. Nous allons distinguer les cas $`\alpha ^2>\frac{1}{1+k}`$, et $`a_p^2>\frac{1}{1+k}`$ pour un $`p[1,k]`$.
Supposons $`\alpha ^2\frac{1}{1+k}`$ :
La forme $`\phi `$ est orthogonale ร $`Ker\mathrm{\Delta }^1`$, cโest-ร -dire ร lโensemble des relevรฉs de formes harmoniques de $`N`$. Par consรฉquent, la forme $`\alpha `$ est elle-mรชme orthogonales aux formes harmoniques de $`N`$, et donc $`\mathrm{d}\alpha ^2+\delta \alpha ^2\lambda _{\mathrm{1,1}}(N,h)`$.
Si $`\delta \alpha ^2\mathrm{d}\alpha ^2`$, alors $`\frac{\delta \alpha ^2}{\alpha ^2}\frac{\lambda _{\mathrm{1,1}}(N,h)}{2}`$. Or, on a les deux inรฉgalitรฉs
$$\alpha ^2\frac{1}{k+1}$$
(7.8)
et
$$\delta \alpha ^2R(\phi )<\epsilon ,$$
(7.9)
dont on peut dรฉduire la majoration
$$\frac{\lambda _{\mathrm{1,1}}(N,h)}{2(k+1)}<\epsilon .$$
(7.10)
On peut choisir $`\epsilon `$ suffisamment petit pour รฉliminer ce cas.
Si $`\mathrm{d}\alpha ^2\delta \alpha ^2`$, alors $`\frac{\mathrm{d}\alpha ^2}{\alpha ^2}\frac{\lambda _{\mathrm{1,1}}(N,h)}{2}`$. On peut dรฉduire de (7.7) que
$$\mathrm{d}\alpha +\underset{i=1}{\overset{k}{}}a_ie_i^2<\epsilon .$$
(7.11)
On veut montrer ร partir de (7.11) que puisque $`\mathrm{d}\alpha `$ est minorรฉ, les $`a_i`$, et donc les $`\mathrm{d}a_i`$ le sont aussi, ce qui contredira le fait que $`R(\phi )<\epsilon `$. Il dรฉcoule de (7.11):
$$\mathrm{d}\alpha ^2+\underset{i=1}{\overset{k}{}}a_ie_i^2+2(\mathrm{d}\alpha ,\underset{i=1}{\overset{k}{}}a_ie_i)<\epsilon $$
(7.12)
et donc
$$\frac{\lambda _{\mathrm{1,1}}(N,h)}{2(k+1)}\epsilon <2(\mathrm{d}\alpha ,\underset{i=1}{\overset{k}{}}a_ie_i),$$
(7.13)
De plus, en utilisant le fait que $`e_i`$ est harmonique, donc cofermรฉe et quโen restriction aux $`1`$-formes on a $`\delta =\mathrm{d}`$, on obtient $`(\mathrm{d}\alpha ,a_ie_i)=(\alpha ,\delta (a_ie_i))=(\alpha ,(\mathrm{d}a_i(e_i)))`$. En appliquant ponctuellement lโinรฉgalitรฉ de Schwarz, on en dรฉduit
$`(\mathrm{d}\alpha ,a_ie_i)`$ $``$ $`{\displaystyle _N}|\alpha ||\mathrm{d}a_i(e_i)|\mathrm{d}v{\displaystyle _N}|\alpha ||\mathrm{d}a_i||e_i|\mathrm{d}v`$ (7.14)
$``$ $`e_i_{\mathrm{}}{\displaystyle _N}|\alpha ||\mathrm{d}a_i|dv=e_i_{\mathrm{}}(|\alpha ||\mathrm{d}a_i|)`$
$``$ $`e_i_{\mathrm{}}\alpha \mathrm{d}a_ie_i_{\mathrm{}}\mathrm{d}a_i`$
et finalement
$$\frac{\lambda _{\mathrm{1,1}}(N,h)}{4}\epsilon <2\underset{i=1}{\overset{k}{}}e_i_{\mathrm{}}\mathrm{d}a_i.$$
(7.15)
Comme selon le lemme 6.8, $`e_i_{\mathrm{}}`$ est majorรฉ en fonction de $`a`$ et $`n`$, on obtient si $`\epsilon `$ est suffisamment petit une minoration de $`_{i=1}^k\mathrm{d}a_i`$, ce qui contredit le fait que $`_{i=1}^k\mathrm{d}a_i^2<\epsilon `$.
En choisissant $`\epsilon `$ suffisamment petit en fonction de $`\lambda _{\mathrm{1,1}}(N,h)`$, $`a`$ et $`n`$, on peut donc รฉcarter le cas $`\alpha ^2\frac{1}{1+k}`$.
Supposons $`a_p^2\frac{1}{1+k}`$ :
Si on note $`\overline{a}_i`$ la valeur moyenne de la fonction $`a_i`$, on peut choisir la base $`(\omega _i)`$ de sorte que $`\overline{a}_i=0`$ pour $`i2`$ (il suffit de choisir comme nouveau $`\omega _1`$ la forme diffรฉrentielle $`_{i=1}^k\overline{a}_i\omega _i`$). On a alors
$$\mathrm{d}a_i^2\lambda _{\mathrm{0,1}}(N,h)a_i^2$$
(7.16)
pour $`i2`$ et
$$\mathrm{d}a_1^2\lambda _{\mathrm{0,1}}(N,h)a_1\overline{a}_1^2.$$
(7.17)
Comme $`\mathrm{d}a_i^2<\epsilon `$ pour tout $`i`$, on a donc
$$a_i_2^2<\frac{\epsilon }{\lambda _{\mathrm{0,1}}(N,h)}$$
(7.18)
pour $`i2`$ et
$$a_1\overline{a}_1_2^2<\frac{\epsilon }{\lambda _{\mathrm{0,1}}(N,h)}.$$
(7.19)
En particulier, si $`\epsilon `$ est suffisamment petit, on a $`p=1`$.
On peut alors รฉcrire
$`\mathrm{d}\alpha +{\displaystyle \underset{i=1}{\overset{k}{}}}a_ie_i^2`$ $`=`$ $`\overline{a}_1e_1^2+\mathrm{d}\alpha +(a_1\overline{a}_1)e_1+{\displaystyle \underset{i=2}{\overset{k}{}}}a_ie_i^2`$ (7.20)
$`+2(\overline{a}_1e_1,\mathrm{d}\alpha +(a_1\overline{a}_1)e_1+{\displaystyle \underset{i=2}{\overset{k}{}}}a_ie_i).`$
On a dโune part
$$(\overline{a}_1e_1,\mathrm{d}\alpha )=\overline{a}_1(\delta e_1,\alpha )=0,$$
et dโautre part
$`(\overline{a}_1e_1,a_ie_i)`$ $``$ $`\overline{a}_1e_1_2a_ie_i_2\overline{a}_1e_1_2e_i_{\mathrm{}}a_i_2`$ (7.21)
$`<`$ $`\overline{a}_1e_1_2{\displaystyle \frac{\sqrt{\epsilon }e_i_{\mathrm{}}}{\lambda _{\mathrm{0,1}}(N,h)}},`$
cette derniรจre inรฉgalitรฉ restant vraie pour $`i=1`$ en remplaรงant $`a_i`$ par $`(\overline{a}_1a_1)`$.
Comme $`\mathrm{d}\alpha +_{i=1}^ka_ie_i^2<\epsilon `$, on a finalement
$$\overline{a}_1e_1^2<2|(\overline{a}_1e_1,\mathrm{d}\alpha +(a_1\overline{a}_1)e_1+\underset{i=2}{\overset{k}{}}a_ie_i)|+\epsilon $$
et donc
$$\overline{a}_1e_1^2\frac{2n\sqrt{\epsilon }e}{\lambda _{\mathrm{0,1}}(N,h)}\overline{a}_1e_1\epsilon <0.$$
(7.22)
On en dรฉduit que $`\overline{a}_1e_1`$ est encadrรฉ par les racines du polynรดme $`X^2\frac{2n\sqrt{\epsilon }e}{\lambda _{\mathrm{0,1}}(N,h)}X\epsilon `$, et en particulier majorรฉ par sa plus grande racine. Comme $`e`$ est majorรฉ en fonction de $`n`$ et de la borne sur la courbure (lemme 6.8), la plus grande racine du polynรดme sโรฉcrit comme une constante dรฉpendant de $`n`$, $`a`$ et $`\lambda _{\mathrm{0,1}}(N,h)`$, multipliรฉe par $`\sqrt{\epsilon }`$. On obtient donc le rรฉsultat souhaitรฉ en prenant comme $`\omega `$ la forme $`\overline{a}_1\omega _1`$.
On va maintenant minorer le spectre de $`e^{}e`$ en fonction du volume de la fibre $`T^k`$.
###### Fait 7.23
Il existe une constante $`c(n,a,(N,h))>0`$ telle que la premiรจre valeur propre de $`e^{}e`$ soit minorรฉe par $`cVol(T^k)^2`$.
Dรฉmonstration : Soient $`\lambda _1,\mathrm{},\lambda _k`$ les valeurs propres de $`e^{}e`$ classรฉes dans lโordre croissant. Comme $`e`$ est injective, ces valeurs propres sont non nulles et la premiรจre vรฉrifie
$$\lambda _1=\frac{_i\lambda _i}{_{i1}\lambda _i}=\frac{Det(e^{}e)}{_{i1}\lambda _i}.$$
(7.24)
Par ailleurs, les valeurs propres de $`e^{}e`$ vรฉrifient $`\lambda _ie^{}e`$, donc
$$\lambda _1\frac{Det(e^{}e)}{e^{}e^{k1}}\frac{Det(e^{}e)}{e^{2k2}}.$$
(7.25)
Lโimage de $`e`$ est un sous-espace de $`Ker\mathrm{\Delta }^2(N)`$ de dimension $`k`$, engendrรฉ par un sous-rรฉseau du rรฉseau des formes harmoniques entiรจres de $`N`$. Si on restreint $`e`$ et $`e^{}`$ ร ce sous-espace, on peut รฉcrire $`Det(e^{}e)=(Dete)^2`$, oรน $`Dete`$ est le dรฉterminant dโune matrice de $`e`$ รฉcrite dans des bases orthonormรฉes de $`๐ข^{}`$ et $`Ime`$, ce qui donne
$$\lambda _1\frac{(Dete)^2}{e^{2k2}}.$$
(7.26)
Le lemme 6.8 donne une majoration de $`e`$ en fonction de $`n`$ et de la borne $`a`$ sur la courbure de $`M`$, il ne reste donc quโร minorer $`Dete`$. Notons $`Det^{}e`$ le dรฉterminant de la matrice de $`e`$ dans la base canonique de $`๐ข^{}=_{}^{k}{}_{}{}^{}`$ et une base orthonormรฉe de $`Ime`$. On a alors $`Dete=(Det^{}e)(VolT^k)`$. Comme les images dans $`Ker\mathrm{\Delta }^2(N)`$ des รฉlรฉments de la base canonique de $`๐ข^{}`$ sont des formes entiรจres, le dรฉterminant $`Det^{}e`$, qui est aussi le volume de $`e([\mathrm{0,1}]^k)`$, est un multiple du volume dโun domaine fondamental du rรฉseau des formes entiรจres dans $`Ime`$. Comme par ailleurs $`Det^{}e`$ est non nul, il sera donc minorรฉ par le volume de ce domaine fondamental. Si on note $`\rho `$ le minimum des normes des $`2`$-formes harmoniques entiรจres non nulles, ce volume est minorรฉ par le volume dโune boule de rayon $`\frac{\rho }{2}`$ dans $`Ime`$, et donc minorรฉ par une constante ne dรฉpendant que de $`n`$ et de la mรฉtrique $`h`$ de $`N`$. On peut donc bien รฉcrire
$$\lambda _1c(n,a,(N,h))Vol(T^k)^2.$$
Nous allons maintenant supposer que $`e`$ nโest pas injective. Notons $`l`$ la dimension de son noyau. Le premier nombre de Betti de $`M`$ est alors $`b_1(N)+l`$. En effet, on a vu que si une $`1`$-forme $`\phi =\pi ^{}(\alpha )+_{i=1}^k\pi ^{}(a_i)\omega _i`$ est harmonique, cela signifie, dโaprรจs (7.4) et (7.5) :
$$\mathrm{\Delta }\alpha =0,\mathrm{d}a_i=0\text{ pour tout }i\text{, et }\underset{i=1}{\overset{k}{}}a_ie_i=0.$$
(7.27)
Comme les fonctions $`a_i`$ sont constantes. Lโensemble des $`a_i`$ tels que $`_{i=1}^ka_ie_i=0`$ est exactement le noyau de $`e`$. Lโespace des formes harmoniques de $`M`$ est donc lโespace engendrรฉ par les relevรฉs des formes harmoniques de $`N`$ et les formes verticales induites par les รฉlรฉments de noyau de $`e`$.
On peut reprendre la dรฉmonstration prรฉcรฉdente en prenant pour $`(\omega _i)_i`$ une base de $`๐ข^{}`$ telle que $`\omega _{kl+1},\mathrm{},\omega _k`$ soit une base de $`Kere`$ (le fait que la forme $`\phi `$ est orthogonale aux formes harmoniques se traduit par le fait que $`a_{kl+1},\mathrm{},a_k=0`$) et en รฉtudiant $`e^{}e`$ restreint ร lโorthogonal de $`Kere`$. On obtient de la mรชme faรงon le rรฉsultat du fait 7.2, ร savoir que la premiรจre valeur propre du laplacien sur $`M`$ est minorรฉe ร une constante multiplicative prรจs par la premiรจre valeur propre de $`(e^{}e)_{|(Kere)^{}}`$.
Pour minorer le spectre de $`(e^{}e)_{|(Kere)^{}}`$, on doit รชtre un peu plus attentif dans la manipulation des bases de $`๐ข^{}`$.
Soit $`=(\omega _1,\mathrm{},\omega _k)`$ une base orthonormรฉe de $`๐ข^{}`$ et $`^{}=(\omega _1^{},\mathrm{},\omega _k^{})`$ une base du rรฉseau des entiers de $`๐ข^{}`$, telles que $`(\omega _1,\mathrm{},\omega _l)`$ et $`(\omega _1^{},\mathrm{},\omega _l^{})`$ soient des bases de $`Kere`$ (comme lโimage du rรฉseau des entiers de $`๐ข^{}`$ est contenue dans un rรฉseau, le noyau de $`e`$ est effectivement engendrรฉ par des รฉlรฉments entiers). La matrice de passage de $``$ ร $`^{}`$ est de la forme
$$P=\left(\begin{array}{cc}P_1& P_2\\ 0& P_3\end{array}\right),$$
$`P_1`$ est un bloc carrรฉ de taille $`l`$. Si se donne une base orthonormรฉe de $`Ime`$, la matrice de $`e`$ sโรฉcrit sous la forme $`(0,A)`$ dans la base $``$ et $`(0,A^{})`$ dans la base $`^{}`$, oรน $`A`$ et $`A^{}`$ sont des blocs carrรฉs de taille $`kl`$ et vรฉrifient $`A^{}=AP_3`$.
Le spectre de $`(e^{}e)_{|(Kere)^{}}`$ est celui de $`A^{}A`$. On peut รฉcrire, comme dans la dรฉmonstration du fait 7.23 :
$$\lambda _1\frac{DetA^{}A}{A^{}A^{kl1}}\frac{(DetA)^2}{A^{2(kl1)}},$$
(7.28)
$`\lambda _1`$ est la premiรจre valeur propre non nulle de $`e^{}e`$. De plus, on a $`DetA^{}=DetADetP_2`$ et donc
$$DetA=\frac{DetA^{}}{DetP_3}=DetA^{}\frac{DetP_1}{DetP}.$$
(7.29)
Le dรฉterminant de $`A^{}`$ est, comme prรฉcรฉdemment, minorรฉ par le covolume du rรฉseau des formes entiรจres dans $`Ime`$, et $`DetP`$ sโinterprรจte gรฉomรฉtriquement comme lโinverse du volume de $`T^k`$.
Il reste ร minorer $`DetP_1`$. Comme $`Kere`$ est engendrรฉ par des รฉlรฉments entiers de $`๐ข^{}`$, lโorthogonal de $`Kere`$ pour la dualitรฉ dรฉfinit un sous-tore $`T^{kl}`$ de $`T^k`$. De plus, le dual de de lโalgรจbre de Lie $`๐ข(T^k/T^{kl})`$ du quotient $`T^k/T^{kl}`$ est isomorphe ร $`Kere`$. La matrice $`P_1`$ est donc la matrice de passage dโune base orthonormรฉe de $`๐ข^{}(T^k/T^{kl})`$ dans une base du base orthonormรฉe de $`๐ข^{}(T^k/T^{kl})`$ dans une base du rรฉseau des entiers de $`๐ข^{}(T^k/T^{kl})`$, et par consรฉquent $`DetP_1`$ est lโinverse du volume de $`T^k/T^{kl}`$ pour la mรฉtrique quotient. Le diamรจtre de $`T^k/T^{kl}`$ est majorรฉ par $`\epsilon `$, comme celui de $`T^k`$, et par consรฉquent son volume aussi.
### 7.2 Petites valeurs propres et norme des $`2`$-formes harmoniques entiรจres
Nous allons ici discuter de la valeur de la constante $`c(n,a,d,(N,h))`$ du thรฉorรจme 7.1, et en particulier de la maniรจre dont elle dรฉpend de la mรฉtrique $`h`$ sur $`N`$.
Dans la dรฉmonstration du thรฉorรจme, la gรฉomรฉtrie de $`N`$ intervient quatre fois: on a besoin de contrรดler sa courbure pour appliquer la formule de OโNeill et majorer les $`e_i`$; en 7.2 apparaissent les valeurs propre $`\lambda _{\mathrm{1,1}}(N,h)`$ et $`\lambda _{\mathrm{0,1}}(N,h)`$; enfin on fait intervenir en 7.23 le minimum des normes des 2-formes harmoniques non nulles dont la classe de cohomologie est entiรจre.
On peut noter que dans la dรฉmonstration du fait 7.2, on a seulement besoin dโune minoration des deux valeurs propres $`\lambda _{\mathrm{1,1}}(N,h)`$ et $`\lambda _{\mathrm{0,1}}(N,h)`$. Lโidรฉe est en fait de sโassurer que le spectre de $`(N,h)`$ nโinterfรจre pas dans la recherche des petites valeurs propres de $`M`$. On sait par ailleurs que $`\lambda _{\mathrm{0,1}}(N,h)`$ peut รชtre minorรฉ en fonction du diamรจtre et de la courbure de $`N`$ (cf. thรฉorรจme 4), et $`\lambda _{\mathrm{1,1}}(N,h)`$ en fonction du diamรจtre, de la courbure et du rayon dโinjectivitรฉ de $`N`$ (thรฉorรจme 5). En outre, une borne sur la courbure est aussi suffisante pour appliquer la formule de OโNeill.
La dรฉmonstration du fait 7.23 introduit quand ร elle la constante
$$\rho (N,h)=\underset{\stackrel{\alpha ^2(N,h)\backslash \left\{0\right\}}{[\alpha ]H^2(N,)}}{inf}\alpha _2$$
(7.30)
dans la minoration de la premiรจre valeur propre du laplacien. Il est naturel de se demander si lโon peut contrรดler $`\rho (N,h)`$ ร lโaide des mรชmes invariants gรฉomรฉtriques que $`\lambda _{\mathrm{0,1}}(N,h)`$ et $`\lambda _{\mathrm{1,1}}(N,h)`$:
###### Question 7.31
Existe-t-il une constante $`c(n,a,d,r)>0`$ telle que si $`(N^n,h)`$ est une variรฉtรฉ riemannienne vรฉrifiant $`diam(N,h)d`$ et $`|K(N,h)|a`$ et $`inj(N,h)r`$, alors $`\rho (N,h)c`$ ?
Un argument de compacitรฉ permet de montrer quโavec lโhypothรจse de rayon dโinjectivitรฉ minorรฉe, on peut rรฉpondre affirmativement ร la question 7.31.
###### Proposition 7.32
Pour tout rรฉels $`a,d,r>0`$ et tout entier $`n^{}`$, il existe une constante $`c(n,a,d,r)>0`$ (non explicite) telle que si $`(N^n,h)`$ est une variรฉtรฉ riemannienne vรฉrifiant $`diam(N,h)d`$ et $`|K(N,h)|a`$ et $`inj(N,h)r`$, alors $`\rho (N,h)c`$.
Dรฉmonstration : On considรจre le tore $`T=T^{b_2(N)}`$, vu comme quotient de $`^2(N,h)`$ par le rรฉseau des formes harmoniques entiรจres. La mรฉtrique $`h`$ sur $`N`$ induit une norme euclidienne sur $`^2(N,h)`$ qui passe au quotient sur $`T`$ en une mรฉtrique plate $`\overline{h}`$. Minorer $`\rho (N,h)`$ revient ร minorer le rayon dโinjectivitรฉ de $`(T,\overline{h})`$.
On sait (\[AC92\]) que lโespace des mรฉtriques $`h`$ sur $`N`$ telles que $`diam(N,h)d`$ et $`|K(N,h)|a`$ et $`inj(N,h)r`$ est relativement compact pour la topologie $`C^\alpha `$. De plus, la mรฉtrique $`\overline{h}`$ dรฉpend continument de $`h`$. En effet, si on se donne un rรฉel $`\epsilon >0`$ et une mรฉtrique $`h`$ sur $`N`$, on aura, pour toute mรฉtrique $`h^{}`$ suffisamment proche de $`h`$ et toute $`2`$-forme $`\alpha `$ harmonique pour la mรฉtrique $`h`$, $`\left|\alpha _2\alpha _2^{}\right|\epsilon \alpha _{\mathrm{}}`$, oรน $`_p`$ et $`_p^{}`$ dรฉsignent les normes pour les mรฉtriques $`h`$ et $`h^{}`$ respectivement. Comme ร courbure et diamรจtre bornรฉs, la norme $`L^{\mathrm{}}`$ des formes harmoniques est contrรดlรฉe par leur norme $`L^2`$ (cf. \[Li80\] et inรฉgalitรฉ (6.43)), on peut รฉcrire $`\left|\alpha _2\alpha _2^{}\right|\tau (n,a,d)\epsilon \alpha _2`$. Si $`\alpha ^{}`$ est le reprรฉsentant harmonique pour $`h^{}`$ de la classe de cohomologie de $`\alpha `$, on peut donc finalement trouver un voisinage $`๐ฑ`$ de $`h`$ tel que pour toute mรฉtrique $`h^{}`$ dans $`๐ฑ`$,
$$\alpha ^{}_2^{}\alpha _2^{}(1+\epsilon )\alpha _2.$$
(7.33)
Quitte ร restreindre le voisinage $`๐ฑ`$, on a rรฉciproquement
$$\alpha _2(1+\epsilon )\alpha ^{}_2^{},$$
(7.34)
ce qui implique bien la continuitรฉ de $`h\overline{h}`$. Lโensemble dรฉcrit par $`\overline{h}`$ quand $`h`$ varie est donc relativement compact dans lโespace des mรฉtriques plates du tore. Il existe par consรฉquent une mรฉtrique sur $`T`$ qui rรฉalise la borne infรฉrieure du rayon dโinjectivitรฉ de $`T`$ quand $`h`$ varie. En particulier, cette borne infรฉrieure est non nulle
On peut se demander si le rรฉsultat reste vrai avec des hypothรจses plus faibles :
###### Question 7.35
Existe-t-il une constante $`c(n,a,d)>0`$ telle que si $`(N^n,h)`$ est une variรฉtรฉ riemannienne vรฉrifiant $`diam(N,h)d`$ et $`|K(N,h)|a`$ et alors $`\rho (N,h)c`$ ?
Il faut noter par ailleurs quโune minoration non explicite ne permet pas dโamรฉliorer les minorations dรฉjร connues de la premiรจre valeur propre du spectre. On a besoin dโestimations prรฉcises :
###### Question 7.36
Si $`diam(N,h)d`$, $`|K(N,h)|a`$ et $`inj(N,h)r`$, peut-on minorer $`\rho (N,h)`$ par une constante explicite $`c(n,a,d,r)>0`$ ? Plus prรฉcisรฉment, peut-on trouver une constante $`c`$ de la forme $`c^{}(n,a,d)Vol(N,h)^{\alpha (n)}`$ ou $`c^{}(n,a,d)inj(N,h)^{\alpha (n)}`$ ?
On peut noter quโexpliciter le rรดle du volume de $`N`$ dans cette minoration permet dโobtenir dans le thรฉorรจme 7.1 une minoration de $`\lambda _{\mathrm{1,1}}(M,g)`$ en fonction du volume de $`(M,g)`$.
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# Quantization of complex Lagrangian submanifoldsTo appear in: Adv. Math.2000 AMS Mathematics Subject Classification(s): 46L65, 14A20, 32C38
## 1 Introduction
Let $`๐`$ be a complex contact manifold. A local model for $`๐`$ is an open subset of the projective cotangent bundle $`P^{}Y`$ to a complex manifold $`Y`$. The manifold $`P^{}Y`$ is endowed with the sheaf $`_Y`$ of microdifferential operators of . In , Kashiwara proves the existence of a canonical stack $`\mathrm{๐ฌ๐๐ฝ}(๐ค_๐)`$ on $`๐`$, locally equivalent to the stack of $`_Y`$-modules. Let $`\mathrm{\Lambda }๐`$ be a smooth Lagrangian submanifold. In the same paper, Kashiwara states that there exists a globally defined holonomic system simple along $`\mathrm{\Lambda }`$ in the stack $`\mathrm{๐ฌ๐๐ฝ}(๐ค_๐|_\mathrm{\Lambda })`$ twisted by half-forms on $`\mathrm{\Lambda }`$.
Now, let $`๐`$ be a complex symplectic manifold. A local model for $`๐`$ is an open subset of the cotangent bundle $`T^{}X`$ to a complex manifold $`X`$. The manifold $`T^{}X`$ is endowed with the sheaf $`๐ฒ_X`$ of WKB-differential operators, similar to $`_X`$, but with an extra central parameter, a substitute to the lack of homogeneity. (Note that, in the literature, $`๐ฒ_X`$ is also called a deformation-quantization ring, or a ring of semi-classical differential operators. See for a precise description of the ring $`๐ฒ_X`$ and its links with $`_X`$.) A stack $`\mathrm{๐ฌ๐๐ฝ}(๐ถ_๐)`$ on $`๐`$ locally equivalent to the stack of $`๐ฒ_X`$-modules has been constructed in the formal case (in the general setting of Poisson manifolds) by and in the analytic case (and by a different method, similar to ) by . (See also for papers closely related to this subject.)
In this paper, we prove that if $`\mathrm{\Lambda }`$ is a smooth Lagrangian submanifold of the complex symplectic manifold $`๐`$, there exists a globally defined simple holonomic module along $`\mathrm{\Lambda }`$ in the stack $`\mathrm{๐ฌ๐๐ฝ}(๐ถ_๐|_\mathrm{\Lambda })`$ twisted by half-forms on $`\mathrm{\Lambda }`$. As a by-product, we prove that there is an equivalence of stacks between that of twisted regular holonomic modules along $`\mathrm{\Lambda }`$ and that of local systems on $`\mathrm{\Lambda }`$. The local model for our theorem is given by $`๐=T^{}X`$ and $`\mathrm{\Lambda }=T_X^{}X`$, the zero-section of $`T^{}X`$. In this case, a simple module is the sheaf $`๐ช_X^\tau `$ whose sections are series $`_{\mathrm{}<jm}f_j\tau ^j`$ ($`m`$), where the $`f_j`$โs are sections of $`๐ช_X`$ and the family $`\{f_j\}_j`$ satisfies certain growth conditions on compact subsets of $`X`$. The problem we solve here is how to patch together these local models.
Our proof consists in showing that if $`๐`$ is a complex symplectic manifold and $`\mathrm{\Lambda }`$ a Lagrangian submanifold, then there exists a โcontactificationโ $`๐`$ of $`๐`$ in a neighborhood of $`\mathrm{\Lambda }`$. Local models for $`๐`$ and $`๐`$ are an open subset of the cotangent bundle $`T^{}X`$ to a complex manifold $`X`$ and an open subset of the projective cotangent bundle $`P^{}(X\times )`$, respectively. With the same techniques as in , we construct a stack $`\mathrm{๐ฌ๐๐ฝ}(๐ค_{๐,\widehat{t}})`$ on $`๐`$ locally equivalent to the stack of modules over the ring $`_{X\times ,\widehat{t}}`$ of microdifferential operators commuting with $`/_t`$, where $`t`$ is the coordinate on $``$. We then apply Kashiwaraโs existence theorem for simple modules along Lagrangian manifolds in the contact case to deduce the corresponding result in the symplectic case. In fact, the technical heart of this paper is devoted to giving a detailed proof, based on the theory of symbols of simple sections of holonomic modules, of Kashiwaraโs result stated in .
Acknowledgements We would like to thank Louis Boutet de Monvel, Masaki Kashiwara, and Pietro Polesello for their useful comments and insights.
## 2 Stacks
Stacks were invented by Grothendieck and Giraud and we refer to for an exposition. Roughly speaking, a prestack (resp. a stack) is a presheaf (resp. a sheaf) of categories, as we shall see below.
In sections 2 and 3, we denote by $`X`$ a topological space. However, all definitions and results easily extend when replacing $`X`$ with a site, that is, a small category $`๐_X`$ endowed with a Grothendieck topology.
###### Definition 2.1.
1. A prestack $`๐ฒ`$ on $`X`$ is the assignment of a category $`๐ฒ(U)`$ for every open subset $`UX`$, a functor $`\rho _{VU}:๐ฒ(U)\stackrel{}{}๐ฒ(V)`$ for every open inclusion $`VU`$, and an isomorphism of functors $`\lambda _{WVU}:\rho _{WV}\rho _{VU}\rho _{WU}`$ for every open inclusion $`WVU`$, such that $`\rho _{UU}=\mathrm{id}_{๐ฒ(U)}`$, $`\lambda _{UUU}=\mathrm{id}_{\mathrm{id}_{๐ฒ(U)}}`$, and the following diagram of isomorphisms of functors from $`๐ฒ(U)`$ to $`๐ฒ(Y)`$ commutes for every open inclusion $`YWVU`$
For $`F๐ฒ(U)`$ and $`VU`$, we will write $`F|_V`$ for short instead of $`\rho _{VU}(F)`$.
2. A separated prestack is a prestack $`๐ฒ`$ such that for any $`F,G๐ฒ(U)`$, the presheaf $`om_{๐ฒ|_U}(F,G)`$, defined by $`V\mathrm{Hom}_{๐ฒ(V)}(F|_V,G|_V)`$, is a sheaf.
A stack is a separated prestack satisfying suitable glueing conditions, which may be expressed in terms of descent data.
###### Definition 2.2.
Let $`U`$ be an open subset of $`X`$, $`๐ฐ=\{U_i\}_{iI}`$ an open covering of $`U`$ and $`๐ฒ`$ a separated prestack on $`X`$.
1. A descent datum on $`๐ฐ`$ for $`๐ฒ`$ is a pair
$$(\{_i\}_{iI},\{\theta _{ij}\}_{i,jI}),\text{ with }_i๐ฒ(U_i),\theta _{ij}:_j|_{U_{ij}}\stackrel{}{}_i|_{U_{ij}}$$
(2.1)
such that the following diagram of isomorphisms in $`๐ฒ(U_{ijk})`$ commutes
2. The descent datum (2.1) is called effective if there exist $`๐ฒ(U)`$ and isomorphisms $`\theta _i:|_{U_i}\stackrel{}{}_i`$ in $`๐ฒ(U_i)`$ satisfying the natural compatibility conditions with the $`\theta _{ij}`$โs and $`\lambda `$โs.
Note that if the descent datum (2.1) is effective, then $``$ is unique up to unique isomorphism.
###### Definition 2.3.
A stack is a separated prestack such that for any open subset $`U`$ of $`X`$ and any open covering $`๐ฐ=\{U_i\}_{iI}`$ of $`U`$, the descent datum is effective.
To end this section, let us go up one level, and recall the glueing conditions for stacks.
###### Definition 2.4.
Let $`U`$ be an open subset of $`X`$, $`๐ฐ=\{U_i\}_{iI}`$ an open covering of $`U`$.
1. A descent datum for stacks on $`๐ฐ`$ is a triplet
$$(\{๐ฒ_i\}_{iI},\{\phi _{ij}\}_{i,jI},\{\alpha _{ijk}\}_{i,j,kI}),$$
(2.2)
where the $`๐ฒ_i`$โs are stacks on $`U_i`$, $`\phi _{ij}:๐ฒ_j|_{U_{ij}}\stackrel{}{}๐ฒ_i|_{U_{ij}}`$ are equivalences of stacks, and $`\alpha _{ijk}:\phi _{ij}\phi _{jk}\phi _{ik}:๐ฒ_k|_{U_{ijk}}\stackrel{}{}๐ฒ_i|_{U_{ijk}}`$ are isomorphisms of functors such that for any $`i,j,k,lI`$ the following diagram of isomorphisms of functors from $`๐ฒ_l|_{U_{ijkl}}`$ to $`๐ฒ_i|_{U_{ijkl}}`$ commutes
(2.3)
2. The descent datum (2.2) is called effective if there exist a stack $`๐ฒ`$ on $`U`$, equivalences of stacks $`\phi _i:๐ฒ|_{U_i}\stackrel{}{}๐ฒ_i`$ and isomorphisms of functors $`\alpha _{ij}:\phi _{ij}\phi _j\phi _i:๐ฒ|_{U_{ij}}\stackrel{}{}๐ฒ_i|_{U_{ij}}`$, satisfying the natural compatibility conditions.
Note that if the descent datum (2.3) is effective, then $`๐ฒ`$ is unique up to equivalence and such an equivalence is unique up to unique isomorphism.
In the language of $`2`$-categories, the following theorem asserts that the $`2`$-prestack of stacks is a $`2`$-stack.
###### Theorem 2.5.
(cf ) Descent data for stacks are effective.
Denote by $`๐ฒ`$ the stack associated with the descent datum (2.3). Objects of $`๐ฒ(U)`$ can be described by pairs
$$(\{_i\}_{iI},\{\xi _{ij}\}_{i,jI}),$$
(2.4)
where $`_i๐ฒ_i(U_i)`$ and $`\xi _{ij}:\phi _{ij}(_j|_{U_{ij}})\stackrel{}{}_i|_{U_{ij}}`$ are isomorphisms in $`๐ฒ_i(U_{ij})`$ such that for $`i,j,kI`$ the following diagram in $`๐ฒ_i(U_{ijk})`$ commutes
(2.5)
## 3 Twisted modules
Let us now recall how stacks of twisted modules are constructed in .
Let $`๐ค`$ be a commutative unital ring and $`๐`$ a sheaf of $`๐ค`$-algebras on $`X`$. Denote by $`\mathrm{Mod}(๐)`$ the category of left $`๐`$-modules, and by $`\mathrm{๐ฌ๐๐ฝ}(๐)`$ the associated stack $`XU\mathrm{Mod}(๐|_U)`$.
Consider an open covering $`๐ฐ=\{U_i\}_{iI}`$ of $`X`$, a family of $`๐ค`$-algebras $`๐_i`$ on $`U_i`$ and $`๐ค`$-algebra isomorphisms $`f_{ij}:๐_j|_{U_{ij}}\stackrel{}{}๐_i|_{U_{ij}}`$. The existence of a sheaf of $`๐ค`$-algebras locally isomorphic to $`๐_i`$ requires the condition $`f_{ij}f_{jk}=f_{ik}`$ on triple intersections. The weaker conditions (3.2) and (3.3) below are needed for the existence of a $`๐ค`$-additive stack locally equivalent to $`\mathrm{๐ฌ๐๐ฝ}(๐_i)`$.
###### Definition 3.1.
A $`๐ค`$-algebroid descent datum $`๐ `$ on $`๐ฐ`$ is a triplet
$$๐ =(\{๐_i\}_{iI},\{f_{ij}\}_{i,jI},\{a_{ijk}\}_{i,j,kI}),$$
(3.1)
where $`๐_i`$ is a $`๐ค`$-algebra on $`U_i`$, $`f_{ij}:๐_j|_{U_{ij}}\stackrel{}{}๐_i|_{U_{ij}}`$ is a $`๐ค`$-algebra isomorphism, $`a_{ijk}๐_i^\times (U_{ijk})`$ is an invertible section, and (3.2) and (3.3) below are satisfied:
$`f_{ij}f_{jk}`$ $`=\mathrm{Ad}(a_{ijk})f_{ik}\text{ as }๐ค\text{-algebra isomorphisms }๐_k|_{U_{ijk}}\stackrel{}{}๐_i|_{U_{ijk}},`$ (3.2)
$`a_{ijk}a_{ikl}`$ $`=f_{ij}(a_{jkl})a_{ijl}\text{ in }๐_i^\times (U_{ijkl}).`$ (3.3)
(Here $`\mathrm{Ad}(a_{ijk})`$ denotes the automorphism of $`๐_i|_{U_{ijk}}`$ given by $`aa_{ijk}aa_{ijk}^1`$.)
###### Remark 3.2.
The notion of an algebroid stack exists intrinsically, without using coverings or descent data. It has been introduced by and developed in . In this paper, we shall restrict ourselves to algebroids presented by descent data.
###### Remark 3.3.
Let $`_i๐_i^\times `$ be multiplicative subgroups, invariant by $`f_{ij}`$, and such that for any $`b_i,b_i^{}_i`$, the equality $`\mathrm{Ad}(b_i)=\mathrm{Ad}(b_i^{})`$ implies $`b_i=b_i^{}`$. Assume that $`a_{ijk}_i`$. Then, as noticed e.g., in \[9, pag. 2\], condition (3.3) follows from (3.2).
Let us recall how to define the stack of โ$`๐ `$-modulesโ in terms of local data.
A $`๐ค`$-algebra morphism $`f:\stackrel{}{}๐`$ induces a functor
$$\stackrel{~}{f}:\mathrm{๐ฌ๐๐ฝ}(๐)\stackrel{}{}\mathrm{๐ฌ๐๐ฝ}()$$
defined by $`{}_{f}{}^{}`$, where $`{}_{f}{}^{}`$ denotes the sheaf of $`๐ค`$-vector spaces $``$, endowed with the $``$-module structure given by $`bm:=f(b)m`$ for $`b`$ and $`m`$.
For $`a๐^\times `$ an invertible section, the automorphism $`\mathrm{Ad}(a)`$ induces the functor $`\stackrel{~}{\mathrm{Ad}(a)}`$ between $`\mathrm{๐ฌ๐๐ฝ}(๐)`$ and itself, and we denote by
$$\stackrel{~}{a}:\stackrel{~}{\mathrm{Ad}(a)}\mathrm{id}_{\mathrm{๐ฌ๐๐ฝ}(๐)}$$
the isomorphism of functors given by $`\stackrel{~}{a}():{}_{\mathrm{Ad}(a)}{}^{}\stackrel{}{}`$, $`ua^1u`$, for $`u\mathrm{๐ฌ๐๐ฝ}(๐)`$. (Note that $`\stackrel{~}{a}()(a^{}u)=a^1aa^{}a^1u=a^{}\stackrel{~}{a}()(u)`$.)
###### Definition 3.4.
* The stack of twisted modules associated to the $`๐ค`$-algebroid descent datum $`๐ `$ on $`๐ฐ`$ in (3.1) is the stack defined (using Theorem 2.5) by the descent datum
$$\mathrm{๐ฌ๐๐ฝ}(๐ )=(\{\mathrm{๐ฌ๐๐ฝ}(๐_i)\}_{iI},\{\stackrel{~}{f}_{ji}\}_{i,jI},\{\stackrel{~}{a}_{kji}\}_{i,j,kI}).$$
(3.4)
* One sets $`\mathrm{Mod}(๐ ):=\mathrm{๐ฌ๐๐ฝ}(๐ )(X)`$. Objects of the category $`\mathrm{Mod}(๐ )`$ are called twisted modules.
According to (2.4), objects of $`\mathrm{Mod}(๐ )`$ are described by pairs
$$=(\{_i\}_{iI},\{\xi _{ij}\}_{i,jI}),$$
where $`_i`$ are $`๐_i`$-modules and $`\xi _{ij}:{}_{f_{ji}}{}^{}_{j}^{}|_{U_{ij}}\stackrel{}{}_i|_{U_{ij}}`$ are isomorphisms of $`๐_i`$-modules, such that for any $`u_k_k`$ one has
$$\xi _{ij}({}_{f_{ji}}{}^{}\xi _{jk}^{}(u_k))=\xi _{ik}(a_{kji}^1u_k)$$
(3.5)
as morphisms $`{}_{f_{kj}f_{ji}}{}^{}_{k}^{}\stackrel{}{}_i`$. Indeed, (3.5) translates the commutativity of (2.5).
###### Example 3.5.
Let $`X`$ be a complex manifold, and denote by $`\mathrm{\Omega }_X`$ the sheaf of holomorphic forms of maximal degree. Take an open covering $`๐ฐ=\{U_i\}_{iI}`$ of $`X`$ such that there are nowhere vanishing sections $`\omega _i\mathrm{\Omega }_{U_i}`$. Let $`t_{ij}๐ช_{U_{ij}}^\times `$ be the transition functions given by $`\omega _j|_{U_{ij}}=t_{ij}\omega _i|_{U_{ij}}`$. Choose determinations $`s_{ij}๐ช_{U_{ij}}^\times `$ for the multivalued functions $`t_{ij}^{1/2}`$. Since $`s_{ij}s_{jk}`$ and $`s_{ik}`$ are both determinations of $`t_{ik}^{1/2}`$, there exists $`c_{ijk}\{1,1\}`$ such that
$$s_{ij}s_{jk}=c_{ijk}s_{ik}.$$
(3.6)
We thus get a $``$-algebroid descent datum
$$_{\sqrt{\mathrm{\Omega }_X}}:=(\{_{U_i}\}_{iI},\{\mathrm{id}_{_{U_{ij}}}\}_{i,jI},\{c_{ijk}\}_{i,j,kI}).$$
Note that, since $`c_{ijk}^2=1`$, there is an equivalence $`\mathrm{Mod}(_{\sqrt{\mathrm{\Omega }_X}})\mathrm{Mod}(_{\sqrt{\mathrm{\Omega }_X^1}})`$.
Recall from (see also \[5, ยง1\]), that there is an equivalence
$$\mathrm{Mod}(_{\sqrt{\mathrm{\Omega }_X}})\mathrm{Mod}(_X)$$
(3.7)
if and only if the cohomology class $`[c_{ijk}]H^2(X;_X^\times )`$ is trivial. Consider the long exact cohomology sequence
$$H^1(X;_X^\times )\stackrel{๐ผ}{}H^1(X;๐ช_X^\times )\stackrel{๐ฝ}{}H^1(X;d๐ช_X)\stackrel{๐พ}{}H^2(X;_X^\times )$$
associated with the short exact sequence
$$1\stackrel{}{}_X^\times \stackrel{}{}๐ช_X^\times \stackrel{d\mathrm{log}}{}d๐ช_X\stackrel{}{}0.$$
One has $`[c_{ijk}]=\gamma (\frac{1}{2}\beta ([\mathrm{\Omega }_X]))`$, so that $`[c_{ijk}]=1`$ if and only if there exists a line bundle $``$ such that $`\frac{1}{2}\beta ([\mathrm{\Omega }_X])=\beta ([])`$, i.e. such that $`\beta ([\mathrm{\Omega }_X_๐ช^2])=0`$. This last condition holds if and only if there exists a local system of rank one $`L`$ such that $`[\mathrm{\Omega }_X_๐ช^2]=\alpha ([L])`$. Summarizing, (3.7) holds if and only if there exist $``$ and $`L`$ as above, such that
$$\mathrm{\Omega }_XL_{}^2.$$
The twisted sheaf of half-forms in $`\mathrm{Mod}(_{\sqrt{\mathrm{\Omega }_X}})`$ is given by
$$\sqrt{\mathrm{\Omega }_X}=(\{๐ช_{U_i}\}_{iI},\{s_{ij}\}_{i,jI}).$$
We denote by $`\sqrt{\omega _i}`$ the section corresponding to $`1๐ช_{U_i}`$. Hence, on $`U_{ij}`$ we have
$$\sqrt{\omega _j}=s_{ij}\sqrt{\omega _i}.$$
Denote by $`[s_{ij}]`$ the equivalence class of $`s_{ij}`$ in $`๐ช_{U_{ij}}^\times /_{U_{ij}}^\times `$. Since the $`s_{ij}`$โs satisfy (3.6), we notice that
$$\sqrt{\mathrm{\Omega }_X^\times }/_X^\times =(\{๐ช_{U_i}^\times /_{U_i}^\times \}_{iI},\{[s_{ij}]\}_{i,jI})\mathrm{Mod}(_X)$$
(3.8)
is a (usual, i.e. not twisted) sheaf.
Let $`๐ =(\{๐_i\}_{iI},\{f_{ij}\}_{i,jI},\{a_{ijk}\}_{i,j,kI})`$ be a $`๐ค`$-algebroid descent datum on $`๐ฐ`$ as in (3.1), and consider a pair
$$=(\{_i\}_{iI},\{\xi _{ij}\}_{i,jI}),$$
where $`_i`$ are $`๐_i`$-modules and $`\xi _{ij}:{}_{f_{ji}}{}^{}_{j}^{}|_{U_{ij}}\stackrel{}{}_i|_{U_{ij}}`$ are isomorphisms of $`๐_i`$-modules which do not necessarily satisfy (3.5). Assume instead that there are isomorphisms
$$om_{๐_i}({}_{f_{ji}}{}^{}_{j}^{}|_{U_{ij}},_i|_{U_{ij}})๐ค_{U_{ij}}.$$
(3.9)
Under this assumption, we will show in Proposition 3.8 below that $``$ makes sense as a global object of $`\mathrm{๐ฌ๐๐ฝ}(๐ _๐ค๐ฒ)`$, for a suitable twist $`๐ฒ`$.
Consider the isomorphisms
$$\varphi _{ij}:om_{๐_i}({}_{f_{ji}}{}^{}_{j}^{}|_{U_{ij}},_i|_{U_{ij}})\stackrel{}{}๐ค_{U_{ij}},$$
(3.10)
defined by $`\xi _{ij}1`$. The dotted arrow defined by the following commutative diagram of isomorphisms is the multiplication by a section $`c_{ijk}`$ of $`๐ค_{U_{ijk}}^\times `$.
(3.11)
Here, the first vertical arrow on the left follows from the equality $`{}_{f_{ij}f_{jk}}{}^{}_{i}^{}={}_{f_{jk}}{}^{}({}_{f_{ij}}{}^{}_{i}^{})`$, and the second one from $`{}_{f_{ij}f_{jk}}{}^{}_{i}^{}={}_{\mathrm{Ad}(a_{ijk})f_{ik}}{}^{}_{i}^{}`$.
###### Lemma 3.6.
The constants $`c_{ijk}`$ defined above satisfy the cocycle condition
$$c_{ijk}c_{ikl}=c_{jkl}c_{ijl}.$$
* Consider morphisms
$$\eta _{ji}:{}_{f_{ij}}{}^{}_{i}^{}\stackrel{}{}_j.$$
These induce morphisms
$${}_{f_{jk}}{}^{}\eta _{ji}^{}:{}_{f_{jk}}{}^{}({}_{f_{ij}}{}^{}_{i}^{})\stackrel{}{}{}_{f_{jk}}{}^{}_{j}^{}.$$
The composition
$${}_{f_{ik}}{}^{}_{i}^{}\stackrel{\stackrel{~}{a_{ijk}}^1}{}{}_{\mathrm{Ad}(a_{ijk})f_{ik}}{}^{}_{i}^{}={}_{f_{jk}}{}^{}({}_{f_{ij}}{}^{}_{i}^{})\stackrel{{}_{f_{jk}}{}^{}\eta _{ji}^{}}{}{}_{f_{jk}}{}^{}_{j}^{}$$
is given by
$$u_ia_{ijk}u_i{}_{f_{jk}}{}^{}\eta _{ji}^{}(a_{ijk}u_i).$$
Hence, the composition of the vertical isomorphisms in the left column of (3.11) is given by
$$\eta _{ji}_{}\eta _{kj}\eta _{kj}({}_{f_{jk}}{}^{}\eta _{ji}^{}(a_{ijk})).$$
One then has
$$\phi _{ji}(\eta _{ji})\phi _{kj}(\eta _{kj})=\phi _{ki}(\eta _{kj}({}_{f_{jk}}{}^{}\eta _{ji}^{}(a_{ijk})))c_{ijk}.$$
(3.12)
Using (3.12), we have, on one hand
$`\phi _{ji}(\eta _{ji})\phi _{kj}(\eta _{kj})\phi _{lk}(\eta _{lk})`$ $`=\phi _{ji}(\eta _{ji})\phi _{lj}(\eta _{lk}({}_{f_{kl}}{}^{}\eta _{kj}^{}(a_{jkl})))c_{jkl}`$
$`=\phi _{li}(\eta _{lk}({}_{f_{kl}}{}^{}\eta _{kj}^{}(a_{jkl}{}_{f_{jl}}{}^{}\eta _{ji}^{}(a_{ijl}))))c_{ijl}c_{jkl}`$
$`=\phi _{li}(\eta _{lk}({}_{f_{kl}}{}^{}\eta _{kj}^{}({}_{f_{jk}f_{kl}}{}^{}\eta _{ji}^{}(f_{ij}(a_{jkl})a_{ijl}))))c_{ijl}c_{jkl},`$
and on the other hand
$`\phi _{ji}(\eta _{ji})\phi _{kj}(\eta _{kj})\phi _{lk}(\eta _{lk})`$ $`=\phi _{ki}(\eta _{kj}({}_{f_{jk}}{}^{}\eta _{ji}^{}(a_{ijk})))\phi _{lk}(\eta _{lk})c_{ijk}`$
$`=\phi _{li}(\eta _{lk}({}_{f_{kl}}{}^{}\eta _{kj}^{}({}_{f_{jk}f_{kl}}{}^{}\eta _{ji}^{}(a_{ikl}a_{ijk}))))c_{ikl}c_{ijk}.`$
The conclusion follows using (3.3). โ
###### Remark 3.7.
Lemma 3.6 is a particular case of a general result which asserts that equivalence classes of locally trivial $`๐ค`$-algebroids on $`X`$ are in one-to-one correspondence with $`H^2(X;๐ค^\times )`$. The analogue result for gerbes is discussed in , and we refer to for the formulation in terms of algebroids.
Let us recall our setting. Given a $`๐ค`$-algebroid descent datum
$$๐ =(\{๐_i\}_{iI},\{f_{ij}\}_{i,jI},\{a_{ijk}\}_{i,j,kI}),$$
consider a pair
$$=(\{_i\}_{iI},\{\xi _{ij}\}_{i,jI}),$$
where $`_i`$ are $`๐_i`$-modules and $`\xi _{ij}:{}_{f_{ji}}{}^{}_{j}^{}|_{U_{ij}}\stackrel{}{}_i|_{U_{ij}}`$ are isomorphisms of $`๐_i`$-modules. Assuming (3.9), Lemma 3.6 guarantees that the constants $`c_{ijk}`$ defined by (3.11) satisfy the cocycle condition. We can thus consider the $`๐ค`$-algebroid descent datum
$$๐ฒ=(\{๐ค_{U_i}\}_{iI},\{\mathrm{id}_{๐ค_{U_{ij}}}\}_{i,jI},\{c_{ijk}\}_{i,j,kI}).$$
Set
$$๐ _๐ค๐ฒ=(\{๐_i\}_{iI},\{f_{ij}\}_{i,jI},\{a_{ijk}c_{ijk}\}_{i,j,kI}).$$
###### Proposition 3.8.
Let $`=(\{_i\}_{iI},\{\xi _{ij}\}_{i,jI})`$ be as above, and assume (3.9), then
$$=(_i,\xi _{ij})\mathrm{Mod}(๐ _๐ค๐ฒ).$$
* By (3.5), it is enough to show that
$$\xi _{ik}=\xi _{ij}({}_{f_{ji}}{}^{}\xi _{jk}^{}(a_{kji}c_{kji})).$$
We have $`\phi _{ik}(\xi _{ik})=1=\phi _{jk}(\xi _{jk})\phi _{ij}(\xi _{ij})=\phi _{ik}(\xi _{ij}({}_{f_{ji}}{}^{}\xi _{jk}^{}(a_{kji}c_{kji})))`$ by (3.12). The conclusion follows since $`\phi _{ik}`$ is an isomorphism. โ
## 4 Microdifferential modules
Here we review a few notions from the theory of microdifferential modules. References are made to and also to for complementary results. See for an exposition.
Let $`Y`$ be a complex analytic manifold, and $`\pi :T^{}Y\stackrel{}{}Y`$ its cotangent bundle. The sheaf $`_Y`$ of microdifferential operators on $`T^{}Y`$ is a $``$-central algebra endowed with a $``$-filtration by the order. Denote by $`_Y(m)`$ its subsheaf of operators of order at most $`m`$, and by $`๐ช_{T^{}Y}(m)`$ the sheaf of functions homogeneous of degree $`m`$ in the fiber of $`\pi `$. Denote by $`\mathrm{eu}`$ the Euler vector field, i.e. the infinitesimal generator of the action of $`^\times `$ on $`T^{}Y`$. Then $`f๐ช_{T^{}Y}(m)`$ if and only if $`\mathrm{eu}f=mf`$. In a homogeneous symplectic local coordinate system $`(x;\xi )`$ on $`UT^{}Y`$, a section $`P\mathrm{\Gamma }(U;_Y(m))`$ is written as a formal series
$$P=\underset{jm}{}p_j(x;\xi ),p_j\mathrm{\Gamma }(U;๐ช_{T^{}Y}(j)),$$
(4.1)
with the condition that for any compact subset $`K`$ of $`U`$ there exists a constant $`C_K>0`$ such that $`\underset{K}{sup}|p_j|C_K^j(j)!`$ for all $`j<0`$.
If $`Q=_{jn}q_j`$ is another section, the product $`PQ=R=_{jm+n}r_j`$ is given by the Leibniz rule
$$r_k=\underset{k=i+j|\alpha |}{}\frac{1}{\alpha !}(_\xi ^\alpha p_i)(_x^\alpha q_j).$$
The symbol map
$$\sigma _m:_Y(m)\stackrel{}{}๐ช_{T^{}Y}(m),Pp_m$$
does not depend on the choice of coordinates and induces the symbol map
$$\sigma :_Y\stackrel{}{}gr_Y\stackrel{}{}\underset{m}{}๐ช_{T^{}Y}(m).$$
The formal adjoint of $`P=_{jm}p_j`$ is defined by
$$P^{}=\underset{jm}{}p_j^{},p_j^{}(x;\xi )=\underset{j=k|\alpha |}{}\frac{(1)^{|\alpha |}}{\alpha !}_\xi ^\alpha _x^\alpha p_k(x;\xi ).$$
It depends on the choice of coordinates, and more precisely on the choice of the top degree form $`dx_i\mathrm{}dx_n\mathrm{\Omega }_Y`$. One thus considers the twist of $`_Y`$ by half-forms
$$_Y^\sqrt{v}:=\pi ^1\sqrt{\mathrm{\Omega }_Y}_{\pi ^1๐ช_Y}_Y_{\pi ^1๐ช_Y}\pi ^1\sqrt{\mathrm{\Omega }_Y^1}.$$
This is a sheaf of filtered $``$-algebras endowed with a canonical anti-isomorphism
$$:_Y^\sqrt{v}\stackrel{}{}a_{}_Y^\sqrt{v},$$
where $`a`$ denotes the antipodal map on $`T^{}Y`$. There is a subprincipal symbol
$$\sigma _{m1}^{}:_Y^\sqrt{v}(m)\stackrel{}{}๐ช_{T^{}Y}(m1),P\frac{1}{2}\sigma _{m1}(P(1)^mP^{}).$$
In local coordinates, $`\sigma _{m1}^{}(P)=p_{m1}\frac{1}{2}_i_{x_i}_{\xi _i}p_m`$ for $`P`$ as in (4.1).
The following definition is adapted from .
###### Definition 4.1.
Let $`\mathrm{\Lambda }`$ be a smooth, locally closed, $`^\times `$-conic submanifold of $`T^{}Y`$. Let $``$ be a coherent $`_Y`$-module supported by $`\mathrm{\Lambda }`$.
1. One says that $``$ is regular (resp. simple) along $`\mathrm{\Lambda }`$ if there locally exists a coherent sub-$`_Y(0)`$-module $`_0`$ of $``$ which generates it over $`_Y`$, and such that $`_0/_Y(1)_0`$ is an $`๐ช_\mathrm{\Lambda }(0)`$-module (resp. a locally free $`๐ช_\mathrm{\Lambda }(0)`$-module of rank one).
2. Let $``$ be simple along $`\mathrm{\Lambda }`$. A section $`u`$ is called a simple generator if $`_0=_Y(0)u`$ satisfies the conditions in (a), and the image of $`u`$ in $`_0/_Y(1)_0`$ generates this module over $`๐ช_\mathrm{\Lambda }(0)`$.
Set
$$_\mathrm{\Lambda }=\{P_Y(1)|_\mathrm{\Lambda };\sigma _1(P)|_\mathrm{\Lambda }=0\},$$
(4.2)
and denote by $`_\mathrm{\Lambda }`$ the sub-algebra of $`_Y|_\mathrm{\Lambda }`$ generated by $`_\mathrm{\Lambda }`$.
###### Remark 4.2.
Let $`u`$ be a generator of a coherent $`_Y`$-module $``$. Then $`_Y/`$, where $`=\{P_Y:Pu=0\}`$. Set
$$\overline{}=\{\sigma _m(P);m,P_Y(m)\},$$
and note that $`\mathrm{supp}=\mathrm{supp}(๐ช_{T^{}Y}/\overline{})`$. Then $``$ is simple if and only if there locally exists a generator $`u`$ such that the ideal $`\overline{}`$ is reduced. Moreover, such a section $`u`$ is a simple generator and the sub-$`_Y(0)`$-module $`_0`$ generated by $`u`$ satisfies
$$_\mathrm{\Lambda }_0_0.$$
Indeed, in a homogeneous symplectic local coordinate system we may write $`P_\mathrm{\Lambda }`$ as $`P=P^{}+Q`$ with $`Q`$ of order $`0`$ and $`P^{}u=0`$.
### Symbol of sections of simple systems
Let us recall the notion of symbol for simple generators. References are made to .
For a vector field $`v\mathrm{\Theta }_\mathrm{\Lambda }`$ on $`\mathrm{\Lambda }`$, denote by $`_v^{1/2}`$ its Lie derivative action on the twisted sheaf $`\sqrt{\mathrm{\Omega }_\mathrm{\Lambda }}`$. Then $`_v^{1/2}`$ is an operator of order one in the ring $`๐_\mathrm{\Lambda }^\sqrt{v}=\sqrt{\mathrm{\Omega }_\mathrm{\Lambda }}_{๐ช_\mathrm{\Lambda }}๐_\mathrm{\Lambda }_{๐ช_\mathrm{\Lambda }}\sqrt{\mathrm{\Omega }_\mathrm{\Lambda }^1}`$
of differential operators acting on $`\sqrt{\mathrm{\Omega }_\mathrm{\Lambda }}`$.
Define $`_\mathrm{\Lambda }^\sqrt{v}`$ and $`_\mathrm{\Lambda }^\sqrt{v}`$ as in (4.2), replacing $`_Y`$ with $`_Y^\sqrt{v}`$, and denote by $`H_f`$ the Hamiltonian vector field of $`f๐ช_{T^{}Y}`$. Note that $`H_fT\mathrm{\Lambda }`$ if $`f|_\mathrm{\Lambda }=0`$. For $`P_\mathrm{\Lambda }^\sqrt{v}`$, consider the transport operator
$$\mathrm{L}(P)=_{H_{\sigma _1(P)}|_\mathrm{\Lambda }}^{1/2}+\sigma _0^{}(P)|_\mathrm{\Lambda }.$$
One checks that $`\mathrm{L}`$ satisfies the relations $`\mathrm{L}(AP)=\sigma _0(A)\mathrm{L}(P)`$, $`\mathrm{L}(PA)=\mathrm{L}(P)\sigma _0(A)`$, and $`\mathrm{L}([P,Q])=[\mathrm{L}(P),\mathrm{L}(Q)]`$, for $`P,Q_\mathrm{\Lambda }^\sqrt{v}`$ and $`A_\mathrm{\Lambda }^\sqrt{v}(0)`$ (see e.g. \[10, ยง8.3\]). It follows that $`\mathrm{L}`$ extends as a $``$-algebra morphism
$$\mathrm{L}:_\mathrm{\Lambda }^\sqrt{v}\stackrel{}{}๐_\mathrm{\Lambda }^\sqrt{v}$$
(4.3)
by setting $`\mathrm{L}(P_1\mathrm{}P_r)=\mathrm{L}(P_1)\mathrm{}\mathrm{L}(P_r)`$, for $`P_i_\mathrm{\Lambda }^\sqrt{v}`$.
Let $``$ be a simple $`_Y^\sqrt{v}`$-module along $`\mathrm{\Lambda }`$, and $`u`$ a simple generator. The twisted subsheaf of $`\sqrt{\mathrm{\Omega }_\mathrm{\Lambda }}`$ defined by
$$\{\sigma \sqrt{\mathrm{\Omega }_\mathrm{\Lambda }};\mathrm{L}(P)\sigma =0P_\mathrm{\Lambda }^\sqrt{v},Pu=0\}$$
(4.4)
is locally a free sheaf of rank one over $``$.
###### Definition 4.3.
Let $``$ be a simple $`_Y^\sqrt{v}`$-module along $`\mathrm{\Lambda }`$ and let $`u`$ be a simple generator. The symbol of $`u`$ is defined by
$$\sigma _\mathrm{\Lambda }(u)=[\sigma ]\sqrt{\mathrm{\Omega }_\mathrm{\Lambda }^\times }/_\mathrm{\Lambda }^\times .$$
for $`\sigma `$ as in (4.4).
The Euler vector field $`\mathrm{eu}`$ acts on $`\sqrt{\mathrm{\Omega }_\mathrm{\Lambda }}`$ by $`_{\mathrm{eu}}^{1/2}`$, and one says that $`\sigma \sqrt{\mathrm{\Omega }_\mathrm{\Lambda }}`$ is homogeneous of degree $`\lambda `$ if $`\mathrm{eu}\sigma =\lambda \sigma `$. Hence, the notion of homogeneous section of $`\sqrt{\mathrm{\Omega }_\mathrm{\Lambda }^\times }/_\mathrm{\Lambda }^\times `$ makes sense. One calls order of $`u`$ the homogeneous degree of $`\sigma _\mathrm{\Lambda }(u)`$. Then the equivalence class of $`\lambda `$ in $`/`$ does not depend on $`u`$, and is called the order of $``$.
If $`P_Y^\sqrt{v}(m)`$ is such that $`\sigma _m(P)|_\mathrm{\Lambda }`$ never vanishes, then
$$\sigma _\mathrm{\Lambda }(Pu)=\sigma _m(P)\sigma _\mathrm{\Lambda }(u).$$
(4.5)
Also recall from loc.cit. that simple modules of the same order are locally isomorphic.
## 5 Quantization of contact manifolds
Here we review Kashiwaraโs construction of the stack of microdifferential modules on a contact manifold.
Let $`Y`$ be a complex analytic manifold, $`\pi :T^{}Y\stackrel{}{}Y`$ its cotangent bundle, $`\dot{T}^{}Y=T^{}YY`$ the complementary of the zero-section, $`\varpi :P^{}Y\stackrel{}{}Y`$ the projective cotangent bundle, and $`\gamma :\dot{T}^{}Y\stackrel{}{}P^{}Y`$ the projection. The sheaf of microdifferential operators on $`P^{}Y`$ is given by $`\gamma _{}(_Y|_{\dot{T}^{}X})`$, and we still denote it by $`_Y`$ for short. Since the antipodal map induces the identity on $`P^{}Y`$, the anti-involution $``$ is well defined on the sheaf $`_Y^\sqrt{v}`$.
Let $`\chi _{ij}:P^{}Y_iV_i\stackrel{}{}V_jP^{}Y_j`$ be a contact transformation. Recall that there locally exists a quantized contact transformation (QCT for short) above $`\chi _{ij}`$. This is an isomorphism of filtered $``$-algebras $`\mathrm{\Phi }_{ij}:\chi _{ij}^1_{Y_j}^\sqrt{v}|_{V_j}\stackrel{}{}_{Y_i}^\sqrt{v}|_{V_i}`$. Moreover, one can ask that $`\mathrm{\Phi }_{ij}`$ is $``$-preserving. Such a quantization is not unique, but a key remark by Kashiwara is that $``$-preserving filtered automorphisms of $`_Y^\sqrt{v}`$ are of the form $`\mathrm{Ad}(Q)`$ for a unique operator $`Q_Y^\sqrt{v}(0)`$ satisfying
$$QQ^{}=1,\sigma _0(Q)=1.$$
(5.1)
###### Definition 5.1.
A complex contact manifold $`๐=(๐,๐ช_๐(1),\alpha )`$ is a complex manifold $`๐`$ of dimension $`2n+1`$ endowed with a line bundle $`๐ช_๐(1)`$ and a 1-form $`\alpha \mathrm{\Gamma }(๐,\mathrm{\Omega }_๐^1_๐ช๐ช_๐(1))`$, such that $`\alpha (d\alpha )^n`$ is a non-degenerate section of $`\mathrm{\Omega }_๐^{2n+1}_๐ช๐ช_๐(n+1)`$.
(Here we set $`๐ช_๐(k)=๐ช_๐(1)^k`$, and we use the fact that $`\alpha (d\alpha )^r`$ is a well-defined section of $`\mathrm{\Omega }_๐^{2r+1}_๐ช๐ช_๐(r+1)`$ for $`0rn`$.)
There is an open covering $`๐ฑ=\{V_i\}_{iI}`$ of $`๐`$ and contact embeddings $`\chi _i:V_iP^{}Y_i`$. Up to refining the covering (we still denote it by $`๐ฑ`$), the induced contact transformations $`\chi _{ij}:\chi _i(V_{ij})\stackrel{}{}\chi _j(V_{ij})`$ can be quantized to a $``$-preserving filtered $``$-algebra isomorphism
$$\mathrm{\Phi }_{ij}:\chi _{ij}^1(_{Y_j}^\sqrt{v}|_{\chi _j(V_{ij})})\stackrel{}{}_{Y_i}^\sqrt{v}|_{\chi _i(V_{ij})}.$$
The composition $`\mathrm{\Phi }_{ij}\mathrm{\Phi }_{jk}\mathrm{\Phi }_{ik}^1`$ is a $``$-preserving automorphism of $`_{Y_i}^\sqrt{v}|_{\chi _i(V_{ijk})}`$, and hence is equal to $`\mathrm{Ad}(Q_{ijk})`$ for a unique $`Q_{ijk}_{Y_i}^\sqrt{v}(V_{ijk})`$ satisfying (5.1). This proves the theorem below, thanks to Remark 3.3.
###### Theorem 5.2.
(cf \[9, Theorem 2\]) The triplet
$$๐ค_๐=(\{\chi _i^1_{Y_i}^\sqrt{v}|_{V_i}\}_{iI},\{\chi _i^1(\mathrm{\Phi }_{ij})\}_{i,jI},\{\chi _i^1(Q_{ijk})\}_{i,j,kI})$$
(5.2)
is a $``$-algebroid descent datum over $`๐`$. In particular, there is an associated stack $`\mathrm{๐ฌ๐๐ฝ}(๐ค_๐)`$ on $`๐`$ locally equivalent to the stack $`\mathrm{๐ฌ๐๐ฝ}(_{Y_i}^\sqrt{v})`$ of microdifferential modules.
### Good modules
Any local notion, such as that of being coherent, simple, or regular, immediately extends to the category $`\mathrm{Mod}(๐ค_๐)`$. Here, we will discuss the non local notion of being good (cf ).
Remark first that the construction in Theorem 5.2 also applies when replacing $`_{Y_i}`$ with $`_{Y_i}(0)`$. One thus gets a $``$-algebroid descent datum $`๐ค_๐(0)`$ as well as a $``$-linear functor $`๐ค_๐(0)\stackrel{}{}๐ค_๐`$. This induces a forgetful functor
$$\mathrm{Mod}(๐ค_๐)\stackrel{๐๐๐}{}\mathrm{Mod}(๐ค_๐(0)).$$
This functor locally admits a left adjoint, hence it has a left adjoint
$$\mathrm{Mod}(๐ค_๐(0))\stackrel{๐๐ฅ๐ก}{}\mathrm{Mod}(๐ค_๐).$$
###### Definition 5.3.
A coherent module $`\mathrm{Mod}(๐ค_๐)`$ is good if for any relatively compact open subset $`U๐`$ there exists $`_0\mathrm{Mod}(๐ค_๐(0)|_U)`$ such that $`|_U\mathrm{๐๐ฅ๐ก}(_0)`$.
### Quantization with parameters
Assume now that the bundle $`๐ช_๐(1)`$ is trivial, i.e. that there exists a nowhere vanishing section $`\tau \mathrm{\Gamma }(๐;๐ช_๐(1))`$. Consider an open covering $`\{V_i\}_{iI}`$ of $`๐`$ and contact embeddings $`\chi _i:V_iP^{}Y_i`$ to which is attached the $``$-algebroid descent datum (5.2).
###### Lemma 5.4.
Up to refining the covering, there exist $``$-preserving QCTโs $`\mathrm{\Phi }_{ij}`$ above $`\chi _{ij}`$, and sections $`T_i\mathrm{\Gamma }(\chi _i(V_i);_{Y_i}^\sqrt{v}(1))`$, such that
$$\sigma (T_i)\chi _i=\tau ,T_i^{}=T_i,\mathrm{\Phi }_{ij}(T_j)=T_i.$$
* Let $`(t,x_1,\mathrm{},x_n,u_1,\mathrm{},u_n)`$ be a local coordinate system on $`V_i`$ such that the contact form is given by $`dt+_iu_idx_i`$ and $`\tau =\sigma (_t)\chi _i`$. Then set $`T_i=_t`$.
Let $`\mathrm{\Phi }_{ij}^{}`$ be a QCT above $`\chi _{ij}`$ such that $`\mathrm{\Phi }_{ij}^{}(T_j)=T_i`$. Then the proof goes as that of \[19, Lemma 5.3 (iii)\]. Consider the QCT above the identity given by
$$\mathrm{\Phi }_{ij}^{}:=\mathrm{\Phi }_{ij}^1\mathrm{\Phi }_{ij}^{}.$$
There exists $`Q_{Y_j}^\sqrt{v}`$ such that $`\mathrm{\Phi }_{ij}^{}=\mathrm{Ad}(QQ^{})`$. Hence $`\mathrm{\Phi }_{ij}:=\mathrm{\Phi }_{ij}^{}\mathrm{Ad}(Q)`$ is $``$-preserving QCT above $`\chi _{ij}`$, and $`\mathrm{\Phi }_{ij}(T_j)=T_i`$. โ
Denote by $`_{Y_i,T_i}^\sqrt{v}`$ the subalgebra of $`_{Y_i}^\sqrt{v}`$ of operators commuting with $`T_i`$. As above, denote by $`Q_{ijk}_{Y_i}^\sqrt{v}(\chi _i(V_{ijk}))`$ the unique operator satisfying (5.1) such that $`\mathrm{\Phi }_{ij}\mathrm{\Phi }_{jk}\mathrm{\Phi }_{ik}^1=\mathrm{Ad}(Q_{ijk})`$. Since $`\mathrm{Ad}(Q_{ijk})(T_i)=T_i`$, it follows that $`Q_{ijk}`$ is a section of $`_{Y_i,T_i}^\sqrt{v}`$.
###### Proposition 5.5.
Let $`๐`$ be a complex contact manifold. Assume that there exists a nowhere vanishing section $`\tau \mathrm{\Gamma }(๐;๐ช_๐(1))`$. Then the triplet
$$๐ค_{๐,\tau }=(\{\chi _i^1_{Y_i,T_i}^\sqrt{v}|_{\chi _i(V_i)}\}_{iI},\{\chi _i^1(\mathrm{\Phi }_{ij})\}_{i,jI},\{\chi _i^1(Q_{ijk})\}_{i,j,kI})$$
(5.3)
is a $``$-algebroid descent datum on $`๐ฑ`$. In particular, there is an associated stack $`\mathrm{๐ฌ๐๐ฝ}(๐ค_{๐,\tau })`$ on $`๐`$ locally equivalent to the stack $`\mathrm{๐ฌ๐๐ฝ}(_{Y_i,T_i}^\sqrt{v})`$.
## 6 Simple holonomic modules on contact manifolds
Here, we give a proof of a result of Kashiwara on the existence of twisted simple holonomic modules along smooth Lagrangian submanifolds (also called Legendrian in the literature) of complex contact manifolds.
Let $`๐`$ be a complex contact manifold. Recall from Theorem 5.2 that there is an algebroid descent datum $`๐ค_๐`$ of microdifferential operators on $`๐`$. Let $`\mathrm{\Lambda }๐`$ be a Lagrangian submanifold. With notations as in Example 3.5, consider the stack $`\mathrm{๐ฌ๐๐ฝ}(๐ค_๐|_\mathrm{\Lambda }_{}_{\sqrt{\mathrm{\Omega }_\mathrm{\Lambda }}})`$ of twisted microdifferential modules on $`\mathrm{\Lambda }`$.
Let us recall Kashiwaraโs theorem announced in .
###### Theorem 6.1.
Let $`๐`$ be a complex contact manifold and let $`\mathrm{\Lambda }๐`$ be a Lagrangian submanifold. There exists $`\mathrm{Mod}(๐ค_๐|_\mathrm{\Lambda }_{}_{\sqrt{\mathrm{\Omega }_\mathrm{\Lambda }}})`$ which is simple along $`\mathrm{\Lambda }`$.
###### Remark 6.2.
It will follow from the proof that moreover $``$ is good.
Here, we give a proof of this result using the notion of symbol for simple sections of holonomic modules recalled in Section 4.
* Let us denote for short by
$$๐ค_๐=(\{๐_i\}_{iI},\{f_{ij}\}_{i,jI},\{a_{ijk}\}_{i,j,kI})$$
the $``$-algebroid descent datum (5.2) attached to an open covering $`๐=_{iI}V_i`$. Up to a refinement, we may assume that the $``$-algebroid descent datum $`_{\sqrt{\mathrm{\Omega }_\mathrm{\Lambda }}}`$ is attached to the same covering. More precisely, using notations as in Example 3.5 for $`X=\mathrm{\Lambda }`$ and $`U_i=\mathrm{\Lambda }_i=\mathrm{\Lambda }V_i`$, we assume that there are nowhere vanishing sections $`\omega _i\mathrm{\Omega }_{\mathrm{\Lambda }_i}`$, and $`s_{ij}๐ช_{\mathrm{\Lambda }_{ij}}^\times `$ with $`\sqrt{\omega _j}=s_{ij}\sqrt{\omega _i}`$, such that
$$_{\sqrt{\mathrm{\Omega }_X}}=(\{_{\mathrm{\Lambda }_i}\}_{iI},\{\mathrm{id}_{_{\mathrm{\Lambda }_{ij}}}\}_{i,jI},\{c_{ijk}\}_{i,j,kI}),$$
where the $`c_{ijk}`$โs are defined by
$$s_{ij}s_{jk}=c_{ijk}s_{ik}.$$
Up to a further refinement of the covering, we may assume that there exist simple $`๐_i`$-modules $`_i`$ of order $`0`$ along $`\mathrm{\Lambda }_i`$ and simple generators $`u_i`$ of $`_i`$ such that
$$\sigma _{\mathrm{\Lambda }_i}(u_i)=[\sqrt{\omega _i}](\sqrt{\mathrm{\Omega }_\mathrm{\Lambda }^\times }/_\mathrm{\Lambda }^\times )|_{\mathrm{\Lambda }_i}.$$
Since $`_i`$ and $`{}_{f_{ji}}{}^{}_{j}^{}`$ are simple $`๐_i`$-modules along $`\mathrm{\Lambda }_{ij}`$ of the same order, they are isomorphic, and one has
$$om_{๐_i}({}_{f_{ji}}{}^{}_{j}^{},_i)_{\mathrm{\Lambda }_{ij}}.$$
Let us describe explicitly such an isomorphism. For $`\stackrel{~}{\xi }_{ij}om_{๐_i}({}_{f_{ji}}{}^{}_{j}^{},_i)`$, let $`\stackrel{~}{b}_{ij}`$ is a section of $`๐_i`$ satisfying $`\stackrel{~}{\xi }_{ij}(u_j)=\stackrel{~}{b}_{ij}u_i`$. One has
$`\sigma _\mathrm{\Lambda }(\stackrel{~}{\xi }_{ij}(u_j))`$ $`=\sigma _\mathrm{\Lambda }(u_j)=[\sqrt{\omega _j}]=[s_{ij}\sqrt{\omega _i}],`$
$`\sigma _\mathrm{\Lambda }(\stackrel{~}{b}_{ij}u_i)`$ $`=\sigma (\stackrel{~}{b}_{ij})\sigma _\mathrm{\Lambda }(u_i)=[\sigma (\stackrel{~}{b}_{ij})\sqrt{\omega _i}],`$
where the fourth equality follows from (4.5). In particular, $`s_{ij}^1\sigma (\stackrel{~}{b}_{ij})_{\mathrm{\Lambda }_{ij}}^\times `$. We can thus consider the isomorphism
$`\phi _{ij}:om_{๐_i}({}_{f_{ji}}{}^{}_{j}^{},_i)`$ $`\stackrel{}{}_{\mathrm{\Lambda }_{ij}}`$
$`\stackrel{~}{\xi }_{ij}`$ $`s_{ij}^1\sigma (\stackrel{~}{b}_{ij}).`$
Set
$$\xi _{ij}=\phi _{ij}^1(1):{}_{f_{ji}}{}^{}_{j}^{}\stackrel{}{}_i,$$
and let $`b_{ij}๐_i`$ satisfy $`\xi _{ij}(u_j)=b_{ij}u_i`$. As $`\phi _{ij}(\xi _{ij})=1`$, one has
$$\sigma (b_{ij})=s_{ij}.$$
(6.1)
By Proposition 3.8, we get a twisted module
$$=(_i,\xi _{ij})\mathrm{Mod}(๐ค_๐|_\mathrm{\Lambda }_{}๐ฒ),$$
where
$$๐ฒ=(\{_{\mathrm{\Lambda }_i}\}_{iI},\{\mathrm{id}_{_{\mathrm{\Lambda }_{ij}}}\}_{i,jI},\{\stackrel{~}{c}_{ijk}\}_{i,j,kI}),$$
the $`\stackrel{~}{c}_{ijk}`$โs being given by (3.11). Since $`\phi _{ij}(\xi _{ij})=1=\phi _{jk}(\xi _{jk})`$, by (3.12) we have
$$\stackrel{~}{c}_{ijk}=\phi _{ik}(\xi _{ij}({}_{f_{ji}}{}^{}\xi _{jk}^{}(a_{kji}))).$$
To show that $`๐ฒ=_{\sqrt{\mathrm{\Omega }_\mathrm{\Lambda }}}`$, it thus remains to prove that $`\stackrel{~}{c}_{ijk}=c_{ijk}`$, i.e. that
$$s_{ij}s_{jk}=\stackrel{~}{c}_{ijk}s_{ik}.$$
One has
$$\begin{array}{cc}\hfill \xi _{ij}({}_{f_{ji}}{}^{}\xi _{jk}^{}(a_{kji}u_k))& =f_{ij}f_{jk}(a_{kji})\xi _{ij}({}_{f_{ji}}{}^{}\xi _{jk}^{}(u_k))\hfill \\ & =f_{ij}f_{jk}(a_{kji})\xi _{ij}(b_{jk}u_j)\hfill \\ & =f_{ij}f_{jk}(a_{kji})f_{ij}(b_{jk})\xi _{ij}(u_j)\hfill \\ & =f_{ij}f_{jk}(a_{kji})f_{ij}(b_{jk})b_{ij}u_i.\hfill \end{array}$$
Then
$$\begin{array}{cc}\hfill \stackrel{~}{c}_{ijk}& =\phi _{ik}(\xi _{ij}({}_{f_{ji}}{}^{}\xi _{jk}^{}(a_{kji})))\hfill \\ & =s_{ik}^1\sigma (f_{ij}f_{jk}(a_{kji})f_{ij}(b_{jk})b_{ij})\hfill \\ & =s_{ik}^1\sigma (b_{jk})\sigma (b_{ij})\hfill \\ & =s_{ik}^1s_{jk}s_{ij},\hfill \end{array}$$
where the third equality follows from the fact that $`\sigma (a_{kji})=1`$, and the fourth one from (6.1). โ
###### Definition 6.3.
A complex homogeneous symplectic manifold $`=(,\omega ,v)`$ is a complex symplectic manifold $`(,\omega )`$ endowed with a vector field $`v`$ satisfying $`_v\omega =\omega `$.
Corollary 6.4 below was announced in . Although we shall not use it, we give a proof for the readerโs convenience. Also note that our statement corrects that in loc. cit., following a private communication with M. Kashiwara.
Let $`๐=(๐,๐ช_๐(1),\alpha )`$ be a contact manifold and denote by $`\gamma :\stackrel{~}{๐}\stackrel{}{}๐`$ the total space of the $`^\times `$-principal bundle associated with the dual of the line bundle $`๐ช_๐(1)`$. Then $`\stackrel{~}{๐}`$ is a homogeneous symplectic manifold and there exists a covering $`\{V_i\}_{iI}`$ of $`๐`$ by contact charts $`\chi _i:V_iP^{}Y_i`$ such that, setting $`\stackrel{~}{V}_i=\gamma ^1V_i`$, $`\{\stackrel{~}{V}_i\}_{iI}`$ is a covering of $`\stackrel{~}{๐}`$ and there are commutative diagrams
where $`\gamma _i`$ is the the projection $`\dot{T}^{}Y_i\stackrel{}{}P^{}Y_i`$ and $`\stackrel{~}{\chi }_i:\stackrel{~}{V}_i\dot{T}^{}Y_i`$ are homogeneous symplectic maps.
We denote by $`\mathrm{๐ฌ๐๐ฝ}_{\text{loc-sys}}(_{\gamma ^1\mathrm{\Lambda }})`$ the full substack of $`\mathrm{๐ฌ๐๐ฝ}(_{\gamma ^1\mathrm{\Lambda }})`$ consisting of local systems, and by $`\mathrm{๐ฌ๐๐ฝ}_{\text{reg-}\mathrm{\Lambda }}(๐ค_๐|_\mathrm{\Lambda })`$ the full substack of $`\mathrm{๐ฌ๐๐ฝ}(๐ค_๐|_\mathrm{\Lambda })`$ consisting of modules regular along $`\mathrm{\Lambda }`$
###### Corollary 6.4.
(cf ) There is an equivalence of stacks
$$\mathrm{๐ฌ๐๐ฝ}_{\text{reg-}\mathrm{\Lambda }}(๐ค_๐|_\mathrm{\Lambda }_{}_{\sqrt{\mathrm{\Omega }_\mathrm{\Lambda }}})\gamma _{}\mathrm{๐ฌ๐๐ฝ}_{\text{loc-sys}}(_{\gamma ^1\mathrm{\Lambda }}).$$
* The sheaf of rings $`_Y`$ on $`T^{}Y`$ is a subsheaf of the sheaf of rings $`_Y^{}`$ of , this last sheaf being the microlocalization along the diagonal of the sheaf $`๐ช_{Y\times Y}`$ (up to a shift and tensorizing by holomorphic forms), see \[13, chapter 11\] for a detailed construction. One defines the sheaf $`_Y^{\sqrt{v},}`$ similarly as we have defined $`_Y^\sqrt{v}`$.
With notations as in (5.2), the isomorphisms $`\mathrm{\Phi }_{ij}`$ are induced by sections of a microdifferential module supported by the graph of $`\chi _{ij}`$, and hence extend to isomorphisms $`\mathrm{\Phi }_{ij}^{}`$. Considering $`Q_{ijk}`$ as sections of $`_{Y_i}^{\sqrt{v},}`$, we get a $``$-algebroid descent datum
$$๐ค_{\stackrel{~}{๐}}^{}=(\{\stackrel{~}{\chi }_i^1_{Y_i}^{\sqrt{v},}|_{\stackrel{~}{\chi }_i(\stackrel{~}{V}_i)}\}_{iI},\{\stackrel{~}{\chi }_i^1(\mathrm{\Phi }_{ij}^{})\}_{i,jI},\{\stackrel{~}{\chi }_i^1(Q_{ijk})\}_{i,j,kI}),$$
and the associated stack $`\mathrm{๐ฌ๐๐ฝ}(๐ค_{\stackrel{~}{๐}}^{})`$ on $`\stackrel{~}{๐}`$. The inclusion $`\gamma ^1_Y^\sqrt{v}_Y^{}`$ induces a functor of extension of scalars $`\mathrm{๐ฌ๐๐ฝ}(๐ค_๐)\stackrel{}{}\gamma _{}\mathrm{๐ฌ๐๐ฝ}(๐ค_{\stackrel{~}{๐}}^{})`$. Let us denote by $`^{}`$ this functor.
By Theorem 6.1 there is a simple system $``$ in $`\mathrm{Mod}(๐ค_๐|_\mathrm{\Lambda }_{}_{\sqrt{\mathrm{\Omega }_\mathrm{\Lambda }}})`$. Consider the functor
$$\mathrm{\Psi }:\mathrm{๐ฌ๐๐ฝ}_{\text{reg-}\mathrm{\Lambda }}(๐ค_๐|_\mathrm{\Lambda }_{}_{\sqrt{\mathrm{\Omega }_\mathrm{\Lambda }}})\stackrel{}{}\gamma _{}\mathrm{๐ฌ๐๐ฝ}_{\text{loc-sys}}(_{\gamma ^1\mathrm{\Lambda }})$$
given by $`om_{๐ค_{\stackrel{~}{๐}}^{}}(^{},^{})`$. The functor $`\mathrm{\Psi }`$ is locally an equivalence by the results of , and hence is an equivalence. โ
###### Example 6.5.
Let $`๐=P^{}`$, $`\mathrm{\Lambda }=P_{\{0\}}^{}=\{\mathrm{pt}\}`$, $`\gamma ^1\mathrm{\Lambda }=^\times `$. In this case, simple objects of rank one are classified by $`^\times /`$.
## 7 WKB-modules
The relationship between microdifferential operators on a complex contact manifold and WKB-differential operators on a complex symplectic manifold is classic, and is discussed e.g., in in the case of cotangent bundles. This study (including the analysis of the action of quantized contact transformations) is systematically performed in , and here we follow their presentation.
Let $`X`$ be a complex manifold, $`t`$ the coordinate, and set
$$_{X\times ,\widehat{t}}=_{X\times ,_t}=\{P_{X\times };[P,/_t]=0\}.$$
Set $`\dot{P}^{}(X\times )=\{\tau 0\}`$, and consider the projection
$$\rho :\dot{P}^{}(X\times )\stackrel{}{}T^{}X,$$
(7.1)
given in local coordinates by $`\rho (x,t;\xi ,\tau )=(x;\xi /\tau )`$. The ring of WKB-operators on $`T^{}X`$ is defined by
$$๐ฒ_X:=\rho _{}(_{X\times ,\widehat{t}}).$$
We similarly set $`๐ฒ_X^\sqrt{v}:=\rho _{}(_{X\times ,\widehat{t}}^\sqrt{v})`$. In a local symplectic coordinate system $`(x;\xi )`$ on $`T^{}X`$, a section $`P๐ฒ_X(U)`$ is written as a formal series
$$P=\underset{jm}{}p_j(x;\xi )\tau ^j,p_j๐ช_{T^{}X}(U),m,$$
with the condition that for any compact subset $`K`$ of $`V`$ there exists a constant $`C_K>0`$ such that $`\underset{K}{sup}|p_j|C_K^j(j)!`$ for all $`j<0`$.
One sets
$$๐ค:=๐ฒ_{\{\mathrm{pt}\}}.$$
(7.2)
Hence, an element $`a๐ค`$ is written as a formal series
$$a=\underset{jm}{}a_j\tau ^j,a_j,m,$$
with the condition that there exist $`C>0`$ with $`|a_j|C^j(j)!`$ for all $`j<0`$.
Note that $`๐ฒ_X`$ is $``$-filtered $`๐ค`$-central algebra, and the principal symbol map $`\sigma _m:๐ฒ_X(m)\stackrel{}{}๐ช_{T^{}X}`$ induces an isomorphism of graded algebras $`gr๐ฒ_X\stackrel{}{}๐ช_{T^{}X}[\tau ^1,\tau ]`$. Note also that $`\pi ^1๐_X`$ is a subring of $`๐ฒ_X`$.
We can now mimic Definition 4.1 for $`๐ฒ_X`$-modules.
###### Definition 7.1.
Let $`\mathrm{\Lambda }`$ be a smooth Lagrangian submanifold of $`T^{}X`$. Let $``$ be a coherent $`๐ฒ_X`$-module supported by $`\mathrm{\Lambda }`$. One says that $``$ is regular (resp. simple) along $`\mathrm{\Lambda }`$ if there locally exists a coherent sub-$`๐ฒ_X(0)`$-module $`_0`$ of $``$ which generates it over $`๐ฒ_X`$, and such that $`_0/๐ฒ_X(1)_0`$ is an $`๐ช_\mathrm{\Lambda }`$-module (resp. a locally free $`๐ช_\mathrm{\Lambda }`$-module of rank one).
Note that, as follows e.g. from Corollary 9.2 below, if $``$ is regular then it is locally a finite direct sum of simple modules.
###### Notation 7.2.
Let $`X`$ be a complex manifold. We denote by $`๐ช_X^\tau `$ the simple $`๐ฒ_X`$-module along the zero-section $`T_X^{}X`$ defined by $`๐ช_X^\tau =๐ฒ_X/`$, where $``$ is the left ideal generated by the vector fields on $`X`$.
###### Proposition 7.3.
1. Any two $`๐ฒ_X`$-modules simple along $`\mathrm{\Lambda }`$ are locally isomorphic. In particular, any simple module along $`T_X^{}X`$ is locally isomorphic to $`๐ช_X^\tau `$.
2. If $``$, $`๐ฉ`$ are simple $`๐ฒ_X`$-modules along $`\mathrm{\Lambda }`$, then $`Rom_{๐ฒ_X}(,๐ฉ)`$ is a $`๐ค`$-local system of rank one on $`\mathrm{\Lambda }`$.
* Since both statements are local on $`\mathrm{\Lambda }`$, we may assume that $`X`$ is endowed with a local coordinate system $`(x_1,\mathrm{},x_n)`$ and that $`\mathrm{\Lambda }`$ is the zero-section $`T_X^{}X`$ of $`T^{}X`$.
(i) Any $`๐ฒ_X`$-module simple along $`T_X^{}X`$ is locally isomorphic to $`๐ช_X^\tau `$. Indeed, the proof of the theorem (due to ) which asserts that if $`Y`$ is a complex manifold, then simple $`_Y`$-modules along smooth regular involutive submanifolds of $`P^{}Y`$ are locally isomorphic applies when replacing $`_{X\times }`$ with $`_{X\times ,\widehat{t}}`$ (for an exposition, see \[21, Ch 1, Th 6.2.1\]).
(ii) By (i) we may assume that $`=๐ฉ=๐ช_X^\tau `$. The result easily follows, representing $`Rom_{๐ฒ_X}(๐ช_X^\tau ,๐ช_X^\tau )`$ by the Koszul complex $`K^{}(๐ช_X^\tau ,(_{x_1},\mathrm{},_{x_n}))`$ associated with the sequence $`(_{x_1},\mathrm{},_{x_n})`$ acting on $`๐ช_X^\tau `$. โ
Recall the projection $`\rho `$ in (7.1) and note that $`\rho ^1\mathrm{\Lambda }`$ is an involutive submanifold of $`\dot{P}^{}(X\times )`$.
###### Proposition 7.4.
Let $`\mathrm{\Lambda }`$ be a smooth Lagrangian submanifold of $`T^{}X`$.
1. Locally, there exists a Lagrangian submanifold $`\mathrm{\Lambda }^0\rho ^1(\mathrm{\Lambda })`$ on which $`\rho `$ induces an isomorphism $`\mathrm{\Lambda }^0\mathrm{\Lambda }`$.
2. Let $``$ be a simple $`๐ฒ_X`$-module along $`\mathrm{\Lambda }`$. Then, locally there exists a simple $`_{X\times }`$-module $`^0`$ along $`\mathrm{\Lambda }^0`$ such that $`\mathrm{๐๐๐}(^0)`$, where $`\mathrm{๐๐๐}`$ denotes the forgetful functor $`\mathrm{Mod}(_{X\times })\stackrel{}{}\mathrm{Mod}(_{X\times ,\widehat{t}})`$.
* Since the problem is local, we may assume that $`\mathrm{\Lambda }=T_X^{}X`$ in a system of local symplectic coordinates $`(x;u)T^{}X`$. In the corresponding system of homogeneous coordinates $`(x,t;\xi ,\tau )P^{}(X\times )`$, one has $`\rho ^1\mathrm{\Lambda }=\{\xi =0\}`$. This set is foliated by the Lagrangian submanifolds $`\mathrm{\Lambda }^c=\{\xi =0,t=c\}`$, for $`c`$. For $`ZX`$ a closed submanifold, denote by $`๐_{Z|X}`$ the sheaf of finite order holomorphic microfunctions on $`P_Z^{}X`$. A simple $`_X`$-module along $`\mathrm{\Lambda }^0`$ is given by $`^0=๐_{X\times \{0\}|X\times }`$. One then immediately checks that $`^0๐ช_X^\tau `$. โ
## 8 Contactification of symplectic manifolds
In this section, we recall some well-known facts from the specialists on contact and symplectic geometry.
###### Definition 8.1.
A contactification of a symplectic manifold $`(๐,\omega )`$ is a complex contact manifold $`(๐,๐ช_๐(1),\alpha )`$, and a morphism of complex manifolds $`\rho :๐\stackrel{}{}๐`$ such that the following conditions are satisfied:
1. the line bundle $`๐ช_๐(1)`$ has a nowhere vanishing section $`\tau `$,
2. there exists an open covering $`๐=_{iI}U_i`$, holomorphic functions $`t_i`$ on $`\rho ^1(U_i)`$ and primitives $`\sigma _i\mathrm{\Omega }_{U_i}^1`$ of $`\omega |_{U_i}`$ such that $`\chi _i:=(\rho ,t_i)`$ gives an isomorphism $`\chi _i:\rho ^1(U_i)\stackrel{}{}U_i\times `$, and $`dt_i+\rho ^{}(\sigma _i)=(\alpha /\tau )|_{\rho ^1(U_i)}`$.
###### Lemma 8.2.
A contactification $`\rho :๐\stackrel{}{}๐`$ has a structure of $``$-principal bundle.
* With notations as in Definition 8.1, set $`\chi _{ij}=\chi _j\chi _i^1|_{U_{ij}}:U_{ij}\times \stackrel{}{}U_{ij}\times `$. We have $`\chi _{ij}=(\mathrm{id}_{U_{ij}},\theta _{ij})`$ for some transition functions $`\theta _{ij}:U_{ij}\times \stackrel{}{}`$. Since $`d(t_it_j)=\rho ^{}(\sigma _i\sigma _j)`$, up to shrinking the covering we may assume that $`\sigma _i\sigma _j=df_{ij}`$ for some functions $`f_{ij}๐ช_{U_{ij}}`$. The transition functions are the translations $`\theta _{ij}(p,t)=t+f_{ij}(p)+c_{ij}`$ for some $`c_{ij}`$. โ
###### Lemma 8.3.
A symplectic manifold $`(๐,\omega )`$ admits a contactification if and only if the de Rham cohomology class $`[\omega ]H^2(๐;_๐)`$ vanishes.
* (a) Let $`๐=_{iI}U_i`$ be a covering such that $`\omega |_{U_i}`$ has a primitive $`\sigma _i\mathrm{\Omega }_{U_i}^1`$. Up to shrinking the covering we may assume that $`\sigma _i\sigma _j=df_{ij}`$ for some functions $`f_{ij}๐ช_{U_{ij}}`$. Then $`d(f_{ij}+f_{jk}f_{ik})=0`$, and $`_{U_{ijk}}c_{ijk}=f_{ij}+f_{jk}f_{ik}`$ is a Cech cocycle representing $`[\omega ]H^2(๐;_๐)`$. Since $`[\omega ]=0`$, we have $`c_{ijk}=c_{ij}+c_{jk}c_{ik}`$. Setting $`\stackrel{~}{f}_{ij}=f_{ij}c_{ij}`$, we have $`\stackrel{~}{f}_{ij}+\stackrel{~}{f}_{jk}\stackrel{~}{f}_{ik}=0`$. Endow $`U_i\times `$ with the contact form $`dt+\sigma _i`$. A contactification of $`๐`$ is thus given by the $``$-principal bundle $`๐\stackrel{}{}๐`$ with local charts $`U_i\times `$, and transition functions over $`U_{ij}`$ given by $`(p,t)(p,t+\stackrel{~}{f}_{ij}(p))`$.
(b) Let $`\rho :๐\stackrel{}{}๐`$ be a contactification, and use notations as in Definition 8.1. We have $`d(\alpha /\tau )=\rho ^{}\omega `$ and hence $`[\rho ^{}\omega ]=0`$ in $`H^2(๐;_๐)`$. Then $`[\omega ]=0`$ in $`H^2(๐;_๐)\stackrel{}{}H^2(๐;_๐)`$. โ
###### Lemma 8.4.
Let $`(๐,\omega )`$ be a symplectic manifold and $`\mathrm{\Lambda }`$ a Lagrangian submanifold. After replacing $`๐`$ with a neighborhood of $`\mathrm{\Lambda }`$ there exists a contactification $`\rho :๐\stackrel{}{}๐`$ and a Lagrangian submanifold $`\mathrm{\Lambda }^0๐`$ on which $`\rho `$ induces an isomorphism $`\mathrm{\Lambda }^0\mathrm{\Lambda }`$.
* The restriction map $`H^2(๐,_๐)\stackrel{}{}H^2(\mathrm{\Lambda };_\mathrm{\Lambda })`$ is given by $`[\omega ][j^{}\omega ]`$, where $`j:\mathrm{\Lambda }๐`$ is the embedding. Since $`\omega |_\mathrm{\Lambda }=0`$, there exists an open neighborhood $`U\mathrm{\Lambda }`$ such that $`[\omega |_U]=0`$. Thus, up to replacing $`๐`$ with $`U`$, we can assume that $`[\omega ]=0`$.
Let us adapt the arguments in part (a) of the proof of Lemma 8.3 above. Let $`๐=_{iI}U_i`$ be a covering such that $`\omega |_{U_i}`$ has a primitive $`\sigma _i\mathrm{\Omega }_{U_i}^1`$. We may assume that $`\sigma _i|_\mathrm{\Lambda }=0`$. Let $`f_{ij}๐ช_{U_{ij}}`$ satisfy $`\sigma _i\sigma _j=df_{ij}`$. We may assume moreover that $`f_{ij}|_\mathrm{\Lambda }=0`$. Since $`f_{ij}+f_{jk}f_{ik}`$ is locally constant on $`\mathrm{\Lambda }`$, it must vanish on $`\mathrm{\Lambda }`$. Thus, the $``$-principal bundle $`\rho :๐\stackrel{}{}๐`$ is described by the local charts $`U_i\times `$, with transition functions $`(p,t)(p,t+f_{ij}(p))`$ and contact form $`dt+\sigma _i`$. We define $`\mathrm{\Lambda }^0`$ by $`\mathrm{\Lambda }^0|_{U_i\times }=\mathrm{\Lambda }|_{U_i}\times \{0\}`$. โ
## 9 Simple holonomic modules on symplectic manifolds
Let us start by recalling the construction of of the stack of WKB-modules on a symplectic manifold, in the special case where there exists a contactification.
Let $`\rho :๐\stackrel{}{}๐`$ be a contactification of a complex symplectic manifold. Denote by
$$๐ค_{๐,\tau }=(\{๐_i\}_{iI},\{f_{ij}\}_{i,jI},\{a_{ijk}\}_{i,j,kI})$$
the $``$-algebroid descent datum on $`๐`$ given by Proposition 5.5. Consider the $``$-algebroid descent datum on $`๐`$
$$๐ถ_๐=(\{\rho _{}๐_i\}_{iI},\{\rho _{}f_{ij}\}_{i,jI},\{\rho _{}a_{ijk}\}_{i,j,kI}).$$
The stack $`\mathrm{๐ฌ๐๐ฝ}(๐ถ_๐)`$ associated with this $``$-algebroid descent datum is the stack of WKB-modules in .
One defines the notion of good $`๐ถ_๐`$-module similarly as for $`๐ค_๐`$-modules.
Here, we prove the existence of twisted simple holonomic modules along smooth Lagrangian submanifolds of complex symplectic manifolds.
###### Theorem 9.1.
Let $`๐`$ be a complex symplectic manifold and let $`\mathrm{\Lambda }๐`$ be a Lagrangian submanifold. There exists a module $`\mathrm{Mod}(๐ถ_๐|_\mathrm{\Lambda }_{}_{\sqrt{\mathrm{\Omega }_\mathrm{\Lambda }}})`$ which is simple along $`\mathrm{\Lambda }`$. Moreover, $``$ is good.
* By Lemma 8.3 there exist a contactification $`\rho :๐\stackrel{}{}๐`$ and a Lagrangian submanifold $`\mathrm{\Lambda }^0๐`$ on which $`\rho `$ induces an isomorphism $`\mathrm{\Lambda }^0\mathrm{\Lambda }`$. By Theorem 6.1 there exists $`^0\mathrm{Mod}(๐ค_๐|_{\mathrm{\Lambda }^0}_{}_{\sqrt{\mathrm{\Omega }_{\mathrm{\Lambda }^0}}})`$ which is simple along $`\mathrm{\Lambda }^0`$. As in Proposition 7.4 we then set $`=\mathrm{๐๐๐}(^0)`$, where
$`\mathrm{๐๐๐}:\mathrm{Mod}(๐ค_๐|_{\mathrm{\Lambda }^0}_{}_{\sqrt{\mathrm{\Omega }_{\mathrm{\Lambda }^0}}})`$ $`\stackrel{}{}\mathrm{Mod}(๐ค_{๐,\widehat{t}}|_{\mathrm{\Lambda }^0}_{}_{\sqrt{\mathrm{\Omega }_{\mathrm{\Lambda }^0}}})`$
$`\mathrm{Mod}(๐ถ_๐|_\mathrm{\Lambda }_{}_{\sqrt{\mathrm{\Omega }_\mathrm{\Lambda }}})`$
is the natural forgetful functor. โ
###### Corollary 9.2.
There is a $`๐ค`$-equivalence of stacks
$$\mathrm{๐ฌ๐๐ฝ}_{\text{reg-}\mathrm{\Lambda }}(๐ถ_๐|_\mathrm{\Lambda }_{}_{\sqrt{\mathrm{\Omega }_\mathrm{\Lambda }}})\mathrm{๐ฌ๐๐ฝ}_{\text{loc-sys}}(๐ค_\mathrm{\Lambda }).$$
(9.1)
* By Theorem 9.1, there exists a simple module $``$ in $`\mathrm{Mod}(๐ถ_๐|_\mathrm{\Lambda }_{}_{\sqrt{\mathrm{\Omega }_\mathrm{\Lambda }}})`$. By Proposition 7.3, a functor (9.1) is given by $`om_{๐ถ_๐}(,)`$. Proving that it is an equivalence is a local problem, and so we may assume $`๐=T^{}X`$, $`\mathrm{\Lambda }=T_X^{}X`$. Then a quasi inverse is given by $`FF_๐ค๐ช_X^\tau `$. โ
###### Remark 9.3.
Note that in the real setting, a link between simple holonomic modules on real Lagrangian submanifolds and Fourier distributions, Maslov index, etc. is investigated in .
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# Notes on ๐ถ-graded modules over an affine semigroup ring ๐พโข[๐ถ]
## 1. Introduction
First, we fix the notation used throughout this paper. Let $`C^d^d`$ be an affine semigroup (i.e., $`C`$ is a finitely generated additive submonoid of $`^d`$), and $`R:=K[๐ฑ^๐๐C]K[x_1^{\pm 1},\mathrm{},x_d^{\pm 1}]`$ the semigroup ring of $`C`$ over a field $`K`$. Here $`๐ฑ^๐`$ for $`๐=(c_1,\mathrm{},c_d)C`$ denotes the monomial $`_{i=1}^dx_i^{c_i}`$. We always assume that $`C=^d`$ and $`C(C)=\{0\}`$. Thus $`dimR=d`$ and $`๐ช:=(๐ฑ^๐0๐C)`$ is a maximal ideal of $`R`$.
Of course, $`R=_{๐C}K๐ฑ^๐`$ is a $`^d`$-graded ring. We say a $`^d`$-graded ideal of $`R`$ is a monomial ideal. Let $`\mathrm{Mod}R`$ be the category of $`^d`$-graded $`R`$-modules and their degree preserving $`R`$-homomorphisms, and $`\mathrm{mod}R`$ its full subcategory consisting of finitely generated modules. As usual, for $`M\mathrm{Mod}R`$ and $`๐^d`$, $`M_๐`$ denotes the degree $`๐`$ component of $`M`$, and $`M(๐)`$ denotes the shifted module of $`M`$ with $`M(๐)_๐=M_{๐+๐}`$. We say $`M\mathrm{Mod}R`$ is $`C`$-graded, if $`M_๐=0`$ for all $`๐C`$. A monomial ideal $`IR`$ and the quotient ring $`R/I`$ are $`C`$-graded modules. Let $`\mathrm{mod}_CR`$ be the full subcategory of $`\mathrm{mod}R`$ consisting of $`C`$-graded modules.
Miller proved that $`\mathrm{mod}_CR`$ has enough injectives and any object has a minimal injective resolution in this category, which is unique up to isomorphism and has finite length. This resolution is called a minimal irreducible resolution, since an indecomposable injective in $`\mathrm{mod}_CR`$ corresponds to a monomial irreducible ideal.
In ยง2, under the assumption that $`R`$ is Cohen-Macaulay and simplicial, we show that information on $`M\mathrm{mod}_CR`$ such as depth and Cohen-Macaulay property can be read off from numerical invariants of the minimal irreducible resolution of $`M`$ (something analogous to โBass numbersโ). One might think these results should be a variant of the fact that depth and related conditions can be characterized by Bass numbers. But this insight is not quite correct. Philosophically, our result is rather closer to a theorem of Eagon and Reiner stating that the Stanley-Reisner ring of a simplicial complex $`\mathrm{\Delta }`$ is Cohen-Macaulay if and only if that of the Alexander dual $`\mathrm{\Delta }^{}`$ has a linear free resolution, and Millerโs generalization of this result to finitely generated $`^d`$-graded modules over a polynomial ring $`K[x_1,\mathrm{},x_d]`$ (see ). In fact, in the polynomial ring case, the โAlexander dualโ of our Theorem 2.6 corresponds to his \[13, Theorem 4.20\]. But the proofs are not similar.
In ยงยง3-6, we assume that $`R`$ is normal (but not necessarily simplicial). In ยง3, we study the full subcategory $`\mathrm{Sq}`$ of $`\mathrm{mod}_CR`$ consisting of squarefree modules. The notion of squarefree modules over a normal semigroup ring was introduced by the author in . A monomial ideal $`IR`$ is squarefree if and only if it is a radical ideal (i.e., $`I=\sqrt{I}`$). The category $`\mathrm{Sq}`$ behaves much nicer than $`\mathrm{mod}_CR`$ as shown in . In this paper, we show that a squarefree module $`M`$ is sequentially Cohen-Macaulay (a non-pure generalization of the Cohen-Macaulay property) if and only if the linear strand of its minimal irreducible resolution is acyclic.
In ยง4, assuming that $`R`$ is normal, we study the quotient ring $`R/I`$ by a radical monomial ideal $`I`$, focusing on the problem when $`R/I`$ is sequentially Cohen-Macaulay. When $`R`$ is a polynomial ring, $`R/I`$ is usually called a Stanley-Reisner ring. As a Stanley-Reisner ring is associated with a simplicial complex, our $`R/I`$ is associated with a polyhedral complex $`\mathrm{\Delta }`$ contained in the polyhedral cone $`_0C^d`$. So we denote it by $`K[\mathrm{\Delta }]`$. Our $`K[\mathrm{\Delta }]`$ is a special case of the rings Stanley constructed from more general polyhedral complexes in \[18, ยง4\], but still an interesting class. On the other hand, the sequentially Cohen-Macaulay property has become important in the theory of Stanley-Reisner rings, since it is closely related to non-pure shellability and shifting of simplicial complexes (c.f. ). Among other things, we show that the sequentially Cohen-Macaulay property of $`K[\mathrm{\Delta }]`$ is a topological property of the โgeometric realizationโ $`|\mathrm{\Delta }|B`$ of the complex $`\mathrm{\Delta }`$. Here $`B`$ is a $`(d1)`$-dimensional polytope which is the intersection of the cone $`_0C`$ and a hyperplane $`H^d`$.
With the above notation, a squarefree module $`M`$ gives a $`K`$-constructible sheaf $`M^+`$ on $`B`$. For example, $`K[\mathrm{\Delta }]^+`$ is the $`K`$-constant sheaf on the geometric realization $`|\mathrm{\Delta }|B`$. As shown in , we can connect the local cohomology of $`M`$ and the sheaf cohomology of $`M^+`$. More precisely, we have $`[H_๐ช^{i+1}(M)]_0H^i(B;M^+)`$ for all $`i1`$ and an exact sequence $`0[H_๐ช^0(M)]_0M_0H^0(B;M^+)[H_๐ช^1(M)]_00`$. We say $`M\mathrm{Sq}`$ is regular, if $`[H_๐ช^0(M)]_0=[H_๐ช^1(M)]_0=0`$. For all $`M\mathrm{Sq}`$, there is a unique $`\overline{M}\mathrm{Sq}`$ which is regular and $`\overline{M}^+M^+`$. In ยง5, we show that if $`M`$ is regular then the (sequentially) Cohen-Macaulay property of $`M`$ depends only on the sheaf $`M^+`$. This result can imply the main result of ยง4.
Assume that $`R`$ is normal and $`IR`$ is a (not necessarily monomial) ideal. In ยง6, generalizing a result of Herzog, Takayama and Terai , we show that if $`R/I`$ is Cohen-Macaulay and $`\sqrt{I}`$ is a monomial ideal, then $`R/\sqrt{I}`$ is Cohen-Macaulay again. (If $`\sqrt{I}`$ is not a monomial ideal, this statement does not hold. Otherwise, if $`IS:=K[x_1,\mathrm{},x_d]`$ defines a set theoretic complete intersection subscheme of $`๐ธ^d`$, then $`S/\sqrt{I}`$ must be Cohen-Macaulay. This is clearly strange.)
Let $`S=K[x_1,\mathrm{},x_d]`$ be a polynomial ring, and $`IS`$ a monomial ideal. Recently, Takayama showed that the range $`\{๐^d[H_๐ช^i(S/I)]_๐0\}`$ is controlled by the degrees of minimal generators of $`I`$ (especially, when $`H_๐ช^i(S/I)`$ has finite length). In ยง6, after giving a simple new proof of this result, we will (partially) extend it to affine semigroup rings. For example, we prove the following: Assume that $`R`$ is normal and simplicial. For $`M\mathrm{mod}_CR`$, $`H_๐ช^i(M)`$ has finite length if and only if $`H_๐ช^i(M)`$ is $`C`$-graded. A similar result holds for arbitrary affine semigroup rings, but some modification is necessary.
## 2. Irreducible resolutions over Cohen-Macaulay simplicial semigroup rings
We use the same notation as in the introduction. Let $`_0:=\{rr0\}`$ be the set of non-negative real numbers, and $`๐:=\{\gamma _i๐_i\gamma _i_0,๐_iC\}^d`$ the polyhedral cone spanned by $`C`$. Let $`๐`$ be the set of (non-empty) faces of $`๐`$. Note that $`๐`$ itself and $`\{0\}`$ belong to $`๐`$. For $`F๐`$, $`P_F`$ denotes the monomial ideal $`(๐ฑ^๐๐CF)`$ of $`R`$. Then $`P_F`$ is a prime ideal. Conversely, any monomial prime ideal of $`R`$ is of the form $`P_F`$ for some $`F๐`$.
It is well-known that $`\mathrm{Mod}R`$ has enough injectives (c.f. \[15, Chapter 11\]). We denote the injective hull of $`R/P_F`$ in $`\mathrm{Mod}R`$ by $`{}_{}{}^{}E(R/P_F)`$. Assume that $`M\mathrm{Mod}R`$ is indecomposable. Then $`M`$ is injective in $`\mathrm{Mod}R`$ if and only if there are some $`๐^d`$ and $`F๐`$ such that $`M{}_{}{}^{}E(R/P_F)(๐)`$. Recall that $`M\mathrm{Mod}R`$ has a minimal injective resolution in $`\mathrm{Mod}R`$, which is unique up to isomorphism. For a monomial prime ideal $`P_F`$, set $`\mu _i(P_F,M):=dim_{K(F)}(\mathrm{Ext}_R^i(R/P_F,M)R_{P_F})`$, where $`K(F)`$ is the quotient field of $`R/P_F`$. It is well-known that if $`J^{}`$ is a minimal injective resolution of $`M`$ in $`\mathrm{Mod}R`$, then we have $`J^i{}_{}{}^{}E(R/P_F)^{\mu _i(P_F,M)}`$ for all $`i`$ as underlying $`R`$-modules (i.e., if we forget the grading of the modules).
We say an ideal $`IR`$ is irreducible if every expression $`I=I_1I_2`$ of $`I`$ as an intersection of two ideals satisfies that $`I=I_1`$ or $`I=I_2`$. An irreducible ideal is always a primary ideal, while the converse is not true. In particular, the radical $`\sqrt{I}`$ of an irreducible ideal $`I`$ is a prime ideal. According to Miller , we use the letter $`W`$ to denote a monomial irreducible ideal. Then we have $`\sqrt{W}=P_F`$ for some $`F๐`$. In this case, $`dim(R/W)`$ equals the dimension of $`F`$ as a polyhedral cone. When $`R`$ is a polynomial ring $`K[x_1,\mathrm{},x_d]`$, $`I`$ is a monomial irreducible ideal if and only if it is a complete intersection ideal generated by powers of variables.
For $`M\mathrm{Mod}R`$, we say the submodule $`_{๐C}M_๐`$ is the $`C`$-graded part of $`M`$, and we denote it by $`M_C`$.
###### Proposition 2.1 (Miller \[14, ยง2\]).
We have the following.
* Let $`IR`$ be a monomial ideal. Then $`I`$ is irreducible if and only if there are some $`๐C`$ and $`F๐`$ such that $`R/I[{}_{}{}^{}E(R/P_F)(๐)]_C`$.
* The category $`\mathrm{mod}_CR`$ has enough injectives. An indecomposable module $`M`$ is injective in $`\mathrm{mod}_CR`$ if and only if $`MR/W`$ for some monomial irreducible ideal $`WR`$.
We say an injective resolution in $`\mathrm{mod}_CR`$ is an irreducible resolution.
###### Theorem 2.2 (Miller \[14, Theorem 2.4\]).
Let $`M\mathrm{mod}_CR`$ be a $`C`$-graded module, and $`J^{}`$ a minimal injective resolution of $`M`$ in $`\mathrm{Mod}R`$. Then the $`C`$-graded part $`[J^{}]_C`$ of $`J^{}`$ is an irreducible resolution of $`M`$, and has finite length.
We say the irreducible resolution given in Theorem 2.2 is a minimal irreducible resolution. Every $`M\mathrm{mod}_CR`$ has a minimal irreducible resolution, and this is unique up to isomorphism. Any irreducible resolution is a direct sum of a minimal one and an exact sequence. It is noteworthy that Helm and Miller gave an algorithm to compute minimal irreducible resolutions.
Let $`I^{}`$ be a minimal irreducible resolution of $`M\mathrm{mod}_CR`$. Then, for a monomial irreducible ideal $`W`$ and an integer $`i0`$, we have $`\nu _i(W,M)`$ satisfying
$$I^i(R/W)^{\nu _i(W,M)}.$$
Note that
(2.1)
$$\mu _i(P_F,M)\underset{\sqrt{W}=P_F}{}\nu _i(W,M)$$
for all $`F๐`$. Set $`\mathrm{ir}\mathrm{dim}M:=\mathrm{max}\{iI^i0\}`$. Theorem 2.2 states that $`\mathrm{ir}\mathrm{dim}M<\mathrm{}`$ for all $`M`$. But, if $`R`$ is not a polynomial ring, then $`sup\{\mathrm{ir}\mathrm{dim}MM\mathrm{mod}_CR\}=\mathrm{}`$. In fact, for a given integer $`n`$, we have $`\mathrm{ir}\mathrm{dim}((R/๐ช)(๐))>n`$ for sufficiently โlargeโ $`๐C`$.
###### Lemma 2.3.
Let $`M\mathrm{mod}_CR`$. If $`\nu _0(W,M)0`$, then $`\sqrt{W}`$ is an associated prime of $`M`$. Conversely, if $`P_F`$ is an associated prime of $`M`$, then there is a monomial irreducible ideal $`W`$ with $`\sqrt{W}=P_F`$ and $`\nu _0(W,M)0`$. In particular, $`dimM=\mathrm{max}\{dim(R/W)\nu _0(W,M)0\}`$.
###### Proof.
The first assertion follows from (2.1). Since $`M`$ is a submodule of $`I^0:=(R/W^{})^{v_0(W^{},M)}`$ and an irreducible ideal is a primary ideal, we have $`\mathrm{Ass}(M)\mathrm{Ass}(I^0)=\{\sqrt{W^{}}v_0(W^{},M)0\}`$. So the second statement follows. The last assertion is easy. โ
For further information on irreducible resolutions, consult \[15, Chapter 11\].
We say $`R`$ is simplicial if there are $`๐_1,\mathrm{},๐_dC`$ such that $`๐=\{_{i=1}^d\gamma _i๐_i\gamma _i_0\}`$. In this case, we have $`C=^d\{_{i=1}^d\alpha _i๐_i\alpha _i_0\}`$, and $`\{๐_1,\mathrm{},๐_d\}`$ is a basis of $`^d`$ (and $`^d`$). A polynomial ring $`K[x_1,\mathrm{},x_d]`$ is a primary example of a simplicial semigroup ring.
###### Lemma 2.4.
Assume that $`R`$ is Cohen-Macaulay and simplicial. Then $`R/W`$ is Cohen-Macaulay for all monomial irreducible ideal $`W`$.
###### Proof.
By Proposition 2.1 (1), there is some $`๐=_{i=1}^d\alpha _i๐_iC`$ ($`\alpha _i_0`$) such that $`[{}_{}{}^{}E(R/P_F)(๐)]_CR/W,`$ where $`F`$ is the face of $`๐`$ with $`P_F=\sqrt{W}`$. We may assume that $`F`$ is spanned by $`๐_1,๐_2,\mathrm{},๐_n`$ as a polyhedral cone. Set $`H:=CF`$ to be a submonoid of $`C`$, and set
$$U:=\{\underset{i=1}{\overset{d}{}}\beta _i๐_iC\beta _i0\text{ for all }i\text{ and }\beta _i\alpha _i\text{ for all }i>n\}.$$
Then $`U`$ is an $`H`$-ideal in the sense of , that is, $`H+UU`$. (In , they assumed that $`U`$ is contained in the group generated by $`H`$. But their results hold without this assumption.) Note that $`A:=K[H]=_{๐H}K๐ฑ^๐`$ is a simplicial affine semigroup ring. Clearly, $`AR/P_F`$. But here we regard $`A`$ as a subring of $`R`$ via the inclusion $`HC`$. As an $`A`$ module, $`R/W`$ is isomorphic to $`K[U]`$, and it is finitely generated. Since $`R=K[C]`$ is Cohen-Macaulay, $`U`$ satisfies the condition (iii) of \[5, Theorem 2,2\] as an $`H`$-ideal. Thus $`R/W`$ is a Cohen-Macaulay module over $`A`$ (thus over $`R`$). โ
###### Example 2.5.
(1) In the non-simplicial case, a monomial irreducible ideal need not be Cohen-Macaulay, even if the ring is normal. Consider the affine semigroup $`C`$ in $`^3`$ generated by $`(0,0,1)`$, $`(0,1,1)`$, $`(1,1,1)`$ and $`(1,0,1)`$. Then the semigroup ring $`R`$ of $`C`$ is normal but not simplicial. We denote the monomials in $`R`$ with degrees $`(0,0,1)`$, $`(0,1,1)`$, $`(1,1,1)`$ and $`(1,0,1)`$ by $`x`$, $`y`$, $`z`$ and $`w`$ respectively. Note that $`R=K[x,y,z,w]/(xzyw)`$. Set $`W=(y^2,yz,z^2)`$. Then $`W`$ is an irreducible ideal, since $`R/W[{}_{}{}^{}E(R/(y,z))((1,1,1))]_C`$. But computation by Macaulay 2 shows that $`R/W`$ is not Cohen-Macaulay.
(2) If $`R`$ is normal (but not necessarily simplicial), then $`R/P_F`$ is a normal semigroup ring (in particular, Cohen-Macaulay) for any monomial prime ideal $`P_F`$. If $`R`$ is simplicial and Cohen-Macaulay, then so is $`R/P_F`$ by Lemma 2.4. On the other hand, if $`R`$ is Cohen-Macaulay but not simplicial, then $`R/P_F`$ need not be Cohen-Macaulay. The following example is due to Professor Ngo Viet Trung. The affine semigroup ring $`R:=k[s^4,s^3t,st^3,t^4,s^4u,t^4u]`$ is Cohen-Macaulay. But $`P=(s^4u,t^4u)`$ is a prime ideal and $`R/Pk[s^4,s^3t,st^3,t^4]`$ is not Cohen-Macaulay.
When $`R`$ is Cohen-Macaulay, $`\omega _R`$ denotes the canonical module of $`R`$.
###### Theorem 2.6.
If $`R`$ is Cohen-Macaulay and simplicial, and $`M\mathrm{mod}_CR`$, then we have
(2.2)
$$\mathrm{depth}_RM=\mathrm{min}\{dim(R/W)+i\nu _i(W,M)0\}.$$
###### Proof.
Let $`I^{}:0I^0\stackrel{d^0}{}I^1\stackrel{d^1}{}\mathrm{}`$ be a minimal irreducible resolution of $`M`$, and set $`\mathrm{\Omega }^i(M):=\mathrm{ker}(d^i)`$. Of course, $`\mathrm{\Omega }^0(M)=M`$. Set $`\delta `$ to be the right hand side of (2.2).
To prove $`\mathrm{depth}_RM\delta `$, we will show that $`\mathrm{depth}(\mathrm{\Omega }^i(M))\delta i`$ for all $`i`$ by backward induction on $`i`$. (Here we set the depth of the 0 module to be $`+\mathrm{}`$.) If $`id`$, there is nothing to prove. Assume that $`\mathrm{depth}_R(\mathrm{\Omega }^{i+1}(M))\delta i1`$. Consider the short exact sequence
(2.3)
$$0\mathrm{\Omega }^i(M)I^i\mathrm{\Omega }^{i+1}(M)0.$$
Recall that $`I^i=(R/W)^{\nu _i(W,M)}`$ and each $`R/W`$ is Cohen-Macaulay by Lemma 2.4. By the assumption, if $`\nu _i(W,M)0`$ then $`dim(R/W)\delta i`$. So $`\mathrm{depth}_RI^i\delta i`$. By (2.3), we have $`\mathrm{depth}_R(\mathrm{\Omega }^i(M))\delta i`$.
We can take $`n`$ such that $`\nu _n(W,M)0`$ for some $`W`$ with $`dim(R/W)=\delta n`$. To prove $`\mathrm{depth}_RM=\delta `$, we will show that $`\mathrm{depth}_R(\mathrm{\Omega }^i(M))=\delta i`$ for all $`in`$ by backward induction on $`i`$. Since $`\sqrt{W}`$ is an associated prime of $`\mathrm{\Omega }^n(M)`$ and $`dim(R/\sqrt{W})=\delta n`$, we have $`\mathrm{depth}_R(\mathrm{\Omega }^n(M))\delta n`$. (In fact, we have $`H_๐ช^{\delta n}(\mathrm{\Omega }^n(M))^{}\mathrm{Ext}_R^{d\delta +n}(\mathrm{\Omega }^n(M),\omega _R)0`$ by an argument similar to \[1, Theorem 8.1.1\].) But we have seen that $`\mathrm{depth}_R(\mathrm{\Omega }^n(M))\delta n`$. So $`\mathrm{depth}_R(\mathrm{\Omega }^n(M))=\delta n`$. Assume that $`\mathrm{depth}_R(\mathrm{\Omega }^{i+1}(M))=\delta i1`$. Since $`\mathrm{depth}I^i\delta i`$, we have $`\mathrm{depth}_R(\mathrm{\Omega }^i(M))=\delta i`$ by (2.3). โ
The following is an easy consequence of the theorem.
###### Corollary 2.7.
Assume that $`R`$ is Cohen-Macaulay and simplicial. Then $`M\mathrm{mod}_CR`$ is a Cohen-Macaulay module of dimension $`p`$ if and only if $`\nu _i(W,M)0`$ implies $`pidim(R/W)p`$ for all $`i`$.
We can also characterize Serreโs condition $`(S_n)`$ in our context.
###### Theorem 2.8.
Assume that $`R`$ is Cohen-Macaulay and simplicial. If $`IR`$ is a monomial ideal with $`dim(R/I)=p`$, then the following are equivalent for an integer $`n2`$.
* $`R/I`$ satisfies Serreโs condition $`(S_n)`$.
* If $`in1`$ and $`dim(R/W)<pi`$, then $`\nu _i(W,R/I)=0`$.
To prove the theorem, we need the following (more or less) well-known facts. Note that these facts hold in much wider context (e.g., over a noetherian local ring admitting a dualizing complex).
###### Lemma 2.9.
For a monomial ideal $`IR`$ and a module $`M\mathrm{mod}R`$, we have the following.
* We always have $`dim(\mathrm{Ext}_R^j(M,\omega _R))dj`$. And the equality holds if and only if there is an associated prime $`P`$ of $`M`$ with $`dim(R/P)=dj`$.
* If $`R/I`$ satisfies the $`(S_n)`$ condition for some $`n2`$, then all associated primes of $`R/I`$ have the same dimension.
* Let $`n2`$. Then $`R/I`$ satisfies $`(S_n)`$ if and only if $`dim(\mathrm{Ext}_R^j(R/I,\omega _R))djn`$ for all $`j>ddim(R/I)`$.
###### Proof.
(1) Essentially same as \[1, Theorem 8.1.1\].
(2) See \[7, Remark 2.4.1\].
(3) Follows from (2) and the local duality. โ
The parts (2) and (3) of Lemma 2.9 do not hold for a module $`M`$. So Theorem 2.8 only concerns ideals. But if all minimal primes of $`M\mathrm{mod}_CR`$ have the same dimension, then Theorem 2.8 holds for $`M`$.
Proof of Theorem 2.8. Set $`A:=R/I`$. As before, let $`I^{}:0I^0\stackrel{d^0}{}I^1\stackrel{d^1}{}\mathrm{}`$ be a minimal irreducible resolution of $`A`$, and set $`\mathrm{\Omega }^i(A):=\mathrm{ker}(d^i)`$.
Assume that $`A`$ satisfies the condition (b). Then $`\mathrm{Ext}_R^j(I^i,\omega _R)=0`$ for all $`i<n`$ and $`j>dp+i`$ by Lemma 2.4. The short exact sequence (2.3) yields
$$\mathrm{Ext}_R^j(\mathrm{\Omega }^i(A),\omega _R)\mathrm{Ext}_R^{j+1}(\mathrm{\Omega }^{i+1}(A),\omega _R)$$
for all $`i<n1`$ and $`j>dp+i`$. Using this, we get
$$\mathrm{Ext}_R^j(A,\omega _R)\mathrm{Ext}_R^{j+n1}(\mathrm{\Omega }^{n1}(A),\omega _R)$$
for all $`j>dp`$. In this case, the dimension of any associated prime of $`\mathrm{\Omega }^{n1}(A)`$ is at least $`pn+1`$ by (b), and we have $`dim(\mathrm{Ext}_R^{j+n1}(\mathrm{\Omega }^{n1}(A),\omega _R))djn`$ by Lemma 2.9 (1) (note that $`j+n1>dp+n1`$). Thus $`A`$ satisfies the $`(S_n)`$ condition by Lemma 2.9 (3).
Conversely, assume that $`A`$ satisfies $`(S_n)`$ for some $`n2`$ but $`\nu _i(W,A)0`$ for some $`in1`$ and some $`W`$ with $`dim(R/W)pi1`$. Set $`P=\sqrt{W}`$. Then $`dimA_Pp(pi1)=i+1`$. But, we have $`\mathrm{Ext}_R^i(R/P,A)0`$ by (2.1). Hence $`\mathrm{depth}_{R_P}A_Pi<\mathrm{min}\{n,dimA_P\}`$. It contradicts $`(S_n)`$. โ
###### Definition 2.10.
We say $`M\mathrm{mod}R`$ is a generalized Cohen-Macaulay module if $`H_๐ช^i(M)`$ has finite length (i.e., $`dim_KH_๐ช^i(M)<\mathrm{}`$) for all $`i<dimM`$. We can also define this concept over a noetherian local ring in a similar way.
By Lemma 2.9 (1), all minimal primes of a generalized Cohen-Macaulay module have the same dimension.
###### Proposition 2.11.
Assume that $`R`$ is Cohen-Macaulay and simplicial. For $`M\mathrm{mod}_CR`$ with $`dimM=p`$, $`M`$ is a generalized Cohen-Macaulay module if and only if $`\nu _i(W,M)=0`$ for all $`i`$ and all $`W`$ with $`0<dim(R/W)<pi`$.
###### Proof.
(Sufficiency) By argument similar to the proof of Theorem 2.6, $`H_๐ช^j(\mathrm{\Omega }^i(M))`$ has finite length for all $`j<pi`$. Since $`\mathrm{\Omega }^0(M)=M`$, we are done.
(Necessity) We will prove the contrapositive. Assume that $`\nu _n(W,M)0`$ for some $`n0`$ and some $`W`$ with $`0<dim(R/W)<pn`$. We may assume that $`n`$ is minimum among such integers, and set $`q:=dim(R/W)`$. We can prove that $`H_๐ช^{q+ni}(\mathrm{\Omega }^i(M))`$ does not have finite length for all $`in`$ by backward induction. Since $`\mathrm{\Omega }^0(M)=M`$ and $`q+n<p`$, $`M`$ is not generalized Cohen-Macaulay. โ
The notion of Buchsbaum modules is an intermediate concept between Cohen-Macaulay modules and generalized Cohen-Macaulay modules. is a nice introduction of this theory.
###### Proposition 2.12.
Assume that $`R`$ is Cohen-Macaulay and simplicial. Let $`M\mathrm{mod}_CR`$ be a $`C`$-graded module of dimension $`p`$. If $`\nu _i(W,M)=0`$ for all $`i0`$ and all monomial irreducible ideal $`W`$ with $`W๐ช`$ and $`dim(R/W)<pi`$, then $`M`$ is Buchsbaum.
###### Proof.
By an argument similar to the above, we can see that $`[H_๐ช^i(M)]_๐=0`$ for all $`i<p`$ and all $`0๐^d`$. So the assertion follows from \[23, Corollary 4.6\]. โ
## 3. Squarefree modules over a normal semigroup ring
For the results in the previous section, the assumption that $`R`$ is simplicial is really necessary. But when we consider radical monomial ideals in a normal semigroup ring, we can remove this assumption.
Throughout this section, we assume that $`R`$ is normal.
For a point $`u๐(=_0C)`$, we always have a unique face $`F๐`$ whose relative interior contains $`u`$. In this notation, we denote $`s(u)=F`$.
###### Definition 3.1 ().
Assume that $`R`$ is normal. We say an $`R`$-module $`M`$ is squarefree, if $`M\mathrm{mod}_CR`$ and the multiplication map $`M_๐y๐ฑ^๐yM_{๐+๐}`$ is bijective for all $`๐,๐C`$ with $`s(๐+๐)=s(๐)`$.
For a monomial ideal $`IR`$, $`I`$ (or $`R/I`$) is a squarefree module if and only if $`I`$ is a radical ideal (i.e., $`I=\sqrt{I}`$). Since the canonical module $`\omega _R`$ of $`R`$ is isomorphic to the ideal $`(๐ฑ^๐๐C\text{ with }s(๐)=๐)`$, it is also squarefree. It is easy to check that if $`M`$ is squarefree, we have $`dim_KM_๐=dim_KM_๐`$ for all $`๐,๐C`$ with $`s(๐)=s(๐)`$.
Let us recall basic properties of squarefree modules. See for detail. $`\mathrm{Sq}`$ denotes the full subcategory of $`\mathrm{Mod}R`$ consisting of squarefree modules. Then $`\mathrm{Sq}`$ is closed under kernels, cokernels and extensions in $`\mathrm{Mod}R`$. Hence $`\mathrm{Sq}`$ is an abelian category. Moreover, it admits enough projectives and injectives, and an indecomposable injective object is isomorphic to $`R/P_F`$ for some $`F๐`$.
If $`M\mathrm{mod}R`$ and $`N\mathrm{Mod}R`$, then the $`R`$-module $`\mathrm{Hom}_R(M,N)`$ has a natural $`^d`$-grading with $`\mathrm{Hom}_R(M,N)_๐=\mathrm{Hom}_{\mathrm{Mod}R}(M,N(๐))`$. Similarly, $`\mathrm{Ext}_R^i(M,N)`$ has a natural $`^d`$-grading.
###### Lemma 3.2 ().
Let $`M`$ be a squarefree $`R`$-module, and $`I^{}`$ its minimal irreducible resolution. Then we have the following.
* If $`\nu _i(W,M)0`$ for some $`i0`$, then $`W`$ is a prime ideal. Moreover, $`I^{}`$ is a minimal injective resolution of $`M`$ in $`\mathrm{Sq}`$.
* $`dimI^i>dimI^{i+1}`$ for all $`i`$. In particular, $`\mathrm{ir}\mathrm{dim}Md`$.
* $`\mathrm{Ext}_R^i(M,\omega _R)`$ is squarefree for all $`i`$.
Let $`M`$ be a squarefree module. For each $`F๐`$, take some $`๐(F)C\mathrm{rel}\mathrm{int}(F)`$ (i.e., $`s(๐(F))=F`$). If $`F,G๐`$ and $`GF`$, \[22, Theorem 3.3\] gives a $`K`$-linear map $`\phi _{G,F}^M:M_{๐(F)}M_{๐(G)}`$. These maps satisfy $`\phi _{F,F}^M=\mathrm{Id}`$ and $`\phi _{H,G}^M\phi _{G,F}^M=\phi _{H,F}^M`$ for all $`HGF`$.
###### Theorem 3.3 (\[22, Theorem 4.15\]).
For $`M\mathrm{Sq}`$ and $`F๐`$ with $`dimF=t`$, we have
(3.1)
$$\nu _i(P_F,M)=dim_K[\mathrm{Ext}_R^{dit}(M,\omega _R)]_{๐(F)}.$$
Since $`\mathrm{Ext}_R^{dit}(M,\omega _R)`$ is squarefree, the value of the right side of the equality (3.1) does not depend on the choice of $`๐(F)`$.
Propositions 3.4, 3.5 and 3.6, which follow from Theorem 3.3, are the squarefree modules version of results in the previous section. While (the latter part of) Proposition 3.4 has been stated in , we state it here for the readerโs convenience. Since $`R/P_F`$ is always Cohen-Macaulay whenever $`R`$ is normal, we can also prove these results by arguments similar to the previous section.
###### Proposition 3.4 ().
If $`M\mathrm{Sq}`$, then $`\mathrm{depth}_RM=\mathrm{min}\{dim(R/P_F)+i\nu _i(P_F,M)0\}.`$ In particular, $`M`$ is a Cohen-Macaulay module of dimension $`p`$ if and only if $`\nu _i(P_F,M)=0`$ for all $`i`$ and all $`P_F`$ with $`dim(R/P_F)pi`$ .
###### Proposition 3.5.
If $`IR`$ is a radical monomial ideal with $`dim(R/I)=p`$, then the following are equivalent for an integer $`n2`$.
* $`R/I`$ satisfies Serreโs condition $`(S_n)`$.
* If $`in1`$ and $`dim(R/P_F)pi`$, then $`\nu _i(P_F,R/I)=0`$.
For squarefree modules, we can prove the converse of Proposition 2.12 by virtue of \[23, Corollary 4.6\].
###### Proposition 3.6.
Assume that $`R`$ is normal. Let $`M`$ be a squarefree $`R`$-module of dimension $`p`$. Then $`M`$ is Buchsbaum, if and only if it is generalized Cohen-Macaulay, if and only if $`\nu _i(P_F,M)=0`$ for all $`i`$ and all monomial prime ideal $`P_F`$ with $`dim(R/P_F)0,pi`$.
The linear strand of a minimal free resolution of a finitely generated graded module over a polynomial ring is an important notion introduced by Eisenbud (c.f. ). Here we introduce the analog of this concept for an irreducible resolution of squarefree modules. (Miller also studied this concept implicitly. See \[14, ยง3\].)
Let $`M`$ be a squarefree $`R`$-module, and $`I^{}`$ its minimal irreducible resolution. For each $`l`$, we define the $`l`$-linear strand $`\mathrm{lin}_l(M)`$ of (the minimal irreducible resolution of) $`M`$ as follows: The term $`\mathrm{lin}_l(M)^i`$ of cohomological degree $`i`$ is
$$\underset{dim(R/P_F)=li}{}(R/P_F)^{\nu _i(P_F,M)},$$
which is a direct summand of $`I^i=(R/P_F)^{\nu _i(P_F,M)}`$, and the differential $`\mathrm{lin}_l(M)^i\mathrm{lin}_l(M)^{i+1}`$ is the corresponding component of the differential $`I^iI^{i+1}`$ of $`I^{}`$. By Proposition 3.4, $`M`$ is a Cohen-Macaulay module of dimension $`p`$ if and only if $`I^{}\mathrm{lin}_p(M)`$. Set $`\mathrm{lin}(M):=_{0ld}\mathrm{lin}_l(M)`$.
Using spectral sequence argument, we can construct $`\mathrm{lin}(M)`$ from a (not necessarily minimal) irreducible resolution $`J^{}`$ of $`M`$. For each $`i`$, let $`J^i_{j0}J^{i,j}`$ be a decomposition with $`dim(J^{i,j})=j`$ or $`J^{i,j}=0`$. Consider the filtration $`J^{}=F_0J^{}F_1J^{}\mathrm{}F_dJ^{}=0`$ with $`F_pJ^i=_{jdp}J^{i,j}`$, and the associated spectral sequence $`\{E_r^,,d_r\}`$. Since $`F_pJ^i`$ is a direct summand of $`J^i`$, we have $`E_0^{p,q}=(F_pJ^{p+q}/F_{p+1}J^{p+q})J^{p+q,dp}`$ and $`J_0^t:=_{p+q=t}E_0^{p,q}J^t`$. The maps $`d_0^{p,q}:E_0^{p,q}E_0^{p,q+1}`$ make $`J_0^{}`$ a cochain complex. On the other hand, we always have a decomposition $`J^{}=I^{}T^{}`$ such that $`I^{}`$ is minimal and $`T^{}`$ is exact. If we identify $`J_0^t`$ with $`J^t=I^tT^t`$, the differential $`d_0`$ of $`J_0^{}`$ is given by $`(0,d_T^{})`$. As before, let $`I^i_{j0}I^{i,j}`$ be a decomposition with $`dim(I^{i,j})=j`$ or $`I^{i,j}=0`$. Then we have $`E_1^{p,q}F_pI^{p+q}/F_{p+1}I^{p+q}I^{p+q,dp}`$ and $`J_1^t:=_{p+q=t}E_1^{p,q}I^t`$. The maps $`d_1^{p,q}:E_1^{p,q}E_1^{p+1,q}`$ make $`J_1^{}`$ a cochain complex, which is isomorphic to $`\mathrm{lin}(M)`$.
A minimal injective resolution of $`\omega _R`$ in $`\mathrm{Mod}R`$ is of the form
(3.2)
$$I^{}:0I^0I^1\mathrm{}I^d0,$$
$$I^i=\underset{\begin{array}{c}F๐\\ dimF=di\end{array}}{}{}_{}{}^{}E(R/P_F)$$
and the differential is composed of the maps $`(\pm 1)\mathrm{nat}:{}_{}{}^{}E(R/P_F){}_{}{}^{}E(R/P_F^{})`$ for $`F,F^{}๐`$ with $`FF^{}`$ and $`dimF=dimF^{}+1`$, where $`\mathrm{nat}:{}_{}{}^{}E(R/P_F){}_{}{}^{}E(R/P_F^{})`$ is induced by the natural surjection $`R/P_FR/P_F^{}`$, and the sign $`\pm `$ is given by an incidence function (c.f. \[1, ยง6.2\]) of $`๐`$.
Let $`D^b(\mathrm{Sq})`$ be the bounded derived category of $`\mathrm{Sq}`$. By \[23, Lemma 2.3\], we have $`D^b(\mathrm{Sq})D_{\mathrm{Sq}}^b(\mathrm{Mod}R)`$, so we will identify these categories. For a complex $`M^{}`$ and an integer $`n`$, let $`M^{}[n]`$ be the $`n^{\mathrm{th}}`$ translation of $`M^{}`$. That is, $`M^{}[n]`$ is a complex with $`M^i[n]=M^{i+n}`$. Since $`\mathrm{Ext}_R^i(M,\omega _R)\mathrm{Sq}`$ for all $`M\mathrm{Sq}`$,
$$๐():=\mathrm{R}\mathrm{Hom}_R(,\omega _R)$$
gives a contravariant functor from $`D^b(\mathrm{Sq})`$ to itself satisfying $`๐๐\mathrm{Id}_{D^b(\mathrm{Sq})}`$.
For $`M\mathrm{Mod}R`$, we say the submodule $`_{๐C}M_๐`$ is the $`C`$-graded part of $`M`$. Since $`R`$ is normal now, the $`C`$-graded part of $`{}_{}{}^{}E(R/P_F)`$ is isomorphic to $`R/P_F`$ (c.f. \[15, Chapter 13\]). Moreover, we have the following.
###### Lemma 3.7.
Assume that $`R`$ is normal. For $`M\mathrm{Sq}`$ and $`F๐`$, the $`C`$-graded part of $`\mathrm{Hom}_R(M,{}_{}{}^{}E(R/P_F))`$ is isomorphic to $`(M_{๐(F)})^{}_K(R/P_F)`$.
###### Proof.
When $`R`$ is a polynomial ring, this result was given in \[21, Lemma 3.20\]. The general case can be proved by essentially same argument. But we give a precise proof here for the readerโs convenience.
Let $`M^{}`$ be the submodule of $`M`$ generated by $`M_{๐(F)}`$. By the injectivity of $`{}_{}{}^{}E(R/P_F)`$, $`f\mathrm{Hom}_{\mathrm{Mod}R}(M^{},{}_{}{}^{}E(R/P_F)(๐))`$ for $`๐C`$ can be extended to $`M{}_{}{}^{}E(R/P_F)(๐)`$. On the other hand, if $`๐^{}CF`$, then $`[{}_{}{}^{}E(R/P_F)]_{๐+๐^{}}=0`$. Hence, for $`0g\mathrm{Hom}_{\mathrm{Mod}R}(M,{}_{}{}^{}E(R/P_F)(๐))`$, there is some $`๐^{}F`$ such that $`g(y)0`$ for some $`yM_๐^{}`$. In this situation, we have $`g(๐ฑ^{๐(F)}y)=๐ฑ^{๐(F)}g(y)0`$. Since $`M`$ is squarefree and $`s(๐(F))=s(๐(F)+๐^{})`$, there is $`zM_{๐(F)}`$ such that $`๐ฑ^๐^{}z=๐ฑ^{๐(F)}y`$. Clearly, $`g(z)0`$. Hence the restriction $`g|_M^{}`$ of $`g\mathrm{Hom}_{\mathrm{Mod}R}(M,{}_{}{}^{}E(R/P_F)(๐))`$ to $`M^{}`$ is 0 if and only if $`g=0`$. Combining these observations, we have
$$[\mathrm{Hom}_R(M^{},{}_{}{}^{}E(R/P_F))]_C[\mathrm{Hom}_R(M,{}_{}{}^{}E(R/P_F))]_C.$$
Since $`\mathrm{ann}(y)P_F`$ for all $`0yM_{๐(F)}`$, $`f\mathrm{Hom}_K(M_{๐(F)},{}_{}{}^{}E(R/P_F)(๐))`$ can be extended to $`\stackrel{~}{f}\mathrm{Hom}_{\mathrm{Mod}R}(M^{},{}_{}{}^{}E(R/P_F)(๐))`$. Note that $`{}_{}{}^{}E(R/P_F)_{๐+๐(F)}0`$ if and only if $`{}_{}{}^{}E(R/P_F)_{๐+๐(F)}K`$ if and only if $`๐F`$. Hence we have $`[\mathrm{Hom}_R(M^{},{}_{}{}^{}E(R/P_F))]_C(M_{๐(F)})^{}_KR/P_F`$ as graded $`K`$-vector spaces. But it is easy to see that this is also an isomorphism of $`R`$-modules. โ
###### Lemma 3.8.
Assume that $`R`$ is normal. If $`M\mathrm{Sq}`$, then $`๐(M)`$ is quasi-isomorphic to the complex $`D^{}:0D^0D^1\mathrm{}D^d0`$ with
(3.3)
$$D^i=\underset{\begin{array}{c}F๐\\ dimF=di\end{array}}{}(M_{๐(F)})^{}_K(R/P_F).$$
Here the differential is the sum of the maps
$$(\pm \phi _{F,F^{}}^M)^{}\mathrm{nat}:(M_{๐(F)})^{}_KR/P_F(M_{๐(F^{})})^{}_KR/P_F^{}$$
for $`F,F^{}๐`$ with $`FF^{}`$ and $`dimF=dimF^{}+1`$, and $`\mathrm{nat}`$ denotes the natural surjection $`R/P_FR/P_F^{}`$.
Moreover, if $`M^{}D^b(\mathrm{Sq})`$, then $`๐(M^{})`$ is quasi-isomorphic to the total complex of the double complex $`W^{i,j}(M^{})=D^i(M^j)`$. Here the differential $`W^{i,j}=D^i(M^j)W^{i+1,j}=D^{i+1}(M^j)`$ is the $`i^{\mathrm{th}}`$ differential of $`D^{}(M^j)`$ and the differential
$$W^{i,j}=\underset{\begin{array}{c}F๐\\ dimF=di\end{array}}{}(M_{๐(F)}^j)^{}_K(R/P_F)W^{i,j+1}=\underset{\begin{array}{c}F๐\\ dimF=di\end{array}}{}(M_{๐(F)}^{j1})^{}_K(R/P_F)$$
is induced by the differential $`_M^{}^{j1}:M^{j1}M^j`$ of $`M^{}`$.
###### Proof.
When $`R`$ is a polynomial ring, the assertion was given in \[24, ยง3\]. The general case can be proved by a similar argument to . Here we only remark that $`\mathrm{R}\mathrm{Hom}_R(M^{},\omega _R)`$ is quasi-isomorphic to $`\mathrm{Hom}_R^{}(M^{},I^{})`$ where $`I^{}`$ is the minimal injective resolution of $`\omega _R`$ described in (3.2), and the isomorphism (3.3) clearly follows from Lemma 3.7. โ
Convention. In the sequel, as an explicit complex, $`๐(M^{})`$ for $`M^{}D^b(\mathrm{Sq})`$ means the complex described in Lemma 3.8.
The next result refines Theorem 3.3.
###### Theorem 3.9.
Assume that $`R`$ is normal. If $`M\mathrm{Sq}`$, then we have $`\mathrm{lin}_i(M)๐(\mathrm{Ext}_R^{di}(M,\omega _R))[di]`$ for all $`i`$.
###### Proof.
Note that $`๐๐(M)`$ is a complex of injective objects in $`\mathrm{Sq}`$, and quasi-isomorphic to $`M`$. Hence it is a (non-minimal) irreducible resolution of $`M`$. Set $`J^{}:=๐๐(M)`$ and $`N^{}:=๐(M)`$. Recall that $`J^{}=๐(N^{})`$ is the total complex of the double complex $`W^{}`$ defined in Lemma 3.8. We use the same notation as in the spectral sequence construction of $`\mathrm{lin}(M)`$. Recall that $`J^iJ_0^i`$ for all $`i`$. The construction of $`J_0^{}`$ cancels the horizontal differential of $`W^{}`$ (i.e., the differential which comes from that of $`๐(N^i)`$). Hence the differential of $`J_0^{}`$ is just induced by that of $`N^{}`$. So if we set $`(N^{})`$ to be the complex such that $`(N^{})^i=H^i(N)`$ for all $`i`$ and the differential maps are zero, then $`J_1^{}`$ is isomorphic to $`๐((N^{}))`$. Since the $`i^{\mathrm{th}}`$ cohomology of $`(N^{})`$ is $`H^i(N)=\mathrm{Ext}_R^i(M,\omega _R)`$, the assertion follows. โ
###### Definition 3.10 (Stanley, ).
Let $`M\mathrm{mod}R`$. We say $`M`$ is sequentially Cohen-Macaulay if there is a finite filtration
$$0=M_0M_1\mathrm{}M_r=M$$
of $`M`$ by graded submodules $`M_i`$ satisfying the following conditions.
* Each quotient $`M_i/M_{i1}`$ is Cohen-Macaulay.
* $`dim(M_i/M_{i1})<dim(M_{i+1}/M_i)`$ for all $`i`$.
###### Corollary 3.11.
Let $`M\mathrm{Sq}`$. Then $`M`$ is sequentially Cohen-Macaulay if and only if $`\mathrm{lin}(M)`$ is acyclic (i.e., $`H^i(\mathrm{lin}(M))=0`$ for all $`i0`$).
###### Proof.
For $`0N\mathrm{Sq}`$, $`H^i(๐(N))=0`$ for all $`in`$ if and only if $`N`$ is Cohen-Macaulay and $`dimN=dn`$. By Theorem 3.9, $`\mathrm{lin}(M)`$ is acyclic if and only if $`\mathrm{Ext}_R^{di}(M,\omega _R)`$ is a Cohen-Macaulay module of dimension $`i`$ for all $`i`$. So the assertion follows from \[17, Theorem III.2.11\]. โ
###### Remark 3.12.
Assume that $`R`$ is a polynomial ring. It is easy to see that $`\mathrm{lin}(M)`$ of $`M\mathrm{Sq}`$ is acyclic if and only if the Alexander dual of $`M`$ is componentwise linear (see, for example, for this concept). So Corollary 3.11 generalizes \[16, Theorem 4.5\].
## 4. Sequentially Cohen-Macaulay face rings
In this section, we always assume that $`R`$ is normal.
We have a hyperplane $`H^d`$ such that $`B:=H๐`$ is a $`(d1)`$-dimensional polytope. Clearly, $`B`$ has the topology induced by the Euclidean space $`^d`$, and $`B`$ is homeomorphic to a closed ball of dimension $`d1`$. For a face $`F๐`$, set $`|F|`$ to be the relative interior of $`FH`$. Regarding $`๐`$ as a partially ordered set by inclusions, we say $`\mathrm{\Delta }๐`$ is an order ideal, if $`\mathrm{\Delta }`$ is a non-empty subset such that; $`F\mathrm{\Delta }`$, $`G๐`$ and $`FG`$ $``$ $`G\mathrm{\Delta }`$. If $`\mathrm{\Delta }`$ is an order ideal, then $`|\mathrm{\Delta }|:=_{F\mathrm{\Delta }}|F|`$ is a closed subset of $`B`$, and $`_{F\mathrm{\Delta }}|F|`$ is a regular cell decomposition (c.f. \[1, ยง6.2\]) of $`|\mathrm{\Delta }|`$. Up to homeomorphism, (the regular cell decomposition of) $`|\mathrm{\Delta }|`$ does not depend on the particular choice of the hyperplane $`H`$. The dimension $`dim|\mathrm{\Delta }|`$ of $`|\mathrm{\Delta }|`$ is given by $`\mathrm{max}\{dim|F|F\mathrm{\Delta }\}`$. Here $`dim|F|`$ denotes the dimension of $`|F|`$ as a cell (we set $`dim\mathrm{}=1`$), that is, $`dim|F|=dimF1`$.
We assign an order ideal $`\mathrm{\Delta }๐`$ to the ideal $`I_\mathrm{\Delta }:=(๐ฑ^๐๐C\text{ and }s(๐)\mathrm{\Delta })`$ of $`R`$. (For the definition of $`s(๐)`$, see the beginning of ยง3.) Note that $`I_\mathrm{\Delta }`$ is a radical ideal, and any radical monomial ideal of $`R`$ is of the form $`I_\mathrm{\Delta }`$ for some $`\mathrm{\Delta }`$. Set $`K[\mathrm{\Delta }]:=R/I_\mathrm{\Delta }`$. Clearly,
$$K[\mathrm{\Delta }]_๐\{\begin{array}{cc}K\hfill & \text{if }๐C\text{ and }s(๐)\mathrm{\Delta }\text{,}\hfill \\ 0\hfill & \text{otherwise.}\hfill \end{array}$$
In particular, if $`\mathrm{\Delta }=๐`$ (resp. $`\mathrm{\Delta }=\{\{0\}\}`$), then $`I_\mathrm{\Delta }=0`$ (resp. $`I_\mathrm{\Delta }=๐ช`$) and $`K[\mathrm{\Delta }]=R`$ (resp. $`K[\mathrm{\Delta }]=K`$). We have $`dimK[\mathrm{\Delta }]=dim|\mathrm{\Delta }|+1`$. When $`R`$ is a polynomial ring, $`K[\mathrm{\Delta }]`$ is nothing other than the Stanley-Reisner ring of a simplicial complex $`\mathrm{\Delta }`$. (If $`R`$ is simplicial, then $`B`$ is a simplex and $`\mathrm{\Delta }`$ is a simplicial complex.) Clearly, $`I_\mathrm{\Delta }`$ and $`K[\mathrm{\Delta }]`$ are squarefree modules.
Let $`\mathrm{\Delta },\mathrm{\Sigma }๐`$ be order ideals with $`\mathrm{\Delta }\mathrm{\Sigma }`$. When we consider such a pair, we assume that $`\mathrm{\Sigma }\{\{0\}\}`$, but allow the case $`\mathrm{\Sigma }=\mathrm{}`$. Thus, $`|\mathrm{\Sigma }|=\mathrm{}`$ if and only if $`\mathrm{\Sigma }=\mathrm{}`$. We set $`I_{\mathrm{}}=R`$. We always have $`I_\mathrm{\Sigma }I_\mathrm{\Delta }`$, and $`I_{\mathrm{\Delta }/\mathrm{\Sigma }}:=I_\mathrm{\Sigma }/I_\mathrm{\Delta }`$ is a squarefree $`R`$-module. Note that $`I_{\mathrm{\Delta }/\mathrm{}}=K[\mathrm{\Delta }]`$. We say a pair $`(\mathrm{\Delta },\mathrm{\Sigma })`$ is Cohen-Macaulay if so is the module $`I_{\mathrm{\Delta }/\mathrm{\Sigma }}`$. (The Cohen-Macaulay property of $`(\mathrm{\Delta },\mathrm{\Sigma })`$ depends on $`\mathrm{char}(K)`$. But, in this paper, we fix the base field $`K`$. So we omit the phrase โover $`K`$โ.) Clearly, the Cohen-Macaulay property of $`(\mathrm{\Delta },\mathrm{\Sigma })`$ depends only on $`\mathrm{\Delta }\mathrm{\Sigma }`$. For a topological space $`X`$, $`\mathrm{Sh}(X)`$ denotes the category of sheaves of $`K`$-vector spaces on $`X`$. Recall that $`Z:=|\mathrm{\Delta }||\mathrm{\Sigma }|`$ admits Verdierโs dualizing complex (over $`K`$) $`๐_Z^{}D^b(\mathrm{Sh}(Z))`$. See \[12, V. ยง2\]. Recall a few results from , some parts of which have been obtained by Stanley \[18, ยงยง4,5\] (this does not mean that the next theorem contains corresponding results in \[18, ยงยง4,5\], since the rings treated there are more general than ours).
###### Theorem 4.1 (\[23, Theorem 4.9, Proposition 4.10\]).
We have the following.
* $`K[\mathrm{\Delta }]`$ is Cohen-Macaulay if and only if $`^i(๐_{|\mathrm{\Delta }|}^{})=0`$ and $`\stackrel{~}{H}^i(|\mathrm{\Delta }|;K)=0`$ for all $`idim|\mathrm{\Delta }|`$.
* Let $`(\mathrm{\Delta },\mathrm{\Sigma })`$ be a pair of order ideals of $`๐`$ with $`\mathrm{\Sigma }\mathrm{}`$, and $`h`$ the embedding map from $`Z:=|\mathrm{\Delta }||\mathrm{\Sigma }|`$ to its closure $`\overline{Z}=|\mathrm{\Delta }|`$. Then $`(\mathrm{\Delta },\mathrm{\Sigma })`$ is Cohen-Macaulay if and only if $`\mathrm{R}^ih_{}๐_Z^{}=0`$ and $`H_c^i(Z;K)=0`$ for all $`idimZ`$. Here $`H_c^i()`$ denotes the cohomology with compact support.
Consequently, either $`\mathrm{\Sigma }=\mathrm{}`$ or $`\mathrm{\Sigma }\mathrm{}`$, the Cohen-Macaulay property of $`(\mathrm{\Delta },\mathrm{\Sigma })`$ depends only on the pair of topological spaces $`(|\mathrm{\Delta }|,|\mathrm{\Sigma }|)`$ (or even $`(\overline{Z},Z)`$).
###### Remark 4.2.
We can check the Cohen-Macaulay property of a pair $`(\mathrm{\Delta },\mathrm{\Sigma })`$ explicitly. Recall the combinatorial description of $`๐(I_{\mathrm{\Delta }/\mathrm{\Sigma }})=\mathrm{R}\mathrm{Hom}_R(I_{\mathrm{\Delta }/\mathrm{\Sigma }},\omega _R)D^b(\mathrm{Sq})`$ given in Lemma 3.8. For each $`F๐`$, the complex $`[๐(I_{\mathrm{\Delta }/\mathrm{\Sigma }})]_{๐(F)}`$ of $`K`$-vector spaces is isomorphic to the complex $`D_F^{}`$ defined by
$$D_F^i:=\underset{\begin{array}{c}G\mathrm{\Delta }\mathrm{\Sigma },GF\\ dimG=di\end{array}}{}Ke_G,$$
$$:e_G\underset{\begin{array}{c}G^{}\mathrm{\Delta }\mathrm{\Sigma },GG^{}F\\ dimG^{}=dimG1\end{array}}{}\pm e_G^{},$$
where $`e_G`$ is a basis element, and the sign $`\pm `$ is given by an incidence function of $`๐`$. Hence $`(\mathrm{\Delta },\mathrm{\Sigma })`$ is Cohen-Macaulay if and only if $`H^i(D_F^{})=0`$ for all $`iddimI_{\mathrm{\Delta }/\mathrm{\Sigma }}`$ and all $`F\mathrm{\Delta }`$.
For $`F\mathrm{\Delta }`$, set $`\delta (F):=\mathrm{max}\{dim|G|GF,G\mathrm{\Delta }\}`$. For $`i`$ with $`idim|\mathrm{\Delta }|`$, we call $`\mathrm{\Delta }^{(i)}:=\{F\mathrm{\Delta }dim|F|i\}`$ the $`i`$-skeleton of $`\mathrm{\Delta }`$, and $`\mathrm{\Delta }^{[i]}:=\{F\mathrm{\Delta }^{(i)}\delta (F)i\}`$ the pure $`i`$-skeleton of $`\mathrm{\Delta }`$. Clearly, these are order ideals of $`๐`$ again. It is easy to see that
(4.1)
$$\mathrm{\Delta }^{[i]}(\mathrm{\Delta }^{[i+1]})^{(i)}=\{F\mathrm{\Delta }\delta (F)=i\}$$
for all $`idim|\mathrm{\Delta }|=:r`$. Here we set $`(\mathrm{\Delta }^{[r+1]})^{(r)}=\mathrm{}`$.
For a finitely generated graded $`R`$-module $`M`$ and an integer $`i`$, we set
$$M_{(i)}:=\{yMdim(Ry)i\}.$$
Then $`M_{(i)}`$ is a submodule of $`M`$ with $`M=M_{(d)}M_{(d1)}\mathrm{}M_{(0)}M_{(1)}=0`$. It is easy to see that $`M`$ is sequentially Cohen-Macaulay if and only if $`M_{(i)}/M_{(i1)}`$ is Cohen-Macaulay (of course, we allow the case $`M_{(i)}/M_{(i1)}=0`$) for all $`i`$.
Since $`K[\mathrm{\Delta }]_{(i+1)}/K[\mathrm{\Delta }]_{(i)}I_{\mathrm{\Delta }^{[i]}/(\mathrm{\Delta }^{[i+1]})^{(i)}}`$, we have the following.
###### Lemma 4.3 (c.f. \[17, Proposition 2.10\], \[3, Remark in p.4\]).
With the above notation, $`K[\mathrm{\Delta }]`$ is sequentially Cohen-Macaulay if and only if the pair $`(\mathrm{\Delta }^{[i]},(\mathrm{\Delta }^{[i+1]})^{(i)})`$ is Cohen-Macaulay for all $`idim|\mathrm{\Delta }|`$.
When $`R`$ is a polynomial ring (i.e., $`K[\mathrm{\Delta }]`$ is a Stanley-Reisner ring), the next result has been obtained by Duval .
###### Theorem 4.4 (c.f. Duval \[3, Theorem 3.3\]).
For an order ideal $`\mathrm{\Delta }๐`$, $`K[\mathrm{\Delta }]`$ is sequentially Cohen-Macaulay if and only if $`K[\mathrm{\Delta }^{[i]}]`$ is Cohen-Macaulay for all $`idim|\mathrm{\Delta }|`$.
###### Proof.
The proof of \[3, Theorem 3.3\] for the Stanley-Reisner ring case also works here. This proof uses Lemma 4.3 and Lemma 4.5 below. The latter, which is a well-known result in the Stanley-Reisner ring case, also holds in our context. โ
###### Lemma 4.5.
If $`K[\mathrm{\Delta }]`$ is Cohen-Macaulay, then so is $`K[\mathrm{\Delta }^{(i)}]`$ for all $`idim|\mathrm{\Delta }|`$.
###### Proof.
We may assume that $`i+1=dim|\mathrm{\Delta }|=:r`$. For each $`F\mathrm{\Delta }`$, the canonical module $`\omega _{K[F]}`$ of $`K[F]:=R/P_F`$ is a squarefree $`R`$-module with
$$[\omega _{K[F]}]_๐=\{\begin{array}{cc}K\hfill & \text{if }s(๐)=F\text{,}\hfill \\ 0\hfill & \text{otherwise,}\hfill \end{array}$$
for all $`๐C`$. We have a short exact sequence
(4.2)
$$0\underset{\begin{array}{c}F\mathrm{\Delta }\\ dim|F|=r\end{array}}{}\omega _{K[F]}K[\mathrm{\Delta }]K[\mathrm{\Delta }^{(r1)}]0.$$
Since $`\omega _{K[F]}`$ for $`F`$ with $`dim|F|=r`$ is a Cohen-Macaulay module of dimension $`r+1`$ $`(=dimK[\mathrm{\Delta }])`$ and $`dimK[\mathrm{\Delta }^{(r1)}]=r`$, $`K[\mathrm{\Delta }^{(r1)}]`$ is Cohen-Macaulay. โ
###### Remark 4.6.
In the Stanley-Reisner ring case, Hibi (\[10, ยง2\]) showed that if $`i<dim|\mathrm{\Delta }|`$ then the canonical module of $`K[\mathrm{\Delta }^{(i)}]`$ is generated by its degree 0 part. This also holds in our context. To see this, we may assume that $`i+1=dim|\mathrm{\Delta }|=:r`$. Then the long exact sequence of $`\mathrm{Ext}_R^i(,\omega _R)`$ derived from the sequence (4.2) gives
$$\underset{\begin{array}{c}F\mathrm{\Delta }\\ dim|F|=r\end{array}}{}K[F]\mathrm{Ext}_R^{dr}(K[\mathrm{\Delta }^{(i)}],\omega _R)0\text{(exact)}.$$
The module of right side, which is isomorphic to the canonical module of $`K[\mathrm{\Delta }^{(i)}]`$, is clearly generated by its degree 0 part.
Since $`|\mathrm{\Delta }^{[i]}|`$ depends on the cell decomposition $`|\mathrm{\Delta }|=_{F\mathrm{\Delta }}|F|`$, Theorem 4.4 does not implies that the sequentially Cohen-Macaulay property of $`K[\mathrm{\Delta }]`$ is a topological property of $`|\mathrm{\Delta }|`$. But we can prove this using Lemma 4.3 directly.
###### Theorem 4.7.
The sequentially Cohen-Macaulay property of $`K[\mathrm{\Delta }]`$ only depends on the topological space $`|\mathrm{\Delta }|`$.
###### Proof.
Set $`X:=|\mathrm{\Delta }|`$. Note that the subspace
$$\underset{F\mathrm{\Delta }^{[i]}}{}|F|\underset{F(\mathrm{\Delta }^{[i+1]})^{(i)}}{}|F|$$
of $`X`$ does not depend on the particular cell decomposition of $`X`$ by (4.1). In fact, it coincides with $`\{xXdim_xX=i\}`$, where the dimension $`dim_xX`$ of $`X`$ at $`x`$ is the one defend in \[12, III, Definition 9.10\]. So the assertion follows from Theorem 4.1 and Lemma 4.3. โ
Even in the Stanley-Reisner ring case, it is well-known that the Gorenstein property of $`K[\mathrm{\Delta }]`$ is not a topological property of $`|\mathrm{\Delta }|`$ (i.e., depends on the simplicial decomposition). In the normal semigroup ring case, we have another problem. A normal semigroup ring $`R`$ is always Cohen-Macaulay, but not necessarily Gorenstein, even if it is simplicial. Note that $`K[\mathrm{\Delta }]=R`$ if $`\mathrm{\Delta }=๐`$. So the Gorenstein property can not be determined by the poset structure of $`\mathrm{\Delta }`$ (i.e., sensitive to the semigroup $`C`$). But we can prove that the Gorenstein\* property is topological. Recently, Ichim and Rรถmer also studied the Gorenstein (or Gorenstein\*) property of a toric face ring, which is a notion containing our $`K[\mathrm{\Delta }]`$. Their \[11, Corollary 6.9\] is closely related to Theorem 4.9 below.
###### Definition 4.8 (Stanley, \[17, p.67\]).
We say $`K[\mathrm{\Delta }]`$ is Gorenstein\*, if it is Gorenstein and the canonical module $`\omega _{K[\mathrm{\Delta }]}`$ of $`K[\mathrm{\Delta }]`$ is generated by its degree 0 part (equivalently, $`\omega _{K[\mathrm{\Delta }]}K[\mathrm{\Delta }]`$ as graded modules).
###### Proposition 4.9.
Set $`X:=|\mathrm{\Delta }|`$, and let $`๐_X^{}D^b(\mathrm{Sh}(X))`$ be the dualizing complex of $`X`$. Assume that $`r:=dimX1`$ and $`K[\mathrm{\Delta }]`$ is Cohen-Macaulay. Then $`K[\mathrm{\Delta }]`$ is Gorenstein\* if and only if $`^r(๐_X^{})_x=K`$ for all $`xX`$ and $`H^r(X;K)0`$. Thus the Gorenstein\* property of $`K[\mathrm{\Delta }]`$ only depends on the topological space $`X`$.
Under the assumption of Proposition 4.9, we have $`^r(๐_X^{})[r]๐_X^{}`$ in $`D^b(\mathrm{Sh}(X))`$. When $`X`$ is a manifold with or without boundary (but $`K[\mathrm{\Delta }]`$ is not necessarily Cohen-Macaulay), $`^r(๐_X^{})`$ is called the orientation sheaf of $`X`$ (over $`K`$). Hence, when $`X`$ is a manifold and $`K[\mathrm{\Delta }]`$ is Cohen-Macaulay, $`K[\mathrm{\Delta }]`$ is Gorenstein\* if and only if $`X`$ is an orientable manifold over $`K`$ by \[23, Theorem 4.2\] (quoted as Theorem 5.7 below). In particular, if $`X`$ is homeomorphic to an $`r`$-dimensional sphere, then $`K[\mathrm{\Delta }]`$ is Gorenstein\*.
###### Proof.
The last statement easily follows from the first one. In fact, the assumption that $`r1`$ is not essential. If $`r=1`$ (i.e., $`X=\mathrm{}`$), then $`K[\mathrm{\Delta }]`$ is Gorenstein\*. When $`r=0`$, $`K[\mathrm{\Delta }]`$ is Gorenstein\* if and only if $`X`$ consists of two points. On the other hand, we have seen that the Cohen-Macaulay property of $`K[\mathrm{\Delta }]`$ is a topological property of $`X`$. So it suffices to prove the first statement. Recall that $`\omega _{K[\mathrm{\Delta }]}`$ is a squarefree module with $`\mathrm{Ass}(K[\mathrm{\Delta }])=\mathrm{Ass}(\omega _{K[\mathrm{\Delta }]})`$. By \[23, Theorem 4.2\], $`^r(๐_X^{})_x=K`$ for all $`xX`$ if and only if $`[\omega _{K[\mathrm{\Delta }]}]_{๐(F)}=K`$ for all $`F\mathrm{\Delta }\{0\}`$. Hence, under the assumption that $`^r(๐_X^{})_x=K`$ for all $`xX`$, we have $`\omega _{K[\mathrm{\Delta }]}K[\mathrm{\Delta }]`$ as graded modules, if and only if $`\omega _{K[\mathrm{\Delta }]}`$ is generated by its degree 0 part, if and only if $`[\omega _{K[\mathrm{\Delta }]}]_00`$ (equivalently, $`[\omega _{K[\mathrm{\Delta }]}]_0=K`$). By \[23, Theorem 3.3\] (quoted as Theorem 5.2 below) and the local duality, we have $`H^r(X;K)[H_๐ช^{r+1}(K[\mathrm{\Delta }])]_0[\omega _{K[\mathrm{\Delta }]}]_0^{}`$. So we are done. โ
###### Example 4.10.
In Proposition 4.9, that $`^r(๐_X^{})_x=K`$ for all $`xX`$ is not a sufficient condition for $`K[\mathrm{\Delta }]`$ to be Gorenstein\*. In fact, if $`X`$ is a manifold without boundary, then $`^r(๐_X^{})_x=K`$ for all $`xX`$. But, of course, $`X`$ need not be orientable. For example, if $`X`$ is homeomorphic to a real projective plane and $`\mathrm{char}(K)2`$, then $`K[\mathrm{\Delta }]`$ is Cohen-Macaulay but not Gorenstein\*.
## 5. Constructible sheaves associated with squarefree modules
Throughout this section, $`R`$ is normal, and $`\mathrm{\Delta }`$ is an order ideal of $`๐`$ with $`X:=|\mathrm{\Delta }|`$.
If $`F\mathrm{\Delta }`$, then $`U_F:=X(_{F^{}F}|F^{}|)`$ is an open set of $`X`$. Note that $`\{U_F\{0\}F\mathrm{\Delta }\}`$ is an open covering of $`X`$. If $`M`$ is a squarefree $`K[\mathrm{\Delta }]`$-module (i.e., $`M\mathrm{Sq}`$ and $`\mathrm{ann}(M)I_\mathrm{\Delta }`$), then we can construct a sheaf $`M^+`$ on $`X`$ as in (see also Remark 5.3 (1) below). More precisely, the assignment $`\mathrm{\Gamma }(U_F,M^+)=M_{๐(F)}`$ for each $`\{0\}F๐`$ and the map $`\phi _{F,F^{}}^M:M^+(U_F^{})=M_{๐(F^{})}M_{๐(F)}=\mathrm{\Gamma }(U_F,M^+)`$ for $`\{0\}F,F^{}\mathrm{\Delta }`$ with $`FF^{}`$ (equivalently, $`U_F^{}U_F`$) actually defines a sheaf. It is easy to see that $`M^+`$ is a constructible sheaf with respect to the cell decomposition $`X=_{F\mathrm{\Delta }}|F|`$. In fact, for all $`\{0\}F\mathrm{\Delta }`$, the restriction $`M^+|_{|F|}`$ of $`M^+`$ to $`|F|X`$ is a constant sheaf with coefficients in $`M_{๐(F)}`$. Thus, if $`p|F|X`$, then the stalk $`(M^+)_p`$ at $`p`$ is isomorphic to $`M_{๐(F)}`$. If $`M`$ is a squarefree module, then so is the submodule $`M_{>0}:=_{0๐C}M_๐`$. For squarefree $`K[\mathrm{\Delta }]`$-modules $`M`$ and $`N`$, it is easy to see that $`M^+N^+`$ if and only if $`M_{>0}N_{>0}`$. In other words, $`M_0`$ is โirrelevantโ to $`M^+`$.
Let $`\mathrm{Sq}(\mathrm{\Delta })`$ be the category of squarefree $`K[\mathrm{\Delta }]`$-modules (hence $`\mathrm{Sq}(\mathrm{\Delta })`$ is a full subcategory of $`\mathrm{Sq}`$). The above operation gives the exact functor $`()^+:\mathrm{Sq}(\mathrm{\Delta })\mathrm{Sh}(X)`$.
###### Example 5.1.
Note that $`K[\mathrm{\Delta }]`$ itself is a squarefree $`K[\mathrm{\Delta }]`$-module, and we have $`K[\mathrm{\Delta }]^+\underset{ยฏ}{K}`$, where $`\underset{ยฏ}{K}`$ is the constant sheaf on $`X:=|\mathrm{\Delta }|`$ with coefficients in $`K`$. For a face $`\{0\}F\mathrm{\Delta }`$, we set $`K[F]:=K[\mathrm{\Delta }]/P_F`$. Then we have $`K[F]^+j_{}\underset{ยฏ}{K}_{|F|}\mathrm{Sh}(X)`$, where $`j:|F|X`$ is the embedding map and $`\underset{ยฏ}{K}_{|F|}`$ is the constant sheaf on $`|F|`$ (note that $`j_{}\underset{ยฏ}{K}_{|F|}`$ is essentially the constant sheaf on the closure $`\overline{|F|}`$ of $`|F|`$, not on $`|F|`$ itself). Recall that the canonical module $`\omega _R`$ is a squarefree $`R`$-module, and $`R=K[\mathrm{\Delta }]`$ if $`\mathrm{\Delta }=๐`$. Then $`\omega _R^+h_!\underset{ยฏ}{K}_B^{}\mathrm{Sh}(B)`$, where $`B^{}`$ is the interior of $`B`$, $`h:B^{}B`$ is the embedding map, and $`\underset{ยฏ}{K}_B^{}`$ is the constant sheaf on $`B^{}`$. It is noteworthy that $`h_!\underset{ยฏ}{K}_B^{}`$ is the orientation sheaf of $`B`$ (over $`K`$). This is a key point for Theorem 5.7 below.
###### Theorem 5.2 (\[23, Theorem 3.3\]).
For $`M\mathrm{Sq}(\mathrm{\Delta })`$, we have an isomorphism
$$H^i(X;M^+)[H_๐ช^{i+1}(M)]_0\text{for all }i1,$$
and an exact sequence
(5.1)
$$0[H_๐ช^0(M)]_0M_0H^0(X;M^+)[H_๐ช^1(M)]_00.$$
In particular, we have $`[H_๐ช^{i+1}(k[\mathrm{\Delta }])]_0\stackrel{~}{H}^i(X;K)`$ for all $`i0`$, where $`\stackrel{~}{H}^i(X;K)`$ denotes the $`i^{th}`$ reduced cohomology of $`X`$ with coefficients in $`K`$.
###### Remark 5.3.
(1) In , we construct the sheaf $`M^+`$ regarding $`M`$ as an $`R`$-module (not a $`K[\mathrm{\Delta }]`$-module). Thus, $`M^+`$ is always a sheaf on $`B`$ there. But this is not a problem. In fact, if $`\mathrm{Sh}(B)`$ is the sheaf constructed from $`M\mathrm{Sq}(\mathrm{\Delta })`$ in the style of , and $`M^+\mathrm{Sh}(X)`$ is the sheaf constructed in the style of the present paper, then we have $`i_{}M^+`$. Here $`i:XB`$ is the embedding map. Since $`H^i(X;M^+)H^i(B;)`$ for all $`i`$, \[23, Theorem 3.3\] implies Theorem 5.2.
(2) When $`R`$ is a polynomial ring (i.e., $`K[\mathrm{\Delta }]`$ is a Stanley-Reisner ring), the last statement of Theorem 5.2 is (a part of) a famous formula of Hochster (\[17, Theorem II.4.1\]). We can also compute $`[H_๐ช^i(M)]_๐`$ for any $`๐^d`$ in terms of the sheaf $`M^+`$ (see \[23, Theorem 3.5\]), while we do not use this result here.
(3) For $`yM_0`$, the sections $`\phi _{F,\{0\}}^M(y)M_{๐(F)}=\mathrm{\Gamma }(U_F;M^+)`$ for all $`\{0\}F๐`$ define the unique global section $`\stackrel{~}{y}\mathrm{\Gamma }(X;M^+)`$. The map $`M_0H^0(X,M^+)`$ in the sequence (5.1) is given by $`y\stackrel{~}{y}`$.
###### Definition 5.4.
A squarefree module $`M`$ is regular, if $`[H_๐ช^0(M)]_0=[H_๐ช^1(M)]_0=0`$.
###### Lemma 5.5.
Let $`M`$ be a squarefree $`K[\mathrm{\Delta }]`$-module. There is a unique (up to isomorphism) squarefree $`K[\mathrm{\Delta }]`$-module $`\overline{M}`$ which is regular and $`\overline{M}^+M^+`$.
###### Proof.
We put $`\overline{M}_{>0}=M_{>0}`$ and $`\overline{M}_0=\mathrm{\Gamma }(X;M^+)`$. It suffices to define $`\phi _{F,\{0\}}^{\overline{M}}:\overline{M}_0\overline{M}_{๐(F)}`$ for each $`\{0\}F๐`$. Under the identification $`\overline{M}_0=\mathrm{\Gamma }(X;M^+)`$ and $`\overline{M}_{๐(F)}=M_{๐(F)}=\mathrm{\Gamma }(U_F;M^+)`$, $`\phi _{F,\{0\}}^{\overline{M}}`$ is given by the restriction map $`\mathrm{\Gamma }(X;M^+)\mathrm{\Gamma }(U_F;M^+)`$. Then $`\overline{M}`$ satisfies the expected condition by the sequence (5.1). โ
###### Remark 5.6.
If $`X`$ is connected, then we have $`K[\mathrm{\Delta }]\overline{K[\mathrm{\Delta }]}`$ in the notation of Lemma 5.5. But, if $`X`$ is not connected, then $`\overline{K[\mathrm{\Delta }]}_0=K^n`$, where $`n`$ is the number of the connected components of $`X`$.
Set $`r:=dimX`$. Consider the complex
$$\omega _{K[\mathrm{\Delta }]}^{}:0\omega _{K[\mathrm{\Delta }]}^r\omega _{K[\mathrm{\Delta }]}^{r+1}\mathrm{}\omega _{K[\mathrm{\Delta }]}^0\omega _{K[\mathrm{\Delta }]}^10,$$
$$\omega _{K[\mathrm{\Delta }]}^i:=\underset{\begin{array}{c}F\mathrm{\Delta }\\ dim|F|=i\end{array}}{}K[F]$$
of squarefree $`K[\mathrm{\Delta }]`$-modules. The translated complex $`\omega _{K[\mathrm{\Delta }]}^{}[1]`$ of $`\omega _{K[\mathrm{\Delta }]}^{}`$ is (quasi-isomorphic to) the normalized dualizing complex of $`K[\mathrm{\Delta }]`$. Hence we have
$$\mathrm{Ext}_{K[\mathrm{\Delta }]}^i(M^{},\omega _{K[\mathrm{\Delta }]}^{})\mathrm{Ext}_{K[\mathrm{\Delta }]}^{i1}(M^{},\omega _{K[\mathrm{\Delta }]}^{}[1])H_๐ช^{i+1}(M^{})^{}$$
for all $`M^{}D^b(\mathrm{Sq}(\mathrm{\Delta }))`$ and all $`i`$. But here we prefer $`\omega _{K[\mathrm{\Delta }]}^{}`$ to $`\omega _{K[\mathrm{\Delta }]}^{}[1]`$, since $`(\omega _{K[\mathrm{\Delta }]}^{})^+D^b(\mathrm{Sh}(X))`$ is quasi-isomorphic to the dualizing complex $`๐_X^{}`$ of $`X`$ as shown in \[23, corollary 4.3\].
If $`M^{}D^b(\mathrm{Sq}(\mathrm{\Delta }))`$, then we have $`\mathrm{R}\mathrm{Hom}_{K[\mathrm{\Delta }]}(M^{},\omega _{K[\mathrm{\Delta }]}^{})D_{\mathrm{Sq}(\mathrm{\Delta })}^b(\mathrm{Mod}R)D^b(\mathrm{Sq}(\mathrm{\Delta }))`$. So we can define $`\mathrm{R}\mathrm{Hom}_{K[\mathrm{\Delta }]}(M^{},\omega _{K[\mathrm{\Delta }]}^{})^+D^b(\mathrm{Sh}(X))`$. Moreover, the following result holds.
###### Theorem 5.7 (\[23, Theorem 4.2\]).
For $`M^{}D^b(\mathrm{Sq}(\mathrm{\Delta }))`$, we have
$$\mathrm{R}\mathrm{Hom}_{K[\mathrm{\Delta }]}(M^{},\omega _{K[\mathrm{\Delta }]}^{})^+\mathrm{R}om((M^{})^+,๐_X^{})$$
in $`D^b(\mathrm{Sh}(X))`$. In particular, $`\mathrm{Ext}_{K[\mathrm{\Delta }]}^i(M^{},\omega _{K[\mathrm{\Delta }]}^{})^+xt^i((M^{})^+,๐_X^{})`$.
The next two results easily follow from Theorem 5.7 and the local duality theorem for $`K[\mathrm{\Delta }]`$.
###### Proposition 5.8.
If $`M`$ is a squarefree $`K[\mathrm{\Delta }]`$-module with $`dimM1`$, then
$$\mathrm{min}\{ixt^i(M^+,๐_X^{})0\}=dim\mathrm{supp}(M^+)=dimM1.$$
Here $`\mathrm{supp}(M^+):=\{xX(M^+)_x0\}`$.
###### Theorem 5.9.
Let $`M`$ be a squarefree $`K[\mathrm{\Delta }]`$-module which is regular. Set $`t:=dim\mathrm{supp}(M^+)=dimM1`$. Then $`M`$ is Buchsbaum if and only if $`xt^i(M^+,๐_X^{})=0`$ for all $`it`$. Similarly, $`M`$ is Cohen-Macaulay if and only if $`xt^i(M^+,๐_X^{})=0`$ for all $`it`$ and $`H^i(X;M^+)=0`$ for all $`i0,t`$. In particular, the Buchsbaum property and the Cohen-Macaulay property of $`M`$ depend only on the sheaf $`M^+`$.
###### Remark 5.10.
In Theorem 5.9 (and Theorem 5.12 below), even $`X`$ is somewhat superfluous, and the closure of $`\mathrm{supp}(M^+)`$ is essential. If we set $`T:=\{F\mathrm{\Delta }M_{๐(F)}0\}`$, then $`_{FT}|F|=\mathrm{supp}(M^+)`$. Let $`\mathrm{\Sigma }`$ be the order ideal of $`๐`$ generated by $`T`$. Then $`Y:=_{F\mathrm{\Sigma }}|F|`$ coincides with the closure of $`\mathrm{supp}(M^+)`$. Note that $`\mathrm{\Sigma }\mathrm{\Delta }`$ and $`M`$ can be regarded as a squarefree $`K[\mathrm{\Sigma }]`$-module. Let $`\mathrm{Sh}(Y)`$ be the sheaf associated with $`M`$ as a $`K[\mathrm{\Sigma }]`$-module. Then $`M^+|_Y`$, $`H^i(Y;)H^i(X;M^+)`$, and $`xt^i(,๐_Y^{})xt^i(M^+,๐_X^{})`$ for all $`i`$. Of course, the Cohen-Macaulay (resp. Buchsbaum) property of $`M`$ does not depend on whether we regard $`M`$ as a $`K[\mathrm{\Delta }]`$-module or a $`K[\mathrm{\Sigma }]`$-module.
Let $`\mathrm{Sh}(X)`$, and set $`๐ข^{}:=\mathrm{R}om(,๐_X^{})D^b(\mathrm{Sh}(X))`$. Recall that $`xt^i(,๐_X^{})=^i(๐ข^{})`$. More precisely, $`xt^i(,๐_X^{})`$ is the sheaf associated with the presheaf defined by $`UH^i(\mathrm{\Gamma }(U;๐ข^{}))`$ for an open subset $`U`$ of $`X`$. Hence the element of $`\mathrm{Ext}^i(,๐_X^{})=H^i(\mathrm{\Gamma }(X;๐ข^{}))`$ gives a global section of $`xt^i(,๐_X^{})`$, that is, we have a natural map $`\mathrm{Ext}^i(,๐_X^{})\mathrm{\Gamma }(X;xt^i(,๐_X^{}))`$.
###### Lemma 5.11.
If $`M\mathrm{Sq}(\mathrm{\Delta })`$, then $`[\mathrm{Ext}_{K[\mathrm{\Delta }]}^i(M,\omega _{K[\mathrm{\Delta }]}^{})]_0\mathrm{Ext}^i(M^+,๐_X^{})`$ for all $`i<0`$. Via this isomorphism, the natural map $`\mathrm{Ext}^i(M^+,๐_X^{})\mathrm{\Gamma }(X;xt^i(M^+,๐_X^{}))`$ coincides with the middle map $`[\mathrm{Ext}_{K[\mathrm{\Delta }]}^i(M,\omega _{K[\mathrm{\Delta }]}^{})]_0\mathrm{\Gamma }(X;\mathrm{Ext}_{K[\mathrm{\Delta }]}^i(M,\omega _{K[\mathrm{\Delta }]}^{})^+)`$ of the complex (5.1) in Theorem 5.2.
###### Proof.
Set $`N^{}:=\mathrm{R}\mathrm{Hom}_{K[\mathrm{\Delta }]}(M,\omega _{K[\mathrm{\Delta }]}^{})D^b(\mathrm{Sq}(\mathrm{\Delta }))`$. By the same argument as Lemma 3.8, we may assume that
$$N^i=\underset{\begin{array}{c}F\mathrm{\Delta }\\ dim|F|=i\end{array}}{}(M_{๐(F)})^{}_KK[F].$$
Thus $`(N^i)^+`$ is a direct sum of copies of the sheaf $`K[F]^+`$ for various $`F\mathrm{\Delta }`$ with $`dim|F|=i`$. While $`K[F]^+`$ is not an injective object in $`\mathrm{Sh}(X)`$, it is a constant sheaf over the closed ball $`|F|`$ and $`H^i(X;K[F]^+)=H^i(|F|;K)=0`$ for all $`i>0`$. From this fact and that $`(N^{})^+\mathrm{R}om(M^+,๐_X^{})`$, $`\mathrm{Ext}^i(M^+,๐_X^{})`$ is the $`i^{\mathrm{th}}`$ cohomology of the complex $`\mathrm{\Gamma }(X;(N^{})^+)`$. On the other hand, $`\mathrm{\Gamma }(X;(N^i)^+)=[N^i]_0`$ for all $`i0`$. Hence $`\mathrm{Ext}^i(M^+,๐_X^{})=[\mathrm{Ext}_{K[\mathrm{\Delta }]}^i(M,\omega _{K[\mathrm{\Delta }]}^{})]_0`$ for all $`i<0`$. Since $`H^i(N^{})^+xt^i(M^+,\omega _{K[\mathrm{\Delta }]}^{})`$, the assertion can be checked easily. โ
###### Theorem 5.12.
Assume that $`M\mathrm{Sq}(\mathrm{\Delta })`$ is regular. Then $`M`$ is sequentially Cohen-Macaulay if and only if the following conditions are satisfied.
* $`xt^i(xt^j(M^+,๐_X^{}),๐_X^{})=0`$ for all $`i,j`$ with $`ij`$ and $`j<0`$.
* $`H^i(X;xt^j(M^+,๐_X^{}))=0`$ for all $`i0,j`$ and all $`j<0`$.
* The natural map $`\mathrm{Ext}^i(M^+,๐_X^{})\mathrm{\Gamma }(X;xt^i(M^+,๐_X^{}))`$ is bijective for all $`i<0`$.
In particular, the sequentially Cohen-Macaulay property of $`M`$ only depends on the sheaf $`M^+`$.
###### Proof.
By argument similar to \[17, Theorem III.2.11\], we can see that $`M`$ is sequentially Cohen-Macaulay, if and only if $`\mathrm{Ext}_{K[\mathrm{\Delta }]}^i(M,\omega _{K[\mathrm{\Delta }]}^{})`$ is either the 0 module or a Cohen-Macaulay module of dimension $`1i`$ for all $`i`$. Since $`\mathrm{Ext}_{K[\mathrm{\Delta }]}^i(M,\omega _{K[\mathrm{\Delta }]}^{})=0`$ for $`i>1`$ and $`\mathrm{Ext}_{K[\mathrm{\Delta }]}^1(M,\omega _{K[\mathrm{\Delta }]}^{})`$ is an artinian module, we do not have to check $`\mathrm{Ext}_{K[\mathrm{\Delta }]}^i(M,\omega _{K[\mathrm{\Delta }]}^{})`$ for $`i1`$. Moreover, the case when $`i=0`$ is also unnecessary in our situation. In fact, since $`M`$ is regular, we have $`[H_๐ช^1(M)^{}]_0=[\mathrm{Ext}_{K[\mathrm{\Delta }]}^0(M,\omega _{K[\mathrm{\Delta }]}^{})]_0=0`$. From this fact and that $`dim\mathrm{Ext}_{K[\mathrm{\Delta }]}^0(M,\omega _{K[\mathrm{\Delta }]}^{})1`$, $`\mathrm{Ext}_{K[\mathrm{\Delta }]}^0(M,\omega _{K[\mathrm{\Delta }]}^{})`$ is Cohen-Macaulay.
If the condition (c) is satisfied, $`\mathrm{Ext}_{K[\mathrm{\Delta }]}^i(M,\omega _{K[\mathrm{\Delta }]}^{})`$ is regular by Lemma 5.11. Conversely, if $`M`$ is sequentially Cohen-Macaulay, then $`\mathrm{Ext}_{K[\mathrm{\Delta }]}^i(M,\omega _{K[\mathrm{\Delta }]}^{})`$ must be regular for $`i<0`$ and (c) is satisfied.
Since $`\mathrm{Ext}_{K[\mathrm{\Delta }]}^i(M,\omega _{K[\mathrm{\Delta }]}^{})^+xt^i(M^+,๐_X^{})`$, the assertion follows from the above observation and Theorem 5.9 (and Proposition 5.8). โ
###### Remark 5.13.
If $`X`$ is not connected, then $`K[\mathrm{\Delta }]`$ is not regular as a squarefree module. So Theorem 5.12 does not imply Theorem 4.7 directly. But, by the following observation, Theorem 4.7 can be reduced to Theorem 5.12.
* 0 dimensional components of $`X`$ (i.e., 1 dimensional components of $`K[\mathrm{\Delta }]`$) are irrelevant to the sequentially Cohen-Macaulay property of $`K[\mathrm{\Delta }]`$. So we can remove them.
* If $`X`$ does not have a 0 dimensional component and $`K[\mathrm{\Delta }]`$ is sequentially Cohen-Macaulay, then $`X`$ is connected.
## 6. Ideals whose radicals are monomial ideals
In this brief section, we generalize a result of Herzog, Takayama and Terai . Let $`IR`$ be a (non-monomial) ideal. Even if $`R/I`$ is Cohen-Macaulay, $`R/\sqrt{I}`$ is not Cohen-Macaulay in general. See the introduction of for an explicit example. But the next theorem states that if $`\sqrt{I}`$ is a monomial ideal then such an example does not exist. When $`R`$ is a polynomial ring, this result was obtained in .
###### Theorem 6.1 (c.f. \[9, Theorem 2.6\]).
Assume that $`R`$ is normal. Let $`I`$ be a (not necessarily graded) ideal whose radical $`\sqrt{I}`$ is a monomial ideal. Then we have the following
(1)If $`R/I`$ is Cohen-Macaulay (more generally, the localization $`(R/I)_๐ช`$ is Cohen-Macaulay), then $`R/\sqrt{I}`$ is also.
(2) If the localization $`(R/I)_๐ช`$ is generalized Cohen-Macaulay, then $`R/\sqrt{I}`$ is Buchsbaum, in particular, it is generalized Cohen-Macaulay again.
The idea of the proof is same as the one given in \[9, Remark 2.7\].
###### Proof.
Let me introduce the notation and facts used throughout this proof. Set $`p:=dim(R/I)=dim(R/\sqrt{I})`$. Since $`\sqrt{I}`$ is a monomial ideal, $`\mathrm{Ext}_R^i(R/\sqrt{I},\omega _R)`$ has a natural $`^d`$-grading. In particular, $`\mathrm{Ext}_R^i(R/\sqrt{I},\omega _R)_RR_๐ช=0`$ implies $`\mathrm{Ext}_R^i(R/\sqrt{I},\omega _R)=0`$. Similarly, if $`\mathrm{Ext}_R^i(R/\sqrt{I},\omega _R)_RR_๐ช`$ has finite length, then so does $`\mathrm{Ext}_R^i(R/\sqrt{I},\omega _R)`$.
(1) It suffices to show that $`\mathrm{Ext}_R^i(R/\sqrt{I},\omega _R)=0`$ for all $`idp`$. Recall that the natural map
$$\mathrm{Ext}_R^i(R/\sqrt{I},\omega _R)H_\sqrt{I}^i(\omega _R)$$
factors through the map $`f:\mathrm{Ext}_R^i(R/\sqrt{I},\omega _R)\mathrm{Ext}_R^i(R/I,\omega _R)`$ induced by the surjection $`R/IR/\sqrt{I}`$. But the map $`\mathrm{Ext}_R^i(R/\sqrt{I},\omega _R)H_\sqrt{I}^i(\omega _R)`$ is an injection by \[22, Theorem 5.9\]. So the map $`f`$ and its localization $`f_RR_๐ช`$ are injective. Since $`(R/I)_๐ช`$ is Cohen-Macaulay, we have $`\mathrm{Ext}_R^i(R/I,\omega _R)_RR_๐ช=0`$ for all $`idp`$. So we are done.
(2) By the assumption, $`\mathrm{Ext}_R^i(R/I,\omega _R)_RR_๐ชH_๐ช^{di}((R/I)_๐ช)^{}`$ has finite length for all $`idp`$. Recall that the map $`f:\mathrm{Ext}_R^i(R/\sqrt{I},\omega _R)\mathrm{Ext}_R^i(R/I,\omega _R)`$ and the localization $`fR_๐ช`$ are injective. Hence if $`\mathrm{Ext}_R^i(R/\sqrt{I},\omega _R)`$ does not have finite length, then $`\mathrm{Ext}_R^i(R/I,\omega _R)_RR_๐ช`$ also. So $`\mathrm{Ext}_R^i(R/\sqrt{I},\omega _R)`$ has finite length for all $`idp`$. By \[23, Corollary 4.6\], $`R/\sqrt{I}`$ is Buchsbaum. โ
## 7. Local cohomology modules of finite length
Let $`S=K[x_1,\mathrm{},x_d]`$ be a polynomial ring, and $`IS`$ a monomial ideal. Recently, Takayama gave an interesting observation that the range
$$\{๐^d[H_๐ช^i(S/I)]_๐0\}$$
is controlled by the degrees of minimal generators of $`I`$ (especially, when $`H_๐ช^i(S/I)`$ has finite length). For this result, he used a combinatorial formula on $`[H_๐ช^i(S/I)]_๐`$. But there is an easy and conceptual proof, and we can generalize his result to finitely generated $`^d`$-graded $`S`$-modules. We regard $`^d`$ as a partially ordered set with $`๐๐`$ $`\stackrel{\text{def}}{}`$ $`a_ib_i`$ for all $`i`$. Here $`๐=(a_1,\mathrm{},a_d)`$ and $`๐=(b_1,\mathrm{},b_d)`$.
###### Definition 7.1 (Miller ).
Let $`๐^d`$. We say an $`S`$-module $`M`$ is positively $`๐`$ determined (โ$`๐`$-p.d.โ, for short) if $`M`$ is $`^d`$-graded, finitely generated, and the multiplication map $`M_๐yx_iyM_{๐+๐_i}`$ is bijective for all $`๐^d`$ and all $`1id`$ with $`b_ia_i`$. Here $`๐_i=(0,\mathrm{},1,\mathrm{},0)^d`$ is the vector with 1 at the $`i^{\mathrm{th}}`$ position.
###### Remark 7.2.
(1) Any finitely generated $`^d`$-graded $`S`$-module is an $`๐`$-p.d.module for sufficiently large $`๐^d`$.
(2) Set $`\mathrm{๐}:=(1,\mathrm{},1)^d`$. Then $`\mathrm{๐}`$-p.d.modules are nothing other than squarefree modules (recall Definition 3.1) over $`S`$.
(3) Let $`I=(๐ฑ^{๐_1},\mathrm{},๐ฑ^{๐_r})`$ be a monomial ideal in $`S`$ with $`๐_j=(b_1^j,\mathrm{},b_d^j)`$. Set $`a_i:=\mathrm{max}\{b_i^j1jr\}`$ for each $`1id`$, and set $`๐:=(a_1,\mathrm{},a_d)`$ (in other words, $`๐ฑ^๐=\mathrm{lcm}(๐ฑ^{๐_1},\mathrm{},๐ฑ^{๐_r})`$). Then $`I`$ and $`S/I`$ are $`๐`$-p.d.modules.
As shown in , $`๐`$-p.d.modules enjoy many interesting properties. Some of them are used in the proof of the next result.
###### Proposition 7.3 (c.f. \[20, Corollary 2\]).
Let $`M`$ be an $`S`$-module which is $`๐`$-p.d. If $`[H_๐ช^i(M)]_๐0`$ for some $`i`$, then $`๐๐\mathrm{๐}`$.
###### Proof.
Let $`P^{}`$ be a $`^d`$-graded minimal free resolution of $`M`$. Then each $`P^i`$ is an $`๐`$-p.d.module again. That is, if $`S(๐)`$ appears in $`P^i`$ as a direct summand, then $`\mathrm{๐}๐๐`$. Let $`\omega _S=S(\mathrm{๐})`$ be the canonical module of $`S`$. If $`S(๐)`$ is $`๐`$-p.d., then
$$\mathrm{Hom}_S(S(๐),\omega _S)(๐+\mathrm{๐})\mathrm{Hom}_S(S(๐),S(๐))S(๐+๐)$$
is also. Hence $`[\mathrm{Hom}_S(P^{},\omega _S)](๐+\mathrm{๐})`$ is a complex of $`๐`$-p.d.modules, and its $`i^{\mathrm{th}}`$ cohomology $`[\mathrm{Ext}_S^i(M,\omega _S)](๐+\mathrm{๐})`$ is also $`๐`$-p.d. Since $`H_๐ช^{di}(M)_๐[\mathrm{Ext}_S^i(M,\omega _S)]_๐^{}`$ by the local duality, we are done. โ
Since $`[\mathrm{Ext}_S^i(M,\omega _S)](๐+\mathrm{๐})`$ is an $`๐`$-p.d.module as shown in the proof of Proposition 7.3, we have the following.
###### Proposition 7.4 (c.f. \[20, Proposition 1\]).
Let $`M`$ be a finitely generated $`^d`$-graded $`S`$-module. Then $`H_๐ช^i(M)`$ has finite length if and only if $`H_๐ช^i(M)`$ is $`^d`$-graded.
To extend Proposition 7.4 to semigroup rings, we have to introduce the notion $`\mathrm{supp}_+(u)`$ for $`u^d`$. There are (the โdefining equationsโ of) hyperplanes $`h_1,\mathrm{},h_t(^d)^{}`$ such that $`๐(=_0C)=\{u^dh_i(u)0\text{ for all }i\text{ }\}`$. We may assume that $`h_1,\mathrm{},h_t`$ form a minimal system defining $`๐`$ (equivalently, the number of $`d1`$ dimensional faces of $`๐`$ is $`t`$). For $`u^d`$, set
$$\mathrm{supp}_+(u):=\{ih_i(u)>0\}\{1,\mathrm{},t\}.$$
For $`u,v๐`$, $`\mathrm{supp}_+(u)=\mathrm{supp}_+(v)`$ if and only if $`s(u)=s(v)`$.
So let be $`h_1,\mathrm{},h_t(^d)^{}`$ the โdefining equationsโ of the cone $`๐^d`$. Recall that, for $`u^d`$, we have $`\mathrm{supp}_+(u)=\{ih_i(u)>0\}\{1,\mathrm{},t\}`$. Set $`\overline{C}:=^d๐`$. Note that $`R`$ is normal if and only if $`C=\overline{C}`$. It is easy to see that, for $`๐^d`$, $`\mathrm{supp}_+(๐)=\mathrm{}`$ if and only if $`๐\overline{C}`$. We say $`M\mathrm{Mod}R`$ is $`\overline{C}`$-graded, if $`M_๐=0`$ for all $`๐\overline{C}`$. Clearly, a $`C`$-graded module is always $`\overline{C}`$-graded, and the converse is true if $`R`$ is normal. We also set
$$\mathrm{\Psi }_C:=\{๐^d\mathrm{supp}_+(๐)\mathrm{supp}_+(๐)\text{ for some }0๐C\}.$$
###### Theorem 7.5.
Assume that $`M\mathrm{mod}R`$ is $`\overline{C}`$-graded. Then $`H_๐ช^i(M)`$ has finite length if and only if $`[H_๐ช^i(M)]_๐=0`$ for all $`๐\mathrm{\Psi }_C`$.
###### Corollary 7.6.
Assume that $`R`$ is simplicial and $`M\mathrm{mod}R`$ is a $`\overline{C}`$-graded module. Then $`H_๐ช^i(M)`$ has finite length if and only if $`H_๐ช^i(M)`$ is $`\overline{C}`$-graded.
###### Proof.
Since $`R`$ is simplicial, we have
$$\mathrm{\Psi }_C=\{๐^d\mathrm{supp}_+(๐)\mathrm{}\}=\{๐^d๐\overline{C}\}=^d\overline{C}.$$
So the assertion follows from Theorem 7.5. โ
###### Example 7.7.
If $`C`$ is not simplicial, Corollary 7.6 does not hold. Let $`R`$ be the ring given in Example 2.5 (1), and set $`M=R/(x^3,x^2y^2,y^3)`$. Then computation by Macaulay 2 shows that $`H_๐ช^1(M)`$ has finite length and $`H_๐ช^1(M)_{(2,0,0)}0`$.
To prove the theorem, we consider the following condition for a module $`N\mathrm{Mod}R`$.
* If $`๐^d`$ and $`0๐C`$ satisfy $`\mathrm{supp}_+(๐)\mathrm{supp}_+(๐)`$, then the multiplication map $`N_{๐๐}y๐ฑ^๐yN_๐`$ is bijective.
###### Lemma 7.8.
Assume that a $`^d`$graded $`R`$-module $`N`$ has finite length and satisfies the condition $`()`$. Then $`N_๐=0`$ for all $`๐\mathrm{\Psi }_C`$.
###### Proof.
If $`๐\mathrm{\Psi }_C`$, then we can take $`0๐C`$ such that $`\mathrm{supp}_+(๐)\mathrm{supp}_+(๐)`$. Since $`\mathrm{supp}_+(๐+n๐)=\mathrm{supp}_+(๐)`$ for all $`n0`$, we have $`N_๐N_{๐๐}N_{๐2๐}\mathrm{}`$ by the condition $`()`$. So $`N_๐`$ must be 0. โ
###### Lemma 7.9.
The full subcategory of $`\mathrm{Mod}R`$ consisting of modules satisfying $`()`$ is closed under kernel and cokernels.
###### Proof.
Follows from the five-lemma. โ
Proof of Theorem 7.5. (Necessity) By Lemma 7.8, it suffices to prove that $`H_๐ช^i(M)`$ satisfies the condition $`()`$. To show this, we use the โฤech complexโ $`L^{}`$ of \[1, ยง6.2\] (see also \[15, Chapter 13\]. There this complex is called โIshidaโs complexโ), which satisfies $`H_๐ช^i(M)H^i(M_RL^{})`$ for all $`i`$. For a face $`F๐`$, $`R_F`$ denotes the localization of $`R`$ at the multiplicatively closed set $`\{๐ฑ^๐๐FC\}`$. Then we have $`L^i=_{dimF=i}R_F`$ for each $`i`$.
By Lemma 7.9, it suffices to prove that $`M_F:=M_RR_F`$ satisfies $`()`$ for all $`F๐`$. Of course, $`M_F`$ is the localization of $`M`$ at $`\{๐ฑ^๐๐FC\}`$. We may assume that $`F=\{u^dh_1(u)=h_2(u)=\mathrm{}=h_n(u)=0\}`$.
Let $`๐\mathrm{\Psi }_C`$, and take $`0๐C`$ with $`\mathrm{supp}_+(๐)\mathrm{supp}_+(๐)`$. If $`๐F`$, $`M_F`$ is a module over $`R[๐ฑ^๐]`$. Thus the multiplication by $`๐ฑ^๐`$ gives a bijection $`(M_F)_{๐๐}(M_F)_๐`$. So we may assume that $`๐F`$. If $`[M_F]_๐0`$, then we have $`๐=๐๐`$ for some $`๐\overline{C}`$ and $`๐F^d`$. Thus $`\mathrm{supp}_+(๐)\mathrm{supp}_+(๐)\{n+1,n+2,\mathrm{},t\}`$. On the other hand, since $`๐F`$, we have $`i\mathrm{supp}_+(๐)\mathrm{supp}_+(๐)`$ for some $`in`$. This is a contradiction. Hence $`[M_F]_๐=0`$. Similarly, we can see that $`[M_F]_{๐๐}=0`$. So any map $`[M_F]_{๐๐}[M_F]_๐`$ is bijective.
(Sufficiency) By the local duality, the graded Matlis dual $`H_๐ช^i(M)^{}`$ of $`H_๐ช^i(M)`$ is finitely generated. For any $`๐^d`$ and $`0๐C`$, we have $`(๐+n๐)\mathrm{\Phi }_C`$ for sufficiently large $`n`$. By the assumption, if $`[H_๐ช^i(M)^{}]_๐0`$ then $`๐\mathrm{\Psi }_C`$. Hence $`๐ช^nH_๐ช^i(M)=0`$ for sufficiently large $`n`$, and $`H_๐ช^i(M)`$ has finite length. โ
## Acknowledgments
The author is grateful to Professor Ngo Viet Trung for telling him the ring given in Example 2.5 (2). He also thanks Professors Ezra Miller and Tim Rรถmer for careful reading and comments on an earlier version of this paper, and Professor Mitsuyasu Hashimoto for useful comment around Lemma 5.11.
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# The Vertical Structure of Stars in Edge-on Galaxies
## 1. Introduction
The age, morphology and composition of the disk, bulge and halo components found in typical spiral galaxies all provide significant constraints on the mechanisms of galaxy formation. In the Milky Way, the disk of the galaxy has been found to be a complex object, with a thin disk whose scale height depends on stellar age and a separate thick disk component first discovered by Gilmore & Reid (1983) using star counts. The Milky Way thick disk is dominated by old, moderately metal-poor stars (Wyse & Gilmore 1995) and has recently been shown to be chemically distinct from the thin disk (e.g. Bensby et al. 2005).
Outside the Milky Way, very little detailed information exists about disk morphologies in other galaxies. Integrated surface brightness profiles of edge-on spirals have shown that many galaxies appear to have thick disks (e.g. Burstein 1979; Dalcanton & Bernstein 2002; Pohlen et al. 2004), while recent studies in a handful of nearby edge-on galaxies using resolved stars have confirmed that extraplanar stars are old and moderately metal-poor (Mould 2005; Tikhonov et al. 2005).
Many different mechanisms have been used to explain the creation of thin and thick disks. These mechanisms fall into three categories: (1) Creation of thicker components by dynamical heating from a thin disk by molecular clouds, spiral arms, galaxy interaction etc. (e.g. Spitzer & Schwarzschild 1951; Lacey 1991; Gnedin 2003). (2) the slow collapse of the proto-Galaxy forming the thick and thin disk (e.g. Eggen et al. 1962; Gilmore 1984), and (3) the formation of a thick disk from mergers by direct accretion of stars or by in situ formation from accreted gas (e.g. Bekki & Chiba 2001; Gilmore et al. 2002; Abadi et al. 2003; Brook et al. 2004). Testing the applicability of these theories beyond the Milky Way requires detailed observations of galaxies with a range of masses and types. In this paper we focus on the disks of edge-on, low mass, late type spiral galaxies.
The work presented here is a summary of our recently published results, Seth et al. (2005a) (Paper I) and Seth et al. (2005b) (Paper II). A more in depth treatment of this work is presented there. In ยง2 we present an overview of our galaxy sample and describe out HST/ACS observations. We present our primary results on the vertical distribution of stars and the metallicity of extraplanar stars in our sample galaxies in ยง3. This is followed by a discussion in ยง4.
## 2. Observations and Sample
We obtained HST/ACS images of 16 edge-on ($`i>80^{}`$), late type galaxies as part of a Cycle 12 snapshot proposal. A subsample of eight fields in six galaxies were near enough to enable distance measurements using the tip of the red giant branch (RGB). These six galaxies and their distances, morphological types, circular velocities, $`K_s`$-band scale lengths and scale heights<sup>1</sup><sup>1</sup>1We use a functional form $`\mathrm{\Sigma }(z)sech^2(z/z_0)`$ to fit our surface density profiles. Scale heights refer to the $`z_0`$ in this equation. as determined from fits to 2MASS data (Paper I), and the number of ACS fields/pointings are listed in Table 1. Because these galaxies are nearby, the ACS field-of-view does not cover the entire galaxy. For most galaxies the fields are roughly centered on the galaxy, however, for those galaxies with two fields (NGC 55 and NGC 4631), the second field (denoted with a โ-DISKโ suffix) is located further out in the disk. The six galaxies are all late type (Sc+), none have a significant bulge, and all have luminosities and inferred masses significantly lower than the Milky Way.
The ACS observations consisted of 676 seconds of exposure time in the F606W filter and 700 seconds in the F814W filter. Our photometric pipeline is described in detail in Paper I and yielded 40,000 to 280,000 stars per field. Extensive artificial star tests ($``$5 million stars per field) enabled determination of the completeness as a function of magnitude, color and local surface brightness allowing us to study the number of stars as a function of distance from the disk midplane (โdisk heightโ). To enable adequate correction of the number counts of stars, we restrict all analysis presented here to stars brighter than the 20% completeness level in the highest surface brightness areas (i.e. the midplane).
Figure 1 shows the color-magnitude diagrams (CMDs) for each of the fields analyzed here. In the nearest galaxies, the red giant branch (RGB), Asymptotic Giant Branch (AGB), upper Main Sequence (MS) and red and blue Helium burning star sequences are clearly visible. We use these different stellar populations to separate the stars we see into three broad age categories. Three boxes are shown on each CMD in Fig. 1 for the RGB, AGB, and MS (including helium burning) stars. To determine how well we could isolate different stellar ages we generated synthetic CMDs using the MATCH code (Dolphin 2002) and assuming a constant star formation rate starting 13 Gyr ago. The synthetic CMD was then convolved with the errors obtained from the artificial star tests for each galaxy (See Paper II for more details). The resulting age distribution for stars falling in the MS, AGB, and RGB boxes in NGC 4144 is shown in Figure 2. From this figure it is clear that the MS, AGB, and RGB boxes roughly correspond to young ($`100`$ Myr), intermediage age ($``$1 Gyr), and old ($`10`$ Gyr) stellar populations. However there is some overlap between the different boxes, particularly between the AGB and RGB (see Paper II for more details).
## 3. Results
In this section we present the main results of our work. This includes an analysis of the scale heights of the resolved stellar populations and a determination of the metallicity of the old extraplanar stellar population.
### 3.1. Vertical Distribution of Stars
Before dividing up the stellar population in different age bins, we first determine the vertical distribution of all the observed stars above the 20% completeness limit. Figure 3 shows the vertical stellar density profile for each field in our sample. This figure was created by determining the number of stars in each galaxy in bins at different disk heights. This number of stars was then corrected for incompleteness and was divided by the image area in each bin to obtain a surface density. The two dashed lines denote the profiles from the two NGC 4631 fields on the side where they are contaminated by stars from nearby companion NGC 4627. The disk heights on the x-axis are scaled by $`z_{1/2}`$, the disk height containing 50% of the $`K_s`$ band light ($`z_{1/2}=0.5493z_0`$) in each galaxy. The two dashed lines show the profile expected for a stellar population with a scale height 1 and 2 times the $`K_s`$ band light.
The most notable result from Figure 3 is that we are able to trace the resolved stellar population far above the midplane, out to 15-20 $`z_{1/2}`$ ($``$3-4 kpc). The fact that the star counts continue to decrease all the way to the edge of the images suggests there are few background and foreground contaminants. Assuming a stellar population at large scale heights similar to Galactic globular clusters (using data taken from Buonanno et al. 1994; Kravtsov et al. 1997), we find that we reach a limiting surface brightness of $``$28 mag arcsec<sup>-2</sup>. We determine the scale height of the full stellar population at large disk heights by fitting the profiles between 5 and 10 $`z_{1/2}`$ above and below the disk to a sech<sup>2</sup> function. The derived scale heights listed in Table 2 (in columns 3 & 4) are typically $``$2 times the 2MASS $`K_s`$ band scale heights. We note that the 2MASS images are not very deep and only trace the highest surface brightness population near the midplane. This therefore suggests that the population at large disk heights has a larger scale height than the population near the midplane.
We now analyze the scale height of our sample of galaxies as a function of stellar age. Using the MS, AGB, and RGB boxes shown in Figure 1, we are able to isolate young, intermediate, and old stellar populations. Figure 4 shows the completeness corrected surface density profiles for the three different stellar population boxes in each field. The youngest MS stars are shown with a solid line, while the intermediate-age AGB and old RGB stars are shown with dotted and dashed lines, respectively. The stellar density profiles have all been normalized to integrate to one. Each field clearly shows that the MS component is the narrowest while the RGB is in all cases the thickest component. We quantified this trend by fitting a sech<sup>2</sup> function to each component excluding the central portions of the profile from the fit. The central portions of the profiles show dips in the star counts which likely result from dust absorption (see Paper II). The scale heights found from the sech<sup>2</sup> fits are shown in columns 5-7 of Table 2. A number of interesting trends are visible. First, in each galaxy/field, the MS component has the smallest scale height with the AGB and RGB having succesively larger scale heights. In other words, the scale height increases with increasing stellar age in all six galaxies. Second, the scale heights of the MS component are larger than (200-500 pc) than the scale height of young stars in the Milky Way ($``$200 pc; Schmidt 1963) despite being less luminous and having shorter scale lengths. This thicker MS component may be related to the thicker dust distribution seen in low mass galaxies by Dalcanton et al. (2004). Finally, the axial ratio of the RGB component ($`h_r/z_{0,\mathrm{RGB}}`$) shown in the last column of Table 2 ranges from 1.0 to 3.3. This axial ratio is similar to the Milky Way thick disk (using data from e.g Chen et al. 2001).
### 3.2. Metallicity
Based on the scale height results, we expect most of the stars at large disk heights to be old. The morphology of the RGB population can provide a metallicity diagnostic for old stellar populations (e.g. Da Costa & Armandroff 1990; Armandroff et al. 1993; Frayn & Gilmore 2002). In this section we use the color of the RGB stars to determine the metallicity of the extraplanar stellar population.
The left panel of Figure 5 shows the composite CMD of stars above 4$`z_{1/2}`$ in all the galaxies excepting the NGC 4631 fields (which have completeness limits that are too high for RGB metallicity analysis). Overplotted are 10 Gyr Padova isochrones at \[Fe/H\] values ranging between -2.3 and 0.0. The peak of the distribution lies between the isochrones with \[Fe/H\] of -1.3 and -0.7 indicating that the dominant population in these galaxies at high latitudes is moderately metal-poor. We determined a metallicity distribution function for each galaxy using three bins in absolute magnitude shown with dotted lines in the left panel of Figure 5. In each bin we determined the metallicity of individual stars by linearly interpolating between the 10 Gyr Padova isochrones. The right panel of Figure 5 shows the metallicity distribution functions for NGC 4244, with the shaded grey region indicating the range of results for the three different bins. The peak metallicity is at \[Fe/H\] of -0.9 with a long tail towards lower metallicities and a sharper cutoff at higher metallicities. The other four galaxies have similar MDFs, with peaks ranging between \[Fe/H\] of -0.7 and -1.1. This peak metallicity also agrees with metallicities found for extraplanar stars in other studies of galaxies in the same mass range (Cole et al. 2000; Davidge 2003; Tiede et al. 2004; Davidge 2005). This metallicity is slightly more metal-poor than the Milky Way thick disk metallicity (\[Fe/H\]$``$-0.6) but significantly more metal-rich than the Milky Way halo ($`[\mathrm{Fe}/\mathrm{H}]`$-1.7 Wyse & Gilmore 1995). We note that the current metallicity of these galaxies is somewhat subsolar (Garnett 2002), furthering the idea that the peak metallicites of \[Fe/H\]$``$-1 are consistent with the presence of a thick disk similar to the Milky Wayโs.
Overplotted in the right panel of Figure 5 is a closed-box โsimpleโ chemical evolution model (Eq. 20 of Pagel 1997) scaled to the peak of the MDF. The overall shape is similar to observed MDF, however our data appears to have a sharper cutoff at the high metallicity end, and a deficit of stars at the low metallicity end. The high metallicity cutoff suggests that the star formation trucated before all the gas was consumed, as would be expected for any of the thick disk scenarios discussed in the introduction. At low metallicities, the deficit of stars may be a manifestation of the G-dwarf problem. However, our results at the low metallicity end are quite uncertain since the close spacing of the isochrones means photometric errors introduce large uncertainties in the metallicity.
Models for the origin of thin and thick disks predict different degrees of variation in the stellar metallicity with height above the plane. We can observe any such trend in our galaxies by examining the median color of the RGB stars as a function of height above the midplane. Figure 6 shows the median RGB color vs. disk height for each of the galaxies in our sample (except NGC 4631). At low disk heights, internal reddening may impact the color of stars - this is shown as the hatched region. Above this height, from $``$3-10$`z_{1/2}`$ ($``$1-2 kpc), the galaxies show very little variation in RGB color suggesting nearly uniform metallicity with increasing disk height. Recent observations by (Davidge 2005; Tikhonov et al. 2005; Mould 2005) of extraplanar stars in similar galaxies also find a lack of strong metallicity gradients. We discuss the implications of this further in the next section.
## 4. Discussion and Conclusions
We have shown above that in our sample of six edge-on, low mass, late type galaxies:
* Stars exist at large heights above the disk plane.
* The scale height of a stellar population increases with age in each galaxy. The young (MS) stellar population has a larger scale height than young stars in the Milky Way, while the old RGB population has an axial ratio similar to the Milky Way thick disk.
* Extraplanar RGB stars have a peak metallicity of \[Fe/H\]$``$-1. In addition there appears to be little gradient in the metallicity at moderate (1-2 kpc) disk heights.
In the Milky Way, it is widely accepted that the age-dependent scale height in the thin disk is a result of slow disk heating, while the thick disk is thought to be formed by some type of merger event. We now consider our results in light of this scenario.
Disk Heating: The steady increase in scale height with stellar population age suggests that some kind of disk heating is occurring in our sample galaxies. We can use our measured scale heights to constrain the rate of disk heating in these galaxies. For an isothermal sech<sup>2</sup> profile, the scale height of a stellar population is proportional to the square of the vertical velocity dispersion ($`z_0\sigma _z^2`$). In the Milky Way, the rate of disk heating is usually expressed as a power law, $`\sigma _zt^\alpha `$, with $`\alpha `$ ranging between 0.3 and 0.6 (Hรคnninen & Flynn 2002). Using the ages for our stellar population boxes derived from the constant star-formation rate synthetic CMDs, we find that the rate of disk heating is much lower, with $`\alpha `$ being no larger than 0.15 (see Paper II for more details). This is not surprising given our galaxies likely have fewer molecular clouds and weaker spiral arms. Based on this slow heating rate, it is plausible that disk heating can account for the observed variation of scale height. However, a model in which steady disk heating accounts for all of the extraplanar stars would predict the existence of intermediate-age RGB stars with higher metallicities and lower scale heights leading to a decreasing metallicity with increasing disk height. The lack of such a trend in our galaxies suggests that a majority of RGB stars at all disk heights formed early and with a well-mixed metallicity distribution. This strongly suggests that mergers or interactions played a role in the formation of these RGB stars.
Merger & Interactions: The thick disk in the Milky Way is thought to have formed in a merger event (Abadi et al. 2003; Brook et al. 2004, e.g.), and recent dynamical observations of a counter-rotating thick disk in FGC 227 by Yoachim & Dalcanton (2005) also strongly suggests a merger origin. In our sample galaxies, the lack of metallicity gradient and the similarity of the axial ratio and metallicity of the RGB population to the Milky Way thick disk suggests that mergers or interactions likely have played an important role in the early history of these galaxies. Further evidence to support this claim comes from the heterogeneity of scale heights we observe. Although all the galaxies in our sample show a similar trend towards increasing scale height with age, in some galaxies the RGB component is significantly thicker than in others despite having similar circular velocity/mass. This variation would be a natural consequence of the stochastic merging process.
Overall, our results require that some disk heating occurs at intermediate ages (to puff up the AGB component), but that events at earlier times (interactions or mergers) created a majority of the RGB stars over a short period of time. In future work, we will focus on the globular cluster systems of our sample galaxies which will provide us with additional evidence on their history and assembly.
#### Acknowledgments.
The authors would like to thank Andrew Dolphin and Antonio Aparicio for their help in generating synthetic CMDs and Leo Girardi for supplying us with isochrones. This work was supported by HST-GO-09765, the Sloan foundation, and NSF Grant CAREER AST-0238683.
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# An extremal problem related to negative refraction
## 1. Introduction
We consider in this note an extremal problem that arose in investigations of certain electromagnetic parameters of artificial materials (metamaterials). The physical interpretation of our solution in terms of bounds for refractive indices and theoretical limitations for the design of metamaterials is described elsewhere ; the purpose of the present work is to give an account of the underlying mathematical problem, which seems to be of some independent interest.
We will be dealing with Hardy spaces $`H^p`$ of the upper half-plane $`\{z=x+iy:y>0\}`$. We are primarily interested in $`H^2`$, but it is convenient to have at our disposal the whole range of spaces corresponding to $`0<p\mathrm{}`$. For $`0<p<\mathrm{}`$, $`H^p`$ consists of those analytic functions $`f`$ in the upper half-plane for which
$$f_p^p=\underset{y>0}{sup}_{\mathrm{}}^{\mathrm{}}|f(x+iy)|^p๐x<\mathrm{};$$
$`H^{\mathrm{}}`$ is the space of bounded analytic functions. A function $`f`$ in $`H^p`$ has a nontangential boundary limit at almost every point of the real axis, and the corresponding limit function, also denoted $`f`$, is in $`L^p=L^p()`$. Indeed, the $`L^p`$ norm of the boundary limit function coincides with the $`H^p`$ norm introduced above. Thus we may view $`H^p`$ as a subspace of $`L^p`$. We refer to for these and other basic facts about $`H^p`$, as well as the twin theory of $`H^p`$ of the unit disk. (We will make a reference to the disk setting at one point.)
The Hilbert space $`H^2`$ is the image of $`L^2(^+)`$ under the Fourier transform. In practice, it is quite common that one considers functions in $`H^2`$ that are Fourier transforms of real-valued functions in $`L^2`$. This leads to the following symmetry condition: $`f(x)=\overline{f(x)}`$. Functions $`f`$ satisfying this condition will be referred to as Hermitian functions. Thus Hermitian functions have even real parts and odd imaginary parts.
The Hilbert transform of a function $`u`$ in $`L^p`$ ($`1p<\mathrm{}`$) is defined as
$$\stackrel{~}{u}(x)=\text{p.v.}\frac{1}{\pi }_{\mathrm{}}^{\mathrm{}}\frac{u(t)}{xt}๐t.$$
It acts boundedly on $`L^p`$ for $`1<p<\mathrm{}`$ and isometrically on $`L^2`$. If $`u`$ is a real-valued function in $`L^p`$ for $`1<p<\mathrm{}`$, then $`u+i\stackrel{~}{u}`$ is in $`H^p`$, and so the role of the Hilbert transform is to link the real and imaginary parts of functions in $`H^p`$. We will only work with Hermitian functions, and we will be interested in computing real parts from imaginary parts. For this reason, it will be convenient for us to consider the following Hilbert operator:
$$v(x)=\text{p.v.}\frac{1}{\pi }_0^{\mathrm{}}v(t)\left(\frac{1}{tx}+\frac{1}{t+x}\right)๐t,$$
acting on functions in $`L^p(^+)`$. Provided $`1<p<\mathrm{}`$, the function $`v+iv`$ will then be in $`H^p`$, with the presumption that $`v`$ is an odd function.
A natural type of problem in $`H^2`$ is that of approximating a given Hermitian function supported on two symmetric intervals. This means that one specifies a desired behavior in a certain frequency band and attempts to find, in some appropriate sense, an optimal approximation in $`H^2`$. Without further restrictions, such a problem makes little sense, because it is easy to see that approximations can be made with arbitrary precision in $`L^2`$ norm. It may be reasonable to prescribe bounds for the norm of the approximating function; see for instance the work of M. G. Kreฤญn and P. Ya. Nudelโman , for interesting results along such lines. In the present note, we take a different route. We shall require the imaginary part of the function to be nonnegative for nonnegative arguments.<sup>1</sup><sup>1</sup>1This reflects the passivity condition for our electromagnetic medium. We are interested in a specific problem of this general kind; it turns out to have an explicit and remarkably simple solution.
## 2. Results
We turn to the statement of the problem. For a finite interval $`I=[a,b]`$ ($`0<a<b`$) and every real number $`\alpha `$ we define the family of functions
$$K_\alpha (I)=\{vL^2(^+):v(t)0\text{for}t>0,v(t)=\alpha \text{for}tI\}.$$
(Here and elsewhere we suppress the obvious โalmost everywhereโ provisions needed when considering pointwise restrictions.) We think of functions in $`K_\alpha (I)`$, or more generally functions in $`L^2(^+)`$, as the imaginary parts of Hermitian functions, and we view them therefore as odd functions on $``$.
Our purpose is to give a parametrization of $`K_\alpha (I)`$ and to solve the extremal problem
$$\lambda =\underset{vK_1(I)}{inf}\chi _Iv_{\mathrm{}},$$
where $`\chi _I`$ denotes the characteristic function of $`I`$. We will show that the extremal problem has the following explicit solution:
$$\lambda =\frac{b^2a^2}{2ab}.$$
We note that the quantity on the right is invariant under dilations $`sI=[sa,sb]`$. This is as it should be since $`v(t)`$ is in $`K_1(I)`$ if and only if $`v(t/s)`$ is in $`K_1(sI)`$ for $`s>0`$. It will become clear that the extremal value $`\lambda `$ is not attained by any function in $`K_1(I)`$. We will also see that the problem is insensitive to which $`L^p`$ norm we choose to minimize.
Clearly, the corresponding extremal problem for $`K_\alpha (I)`$ has solution $`|\alpha |\lambda `$ when $`\alpha <0`$. However, if $`\alpha 0`$, the extremal problem is uninteresting and has solution $`0`$. Thus the sign in the relation $`v(t)=\alpha `$ for $`tI`$ matters in a decisive way.<sup>2</sup><sup>2</sup>2The case $`\alpha <1`$ corresponds to the interesting physical phenomenon of negative refraction, which has received considerable attention in recent years. Artificial, negatively refracting materials, called metamaterials, have been realized in the microwave range , building on previous theoretical ideas . As explained in , the solution to our extremal problem provides a bound for the loss of negatively refracting materials when the real part of the refractive index is constant in a finite bandwidth.
The following lemma is basic for our parametrization of $`K_\alpha (I)`$.
###### Lemma 1.
A real-valued function in $`L^2(^+I)`$ is the restriction to $`^+I`$ of at most one real-valued function $`v`$ in $`L^2(^+)`$ such that $`v(t)`$ is constant on $`I`$.
###### Proof.
We assume two real-valued functions $`v_1`$ and $`v_2`$ in $`L^2(^+)`$ coincide off $`I`$ and are such that both $`u_1(t)=v_1(t)`$ and $`u_2(t)=v_2(t)`$ are constant on $`I`$. If we set $`c=u_2(t)u_1(t)`$ for $`t`$ in $`I`$, then the function $`h=[(u_1u_2+i(v_1v_2)+c]^2`$ will be real for real arguments. A change of variables argument shows that then $`h(i(1+z)/(1z))`$ belongs to $`H^1`$ of the unit disk. But a function in $`H^1`$ can be real only if it is a constant. Clearly, $`h`$ can be a constant only if $`u_1u_2+i(v_1v_2)=0`$. โ
We note that the assumption that $`v`$ is in $`L^2`$ is essential for this lemma; the proof would break down if we assumed, say, that $`v`$ belonged to some $`L^p`$ for $`p<2`$.
The following function will play an essential role in what follows:
$$\sigma (z)=\frac{1}{\sqrt{z^2b^2}\sqrt{z^2a^2}}.$$
This function, which is taken to be positive for real arguments $`x>b`$, is analytic in the slit plane $`([b,a][a,b])`$. For real arguments $`a<|x|<b`$ we define $`\sigma (x)`$ by extending it continuously from the upper half-plane. Thus $`\sigma (x)`$ takes values on the negative imaginary half-axis when $`x`$ is in $`(a,b)`$ and on the positive imaginary half-axis when $`x`$ is in $`(a,b)`$, and otherwise it is real for real arguments. The key point, besides the symmetry $`\sigma (x)=\overline{\sigma (x)}`$, is that $`\sigma `$ provides a means for switching between real and imaginary when switching off and on $`I`$.
The following is our main result.
###### Theorem.
A nonnegative function $`v`$ in $`L^2(^+)`$ is in $`K_\alpha (I)`$ if and only if the following three conditions hold:
(1)
$$_{^+I}v(t)|\sigma (t)|๐t<\mathrm{}$$
(2)
$$\frac{2}{\pi }_{^+I}tv(t)\sigma (t)๐t=\alpha $$
(3)
$$v(x)=\left((1\chi _I)\sigma v\right)(x)/|\sigma (x)|,xI.$$
Some remarks are in order before we give the proof of the theorem.
The integrability condition (1) is merely a slight growth condition at the endpoints of $`I`$; we may write it more succinctly as
$$_0^a\left[v(at)+v(b+t)\right]\frac{dt}{\sqrt{t}}<\mathrm{}.$$
This condition ensures that the integral in (2) and the Hilbert transform appearing in (3) are both well-defined.
At first sight, the theorem may not seem to give an explicit parametrization of $`K_\alpha (I)`$. However, the Hilbert transform appearing in (3) is given by
$$\left((1\chi _I)\sigma v\right)(x)=\frac{1}{\pi }_{^+I}v(t)\sigma (t)\frac{2t}{t^2x^2}๐t,$$
and we observe that the integrand on the right is nonnegative whenever $`v(t)`$ is nonnegative. Hence $`v(x)0`$ for $`x`$ off $`I`$ implies $`v(x)0`$ for $`x`$ in $`I`$. This small miracle implies that $`K_\alpha (I)`$ is parameterized by those nonnegative functions $`v`$ in $`L^2(^+I)`$ for which (1) and (2) hold and such that
$$_I|\left((1\chi _I)\sigma v\right)(x)|^2|\sigma (x)|^2๐x<\mathrm{}.$$
By rephrasing this condition in more explicit terms (see Lemma 3 below), we arrive at the following corollary.
###### Corollary.
A nonnegative function $`\nu `$ in $`L^2(^+I)`$ has an extension to a function in some class $`K_\alpha (I)`$ if and only if the following condition holds:
$$_0^a_0^a\left[\nu (at)\nu (a\tau )+\nu (b+t)\nu (b+\tau )\right]\frac{|\mathrm{log}(t+\tau )|}{\sqrt{t\tau }}๐t๐\tau <\mathrm{}.$$
The difference between (1) and the condition above is the logarithmic factor, which means that the condition of the corollary is only a very slight strengthening of (1). It is clear that for instance boundedness of $`v`$ near the endpoints of $`I`$ is more than enough.
We note that the integrand in (2) is negative to the left of $`I`$ and positive to the right of $`I`$. This means that if $`\alpha `$ is negative, then
$$|\alpha |\frac{2}{\pi }_0^atv(t)|\sigma (t)|๐t,$$
with equality holding if $`v`$ vanishes to the right of $`I`$. It follows that
$$\left((1\chi _I)\sigma v\right)(x)\frac{1}{\pi }_0^av(t)\sigma (t)\frac{2t}{t^2x^2}๐t\frac{|\alpha |}{x^2};$$
we may come as close as we wish to this lower bound by choosing any suitable $`v`$ supported on a small set sufficiently close to $`0`$. Hence our extremal problem (corresponding to $`\alpha =1`$) has solution
$$\lambda =\underset{xI}{\mathrm{max}}\frac{\sqrt{(b^2x^2)(x^2a^2)}}{x^2}=\frac{b^2a^2}{2ab},$$
as proclaimed above. We also observe that the same function $`1/(x^2|\sigma (x)|)`$ would give the infimum for the $`L^p`$ norm over $`I`$ for any other value of $`p>0`$.
If, on the other hand, $`\alpha `$ is positive, we have instead
$$\alpha \frac{2}{\pi }_b^{\mathrm{}}tv(t)|\sigma (t)|๐t,$$
with equality holding if $`v`$ vanishes to the left of $`I`$. In this case, arguing in the same fashion as above, we find that we can get $`\left((1\chi _I)\sigma v\right)(x)`$ as small as we please by letting $`v`$ be supported on a set sufficiently far to the right of $`I`$.
## 3. Proofs
We now turn to the proof of the theorem and its corollary. We will rely on Lemma 1 and two additional lemmas.
###### Lemma 2.
For every $`t`$ in $`(0,a)(b,\mathrm{})`$ the function
$$f_t(x)=\frac{2t}{t^2x^2}\left(1\frac{\sigma (t)}{\sigma (x)}\right)2t\sigma (t)$$
is in $`H^p`$ for $`p>1/2`$, and the following estimates hold
$$f_t_1\frac{C_1}{\sqrt{|(at)(bt)|}},t<2b,$$
$$_{2b}^{\mathrm{}}|f_t(x)|^2๐t\frac{C_2}{|x|+1},$$
where the constants $`C_1`$ and $`C_2`$ only depend on $`a`$ and $`b`$.
###### Proof.
It is immediate that $`f_t`$ belongs to $`H^p`$ for $`p>1/2`$ because the isolated singularities $`\pm t`$ are removable and $`f_t(z)=O(z^2)`$ when $`z\mathrm{}`$. The norm estimates follow from elementary calculations. โ
###### Lemma 3.
A nonnegative function $`\nu `$ in $`L^2(^+I)`$ satisfies
$$_I|((1\chi _I)\sigma \nu )(x)|^2|\sigma (x)|^2๐x<\mathrm{}$$
if and only if the following condition holds:
$$_0^a_0^a\left[\nu (at)\nu (a\tau )+\nu (b+t)\nu (b+\tau )\right]\frac{|\mathrm{log}(t+\tau )|}{\sqrt{t\tau }}๐t๐\tau <\mathrm{}.$$
###### Proof.
The necessary and sufficient condition for square-integrability at the left end-point of $`I`$ is that
$$_0^a\left(_0^a\frac{\nu (at)}{\sqrt{t}(t+x)}๐t\right)^2x๐x<\mathrm{}.$$
By Fubiniโs theorem, we may interchange the order of integration so that this condition becomes
$$_0^a_0^a\nu (at)\nu (a\tau )\frac{|\mathrm{log}(t+\tau )|}{\sqrt{t\tau }}๐t๐\tau <\mathrm{}.$$
Combining this with the corresponding condition at the right end-point of $`I`$, we arrive at the condition of the lemma. โ
The theorem is now proved in the following way. We assume first that we are given a nonnegative function $`v`$ in $`L^2(^+)`$ satisfying (1), (2), and (3). We claim that the function
$$f_1(x)=\frac{1}{\pi }_{(0,a)(b,2b)}v(t)f_t(x)๐t$$
is in $`H^1`$. Indeed, by Lemma 2 and Fubiniโs theorem,
$$f_1_1\frac{C_1}{\pi }_{(0,a)(b,2b)}\frac{v(t)}{\sqrt{|(at)(bt)|}}๐t,$$
and the integral on the right is bounded thanks to (1). On the other hand,
$$f_2(x)=\frac{1}{\pi }_{2b}^{\mathrm{}}v(t)f_t(x)๐t$$
is in $`H^p`$ for $`p>2`$, because by the CauchyโSchwarz inequality and Lemma 2 we have
$$_{\mathrm{}}^{\mathrm{}}|f_2(x)|^p๐x\frac{C_2^{p/2}}{\pi ^p}v_2^p_{\mathrm{}}^{\mathrm{}}\frac{1}{(|x|+1)^{p/2}}๐x.$$
The imaginary part of $`f_1+f_2`$ is supported by $`I`$ and equals $`v`$ there, in view of (3). Since $`v`$ is assumed to be in $`L^2`$, it follows that $`f_1+f_2`$ is in fact in $`H^2`$.
We set $`f=((1\chi _I)v)+i(1\chi _I)v`$, which is a function in $`H^2`$. We observe that the imaginary part of $`f(f_1+f_2)`$ equals $`v`$ and that its real part equals $`\alpha `$ on $`I`$, when taking into account (2). So we have proved that the given $`v`$ is indeed in $`K_\alpha (I)`$.
We now prove the necessity of the three conditions of the theorem. So assume we are given a nonnegative function $`v`$ in $`L^2(^+)`$ belonging to some class $`K_\alpha (I)`$. Setting $`I_\epsilon =[a+\epsilon ,b\epsilon ]`$, we see that $`v`$ also belongs to $`K_\alpha (I_\epsilon )`$ whenever $`0<\epsilon <(ba)/2`$. But since the Hilbert transform of $`v`$ is constant near the endpoints of $`I_\epsilon `$, it follows that
$$_{^+I_\epsilon }v(t)|\sigma _\epsilon (t)|๐t<\mathrm{},$$
where now
$$\sigma _\epsilon (t)=\frac{1}{\sqrt{t^2(b\epsilon )^2}\sqrt{t^2(a+\epsilon )^2}}.$$
We claim that this means that
(4)
$$v(x)=\left((1\chi _{I_\epsilon })\sigma _\epsilon v\right)(x)/|\sigma _\epsilon (x)|,$$
provided $`x`$ is in $`(a+\epsilon ,b\epsilon )`$. Indeed, by Lemma 1, it is enough to verify that
$$(1\chi _{I_\epsilon }(x))v(x)+\chi _{I_\epsilon }(x)\left((1\chi _{I_\epsilon })\sigma _\epsilon v\right)(x)/|\sigma _\epsilon (x)|$$
is in $`K_\alpha (I_\epsilon )`$ for some $`\alpha `$. Since, in view of Lemma 3, the function on the right-hand side of (4) is square-integrable on $`I_\epsilon `$, the claim follows by repeating the argument in the first part of the proof.
We may view $`\left((1\chi _{I_\epsilon })\sigma _\epsilon v\right)(x)`$ as the $`L^1`$ norm of the function
$$h_{x,\epsilon }(t)=\frac{1}{\pi }\left((1\chi _{I_\epsilon }(t))\sigma _\epsilon (t)v(t)\right)\frac{2t}{t^2x^2}.$$
Then (4) says that $`h_{x,\epsilon }_1v(x)|\sigma (x)|`$ when $`\epsilon 0`$. Since we also have that
$$h_{x,\epsilon }(t)\frac{1}{\pi }\left((1\chi _I(t))\sigma (t)v(t)\right)\frac{2t}{t^2x^2}$$
for every $`t`$, we obtain
$$v(x)=\left((1\chi _I)\sigma v\right)(x)/|\sigma (x)|$$
for every $`x`$ in $`(a,b)`$. By a similar argument, we find that
$$\alpha =\underset{\epsilon 0}{lim}\frac{2}{\pi }_{^+I_\epsilon }tv(t)\sigma _\epsilon (t)๐t=\frac{2}{\pi }_{^+I}tv(t)\sigma (t)๐t.$$
The necessity of (1) has already been observed; without it we would reach the contradictory conclusion that $`v(x)=\mathrm{}`$ for almost every $`x(a,b)`$.
We finally note that the corollary is an immediate consequence of the theorem and lemmas 1 and 3.
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# Controllable soliton emission from a Bose-Einstein condensate.
## Abstract
We demonstrate, through numerical simulations, the controllable emission of matter-wave bursts from a Bose-Einstein Condensate in a shallow optical dipole trap. The process is triggered by spatial variations of the scattering length along the trapping axis. In our approach, the outcoupling mechanism are atom-atom interactions and thus, the trap remains unaltered. Once emitted, the matter wave forms a robust soliton. We calculate analytically the parameters for the experimental implementation of this atomic soliton machine gun.
Introduction.- After the remarkable experimental realization of Bose-Einstein condensates (BEC) with alkali atomsAnderson95 it was soon realized that the coherent behavior of the atom cloud could be used as the basis for an โatom laserโ. The first realization of such device used short radio-frequency pulses as an outcoupling mechanism, flipping the spins of some of the atoms and releasing them from the trap Mewes97 . Later, other atom lasers were built with different configurations leading to pulsed, semi-continuous or single-atom coherent sourcesBloch ; Hagley ; Bloch2 ; Pepe ; Ultimo .
There are important differences between these devices and photon pulsed lasers. The nature of the waves and the lack of a population inversion mechanism are the most evident. However, from the practical point of view, one of the most remarkable is that a pulsed BEC laser will show a significant spreading of the atom cloud if the number of atoms per pulse is large.
A way to overcome this difficulty would be to use atomic bright solitons solitons1 ; solitons2 as matter wave pulses. This idea was explored in Ref. Carr04 to generate a train of solitons by the mechanism of modulational instability (MI) and was referred to as an atomic soliton laser. However, although one could extract a few coherent solitons from such a system this atomic soliton laser would be very limited since: (i) the final number of atoms per pulse after MI is only a small fraction of the initial number of atoms in the condensate due to collapse processes, (ii) the number of solitons generated is not large and half of them would be directed backward (iii) the trap should be destroyed for outcoupling and (iv) the pulses will travel at different speeds once the trap is removed.
Thus, it is important to discuss new outcoupling mechanisms for atom lasers. This is specially interesting since the techniques for generating BEC with growing number of particles and their physical properties are nowadays well established and the current challenges in the field face the applications of coherent matter waves to the design of practical devices chip .
In this letter, we show how a highly controllable train of up to several hundreds of matter-wave solitons can be extracted from a BEC without altering the trap properties but instead acting on the scattering length, by changing it along the atom cloud.
System configuration and theoretical model.- Let us assume that a large BEC is strongly trapped in the transverse directions ($`x,y`$) and weakly confined in the longitudinal one ($`z`$). We now consider the effect of a sharp variation along $`z`$ of the scattering length, which is changed from positive (or zero) to negative, thus making it inhomogeneous. This can be done with magnetic FB1 or opticalFB2 techniques and therefore this region of negative scattering length can be displaced along the condensate. Thus, when it is located close to one edge of the trap, overlapping the wing of the cloud, the tail of the condensate may be able to form a single soliton which, because of its higher internal energy will be outcoupled from the cloud and thus emitted outward. When the condensate refills the gap left out by the outgoing pulse a new soliton would be emitted. This process would continue while there is a large enough remnant of atoms in the trap and would lead to a soliton burst escaping from the BEC.
Thus, we will describe our system of $`N`$ weakly interacting bosons of mass $`m`$, trapped in a potential $`V(\stackrel{}{r})`$ is the mean field limit by a Gross-Pitaevskii equation of the form
$$i\mathrm{}\frac{\mathrm{\Psi }}{t}=\frac{\mathrm{}^2}{2m}^2\mathrm{\Psi }+V(\stackrel{}{r})\mathrm{\Psi }+U(z)|\mathrm{\Psi }|^2\mathrm{\Psi },$$
(1)
where $`\mathrm{\Psi }`$ is the condensate wavefunction, $`N=|\mathrm{\Psi }|^2d^3๐ซ`$ The coefficient $`U(z)=4\pi \mathrm{}^2a(z)/m`$ characterizes the 2-body interaction and since we will consider spatially inhomogeneous systems it will be a function of $`z`$.
In this paper, we consider a BEC tightly confined in ($`x,y`$) by a harmonic potential $`V_{}`$ and more relaxed along $`z`$ where there is and optical dipole trap $`V_z`$Stamper98 ; Martikainen99 . Thus, we have
$$V(\stackrel{}{r})=V_{}+V_z=\frac{m\nu _{}^2}{2}\left(x^2+y^2\right)+V_0\left[1\mathrm{exp}\left(\frac{z^2}{L^2}\right)\right],$$
(2)
where $`V_0`$ is the depth of the shallow optical dipole potential and $`L`$ is the characteristic width of the trap. We will consider situations in which the ground state of the optical dipole trap is much larger than the ground state of the transverse harmonic potential.
In Fig.1 we have plotted the geometry of the system The choice of a shallow Gaussian trap is very important for our model, since we are interested on studying the outcoupling of solitons along the $`z`$ axis. Thus, we need a potential barrier that can be overcome by the self-interaction effects. In this situation we can describe the dynamics of the condensate in the quasi-one dimensional limit as given by a factorized wavefunction of the form Perezgarcia98 $`\mathrm{\Psi }(๐,t)=\mathrm{\Phi }_0(x,y)\psi (z,t)`$, satisfying
$$i\frac{\psi }{\tau }=\frac{r_{}^2}{2}\frac{^2\psi }{z^2}+f(z)\psi +g|\psi |^2\psi ,$$
(3)
where $`r_{}=\sqrt{\mathrm{}/m\nu _{}}`$ is the transverse size of the cloud, $`f(\eta )=V_z/(\mathrm{}\nu _{}/2)`$, $`\tau =\nu _{}t`$ is the time measured in units of the inverse of the radial trapping frequency and $`g(z)=\sqrt{8}\pi r_{}^2a(z)`$ is the effective interaction coefficient.
In Fig. 1 we show a sketch of the setup to be considered in this paper. We show the axial optical dipole potential of depth $`V_0`$ and half width $`L`$. The shaded region ($`z>D`$) has negative scattering length. A spatial variation along $`z`$ of the scattering length $`a`$, can be achieved by magnetically tunning the Feschbach resonances FB1 or by their optical manipulation by means of an additional laser beam FB2 . In this paper we choose a step model for $`a(z)`$ in order to present the basic mechanism. We must stress that our ideas apply to smoother variations of $`a`$. Thus, we will use a dependence of the form
$$a(z)=\{\begin{array}{cc}0,\hfill & z<D\hfill \\ a_{}<0,\hfill & z>D.\hfill \end{array}$$
(4)
The continuous lines in Fig. 1 are the numerically calculated profiles of eigenstates of Eq. (3) for three different values of $`a_{}`$. The lower curve corresponds to he linear case ($`a_{}=0`$), whereas lines b) and c) display the shape of the cloud for $`a_{}=a_{cr}`$ and $`a_{}=1.2a_{cr}`$, respectively. The value $`a_{cr}`$ is the (negative) scattering length needed for emitting one soliton from the cloud. We will calculate it by means of an approximate analytical technique.
Single soliton emission.- In Fig. 1 (black region), it is shown the overlapping of the tail of the trapped cloud and the zone where $`a=a_{}`$. Due to the atom-atom interaction, the trapping changes in this region. Thus, if $`a_{}`$ is negative enough, the minimum of the effective potential will be displaced toward this region and the cloud will move to the shaded region of Fig. 1.
To study this possibility we will use an averaged Lagrangian formalism AL for studying localized solutions of Eq. (3). We will use a Gaussian ansatz
$$\psi (z,\tau )=Ae^{(zz_0(\tau ))^2/2w^2}e^{i\left(v(\tau )z\right)},$$
(5)
where $`z_0(\tau )`$ accounts for the motion of the center the cloud with velocity proportional to $`v`$. The standard calculations lead to the equations
$`\ddot{z_0}`$ $`=`$ $`{\displaystyle \frac{d\mathrm{\Pi }}{dz_0}},`$ (6a)
$`\mathrm{\Pi }(z_0)`$ $`=`$ $`\sqrt{\pi }r_{}^2\{{\displaystyle \frac{V_0}{\mathrm{}\nu _{}}}(1{\displaystyle \frac{e^{\frac{z_0^2}{w^2+L^2}}}{\sqrt{1+\frac{w^2}{L^2}}}})`$ (6b)
$`+{\displaystyle \frac{1}{\sqrt{8\pi }}}{\displaystyle \frac{aN}{w}}\text{erfc}\left[{\displaystyle \frac{\sqrt{2}\left(Dz_0\right)}{w}}\right]\},`$
where $`\text{erfc}(u)=\frac{2}{\sqrt{\pi }}_u^{\mathrm{}}\mathrm{exp}(v^2)๐v`$ is the complementary error function. Thus, $`z_0`$ evolves like a classical particle under the effect of a potential $`\mathrm{\Pi }(z_0)`$.This provides a qualitative understanding of the soliton emission: for the linear case ($`a=0`$) the center of the cloud is located at the bottom of the Gaussian trap, which corresponds to its fundamental eigenstate \[Fig. 1 (a)\]; as $`a_{}`$ takes more negative values the effective trapping of the cloud is deformed and the minimum of the equivalent potential moves to the region with $`z>0`$. Fig. 2 shows the equivalent potentials $`\mathrm{\Pi }`$ given by Eq. (6b) for different values of $`a_{}<0`$. Fig. 2(a) shows the effective potential for $`a_{}=0`$. As $`a_{}`$ is decreased there is a limiting value $`a_{cr}`$ for which the potential $`\mathrm{\Pi }(0)=\mathrm{\Pi }(\mathrm{})`$ \[Fig. 2(b)\], thus if the atom cloud is initially placed at $`z_0=0`$ it will oscillate around the minimum and escape $`z_0(\tau )\mathrm{}`$ for $`\tau \mathrm{}`$, a phenomenon which is called soliton emissionemision . The critical value of $`a_{}`$ that corresponds to the threshold for soliton emission can be obtained within our formalism from the condition, $`\mathrm{\Pi }(0)=\mathrm{\Pi }(\mathrm{})`$ which yields to:
$`Na_{cr}`$ $`=`$ $`{\displaystyle \frac{\sqrt{8\pi }V_0}{\mathrm{}\nu _{}}}{\displaystyle \frac{wL}{\sqrt{L^2+w^2}}}\left[\text{erfc}\left({\displaystyle \frac{\sqrt{2}D}{w}}\right)2\right]^1`$ (7a)
$``$ $`{\displaystyle \frac{\sqrt{8\pi }V_0}{\mathrm{}\nu _{}}}{\displaystyle \frac{wL}{\sqrt{L^2+w^2}}}5.0\mathrm{\Delta }{\displaystyle \frac{L}{w}},`$ (7b)
the first approximation is valid for $`Dw`$ and the second for $`V_0=\mathrm{\Delta }\mathrm{}\nu _{}`$, with $`\mathrm{\Delta }<1`$, which provides a shallow trap with $`wL`$. We have found that the ratio between the exact numerical value and the prediction from Eq. (7a) is a factor $`1.15`$ which is essentially due to the specific choice of the ansatz. In fact, once the soliton is emitted the wavefunction takes a hyperbolic secant shape and its width differs from the Gaussian by a similar factor. However, taking a sech profile as variational ansatz does not yield to analytical results with the potentials we have considered.
For more negative values of $`a<a_{cr}`$ not only the soliton is outcoupled from the system but also it propagates with a finite asymptotic speed in the nonlinear medium. An example of the effective potential for $`a=1.2a_{cr}`$ is shown in Fig. 2(c) and the corresponding soliton emission in Fig. 2(d), for a <sup>7</sup>Li condensate such as that of Ref. solitons1
Partial outcoupling and multisoliton emission.- A deeper numerical exploration based on Eq. (3) reveals more interesting effects which are beyond the averaged lagrangian description. An example is shown in Fig. 3 for parameter values $`V_0=\mathrm{}\nu _{}/2`$, $`\nu _{}=1KHz`$, $`L=4r_{}`$, $`D=2.5L`$, $`N=310^5`$, $`w=5.4r_{}`$. The vertical axis in each figure is time from $`t=0`$ to $`t=1s`$ and the horizontal width of each window is $`100r_{}300\mu m`$. When $`a_{}=0.9a_{cr}`$ the atom cloud is only slightly deviated to the region with $`a<0`$ \[Fig. 3 (a) \]. Decreasing $`a_{}`$ down to $`a_{}=1.95a_{cr}`$ leads to the phenomenon of soliton emission with some remnant staying on the optical dipole trap \[Fig. 3(b)\]. However, as $`a_{}`$ is made more and more negative we obtain the emission of an integer number of solitons.
The phenomenon of partial emission can be qualitatively understood with the variational method. As it can be seen in Fig. 2 the effective potentials $`\mathrm{\Pi }`$ have only one minimum for $`D=2.0L`$. However if the value of $`D`$ is increased to say $`D=2.5L`$, the potential $`\mathrm{\Pi }`$ for values close to $`a_{cr}`$ presents one maximum between two adjacent minima. This configuration of the effective potential causes the split of the matter wave in two parts: one remains in the trap and the other is emitted as a soliton.
From our numerical exploration we have found some regularities that are noticeable. In first place, as it can be appreciated from Fig. 3, the velocities of the outgoing pulses are almost equal. Secondly, the intervals of emission are very regular and decrease with the effective nonlinearity with a law of the form $`\tau (a_{}/a_{cr})^2`$. Finally, the values of the scattering length required for outcoupling several solitons are almost exact multiples of $`a_{cr}N`$ as it is clear by the approximate linear dependence of the number of solitons emitted in the inset of Fig. 5. This can be interpreted as the requirement of a critical effective nonlinearity remaining in the trapped cloud for the formation of every soliton. Finally, we must comment from Fig. 4 that the amplitude of the emitted solitons decays from the first one of the burst to the last approximately in a 30%, depending on the number of pulses.
Discussion and conclusions.- These regular and controllable soliton trains could be realized with any atomic specie for which the scattering length could be made negative on a region of the space. This has been done with several atomic species such as <sup>7</sup>Li, <sup>85</sup>Rb and <sup>133</sup>Cs. Taking for instance the experimental parameters of Ref. solitons1 for Lithium and assuming $`3\times 10^5`$ atoms in the optical dipole trap we obtain that a few (2-3) solitons could be emitted for a depth of $`V_0=\mathrm{}\nu _{}/2`$. However, increasing the initial particle number in the zero-scattering length region could allow the observation of soliton trains with a higher number of solitons, e.g. for $`3\times 10^6`$ we estimate up to $`N_{sol}20`$. For <sup>85</sup>Rb and taking parameter data from Ref. 85Rb we estimate about $`25`$. Even more interesting effects could be obtained with Cesium, for which experimental data of Ref. 133Cs lead to an estimate of $`N_{sol}200`$. In all cases these numbers can be raised or lowered by acting on the parameters: $`a_{},N,D`$ and $`V_0`$.
In summary, we have proposed a novel mechanism for outcoupling coherent matter wave pulses from a Bose-Einstein Condensate. Our system is able to perform a regular and controllable emission of atomic soliton bursts that are easily extracted by an adequate choice of the control parameters. Using this mechanism a train of even several hundred of solitons could be coherently outcoupled from a condensate. As the techniques for coherently feeding the remaining condensate progress our idea could provide an outcoupling mechanism for a continuous atomic soliton laser.
This work was supported by Ministerio de Educaciรณn y Ciencia, Spain (projects FIS2004-02466, BFM2003-02832, network FIS2004-20188-E) and by Xunta de Galicia (project PGIDIT04TIC383001PR).
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# Acceleration disturbances and requirements for ASTROD I
## 1 Introduction
ASTRODynamical Space Test of Relativity using Optical Devices (ASTROD) aims at testing relativistic gravity, measuring the solar-system parameters with high precision and detecting gravitational waves from massive black holes and galactic binary stars. The concept of ASTROD is to put two spacecraft in separate solar orbits and carry out laser interferometic ranging with Earth reference stations (e.g. a spacecraft at the Earth-Sun L1/L2 points). A simple version of ASTROD, ASTROD I<sup>1</sup><sup>1</sup>1ASTROD I was previously referred to as Mini-ASTROD (for instance in )., has been studied as the first step to ASTROD . ASTROD I employs one spacecraft in a solar orbit and carries out interferometric ranging and pulse ranging with ground stations. The main scientific goals of ASTROD I are to test relativistic gravity and the fundamental laws of spacetime with three-order-of-magnitude improvement in sensitivity and to improve the solar, planetary and asteroid parameter determination by 1 to 3 orders of magnitude. The technological goal of ASTROD I is to prepare for the ASTROD mission.
The acceleration disturbance goal of the ASTROD I proof mass is $`10^{13}`$ m s<sup>-2</sup> Hz<sup>-1/2</sup> at frequency $`\nu `$ of 0.1 mHz. The power spectral density of the allowed level of the acceleration noise is shown in figure 1. Assuming a 10 ps one-way timing accuracy (3 mm ranging accuracy) and the acceleration noise of $`10^{13}`$ m s<sup>-2</sup> Hz<sup>-1/2</sup> at frequency of about 0.1 mHz, a simulation for 400 days (350โ750 days after launch) showed that ASTROD I could determine the relativistic parameters $`\gamma `$ and $`\beta `$, and the solar quadrupole parameter $`J_2`$ to levels of 10<sup>-7</sup>, 10<sup>-7</sup> and 10<sup>-8</sup>, respectively . In the simulation, (i) the timing noise is modeled as Gaussian random noise; (ii) unknown acceleration noise is modeled to have Gaussian random magnitude with zero mean and with standard deviation 10<sup>-15</sup> m s<sup>-2</sup> and to have its direction changed randomly every 4 h (equivalent to 10<sup>-13</sup> m s<sup>-2</sup> Hz<sup>-1/2</sup> for $`\nu `$ $``$ 0.1 mHz assumed as the requirement of the drag-free system) and (iii) five range points are taken each day (at 0.2 d interval). Longer term systematic effects will be studied in a future paper. This simulation agrees with the scientific goals of ASTROD I. The timing uncertainty of event timer reaches 3 ps in satellite laser ranging at present. Space qualified versions of similar accuracy are under development. For a ranging uncertainty of 3 mm in a distance of 3 $`\times `$ 10<sup>11</sup> m (2 AU), the laser/clock frequency needs to be known to one part in 10<sup>14</sup>. This can be set as a requirement of the space laser/clock or a requirement for laser frequency monitoring through ground clock and modeling. As to ground station jitter, monitoring to an accuracy of 3 mm is required and can be achieved. The atmospheric effects on laser propagation will be monitored and subtracted to mm-level by using 2-color (2-wavelength) ranging (one color for pulse ranging and one for interferometric ranging). These measurement uncertainties are not cumulative in the range determination while the acceleration disturbances accumulate in time in the geodesic deviations. In order to achieve the acceleration disturbance goal, a drag-free control system using capacitive sensors will be employed.
In order to decide on detailed designs of the accelerometer for ASTROD I, we have to know the sources and magnitude of acceleration disturbances that could arise in the accelerometer. In this paper, we carry out analyses, mainly based on existing literature on acceleration disturbances for other gravitational missions (LISA and the LISA Pathfinder ), to give a preliminary overview of the acceleration disturbance to the ASTROD I proof mass. Based on the analyses, we infer some of the requirements for the designs of the ASTROD I payload and spacecraft. Also, we compare parameter values we have assumed in the analyses with those for LISA to confirm the feasibility of the requirements for ASTROD I.
First, we will give an overview of the ASTROD I configuration (section 2) and the control-loop model we assumed (section 3). Then, we will estimate the magnitude of acceleration disturbances and requirements for ASTROD I in sections 4 to 8, and compare the requirements of ASTROD I with LISA in section 9.
## 2 ASTROD I spacecraft configuration
The ASTROD I spacecraft has a cylindrical shape with diameter 2.5 m and height 2 m. Its cylindrical side is covered by solar panels. The cylindrical axis is perpendicular to the orbit plane and the telescope is set to point toward a ground laser station. The effective area of receiving sunlight is about 5 m<sup>2</sup> and it can generate power that is larger than 500 W. The total mass of the spacecraft is about 350 kg and that of payload is 100โ120 kg (see for more detailed descriptions of the configuration). The orbit distance from the Sun varies from about 0.5 AU to 1 AU (figure 2 of ).
The proof mass ($`m_p`$ = 1.75 kg) is a rectangular parallelepiped (50 $`\times `$ 50 $`\times `$ 35 mm<sup>3</sup>)<sup>2</sup><sup>2</sup>2This is the current design of the proof mass. A cylindrical shape is also considered as an alternative design of the proof mass. made from Au-Pt alloy (density $`\rho `$ = 2 $`\times `$ 10<sup>4</sup> kg m<sup>-3</sup>). The six sides of the proof mass are surrounded by electrodes mounted on the housing for capacitive sensing and control. The gap between each side of the proof mass and the opposing electrode is 2 mm. Assumed values for the capacitance and voltages are listed in table 1.
## 3 Control-loop model
Various acceleration disturbances would act on the proof mass in different ways. In order to infer how these different kinds of acceleration disturbances would contribute to the total acceleration disturbance of the proof mass, we tentatively assume a simple single-mass and single-axis control-loop model . The diagram is shown in figure 2. The relative difference between the displacement disturbance amplitudes of the proof mass and of the spacecraft, $`X_{ps}=X_pX_s`$, with position readout noise $`X_{nr}`$, is measured by the position displacement sensor. The output of this sensor is converted to acceleration disturbance $`f_r`$, by a transfer function $`R`$ ($`\omega _R^2`$). This acceleration is supplied to the thruster and the output acceleration disturbance with thruster noise $`N_t`$ is applied to the spacecraft. The spacecraft also experiences acceleration disturbances by coupling to the proof mass with a coupling constant $`K`$ and external environmental disturbances ($`f_{ns}`$) that work directly on the spacecraft. The total acceleration of the spacecraft $`f_s`$ is converted to the position noise $`X_s`$ with a transfer function $`S`$.
We consider only the sensitive axis of the proof mass without actuation. The proof mass would experience disturbances by spacecraft-proof mass coupling and environmental disturbances $`f_{np}`$ (see section 5). The total acceleration disturbance of the proof mass $`f_p`$ is converted to the displacement disturbance $`X_p`$ with a transfer function $`P`$.
From this control-loop model, we obtain the following linear loop equations :
$`f_r`$ $`=`$ $`R(X_{ps}+X_{nr})`$ (1)
$`f_t`$ $`=`$ $`T(f_r+N_t)`$ (2)
$`f_s`$ $`=`$ $`f_{ns}+f_t(m_p/M_{sc})f_{ps}`$ (3)
$`f_p`$ $`=`$ $`f_{np}+f_{ps}`$ (4)
$`f_{ps}`$ $`=`$ $`KX_{ps}`$ (5)
$`X_s`$ $`=`$ $`Sf_s`$ (6)
$`X_p`$ $`=`$ $`Pf_p`$ (7)
By solving these equations for $`X_p`$, and assuming $`S=P=\omega ^2`$ and $`u`$ $``$ $`STR`$ (= $`T\omega _R^2\omega ^2`$; $`T`$ is nominally 1 and its effects are absorbed in $`R`$ and $`N_t`$), we obtain (e.g. ):
$$f_pX_{nr}(K)+f_{np}+(f_{ns}+TN_t)K\omega ^2u^1$$
(8)
where $`\omega =2\pi \nu `$. This acceleration disturbance has to be less than the acceleration noise goal of $`10^{13}`$ m s<sup>-2</sup> Hz<sup>-1/2</sup> at $`\nu `$ = 0.1 mHz.
We will estimate the values of the direct acceleration disturbances of the spacecraft ($`f_{ns}`$ and the thruster noise $`TN_t`$) in section 4 and of the proof mass ($`f_{np}`$) in sections 5 and 6, and the stiffness $`K`$ in section 7. We will discuss the requirements for $`X_{nr}`$ and $`u`$ in section 8.
## 4 Direct acceleration disturbances of the spacecraft
The spacecraft would be affected by environmental disturbances that stem from, for example, solar radiation pressure, solar wind and micrometeorite impacts. Among these sources of disturbances, solar radiation pressure is considered to be the major contributor to the acceleration disturbances (section 7 of ). The contribution from solar wind might be comparable to radiation pressure, but the spectral behavior of the solar wind is not well known.
By assuming a perfectly reflecting surface of the spacecraft, acceleration noise caused by fluctuation in solar irradiance $`\delta W_0`$ is $`f_{ns,srp}=\frac{4}{3}(A_{sc}\delta W_0/M_{sc}c)`$, where $`A_{sc}`$ is the area of the spacecraft facing the Sun, $`M_{sc}`$ is the mass of the spacecraft and $`c`$ is the speed of light in vacuum. From the data of the VIRGO experiment on SOHO , fractional fluctuation in solar irradiance is $`\delta W_0/W_02.8\times 10^3`$ Hz<sup>-1/2</sup> at 0.1 mHz (figure 6 of ). Assuming that $`\delta W_0/W_0`$ is the same at 0.5 AU, and taking a total irradiance of 5500 W m<sup>-2</sup>, the area of 5 m<sup>2</sup> and a mass of 350 kg, we obtain $`f_{ns,srp}`$ = 9.8 $`\times `$ 10<sup>-10</sup> m s<sup>-2</sup> Hz<sup>-1/2</sup>.
The impact rate of 1-ng meteorites on the ASTROD I spacecraft (its surface area is 25.5 m<sup>2</sup>) is about 7 events per day or about 0.08 mHz<sup>3</sup><sup>3</sup>3This rate was tentatively estimated from the meteorite flux/mass distribution used by the LISA study team (figure 7.2-29 of ).. The average velocity of meteorites is 18 km s<sup>-1</sup> in an Earth frame near the Earth (figure 7.2-30 of ). An impact of 1-ng meteorite with the average velocity of 18 km s<sup>-1</sup> on the surface of the spacecraft (350 kg) without reflection produces a linear velocity increment of about 5 $`\times `$ 10<sup>-11</sup> m s<sup>-1</sup>. Smaller meteorites have larger flux (figure 7.2-29 of ), but their impacts on velocity changes are smaller. Therefore, the contribution to the acceleration disturbances at 0.1 mHz seems insignificant, in comparison with the effect of solar radiation pressure. However, because such discrete changes could directly affect the ASTROD I experiment, the impact effects have to be studied carefully in detail.
In addition to the acceleration noise from the environmental disturbances, the spacecraft would suffer from the thruster noise. A force fluctuation of 10 $`\mu `$N Hz<sup>-1/2</sup> in thruster corresponds to the acceleration disturbance of $`TN_t`$ = 2.8 $`\times `$ 10<sup>-8</sup> m s<sup>-2</sup> Hz<sup>-1/2</sup>.
Therefore, the rss (root-sum-square) of these acceleration disturbances is:
$$f_{ns}+TN_t2.8\times 10^8\mathrm{m}\mathrm{s}^2\mathrm{Hz}^{1/2}$$
(9)
at 0.1 mHz, dominated by the thruster noise. As per , LISA requires a force fluctuation to be less than a few $`\mu `$N Hz<sup>-1/2</sup>. However, recent studies for the LISA Pathfinder indicate that force noise up to 0.1 mN Hz<sup>-1/2</sup> at 0.1 mHz can be tolerated by increasing the gain ; the LISA study team is designing high gain loops and the requirement has been relaxed<sup>4</sup><sup>4</sup>4This is a comment from one of the referees..
## 5 Direct acceleration disturbances of the proof mass
Direct proof-mass acceleration disturbances can be classified into two categories depending on their origins : environmental disturbances ($`f_{nep}`$) and proof-mass sensor back-action acceleration disturbances ($`f_{nbp}`$). The former includes disturbances related to magnetic effects ($`f_{m1}`$, $`f_{m2}`$, $`f_{m3}`$, $`f_{L1}`$ and $`f_{L2}`$), impact effects (by cosmic ray ($`f_c`$) and residual gas ($`f_{rg}`$)), temperature dependent effects (radiometric and outgassing effects ($`f_{re}`$ and $`f_{og}`$) and thermal radiation pressure ($`f_{tr}`$)) and gravity gradients caused by thermal distortion of the spacecraft ($`f_{gg}`$). The latter is originated from voltage fluctuations ($`f_{b1}`$ and $`f_{b2}`$), charge fluctuations ($`f_{b3}`$ and $`f_{b4}`$), readout electronics ($`f_{ba}`$), patch field voltages ($`f_{pe}`$) and thermal voltage noise by dielectric losses ($`f_{dl}`$).
Parameter values and physical constants used for the estimations are listed in tables 1 and 2, respectively. Table 3 provides a summary of the expressions used to estimate the direct proof-mass acceleration disturbances and the estimated values. The disturbances noted as $`f_{m1}`$, $`f_{m2}`$, $`f_{L1}`$, $`f_{L2}`$, $`f_c`$, $`f_{rg}`$, $`f_{re}`$, $`f_{tr}`$, $`f_{gg}`$ and $`f_{b1}`$-$`f_{b4}`$ in table 3 correspond to $`A_1`$-$`A_6`$ and $`A_8`$-$`A_{14}`$ in , respectively<sup>5</sup><sup>5</sup>5The acceleration disturbance due to laser photon radiation pressure, noted as $`A_7`$ in , would not arise in ASTROD I, as the laser beam is not to be injected directly on the surface of the ASTROD I proof mass.. We will briefly explain each disturbance below.
### 5.1 Magnetostatic interaction
The lowest order force on a proof mass with a magnetic moment $`\stackrel{}{M_p}`$ in an external magnetic field $`\stackrel{}{B}`$ is given by $`\stackrel{}{F_m}`$ = $`\stackrel{}{}(\stackrel{}{M_p}\stackrel{}{B})`$ . The magnetic moment of the proof mass is a vector sum of the remanent moment $`\stackrel{}{M_r}`$ and the induced moment : $`\stackrel{}{M_p}=\stackrel{}{M_r}+(\chi _mV_p/\mu _0)\stackrel{}{B}`$, where $`\chi _m`$ and $`V_p`$ are magnetic susceptibility and the volume of the proof mass, respectively; $`\mu _0=1.26\times 10^6`$ N A<sup>-2</sup> is the permeability of vacuum. The external magnetic field would be given by the superposition of the interplanetary magnetic field $`\stackrel{}{B_{ip}}`$ and a local magnetic field $`\stackrel{}{B_{sc}}`$. Dominant terms in acceleration of the proof mass due to the induced magnetic moment produce acceleration disturbances noted $`f_{m1}`$ and $`f_{m2}`$ in table 3, where $`\xi _m`$ is a scaling factor for possible suppression by magnetic shielding.
The average interplanetary magnetic field at 1 AU from the Sun varies from $`10^9`$ to 3.7 $`\times `$ 10<sup>-8</sup> T . Ulysses data obtained near 1 AU from the Sun (figure 9 of ) showed a $`\nu ^{2/3}`$ dependence of the variation in the interplanetary magnetic field and $`\delta B_{ip}`$ can be inferred to be about $`10^7`$ T Hz<sup>-1/2</sup> at 0.1 mHz. As the behavior of $`\delta B_{ip}`$ at 0.5 AU from the Sun is uncertain, we use a somewhat higher value of $`\delta B_{ip}=4\times 10^7`$ T Hz<sup>-1/2</sup> (see footnote in section 5.2). According to the formulation studies for LISA and the implementation work of the LISA Pathfinder, batteries and micro-thrusters are the principal suspects of the origins of the local magnetic field in LISA<sup>6</sup><sup>6</sup>6This is a comment from one of the referees.. We need elaborate modeling works to estimate the magnitude of local magnetic field. Here, we use the same values used in analyses for LISA : $`B_{sc}8\times 10^7`$ T, $`\delta B_{sc}`$ = 10<sup>-7</sup> T Hz<sup>-1/2</sup> and $`|\stackrel{}{}B_{sc}|`$ $``$ 3 $`B_{sc}r_m^1`$ $``$ 3 $`\times `$ 10<sup>-6</sup> T m<sup>-1</sup> (where $`r_m`$ = 0.75 m). Silvestri et al. have reported that the magnetic susceptibility of five samples ranged from $``$2.8 $`\times `$ 10<sup>-5</sup> to $``$2.1 $`\times `$ 10<sup>-5</sup> for two of them without traceable iron contamination and from +1.1 $`\times `$ 10<sup>-5</sup> to +8.8 $`\times `$ 10<sup>-5</sup> for the rest samples with a trace of iron . Lower susceptibility may be achievable by controlling the manufacturing process of the alloy. The requirement for LISA is $`\chi _m`$ = 3 $`\times `$ 10<sup>-6</sup> . We use a somewhat moderate value of $`\chi _m`$ = 5 $`\times `$ 10<sup>-5</sup>. Assuming $`\xi _m`$ = 1 as per , we obtain $`f_{m1}`$ = 1.2 $`\times `$ 10<sup>-15</sup> m s<sup>-2</sup> Hz<sup>-1/2</sup> and $`f_{m2}`$ = 3.4 $`\times `$ 10<sup>-15</sup> m s<sup>-2</sup> Hz<sup>-1/2</sup> at 0.1 mHz.
Acceleration disturbance due to the magnetic remanent moment is given by $`f_{m3}`$ listed in table 3 (equation 4(c) of ), where $`\mathrm{\Delta }(\delta B)`$ is time-varying magnetic field gradients. According to measurements by Gill et al., magnetic remanent moment of three 4.2-cm 70/30 Au-Pt alloy cubes (one of them was 73/27 Au-Pt alloy) ranged from 3.5 $`\times `$ 10<sup>-9</sup> A m<sup>2</sup> kg<sup>-1</sup> to 3.1 $`\times `$ 10<sup>-8</sup> A m<sup>2</sup> kg<sup>-1</sup> . By scaling the largest measured remanent moment by weight, it is $`5.4\times 10^8`$ A m<sup>2</sup> for a 1.75-kg proof mass. By using a somewhat relaxed value of 1 $`\times `$ 10<sup>-7</sup> A m<sup>2</sup> and $`|(\delta B)|=4\times 10^8`$ T m<sup>-1</sup> Hz<sup>-1/2</sup> , which is a factor of four higher than the assumed requirement for LISA in , we obtain $`f_{m3}=1.6\times 10^{15}`$ m s<sup>-2</sup> Hz<sup>-1/2</sup> at 0.1 mHz. In the error estimates for LISA , magnetic remanent moment of 2 $`\times `$ 10<sup>-8</sup> A m<sup>2</sup> is used.
### 5.2 Lorentz force
The proof mass in orbit would be charged up by cosmic-ray impacts . As the charged proof mass moves through the interplanetary magnetic field ($`\stackrel{}{B_{ip}}`$) with a velocity $`v`$ of about 4 $`\times `$ 10<sup>4</sup> m s<sup>-1</sup>, it experiences the Lorentz force: $`\stackrel{}{F_L}`$ = $`q\stackrel{}{v}\times \stackrel{}{B_{ip}}`$, where $`q`$ is the built-up charge. Acceleration disturbances due to the fluctuation of the charge build-up ($`\delta q`$) in the proof mass and of the average interplanetary magnetic field are given by $`f_{L1}`$ and $`f_{L2}`$ (see table 3), respectively, where $`\xi _e`$ is an electrostatic shielding factor; in the rest frame of the proof mass, the motion of the proof mass through the interplanetary magnetic field generates an electrostatic field.
By assuming the Poisson distribution of cosmic-ray impacts, the average charge fluctuation spectral density can be defined as $`\delta q(\omega )\sqrt{2e\dot{q}}/\omega `$. For the frequency ($`\nu =(2\pi )^1\omega `$) of 0.1 mHz and the effective charging rate ($`\dot{q}`$) of 288 $`+\mathrm{e}`$ s<sup>-1</sup>, which was estimated for a LISA proof-mass (46-mm cube) by a simulation using GEANT 4 toolkit , we obtain $`\delta q(\omega )`$ = 6.1 $`\times `$ 10<sup>-15</sup> C Hz<sup>-1/2</sup>. Using this value, $`\xi _e`$ = 10 and $`B_{ip}`$ = 1.2 $`\times `$ 10<sup>-7</sup> T, we obtain $`f_{L1}`$ = 1.7 $`\times `$ 10<sup>-18</sup> m s<sup>-2</sup> Hz<sup>-1/2</sup> at 0.1 mHz. The maximum value of $`B_{ip}`$ at 1 AU is about 3 $`\times `$ 10<sup>-8</sup> T. We tentatively use four times the maximum value as $`B_{ip}`$ at 0.5 AU, by assuming 1/(distance)<sup>2</sup> dependence of $`B_{ip}`$<sup>7</sup><sup>7</sup>7The solar dipole field at 0.5 AU is less than $``$ 10<sup>-10</sup> T. The main magnetic field at 0.5 AU to 1 AU is due to the influence of solar winds which attenuate with (distance)<sup>2</sup>.. Because the volume of the ASTROD I proof-mass is about 10 % smaller than the LISA proof mass, the charging rate for ASTROD I might be smaller than the value for LISA. Bao et al. are working on simulations to estimate the charging rates for ASTROD I . Stebbins et al. use $`\xi _e`$ = 100, which is one order of magnitude larger than the value we used here, in the current error estimates for LISA . Taking a nominal maximum build-up charge $`q`$ = 10<sup>-12</sup> C, which is one order of magnitude larger than the value used in the error estimates for LISA , we obtain $`f_{L2}`$ = 9.1 $`\times `$ 10<sup>-16</sup> m s<sup>-2</sup> Hz<sup>-1/2</sup>.
### 5.3 Cosmic-ray impacts
Some cosmic rays get stopped in the proof mass and deposit momentum . Assuming the Poisson distribution of cosmic ray impact, spectral density of momentum transfer ($`p`$) is $`2p^2\lambda `$, where $`\lambda `$ is the fluctuation in the impact rate . The acceleration disturbance due to the fluctuation in the impact rate is given by $`f_c`$, listed in table 3. The impact rate was inferred from simulations done for LISA; by adding the effects of all stopped particles (protons and helium) and taking into account their directions, the acceleration disturbance by momentum transfer was estimated to be $``$ 2 $`\times `$ 10<sup>-18</sup> m s<sup>-2</sup> Hz<sup>-1/2</sup> for a LISA proof mass . This corresponds to a disturbance due to momentum transfer by protons (mass $`m=`$1.7 $`\times `$ 10<sup>-27</sup> kg), with an impact rate of $``$ 30 s<sup>-1</sup>, at incident energy $`E_d`$ = 200 MeV (= 3.2 $`\times `$ 10<sup>-11</sup> J). Using these values, the acceleration disturbance becomes $`f_c`$ = 1.5 $`\times `$ 10<sup>-18</sup> m s<sup>-2</sup> Hz<sup>-1/2</sup>.
### 5.4 Residual-gas impacts
From the kinetic theory, the number of residual-gas molecules (assumed as ideal gas) that pass an area ($`A_p`$) of the proof mass per second is given by $`\varpi `$ = $`nA_P\overline{v}/6`$, where $`n`$ = $`P(k_BT_P)^1`$ is the number density of the molecules and $`\overline{v}`$ = $`\sqrt{3k_BT_Pm_N^1}`$ is the average thermal velocity ($`P`$ is the pressure of residual gas, $`k_B`$ = 1.38 $`\times `$ 10<sup>-23</sup> J K<sup>-1</sup> is the Boltzmann constant, $`T_P`$ is the temperature of the proof-mass housing and $`m_N`$ = 4.65 $`\times `$ 10<sup>-26</sup> kg is the mass of nitrogen molecules). Assuming the Poisson distribution of the impact rate, we define the spectral density of fluctuations in $`\varpi `$ as $`\delta \varpi (\omega )`$ $``$ $`\sqrt{2\varpi }`$.
Acceleration due to the residual gas impacts is given by $`2m_N\varpi \overline{v}m_p^1`$ and acceleration disturbance due to fluctuations in the impact rate of residual gas is given by $`f_{rg}`$ listed in table 3. For $`P`$ = 10<sup>-5</sup> Pa, $`A_P`$ = 1.75 $`\times `$ 10<sup>-3</sup> m<sup>2</sup> and $`T_P`$ = 293 K, we obtain $`f_{rg}`$ = 7.4 $`\times `$ 10<sup>-16</sup> m s<sup>-2</sup> Hz<sup>-1/2</sup>. Stebbins et al. use 3$`\times `$10<sup>-6</sup> Pa as the residual gas pressure around the proof mass in the error estimates for LISA .
### 5.5 Radiometric effect
Acceleration disturbance due to the radiometric effect (e.g. ) is given by $`f_{re}`$ listed in table 3, where $`\delta T_d`$ is fluctuation in temperature difference across the proof mass housing. This value has to be estimated by carrying out thermal modeling. According to thermal analysis for LISA, temperature fluctuation on the optical bench, to which the proof mass housing is mounted, due to power dissipation of amplifiers is about 3.0 $`\times `$ 10<sup>-5</sup> K Hz<sup>-1/2</sup> at 1 mHz (table 6.2-28 of ). By assuming that the fluctuation rises as $`1/\nu `$, the temperature fluctuation of the optical bench is $`\delta T_{ob}=3.0\times 10^4`$ K Hz<sup>-1/2</sup> at 0.1 mHz. At the frequency of 0.1 mHz and higher frequencies, the temperature fluctuation on the optical bench would be dominated by the fluctuation in the power dissipation ; fluctuation due to solar irradiance at 0.1 mHz is $`1.1\times 10^6`$ K Hz<sup>-1/2</sup> (table 6.2-16 of ). The ratio $`\delta T_{ob}`$/$`\delta T_d`$ would range from 30 to 100 . By using a value of 30 for the ratio, we obtain $`\delta T_d=1.0\times 10^5`$ K Hz<sup>-1/2</sup> at 0.1 mHz. By using this value, we obtain $`f_{re}`$ = 1.7 $`\times `$ 10<sup>-16</sup> m s<sup>-2</sup> Hz<sup>-1/2</sup> at 0.1 mHz.
### 5.6 Temperature dependent outgassing effect
Outgassing from walls of the sensor cage is thought to produce greater acceleration noise than the radiometric effect . An analysis done for the LISA Technology Package (LTP), assuming a simple model of flow circuit with a linear approximation, shows that the outgassing effect is nearly 10 times the radiometric effect . By using this estimate and the estimate we made for the radiometric effect in the previous section, we obtain $`f_{og}`$ = 1.7 $`\times `$ 10<sup>-15</sup> m s<sup>-2</sup> Hz<sup>-1/2</sup> at 0.1 mHz for ASTROD I.
### 5.7 Thermal radiation pressure
By assuming a perfectly reflecting surface of the proof mass, thermal radiation pressure produces acceleration disturbance ($`f_{tr}`$ in table 3) due to fluctuations in the temperature difference across the proof-mass housing. In the expression of $`f_{tr}`$ listed in table 3, $`\sigma `$ = 5.7 $`\times `$ 10<sup>-8</sup> W m<sup>-2</sup> K<sup>-4</sup> is the Stefan-Boltzmann constant, $`A_p`$ is the area of the proof mass and $`T_p`$ is the temperature of the proof-mass housing. A factor of one-third is multiplied, as done in the estimation for LISA by Schumaker , as a margin for the fact that not all of the radiation momentum is normally incident on the proof mass. For the housing temperature of 293 K and the temperature fluctuation of 1.0 $`\times `$ 10<sup>-5</sup> K Hz<sup>-1/2</sup>, we obtain $`f_{tr}`$ = 1.3 $`\times `$ 10<sup>-16</sup> m s<sup>-2</sup> Hz<sup>-1/2</sup> at 0.1 mHz. The same value was used for the temperature fluctuation in the error estimates for LISA .
### 5.8 Gravity gradients due to thermal distortion of the spacecraft
The temperature fluctuation in solar irradiance would cause fluctuation in distortions of the spacecraft: $`\delta xx^1=CTE|\delta T_{sc}|`$, where $`CTE`$ and $`\delta T_{sc}`$ are the coefficient of thermal expansion and temperature fluctuation of the spacecraft, respectively. The inherent fluctuation in solar radiation is $`\delta W_0/W_0=4\delta T/T_p2.8\times 10^3`$ Hz<sup>-1/2</sup> at 0.1 mHz . Therefore, the temperature fluctuation is 0.2 K Hz<sup>-1/2</sup> for $`T_p`$ = 293 K. The gravitational disturbance by a 1-kg mass ($`M_{dis}`$) separated from the proof mass in the sensitive axis by $`x`$ = 0.5 m of aluminium structure would be $`f_{gg}=2.7\times 10^{15}`$ m s<sup>-2</sup> Hz<sup>-1/2</sup> for $`\delta T_{sc}`$ = 0.2 K Hz<sup>-1/2</sup> and $`CTE`$ (of Aluminium) = 2.5 $`\times `$ 10<sup>-5</sup> K<sup>-1</sup>. The same disturbing mass ($`M_{dis}`$ = 1 kg and $`x`$ = 0.5 m) is assumed in . In reality, the mass to be involved in thermal distortion would be much larger, but the influence would be largely canceled because of the axial symmetry in the original spacecraft geometry. The inherent fluctuation in solar radiation could be reduced largely by thermal shielding. In the estimation for LISA, $`\delta T_{sc}`$ = 0.004 K Hz<sup>-1/2</sup> is used . For a more accurate estimate, gravity effects by thermal and non-thermal distortions of the spacecraft and the payload have to be studied by appropriate modeling.
## 6 Proof-mass sensor back-action acceleration disturbances
The total mechanical energy of the capacitive sensing system can be expressed as (e.g. equation (A.3) of ):
$$W=\frac{1}{2}\underset{i}{}C_i(V_iV_s)^2+\frac{1}{2}\frac{q^2}{C}+qV_s$$
(10)
where $`q`$ is a net charge of the proof mass; $`C`$ is the sum of the capacitances due to the applied voltages on the surrounding electrodes $`i`$ and the potential to ground $`g`$ : $`C=_iC_i`$, where $`i=x_1,x_2,y_1,y_2,z_1,z_2,g`$; $`V_s`$ is the voltage induced on the proof mass by the applied voltages on the electrodes: $`V_sC^1_iV_iC_i`$. The first term of $`W`$ is the total energy done on the proof mass by the applied voltages of the surrounding electrodes. The second term is the energy acquired on the proof mass by the image charge on the surrounding electrodes. The third term is the energy stored on the proof mass by the deposit of the free charge on the proof mass.
The $`x`$-component of the force on the proof mass is given by differentiation of equation (10) with respect to $`x`$. For a simplicity, we assume that neither the free charge $`q`$ nor the potentials $`V_i`$ have appreciable gradients along the $`x`$-axis, as per . Acceleration disturbances due to fluctuations in the applied voltages and charge can be given by $`f_{b1}`$ \- $`f_{b4}`$ as listed in table 3. A detailed description on the deviations of $`f_{b1}`$ \- $`f_{b4}`$ is given in Appendix A of . The parameter values used in this section and the estimated values are listed in tables 1 and 3, respectively.
To simplify the analysis, several assumptions were made in the process of deriving the expressions for the four classes of acceleration disturbances ($`f_{b1}`$ \- $`f_{b4}`$) : (1) $`C_i`$ are all comparable with each other in magnitude and on the order of $`C_x`$ $``$ 6 pF (= $`ฯต_0Ad^1`$, where $`\epsilon _0`$ = 8.9 $`\times `$ 10<sup>-12</sup> F m<sup>-1</sup> is the permittivity of vacuum, $`A`$ is the area of each electrode and $`d`$ = 2 mm is the gap); (2) only the capacitances $`C_{x1}`$ and $`C_{x2}`$ have nonzero gradients along the $`x`$-axis: $`C_{x1}^{}`$ = $`C_{x2}^{}`$, and the gradient of the total capacitance $`C`$ $``$ 6 $`C_x`$ is: $`C^{}`$ $``$ $`(C_x\mathrm{\Delta }d)d^2`$, where $`\mathrm{\Delta }d`$ is the gap asymmetry in the $`x`$-direction; (3) the average of the potentials on opposing faces is same for all three axes and expressed as $`V_{x_0}`$ ($`(V_{x1}+V_{x2})/2=`$ 0.1 V); (4) the magnitude of the fluctuation in potential $`\delta V_i`$ is all identical and take the value of the average fluctuation of all the potentials for the three axes and the voltage to ground. We express the fluctuation as $`\delta V_{x_0}`$. It should be noted that we do not consider any cross-talk effects that arise in the sensitive axis due to forces applied to the other degrees of freedom.
### 6.1 Fluctuations in voltage imbalance and charge
$`f_{b1}`$ is associated only with sensing voltages but not the free charge, and its value is $`f_{b1}`$ $``$ 1.4 $`\times `$ 10<sup>-15</sup> m s<sup>-2</sup> Hz<sup>-1/2</sup>, by assuming $`V_{0g}`$ $``$ $`V_{x_0}V_g`$ = 0.05 V. This value is larger than the value used for LISA by a factor of five. Also, we use $`\delta V_d`$ = 1.0 $`\times `$ 10<sup>-4</sup> V Hz<sup>-1/2</sup> as fluctuation in voltage imbalance $`V_d`$ ($``$ $`V_{x1}V_{x2}`$ = 0.01 V). This value ($`\delta V_d`$) is one order of magnitude larger than the value used for LISA. Further, we assume that $`C_g/C`$ $``$ 1/6, as per .
$`f_{b2}`$ and $`f_{b3}`$ arise from fluctuations in the force due to the interaction between the net free charge and applied sensing voltages. $`f_{b2}`$ is due to fluctuation in voltage imbalance and its value is $`f_{b2}`$ $``$ $``$4.8 $`\times `$ 10<sup>-15</sup> m s<sup>-2</sup> Hz<sup>-1/2</sup>, where the net free charge on the proof mass is set to the nominal maximum built-up charge, described earlier. $`f_{b3}`$ is due to fluctuation in charge and its value is $`f_{b3}`$ $``$ $``$2.9 $`\times `$ 10<sup>-15</sup> m s<sup>-2</sup> Hz<sup>-1/2</sup>. The value of $`V_d`$ used here is larger than the one used in the error estimates for LISA by a factor of two. $`f_{b4}`$ is associated only with the free charge and its value is $`f_{b4}`$ $``$ 4.0 $`\times `$ 10<sup>-17</sup> m s<sup>-2</sup> Hz<sup>-1/2</sup>. In the estimation, we assume that the gap asymmetry in the $`x`$-direction is $`\mathrm{\Delta }d`$ = 10 $`\mu `$m, which is larger than by one order of magnitude.
### 6.2 Readout electronics
Readout electronics for the capacitive sensing will be similar to ones studied for LISA (e.g. ). For the standard resonant inductive-bridge scheme discussed in , the sources of disturbances due to readout electronics can be classified into two categories: imperfections in the capacitance bridge and the electric noise in the detecting circuit. The former includes fluctuations of inductance imbalance ($`\mathrm{\Delta }L/L`$), fluctuations of mutual inductance imbalance ($`\mathrm{\Delta }M/M`$), bias oscillator relative amplitude noise ($`\mathrm{\Delta }V/V_{M0}`$, where $`V_{M0}`$ is the 100 kHz bias voltage capacitively applied to the proof mass by injection electrodes on insensitive faces (the $`y`$ and/or $`z`$ faces) of the proof mass; the sensing electrodes are grounded) and bias oscillator phase noise ($`\mathrm{\Delta }\varphi `$). The latter includes current noise and thermal noise.
Using the expressions for these sources of the disturbances given in table 1 of , the main contributions to the acceleration disturbance at 0.1 mHz for ASTROD I are from $`\mathrm{\Delta }V/V_{M0}`$ ($``$ 1.5 $`\times `$ 10<sup>-15</sup> m s<sup>-2</sup> Hz<sup>-1/2</sup>) and $`\mathrm{\Delta }\varphi `$ ($``$ 9.1 $`\times `$ 10<sup>-16</sup> m s<sup>-2</sup> Hz<sup>-1/2</sup>), and the thermal noise ($``$ 3.8 $`\times `$ 10<sup>-17</sup> m s<sup>-2</sup> Hz<sup>-1/2</sup>) and the current noise ($``$ 1.3 $`\times `$ 10<sup>-17</sup> m s<sup>-2</sup> Hz<sup>-1/2</sup>). Contributions from fluctuations of $`\mathrm{\Delta }L/L`$ and $`\mathrm{\Delta }M/M`$ are insignificant. In the estimation, we use the following parameters based on experimental results reported in : $`\mathrm{\Delta }L/L10^4`$, $`\mathrm{\Delta }M/M6\times 10^8`$, $`\mathrm{\Delta }V/V_{M0}10^3`$ Hz<sup>-1/2</sup>, $`\mathrm{\Delta }\varphi 5.7\times 10^4`$ Hz<sup>-1/2</sup> and the turn ratio of the transformer $`n`$ = 1 (40 turns, inductance of 5 mH and a quality factor of 165), and a somewhat relaxed residual imbalance of the bridge $`\rho _{dc}\mathrm{\Delta }d/d=5\times 10^3`$. This relaxation has resulted in the dominant contributions from the bias oscillator relative amplitude noise and the phase noise. In , the thermal noise is dominant as they use $`\rho _{dc}10^4`$ in their estimation. The rss of these disturbances is $`f_{ba}`$ $``$ 1.8 $`\times `$ 10<sup>-15</sup> m s<sup>-2</sup> Hz<sup>-1/2</sup> at 0.1 mHz and listed in table 3.
### 6.3 Patch field voltage
Even when the sensing voltages are not applied on the electrodes, differences in local surface properties of the electrodes and the proof mass could lead a potential difference, patch-field voltage, between them . The charge fluctuations $`\delta q`$ result in acceleration disturbance $`f_{pe}`$ (table 3) through the patch field . By taking the average patch-filed voltage difference between opposing electrodes as $`V_{pe}`$=0.1 V , we obtain $`f_{pe}`$ = 2.9 $`\times `$ 10<sup>-14</sup> m s<sup>-2</sup> Hz<sup>-1/2</sup> at 0.1 mHz. This is the dominant contribution to the total acceleration disturbance of the ASTROD I proof mass. The LISA study team is investigating the possibility of measuring and compensating voltage differences across capacitors, during the mission commission process, to considerably better than 0.01 V .
### 6.4 Dielectric losses
Dielectric losses are thought to stem from surface contamination of electrodes and produces thermal voltage noise :
$$\delta v_{diel}=\sqrt{4k_BT_p\frac{\delta }{\omega C_x}}$$
(11)
where $`\delta `$ is loss angle. The upper limit of $`\delta `$ is reported to be $`10^5`$ for Al electrodes . For $`\delta =10^5`$, this voltage noise is about 6.6 $`\mu `$V Hz<sup>-1/2</sup> at 0.1 mHz and produces acceleration disturbance $`f_{dl}`$ (see table 3 for the expression), through residual dc bias voltage on electrodes, in the sensitive axis . By making the same assumption of the average potential difference $`V_0=0.1`$ V between a given electrode and the proof mass as , we obtain $`f_{dl}`$ $``$ 1.6 $`\times `$ 10<sup>-15</sup> m s<sup>-2</sup> Hz<sup>-1/2</sup>. In the error estimates for LISA , $`V_0=0.01`$ V is used.
### 6.5 Summary of the direct acceleration disturbances of the proof mass
By adding in quadrature, the rss of the proof-mass environmental acceleration disturbances ($`f_{nep}`$) and the sensor back-action acceleration disturbances ($`f_{nbp}`$) are 5.2 $`\times `$ 10<sup>-15</sup> m s<sup>-2</sup> Hz<sup>-1/2</sup> and 3.0 $`\times `$ 10<sup>-14</sup> m s<sup>-2</sup> Hz<sup>-1/2</sup>, respectively. Therefore, the total direct proof-mass acceleration disturbance ($`f_{np}`$) is 3.0 $`\times `$ 10<sup>-14</sup> m s<sup>-2</sup> Hz<sup>-1/2</sup> at 0.1 mHz, dominated by the sensor back-action acceleration disturbances.
## 7 Proof mass-spacecraft coupling
The stiffness $`K`$ is considered to stem from gravity gradients ($`K_{gg}`$), fluctuations in sensing capacitance and capacitance gradients ($`K_{s1}`$, $`K_{s2}`$ and $`K_{s3}`$), bias voltage ($`K_{s4}`$), patch field voltage ($`K_{s5}`$) and magnetic field gradients ($`K_{m1}`$ and $`K_{m2}`$). Table 4 gives a summary of estimated values for these sources of stiffness. The expressions used in the estimations are briefly described below. A detailed description on deviations of the expressions for $`K_{gg}`$, $`K_{s1}`$, $`K_{s2}`$, $`K_{s3}`$ and $`Ks_5`$ is given in Appendix A of <sup>8</sup><sup>8</sup>8$`K_{s5}`$ in text is noted as $`K_{s4}`$ in ..
### 7.1 Gravity gradients
For a given disturbing point mass ($`M_{dis}`$) at a distance $`x`$ from the center of mass of the proof-mass on the sensitive axis, the amplitude of the acceleration disturbance of the proof-mass caused by a positional fluctuation $`X_{ps}`$ is $`a_{gg}K_{gg}X_{ps}`$, where $`K_{gg}`$ is given in table 4.
Making the same assumption of $`M_{dis}`$ = 0.03 kg and $`x`$ = 0.05 m as , we obtain $`K_{gg}`$ $``$ 3.2 $`\times `$ 10<sup>-8</sup> s<sup>-2</sup>. For a more detailed analysis, the identification of the disturbing mass is necessary. Gravitational modeling for ASTROD I is in progress .
This disturbance arises from any positional fluctuation and is different from the gravity gradient caused by thermal distortion or motion ($`f_{gg}`$), which was discussed earlier.
### 7.2 Fluctuations in capacitive sensing
The expressions for $`K_{s1}`$, $`K_{s2}`$ and $`K_{s3}`$ can be obtained in the similar way as $`f_{b1}`$, $`f_{b2}`$, $`f_{b3}`$ and $`f_{b4}`$ in the previous section, under the following assumptions : fluctuations in the capacitances $`C_{x1}`$ and $`C_{x2}`$ and their derivatives produce disturbing forces in the $`x`$-direction, but not the fluctuations in other capacitances or their derivatives; we ignore the cross-coupling effects. Also, we assume $`\delta C_{x1}`$ = $`\delta C_{x2}(C_x/d)\delta x`$ and $`\delta C_{x1}^{}`$ = $`\delta C_{x2}^{}(C_x/d^2)\delta x`$.
$`K_{s1}`$ is due to the fluctuations in the Coulomb interaction between the charged proof mass and the image charges on the surrounding electrodes. $`K_{s2}`$ arises from interaction between the net free charge $`q`$ on the proof-mass and the average electrode voltages. $`K_{s3}`$ is due to the applied voltages across electrodes and the voltage difference across opposite electrodes. The estimated values for $`K_{s1}`$, $`K_{s2}`$ and $`K_{s3}`$ are given in table 4.
### 7.3 Bias voltage
By assuming that the sensing electrodes are grounded, the dominant term of the readout stiffness along the sensitive axis is given by $`K_{s4}`$ , listed in table 4. Assuming, as per , that the proof mass is biased to $`V_{M0}=0.6`$ V, $`K_{s4}`$ $``$ 3.1 $`\times `$ 10<sup>-7</sup> s<sup>-2</sup>.
### 7.4 Patch field voltage
By assuming a nominal patch-field voltage of $`V_{pe}`$ = 0.1 V and an overall multiplicative factor of $`\gamma `$ = 5 as per , the contribution due to the patch field to the proof mass-spacecraft coupling is $`K_{s5}`$ $``$ 1.2 $`\times `$ 10<sup>-9</sup> s<sup>-2</sup> (table 4).
### 7.5 Magnetic field gradients
The magnetic stiffness is given by $`K_m`$ = $`\frac{1}{m_p}\stackrel{}{}[\stackrel{}{}(\stackrel{}{M_p}\stackrel{}{B})]`$ . The expressions for dominant terms of the stiffness due to the induced magnetic moment and the remanent moment are given by $`K_{m1}`$ and $`K_{m2}`$ (as listed in table 4), respectively, where $`|\stackrel{}{}^2B_{sc}|12B_{sc}r_m^2`$ = 1.7 $`\times `$ 10<sup>-5</sup> T m<sup>-2</sup> for $`r_m`$ = 0.75 m. Their contributions to the total stiffness are insignificant.
### 7.6 Summary of the estimated values of stiffness
These contributions to the coupling constant $`K`$ are summarized in table 4. The rss of the coupling constant is 3.1 $`\times `$ 10<sup>-7</sup> s<sup>-2</sup>. This is slightly below the requirement for the total stiffness in LISA (4 $`\times `$ 10<sup>-7</sup> s<sup>-2</sup> ).
## 8 Requirements for the readout sensitivity and spacecraft control-loop gain
We have estimated values for the coupling constant $`K`$, the direct spacecraft acceleration disturbance $`f_{ns}`$ and the direct proof-mass acceleration disturbance $`f_{np}`$. By using the expression for the total acceleration disturbance of $`f_p`$ (equation (8)), we infer the requirements for the readout sensitivity $`X_{nr}`$ and the spacecraft control-loop gain $`u`$. In this process, we allocate an identical magnitude $`f_a`$ to each term of the expression; $`f_a^2=f_p^2/3`$. For ASTROD I, the noise goal is $`10^{13}`$ m s<sup>-2</sup> Hz<sup>-1/2</sup> at 0.1 mHz and, therefore, $`f_a=5.8\times 10^{14}`$ m s<sup>-2</sup> Hz<sup>-1/2</sup>. The second term of the expression is $`f_{np}`$, which was estimated to be 3.0 $`\times `$ 10<sup>-14</sup> m s<sup>-2</sup> Hz<sup>-1/2</sup> (table 3); it is about a factor of 2 smaller than the allocated requirement. With the estimated total stiffness $`K=3.1\times 10^7`$ s<sup>-2</sup>, we obtain $`X_{nr}1.9\times 10^7`$ m Hz<sup>-1/2</sup> from the first term and $`u3.8\times 10^5`$ from the last term of the expression.
From figure 1, one can see that the acceleration noise spectral density requirements for ASTROD I take its lowest value of about 0.4 $`\times `$ 10<sup>-13</sup> m s<sup>-2</sup> Hz<sup>-1/2</sup> at 0.3 mHz. At this frequency, $`f_a`$ becomes 2.3 $`\times `$ $`10^{14}`$ m s<sup>-2</sup> Hz<sup>-1/2</sup>. Therefore, the requirement for $`X_{nr}`$ becomes more stringent at 0.3 mHz: $`X_{nr}7.4\times 10^8`$ m Hz<sup>-1/2</sup>. As for the second term of $`f_p`$, $`f_{np}`$ is smaller at higher frequencies . Our estimate of $`f_{np}`$ is dominated by the contribution from $`f_{pe}`$, which scales as $`\nu ^1`$. Therefore, at 0.3 mHz, $`f_{np}`$ would be $``$ 1 $`\times `$ 10<sup>-14</sup> m s<sup>-2</sup> Hz, which is smaller than $`f_a`$ at 0.3 mHz. The last term scales as $`\omega ^2`$, and $`f_{ns}`$ is expected to be smaller at higher frequencies because of the $`\nu ^{1/3}`$ dependence of the fractional fluctuation in solar irradiance (figure 6 of ). Therefore, by assuming the same level of the thruster noise, the requirement for $`u`$ at 0.3 mHz would be less stringent than that at 0.1 mHz by a factor of $``$ 4.
In summary, the requirements for the readout sensitivity and the control loop gain for ASTROD I are $`X_{nr}7.4\times 10^8`$ m Hz<sup>-1/2</sup> and $`u3.8\times 10^5`$, respectively.
## 9 Comparison with LISA
Main relaxed parameter-values are listed in table 5 in comparison with LISA. The values for LISA are quoted from the current error estimates by Stebbins et al. , except the thruster noise quoted from . Recently the LISA requirements for the thruster noise and the residual gas pressure have been relaxed further from the values given in table 5 (this is a comment from one of the referees).
## 10 Summary, discussion and conclusions
We have estimated the spacecraft acceleration disturbance $`f_{ns}`$ (section 4), the proof-mass acceleration disturbances $`f_{np}`$ (sections 5 and 6) and the stiffness $`K`$ between the spacecraft and the proof mass (section 7). By using the expression (8), we have inferred the requirements for the readout sensitivity $`X_{nr}`$ and the control-loop gain $`u`$ (section 8).
Table 6 provides a summary of the estimated acceleration disturbances and the stiffness (section (b)), and the requirements (section (c)). The estimated total acceleration disturbance $`f_p`$ at 0.1 mHz (section (a) of table 6) is about 13 % less than the noise goal of $`10^{13}`$ m s<sup>-2</sup> Hz<sup>-1/2</sup>. We have compared the parameter values used in the estimation with LISA in section 9.
The total direct acceleration disturbance of the proof mass ($`f_{np}`$) at 0.1 mHz was estimated to be nearly a factor of two smaller than the requirement. This $``$ 50 % margin may be allocated for unknown disturbances or disturbances that would arise but have not been studied yet. These unestimated disturbances would be originated from, for instance, cross-talks in the capacitive sensing and magnetic damping of the proof mass. An estimate of acceleration disturbance due to magnetic damping for LISA is about 2 $`\times `$ 10<sup>-16</sup> m s<sup>-2</sup> Hz<sup>-1/2</sup> at 0.1 mHz . The contribution from the magnetic damping effect to ASTROD I would be in the similar order and insignificant.
The sensor back-action acceleration disturbances can be reduced by increasing the magnitude of the gap $`d`$. The total stiffness $`K`$ would be also reduced, for example, by a factor of four by changing the gap to 4 mm. The optimum design for the capacitive sensing is to be discussed based on results from the ongoing laboratory torsion balance experiment for ASTROD I .
Parameter values we used in this paper are mainly based on the results of studies done for LISA and LISA Pathfinder. This may be sufficient for the preliminary estimation. More accurate estimation would be obtained by carrying out the following works dedicated for ASTROD I: (a) modeling local magnetic fields of the spacecraft, (b) estimating the effective charging rate of the proof-mass, (c) estimating the cosmic-ray impact rate of the proof-mass, (d) estimating the micrometeorite impact effects, (e) thermal modeling of the proof-mass housing and the spacecraft, (f) gravitational modeling that includes thermal and non-thermal deformation of the spacecraft and the payload, (g) electrostatic modeling for the capacitive sensors and (h) estimating environmental factors (such as the interplanetary magnetic field, solar wind, solar radiation and cosmic rays) in the varying orbit (0.5 AU to 1 AU). Simulations to estimate charging rates for ASTROD I are in progress .
We have tentatively estimated acceleration disturbances for ASTROD I. This work has allowed us to set preliminary requirements for ASTROD I. To improve the current estimation, the disturbances that have not been studied yet have to be included and more detailed modeling works are necessary for ASTROD I. In comparison with LISA, requirements for ASTROD I can be largely relaxed. This will make the technological developments for ASTROD I less demanding to meet the drag-free requirements.
This work was funded by the National Science Council and the Foundation of Minor Planets of Purple Mountain Observatory. We thank A. Rรผdiger, S. Vitale, D. K. Gill, A. Pulido Patรณn, and, especially, the referees for useful information on acceleration disturbances and helpful comments on the manuscript.
## References
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# Electronic and magnetic properties of a hexanuclear ferric wheel
## I Introduction
Molecular magnetism has received great attention in the past few years, especially after a series of molecules has been discovered which show magnetism due to a molecular origin (see the review articlescaneschi1999 ; pilawa1999 ). One particular type of molecular magnets are the so-called ferric wheels where iron atoms are arranged ring-like. So far, wheels with 6 up to 18 iron atoms have been synthesizedcaneschi1999 . This has also triggered interest among theorists. The description of magnetism still poses challenges, and only in the last years ab-initio calculations of large molecular magnets have become feasible postnikovpsik .
The early studies employed wave-function based methods: for example, the exchange interaction in transition metal oxides was investigated at the level of periodic Hartree-Fock calculations Towler1994 . Usually, it is found that the Hartree-Fock approximation strongly underestimates the exchange coupling (e.g. Towler1994 ; ricart1995 ; catti1995 ). This was also demonstrated for smaller molecules, where demanding configuration interaction calculations were feasible: it was shown, that for a proper description of the magnetic interactions, highly correlated wave functions are necessaryFink1 ; Fink2 .
With the help of an appropriate embedding, the bulk could be approximatively described with clusters. Then, wave-function based correlation methods could be applied and a value in reasonable agreement with the experimental value of the exchange coupling was deduced degraaf ; moedl .
In the past few years, density functional calculations have become more and more widespread, and have also been applied to molecular magnets (e.g. postnikovpsik ; ruiz2003 ; zeng1 ). They have the advantage of including approximatively the effects of electronic correlation, at a very low computational expense. However, it is often observed that exchange couplings, computed at the density functional level (especially the local density approximation), strongly overestimate the experimental value. In a comparative study, it was shown that hybrid functionals, which include the exact (Fock) exchange, improve this and provide surprisingly accurate values iberio . This was also observed earlier in cluster calculations which were used to model the solid MartinIllas1997 .
The ferric wheel with six iron atoms is one of the smaller molecular magnets and thus ab-initio calculations can be performed at a relatively low expense, besides the approach using model Hamiltonians Normand ; Honecker . Various molecules have been synthesized Caneschi1995angew ; Caneschi1996 ; Abbati1997 ; Lascialfari1997 ; Pilawa1997 ; Saalfrank1997 ; Waldmann1999 ; Affronte1999 ; Cornia1999 ; Affronte2002 ; Pilawa2003 and techniques such as nuclear magnetic resonance, high-field DC and pulsed-field differential magnetization experiments, susceptibility measurements, high field torque magnetometry, inelastic neutron scattering and heat capacity measurements have been applied.
The experimental values for the exchange couplings, extracted from susceptibility measurements or from the singlet-triplet gap, are typically around $``$20 K caneschi1999 ; Waldmann1999 ; Cornia1999 ; Affronte1999 ; Pilawa2003 , although also larger values have been suggested (e. g. 38 KLascialfari1997 ). As there are still only few ab-initio calculationspostnikov2003 ; postnikov2004 , these systems are therefore an interesting object to study.
This article will deal with one of the six-membered iron ferric wheels Abbati1997 . In previous studies on a similar system with a gradient corrected functional, an overestimation of the calculated exchange coupling by a factor of $``$ 4 was found and it was argued that schemes such as LDA+$`U`$ or self-interaction correction might resolve this problempostnikov2003 ; postnikov2004 . As was argued above, hybrid functionals might be a different way to deal with this problem. We therefore studied this ring at the level of various functionals including hybrid functionals, and the traditional Hartree-Fock approach.
## II Method
All the calculations were done with the code CRYSTAL2003 crystal ; dovesi . We employed the Hartree-Fock (HF) method (in the variant of unrestricted Hartree-Fock), the local density approximation (LDA) with the correlation functional of Perdew and ZungerPZ , the gradient corrected functional of Perdew, Burke and Ernzerhof (PBE)PBE and the hybrid functional B3LYP (a functional with admixtures, amongst others, of functionals by Becke, Lee, Yang and Parr). As described later on in this section, we performed the calculations mainly on a molecular system. Therefore, we also used a basis set library for molecular basis setsBasistext rather then the CRYSTAL basis set library which contains essentially basis sets for periodic systems. We chose a \[5s4p2d\] basis set for iron Fe , where the diffuse exponent 0.041148 was omitted, a \[2s\] basis set for hydrogenditchfield , a \[3s2p\] basis set for carbonhehre where the exponent 0.1687144 was replaced with 0.25 to avoid convergence problems, a \[3s2p\] basis set for oxygenoxy with an additional $`d`$-exponent of 0.8 and a \[3s2p\] basis set for lithium oxy where the exponent 0.2 was omitted. After the mentioned modifications, the final basis sets were thus of the size \[4s3p2d\] (iron), \[3s2p\] (carbon), \[3s2p1d\] (oxygen), \[2s\] (hydrogen) and \[2s1p\] (lithium).
To explore the influence of further basis functions, various tests were performed: calculations with an additional $`sp`$-shell (exponent 0.04) and an additional $`d`$-shell (exponent 0.2) at the iron atom have been carried out, at the B3LYP and HF level. Also calculations with an additional $`sp`$-shell (exponent 0.12) at the oxygen atoms were performed at the B3LYP level. This diffuse exponent led to linear dependence problems at the HF level. However, a HF calculation where the outermost exponent 0.270006 for oxygen was replaced with two exponents (0.5, 0.2) was numerically stable. No significant changes for the computed exchange couplings were observed in any of these aforementioned tests. Therefore the basis sets, as described in the first paragraph of this section, can be considered reliable, and all the results were obtained with these basis sets, if not explicitly stated otherwise.
The properties of the ferromagnetic (FM) state (all spins parallel, total spin 30 $`\mu _B`$) and of the antiferromagnetic (AF) state (spins alternating up and down, total spin 0) were computed with each method. To obtain the local magnetic moments, a Mulliken population analysis was performed.
The calculations were performed on an isolated molecule where the full symmetry of the molecule was exploited (ferromagnet: $`C_{3i}`$, antiferromagnet: $`C_3`$). The geometry was chosen according to the measurements by Abbati et al Abbati1997 and is displayed in figure 1. The isolated molecule was charged with +1, as lithium is ionized and charge neutrality is restored by a PF$`{}_{}{}^{}{}_{6}{}^{}`$ ion for this particular ferric wheel. To assess the influence of the centered lithium ion it was also removed in some of the calculations.
## III Results
The total energies, the difference in total energies of the ferromagnet and the antiferromagnet and the exchange parameters $`J`$ are given in table 4. To obtain the exchange parameters J from the calculation, the spin Hamiltonian of the Ising model was used:
$$\text{H}=J\underset{i=1}{\overset{5}{}}S_iS_{i+1}JS_6S_1$$
The variable $`i`$ corresponds to the site index of the iron atoms, which in this model possess a spin of $`S=5/2`$. With this Hamiltonian, the energy difference between ferromagnet(FM) and antiferromagnet(AF) then corresponds to:
$$E_{tot}^{FM}E_{tot}^{AF}=12JS^2$$
Therefore, the exchange parameter J results from the difference of large numbers. Still the results are numerically stable, as can be seen in table 4. For the particular ferric wheel under consideration, the experimental value for the exchange coupling was determined to be -21 K Abbati1997 ; caneschi1999 .
This can now compared with the computed values. One interesting result is the ferromagnetic exchange parameter obtained at the Hartree-Fock level (+7 K ) which does not reproduce the real antiferromagnetic nature of the molecule. This result was stable $`(\pm 2K)`$ with respect to all possible variations in the basis set as mentioned in section II, and also when the central lithium ion was removed.
Such a problem did not show up with any of the density functionals. The magnitude of $`J`$ of the B3LYP calculation (-31 K) is about 50$`\%`$ larger than the experimental value. In the PBE and LDA calculations, it is even overestimated by a factor of $``$ 5 (PBE: -108 K, LDA: -120 K). This agrees well with an earlier calculation for a similar ferric wheel by Postnikov et al postnikov2003 ; postnikov2004 with the PBE functional, where $`J`$ was also largely overestimated (-80 K). The authors argued, that the $`d`$-orbitals in density functional calculations with functionals such as LDA or PBE, were not sufficiently localized to compare with experiment, and therefore the resulting exchange parameters would overestimate the experimental values. They therefore suggested to use methods such as LDA+$`U`$ or self-interaction corrections to improve this shortcoming. The B3LYP functional may be viewed as another possibility to remedy this problem, as it interpolates between the pure Hartree-Fock approach (with an unscreened Coulomb interaction and the correct Fock exchange, but without electronic correlation), and standard functionals (including electronic correlation, but too delocalized $`d`$-orbitals). This was already demonstrated for bulk NiO iberio , where the B3LYP value for the exchange couplings was also found to be about 50 % larger than the experimental value, and LDA overestimated also by a factor of $`>4`$. Also, the total density of states was calculated at the HF and B3LYP level and is shown in figure 2. Note that these calculations were carried out as bulk calculations. The HOMO-LUMO gap at the HF level is about ten times larger than at the B3LYP level, but there is no significant difference between the ferromagnetic and the antiferromagnetic state at one level of theory. However, the smaller gap will favor hybridization and thus the more delocalized situation.
Comparing the results, the wrong sign for J at the HF level is rather surprising as it was usually observed that the Hartree-Fock result underestimated the experimental value, but with the correct sign, with few exceptionsruiz2003 . A possible source for the discrepancy can be found in the geometry of the molecule. For a system where the magnetic centers are in line with the bonding bridge atom, consideration of correlation leads to a strong increase of the antiferromagnetic exchange parameter as it was found for NiO by de Graaf et al degraaf and for other complexes by Fink et alFink1 . In our case, the Fe-O-Fe bonding angle is in the range of 100 degreesAbbati1997 and for such systems (and similarly for other ferric wheels, see e.g. Waldmann et alwaldmann2001 ) the exchange parameter is supposed to be less antiferromagnetic or even ferromagnetic compared to in-line geometries which is a result of the ordering of the magnetic orbitals according to the Goodenough-Kanamori ruleskahn . In addition, due to the to large HOMO-LUMO gap, the system is too ionic at the HF level and thus the overlap of the orbitals is underestimated, which also has an impact on the magnitude of the computed exchange interaction. Thus the small antiferromagnetic coupling and the underestimation of the exchange parameter due to the neglect of correlation at the HF level may be the reason for the wrong sign.
The central lithium atom has an electrostatic influence to the molecule and the lithium ionic radius affects the diameter of the ferric wheel. Therefore, when the lithium is removed and the geometry input is not changed as in our case, no impact on the exchange parameter is expected.
To analyze the magnetic states, the magnetic moments were computed. They are displayed in table 4. A decrease of the local magnetic moment of the iron atoms from HF over B3LYP and PBE to LDA is observable, where the moment becomes more delocalized and is transferred to the surrounding oxygen atoms. This is particularly obvious in the ferromagnetic state. In the antiferromagnetic state, the iron magnetic moments are virtually the same as in the ferromagnetic state, apart from the sign. However, the bridge oxygens carry virtually no magnetic moment, as there is an overlap of up and down spin density at these sites, originating from two neighboring irons, which as a whole cancels. For the same reason there is no magnetic moment at the centered lithium, and even in the ferromagnetic case the lithium spin is virtually zero.
This feature is reflected in the spin-density plots in figure 3 and in larger resolution in figure 4. For a more detailed study, an one-dimensional plot of the spin-density at the LDA level along a line from O(apical) via Fe, O(bridge) and Fe to O(apical) was performed (figure 5). The iron atoms subsist in a high-spin state $`3d_{}^5\text{d}_{}^0`$, but except for the HF method the moment is closer to $`4\mu _B`$ as it was also obtained by Postnikov et al for a slightly different moleculepostnikov2003 ; postnikov2004 . As can be seen in table 4, the local magnetic moment of the apical oxygen atoms increases from HF over B3LYP and PBE to LDA. In addition, there is a local magnetic moment at the carbon C(1) atom (see figure 1 and table 4) which is about 4 times (HF) to 2 times (LDA) smaller than the local magnetic moment of the apical oxygen atom. Corresponding to the moment of the C(1) atom there is a small magnetic moment at the carbon atoms C(2) and C(3), but with opposite sign and a much smaller value. Further on there is actually a moment at the first carbon atom (next to the carbon atoms C(2) and C(3)) of each $`C_6H_5`$ ring. Its value is of the same order as the value of the C(2) and C(3) carbon atoms moments, but with opposite sign. In the ferromagnetic state, there is a small magnetic moment at the C(4) bridge carbon atoms corresponding to the magnetic moment of the bridge oxygen atoms which is also, as already mentioned above, increasing from HF over B3LYP and PBE to LDA.
The results from the Mulliken population analysis of the iron, oxygen (bridge and apical), lithium and carbon atoms are given in tables 4 and 4. The net charge of iron is between 1.27 and 2.16, at the various levels of theory, and thus far away from a formal chargeAbbati1997 of +3. Again, Hartree-Fock gives a very localized picture with the largest net charge (2.16), and LDA a delocalized picture with an iron charge of only 1.27. Oxygen carries essentially a single negative charge ($``$ -0.8 in LDA, $``$ -1.1 in HF), the carbon atoms which have oxygen as a neighbor all have donated charge (0.4 - 0.8 $`|e|`$, depending on the site and the level of theory). The lithium charge is between 0.4 (LDA) and 0.8 (HF). A plot of the difference of the charge density was performed between the values at the LDA and HF level (figure 6). There is a positive charge density difference at the iron atoms because of the smaller charge at the LDA level compared to HF, the results for the other atoms follow from tables 4 and 4. This plot thus visualizes the more covalent picture obtained at the LDA level. The stronger covalency also becomes obvious in the Mulliken overlap population: for example, the overlap population between Fe and the neighbouring oxygens atoms is 0.11-0.12 at the LDA level, but only 0.07-0.08 at the HF level.
## IV Summary
Ab-initio calculations for a hexanuclear ferric wheel Abbati1997 were performed. The exchange coupling parameter $`J`$ was computed at various levels of theory, with the B3LYP functional providing the best agreement with experiment (-31 K, experiment: -21 K). LDA and PBE were found to grossly overestimate the exchange coupling parameter due to the too large delocalization of the $`d`$-orbitals. Surprisingly, the calculation at the Hartree-Fock level led to a ferromagnetic coupling. The electronic population and the spin densities were calculated. The total spin is distributed over several sites, besides the iron atom all of the neighboring oxygens carry some spin. The delocalization increases from HF over B3LYP, PBE to LDA which shows the strongest delocalization and thus the smallest magnetic moment on the iron site.
## V Acknowledgments
Most of the calculations were performed at the compute-server cfgauss (Compaq ES 45) of the data processing center of the TU Braunschweig. The geometry plot of the molecule was performed with MOLDENmolden .
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# Numerical Analysis of Coherent Many-Body Currents in a Single Atom Transistor
## I Introduction
The recently proposed Single Atom Transistor (SAT) setup sat provides new opportunities to experimentally examine the coupling of a spin-1/2 system with bosonic and fermionic modes. Such couplings form fundamental building blocks in several areas of physics. For example, atoms passing through a cavity can allow the quantum non-demolition (QND) readout of single-photon states in quantum optics cqnd , and in solid state physics, such systems occur in Single Electron Transistors set , in studies of electron counting statistics levitov and in the transport of electrons past impurities such as quantum dots cazmas .
In the SAT setup, which was motivated by the significant experimental advances made recently with cold atoms in 1D esslinger1d ; bloch1d ; porto1d , a single spin-1/2 impurity atom, $`q`$, is used to switch the transport of a gas of cold atoms in a 1D optical lattice (Fig. 1). The impurity atom, which can encode a qubit on two internal spin states, is transparent to a gas of probe atoms in one spin state (the โonโ state), but acts as a single atom mirror in the other (the โoffโ state), prohibiting transport via a quantum interference mechanism (Fig. 1). Observation of probe atoms that are initially situated to one side of the impurity, and which can constitute either a 1D degenerage Bose or Fermi gas, can then be used as a QND measurement qnd of the qubit state of the impurity atom $`|\psi _q=\alpha |+\beta |`$ sat (see Fig. 1).
The long coherence times associated with atoms in optical lattices allow many-body effects to contribute coherently to the transport properties over longer timescales than is observed in other systems where bosonic and fermionic modes couple to a spin 1/2 system. This produces novel physics in which the current of atoms passing the impurity, especially in a regime of weak coupling between probe atoms and impurity, is sensitive to interactions between the probe atoms sat . These effects could be directly observed in experiments, for example, via measurements of the density of probe atoms on each site of the impurity atom as a function of time.
In this article we present a detailed numerical analysis of these currents, making use of recently developed numerical methods vidal to calculate the dynamics of the bosonic probe atoms by directly integrating the many-body Schrรถdinger equation in 1D on an adaptively truncated Hilbert space. When these currents are compared to analytical calculations of transmission coefficients for single particles passing the impurity atom and the related currents for a non-interacting 1D Fermi gas, significant interaction effects are observed, as first discussed in Ref. sat . Here we provide new insight into the time dependence of these currents, and what conclusions can be drawn from our numerical results on different timescales. We then calculate the initial currents for atoms at zero temperature diffusing past the impurity (where the initial mean momentum of the 1D gas, $`\widehat{k}_{t=0}=0`$, with $`\widehat{k}`$ is the operator corresponding to the quasi-momentum in the lowest Bloch band and $`t`$ the time), and explore the effects observed for different interaction strengths of bosonic probe atoms. We then also investigate the currents for fermions and bosons when the probe atoms are initially kicked ($`\widehat{k}_{t=0}0`$). This study is complementary to the analytical study of the SAT that is given in recent article by Micheli et al. satandi .
In section II we discuss the basic physics of the SAT, and give a summary of the dynamics found in sat for single particles and non-interacting fermions. Then we present in detail the numerical techniques that we use to compute the exact time evolution of the many-body 1D system. The time-dependence of the resulting currents is discussed in section III, followed by a presentation of the values of the initial steady state currents, both in the diffusive ($`\widehat{k}_{t=0}=0`$) and kicked ($`\widehat{k}_{t=0}0`$) regimes. The conclusions are then summarised in section IV.
## II Overview
### II.1 The Single Atom Transistor
#### II.1.1 The System
As described in section I, we consider probe atoms $`b`$, which are loaded into an optical lattice hubbardtoolbox ; greiner ; oltheory ; esslingerfermions with strong confinement in two dimensions, so that the atoms are restricted to move along a lattice in 1D. The probe atoms are initially situated to the left of a site containing an impurity atom $`q`$, which is trapped independently (by a species or spin-dependent sdol potential), fixing it to a particular site while the probe atoms are free to move. In order to produce the โonโ and โoffโ states of the SAT, we must appropriately engineer the effective spin-dependent interaction between the probe atoms and the impurity, $`H_{\mathrm{int}}=_\sigma U_{\mathrm{eff},\sigma }\widehat{b}_0^{}\widehat{b}_0\widehat{q}_\sigma ^{}\widehat{q}_\sigma `$. Here, $`\widehat{b}_i^{}`$ and $`\widehat{q}^{}`$ are second-quantised creation operators for the $`b`$ and $`q`$ atoms respectively, obeying the standard commutation (anti-commutation) relations for bosons (fermions) and the site index $`i`$ is chosen so that the impurity is on site $`i=0`$. These interactions can be controlled using either a magnetic juliennelattice ; magfeshbach or optical optfeshbach Feshbach resonance. For simplicity we discuss the case of an optical Feshbach resonance, depicted in Fig. 2. Here, lasers are used to drive a transition from the atomic state $`\widehat{b}_0^{}\widehat{q}_\sigma ^{}|\mathrm{vac}`$ via an off-resonant excited molecular state to a bound molecular state back in the lowest electronic manifold $`\widehat{m}_\sigma ^{}|\mathrm{vac}`$ on the impurity site, $`i=0`$ (see Fig. 2). The two-photon Rabi frequency for this process is denoted $`\mathrm{\Omega }_\sigma `$ and the Raman detuning $`\mathrm{\Delta }_\sigma `$, and throughout this article we use units with $`\mathrm{}=1`$.
#### II.1.2 Single Atoms
We consider initially a single probe atom passing the impurity. If the coupling to the molecular state is far off resonance ($`\mathrm{\Omega }_\sigma \left|\mathrm{\Delta }_\sigma \right|`$), the effect of the Feshbach resonance is to modify the interaction between the $`b`$ and $`q`$ atoms in the familiar manner, with $`U_{\mathrm{eff}}=U_{qb}+\mathrm{\Omega }_\sigma ^2/\mathrm{\Delta }_\sigma `$. This can be used to screen the background interaction between these atoms, $`U_{qb}`$, so that the โonโ state of the SAT ($`U_{\mathrm{eff}}=0`$) can be produced by choosing $`\mathrm{\Delta }_{}=\mathrm{\Omega }_{}^2/U_{qb}`$.
If the coupling is resonant ($`\mathrm{\Delta }_{}=0`$), then the physical mechanism is different, and the passage of a probe atom $`b`$ past the impurity is blocked by quantum interference. The mixing of the unbound atomic state and the molecular state on the impurity site produces two dressed states
$$\frac{1}{\sqrt{2}}\left(\widehat{b}_0^{}\widehat{q}_{}^{}|\mathrm{vac}\pm m_{}^{}|\mathrm{vac}\right),$$
(1)
with energies
$$\epsilon _\pm =\frac{U_{qb}}{2}\pm \left(\frac{U_{qb}^2}{4}+\mathrm{\Omega }_{}^2\right)^{1/2}.$$
(2)
The two resulting paths for a particle of energy $`\epsilon `$ then destructively interfere so that when $`\mathrm{\Omega }_{}J`$, where $`J`$ is the normal tunneling amplitude between neighbouring lattice sites, and $`U_{qb}=0`$, the effective tunnelling amplitude past the impurity (see Fig. 3) is
$$J_{\mathrm{eff}}=\left(\frac{J^2}{\epsilon +\mathrm{\Omega }_{}}\frac{J^2}{\epsilon \mathrm{\Omega }_{}}\right)0.$$
(3)
This is reminiscent of the interference effect which underlies Electromagnetically Induced Transparency eit , and corresponds to the effective interaction $`U_{\mathrm{eff}}\mathrm{}`$ required for the โoffโ state of the SAT.
In Refs. sat ; satandi , the Lippmann-Schwinger equation is solved exactly for scattering from the impurity of an atom $`b`$ with incident momentum $`k>0`$ in the lowest Bloch-band, where the energy of the particle $`\epsilon (k)=2J\mathrm{cos}(ka)`$, with $`a`$ the lattice spacing. The resulting transmission probabilities $`T(p)`$ are in the form of Fano Profiles satfano . For $`\mathrm{\Omega }_\sigma J`$ these have a minimum corresponding to complete reflection for $`\epsilon (k)=\mathrm{\Delta }_\sigma `$ and complete transmission for $`\epsilon (k)=\mathrm{\Delta }_\sigma \mathrm{\Omega }_\sigma ^2/U_{qb}`$. For $`\mathrm{\Omega }_\sigma >4J`$, the transmission coefficients are approximately independent of $`k`$, and so complete transparency of the impurity atom is obtained for $`\mathrm{\Delta }_\sigma =\mathrm{\Omega }_\sigma ^2/U_{qb}`$ and complete blocking of the incident atoms for $`\mathrm{\Delta }=0`$.
#### II.1.3 Many Atoms
The treatment of this system for many atoms is similar to the single atom case, but the motion of the probe atoms in the lattice, except on the impurity site, is governed by a (Bose-) Hubbard Hamiltonian oltheory . As the two spin channels for the impurity atom, $`q`$ can be treated independently, we will consider only a single spin channel $`q_\sigma `$, and drop the subscript in the notation throughout the remainder of the article sat . The Hamiltonian for the system is then given (with $`\mathrm{}1`$) by $`\widehat{H}=\widehat{H}_b+\widehat{H}_0`$, with
$`\widehat{H}_\mathrm{b}`$ $`=`$ $`J{\displaystyle \underset{ij}{}}\widehat{b}_i^{}\widehat{b}_j+{\displaystyle \frac{1}{2}}U_{bb}{\displaystyle \underset{j}{}}\widehat{b}_j^{}\widehat{b}_j\left(\widehat{b}_j^{}\widehat{b}_j1\right),`$
$`\widehat{H}_0`$ $`=`$ $`\mathrm{\Omega }(\widehat{m}^{}\widehat{q}\widehat{b}_0+\mathrm{h}.\mathrm{c})\mathrm{\Delta }\widehat{m}^{}\widehat{m}`$ (4)
$`+U_{qb}\widehat{b}_0^{}\widehat{q}^{}\widehat{q}\widehat{b}_0+U_{bm}\widehat{b}_0^{}\widehat{m}^{}\widehat{m}\widehat{b}_0.`$
Here, $`H_\mathrm{b}`$ gives a Hubbard Hamiltonian for the $`b`$ atoms with tunnelling matrix elements $`J`$, and collisional interactions $`U_{bb}`$. For fermions, $`U_{bb}=0`$, whereas for bosons $`U_{bb}=4\pi \mathrm{}^2a_{bb}d^3๐ฑ\left|\mathrm{w}_j(๐ฑ)\right|^4/m_b`$, with $`\mathrm{w}_j(๐ฑ)`$ the Wannier-function on site $`j`$, and $`a_{bb}`$ and $`m_b`$ the scattering length and mass of $`b`$ atoms respectively. $`H_0`$ describes the dynamics in the presence of the impurity on site 0, where atoms $`b`$ and $`q`$ are converted to a molecular state with effective Rabi frequency $`\mathrm{\Omega }`$ and detuning $`\mathrm{\Delta }`$, and the final two terms describe background interactions, $`U_{\alpha \beta }`$ for two particles $`\alpha ,\beta \{q,b,m\}`$, which are typically weak and will be neglected in our treatment. This single-band model is valid in the limit for $`U_{\alpha \beta },J,\mathrm{\Omega },\mathrm{\Delta }\omega `$, where $`\omega `$ is the energy separation between Bloch bands, an inequality which is fulfilled in current experiments. The robustness of the SAT with respect to loss processes is discussed in sat .
In the rest of this article, we will study the current of atoms past the impurity site that develops as a function of time, and how this current depends on the interaction between probe atoms and on interactions between the probe atoms and the impurity.
### II.2 Atomic Currents through the SAT
To analyse the case of many atoms passing the impurity site, we consider the probe atoms $`b`$ to be prepared initially to the left of the impurity, in a ground state corresponding to a 1D box potential. The current of atoms passing the impurity is $`I(t)=d\widehat{N}_R/dt`$, where $`N_R`$ is the mean number of atoms to the right of the impurity, $`N_R=_{j>0}\widehat{b}_j^{}\widehat{b}_j`$. For a sufficiently large number of atoms in the initial cloud, this current is generally found to rapidly settle into an initial steady state current, $`I_{\mathrm{ss}}`$, on relatively short timescales ($`tJ1`$) (see section III.1 for further discussion of steady state currents for bosons).
For a non-interacting Fermi gas at zero temperature, the currents can be calculated exactly when $`U_{bm}=U_{qb}`$, as the equations of motion are linear. Scattering from the impurity occurs independently for each particle in the initial Fermi sea, and after a short transient period of the order of the inverse tunnelling rate $`1/J`$, a steady state current $`I_{\mathrm{ss}}`$ is established. This can be calculated either by integrating the single-particle transmission probabilities sat ; satandi or by direct numerical integration of the Heisenberg equations.
For a non-interacting and very dilute Bose gas, the situation will be identical to considering a single particle. However, for higher densities, many-boson effects become important, and additionally for non-zero interactions the situation becomes even more complicated. In the limit $`U_{bb}/J\mathrm{}`$ in 1D (the Tonks gas regime) it is usually possible to replace the bosonic operators $`\widehat{b}_i,\widehat{b}_i^{}`$ by fermionic operators $`\widehat{f}_i,\widehat{f}_i^{}`$ using a Jordan-Wigner transformation jordanwigner . However, in this case the resulting Hamiltonian,
$`\widehat{H}`$ $`=`$ $`J{\displaystyle \underset{ij}{}}\widehat{f}_i^{}\widehat{f}_j\mathrm{\Delta }\widehat{m}^{}\widehat{m}+(1)^{\widehat{N}_L}\mathrm{\Omega }(\widehat{m}^{}\widehat{q}\widehat{f}_0+\mathrm{h}.\mathrm{c})`$ (5)
$`+U_{qb}\widehat{f}_0^{}\widehat{q}^{}\widehat{q}\widehat{f}_0+U_{bm}\widehat{f}_0^{}\widehat{m}^{}\widehat{m}\widehat{f}_0,`$
contains a nonlinear phase factor resulting from the coupling on the impurity site, $`(1)^{\widehat{N}_L}`$, where $`\widehat{N}_L=_{j<0}\widehat{f}_j^{}\widehat{f}_j`$ is the operator for the number of atoms to the left of the impurity site. For $`\mathrm{\Omega }=0`$, the boson currents are exactly the same as the currents for noninteracting fermions as $`\widehat{b}_i^{}\widehat{b}_i=\widehat{f}_i^{}\widehat{f}_i`$. For finite $`\mathrm{\Omega }`$ it is not clear what role the phase factor will play in determining the system dynamics, although for sufficiently large $`\mathrm{\Omega }J`$ we again expect very little current to pass the impurity.
Thus, for the intermediate regime $`\mathrm{\Omega }J`$, and for the case of finite interaction strength $`U/J`$ there are no known analytical solutions for the currents. For this reason, we specifically study these regimes in this paper, using near-exact numerical methods.
### II.3 Time-Dependent Numerical Algorithm for 1D Many-Body Systems
The algorithm that we use to compute the time evolution of our many body system for bosonic probe atoms was originally proposed by Vidal vidal . This method allows near-exact integration of the many body Schrรถdinger equation in 1D by an adaptive decimation of the Hilbert space, provided that the Hamiltonian couples nearest-neighbour sites only and that the resulting states are only โslightly entangledโ (this will be explained in more detail below). Recently both this algorithm dissipativevidal , and similar methods proposed by Verstrate and Cirac dissipativecirac have been generalised to the treatment of master equations for dissipative systems and systems at finite temperature, and progress has been made applying the latter method to 2D systems cirac2d .
In 1D, these methods rely on a decomposition of the many-body wavefunction into a matrix product representation of the type used in Density Matrix Renormalisation Group (DMRG) calculations dmrgreview , which had previously been widely applied to find the ground state in 1D systems. The time dependent algorithms have now been incorporated within DMRG codes dmrgvidal , and also been used to study the coherent dynamics of a variety of systems vidalexamples . In our case, we write the coefficients of the wavefunction expanded in terms of local Hilbert spaces of dimension $`S`$,
$$|\mathrm{\Psi }=\underset{i_1i_2\mathrm{}i_M=1}{\overset{S}{}}c_{i_1i_2\mathrm{}i_M}|i_1|i_2\mathrm{}|i_M,$$
(6)
as a product of tensors
$$c_{i_1i_2\mathrm{}i_M}=\underset{\alpha _1\mathrm{}\alpha _{M1}}{\overset{\chi }{}}\mathrm{\Gamma }_{\alpha _1}^{[1]i_1}\lambda _{\alpha _1}^{[1]}\mathrm{\Gamma }_{\alpha _1\alpha _2}^{[2]i_2}\lambda _{\alpha _2}^{[2]}\mathrm{\Gamma }_{\alpha _3\alpha _4}^{[2]i_2}\mathrm{}\mathrm{\Gamma }_{\alpha _{M1}}^{[M]i_M}.$$
(7)
These are chosen so that the tensor $`\lambda _\alpha ^{[l]}`$ specifies the coefficients of the Schmidt decomposition nielsonchuang for the bipartite splitting of the system at site $`l`$,
$$|\psi =\underset{\alpha =1}{\overset{\chi _l}{}}\lambda _\alpha ^{[l]}|\varphi _\alpha ^{[1\mathrm{}l]}|\varphi _\alpha ^{[l+1\mathrm{}M]},$$
(8)
where $`\chi _l`$ is the Schmidt rank, and the sum over remaining tensors specify the Schmidt eigenstates, $`|\varphi _\alpha ^{[1\mathrm{}l]}`$ and $`|\varphi _\alpha ^{[l+1\mathrm{}M]}`$. The key to the method is two-fold. Firstly, for many states corresponding to a low-energy in 1D systems we find that the Schmidt coefficients $`\lambda _\alpha ^{[l]}`$, ordered in decreasing magnitude, decay rapidly as a function of their index $`\alpha `$ (this is what we mean by the state being โslightly entangledโ) vidal . Thus the representation can be truncated at relatively small $`\chi `$ and still provide an inner product of almost unity with the exact state of the system $`|\mathrm{\Psi }`$. Secondly, when an operator acts on the local Hilbert state of two neighbouring sites, the representation can be efficiently updated by changing the $`\mathrm{\Gamma }`$ tensors corresponding to those two sites, a number of operations that scales as $`\chi ^3S^3`$ for sufficiently large $`\chi `$ vidal . Thus, we represent the state on a systematically truncated Hilbert space, which changes adaptively as we perform operations on the state.
In order to simulate the time evolution of a state, we perform a Suzuki-Trotter decomposition trotter of the time evolution operator $`\mathrm{exp}(\mathrm{i}\widehat{H}t)`$, which is applied to each pair of sites individually in small timesteps $`\delta t`$. Initial states can also be found using an imaginary time evolution, i.e., the repeated application of the operator $`\mathrm{exp}(\widehat{H}\delta t)`$, together with renormalisation of the state.
In this paper, results are not only produced using the original algorithm as presented in vidal , but also using an optimised version in which the Schmidt eigenstates are forced to correspond to fixed numbers of particles. This allows us to make use of the total number conservation in the Hamiltonian to substantially increase the speed of the code, and also improve the scaling with $`\chi `$ and $`S`$. With this number conserving code we are able to compute results with much higher values of $`\chi `$, however we also find that for insufficiently large $`\chi `$, the results from this code become rapidly unphysical, in contrast to the original code (see section III.1).
In implementations of this method we vary the value of $`\chi `$ to check that the point at which the representation is being truncated does not affect the final results. A useful indicator for convergence of the method is the sum of the Schmidt coefficients discarded in each time step, although in practice the convergence of calculated quantities (such as the single particle density matrix, $`\widehat{b}_i^{}\widehat{b}_j`$) are normally used. This is also discussed further in section III.1
For bosons on an optical lattice we must also choose the dimension $`S`$ of the local Hilbert space, which corresponds to one more than the maximum number of atoms allowed on one lattice site. For simulation of the SAT, we allow a variable dimension of the local Hilbert space $`S_l`$, as we must consider the state of the molecule on the impurity site in addition to the probe atoms. Allowing such a variable dimension dramatically reduces the simulation time, which scales as $`\chi ^3_lS_l^3`$ when $`\chi S`$, and scales proportional to $`S^4`$ when $`\chi `$ is small. For a Bose gas with finite $`U/J`$ we usually take $`S_l=6`$ away from the impurity site, and $`S_0=12`$ on the impurity site, whereas simulations of a Tonks gas can be performed with $`S_l=2`$ away from the impurity site and $`S_0=4`$ on the impurity site.
## III Numerical Results
In section III.1 we discuss the time dependence of the current for bosons and the applicability of our numerical methods in different regimes. We establish the existence of an initial steady state current, $`I_{SS}`$ that appears on a timescale $`tJ1`$, and discuss the observation of a second steady state current $`I_0`$, observed in some cases on a timescale $`tJ10`$. In sections III.2 and III.3 we then present our numerical results for $`I_{SS}`$ for the case where the initial cloud diffuses past the impurity site, and the case where the initial cloud is kicked respectively.
We are primarily interested in the behaviour of the current through the SAT when it is used in the โoffโ state, i.e., we choose $`\mathrm{\Delta }=0`$. To enhance clarity of the results, we also choose $`U_{bq}=U_{bm}=0`$.
In each case, we considered an initial cloud of between $`N=1`$ and $`N=30`$ atoms, confined on $`M=30`$ lattice sites situated immediately to the left of the impurity site. The initial state used corresponds to the ground state, $`|\varphi _0`$ of a Bose-Hubbard model with a box trap.
Our total grid for the time evolution consisted of 61 lattice sites, with the 30 rightmost sites initially unoccupied, and the results we present are, except for very small systems, independent of the size of the initial cloud and of the grid size. Fermionic results are derived from exact integration of the Heisenberg equations of motion, whereas bosonic results are near-exact simulations as described in section II.3.
### III.1 Time Dependence of the current for bosonic probe atoms
The mean number of probe atoms on the right of the impurity, $`N_R`$ is plotted as a function of time, $`t`$, in Fig. 4 for a Tonks gas ($`U/J\mathrm{}`$) with $`\mathrm{\Omega }/J=1`$ and initial state of density $`n=N/M=1`$. These results were calculated with the original simulation algorithm, and it is clear from the figure that the current settles into an initial steady state value $`I_{SS}`$ on the timescale $`tJ1`$. However, as is typical for bosonic probe atoms with $`n>0.5`$, there exists a knee in the curve at a time $`t_{\mathrm{knee}}(\chi )`$, leading to a new and final steady state current, which we will denote $`I_0`$. The time $`t_{\mathrm{knee}}(\chi )`$ depends on the initial density, $`n`$, and coupling, $`\mathrm{\Omega }`$, and as can be seen from this figure, we require a high value of $`\chi `$ to find the exact time. For $`n=1,\mathrm{\Omega }/J=1`$, $`t_{\mathrm{knee}}(\chi )`$ appears to converge to a value between $`tJ=9`$ and $`tJ=12`$ as $`\chi `$ is increased. It is clear that significant level of correlation, or entanglement between the left and right hand side of the system (in the sense of the number of significant Schmidt eigenvalues for a bipartite splitting) are involved in determining the dynamics leading the to knee. However, the actual value of the steady state current $`I_0`$ appears to converge for much lower values of $`\chi `$ and there is essentially no change in this result from $`\chi =10`$ to $`\chi =70`$.
The interpretation of these results is more complex when they are compared with similar results from the new, number conserving version of our code. In Fig. 5 we observe that the behaviour diverges at the same value of $`t_{\mathrm{knee}}(\chi )`$, and even for $`\chi =300`$, the value of $`t_{\mathrm{knee}}(\chi )`$ has only shifted a little further from where it was observed for $`\chi =70`$ with the original version of the code. This confirms that the dynamics on this timescale are dominated by the significant level of correlation, or entanglement between the left and right hand side of the system.
In contrast to the steady state current $`I_0`$ obtained using the original code, though, the current in the number conserving simulations rapidly approaches $`0`$, even for $`\chi =300`$. As can be seen from the dotted line in Fig. 6, this behaviour occurs when the maximum sum of squares of the Schmidt coefficients being discarded in each timestep, $`\epsilon _\lambda =_{\beta >\chi }\lambda _\beta ^2`$, reaches a steady value on the order of $`10^7`$, indicating that the simulation results from the number conserving code are probably not valid for $`t>t_{\mathrm{knee}}`$. Indeed, we observe the same behaviour from the new simulation code with $`\mathrm{\Omega }=0`$, where we know from Eq. 5 that the time dependent current $`I(t)`$ is equal to that for fermions, and should not decrease in this manner (see currents for fermions in Ref. sat ). Interestingly, the original code, which produces the steady state currents $`I_0`$ at finite $`\mathrm{\Omega }`$ reproduces the known result at $`\mathrm{\Omega }=0`$ exactly even for small values of $`\chi `$, with a steady state current $`I_{SS}`$ and no knee.
Our conclusions from these results are as follows:
(i) We know that up to $`t_{\mathrm{knee}}`$ our simulation results are exact, as they are unchanged in the linear region with current $`I_{SS}`$ for $`\chi =20300`$. As this regime lasts at least until $`tJ10`$, these results would be observable in an experimental implementation of the SAT.
(ii) As an impractically large value of $`\chi `$ would be required to reproduce the results exactly on long timescales, we can not be certain what the final behaviour will be for $`t>t_{\mathrm{knee}}(\chi =300)`$. This depends on clearly interesting phenomena that arise from strong correlations between the left and right sides of the impurity site, and could include settling to a final steady state current $`I_0`$. These effects would also be observable in an experiment.
The expected final steady state values $`I_0`$ are already discussed in Ref. sat , and so in the remainder of this article we investigate the initial steady state currents $`I_{SS}`$ in various parameter regimes.
### III.2 Diffusive evolution, with initial mean momentum $`(\widehat{k}_{t=0}=0)`$
We first consider the motion of atoms past the impurity site in the diffusive regime, where the initial state at $`t=0`$ is the ground state of a Bose-Hubbard model on $`M=30`$ lattice sites in a box trap.
#### III.2.1 Dependence of the current on impurity-probe coupling, $`\mathrm{\Omega }`$
In Fig. 7 we show the initial steady state current $`I_{SS}`$ as a function of $`\mathrm{\Omega }/J`$ for fermionic probe atoms, and for bosonic probe atoms with $`U_{bb}/J=4,10,\mathrm{}`$ and $`\mathrm{\Delta }=0`$. All of these results decrease as expected with increasing $`\mathrm{\Omega }/J`$, and even for a relatively small $`\mathrm{\Omega }=2J`$ the current is minimal in each case. At half filling (Fig. 7a), the results for the Tonks gas are identical to the Fermi results for $`\mathrm{\Omega }=0`$, but become substantially different as $`\mathrm{\Omega }`$ increases, with the currents in this regime greater for the bosons. At weaker interactions the currents are smaller than the Tonks result at all $`\mathrm{\Omega }`$, but for $`\mathrm{\Omega }/J>1`$ the currents for $`U/J=4`$ are larger than for a non-interacting Fermi gas. The variation in the currents for different interaction strengths of bosons appears to be due to the broader initial momentum distributions that occur at larger $`U/J`$. At unit filling (Fig. 7b), $`I_{SS}`$ is less dependent on the interaction strength, with all of the bosonic results very close to one another, currents becoming larger than that for fermions when $`\mathrm{\Omega }/J>1`$.
#### III.2.2 Dependence of the current on interaction strength, $`U/J`$
The dependence of the initial steady state current $`I_{SS}`$ on the interaction strength for bosons is depicted more clearly in Fig. 8, both at unit filling, $`n=1`$, and half filling, $`n=1/2`$ for $`\mathrm{\Omega }=J`$. At half filling the current increases with increasing $`U/J`$, which is due to the broader initial momentum distribution produced by the higher interaction energies. In contrast, at higher densities (here $`n=1`$), the probe atoms are blocked better by the SAT for higher interaction strengths, and $`I_{SS}`$ decreases. The key principle here is that bosons appear to be better blocked when they approach the impurity individually. For high densities this is achieved when large interaction strengths eliminate the higher occupancies of all lattice site including the impurity site. For weaker interactions the bosons can swamp the transistor, with one atom being bound to the impurity, whilst other probe atoms tunnel onto and past the impurity site.
This effect is seen in Fig. 9, where the molecular occupation and average probe atom occupation on the impurity site are shown for (a) $`U/J=4`$ and (b) $`U/J=10`$. We see that as $`n`$ increases, the molecular occupation becomes rapidly higher for $`U/J=10`$ than for $`U/J=4`$, despite the larger occupation of probe atoms on the impurity site for $`U/J=4`$. This indicates that for $`U/J=10`$ atoms arrive individually at the impurity site, where they are coupled with the impurity atom into a molecular state, and their transport is efficiently blocked. For $`U/J=4`$, more than one atom enters the impurity site at once, leading to a larger average probe atom occupation on the impurity site, but a comparatively small molecular occupation.
It is important to note, however, that even when $`U/J=4`$, the resulting currents are only slightly larger than they are for non-interacting fermions. At higher interaction strengths we then see an even stronger suppression of the steady state current for dense, strongly interacting bosons. As $`\mathrm{\Omega }`$ increases, both the molecular occupation and probe atom occupation on the impurity site decrease (Fig. 9) as the probability of even a single atom tunnelling onto the impurity site becomes small. For $`\mathrm{\Omega }>2J`$ the blocking mechanism of the SAT functions extremely well even in the regime where the probe atoms are dense and weakly interacting.
#### III.2.3 Dependence of the current on initial density, $`n`$
In Fig. 10 we show the dependence of the initial steady state current, $`I_{SS}`$ on the initial filling factor $`n`$ with (a) $`\mathrm{\Omega }/J=0.5`$ and (b) $`\mathrm{\Omega }/J=1`$. In both cases, the currents for bosons of different interaction strengths are very similar, with the variations following the patterns discussed in the preceding section. These results also agree well with the results for fermions at small $`n`$ and for $`n1`$, but the plateau observed in fermionic currents near $`n0.5`$ does not occur in the currents for bosons. For fermions, this plateau arises from the transmission profile of the SAT as a function of incoming momentum sat , and occurs when the Fermi momentum is raised past the minimum in this transmission profile. For interacting bosons, this correspondence between the momentum distribution of the gas and the transmission profile is destroyed by many-body effects, and we see instead a smooth increase in the current. This results in the bosonic currents being substantially larger than those for fermions near half filling when $`\mathrm{\Omega }1`$ (as was previously observed in Fig. 7a).
### III.3 Kicked evolution, with initial mean momentum $`(\widehat{k}_{t=0}0)`$
In this section we consider an initial state with a non-zero initial momentum, which is obtained, e.g., by briefly tilting the lattice on a timescale much shorter than that corresponding to dynamics of atoms in the lattice. If the tilt is linear, the resulting state will be given by
$$|\varphi (t=0)=\underset{j}{}\mathrm{exp}(\mathrm{i}p_kj\widehat{b}_j^{}\widehat{b}_j)|\varphi _0,$$
(9)
where $`|\varphi _0`$ is the initial many-body ground state, and the quantity $`p_k`$ is determined by the magnitude and duration of the tilt. The effect of this tilt is to translate the ground state in the periodic quasimomentum space by a momentum $`p_k`$. The final mean momentum $`k`$ then depends both on the value $`p_k`$ and the properties of the initial momentum distribution.
#### III.3.1 Dependence of the current on kick strength $`p_k`$
In the case of fermions, the dependence of the current on $`q`$ for different filling factors $`n=N/M`$ and $`\mathrm{\Omega }`$ can be clearly understood in terms of the SAT transmission profile (see satandi ). In Fig. 11a we see the current $`I_{ss}`$ as a function of $`p_k`$ with $`\mathrm{\Omega }=0`$. The currents are each peaked at $`p_k=\pi /2`$, where the resulting mean velocity of the probe atoms is the largest. For $`N/M=1`$, the whole Bloch band is filled, and the momentum distribution is not changed by the application of the kick, i.e., $`\widehat{k}_{t=0}=0`$. In Fig. 11b the same results are shown, but with $`\mathrm{\Omega }/J=1`$. Here we see that for small filling factors, a minimum appears at $`p_k=\pi /2`$, corresponding to the minimum in the transmission profile of the SAT for this incident momentum sat ; satandi . At higher filling factors, this feature of the transmission profile for $`\mathrm{\Omega }/J=1`$ is not sufficiently broad to overcome the increase current due to higer mean velocities in the initial cloud, and the peak at $`p_k=\pi /2`$ reappears. The currents here are, of course, reduced in comparison with those for $`\mathrm{\Omega }=0`$.
Whilst for all $`p_k`$ the currents with no coupling to the impurity atom, i.e., $`\mathrm{\Omega }=0`$, are the same for the Tonks gas as for fermions (Fig. 11a), the currents for finite interaction strengths are found to be remarkably different. In Fig. 12 these rates are plotted for $`U/J=4,7,10`$ for $`N=5,15,30`$ particles initially situated on $`M=30`$ sites. For the very dilute system with $`N=5`$ (Fig. 12a) we see a peak similar to that observed for fermions which is independent of the interaction strength. Here the currents are essentially those for non-interacting particles, and the currents determined by the initial momentum distribution. For $`N=15`$ (Fig. 12b) we observe the surprising result that the current is peaked at a lower value than is observed for fermions or for the Tonks gas. We have observed this peak consistently for such cases of finite interaction, and note that as $`U/J`$ increases, the peak moves back towards $`p_k=\pi /2`$ as the currents converge to the Tonks gas results. As $`N`$ is further increased, the peak continues to move left, and for $`N=30`$ (Fig. 12c) we see a monotonically decreasing current as $`p_k`$ increases. As $`U/J`$ increases these values tend towards the $`p_k`$ independent result observed for the Tonks gas. These results are surprising, but the trends in the behaviour are clear, and they should be directly verifiable in experiments, even without the presence of the impurity atom.
For non-zero coupling to the impurity atom, the currents as a function of $`p_k`$ are shown in Fig. 13. Again we notice that the current for bosons with finite interaction strength is peaked at much lower values of $`p_k`$ than the fermionic currents and that peaks of all of the bosonic currents, including the Tonks currents, as significantly larger than the fermionic currents at half filling, as was observed for diffusive results ($`p_k=0`$). The most remarkable feature of these plots is that despite a significant reduction in the current, the basic dependence on $`p_k`$ is very similar to the $`\mathrm{\Omega }=0`$ results.
#### III.3.2 Dependence of the current on impurity-probe coupling, $`\mathrm{\Omega }`$
The steady state current $`I_{ss}`$ is shown in Fig. 14 as a function of $`\mathrm{\Omega }`$. We observe the same strong decrease in the current due to the operation of the SAT for all of these curves, with the highest currents corresponding to the $`p_k=\pi /2`$ curve as expected. Note that the value of $`I_{SS}`$ is affected equally for all of the kick strengths, and the ratio in the currents for different values of $`p_k`$ is very similar for $`\mathrm{\Omega }=0`$ and $`\mathrm{\Omega }/J=2`$.
## IV Summary
In summary, the SAT setup provides new experimental opportunities to study coherent transport of many atoms past a spin-1/2 impurity due to the relatively long coherence times that exist for systems of atoms in optical lattices. The resulting coherent many-body effects can be clearly seen in the difference between the atomic currents observed for fermions and bosons, and the non-trivial dependence of the current on interaction strength for bosons with finite interactions. Even stronger dependence on these interactions is observed when the probe atoms are initially accelerated to a non-zero momentum. The initial steady state currents would be directly accessible quantities in the experimental implementation of the SAT, and using recently developed methods for time-dependent calculation of many-body 1D systems, we have made quantitative predictions for the corresponding currents for a wide range of system parameters. We cannot be certain about the values the currents approach at long times, although it is possible that the system will settle eventually into a regime with a different steady state current. The high values of $`\chi `$ needed to reproduce this behaviour in our numerical calculations suggest that the currents in this regime could also be strongly sensitive to the coherence properties of the system, which would be very interesting to investigate in an experiment.
###### Acknowledgements.
The authors would like to thank A. Micheli for advice and stimulating discussions, G. Vidal for helpful discussions on the numerical methods, and A. Kantian for his contributions to producing the number conserving version of the program code. AJD thanks the Clarendon Laboratory and DJ thanks the Institute for Quantum Optics and Quantum Information of the Austrian Academy of Sciences for hospitality during the development of this work. This work was supported by EU Networks and OLAQUI. In addition, work in Innsbruck is supported by the Austrian Science Foundation and the Institute for Quantum Information, and work in Oxford is supported by EPSRC through the QIP IRC (www.qipirc.org) (GR/S82176/01) and the project EP/C51933/1.
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# An HLLC Riemann Solver for Relativistic Flows: I. Hydrodynamics
## 1 Introduction
High energy astrophysical phenomena involve, in many cases, relativistic flows, typical examples are superluminal motion of relativistic jets in extragalactic radio sources, accretion flows around massive compact objects, pulsar winds and Gamma Ray Bursts. The modeling of such phenomena has prompted the search for efficient and accurate numerical formulations of the special relativistic fluid equations (for an excellent review see Martรญ & Mรผller, 2003). There is now a strong consensus that the so-called โhigh-resolution shock-capturingโ schemes provide the necessary tools in developing stable and robust relativistic fluid dynamical codes. One of the fundamental ingredients of such schemes is the exact or approximate solution to the Riemann problem.
The solution to the Riemann problem in relativistic hydrodynamics (RHD henceforth) has been extensively studied in literature, and an exact solution can be found within high degree of accuracy by iterative techniques, see Martรญ & Mรผller (1994), Pons et al. (2000), Rezzolla et al. (2003) and references therein. One of the major differences with the classical counterpart is the velocity coupling introduced by the Lorentz factor and the coupling of the latter with the specific enthalpy. This considerably adds to the computational cost, making the use of an exact solver code prohibitive in a multidimensional Godunov-type code.
From this perspective, approximate solvers based on alternative strategies have been devised: local linearization (Eulderink & Mellema, 1995; Falle & Komissarov, 1996), two-shock approximation (Balsara, 1994; Dai & Woodward, 1997; Mignone et al., 2005), flux-splitting methods (Donat et al., 1998), and so forth, see Martรญ & Mรผller (2003) for a comprehensive review. Most of these solvers, however, rely on rather expensive characteristic decompositions of the Jacobian matrix or involve iterative techniques to solve highly nonlinear equations. Although they usually attain better resolution at discontinuities, some of these methods may produce unphysical states with negative densities or pressures, as it has been shown by Einfeldt et al. (1991) for linearized Riemann solvers in the context of classical hydrodynamics.
The HLL method devised by Harten et al. (1983) for classical gasdynamics belongs to a different class of approximate Riemann solvers and has gained increasing popularity among researchers in the last decade. It has been implemented in the context of the relativistic fluid equations by Schneider et al. (1993) and Duncan & Hughes (1994). The HLL approach does not require a full characteristic decomposition of the equations and is straightforward to implement in a Godunov-type code. Besides the computational efficiency, this class of solvers has the attractive feature of being positively conservative in the sense that preserve initially positive densities, energy and pressures.
The HLL formulation, however, lacks the ability to resolve isolated contact discontinuity and for this reason has a more diffusive character than other more sophisticated algorithms. To compensate for this, Toro et al. (1994) developed an extension of the HLL solver for the Euler equations introducing a two-state HLL-type solver called HLLC (where โCโ stands for contact) that improves the treatment of the contact discontinuity, see also Batten et al. (1997). Recently this approach has been generalized to the magnetohydrodynamic equations (Gurski, 2004; Li, 2005).
In the present work, we extend this approach to the relativistic equation of fluid dynamics. The paper is structured as follows: in ยง2 the relevant equations are given, in ยง3 we describe the new approximate Riemann solver and in ยง4 we asses the strength of the new method with one and two dimensional tests.
## 2 The RHD Equations
The motion of an ideal relativistic fluid is governed by conservation of mass, momentum and energy. The pertaining equations are cast as a hyperbolic system of conservation laws (Landau & Lifshitz, 1959) which, in two dimensions, reads
$$\frac{๐ผ}{t}+\frac{๐ญ^x(๐ผ)}{x}+\frac{๐ญ^y(๐ผ)}{y}=0,$$
(1)
where $`๐ผ=(D,m_x,m_y,E)`$ is the unknown vector of conservative variables, whereas $`๐ญ^x`$ and $`๐ญ^y`$ are, respectively, the fluxes along the $`x`$ and $`y`$ directions:
$$๐ญ^x(๐ผ)=\left(\begin{array}{c}Dv_x\\ \\ m_xv_x+p\\ \\ m_yv_x\\ \\ m_x\end{array}\right),๐ญ^y(๐ผ)=\left(\begin{array}{c}Dv_y\\ \\ m_xv_y\\ \\ m_yv_y+p\\ \\ m_y\end{array}\right).$$
(2)
Generalization to three dimensions is straightforward.
In equations (2), $`p`$ is the thermal pressure, whereas $`D`$, $`๐(m_x,m_y)`$ and $`E`$ are, respectively, the mass, momentum and energy densities relative to the lab frame, where the fluid has velocity $`๐(v_x,v_y)`$. Units are conveniently normalized so that the speed of light is $`c=1`$.
The relation between conserved variables $`๐ผ`$ and physical quantities $`๐ฝ=(\rho ,v_x,v_y,p)`$ is
$$D=\gamma \rho ,๐=Dh\gamma ๐,E=Dh\gamma p,$$
(3)
where $`\rho `$ is the proper rest mass density, $`\gamma =(1๐๐)^{\frac{1}{2}}`$ is the Lorentz factor and $`h`$ is the specific enthalpy. Proper closure is provided by specifying an equation of state in the form $`h=h(p,\rho )`$.
For an ideal gas, the enthalpy has the form $`\rho h=\rho +p\mathrm{\Gamma }/(\mathrm{\Gamma }1)`$ and the sound speed is defined by
$$c_s=\sqrt{\frac{\mathrm{\Gamma }p}{\rho h}}.$$
(4)
with $`\mathrm{\Gamma }`$ being the (constant) specific heat ratio. By letting $`p/\rho \mathrm{}`$, we see that the square of the sound speed has the limiting value $`c_s^2\mathrm{\Gamma }1`$. Since it can be shown (Taub, 1948; Anile, 1989; Mignone et al., 2005) that the specific heat ratio $`\mathrm{\Gamma }`$ cannot exceed $`2`$, one always has $`c_s^2<1`$. This is an important result for the positivity of the HLL and HLLC schemes and will be used in a later section.
Equation (3) gives $`๐ผ`$ in terms of the primitive state vector $`๐ฝ`$. The inverse relation involves the solution of a nonlinear equation for the pressure $`p`$:
$$E+p=D\gamma +\frac{\mathrm{\Gamma }}{\mathrm{\Gamma }1}p\gamma ^2$$
(5)
where $`\gamma =\left[1|๐|^2/(E+p)^2\right]^{\frac{1}{2}}`$. Equation (5) can be solved by any standard root finding algorithm.
### 2.1 The Riemann Problem in RHD
Consider a conservative discretization of (1) along the x-direction:
$$\frac{\overline{๐ผ}_i^{n+1}\overline{๐ผ}_i^n}{\mathrm{\Delta }t^n}=\frac{๐_{i+\frac{1}{2}}๐_{i\frac{1}{2}}}{\mathrm{\Delta }x_i}.$$
(6)
The numerical flux functions $`๐_{i+\frac{1}{2}}`$ follow from the solution of Riemann problems with initial data:
$$๐ผ(x,0)=\{\begin{array}{ccc}๐ผ_{L,i+\frac{1}{2}}& \text{if}& x<x_{i+\frac{1}{2}},\\ \multicolumn{3}{c}{}\\ ๐ผ_{R,i+\frac{1}{2}}& \text{if}& x>x_{i+\frac{1}{2}},\end{array}$$
(7)
where $`๐ผ_{L,i+\frac{1}{2}}`$ and $`๐ผ_{R,i+\frac{1}{2}}`$ are the left and right edge values at zone interfaces.
The solution of the Riemann problem for the special relativistic fluid equations has been investigated by Martรญ & Mรผller (1994), Pons et al. (2000), Rezzolla et al. (2003). It consists of a self-similar three-wave pattern generated by the decay of the initial discontinuity (7). The resulting Riemann fan (Fig. 1) is bounded by two nonlinear waves (representing either shocks or rarefactions) separated by a contact discontinuity moving at the fluid velocity. Across the contact discontinuity, pressure and normal velocity are continuous whereas density and tangential velocities experience jumps. The same holds also in the non-relativistic limit. Contrary to the Newtonian counterpart, however, all variables are discontinuous across a shock wave or change smoothly through a rarefaction fan (Pons et al., 2000). This is a consequence of the velocity coupling introduced by the Lorentz factor $`\gamma `$ and by the coupling of the latter with the specific enthalpy $`h`$.
The resulting wave pattern can be solved to a high degree of precision by iterative techniques and has been implemented for the first time in the one-dimensional Godunov-type code by Martรญ & Mรผller (1996). Nevertheless, when tangential velocities are included, the computational effort increases considerably, and the use of an exact solver in a multidimensional Godunov-type code can become prohibitive.
Here we consider a different approach, based on the original prescription given by Harten et al. (1983) (HLL) for the classical Euler equations and subsequently extended by Toro et al. (1994) (HLLC). Both the HLL and HLLC formulations do not require a field by field decomposition of the relativistic equations, a feature which makes them particularly attractive, specially in multi-dimensional applications.
## 3 The HLL Framework
Harten, Lax and van Leer (Harten et al., 1983) proposed an approximate solution to the Riemann problem where the two states bounded by the two acoustic waves are averaged into a single constant state. In other words, the solution to the Riemann problem on the $`x/t=0`$ axis consists of the three possible constant states:
$$๐ผ(0,t)=\{\begin{array}{ccc}๐ผ_L& \text{if}\hfill & \lambda _L0,\hfill \\ \multicolumn{3}{c}{}\\ ๐ผ^{hll}& \text{if}\hfill & \lambda _L0\lambda _R,\hfill \\ \multicolumn{3}{c}{}\\ ๐ผ_R& \text{if}\hfill & \lambda _R0,\hfill \end{array}$$
(8)
where we dropped, for simplicity, the half integer notation $`i+\frac{1}{2}`$. Harten et al. (1983) noted that the single state $`๐ผ^{hll}`$ could be constructed from an a priori estimate of the fastest and slowest signal velocities $`\lambda _L`$ and $`\lambda _R`$:
$$๐ผ^{hll}=\frac{\lambda _R๐ผ_R\lambda _L๐ผ_L+๐ญ_L๐ญ_R}{\lambda _R\lambda _L},$$
(9)
where $`๐ญ_L=๐ญ^x(๐ผ_L)`$, $`๐ญ_R=๐ญ^x(๐ผ_R)`$. Notice that equation (9) represents the integral average of the solution of the Riemann problem over the wave fan (Toro, 1997).
The corresponding interface numerical flux is defined as:
$$๐=\{\begin{array}{ccc}๐ญ_L& \text{if}\hfill & \lambda _L0,\hfill \\ \multicolumn{3}{c}{}\\ ๐ญ^{hll}& \text{if}\hfill & \lambda _L0\lambda _R,\hfill \\ \multicolumn{3}{c}{}\\ ๐ญ_R& \text{if}\hfill & \lambda _R0,\hfill \end{array}$$
(10)
where
$$๐ญ^{hll}=\frac{\lambda _R๐ญ_L\lambda _L๐ญ_R+\lambda _R\lambda _L(๐ผ_R๐ผ_L)}{\lambda _R\lambda _L}.$$
(11)
Thus, given a wave speed estimate for the fastest and slowest speeds $`\lambda _R`$ and $`\lambda _L`$ (see ยง3.1.1), an approximate solution to the Riemann problem can be constructed and the intercell numerical fluxes for the conservative update (6) are computed using (10).
This approach has been applied for the first time to the one-dimensional relativistic equations by Schneider et al. (1993) and later by Duncan & Hughes (1994) for the multidimensional case.
Although the HLL prescription is computationally inexpensive and straightforward to implement, a major drawback is its inability to resolve contact or tangential waves. On the contrary the HLLC scheme, originally introduced by Toro et al. (1994) in the context of the Euler equations of classical gasdynamics, does not suffer from this loss. In the next section we generalize this approach to the equations of relativistic hydrodynamics.
### 3.1 HLLC Solver
The HLLC scheme restores the full wave structure inside the Riemann fan by replacing the single averaged state defined by (9) with two approximate states, $`๐ผ_L^{}`$ and $`๐ผ_R^{}`$. These two states are separated by a middle contact wave which is assumed to have constant speed $`\lambda ^{}`$, so that the full solution to the Riemann problem now reads
$$๐ผ(0,t)=\{\begin{array}{ccc}๐ผ_L& \text{if}& \lambda _L0,\\ \multicolumn{3}{c}{}\\ ๐ผ_L^{}& \text{if}& \lambda _L0\lambda ^{},\\ \multicolumn{3}{c}{}\\ ๐ผ_R^{}& \text{if}& \lambda ^{}0\lambda _R,\\ \multicolumn{3}{c}{}\\ ๐ผ_R& \text{if}& \lambda _R0,\end{array}$$
(12)
and the corresponding intercell numerical fluxes are:
$$๐=\{\begin{array}{ccc}๐ญ_L& \text{if}& \lambda _L0,\\ \multicolumn{3}{c}{}\\ ๐ญ_L^{}& \text{if}& \lambda _L0\lambda ^{},\\ \multicolumn{3}{c}{}\\ ๐ญ_R^{}& \text{if}& \lambda ^{}0\lambda _R,\\ \multicolumn{3}{c}{}\\ ๐ญ_R& \text{if}& \lambda _R0.\end{array}$$
(13)
The intermediate state fluxes $`๐ญ_L^{}`$ and $`๐ญ_R^{}`$ may be expressed in terms of $`๐ผ_L^{}`$ and $`๐ผ_R^{}`$ from the Rankine-Hugoniot jump conditions:
$$\lambda \left(๐ผ^{}๐ผ\right)=๐ญ^{}๐ญ,$$
(14)
where here and throughout the following, quantities without a suffix $`L`$ or $`R`$ refer indifferently to the left ($`L`$) or right ($`R`$) states. Notice that, in general, $`๐ญ=๐ญ^x(๐ผ)`$ but $`๐ญ^{}๐ญ^x(๐ผ^{})`$.
We remind the reader that the HLL and HLLC solvers differ in the representation of the intermediate states. In the case of supersonic flows ($`\lambda _L>0`$ or $`\lambda _R<0`$), in fact, the two solvers become equivalent. The same result also holds for an exact Riemann solver.
If $`\lambda _L`$ and $`\lambda _R`$ are given (see ยง3.1.1), equation (14) represent a system of $`2n`$ equations (where $`n`$ is the number of components of $`๐ผ`$) for the $`4n+1`$ unknowns $`๐ผ_L^{}`$, $`๐ผ_R^{}`$, $`๐ญ_L^{}`$, $`๐ญ_R^{}`$ and $`\lambda ^{}`$. Three additional constraints come from the requirements that both pressure and normal velocity be continuous across the contact wave (i.e. $`v_{x,R}^{}=v_{x,L}^{}`$, $`p_R^{}=p_L^{}`$) and that $`\lambda ^{}=v_{x,L}^{}=v_{x,R}^{}`$. This, however, yields a total of only $`2n+3`$ equations, still not sufficient to solve the system. In order to reduce the number of unknowns and have a well-posed problem, further assumptions have to be made on the form of the fluxes $`๐ญ^{}`$. Here we assume that the two-dimensional fluxes can be written as
$$๐ญ^{}=\left(\begin{array}{c}D^{}v_x^{}\\ \\ m_x^{}v_x^{}+p^{}\\ \\ m_y^{}v_x^{}\\ \\ m_x^{}\end{array}\right).$$
(15)
In such a way, both $`๐ผ^{}`$ and $`๐ญ^{}`$ are expressed in terms of the five unknowns $`D^{}`$, $`v_x^{}`$, $`m_y^{}`$, $`E^{}`$ and $`p^{}`$. The normal components of momentum in the star region, $`m_{x,L}^{}`$ and $`m_{x,R}^{}`$, are not independent variables since, for consistency, we require that $`m_x^{}=(E^{}+p^{})v_x^{}`$. In the classical case, this assumption becomes equivalent to $`m_x^{}=\rho ^{}\lambda ^{}`$. This yields a total of $`11`$ equations in $`11`$ unknowns.
Writing explicitly equation (14) for the left or the right state yields
$$\begin{array}{ccc}D^{}(\lambda \lambda ^{})\hfill & =& D(\lambda v_x),\hfill \\ \multicolumn{3}{c}{}\\ m_x^{}(\lambda \lambda ^{})\hfill & =& m_x(\lambda v_x)+p^{}p,\hfill \\ \multicolumn{3}{c}{}\\ m_y^{}(\lambda \lambda ^{})\hfill & =& m_y(\lambda v_x),\hfill \\ \multicolumn{3}{c}{}\\ E^{}(\lambda \lambda ^{})\hfill & =& E(\lambda v_x)+p^{}\lambda ^{}pv_x.\hfill \end{array}$$
(16)
If one combines the last of (16) together with the second one, the following expression giving $`\lambda ^{}`$ in terms of $`p^{}`$ may be obtained:
$$(A\lambda p^{})v_x^{}=B+p^{},$$
(17)
where $`A=\lambda Em_x`$, $`B=m_x(\lambda v_x)p`$.
By imposing $`p_{x,L}^{}=p_{x,R}^{}`$ across the contact discontinuity we find the following quadratic equation for $`\lambda ^{}`$:
$$F_E^{hll}\left(\lambda ^{}\right)^2(E^{hll}+F_{m_x}^{hll})\lambda ^{}+m_x^{hll}=0.$$
(18)
In equation (18), $`F_E^{hll}`$ and $`F_{m_x}^{hll}`$ are the energy and momentum components of the HLL flux given by equation (11), whereas $`E^{hll}`$ and $`m_x^{hll}`$ are the energy and normal momentum components of the HLL state vector, equation (9). Of the two roots of equation (18) only the one with the minus sign is physically acceptable, since it lies in the range $`(1,1)`$ and, according to the wave-speed estimate presented in ยง3.1.1, can be interpreted as an average velocity between $`\lambda _L`$ and $`\lambda _R`$. The mathematical proof of this statement is given in the appendix (A). Besides, the same root has the correct classical limit, that is $`\lambda ^{}m_x^{hll}/\rho ^{hll}`$ as $`v/c0`$ and $`h1`$. This wave speed is the same one proposed by Toro (1997) and further discussed in Batten et al. (1997).
Once $`\lambda ^{}`$ is known, $`p^{}`$ is computed from (17) and the components of $`๐ผ^{}`$ are readily obtained from (16).
Finally we notice that the method is consistent, in that the integral average over the Riemann fan automatically satisfies the consistency condition by construction (Toro, 1997):
$$\frac{(\lambda ^{}\lambda _L)๐ผ_L^{}+(\lambda _R\lambda ^{})๐ผ_R^{}}{\lambda _R\lambda _L}=๐ผ^{hll},$$
(19)
or, alternatively,
$$\frac{๐ญ_L^{}\lambda _R(\lambda ^{}\lambda _L)+๐ญ_R^{}\lambda _L(\lambda _R\lambda ^{})}{\lambda _R\lambda _L}=\lambda ^{}๐ญ^{hll}.$$
(20)
Incidentally, we notice that equation (18) could have been obtained by algebraic manipulations of equations (19) and (20).
#### 3.1.1 Wave Speed Estimate
The wave speeds needed in our formulation are estimates for the lower and upper bounds of the signal velocities in the solution to the Riemann problem (Toro, 1997). Here we consider the relativistic generalization of the estimates given by Davis (1988) for the Euler equation of gasdynamics. The same choice has been initially adopted by Schneider et al. (1993), Duncan & Hughes (1994) in their relativistic HLL solver and is commonly used by other authors, see, for example, Del Zanna & Bucciantini (2002). Specifically we set:
$$\begin{array}{c}\lambda _L=\mathrm{min}(\lambda _{}(๐ฝ_R),\lambda _{}(๐ฝ_L)),\\ \\ \lambda _R=\mathrm{max}(\lambda _+(๐ฝ_R),\lambda _+(๐ฝ_L)),\end{array}$$
(21)
where $`\lambda _+`$ and $`\lambda _{}`$ are the maximum and minimum eigenvalues of the Jacobian matrix $`๐ญ/๐ผ`$. They are the roots of the quadratic equation
$$(\lambda v_x)^2=\sigma _s(1\lambda ^2),$$
(22)
with $`\sigma _s=c_s^2/\left(\gamma ^2(1c_s^2)\right)`$, and hence
$$\lambda _\pm (๐ฝ)=\frac{v_x\pm \sqrt{\sigma _s\left(1v_x^2+\sigma _s\right)}}{1+\sigma _s}.$$
(23)
It should be mentioned that the wave speed estimate (21) is not the only possible one and different choices (such as the Roe average) may be considered.
#### 3.1.2 Positivity of the HLLC scheme
Adopting the same notations as in Batten et al. (1997), we denote with $`G`$ the set of physically admissible conservative states:
$$G=\{\left(\begin{array}{c}D\\ \\ m_x\\ \\ m_y\\ \\ E\end{array}\right),\begin{array}{c}D>0\hfill \\ \\ E>\sqrt{m_x^2+m_y^2+D^2}\hfill \end{array}\},$$
(24)
where the second inequality simultaneously guarantees pressure positivity and that the total velocity never exceeds the speed of light.
We remind the reader that the pressure $`p(๐ผ^{})`$ computed from the conservative state $`๐ผ^{}`$ using (5) should not be confused with $`p^{}`$ appearing in the flux definition (15). The two pressures are, in fact, different and the positivity argument should apply to $`p(๐ผ^{})`$ rather than $`p^{}`$, which can take negative values under certain circumstances. Similar considerations hold for the velocity $`\lambda ^{}`$ of the contact discontinuity for which, in general, we have $`\lambda ^{}v_x(๐ผ^{})`$. Thus $`p^{}`$ and $`\lambda ^{}`$ may be more conveniently considered as auxiliary variables.
This is one of the fundamental differences between our relativistic solver and the classical HLLC scheme, for which $`\lambda ^{}=m_x^{hll}/\rho ^{hll}`$ and thus only $`p^{}`$ plays the role of an auxiliary parameter. This behavior is a direct consequence of the relativistic coupling between thermodynamical and kinetical terms, a feature absent in the Newtonian formulation.
The positivity of the HLLC scheme is preserved if each of the two intermediate states $`๐ผ_L^{}`$ and $`๐ผ_R^{}`$ are contained in G.
For the density, the proof is trivial and follows from the inequalities $`\lambda _L\lambda ^{}\lambda _R`$ and $`\lambda _Lv_x(L,R)\lambda _R`$, see Appendix ยงA.
Unfortunately the analytical proof of the second statement presents some algebraic difficulties, since the second of (24) reduces to an inequality for a quartic equation in $`\lambda ^{}`$. However, extensive numerical testing, part of which is presented in ยง4, has shown that the second of (24) is always satisfied for all pair of states $`๐ผ_L`$ and $`๐ผ_R`$ whose wave speeds are computed according to (21) and for which an exact analytical solution to the Riemann problem exists (i.e. no vacuum is created).
## 4 Algorithm Validation
We now provide some numerical examples to test our new HLLC solver. For the test problems considered in this section we closely follow Lucas-Serrano et al. (2004).
### 4.1 Implementation Details
The numerical integration of the relativistic equations (1) proceeds via the conservative update (6). For the first-order HLLC scheme, we compute the inter-cell numerical fluxes $`๐_{i+\frac{1}{2}}`$ using (10) with left and right states given, respectively, by $`๐ผ_i`$ and $`๐ผ_{i+1}`$.
For the second order scheme, the input to the Riemann problem are the states
$$\begin{array}{c}๐ฝ_{i+\frac{1}{2},L}^{n+\frac{1}{2}}=๐ฝ_i^{n+\frac{1}{2}}+\frac{\delta ๐ฝ_i^n}{2},\\ \\ ๐ฝ_{i+\frac{1}{2},R}^{n+\frac{1}{2}}=๐ฝ_{i+1}^{n+\frac{1}{2}}\frac{\delta ๐ฝ_{i+1}^n}{2},\end{array}$$
(25)
where $`๐ฝ_i^{n+\frac{1}{2}}`$ follows from a simple Hancock predictor step,
$$๐ผ_i^{n+\frac{1}{2}}=๐ผ_i^n\frac{\mathrm{\Delta }t^n}{2\mathrm{\Delta }x_i}\left[๐ญ\left(๐ฝ_{i+\frac{1}{2},L}^n\right)๐ญ\left(๐ฝ_{i\frac{1}{2},R}^n\right)\right],$$
(26)
with $`๐ฝ_{i+\frac{1}{2},L}^n`$ and $`๐ฝ_{i\frac{1}{2},R}^n`$ computed from (25) by replacing $`๐ฝ^{n+\frac{1}{2}}`$ with $`๐ฝ^n`$.
The $`\delta ๐ฝ`$โs appearing in equation (25) are computed at the beginning of the time step using the fourth-order limited slopes (Colella, 1985; Saltzman, 1994):
$$\delta ๐ฝ_i=s_i\mathrm{min}(\left|\frac{4}{3}\mathrm{\Delta }_0๐ฝ_i\frac{\overline{\delta }๐ฝ_{i+1}+\overline{\delta }๐ฝ_{i1}}{6}\right|,\mathrm{\Delta }_l๐ฝ_i),$$
(27)
where
$$\mathrm{\Delta }_l๐ฝ_i=\alpha \mathrm{min}(|\mathrm{\Delta }๐ฝ_i|,|\mathrm{\Delta }๐ฝ_{i1}|),$$
(28)
and $`\overline{\delta }๐ฝ_i`$ are the second-order slopes
$$\overline{\delta }๐ฝ_i=s_i\mathrm{min}(\mathrm{\Delta }_l๐ฝ_i,|\mathrm{\Delta }_0๐ฝ_i|),$$
(29)
$$\mathrm{\Delta }๐ฝ_i=๐ฝ_{i+1}๐ฝ_i,\mathrm{\Delta }_0๐ฝ_i=\frac{๐ฝ_{i+1}๐ฝ_{i1}}{2},$$
(30)
$$s_i=\frac{\text{sign}(\mathrm{\Delta }๐ฝ_i)+\text{sign}(\mathrm{\Delta }๐ฝ_{i1})}{2}.$$
(31)
The parameter $`\alpha [1,2]`$ adjusts the limiter compression, with $`\alpha =2`$ ($`\alpha =1`$) yielding a more (less) compressive limiter. Notice that, although the use of fourth-order slopes attains sharper representations of discontinuities, the scheme retains global second-order spatial accuracy.
We do not make use of any artificial steepening algorithm to enhance resolution across a contact wave (Martรญ & Mรผller, 1996; Lucas-Serrano et al., 2004; Mignone et al., 2005) in order to highlight the intrinsic capabilities of our new HLLC solver. In the one-dimensional tests, the computational domain is the interval $`[0,1]`$ and the compression parameter is $`\alpha =2`$. In two dimensions we set $`\alpha =2`$, $`\alpha =1.25`$ and $`\alpha =1`$ for density, velocities and pressure, respectively. Additional shock flattening, computed as in (Martรญ & Mรผller, 1996), is used in ยง4.2 and ยง4.6 to prevent spurious numerical oscillations. Outflow boundary conditions are set in problem $`1`$$`4`$.
Multidimensional integration is achieved via Strang directional splitting (Strang, 1968), that is, by successively applying one-dimensional operators in reverse order from one time step to the next one, i.e. $`๐ผ^{n+2}=_x_y๐ผ^{n+1}`$ and $`๐ผ^{n+1}=_y_x๐ผ^n`$. Here $`_x`$ is the operator corresponding to the conservative update (6) (and similarly for $`_y`$). The same time increment $`\mathrm{\Delta }t`$ should be used for two consecutive time steps.
### 4.2 Problem 1
The first test consists in a Riemann problem with initial data
$$(\rho ,v_x,p)=\{\begin{array}{cc}(1,0.9,1)& \text{for}x<0.5,\\ \multicolumn{2}{c}{}\\ (1,0,10)& \text{for}x>0.5.\end{array}$$
(32)
Integration is carried with $`CFL=0.8`$ until $`t=0.4`$ and an ideal equation of state with $`\mathrm{\Gamma }=4/3`$ is used. The breakup of the discontinuity results in the formation of two shock waves separated by a contact discontinuity.
In Fig. 2 we plot the analytical solution for the rest mass density together with the profiles obtained with the first-order HLLC and HLL schemes on $`100`$ uniform computational zones. The two integrations behave similarly near the shock waves, but differ in the ability to resolve the contact discontinuity. As expected, the HLLC scheme yields a sharper representation of the latter, whereas the HLL solver retains a more diffusive character.
The $`L_1`$ norm errors of density are shown in the top-left panel of Fig. 8 for different resolutions. For the sake of comparison, computations have also been performed using the more sophisticated exact Riemann solver described in the one-dimensional code by Martรญ & Mรผller (1996). The errors obtained with the present HLLC scheme and the exact Riemann solver are comparable at low resolution ($`15.3\%`$ and $`13.6\%`$ respectively on $`100`$ points) and become nearly identical as the number of points increases. Conversely, the errors computed with the relativistic HLL scheme are bigger ($`22.2\%`$ on $`100`$ points) and show that almost twice the resolution is needed to achieve the same accuracy obtained with the HLLC or the exact solver.
Fig. 3 shows the results computed with the second-order HLLC scheme on $`400`$ zones, at the same time. The exact profiles for density, velocity and pressure are plotted as solid lines. Additional slope flattening (Martรญ & Mรผller, 1996) has been used to reduce the spurious numerical oscillations observed behind the shock front. All discontinuities are adequately captured and resolved on few computational cells, $`3`$ for the shocks and $`45`$ for the contact discontinuity (contrary to $`7`$ when the HLL solver is employed).
The error in $`L_1`$ norm is $`2.3\%`$ for $`400`$ grid zones and it has been computed at different resolutions using the HLL, HLLC and exact Riemann solver, see Fig. 9. Not surprisingly, the second-order interpolation considerably reduces the errors and higher convergence rates are expected for all schemes. Nevertheless, the three solvers mostly differ in the resolution at the contact discontinuity and, for $`n800`$ grid points, the HLLC and exact Riemann are practically indistinguishable, while at the maximum resolution employed ($`3200`$ zones), the error computed with the HLL scheme is still $`20\%`$ bigger.
### 4.3 Problem 2
In the second test, we prescribe the initial condition
$$(\rho ,v_x,p)=\{\begin{array}{cc}(1,0.6,10)& \text{for}x<0.5,\\ \multicolumn{2}{c}{}\\ (10,0.5,20)& \text{for}x>0.5,\end{array}$$
(33)
with an ideal equation of state with $`\mathrm{\Gamma }=5/3`$. Integration stops at $`t=0.4`$ and CFL = 0.8 has been used in the integration. The initial discontinuity evolves into left-going and right-going rarefaction waves with a contact discontinuity in the middle.
Results for the first-order HLL and HLLC schemes on a $`100`$-point uniform grid are shown in Fig. 4. Again, notice the sharper resolution of the HLLC scheme in proximity of the contact wave. The smooth rarefaction waves are equally resolved by both schemes.
The behavior of the solution under grid resolution effects is described in the top-right panel of Fig. 8. Since the only discontinuity in the problem is the contact wave, the $`L_1`$ norm reflects mostly the different resolution across the discontinuity. The HLLC and the exact solver perform nearly identically, while the HLL exhibits a slightly slower convergence rate. At the maximum resolution, the error in the HLL scheme is $`4.3\%`$ to be compared to the $`3.0\%`$ and $`3.1\%`$ errors obtained from the other two Riemann solvers.
These differences are again reduced in the second-order HLLC scheme, Fig. 5, for which the convergence rates are similar, as shown in the top-right panel of Fig. 9.
### 4.4 Problem 3
The initial condition for this test is
$$(\rho ,v_x,p)=\{\begin{array}{cc}(10,0,40/3)& \text{for}x<0.5,\\ \multicolumn{2}{c}{}\\ (1,0,0)& \text{for}x>0.5,\end{array}$$
(34)
with $`\mathrm{\Gamma }=5/3`$. For numerical reasons, the pressure in the left state has been set equal to a small value, $`p=2/310^6`$. Integration is carried with CFL = 0.8 on $`400`$ grid points; the final integration time is $`t=0.4`$. The initial configuration results in a mildly relativistic blast wave, with a maximum Lorentz factor $`\gamma _{\mathrm{max}}1.4`$. The Riemann fan consists of a rarefaction wave moving to the left, a shock wave adjacent to a contact discontinuity, both moving to the right, see Fig. 6.
Our HLLC scheme is able to capture discontinuities properly; in particular, the shock is resolved within 2-3 zones and the contact discontinuity smears out over 4-5 zones. We remind again that the interpolation algorithm does not make use of additional artificial compression to enhance resolution across the contact wave, as in Martรญ & Mรผller (1996). Moreover, we repeated the test also with the exact Riemann solver and did not find any noticeable difference. In addition and contrary to the previous two test problems, we did not find strong differences between our HLLC method and the HLL scheme. Resolution effects are given in the bottom left panels of Fig. 8 and 9 for the first-order and second order schemes, respectively. As one can see, the solutions computed with the HLL, HLLC and exact solvers behave nearly in the same way, with the $`L_1`$ norm errors being different by less than $`1\%`$ at low resolution and becoming identical for $`n800`$ grid points. We believe that this might be due to the proximity of the contact and shock waves. The quality of our results is, however, similar and comparable to those obtained in previous studies.
### 4.5 Problem 4
In the fourth shock-tube, we prescribe the following initial discontinuity
$$(\rho ,v_x,p)=\{\begin{array}{cc}(1,0,10^3)& \text{for}x<0.5,\\ \multicolumn{2}{c}{}\\ (1,0,10^2)& \text{for}x>0.5.\end{array}$$
(35)
Again we adopt an ideal equation of state with $`\mathrm{\Gamma }=5/3`$. The resulting pattern is similar to that of problem $`3`$, but the specific enthalpy in the left state is much greater than unity, thus resulting in a more thermodynamically relativistic configuration. The solution computed with the second-order scheme at $`t=0.4`$ is shown in Fig. 7 on $`400`$ computational zones and CFL=0.8.
The high pressure jump produces a strong shock wave and a contact discontinuity very close to each other, moving to the right at almost the same speeds. The higher compression in the shell is due to a more pronounced relativistic length-contraction effect caused by a higher Lorentz factor, $`\gamma _{\mathrm{max}}3.7`$. The smaller thickness of the shell between the shock and the contact wave makes this test numerically challenging and particularly demanding in terms of resolution.
Our relativistic HLLC scheme is able to reproduce the solution within a satisfactory agreement, even without using a contact steepening algorithm. The absolute global error in density is $`6.5\%`$ and the density peak in the thin shell achieves $`81.6\%`$ of the exact value. Our results are therefore similar to previous ones proposed in literature.
It should also be pointed out that, similarly to problem 3, we did not find any improvement in the solution by switching to the exact Riemann solver or using the relativistic HLL scheme. This is confirmed by the resolution study carried for the first and second-order schemes (bottom right panels in Fig. 8 and 9). Again, we suggest that the ability to capture the discontinuities relies on the interpolation properties of the algorithm and has a weaker dependence on the Riemann solver for this particular class of problems.
### 4.6 Relativistic Planar Shock Reflection
The initial configuration for this test problem consists in a cold ($`p=0`$), uniform ($`\rho =1`$) flow impinging on a wall located at $`x=0`$. The flow has constant inflow velocity $`v_{in}`$ and the reflection results in the formation of a strong shock wave. For $`t>0`$ the shock propagates upstream and the solution has an analytic form given by (Martรญ & Mรผller, 1996):
$$\rho (r,t)=\{\begin{array}{cc}1& \text{for}r>v_st,\\ \multicolumn{2}{c}{}\\ \sigma & \text{for}r<v_st,\end{array}$$
(36)
where
$$\sigma =\frac{\mathrm{\Gamma }+1+\mathrm{\Gamma }\left(\gamma _{in}1\right)}{\mathrm{\Gamma }1},v_s=(\mathrm{\Gamma }1)\frac{\gamma _{in}|v_{in}|}{\gamma _{in}+1},$$
(37)
are the compression ratio and the shock velocity, respectively. Behind the shock wave ($`r<v_st`$), the gas is at rest (i.e. $`v=0`$) and the pressure has the constant value $`\rho (r,t)(\gamma _{in}1)(\mathrm{\Gamma }1)`$. Conversely, in front of the shock all of the energy is kinetic and thus $`p=0`$, $`v=v_{in}`$.
For numerical reasons, pressure has been initialized to a small finite value, $`p=ฯต(\mathrm{\Gamma }1)`$, with $`ฯต=10^{10}`$ and $`\mathrm{\Gamma }=4/3`$. The computational domain is covered with $`100`$ computational zones and the initial inflow velocity is $`v_{in}=0.99999`$ corresponding to a Lorentz factor $`\gamma _{in}224`$. Integration is carried with CFL =0.4.
Fig. 10 shows the solution at $`t=1.5`$, after the shock has propagated $`\mathrm{\Delta }x_s0.5`$ from the wall. The relative global errors (defined as $`ฯต(L_1)/_i\rho _{ex}(x_i)\mathrm{\Delta }x_i`$) for density, velocity and pressure are, respectively, $`1.8\%`$, $`1.4\%`$ and $`1.4\%`$. This result is in excellent quantitative agreement with the numerical solutions obtained by other authors (Marquina et al., 1992; Martรญ et al., 1997; Aloy et al., 1999; Del Zanna & Bucciantini, 2002). In this test, similarly to problem $`1`$, shock flattening was employed to prevent numerical oscillations.
Also, from the same figure, we notice that our solver suffers from the wall-heating phenomenon, a common pathology in modern Godunov-type schemes. The degree of this pathology is higher than the HLL scheme but less than the exact Riemann solver. We also point out that the problem may be partially mitigated by a proper fine-tunings of the parameters involved in the reconstruction and steepening algorithms. However, we did not follow that approach in the present work.
### 4.7 Two-Dimensional Riemann Problem
Two dimensional Riemann problems involve the interactions of four elementary waves (either shocks, rarefactions, and contact discontinuities) initially separating four constant states. They have been formulated by Schulz et al. (1993), Lax & Liu (1998) in the context of classical hydrodynamics. Here we consider a relativistic extension, initially proposed by Del Zanna & Bucciantini (2002), where the initial configuration involves two shocks and two tangential discontinuities.
The domain is the square $`[1,1]\times [1,1]`$ covered with $`400^2`$ computational zones. The four quadrants $`NE`$ ($`x,y>0`$), $`NW`$ ($`x<0<y`$), $`SW`$ ($`x,y<0`$), $`SE`$ ($`y<0<x`$) divide the square into four constant-state regions:
$$(\rho ,v_x,v_y,p)=\{\begin{array}{cc}(0.1,0,0,0.01)\hfill & \text{for}x,y>0,\hfill \\ \multicolumn{2}{c}{}\\ (0.1,0.99,0,1)\hfill & \text{for}x<0<y,\hfill \\ \multicolumn{2}{c}{}\\ (0.5,0,0,1)\hfill & \text{for}x,y<0,\hfill \\ \multicolumn{2}{c}{}\\ (0.1,0,0.99,1)\hfill & \text{for}y<0<x.\hfill \end{array}$$
(38)
We use an ideal equation of state with $`\mathrm{\Gamma }=5/3`$. The integration is carried out with CFL=0.4 till $`t=0.8`$.
Notice that the initial condition (38) does not exactly prescribe two simple shock waves at the $`NWNE`$ and $`SENE`$ interface. The correct version of this problem has been considered by Mignone et al. (2005). For the sake of comparison, however, we chose to adopt the same initial condition as in Del Zanna & Bucciantini (2002).
The top and bottom panels in Fig. 11 show, respectively, the solutions computed with the HLLC and HLL solvers. The breakup of the initial discontinuity results in two equal-strength curved shock fronts propagating from regions $`NW`$ and $`SE`$ into the upper right portion of the domain ($`NE`$), top panel of Fig. 11. Region $`SW`$ is bounded by two tangential discontinuities and a jet-like structure emerges along the main diagonal.
The initial steady tangential discontinuities, located at the $`W`$ and $`S`$ interfaces, remain automatically sharp in the HLLC formulation, since they are exactly captured by the approximate Riemann solver. The same results has also been shown by Mignone et al. (2005) who used a two-shock iterative nonlinear solver. We emphasize that this is property pertains to the Riemann solver itself and does not depend on the interpolation algorithm. Indeed the same result holds when the first-order scheme is employed. This feature is absent from the HLL formulation, where tangential discontinuities spread along the cartesian axis due to extra numerical diffusion. This behavior is manifestly evident in the grid-aligned spurious waves visible in the bottom panel of Fig. 11, see also Del Zanna & Bucciantini (2002); Lucas-Serrano et al. (2004).
### 4.8 Axisymmetric Relativistic Jet
Finally, as an example of an astrophysical application, we consider the propagation of a light, axisymmetric relativistic jet in 2-D cylindrical geometry. For the sake of comparison, the parameters of the simulation are the same used by Del Zanna & Bucciantini (2002) and by Lucas-Serrano et al. (2004). The initial condition is prescribed as
$$(\rho ,v_r,v_z,p)=\{\begin{array}{cc}(0.1,\mathrm{\hspace{0.17em}0},\mathrm{\hspace{0.17em}0.99},\mathrm{\hspace{0.17em}10}^2)\hfill & \text{for}r,z<1,\hfill \\ \multicolumn{2}{c}{}\\ (10,\mathrm{\hspace{0.17em}0},\mathrm{\hspace{0.17em}0},\mathrm{\hspace{0.17em}10}^2)\hfill & \text{otherwise}.\hfill \end{array}$$
(39)
The jet is pressure matched and its internal relativistic Mach number is 17.1. We use an ideal equation of state with $`\mathrm{\Gamma }=5/3`$ both for the jet and the ambient medium. The computational domain covers the region $`0r12`$, $`0z35`$, with $`240\times 700`$ grid points, so that we have 20 cells per jet radius. At the symmetry axis, $`r=0`$, we impose reflecting boundary conditions; outflow boundary conditions are set everywhere else, except at the inlet region where we keep the beam constant. The CFL number is $`0.5`$ and the jet evolution is followed until $`t=80`$.
For the sake of comparison, we have also carried the simulation using the relativistic HLL solver. Figure 12 shows two snapshots of the rest mass density at times $`t=40`$ and $`t=80`$. The upper-half of each panel refers to the HLLC integration, whereas the lower portion displays the result obtained with the HLL scheme. We see that all the structural features characteristic of jet propagation can be clearly identified, with good resolution of shock waves and contact discontinuities. It is clear from the figures that the HLLC integration shows a significantly greater amount of small scale structure, that is not visible in the HLL results. This is due to the larger numerical diffusion introduced by the latter in subsonic regions which prevent sharp resolution of shear and tangential waves.
The average advance speed of the jet head is $`0.39`$ (to be compared with a onedimensional theoretical estimate of $`0.44`$; Martรญ et al. (1997)). Moreover, we can observe the absence of the carbuncle problem, that usually appears as an extended nose in front of the jet, on the axis (Quirk, 1994).
## 5 Efficiency Comparison
Previous numerical tests have shown that the quality of solution achieved with the HLLC scheme can be competitive with more complex exact or iterative non-linear Riemann solvers, see for example Martรญ & Mรผller (1996), Mignone et al. (2005). Another aspect which plays in favour of the HLLC formalism is the computational efficiency, particularly crucial in long-term simulations in two and three dimensions.
Table 1 gives the normalized CPU time required by the HLL, HLLC and approximate two-shock nonlinear Riemann solvers (for the latter see Mignone et al., 2005). All three solvers are available in the C authorโs code and have been written performing similar degree of optimizations. On the contrary, the FORTRAN code for the exact solution to the Riemann problem, available from Martรญ & Mรผller (2003), was found to be more than a factor of $`7`$ slower than the HLL solver. We believe that this might be due to a lower degree of optimization and by the time consuming numerical integration across the rarefaction fan (Pons et al., 2000).
For illustrative purposes, we consider the first four one-dimensional tests and the two-dimensional Riemann problem. Integrations have been carried using the first-order scheme on $`4000`$ and $`400^2`$ zones, respectively. No optimization flags were used during the compilation.
From the table, one can easily conclude that the HLLC scheme requires little additional costs with respect to the HLL approach (between $`4\%`$ and $`7\%`$ in one dimension), while the iterative nonlinear solver is certainly more expensive, being by a factor of more than $`30\%`$ slower.
In making the comparison, however, one should keep in mind that HLL-type solvers are iteration-free since the underlying algorithms always require a fixed number of operations, regardless of the initial condition. On the contrary, iterative nonlinear Riemann solvers have a certain degree of adaptability since the number of iterations to achieve convergence depends on the strength of the discontinuity at the zone interface. In smooth regions of the flow, for example, fewer iterations are usually needed. This explains why the fourth test problems is particularly time consuming, since a stronger discontinuity is involved.
In this respect, a direct comparison between different Riemann solvers becomes problem-dependent and can be used only as an order-of-magnitude estimate. Conversely, we do not expect the HLLC/HLL efficiency ratio to change with increasing complexity of the flow patterns. For this reason, we believe that for problems involving rich and complex structures the trade-off between computational efficiency and quality of results is certainly worth the effort.
## 6 Conclusions
We have presented, for the first time, an extension of the HLLC scheme by Toro et al. (1994) to relativistic gas dynamics. The solver is robust, computationally efficient and complete, in that it considers the full wave structure in the solution to the Riemann problem. The solver retains the attractive feature of being positively conservative, typical of the HLL scheme family.
The major improvement over the simple single-state HLL solver is the ability to resolve contact and tangential discontinuities. This property has been demonstrated by direct comparisons in several one- and two-dimensional test problems, where differences are strongly evident. The results indicate that the new HLLC solver attains sharper representation of discontinuities, quantitatively similar to the exact but algebraically and computationally more intensive Riemann solver.
The additional computational cost over the traditional HLL approach is less than $`8\%`$ and we believe that the improved quality of results largely justifies the trade-off between the two approximate Riemann solvers.
Extension to relativistic magnetized flows will be considered in a forthcoming paper. We notice, however, that the HLLC formalism presented in this work can be easily generalized to the case of vanishing normal component of magnetic field. When this degeneracy occurs (as in the propagation of jets with toroidal magnetic field, see for example Leismann et al., 2005), in fact, the solution to the Riemann problem is entirely analogous to the non-magnetized case, since only three waves are actually involved. This extension will be presented in Mignone et al. (2005).
Finally, we mention that the relativistic HLLC scheme does not make any assumption on the equation of state, and efforts to incorporate different equations of state should be minimal.
## Appendix A Proof of $`\lambda _L\lambda ^{}\mathrm{\Lambda }_R`$
In what follows we prove some important results concerning the positivity of the our relativistic HLLC scheme. The proof is given below in A.4. Propositions A.1 through A.3 demonstrate some preliminary results.
We assume that the fastest and slowest signal velocities are computed according to the prescription given in ยง3.1.1, and that $`\lambda _R>0`$ and $`\lambda _L<0`$, which is the case of applicability for the intermediate fluxes (13). Obviously, the initial left and right states are supposed to be physically admissible, i.e. they belong to the set $`G`$ defined in ยง3.1.2.
###### Proposition A.1
$`A_R>0`$, $`A_L<0`$
*
We will only prove $`A_R>0`$, since the proof for $`A_L`$ is similar. For the sake of clarity, we omit the subscript $`R`$. Using the definition of $`A`$ given after equation (17) one has
$$A=(E+p)\left[\lambda \left(1\frac{\sigma }{\mathrm{\Gamma }}\right)v_x\right],$$
(40)
where $`\sigma =c_s^2/\gamma ^2`$ is always positive numbers. Equation (40) is always positive for $`\lambda >\lambda _0`$, where
$$\lambda _0\frac{\mathrm{\Gamma }v_x}{\mathrm{\Gamma }\sigma }.$$
(41)
However, according to the choice given in ยง3.1.1, $`\lambda `$ must satify
$$f(\lambda )=(\lambda v_x)^2\frac{\sigma }{1c_s^2}(1\lambda ^2)0.$$
(42)
with $`\sigma _s=\sigma /(1c_s^2)`$. Equation (42) simply states that $`\lambda `$ must be greater than the root with the positive sign $`\lambda _+`$ (for the left state, $`\lambda _L`$ is always less than the root with the negative sign $`\lambda _{}`$). Direct substitution of $`\lambda _0`$ from (41) into (42) shows, after some algebra, that
$$f(\lambda _0)K_1\left[c_s^4+(v_x^2+2\mathrm{\Gamma })c_s^2\mathrm{\Gamma }^2\right],$$
(43)
where the equality occurs in the limit of zero tangential velocities and $`K_1`$ is always a positive quantity. Since $`1<\mathrm{\Gamma }<2`$ and $`c_s^2`$ has the limiting value $`(\mathrm{\Gamma }1)`$ the expression in square bracket is always negative, which means that $`\lambda \lambda _+>\lambda _0`$. This implies that $`AA_R`$ is always positive with our choice of $`\lambda \lambda _R`$.
Since one can prove, in a similar way, that $`A_L<0`$, we have the important results that $`U_E^{hll}=A_RA_L>0`$.
###### Proposition A.2
$`B_R+A_R>0`$, $`B_LA_L>0`$.
*
Again we only give the proof for the right state, the other case being similar. The function $`A+B`$ (the subscript $`R`$ is omitted), with $`A`$ and $`B`$ defined after equation (17) increases linearly with $`\lambda `$ and is positive for $`\lambda >\lambda _0`$, where
$$\lambda _0=\frac{v_x\gamma ^2\mathrm{\Gamma }(v_x+1)+c_s^2}{\gamma ^2\mathrm{\Gamma }(v_x+1)c_s^2}.$$
(44)
However, direct substitution of $`\lambda _0`$ in (42) shows, after extensive manipulations, that
$$f(\lambda _0)=K_2\left[c_s^4(1+2\mathrm{\Gamma })c_s^2+\gamma ^2\mathrm{\Gamma }^2(1v_x^2)\right],$$
(45)
where $`K_2`$ is always a positive quantity. It can be easily verified that the function in square brackets is always positive if $`c_s^2[0,\mathrm{\Gamma }1]`$ and $`1<\mathrm{\Gamma }<2`$. Thus we must have $`B_R/A_R>1`$.
###### Proposition A.3
$`\lambda _LA_RA_R<0`$, $`\lambda _RA_LB_L>0`$.
*
For the right state we have that
$$\lambda _LAB\lambda _{}AB\left(\frac{2v}{1+\sigma _s}\lambda \right)AB,$$
(46)
where the last inequality follows from the fact that the two roots of equation (42) satisfy
$$\lambda _{}=\frac{2v}{1+\sigma _s}\lambda _+,\text{and}\lambda _+\lambda .$$
(47)
Using the fact that $`\lambda ^2(1+\sigma _s)2\lambda vv^2+\sigma _s`$ and that $`A>0`$, the last expression in equation (46) can be shown to obey the following
$$\left(\frac{2v}{1+\sigma _s}\lambda \right)ABg,$$
(48)
where
$$g=K_3\left[v^2(\mathrm{\Gamma }c_s^21)+1\mathrm{\Gamma }+c_s^22c_s^2v_t^2\right],$$
(49)
with $`K_3`$ being a positive quantity. The expression in square bracket in equation (49) is always negative under the same assumptions used previously. Thus we have $`\lambda _L<B_R/A_R`$ and, similarly, one can prove that $`\lambda _R>B_L/A_L`$.
###### Proposition A.4
$`\lambda _L\lambda ^{}\lambda _R`$.
*
We now show that the choice of eigenvalues given in ยง3.1.1 always guarantees $`\lambda _L\lambda ^{}\lambda _R`$.
The starting point is to note that the quadratic equation (18) can be more conveniently written as
$$(A_L\lambda ^{}B_L)(1\lambda _R\lambda ^{})=(A_R\lambda ^{}B_R)(1\lambda _L\lambda ^{}),$$
(50)
which defines the intersection of two quadratic functions. The parabola on the left hand side vanishes in $`\lambda ^{}=1/\lambda _R>1`$ and $`\lambda ^{}=B_L/A_L<1`$, whereas the parabola on the right hand side in $`\lambda ^{}=1/\lambda _L<1`$ and $`\lambda ^{}=B_R/A_R>1`$. Moreover the two quadratics have the same concavity, since $`\text{sign}(A_L\lambda _R)=\text{sign}(A_R\lambda _L)=1`$. Thus the intersection must necessarily satisfy
$$\mathrm{min}(\frac{B_R}{A_R},\frac{B_L}{A_L})\lambda ^{}\mathrm{max}(\frac{B_R}{A_R},\frac{B_L}{A_L}).$$
(51)
However, for any $`\lambda (1,1)`$ one has
$$\lambda AB=(\lambda v_x)^2(E+p)+p(1\lambda ^2)>0,$$
(52)
which, together with the results previously shown, implies that
$$\begin{array}{c}1>\lambda _R>\mathrm{max}(\frac{B_R}{A_R},\frac{B_L}{A_L}),\hfill \\ \\ 1<\lambda _L<\mathrm{min}(\frac{B_R}{A_R},\frac{B_L}{A_L}).\hfill \end{array}$$
(53)
and hence
$$\lambda _L\lambda ^{}\lambda _R.$$
(54)
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# 1 Diagrams for contributions containing soft divergences.
NRQCD Factorization and Universality of NRQCD Matrix Elements
J.P. Ma
Institute of Theoretical Physics , Academia Sinica, Beijing 100080, China
Department of Physics, Shandong University, Jinan Shandong 250100, China
Z.G. Si
Department of Physics, Shandong University, Jinan Shandong 250100, China
## Abstract
The approach of nonrelativistic QCD(NRQCD) factorization was proposed to study inclusive production of a quarkonium. It is widely used and successful. However, a recent study of gluon fragmentation into a quarkonium at two-loop level shows that the factorization is broken. It is suggested that the color-octet NRQCD matrix elements should be modified by adding a gauge link to restore the factorization. The modified matrix elements may have extra soft-divergences at one-loop level which the unmodified can not have, and this can lead to a violation of the universality of these matrix elements. In this letter, we examine in detail the NRQCD factorization for inclusive quarkonium production in $`e^+e^{}`$ annihilation at one-loop level. Our results show that the factorization can be made without the modification of NRQCD matrix elements and it can also be made for relativistic corrections. It turns out that the suggested gauge link will not lead to nonzero contributions to color-octet NRQCD matrix elements at one-loop level and at any order of $`v`$. Therefore the universality holds at least at one-loop level.
A quarkonium system provide a unique place to study the dynamics of QCD, because a quarkonium mainly consists of a heavy quark pair $`Q\overline{Q}`$ and they move with a small velocity $`v`$. An extensive review of quarkonium physics can be found in . A decade ago the approach of NRQCD factorization was proposed to study inclusive production of a quarkonium. In this approach the production of a $`Q\overline{Q}`$ pair can be studied with perturbative QCD because the mass $`m`$ of $`Q`$ provides a large scale, the formation of the $`Q\overline{Q}`$ pair into a quarkonium is characterized with NRQCD matrix elements by an expansion in $`v`$. In order to have prediction power these matrix elements should be universal, i.e., they do not depend on how the $`Q\overline{Q}`$ pair is produced. This leads to the NRQCD factorization. This approach is widely used and successful. Especially, it can systematically take higher Fock state components of a quarkonium, including those components which contain a heavy quark pair in color-octet, into account. A striking success of the approach is to explain the $`\psi ^{}`$ anomaly at Tevatron by taking color-octet components into account.
Although the approach is successful, a complete proof of the factorization does not exist. Studies of various processes at one-loop level really show that the factorization hold at one-loop level. But a recent study shows that the factorization for gluon fragmentation into a quarkonium at two-loop level is incomplete, indicated by that some uncancelled infrared(I.R.) divergences are not matched by NRQCD matrix elements. To restore the factorization, a gauge link is introduced to modify color-octet NRQCD matrix elements for matching these I.R. divergences. At one-loop level, this gauge link can generate some extra I.R. divergences in the modified matrix elements than those unmodified. This can affect the factorization at one-loop level in the cases studied before and can lead to a violation of the universality of color-octet NRQCD matrix elements.
It is the purpose of the letter to examine if the universality is lost and if the factorization can be done in inclusive production of a quarkonium through $`e^+e^{}`$-annihilation through a color-octet $`Q\overline{Q}`$ pair. Our analysis includes not only the leading contributions in the $`v`$-expansion, but also the contribution from relativistic corrections at order of $`v^2`$. We show at one-loop level in detail how soft divergences are cancelled or matched by color-octet NRQCD matrix elements without the gauge link suggested in . Our study also shows that the relativist correction for a color-octet $`Q\overline{Q}`$ pair can also be factorized in the same way. To our knowledge, there is no known example studied at one-loop level to show that the NRQCD factorization holds.
We consider the inclusive production of a quarkonium $`H`$ in the process
$$e^+e^{}\gamma ^{}H+X,$$
(1)
where the virtual photon is with the momentum $`q_\gamma `$ and $`q_\gamma ^2`$ is much larger than the square of the quarkonium mass. For this process we need to calculate the tensor
$$T^{\mu \nu }(P,q_\gamma ,H)=\underset{X}{}d^4xe^{iq_\gamma x}0|J^\nu (0)|H(P)+XH(P)+X|J^\mu (x)|0$$
(2)
where the quarkonium carries $`P`$ and $`J^\mu `$ is the electric current. The NRQCD factorization in suggests the tensor can be written in a factorized form:
$`T^{\mu \nu }(P,q_\gamma ,H)`$ $`=`$ $`F^{\mu \nu }(^1S_0)0|O_8(^1S_0,H)|0+G^{\mu \nu }(^1S_0)0|P_8(^1S_0,H)|0`$ (3)
$`+F^{\mu \nu }(^3P_0)0|O_8(^3P_0,H)|0+F^{\mu \nu ,ijkl}(^3P_2)0|O_8^{\{ij\},\{kl\}}(^3P_2,H)|0+\mathrm{},`$
where the matrix elements in the right hand side are defined with NRQCD fields and can be found in . We only consider the production through those color-octet channels which can be at the leading order of $`\alpha _s`$, i.e., the channel with the quantum number $`{}_{}{}^{1}S_{0}^{}`$ and $`{}_{}{}^{3}P_{J}^{}`$ with $`J=0,2`$. There is a velocity-power counting rule to determine the relative importance of NRQCD matrix elements for a given quarkonium. In the factorized form the second term is for the relativistic correction of the channel $`{}_{}{}^{1}S_{0}^{}`$. The coefficients in the front of the NRQCD matrix elements can be calculated with perturbative QCD, they are series in $`\alpha _s`$, e.g.,
$$G^{\mu \nu }(^1S_0)=G_0^{\mu \nu }(^1S_0)+G_1^{\mu \nu }(^1S_0)+\mathrm{},$$
(4)
where the subscriber $`0(1)`$ stand for the tree(one-loop) contribution. If the factorization holds, these perturbative coefficients should not contain any I.R. divergence. To determine these coefficients, one replaces the quarkonium $`H`$ with a $`Q\overline{Q}`$ state and calculates $`T^{\mu \nu }`$ and the NRQCD matrix elements. By comparing both sides of Eq.(3) calculated with the $`Q\overline{Q}`$ pair the perturbative coefficients can be extracted. If the factorization holds, soft divergences in $`T^{\mu \nu }`$ will have the same form as those appearing in the matrix elements so that the perturbative coefficients do not contain any soft divergence.
To study the factorization we need to calculate the tensor $`T^{\mu \nu }`$ in perturbative theory after replacing the quarkonium $`H`$ with those $`Q\overline{Q}`$ states:
$$T^{\mu \nu }(P,q_\gamma ,Q\overline{Q})=\underset{X}{}d^4xe^{iq_\gamma x}0|J^\nu (0)|Q(p_1^{})\overline{Q}(p_2^{})+X\overline{Q}(p_2)Q(p_1)+X|J^\mu (x)|0.$$
(5)
It should be noted that the heavy quarks carry different momenta in the amplitude and its complex conjugated. This will enable us to identify different states of the $`Q\overline{Q}`$ pair. These momenta are given as:
$$p_1=\frac{1}{2}P+q,p_2=\frac{1}{2}Pq,p_1^{}=\frac{1}{2}P+q^{},p_2^{}=\frac{1}{2}Pq^{}.$$
(6)
In the rest frame of the $`Q\overline{Q}`$, $`q^\mu =(0,๐ช)`$ and $`q^\mu =(0,๐ช^{})`$ with $`๐ช^2=๐ช^2`$ and
$$P^2=(p_1+p_2)^2=4(m^2+๐ช^2),๐ฏ=\frac{๐ช}{m},๐ฏ^{}=\frac{๐ช^{}}{m},$$
(7)
$`v(v^{})`$ is the velocity of the heavy quark in the rest frame. At tree level, the unobserved state $`X`$ contains only one gluon. By expanding $`๐ฏ`$ and $`๐ฏ^{}`$ and identifying quantum numbers, one can determine four coefficients in Eq.(3), i.e.,
$`T_0^{\mu \nu }(P,q_\gamma ,Q\overline{Q})`$ $`=`$ $`F_0^{\mu \nu }(^1S_0)0|O_8(^1S_0,Q\overline{Q})|0+G_0^{\mu \nu }(^1S_0)0|P_8(^1S_0,Q\overline{Q})|0`$ (8)
$`+F_0^{\mu \nu }(^3P_0)0|O_8(^3P_0,Q\overline{Q})|0+F_0^{\mu \nu ,ijkl}(^3P_2)0|O_8^{\{ij\},\{kl\}}(^3P_2,Q\overline{Q})|0`$
$`+๐ช(v^4).`$
In the expansion in $`๐ฏ`$ and $`๐ฏ^{}`$, the leading terms give contributions to $`F_0^{\mu \nu }(^1S_0)`$ for the $`{}_{}{}^{1}S_{0}^{}`$ state, the next-to-leading terms, which are linear in $`๐ฏ`$ and $`๐ฏ^{}`$ like $`๐ฏ๐ฏ^{}`$, give contributions to $`F_0^{\mu \nu }(^3P_0)`$ for the $`{}_{}{}^{3}P_{0}^{}`$ state and $`F_0^{\mu \nu ,ijkl}(^3P_2)`$ for the $`{}_{}{}^{3}P_{2}^{}`$ state. The next-to-next-to-leading terms are proportional either to the tensor $`v^iv^j`$ or to $`v^iv^j`$. One can decompose the tensor $`v^iv^j`$ or to $`v^iv^j`$ into the component of the $`S`$-wave with $`l=0`$ and of the $`D`$-wave with $`l=2`$. The $`S`$-wave component corresponds to the relativistic correction of the $`{}_{}{}^{1}S_{0}^{}`$ state. At tree level all these coefficients are nonzero and contain no soft divergence, their detailed forms are not important for our purpose because we will show that the soft-divergent correction at one-loop to these coefficients is proportional to the tree-level result $`T_0^{\mu \nu }`$. We will use Feynman gauge in this letter. In this gauge one can clearly see how soft divergences are cancelled or matched in a diagram-by-diagram manner.
One-loop corrections consist of two parts. One is the virtual correction, another is the real correction in which the unobserved state $`X`$ contains two gluons or a light quark pair. Beside corrections from wave-function renormalization there are many Feynman diagrams. Since we are only interested in soft divergences, we do not need to consider all diagrams but only those containing soft divergences. The diagrams with soft divergences are given in Fig.1.. To obtain the soft-divergent parts of these diagrams we employ the eikonal approximation with some modification.
To illustrate the eikonal approximation used here, we consider the contribution from Fig.1a to the matrix element $`R^\mu =\overline{Q}(p_2)Q(p_1),G(k,\epsilon ^{},a)|J^\mu (x)|0`$:
$`R_{1a}^\mu `$ $`=`$ $`{\displaystyle \frac{d^4k_1}{(2\pi )^4}\overline{u}(p_1)(ig_sT^b\gamma ^\rho )\frac{\gamma (p_1+k_1)+m}{(p_1+k_1)^2m^2}(ieQ\gamma ^\mu )\frac{\gamma (k_2p_2)+m}{(p_2+k_2)^2m^2}(ig_sT^c\gamma ^\sigma )v(p_2)}`$ (9)
$`{\displaystyle \frac{1}{k_1^2}}{\displaystyle \frac{1}{k_2^2}}(g_sf^{abc})\left[(kk_1)_\sigma \epsilon _\rho ^{}+(k_1k_2)\epsilon ^{}g_{\rho \sigma }+(k_2+k)_\rho \epsilon _\sigma ^{}\right],`$
with $`k_2=kk_1`$. The soft divergence appears when $`k_1`$ or $`k_2`$ becomes soft, i.e., all components of $`k_1`$ or $`k_2`$ becomes small, and when $`k_1`$, hence also $`k_2`$, is collinear to the momentum $`k`$ of the outgoing gluon. If the gluon with $`k_1`$ is soft, the standard approximation is to neglect all $`k_1`$ in nominators and keep only the leading order in $`k_1`$ in the denominators. Therefore the contribution from the soft region of $`k_1`$ can be written as:
$`R_{1a,s1}^\mu `$ $`=`$ $`ieQg_s^3{\displaystyle \frac{d^4k_1}{(2\pi )^4}\overline{u}(p_1)f^{abc}T^bT^c\frac{2p_1^\rho }{2p_1k_1}\gamma ^\mu \frac{\gamma (p_2k)+m}{2p_2k}\gamma ^\sigma v(p_2)}`$ (10)
$`{\displaystyle \frac{1}{k_1^2}}{\displaystyle \frac{1}{(2k_1k)}}\left[k_\sigma \epsilon _\rho ^{}k\epsilon ^{}g_{\rho \sigma }+2k_\rho \epsilon _\sigma ^{}\right].`$
We use dimensional regularization with $`d=4ฯต`$ to regularize divergences. The pole with $`1/ฯต_I`$ is for I.R. divergence with $`ฯต_I=4d`$. Other poles without the subscriber $`I`$ are for U.V. divergences. Calculating the integral we find that the soft amplitude contains double pole in the form $`ฯต^1(ฯต^1ฯต_I^1)`$, indicating that the approximation may not be convenient here. A convenient approximation we will use is to keep denominators exact. With this approximation, the soft- and collinear-divergent part of the contribution can be written as
$`R_{1a,sc}^\mu `$ $`=`$ $`ieQg_s^3{\displaystyle \frac{d^4k_1}{(2\pi )^4}\frac{1}{(2p_1k_1+i0)(2p_2k_2+i0)(k_1^2+i0)(k_2^2+i0)}}`$ (11)
$`\overline{u}(p_1)f^{abc}T^bT^c\left[(2p_1^\rho +\gamma ^\rho \gamma ^+k_1^{})\gamma ^\mu (2p_2^\sigma k_2^{}\gamma ^+\gamma ^\sigma )\right]v(p_2)`$
$`\left[(kk_1)^{}n_\sigma \epsilon _\rho ^{}+(2k_1k)^{}n\epsilon ^{}g_{\rho \sigma }+(k_2+k)^{}n_\rho \epsilon _\sigma ^{}\right],`$
and $`R_{1a}^\mu R_{1a,sc}^\mu `$ contains no any soft divergence. In the above equation we have taken a frame in which
$$k^\mu =(0,k^{},0,0)=k^{}n^\mu ,q_\gamma ^\mu =(q_\gamma ^+,q_\gamma ^{},0,0).$$
(12)
It should be noted that $`R_{2a,sc}^\mu `$ can be written in a covariant form. Performing a similar analysis for other diagrams and loop-momentum integration we obtain the soft-divergent part of the one-loop correction to $`R^\mu `$. The sum of contributions from Fig.1a, Fig.1d and Fig.1e can be written in a compact form:
$`R_{1a,sc}^\mu +R_{1d,sc}^\mu +R_{1e,sc}^\mu `$ $`=`$ $`{\displaystyle \frac{ieQg_s^3}{(4\pi )^2}}N_c\left\{2\left({\displaystyle \frac{2}{ฯต_I}}\right)^2{\displaystyle \frac{2}{ฯต_I}}\left[2\gamma 2+\mathrm{ln}{\displaystyle \frac{(2p_1k)^2}{4\pi m^2\mu ^2}}+\mathrm{ln}{\displaystyle \frac{(2p_2k)^2}{4\pi m^2\mu ^2}}\right]\right\}`$ (13)
$`\overline{u}(p_1)T^a\left[\gamma ^\nu {\displaystyle \frac{\gamma (p_1+k)+m}{(p_1+k)^2m^2}}\gamma ^\mu +\gamma ^\mu {\displaystyle \frac{\gamma (p_2k)+m}{(p_2+k)^2m^2}}\gamma ^\nu \right]v(p_2),`$
i.e., the soft-divergent part is proportional to the tree-level amplitude. The double pole in $`ฯต_I`$ is from the overlap of collinear- and soft region of the loop momentum. The contribution from Fig.1b and Fig.1c contain not only soft divergences but also Coulomb singularities. With our approximation the sum of these two diagrams gives:
$`R_{2b,cs}^\mu +R_{2c,cs}^\mu `$ $`=`$ $`{\displaystyle \frac{ieQg_s^3}{(4\pi )^2}}{\displaystyle \frac{1+2v^2}{N_c}}\left[{\displaystyle \frac{\pi ^2}{2v}}+{\displaystyle \frac{2}{ฯต_I}}\left(1{\displaystyle \frac{2v^2}{3}}\right)\right]`$ (14)
$`\epsilon _\nu ^{}\overline{u}(p_1)T^a\left[\gamma ^\nu {\displaystyle \frac{\gamma (p_1+k)+m}{(p_1+k)^2m^2}}\gamma ^\mu +\gamma ^\mu {\displaystyle \frac{\gamma (p_2k)+m}{(p_2+k)^2m^2}}\gamma ^\nu \right]v(p_2)+\mathrm{},`$
where $`\mathrm{}`$ stand for terms which are finite with $`ฯต,v0`$, and higher orders in $`v`$.
Putting everything together, we obtain the soft divergent part of the virtual one-loop contribution to $`T^{\mu \nu }`$ in Feynman gauge:
$`T_{1,vir.}^{\mu \nu }(P,q_\gamma ,Q\overline{Q})`$ $`=`$ $`{\displaystyle \frac{\alpha _s}{2\pi }}{\displaystyle \frac{1}{N_c}}\left[{\displaystyle \frac{2}{ฯต_I}}\left(1+{\displaystyle \frac{2}{3}}v^2+{\displaystyle \frac{2}{3}}v^2+๐ช(v^4)\right)+{\displaystyle \frac{\pi ^2}{2v}}(1+๐ช(v^2))\right]T_0^{\mu \nu }(P,q_\gamma ,Q\overline{Q})`$
$`{\displaystyle \frac{\alpha _s}{2\pi }}N_c\left[\left({\displaystyle \frac{2}{ฯต_I}}\right)^2{\displaystyle \frac{2}{ฯต_I}}\left(\gamma 1+L_n\right)\right]T_0^{\mu \nu }(P,q_\gamma ,Q\overline{Q})`$
$`{\displaystyle \frac{\alpha _s}{2\pi }}\left({\displaystyle \frac{2}{ฯต_I}}\right)\left[{\displaystyle \frac{N_c^21}{N_c}}+\left({\displaystyle \frac{5}{6}}N_c{\displaystyle \frac{1}{3}}N_f\right)\right]T_0^{\mu \nu }(P,q_\gamma ,Q\overline{Q})+\mathrm{}`$
$`L_n`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{ln}{\displaystyle \frac{(p_1k)(p_1^{}k)(p_2k)(p_2^{}k)}{(\pi m^2\mu ^2)^2}},`$ (15)
where $`\mathrm{}`$ denote finite parts and the correction terms in the third line come from wave-function renormailization of the heavy quarks and the gluon.
Now we turn to the real correction. The real correction consists of amplitudes of two gluons or a light quark pair. The two gluons can be emitted by heavy quarks and by gluon splitting, the light quark pair can only be generated through gluon splitting. It should be noted that a light quark pair can also be produced in other ways, but the contribution of this is irrelevant for the factorization discussed here, because in this case the heavy quark pair is in states other than those given in Eq.(8) and the contribution at order $`\alpha _s^2`$ is free from soft divergences. The Feynman diagrams for the real correction is given in Fig.2.
We write the amplitude with two gluons in the final state as
$$R_{gg}^\mu =R_{gg,Q}^\mu +R_{gg,g}^\mu ,$$
(16)
where the first term is only for two gluons emitted by heavy quarks and second is for gluon splitting. If two gluons are emitted by heavy quarks, there are only I.R. divergences, when one of the gluons becomes soft. Collinear singularities will not appear here because of the massive quarks. When the gluon with the momentum $`k_1`$ becomes soft, the soft contribution of $`R_{gg,Q}^\mu `$ with the standard eikonal approximation can be written:
$$R_{gg,Q,s}^\mu =\frac{ig_s}{2}f^{a_1a_2a}\left(\frac{p_1\epsilon _1^{}}{p_1k_1}+\frac{p_2\epsilon _1^{}}{p_2k_1}\right)\overline{Q}(p_2)Q(p_1),G(k_2,\epsilon _2,a)|J^\mu (0)|0|_{Treelevel}+\mathrm{},$$
(17)
where we only give those contributions explicitly, which are from the color-octet $`Q\overline{Q}`$ pair with the quantum numbers $`{}_{}{}^{1}S_{0}^{}`$ and $`{}_{}{}^{3}P_{J}^{}`$ with $`J=0,2`$. The contributions from states different than those above are indicated by $`\mathrm{}`$. They are not important for our purpose, but they are important for the complete cancellation of soft divergences in the KLN theorem. We will discuss this later. The contribution from $`R_{gg,Q,s}^\mu `$ to $`T^{\mu \nu }`$ can be calculated easily, we obtain:
$$T_{1,Q,s}^{\mu \nu }(P,q_\gamma ,Q\overline{Q})=\frac{\alpha _s}{2\pi }N_c\left(\frac{2}{ฯต_I}\right)\left[1+\frac{1}{3}v^2+\frac{1}{3}v^2+๐ช(v^4)\right]T_0^{\mu \nu }(P,q_\gamma ,Q\overline{Q})+\mathrm{}$$
(18)
where we have only collected the contributions from relevant $`Q\overline{Q}`$ states. Contributions from other states and higher orders of $`v`$ are represented by $`\mathrm{}`$.
The amplitude from the gluon splitting is given by:
$`R_{gg,g}^\mu `$ $`=`$ $`eQg_s^2f^{aa_1a_2}\overline{u}(p_1)T^a\left[\gamma ^\nu {\displaystyle \frac{\gamma (p_1+k)+m}{(p_1+k)^2m^2}}\gamma ^\mu +\gamma ^\mu {\displaystyle \frac{\gamma (p_2k)+m}{(p_2+k)^2m^2}}\gamma ^\nu \right]v(p_2)`$ (19)
$`\left[(k+k_1)\epsilon _2^{}\epsilon _{1\nu }^{}+(k_1+k_2)_\nu \epsilon _1^{}\epsilon _2^{}+(k_2k)\epsilon _1^{}\epsilon _{2\nu }^{}\right]{\displaystyle \frac{1}{(k_1+k_2)^2}},`$
with $`k=k_1+k_2`$. The contribution from $`R_{gg,g}^\mu `$ to $`T^{\mu \nu }`$ is more difficult to evaluate than that from $`R_{gg,Q}^\mu `$, because there is a overlap between the collinear- and soft region of gluon momenta. However there is a standard method, called phase-space slicing method, discussed in detail in . We will use the method. If the gluon with $`k_1`$ is soft, the amplitude can be approximated by:
$`R_{gg,g,s}^\mu `$ $`=`$ $`eQg_s^2f^{aa_1a_2}\overline{u}(p_1)T^a\left[\gamma ^\nu {\displaystyle \frac{\gamma (p_1+k_2)+m}{(p_1+k_2)^2m^2}}\gamma ^\mu +\gamma ^\mu {\displaystyle \frac{\gamma (p_2k_2)+m}{(p_2+k_2)^2m^2}}\gamma ^\nu \right]v(p_2)`$ (20)
$`{\displaystyle \frac{(2k_2\epsilon _1^{})\epsilon _{2\nu }^{}}{(k_1+k_2)^2}}.`$
This amplitude $`R_{gg,g,s}^\mu `$ interfered with the $`R_{gg,Q,s}^\mu `$ will give contributions with the soft divergences which are exactly those in the contributions from the interference between $`R_{gg,Q}^\mu `$ and $`R_{gg,g}^\mu `$. By the method mentioned above we have the contribution from the interference:
$`T_{1,int.}^{\mu \nu }(P,q_\gamma ,Q\overline{Q})={\displaystyle \frac{\alpha _s}{2\pi }}N_c\left\{\left({\displaystyle \frac{2}{ฯต_I}}\right)^2{\displaystyle \frac{2}{ฯต_I}}\left[\gamma +\mathrm{ln}{\displaystyle \frac{s_{min}^2}{\pi m^2\mu ^2}}\right]\right\}T_0^{\mu \nu }(P,q_\gamma ,Q\overline{Q}),`$ (21)
where $`s_{min}`$ is a cut-off in the phase-space slicing method. Our final result will not depend on it. The contributions only from the gluon splitting into two gluons and a light quark pair can be calculated directly with the phase-space slicing method. They contain only collinear singularities. The result is:
$`T_{1,col}^{\mu \nu }(P,q_\gamma ,Q\overline{Q})={\displaystyle \frac{\alpha _s}{2\pi }}\left({\displaystyle \frac{2}{ฯต_I}}\right)\left\{N_c\left[L_n+\mathrm{ln}{\displaystyle \frac{s_{min}^2}{\pi m^2\mu ^2}}+{\displaystyle \frac{11}{6}}\right]{\displaystyle \frac{N_f}{3}}\right\}T_0^{\mu \nu }(P,q_\gamma ,Q\overline{Q}),`$ (22)
where the same dependence on the cut-off $`s_{min}`$ appears and it will be cancelled by that in $`T_{1,int.}^{\mu \nu }`$. Putting everything together we obtain the infrared divergent part of the real correction as:
$`T_{1,real}^{\mu \nu }(P,q_\gamma ,Q\overline{Q})`$ $`=`$ $`{\displaystyle \frac{\alpha _s}{2\pi }}\left\{N_c\left[\left({\displaystyle \frac{2}{ฯต_I}}\right)^2{\displaystyle \frac{2}{ฯต_I}}\left(\gamma 1+L_n\right)\right]\right\}T_0^{\mu \nu }(P,q_\gamma ,Q\overline{Q})`$ (23)
$`+{\displaystyle \frac{\alpha _s}{2\pi }}{\displaystyle \frac{2}{ฯต_I}}\left\{N_c\left[{\displaystyle \frac{11}{6}}+{\displaystyle \frac{1}{3}}v^2+{\displaystyle \frac{1}{3}}v^2\right]{\displaystyle \frac{N_f}{3}}\right\}T_0^{\mu \nu }(P,q_\gamma ,Q\overline{Q})+\mathrm{},`$
Finally, we obtain the total soft-divergent part of $`T^{\mu \nu }`$ at one loop:
$`T_1^{\mu \nu }(P,q_\gamma ,Q\overline{Q})`$ $`=`$ $`{\displaystyle \frac{\alpha _s}{2\pi }}\left[(v^2+v^2)\left({\displaystyle \frac{2}{ฯต_I}}\right){\displaystyle \frac{N_c^22}{3N_c}}{\displaystyle \frac{1}{N_c}}{\displaystyle \frac{\pi ^2}{2v}}\right]T_0^{\mu \nu }(P,q_\gamma ,Q\overline{Q})+\mathrm{},`$ (24)
where we only give the relevant part in detail. It should be noted that the $`[\mathrm{}]`$ of the part does not contain $`๐ฏ๐ฏ^{}`$ from our calculation at the orders considered here. The $`\mathrm{}`$ denote contributions from other states of the $`Q\overline{Q}`$ pair, which can not be produced at tree-level. These contributions are from $`R_{gg,Q}^\mu `$, represented by $`\mathrm{}`$ in Eq.(18), they contain terms like $`๐ฏ๐ฏ^{}`$ because the $`Q\overline{Q}`$ pair at one-loop order can be in a color-singlet $`P`$-wave states.
Before we turn to our result of NRQCD matrix elements to finalize NRQCD factorization, we briefly discuss here the cancellation of soft divergences in KLN theorem. It should be noted that in $`T^{\mu \nu }`$ defined in Eq.(5) the heavy quark pair can be in an arbitrary state which is allowed to be produced at a given order of $`\alpha _s`$, i.e., we do not sum over all possible states of the heavy quark pair. If we sum all these states and take $`๐ฏ=๐ฏ^{}`$, the sum should not contain any soft divergence, as stated by KLN theorem. The sum contains not only those terms explicitly given in Eq.(24), but also those represented by $`\mathrm{}`$, i.e., the contributions from those states which can not be produced at tree level, these states are produced through emission of two gluons by heavy quarks, whose contributions are represented in Eq.(18) by $`\mathrm{}`$. Taking these contributions into account we have the soft-divergent part of the sum:
$`{\displaystyle T_1^{\mu \nu }(P,q_\gamma ,Q\overline{Q})}|_{๐ฏ=๐ฏ^{}}`$ $`=`$ $`{\displaystyle \frac{\alpha _s}{\pi }}v^2\left({\displaystyle \frac{2}{ฯต_I}}\right){\displaystyle \frac{N_c^22}{3N_c}}{\displaystyle T_0^{\mu \nu }(P,q_\gamma ,Q\overline{Q})}|_{๐ฏ=๐ฏ^{}}`$ (25)
$`{\displaystyle \frac{\alpha _s}{\pi }}v^2\left({\displaystyle \frac{2}{ฯต_I}}\right){\displaystyle \frac{N_c^22}{3N_c}}{\displaystyle T_0^{\mu \nu }(P,q_\gamma ,Q\overline{Q})}|_{๐ฏ=๐ฏ^{}}+\mathrm{},`$
where the second term comes from the contributions of those states which are not given explicitly in Eq.(18). By summing of these states and some manipulation the sum can be written into a compact form as given in the above. We see clearly that the sum is free from any soft divergence as KLN theorem states.
The one-loop correction of those NRQCD matrix elements in Eq.(8) and Eq.(12) can be divided into two parts as a virtual part and a real part. The virtual part consists of corrections from vertex and wave-function renormailization. Examples of Feynman diagrams at one-loop level for the real part, vertex- and wave-function renormailization correction are given in Fig.3. up to the orders of $`v`$ we consider. We calculate these corrections in Feynman gauge. The one loop correction can be written as:
$$0|O_n^{Q\overline{Q}}|0_1=0|O_n^{Q\overline{Q}}|0_0\left[\frac{\pi \alpha _s}{4v}\frac{1}{N_c}(1+๐ช(v^2))+๐ช(v^2)\right]$$
(26)
where $`O_n^{Q\overline{Q}}`$ is any of $`O_8(^1S_0)`$, $`P_8(^1S_0)`$, $`O(^3P_J)`$ with $`J=0,1,2`$, the matrix element with the subscriber $`0(1)`$ is the tree-level(one-loop) result of the operator. The first term represents the Coulomb singularity, the second term starts at the relative order of $`v^2`$, combined with $`0|O_n^{Q\overline{Q}}|0_0`$ they should be written in the form as matrix elements of operators. Using this result we can clearly see that the Coulomb singularity in Eq.(24) is reproduced. Hence, one can conclude that all perturbative coefficients $`F^{}s`$ in Eq.(3) are free from the Coulomb singularity and also from I.R. divergences.
To assess if $`G^{\mu \nu }(^1S_0)`$ contains any I.R. divergence, we need to take the operator mixing between $`O_8(^1S_0)`$ and others into account which is contained in the second term in Eq.(26). At one loop and order $`v^2`$, the operator can be mixed with the color-singlet operator with quantum numbers $`{}_{}{}^{1}P_{1}^{}`$ and the color-octet operator $`P_8(^1S_0,Q\overline{Q})`$. The later is relevant for our case. In Feynman gauge we have
$`0|O_8(^1S_0,Q\overline{Q})|0`$ $`=`$ $`{\displaystyle \frac{\alpha _s}{2\pi }}{\displaystyle \frac{1}{N_c}}\left({\displaystyle \frac{2}{ฯต_I}}\right)(0|O_8(^1S_0,Q\overline{Q})|0_0+{\displaystyle \frac{2}{3m^2}}0|P_8(^1S_0,Q\overline{Q})|0_0+๐ช(v^4))`$
$`{\displaystyle \frac{\alpha _s}{2\pi }}{\displaystyle \frac{N_c^21}{N_c}}\left({\displaystyle \frac{2}{ฯต_I}}\right)0|O_8(^1S_0,Q\overline{Q})|0_0`$
$`+{\displaystyle \frac{\alpha _s}{2\pi }}N_c\left({\displaystyle \frac{2}{ฯต_I}}\right)(0|O_8(^1S_0,Q\overline{Q})|0_0+{\displaystyle \frac{1}{3m^2}}0|P_8(^1S_0,Q\overline{Q})|0_0)`$
$`+\mathrm{},`$
$`0|P_8(^1S_0,Q\overline{Q})|0_0`$ $`=`$ $`m^2(v^2+v^2)0|O_8(^1S_0,Q\overline{Q})|0_0,`$ (27)
where contributions at orders higher than $`v^2`$ or from mixing of other irrelevant operators and those with the Coulomb singularity are represented by $`\mathrm{}`$. In the above the first line comes from the vertex correction, the second from the external quark legs and the third comes from the real gluon exchange. With this result, one can see clearly how the soft divergences are absorbed into the matrix element at leading order of $`v`$ on a diagram-by-diagram basis. At order of $`v^0`$ all I.R. divergences are cancelled. At the next-to-leading order of $`v`$, the net divergence in the above equation exactly matches that in $`T_1^{\mu \nu }`$ and the matching is also on a diagram-by-diagram basis.
Our results clearly show that without the modification of relevant NRQCD matrix elements in our case the NRQCD factorization holds at order of $`v^2`$ and the relativistic correction can also be factorized. Adding a gauge link in the those color-octet NRQCD matrix elements, as suggested in , the modified matrix elements will receive extra contributions. These contributions relevant for our study are given by types of Feynman diagrams in Fig.4. We use Feynman gauge to study these contributions. In other gauges, like Coulomb gauge, a gluon can also be exchanged between gauge links. By an infrared power counting each contribution from Fig.4. is with infrared divergences. The virtual contribution from the type(a) of diagrams in Fig.4, after integrating of the energy of the virtual gluon and neglecting terms which generate power divergences, is proportional to the integral:
$`{\displaystyle \frac{d^3q}{(2\pi )^3}\frac{1}{2|๐ช|}\frac{1}{|๐ช|+q^3}\left[\frac{1+v^3}{|๐ช|+๐ฏ๐ช}+\frac{1v^3}{|๐ช|๐ฏ๐ช}\right]}.`$ (28)
In the above the factor $`1/(|๐ช|+q^3)`$ comes from the eikonal propagator of the gauge link, which is determined by the moving direction of the quarkonium. We take the direction in the $`z`$-direction. It is clearly that the integral is infrared-divergent. But the real contribution from the type(b) of diagrams in Fig.4 is also proportional to this integral and the proportional coefficient has different sign than that of the virtual contribution. The total contribution from the gauge link to the color-octet matrix elements is zero at one loop and at any order of $`v`$. Therefore, at one-loop level, the suggested gauge link will not lead to a violation of the universality of color-octet matrix elements and it will not affect all existing one-loop results.
To summarize: We have examined in detail the NRQCD factorization in inclusive production of a quarkonium through $`e^+e^{}`$-annihilation. Our results show that the factorization can be made for production of a $`Q\overline{Q}`$ pair in color octet with the quantum number $`{}_{}{}^{1}S_{0}^{}`$ and $`{}_{}{}^{3}P_{0,2}^{}`$ and also for relativistic correction to the $`S`$-wave state. The modification of color-octet NRQCD matrix elements with the suggested gauge link will not affect the NRQCD factorization at one loop in cases studied before and the case studied here, because the gauge link does not lead to nonzero contributions to color-octet NRQCD matrix elements at one loop. Because of this the universality of these matrix elements holds at one loop.
Acknowledgements
The authors would like to thank Prof. J.W. Qiu for communications about the recent work and Prof. G. Bodwin for intensive discussions, which greatly help to understand the problem studied here. The warm hospitality of the Taipei Summer Institute of NCTS/TPE at the Physics Department of NTU, where the first draft of the paper is completed, is acknowledged. This work is supported by National Nature Science Foundation of P. R. China.
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# Introduction
## Introduction
It has been shown in many works of the last three decades that combinatorial methods applied to Commutative Algebra and Algebraic Geometry are very effective. Most of these, though, deal with standard-graded polynomial rings, i.e. with polynomial rings where the weight or degree of all the variables involved is $`1`$.
Other than the work \[A\], which is concerned with positively graded algebras having a specific Hilbert function (see \[S\], Chapter 10 for a nice survey) and the introductory work \[BR\], there is a small amount of literature about non-standard graded (or weighted) algebras, where the weight of a variable can be any positive integer.
This lack persuaded us to explore the realm of weighted graded algebras with algebraic and combinatorial tools.
Our work is divided in four sections, each centred on one topic and inspired by some of the more significant results in the standard case: generic initial ideals, Prime Avoidance, Castelnuovo-Mumford regularity and lexicographic ideals.
The first section is dedicated to the definition and main properties of generic initial ideals. We first study and describe the automorphisms of a weighted polynomial ring, which are necessary for the definition, the existence of generic initial ideals is discussed and a Borel-fixed type of property, i.e. fixedness under the action of a subgroup of the group of the automorphisms is proved. Moreover, we define a combinatorial counterpart of this property, i.e. being โweighted strongly stableโ, and prove that generic initial ideals are enriched with it.
In the second section we recover an analogue of the homogeneous prime avoidance Lemma (Lemma 2.1), which grants the existence of an almost-regular form of degree equal to the least common multiples of the weights. This weaker statement is though enough to evince some conclusions which generalize the known fact that depth does not change after taking generic initial ideals with respect to the degree reverse-lexicographic order (Proposition 2.7). This is performed under the assumption that each weight is divisible by the previous ones.
The third section takes into consideration the Castelnuovo-Mumford regularity, which is defined in terms of local cohomology. In the first part we report some technical lemmata, which are useful in what follows as computational tools. Then we prove that the regularity of an ideal can be calculated using its graded Betti numbers, as it can be done in the standard case but with a correction due to the weights and the number of their occurrences. This is achieved in Theorem 3.5 and is essentially a paraphrase of what was performed in \[Be\] for pseudograded algebras. The section concludes with a result which predicts that Castelnuovo-Mumford regularity does not change when taking generic initial ideal (Proposition 3.6). Again, this is proved under the hypothesis that each weight is divisible by the previous ones, and is false in general, as shown in Example 3.7.
In the last section we deal with lexicographic ideals. In a standard graded polynomial ring lexicographic ideals are enhanced with many features and are well understood. As a consequence a certain amount of information about an arbitrary homogeneous ideal can be gathered by studying the associated lexicographic ideal. In other contexts though, the generalization of this notion turned out to be rather complex (cf. \[ADK\], \[DH\] and \[MP\]). In our situation there is a natural way to define lexicographic ideals, but it is difficult to describe them and to give a criterion to decide whether a homogeneous ideal in a given weighted polynomial ring is lexifiable, i.e. admits an associated lexicographic ideal with the same Hilbert function.
First we recollect some known results about Hilbert functions of a positively graded algebra and underline which are the factors which make our analysis difficult by means of some examples. In particular the shadow of a lexsegment is a lexsegment in the standard graded case, but this fact does not hold in general in the non-standard setting. We thus proceed by proving Proposition 4.9, which yields a method to verify whether a given ideal is lexicographic. It is indeed enough to check if finitely many graded components are generated by lexsegments. This is accomplished by means of an invariant $`G(w)`$, which was introduced in \[D\]. From a computational point of view the test is not optimal yet, since $`G(w)`$ not only depends on the weights but grows rapidly with the number of variables having the same weight.
Next, we are interested in describing the Hilbert function of lexicographic ideals, but to give a complete solution (cf. Problem 4.12) is an hard task. Theorem 4.11 provides an exhaustive answer for polynomial rings in two variables. The next topic we handle is expressed by Question 4.15. One would like to know which sets of weights make a polynomial ring Macaulay-Lex, i.e. such all of its ideals are lexifiable. Theorem 4.16 and the subsequent examples provide partial answers to this issue. We would like to observe that Theorem 4.11 and Theorem 4.16 provide a complete description of lexicographic ideals in two variables. Still, even in two variables it is not clear which ideals are lexifiable.
As a final remark, which completes and concludes this survey, we also mention the technique of polarization. In this setting admissible numerical functions need not to be Hilbert functions of lexicographic ideals. Thus completely polarized ideals (which in the standard case characterize lexicographic ideals), might be the right tool for a Theorem ร la Macaulay \[M\].
The calculations underlying many of the examples and the material of the last section were carried out using \[CoCoA\]. We implemented some procedures (for the computation of Hilbert functions, generic initial ideals, polarization and associated lexicographic ideals) which can be obtained by any of the two authors.
## Notation
In this paper we use some non-standard notation (!) which we illustrate here. When we consider polynomial rings with a non-standard grading, we mean that we work over an infinite field $`K`$ of characteristic $`0`$ and assume the degrees of the variables to be positive integers with no further restriction.
We order the variables by increasing degree or weight and often group together those with the same degree. Therefore we denote the polynomial ring by $`R=K[๐_\mathrm{๐},\mathrm{},๐_๐ง]`$, where $`๐_๐ข=(X_{i1},\mathrm{},X_{il_i})`$, $`\mathrm{deg}X_{ij}=q_i`$ for any $`j=1,\mathrm{},l_i`$, and $`q_1<q_2<\mathrm{}<q_n`$.
It is convenient to denote by $`R_{[i]}`$ the polynomial ring $`K[๐_\mathrm{๐},\mathrm{},๐_๐ข]`$. We let $`w`$ be the weight vector $`(\mathrm{deg}X_{11},\mathrm{},\mathrm{deg}X_{nl_n})`$ so that $`(R,w)`$ stands for a polynomial ring with the graduation given by $`w`$. If $`w`$ does not play an explicit role we denote $`(R,w)`$ simply by $`R`$.
If not elsewhere specified we consider term orderings $`>`$ which are degree compatible and assume $`X_{ij}>X_{ik}`$ if $`j<k`$, $`i=1,\mathrm{},n`$.
Since they are often used, it may be convenient to fix some notation for the total numbers of variables and the least common multiple of the weights. Thus we let $`l=_{i=1}^nl_i`$ and $`q=\mathrm{lcm}(q_1,\mathrm{},q_n)`$.
Finally, given a set $`AR_d`$, $`A`$ denotes the $`K`$-vector space spanned by $`A`$. If $`VR_d`$ is a $`K`$-vector space, $`\{V\}`$ denotes its monomial basis.
## 1 Weighted generic initial ideals
Generic initial ideals are monomial ideals attached to homogeneous ideals. It has been shown in many works that generic initial ideals, although simpler in nature, still contain a considerable amount of information about the original geometrical object. In order to extend their definition in our setting we have first of all to understand which are the graded automorphisms of $`R`$.
###### Proposition 1.1.
The assignment
$$\phi (X_{ij})=\underset{h=1}{\overset{l_i}{}}a_{jh}^iX_{ih}+\psi _{ij}(๐_\mathrm{๐},\mathrm{},๐_{๐ข\mathrm{๐}}),$$
where, for all $`i,j`$, $`\psi _{ij}`$ are homogeneous polynomials in $`R_{[i1]}`$ of degree $`q_i`$ and $`A_i(a_{jh}^i)_{j,h=1,\mathrm{},l_i}M_{l_i}(K)`$ are invertible matrices, defines a graded automorphism of $`R`$. Vice versa, any graded automorphism of $`R`$ is of this kind.
###### Proof.
Since $`q_i<q_j`$ if $`i<j`$, the requirement that $`\phi `$ is a graded homomorphism forces $`\psi _{ij}`$ to be polynomials in the first $`i1`$ sets of variables. Thus it is sufficient to show that $`\phi `$ is surjective if and only if $`A_i`$ are invertible for all $`i=1,\mathrm{},n`$. With some abuse of notation we write $`\phi (๐_๐ข)=A_i๐_๐ข+B_{i1}`$, where $`B_{i1}`$ is a $`l_i\times 1`$ matrix with entries in $`R_{[i1]}`$. If $`\phi `$ is surjective then for all $`i=1,\mathrm{},n`$ there exists a $`l_i\times 1`$ matrix $`C_i`$ with entries in $`R_{[i]}`$ such that $`๐_๐ข=\phi (C_i)`$. If we write $`C_i`$ as $`D_i๐_๐ข+E_{i1}`$, where $`D_iM_{l_i}(K)`$ and $`E_{i1}`$ is a $`l_i\times 1`$ matrix with entries in $`R_{[i1]}`$, we get that $`๐_๐ข=D_iA_i๐_๐ข+F_{i1}`$ where $`F_{i1}`$ consists of polynomials in $`๐_\mathrm{๐},\mathrm{},๐_{๐ข\mathrm{๐}}`$. Therefore $`F_{i1}=0`$ and $`D_iA_i=I`$. Vice versa, suppose that $`A_i`$ is invertible for $`i=1,\mathrm{},n`$. Since $`A_1`$ is invertible, $`\phi (A_1^1๐_\mathrm{๐})=๐_\mathrm{๐}`$. If $`i>1`$, we have $`\phi (A_i^1๐_๐ข)=๐_๐ข+A_i^1B_{i1}`$, where $`A_i^1B_{i1}`$ has entries in $`R_{[i1]}`$. Thus there exists a $`l_i\times 1`$ matrix $`C_{i1}`$ with entries in $`R_{[i1]}`$ such that $`\phi (C_{i1})=A_i^1B_{i1}`$ and $`๐_๐ข=\phi (A_i^1๐_๐ข)\phi (C_{i1})=\phi (A_i^1๐_๐ขC_{i1})`$, as required. โ
Let $`T`$ be the subset of upper triangular automorphisms consisting of those graded automorphisms of $`R`$ such that, for all $`i=1,\mathrm{},n`$, $`A_i`$ is an upper triangular invertible matrix. By the previous proposition it is clear that $`T`$ is a group, since upper triangular invertible matrices form a group (which is called Borel group). Let $`U`$ be the set of the elementary upper triangular automorphisms $`\tau _{ij}^{rc}`$, where $`r<j`$ and $`cK`$, determined by the assignment $`\tau _{ij}^{rc}(X_{ij})=X_{ij}+cX_{ir}`$ and $`\tau _{ij}^{rc}(X_{hk})=X_{hk}`$ if $`(h,k)(i,j)`$. Finally, let $`N`$ be the set of elementary non-linear automorphisms $`\eta _{ij}^m`$, where $`m`$ is a term of degree $`q_i`$ in $`R_{[i1]}`$ defined by $`\eta _{ij}^m(X_{ij})=X_{ij}+m`$ and $`\eta _{ij}^m(X_{hk})=X_{hk}`$ if $`(h,k)(i,j)`$.
###### Proposition 1.2.
$`T`$ is generated by the diagonal subgroup, by $`U`$ and by $`N`$.
###### Proof.
The proof is an easy induction on the number of variables $`l=l_1+\mathrm{}+l_n`$. If $`l=1`$, the only graded automorphisms are the diagonal automorphisms. Let $`\phi T`$, we say $`\phi (X_{ij})=_{h=1}^{l_i}a_{jh}^iX_{ih}+\psi _{ij}`$, and $`A_i(a_{jh}^i)_{j,h=1,\mathrm{},l_i}`$ are upper triangular invertible matrices; also, let $`\phi ^{}`$ be defined by $`\phi ^{}(X_{nl_n})=X_{nl_n}`$ and $`\phi ^{}(X_{ij})=\phi (X_{ij})`$ if $`(i,j)(n,l_n)`$. It is clear that $`\phi ^{}`$ is an automorphism and that therefore belongs to $`T`$. We now write the polynomial $`\psi _{nl_n}`$ as sum of monomials $`m_{nl_n}^hK[๐_\mathrm{๐},\mathrm{},๐_{๐ง\mathrm{๐}}]`$, $`h=1,\mathrm{},s`$, of degree $`q_n`$. We want to find a decomposition of $`\phi `$ by means of elementary non-linear, upper triangular, diagonal automorphisms and of $`\phi ^{}`$. This leads to the conclusion by induction, since $`\phi ^{}`$ fixes the last variable and can be thought of as an automorphism of a polynomial ring in $`l1`$ variables. We denote by $`\delta _{ij}^c`$, with $`cK`$, the diagonal automorphism defined by $`\delta _{ij}^c(X_{ij})=cX_{ij}`$ and $`\delta _{ij}^c(X_{hk})=X_{hk}`$ if $`(h,k)(i,j)`$. Moreover, since $`a_{l_nl_n}^n0`$, we may write $`u_hm_{nl_n}^h/a_{l_nl_n}^n`$ for all $`h=1,\mathrm{},s`$ and $`b_k\frac{a_{nk}^n}{a_{l_nl_n}^n}`$, for $`k=1,\mathrm{},l_n1`$. It is now easy to see that
$$\phi =\eta _{nl_n}^{u_1}\mathrm{}\eta _{nl_n}^{u_s}\tau _{nl_n}^{1b_1}\mathrm{}\tau _{nl_n}^{l_n1b_{l_n1}}\delta _{nl_n}^{a_{l_nl_n}^n}\phi ^{},$$
as desired. โ
### 1.1 Existence of the generic initial ideal
In the standard graded case the generic initial ideal $`\mathrm{Gin}(I)`$ of an ideal $`IK[X_1,\mathrm{},X_n]`$ plays a central role in problems regarding Hilbert functions and free resolutions of graded ideals. Since $`\mathrm{Gin}(I)`$ with respect to some assigned term order is defined as the initial ideal of $`gI`$, where $`g`$ is a generic change of coordinates, i.e. that $`g`$ is a matrix chosen out of a Zariski non-empty open set of $`\mathrm{GL}_n(K)`$, one way of computing it is the following. Let the $`n^2`$ entries of $`g`$ be new indeterminates, we say $`Y_{ij}`$, with $`i,j=1,\mathrm{},n`$. Write $`gI`$ explicitly and apply the Buchsbergerโs Algorithm to compute $`\mathrm{in}(I)`$ as an ideal of $`K(Y_{ij})[X_1,\mathrm{},X_n]`$. After finitely many computations of the so-called $`S`$-pairs the process finishes, the output result is the sought after monomial ideal - in the variables $`X_i`$ only - and the Zariski open set consists of all those matrices for which the finitely many polynomial denominators of the $`S`$-pairs are non-zero. If one considers this point of view, it is evident that weights do not play any role in the construction, which is thus also possible in the weighted case. Thus we can talk of generic initial ideals of homogeneous ideals in a non-standard graded algebra.
In the standard case it is well-known that generic initial ideals are Borel-fixed, i.e. fixed under the action of the Borel subgroup of $`\mathrm{GL}_n(K)`$ consisting of the upper triangular invertible matrices.
###### Theorem 1.3.
Let $`I`$ be a homogeneous ideal of a weighted polynomial algebra $`R`$. Then $`\mathrm{Gin}(I)`$ is $`T`$-fixed, i.e. $`\phi (\mathrm{Gin}(I))=\mathrm{Gin}(I)`$ for all $`\phi T`$.
###### Proof.
We only need to observe that by applying a non-linear automorphism to a monomial $`u`$ one obtains a polynomial of the form $`u+v`$ where $`v`$ is bigger than $`u`$ in the chosen term order. By Proposition 1.2, this is enough to argue as in the standard case, see for instance the proof of Theorem 15.20 in \[E\]. โ
### 1.2 Weighted strongly stable ideals
Generic initial ideals in a standard graded polynomial ring are characterized combinatorially, the simplicity of this description depending on the characteristic of the base field. In a weighted polynomial ring over a base field of characteristic $`0`$ the same can be performed, via the following definition.
###### Definition 1.4.
Let $`I`$ be a monomial ideal. $`I`$ is called (strongly) stable if the following holds: for every $`uI`$, if $`X_{ij}u`$ then $`\frac{X_{ih}u}{X_{ij}}I`$, for every $`h<j`$ and $`\frac{vu}{X_{ij}}I`$ for all monomials $`v`$ of degree $`q_i`$ in $`R_{[i1]}`$.
It is not difficult to prove that weighted generic initial ideals are stable according to this definition.
###### Proposition 1.5.
Let $`I`$ be an homogeneous ideal. $`I`$ is $`T`$-fixed if and only if $`I`$ is strongly stable.
###### Proof.
One begins by observing that $`I`$ is fixed by the subgroup of diagonal matrices if and only if $`I`$ is monomial. Let $`m=X_{ij}^tm^{}`$, where $`X_{ij}m^{}`$. The images $`\tau _{ij}^{rc}(m)`$, $`\eta _{ij}^s(m)`$, with $`\mathrm{deg}s=q_i`$ can be written as $`(X_{ij}+cX_{ir})^tm^{}`$ and $`(X_{ij}+s)^tm^{}`$ respectively. If $`I`$ is $`T`$-fixed both polynomials, and so each of their monomials, belong to $`I`$. In particular the conditions which define strongly stable ideals are verified. Conversely, if $`I`$ is strongly stable the same argument shows that $`I`$ is fixed by the action of the generators of $`U`$ and $`N`$, and is therefore $`T`$-fixed. โ
One of the key properties of strongly stable ideals in the standard graded polynomial ring $`K[X_1,\mathrm{},X_n]`$ is that $`I:(X_1,\mathrm{},X_n)=I:X_n`$. In fact, beside the trivial inclusion $`I:(X_1,\mathrm{},X_n)I:X_n`$ one has that $`mI:x_n`$ iff $`mX_nI`$, which, because of the stability property, implies $`mX_iI`$ for all $`i`$. It is quite clear how this property is weakened in the more general case where variables might have different weights. In particular the good property of stable ideals with respect to taking colons with the last variables plays a central role in the construction of the Eliahou-Kervaire resolution \[EK\] of such an ideal. This is a completely described minimal graded free resolution of such an ideal in terms of its minimal set of monomial generators. On the other hand being able to construct such a resolution having no restriction on the weight vector would mean to know how to describe a minimal resolution of any monomial ideal, since given any such ideal $`I`$, one can choose weights so that in the corresponding polynomial ring $`I`$ is stable, as the next example shows.
###### Example 1.6.
Let $`A`$ be a set of monomials in $`n`$ variables. Then, there exist non-negative integers $`q_1,\mathrm{},q_n`$ such that in the weighted polynomial ring $`(K[X_1,\mathrm{},X_n],(q_1,\mathrm{},q_n))`$ the ideal generated by $`A`$ is strongly stable. In fact, one can choose weights in such a way that none of the exchanges which were described in Definition 1.4 is possible. For instance, it is enough to choose $`q_1<q_2<\mathrm{}<q_n`$ so that $`2q_1>q_n`$, we say $`q_1=n+1,q_2=n+2,\mathrm{},q_n=2n`$.
## 2 Prime Avoidance
A simple fact of linear algebra gives rise to a powerful tool when combined with techniques dealing with generic forms. This is known as Homogeneous Prime Avoidance: If $`๐ญ_1,\mathrm{},๐ญ_n`$ are prime ideals strictly contained in the graded maximal ideal of a standard graded algebra over an infinite field then there exists a homogeneous form of degree $`1`$ in $`๐ช_i๐ญ_i`$. It turns out to be essential in many proofs, since avoiding a finite number of primes is an open property.
###### Lemma 2.1 (Weighted Prime Avoidance).
Let $`q=\mathrm{lcm}(q_1,\mathrm{},q_n)`$ and let $`๐ญ_1,\mathrm{},๐ญ_n`$ be prime ideals with $`๐ญ_i๐ช`$. Then $`๐ช_q_i(๐ญ_i)_q\mathrm{}`$.
###### Proof.
Since the prime ideals are strictly contained in the maximal ideal, we have that $`(๐ญ_i)_q๐ช_q`$ for all $`i`$. Else, one would have that $`(๐ญ_i)_q=๐ช_q`$ and $`X_{jk}๐ญ_i`$, for all $`j`$ and $`k`$, since $`X_{jk}^{q/q_j}๐ช_q`$ and $`๐ญ`$ is prime. But the infinite vector space $`๐ช_q`$ cannot be written as a finite union of proper subspaces $`๐ญ_q`$, and the claim follows. โ
The next example shows that in general it is not possible to find such a form in a smaller degree.
###### Example 2.2.
Let $`(R,w)=(K[X,Y],(2,3))`$. If $`๐ญ_1=(X)`$ and $`๐ญ_2=(Y)`$ then the smallest degree $`d`$ such that $`(X,Y)_d(๐ญ_1)_d(๐ญ_2)_d`$ is $`6`$.
In the following we recover some results which are known in the standard case, provided that some condition on the weights is assumed. It may be convenient to state one of these conditions here.
###### Condition 2.3.
$`(R,w)`$ is a weighted polynomial ring with $`q_iq_{i+1}`$ for $`i=1,\mathrm{},n1`$.
###### Lemma 2.4.
Let $`(R,w)`$ be a ring for which Condition 2.3 is satisfied, and let $`IR`$ be a strongly stable ideal. For any $`i=1,\mathrm{},n`$ and $`j=1,\mathrm{},l_i`$, one has
$$I:X_{ij}^{\mathrm{}}=I:(X_{11},\mathrm{},X_{ij})^{\mathrm{}}.$$
###### Proof.
We only have to prove the inclusion $``$ since the other one is obvious. Let $`m`$ be a monomial such that $`mX_{ij}^sI`$ for some $`s`$. Since $`I`$ is strongly stable, $`mX_{ih}^sI`$ for any $`1hj`$; furthermore the assumption on the degrees of the indeterminates implies that $`\mathrm{deg}X_{ij}=q_i=q_h\frac{q_i}{q_h}=\mathrm{deg}X_{hk}^{q_i/q_h}`$, and consequently $`m(X_{hk}^{q_i/q_h})^sI`$ for any $`1hi1`$ and $`1kl_h`$, as desired. โ
As a consequence we obtain the following proposition.
###### Proposition 2.5.
Let $`(R,w)`$ be a ring for which Condition 2.3 is satisfied, let $`I`$ be a strongly stable ideal and let $`X_{hk}`$ be the (lex-)smallest variable which divides some minimal generator of $`I`$. Then $`X_{hk+1},\mathrm{},X_{nl_n}`$ form a maximal regular sequence on $`R/I`$.
###### Proof.
Clearly the elements $`X_{hk+1},\mathrm{},X_{nl_n}`$ form a regular sequence on $`R/I`$. Since in the quotient ring $`\overline{R}=K[X_{11},\mathrm{},X_{hk}]`$ the ideal $`\overline{I}`$ is strongly stable, by the previous lemma $`\overline{I}^{sat}=\overline{I}:X_{hk}^{\mathrm{}}\overline{I}`$, which implies that $`0pt\overline{R}/\overline{I}=0`$. โ
We recall the following theorem \[BS\], which holds independently of the given weights and is needed for the proof of the final result of this section.
###### Theorem 2.6.
Let $`F`$ be a free $`R`$module with basis and consider the degree reverse lexicographic monomial order. Let $`M`$ be a graded submodule of $`F`$. The elements $`X_{nl_n},X_{nl_n1},\mathrm{},X_{ij+1},X_{ij}`$ form a regular sequence on $`F/M`$ if and only if they form a regular sequence on $`F/\mathrm{in}(M)`$.
###### Proof.
See that of Theorem 15.13 in \[E\]. โ
###### Theorem 2.7.
Let $`(R,w)`$ be a ring for which Condition 2.3 is satisfied, and consider the degree reverse lexicographic order. Then, for any homogeneous ideal $`IR`$,
$$0ptR/I=0ptR/\mathrm{Gin}(I).$$
###### Proof.
Since $`0ptR/I0ptR/\mathrm{Gin}(I)`$, we may assume $`0ptR/I>0`$. By Lemma 2.1, a generic form of degree $`q_n`$ is a non-zerodivisor on $`R/I`$. Thus, after a generic change of coordinates, we may assume that $`X_{nl_n},X_{nl_n1},\mathrm{},X_{ij}`$ is a maximal $`R/I`$-regular sequence and $`\mathrm{Gin}(I)=\mathrm{in}(I)`$. By Theorem 1.3 and Proposition 1.5 $`\mathrm{Gin}(I)`$ is strongly stable, and consequently, by Proposition 2.5, there is a maximal $`R/\mathrm{Gin}(I)`$-regular sequence $`X_{nl_n},X_{nl_n1},\mathrm{},X_{hk}`$. Now Theorem 2.6 yields that $`(h,k)=(i,j)`$, from which the conclusion is straightforward. โ
## 3 Regularity
Local cohomology modules of a graded module over a weighted polynomial ring have a graded structure arising from resolutions by graded injective modules or equivalently from the construction of the ฤech complex. The usual definition of Castelnuovo-Mumford regularity by means of local cohomology still works in this context and we recall it here. Let $`H_๐ช^i(M)`$ denote the $`i^{th}`$ graded local cohomology module of the graded $`R`$-module $`M`$ with support on the graded maximal ideal $`๐ช`$.
###### Definition 3.1.
Let $`R`$ be a weighted polynomial ring with graded maximal ideal $`๐ช`$. We let
$$a^i(M)\{\begin{array}{cc}\mathrm{max}\{j:H_๐ช^i(M)_j0\}\hfill & \text{if }H_๐ช^i(M)0\hfill \\ \mathrm{}\hfill & \text{otherwise}\hfill \end{array}$$
denote the end of the $`i^{th}`$ local cohomology module of $`M`$. The Castelnuovo-Mumford regularity of $`M`$ is then $`\mathrm{reg}M=\mathrm{max}_{1idimM}\{a^i(M)+i\}`$.
However, one of the aspects that made the Castelnuovo-Mumford regularity interesting, i.e. its direct interpretation through the Betti numbers of the minimal free resolution by means of the formula
$$\mathrm{reg}M=\underset{i0}{\mathrm{max}}\{b_i(M)i\},$$
(3.1)
where $`b_i(M)\mathrm{max}_j\{\beta _{ij}(M)0\}`$, fails in the general weighted case. In this section we re-prove some results about regularity which still hold in the weighted case, in order to give in Theorem 3.5 a formula that generalizes (3.1). In the last part, we consider the regularity of a generic initial ideal $`\mathrm{Gin}(I)`$ and prove in Proposition 3.6 that under some assumption on weights it does not differ from that of $`I`$. Also, we provide a counterexample that shows that in general there is no analogue of the well-known theorem \[BS\] valid in the standard case.
We start by recalling some lemmata which are useful in order to control regularity in the non-standard case.
###### Lemma 3.2.
Let $`0NMQ0`$ be a short exact sequence of finitely generated graded $`R`$-modules. Then
* $`\mathrm{reg}N\mathrm{max}\{\mathrm{reg}M,\mathrm{reg}Q+1\}`$;
* $`\mathrm{reg}M\mathrm{max}\{\mathrm{reg}N,\mathrm{reg}Q\}`$;
* $`\mathrm{reg}Q\mathrm{max}\{\mathrm{reg}N1,\mathrm{reg}M\}`$;
* If $`N`$ has finite length, then $`\mathrm{reg}M=\mathrm{max}\{\mathrm{reg}N,\mathrm{reg}Q\}`$.
###### Proof.
The proofs of $`(i)(iii)`$ are easy and descend from the use of the long exact sequence in cohomology $`\mathrm{}H_๐ช^{i1}(Q)H_๐ช^i(N)H_๐ช^i(M)\mathrm{}`$.
As for the proof of $`(iv)`$, it is clear that $`\mathrm{reg}N=a^0(N)`$ and $`a^0(M)=\mathrm{max}\{a^0(N),a^0(Q)\}`$. Thus,
$$\begin{array}{cc}\hfill \mathrm{reg}M& \mathrm{max}\{a^0(M),\underset{i>0}{\mathrm{max}}\{a^i(M)+i\}\}\hfill \\ \hfill =& \mathrm{max}\{a^0(N),a^0(Q),\underset{i>0}{\mathrm{max}}\{a^i(Q)+i\}\},\hfill \end{array}$$
as desired. โ
###### Lemma 3.3.
Let $`M`$ be a finitely generated graded $`R`$-module and let $`xR_d`$. If $`x`$ is a non-zerodivisor on $`M`$ then $`\mathrm{reg}M/xM=\mathrm{reg}M+(d1)`$. More generally, if $`x`$ is such that $`(0:_Mx)`$ has finite length, then
$$\mathrm{reg}M=\mathrm{max}\{\mathrm{reg}0:_Mx,\mathrm{reg}M/xM(d1)\}.$$
###### Proof.
From the exact sequence $`0(0:_Mx)(d)M(d)MM/xM0`$ one obtains the two short exact sequences $`0(0:_Mx)(d)M(d)xM0`$ and $`0xMMM/xM0`$ so that the proof follows easily as an application of Lemma 3.2. โ
###### Lemma 3.4.
Let $`xR_d`$ such that $`0:_Mx`$ is of finite length. Then for all $`i0`$
$$a_๐ช^{i+1}(M)+da_๐ช^i(M/xM)\mathrm{max}\{a_๐ช^i(M),a_๐ช^{i+1}(M)+d\}.$$
###### Proof.
From the two short exact sequences contained in the proof of the last lemma we deduce that $`H_๐ช^i(M(d))H_๐ช^i(xM)`$ for all $`i>0`$ and obtain the long exact sequence in cohomology $`\mathrm{}H_๐ช^i(M)H_๐ช^i(M/xM)H_๐ช^{i+1}(xM)H_๐ช^{i+1}(M)\mathrm{}`$. To prove the first inequality, it is enough to observe that, if $`a_๐ช^i(M/xM)<a_๐ช^{i+1}(M)+d`$ the above long exact sequence in degree $`a_๐ช^{i+1}(M)+d`$ would deliver a contradiction. The proof of the second inequality is analogous. โ
Let $`M`$ be a finitely generated $`R`$-module of finite projective dimension $`s`$. For $`i=1,\mathrm{},s`$ let as before $`b_i(M)\mathrm{max}_j\{\beta _{ij}(M)0\}`$.
###### Theorem 3.5.
Let $`R=K[๐_\mathrm{๐},\mathrm{},๐_๐ง]`$ be a graded polynomial ring and let $`M`$ be a finitely generated $`R`$-module with $`\mathrm{proj}\mathrm{dim}M<\mathrm{}`$. Then
$$\mathrm{reg}M=\underset{i0}{\mathrm{max}}\{a_๐ช^i(M)+i\}=\underset{i0}{\mathrm{max}}\{b_i(M)i\}\underset{j=1}{\overset{n}{}}l_j(q_j1).$$
###### Proof.
By virtue of the previous lemma and induction on the number of variables one first proves that
$$b_0(M)\underset{i0}{\mathrm{max}}\{a_๐ช^i(M)+i\}+\underset{j=1}{\overset{n}{}}l_j(q_j1).$$
(3.2)
Moreover, it is easy to verify that, if $`F`$ is a free $`R`$-module,
$$a_๐ช^n(F)=b_0(F)\underset{i=1}{\overset{n}{}}l_iq_i.$$
(3.3)
The assertion follows by the use of (3.2) and (3.3) combined with an induction argument on the projective dimension of $`M`$. The proof is an adaptation of that of Theorem 5.5 in \[Be\] (to which the interested reader is referred) and, therefore, the details are omitted here. โ
###### Proposition 3.6.
Let $`(R,w)`$ be a weighted polynomial ring for which Condition 2.3 is satisfied and consider the degree reverse lexicographic order. If $`I`$ an homogeneous ideal of $`R`$ then
$$\mathrm{reg}R/I=\mathrm{reg}R/\mathrm{Gin}(I).$$
###### Proof.
Since $`q_n=\mathrm{lcm}(q_1,\mathrm{},q_n)`$, by Lemma 2.1 a generic form in $`๐ช_{q_n}`$ does not belong to any associated prime $`๐ญ๐ช`$ of $`I`$. By applying a generic automorphism we may assume that $`\mathrm{Gin}(I)=\mathrm{in}(I)`$ and that $`X_{nl_n}`$ is almost-regular. Therefore $`(I:X_{nl_n})/I`$ and, consequently, $`\mathrm{in}(I:X_{nl_n})/\mathrm{in}(I)(\mathrm{in}(I):X_{nl_n})/\mathrm{in}(I)`$ have finite length. By Lemma 3.3, it is enough to verify that $`\mathrm{reg}(I:X_{nl_n})/I=\mathrm{reg}(\mathrm{in}(I):X_{nl_n})/\mathrm{in}(I)`$ in order to apply induction on the numbers of the variables, since the assumption on the weights still holds for $`R/(X_{nl_n})`$. But this is clear because the above modules coincide with their $`0^{th}`$ local cohomology module and they have the same Hilbert function. โ
Notice that the assumption on the weights is essential in order to have a generic form of the right degree for the induction. The following example shows that the above result cannot be extended for any choice of weights.
###### Example 3.7.
Let $`(R,w)=(K[X,Y,Z],(2,4,5))`$ and $`I=(XY,YZ,X^5)`$ and consider the degree reverse lexicographic order. A computation with \[CoCoA\] shows that $`\mathrm{Gin}(I)=(X^3,X^2Z,XY^2,Y^3Z)`$. By Theorem 3.5, the regularity of $`I`$ and $`\mathrm{Gin}(I)`$ can be computed by the use of the resolutions
$$0R(14)R(11)R(10)R(9)R(6)I0$$
and
$$\begin{array}{cc}\hfill 0R(19)& R(19)R(17)R(14)R(11)\hfill \\ & R(17)R(10)R(9)R(6)\mathrm{Gin}(I)0\hfill \end{array}$$
of $`I`$ and $`\mathrm{Gin}(I)`$ respectively.
This example points out that also Proposition 2.7 is not valid without Condition 2.3.
## 4 Lexicographic ideals
Although the definition and some of the main properties of Hilbert functions are still valid in a non-standard setting, a great deal is still unknown about them. In particular Macaulayโs Theorem, which provides a necessary and sufficient condition for a numerical function to be the Hilbert function of a finitely generated standard graded algebra has no counterpart in the weighted case. The main tool which is involved in this context, lexicographic ideals, can be easily defined in the non-standard case, but they are not so easily investigated, as the following analysis shows.
### 4.1 Hilbert functions
Here we point out some facts about non-standard graded algebras which are relevant for our purposes. We start by recalling the well-known Hilbert-Serre Theorem: Let $`I`$ be a homogeneous ideal in $`(R,w)`$. The Poincare series $`P(R/I,t)`$ of $`R/I`$ is a rational function in $`t`$ of the form $`g(t)/_{i=1}^n(1t^{q_i})^{l_i}`$, where $`g(t)[t]`$. It is known that the Hilbert function of $`R/I`$ is quasi-polynomial. Some more information is provided by the following result to be found in \[B\], Theorem 2.2.
###### Proposition 4.1.
Let $`I`$ be a homogeneous ideal in $`(R,w)`$ and let $`d`$ be the order of the pole of $`P(R/I,t)`$ at the point $`t=1`$. Then there exist $`q=\mathrm{lcm}(q_1,\mathrm{},q_n)`$ polynomials $`p_0,\mathrm{},p_{q1}[t]`$ of degree at most $`d1`$ with coefficients in $`[q^{d1}(d1)!]^1`$ such that, for all $`l0`$,
$$H_{R/I}(l)=p_j(l)\text{ for }lj\mathrm{mod}q$$
It is also worth observing that in general some of the Hilbert polynomials described in the above proposition can be $`0`$. Also in the case $`\mathrm{gcd}(q_1,\mathrm{},q_n)=1`$ the vanishing of the Hilbert function of $`R/I`$ in $`t=t_0`$ does not imply that $`H_{R/I}(t)=0`$ for all $`t>t_0`$. However, this is true for $`H_R(t)`$ if $`t_0`$ is bigger than the Frobenius number of $`q_1,\mathrm{},q_n`$ (cf. \[SS\], Chapter 1, Section 3 for more details about this subject).
###### Remark 4.2.
Let $`(R,w)`$ be a weighted polynomial ring. If $`w_i=q`$ for all $`i`$, then the Hilbert function $`H_{(R,w)}(t)`$ is equal to $`H_{(R,(1,\mathrm{},1))}(t/q)`$ if $`qt`$ and $`0`$ otherwise; this case is thus essentially equivalent to the standard case. The same observation shows that one may assume without loss of generality that the $`\mathrm{gcd}`$ of the weights is $`1`$.
Another pathology of the weighted case is shown in the following example \[BR\].
###### Example 4.3.
Let $`(R,w)=(K[X,Y,Z,T],(1,6,10,15))`$.
The monomial $`XY^4Z^2T`$ has degree $`60`$, but it is not multiple of any monomial of degree $`30`$.
However, one can show that this can only occur in low degrees, as it is shown in \[BR\], Proposition 4B.5. One makes use of an invariant introduced in \[D\], which we denote by $`G(w)`$. For the readerโs sake we recall here the result, omitting the definition of $`G(w)`$ since it is not essential in what follows.
###### Proposition 4.4.
Let $`(R,w)`$ be a weighted polynomial ring and let $`n>G(w)`$. Then every monomial of $`R_{n+hq}`$ is divisible by a monomial in $`R_{hq}`$, for any $`h`$.
One might wonder if the same holds for an arbitrary ideal generated in more than one degree, i.e. if there exists $`l`$ such that for all $`r0`$ one has $`I_r=I_lR_{rl}`$. Unfortunately this is false and it partially explains why the study of lexicographic ideals is complicated.
###### Example 4.5.
Let $`(R,w)=(K[X,Y,Z],(2,2,3))`$ and $`I=(X^\alpha ,XYZ)`$ for some integer $`\alpha >1`$. Let us suppose that there exists some $`l`$ such that $`I_r=I_lR_{rl}`$ for all $`r0`$. Then, for all $`k0`$, $`X^kI`$ and this implies that $`l`$ is even. On the other hand $`XY^kZI`$ for all $`k0`$; thus there exists $`k_0`$ such that $`XY^{k_0}ZI_l`$ and $`l=2+2k_0+3`$, which means that $`l`$ is odd.
### 4.2 Lexifiable ideals
Let us consider a standard graded polynomial ring $`K[X_1,\mathrm{},X_n]`$ with the degree lexicographic order. We recall that a lexsegment (of degree $`d`$) is a set $`L`$ of monomials of degree $`d`$ with the property: if $`uL`$ and $`v>u`$ with $`\mathrm{deg}v=\mathrm{deg}u`$ then $`vL`$. A homogeneous monomial ideal is said to be lexicographic if all its graded components are spanned as a $`K`$-vector space by lexsegments. It is clear that these definitions can be overtaken and used in the weighted case.
Lexicographic ideals play a central role in many results in commutative algebra because of their well understood structure. It would be of great interest to grasp which of their properties also hold in the weighted case. One of the most important facts concerning a lexsegment $`L`$ is that the so-called shadow, i.e. the set of monomials which are obtained by multiplying the monomials of $`L`$ by all the variables, is still a lexsegment. In the weighted case, given a lexsegment $`L`$ it is natural to consider the $`n`$ shadows $`L๐_\mathrm{๐},\mathrm{},L๐_๐ง`$. In general they are not lexsegments, as the following easy example shows.
###### Example 4.6.
Let $`(R,w)=(K[X,Y],(2,3))`$, and consider the monomial $`XY`$. This is the only monomial of degree $`5`$, and $`\{XY\}`$ is obviously a lexsegment. Its shadow in degree $`8`$ is $`\{XY^2\}`$, which is not lexsegment, since $`X^4`$ does not belong to it.
It is thus quite clear that there are strong restrictions also on Hilbert functions of lexsegment ideals generated in one degree. For instance, an ideal $`I`$ generated in degree $`d`$ is a lexicographic ideal only if contains $`X_{11}^k`$ for some $`k>d`$; this is possible only if $`d=\alpha q_1`$.
###### Example 4.7.
Let $`IR=K[๐_\mathrm{๐},\mathrm{},๐_๐ง]`$ with $`n>1`$ be an ideal generated by a lexsegment in degree $`d=\alpha q_1`$ and $`H_I(d)l_1`$. Then $`I`$ is a lexicographic ideal.
In fact $`I`$ is generated by the lexsegment $`\{X_{11}^d,X_{11}^{d1}X_{12},\mathrm{},X_{11}^{d1}X_{1h}\}`$, with $`hl_1`$; if we let $`mI_r`$ be a monomial with $`rd`$ then, since $`m`$ belongs to $`I`$, $`m=X_{11}^{d1}X_{1j}m^{}`$ for some $`j=1,\mathrm{},h`$ and a monomial $`m^{}R_{rd}`$. If $`s`$ is a monomial with $`\mathrm{deg}s=\mathrm{deg}m`$ and $`sm`$, then $`s`$ must be $`X_{11}^{d1}X_{1j^{}}s^{}`$ with $`j^{}<j`$ or $`j^{}=j`$ and $`s^{}m^{}`$. It is thus clear that $`sI_r`$.
###### Example 4.8.
Let $`(R,w)=(K[X,Y,Z],(2,2,3))`$. By the previous example, the ideal $`(X^3,X^2Y)`$ is lexicographic, whereas $`(X^3,X^2Y,XY^2)`$ is not. However the ideal $`(X^3,X^2Y,XY^2,X^2Z^2)`$ is lexicographic, as an easy verification shows.
At this point it is still not evident if it is possible to determine whether an ideal is lexicographic by looking at finitely many of its graded components. The following proposition yields that this is indeed the case.
###### Proposition 4.9.
Let $`I(R,w)`$ be a homogeneous ideal generated in degree $`d`$ and let $`q\mathrm{lcm}(q_1,\mathrm{},q_n).`$ If $`I_i`$ is spanned (as a $`K`$-vector space) by a lexsegment for all $`id+q+G(w)`$, then $`I`$ is a lexicographic ideal.
###### Proof.
The proof is by induction on $`i`$. We only have to prove that $`I_i`$ is spanned by a lexsegment if $`i>d+q+G(w)`$, provided that this is true for $`I_r`$ with $`r<i`$. If $`I_i=\mathrm{}`$ there is nothing to prove. Else, let $`v_i`$ be the smallest monomial in $`I_i`$ and let $`u>v_i`$, with $`\mathrm{deg}u=\mathrm{deg}v_i`$, be a monomial not in $`I`$. Finally let $`X_{jh}`$ denote the (lex-)smallest variable which divides $`u`$. Now we write $`v_i`$ as $`vm`$, where $`v`$ is a minimal generator of $`I`$, we say of degree $`d^{}d`$, and $`m`$ is the smallest monomial in $`R_{id^{}}`$. Since $`id^{}>q+G(w)`$, by Proposition 4.4, we may write $`m=m^{}m^{\prime \prime }`$, where $`m^{}`$ is the smallest monomial of $`R_q`$, which is $`X_{nl_n}^{q/q_n}`$. Thus, $`v_i=vX_{nl_n}^{q/q_n}m^{\prime \prime }`$. If we now let $`w\frac{v_iX_{jh}^{q/q_j}}{X_{nl_n}^{q/q_n}}`$, it is clear that $`uwv_i`$, $`wI_i`$ and $`w/X_{jh}I_{rq_j}`$. But this is a contradiction, since $`u/X_{jh}w/X_{jh}`$ and $`I_{rq_j}`$ is spanned by a lexsegment. โ
In a polynomial ring with two variables and coprime weights, one can expect to have a description of lexicographic ideals, because of the following observation. Given a polynomial ring $`(R,w)`$ according to \[CL\] we call any set of consecutive monomials of the same degree a block. If $`R`$ is a polynomial ring in two variables, any shadow of a block is clearly a block. With this notation, a lexsegment is an initial block, and Example 4.6 shows that in general the shadow of an initial block needs not to be such.
Before proceeding with the characterization of lexicographic ideals of $`K[X,Y]`$, we need to fix some notation. Given any set $`A`$ in $`R`$ of monomials of degree $`d`$, we let $`\mathrm{Shad}_i(A)R_{d+i}`$ denote the set of the elements $`um`$, where $`uA`$ and $`m`$ is a monomial in $`R_i`$. Clearly, the cardinality of $`\mathrm{Shad}_{q_1}(A)`$ equals that of $`A`$. Moreover, $`|\mathrm{Shad}_{q_2}(A)|=|A|`$ if $`q_11`$ and $`|\mathrm{Shad}_{q_2}(A)|=|A|+1`$ if $`q_1=1`$ and $`A`$ is a block.
Finally, if $`d`$ and $`q_1d`$, we let $`\delta `$ denote the smallest integer $`d+\beta q_2`$, $`\beta `$, divisible by $`q_1`$. It is not difficult to see that such a number exists and it is such that $`\beta <q_1`$ since $`q_1`$ and $`q_2`$ are assumed to be coprime.
###### Lemma 4.10.
Let $`L`$ be a lexsegment of degree $`d`$.
* If $`q_1d`$ then $`Shad_i(L)`$ is a lexsegment for all $`i`$.
* If $`q_1d`$ then $`\{X^{\delta /q_1}\}\mathrm{Shad}_{\delta d}(L)`$ is a lexsegment (of degree $`\delta `$).
###### Proof.
The proof of $`(i)`$ is obvious.
To prove $`(ii)`$, let $`X^aY^b`$ be the largest monomial of $`L`$ and $`\delta =d+\beta q_2`$. First observe that $`b<q_1`$. Secondly, notice that the largest monomial of $`\mathrm{Shad}_{\delta d}(L)`$ is $`X^aY^{b+\beta }`$. Since $`b+\beta <2q_1`$ and is a multiple of $`q_1`$, we have $`b+\beta =q_1`$, so that the only monomial of $`R_\delta `$ which is larger is $`X^{\delta /q_1}`$. โ
###### Theorem 4.11.
Let $`I`$ be a monomial ideal minimally generated in degrees $`d_1<d_2<\mathrm{}<d_r`$ such that the monomials of $`I_{d_i}`$ form a lexsegment for all $`i=1,\mathrm{},r`$. Then $`I`$ is a lexicographic ideal if and only if $`q_1d_1`$ or $`q_1d_1`$ and there exists $`1<sr`$ such that $`q_1d_s`$ and $`d_s\mathrm{min}\{\delta _i\}`$.
###### Proof.
We start by proving that if $`q_1d_1`$ then $`I`$ is a lexicographic ideal. For this purpose it is enough to observe that the shadow of a block is a block and use Lemma 4.10.
We show now that if $`q_1d_1`$, the conditions on $`d_s`$ imply that $`I`$ is lexicographic. In fact it is sufficient to verify that $`\mathrm{Shad}_i(\{I_{d_j}\})`$ are lexsegments for all $`i`$ and $`j=1,\mathrm{},r`$. Since the generator $`X^{d_s/q_1}`$ occurs in degree $`\mathrm{min}\{\delta _i\}`$, the conclusion follows again by Lemma 4.10.
Finally, if $`I`$ is lexicographic and $`q_1d_1`$, then $`X^kI_{kq_1}`$ with $`kq_1=d_s`$ for some $`1<sr`$. By Lemma 4.10 $`(ii)`$ it is thus clear that $`d_s\mathrm{min}\{\delta _i\}`$. โ
The conditions in the previous proposition can be easily re-formulated in terms of Hilbert series. In general, it would be interesting to have a solution for the following problem.
###### Problem 4.12.
Find a combinatorial characterization for the Hilbert series of lexicographic ideals.
In the same spirit of \[MP\], we say that an ideal $`I(R,w)`$ is lexifiable if there exists a lexicographic ideal $`L`$ with the same Hilbert function as $`I`$.
Given a subset $`A`$ of monomials in $`R_i`$, $`i`$, we let $`\mathrm{Lex}(A)`$ denote the set of the $`|A|`$ lexicographic largest elements of $`R_i`$. We also let $`L_i\mathrm{Lex}(\{I_i\})`$. Thus, $`I`$ is lexifiable iff $`L`$ is an ideal of $`(R,w)`$. To establish which ideals are lexifiable is not an easy task. The following example shows an ideal which is not lexifiable in any lex-order.
###### Example 4.13.
Let $`(R,w)=(K[X,Y],(2,3))`$. The ideal $`(Y)`$ provides an easy example of an ideal which is lexifiable if $`Y>_{\mathrm{Lex}}X`$ and not lexifiable if $`X>_{\mathrm{Lex}}Y`$.
Let $`I=(X^3Y^3,X^2Y^4)`$. $`I`$ is not lexifiable in both cases $`X>_{\mathrm{Lex}}Y`$ and $`Y>_{\mathrm{Lex}}X`$. If $`X>_{\mathrm{Lex}}Y`$ then the candidate to be the associated lexicographic ideal with $`I`$ is the ideal $`L=(X^8,X^6Y)`$ but $`H_I(18)=1`$ and $`H_L(18)=2`$. If $`Y>_{\mathrm{Lex}}X`$ the candidate is $`L=(Y^5,Y^4X^2)`$, but again $`H_L(18)=2`$.
###### Example 4.14.
Let $`(R,w)=(K[X,Y],(2,7))`$. The monomials of degree $`28`$ and $`35`$ are $`X^{14},X^7Y^2,Y^4`$ and $`X^{14}Y,X^7Y^3,Y^5`$ respectively. Let us consider the ideals which have exactly one minimal generator in these two degrees. These are $`I_1=(X^{14},X^7Y^3)`$, $`I_2=(X^7Y^2,Y^5)`$, $`I_3=(X^{14},Y^5)`$, $`I_4=(X^7Y^2,X^{14}Y)`$, $`I_5=(Y^4,X^7Y^3)`$ and $`I_6=(Y^4,X^{14}Y)`$. According to our definitions $`I_1`$ is a lexicographic ideal, $`I_2`$ is lexifiable and $`I_1`$ is the lexicographic ideal associated with it, $`I_3`$ is lexifiable associated with $`(X^{14},X^7Y^3,Y^7)`$. The ideals $`I_4`$, $`I_5`$ and $`I_6`$ are not lexifiable, as a computation of the Hilbert function in degree $`42`$, $`42`$ and $`56`$ shows.
Again according to \[MP\], we say that a graded polynomial ring $`(R,w)`$ is Macaulay-Lex if every homogeneous ideal in $`(R,w)`$ is lexifiable. Macaulayโs Theorem together with Remark 4.2 says that $`(R,w)`$ with $`w=(a,\mathrm{},a)`$ is Macaulay-Lex, whereas for a general choice of $`w`$ there are many ideals which are not lexifiable. Thus it is natural to ask the following question.
###### Question 4.15.
Which polynomial rings $`(R,w)`$ are Macaulay-Lex?
The results of the last part of this section shed some light on the problem and provide partial answers to the above question.
###### Theorem 4.16.
Let $`I`$ be a homogeneous ideal in $`(R,w)=(K[X,Y],(1,q_2))`$. There exists a unique lexicographic ideal $`L`$ such that $`H_{R/I}(t)=H_{R/L}(t)`$ for any $`t`$.
###### Proof.
Taking in consideration what we have said before Example 4.13, it is sufficient to prove that
$$\mathrm{Shad}_1(\mathrm{Lex}(\{I_d\}))\mathrm{Lex}(\{I_{d+1}\})\text{ and }\mathrm{Shad}_{q_2}(\mathrm{Lex}(\{I_d\}))\mathrm{Lex}(\{I_{d+q_2}\}).$$
Since $`q_1=1`$, $`\mathrm{Shad}_i`$ of an initial block is still an initial block, and therefore we can reason on cardinalities.
The first inclusion is immediate since, for any $`A`$, $`|\mathrm{Shad}_1(A)|=|A|`$, and consequently $`|\mathrm{Shad}_1(\mathrm{Lex}(\{I_d\}))|=|\mathrm{Lex}(\{I_d\})|=|\{I_d\}|=|\mathrm{Shad}_1(\{I_d\})|`$ which is equal to $`|\mathrm{Lex}(\mathrm{Shad}_1(\{I_d\}))|`$.
For the second inclusion, we write $`\{I_d\}`$ as $`_{i=1}^sB_i`$, where $`B_i`$ are maximal blocks. It is easy to see that $`\mathrm{Shad}_{q_2}(B_i)\mathrm{Shad}_{q_2}(B_j)=\mathrm{}`$. Therefore
$$\begin{array}{cc}\hfill |\mathrm{Lex}(\mathrm{Shad}_{q_2}(\{I_d\}))||\mathrm{Lex}(\{I_d\})|& =|\mathrm{Shad}_{q_2}(\{I_d\})||\{I_d\}|\hfill \\ & =\underset{i=1}{\overset{s}{}}|\mathrm{Shad}_{q_2}(B_i)||B_i|\hfill \\ & 1\hfill \\ & =|\mathrm{Shad}_{q_2}(\mathrm{Lex}(\{I_d\}))||\mathrm{Lex}(\{I_d\})|,\hfill \end{array}$$
and the proof is concluded. โ
###### Example 4.17.
This is an example of an ideal $`I`$ which is not lexifiable in a ring for which Condition 2.3 is satisfied.
Consider $`(R,w)=(K[X,Y,Z],(1,2,4))`$ and let $`I=(X^4,Y^2,X^3Y)`$, for which we have that $`H_I(4)=2`$, $`H_I(5)=3`$, $`H_I(6)=4`$, $`H_I(7)=4`$. $`I`$ is not lexifiable, in fact, if we try to construct the associated lexicographic ideal $`L`$, we shall have $`L_4=\{X^4,X^2Y\}`$, $`L_5=\{X^5,X^3Y,XY^2\}`$ and $`L_6=\{X^6,X^4Y,X^2Y^2,X^2Z\}`$, so that $`H_L(7)5`$.
###### Example 4.18.
Let $`(R,w)=(K[X,Y,Z],(1,a,ab))`$. The ideal
$$I=(X^{ab},Y^b,X^{a+1}Y^{b1})$$
is not lexifiable, therefore $`(R,w)`$ is not Macaulay-Lex. Let $`b>2`$ and suppose that $`I`$ is lexifiable with associated lexicographic ideal $`L`$. We first observe that $`I_j`$ does not contain any monomial divisible by $`Z`$ for all $`j<2ab`$. Secondly, we show that $`X^{ab2a}Z`$ is a monomial of $`L`$. Since $`H_I(ab+(\alpha +1)a)H_I(ab+\alpha a)+1`$ for all positive $`\alpha `$, and $`H(ab+a)=5`$, one gets that $`H_I(ab+(b2)a)H(ab+a)+b3=b+2`$. Since the first $`b+2`$ monomials in degree $`ab+\alpha a`$ are $`X^{ab+\alpha a},X^{ab+(\alpha 1)a}Y,\mathrm{},X^{\alpha a}Y^b,X^{\alpha a}Z`$, this proves our claim. It is convenient now to write down all the monomials of degree $`2ab2a`$ which are $`Y^{2b2}`$, which is the smallest monomial of $`I`$ in this degree. These are
$$\begin{array}{cc}& X^{2ab2a},X^{2ab3a}Y,\mathrm{},X^{ab2a}Y^b,X^{ab2a}Z,\hfill \\ & X^{ab3a}Y^{b+1},X^{ab3a}YZ,\mathrm{},X^aY^{2b3},X^aY^{b3}Z,Y^{2b2}.\hfill \end{array}$$
As a consequence, it is easy to compute that $`H_I(2ab2a)`$ is $`(b+1)+(b3+1)=2b1`$. Analogously one gets that $`H_I(2ab1)=2b`$. Furthermore the monomials of $`L_{2ab2a}`$ have a certain number of multiples in degree $`2ab1`$, which we can count. We can thus estimate the cardinality of $`L_{2ab1}`$ as follows:
$$\begin{array}{cc}\hfill |L_{2ab1}|& |\{uL_{2ab1}:Zu\}+|\{uL_{2ab1}:Zu\}|\hfill \\ & =\left[(b+1)+\frac{b3}{2}+1\right]+\left[\left(\frac{b3}{2}+1\right)+1\right]\hfill \\ & =2b+1.\hfill \end{array}$$
This is a contradiction since $`I`$ and $`L`$ have by definition the same Hilbert function.
Finally, if $`b=2`$, an easy computation of the Hilbert function in degree $`3a+1`$ shows that $`I`$ is not lexifiable.
### 4.3 Polarization
In \[P\] it is shown how in the standard case the lexicographic ideal $`L`$ associated with an ideal $`IR`$ can be also obtained as the result of a finite process which consists of three fundamental steps, which are a) polarizing a monomial ideal b) modding out by a sequence of generic linear forms and c) taking initial ideals (with respect to the lexicographic order). In the non-standard case, following step-by-step the original proof and using generic sequences of homogeneous forms (which are not necessarily linear) it is not difficult to prove that the same procedure also terminates and leads to an ideal $`I๐ฉ`$, which we call completely polarized . Since, as we have already seen, not all ideals are lexifiable, one might make use of the $`I๐ฉ`$, which is a strongly stable monomial ideal with the same Hilbert function as $`I`$, instead.
###### Example 4.19.
Let $`(R,w)=(K[X,Y,Z],(1,2,4))`$ and
$$I=(X^8,X^6Y,X^4Y^2,X^2Y^3,Y^4,X^2YZ,X^6Z).$$
One can verify that the ideal
$$L=(X^8,X^6Y,X^4Y^2,X^4Z,X^2Y^3,X^2YZ,X^2Z^2,Y^6)$$
is the lexicographic ideal associated with $`I`$ and, thus, $`I`$ is lexifiable. On the other hand
$$I๐ฉ=(X^8,X^6Y,X^4Y^2,X^4Z,X^2Y^3,X^2Y^2Z,Y^4).$$
This shows that, even in the favourable case when Condition 2.3 is satisfied, one might have $`I๐ฉL`$.
We conclude by posing the following question.
###### Question 4.20.
Is there a combinatorial characterization of completely polarized ideals?
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# Weakly Interacting, Dilute Bose Gases in 2D
## I Introduction
### I.1 Revival of interest in low-dimensional systems
Low-dimensional systems are interesting in general, as their low-temperature physics is governed by strong long-range fluctuations. These fluctuations inhibit the formation of the true long-range order (LRO), which is a key concept of phase transition theory in 3D. Thus, a 2D uniform system of interacting bosons does not undergo Bose-Einstein condensation at finite temperatures. However, this system turns superfluid below a certain temperature $`T_{KT}`$, identified by Berezinskii, and Kosterlitz and Thouless (BKT) in 1971-73, signalling the presence of a so-called topological order. The elementary excitations of the superfluid phase are pairs of vortices with opposite winding numbers.
The experimental realization of such a system was for many years restricted to films of superfluid <sup>4</sup>He on surfaces, which is also an example of a strongly-interacting system. The breakthroughs in experimental physics at the end of the last century have changed the situation drastically. The combination of laser cooling (S. Chu, C. Cohen-Tannoudji, W. D. Phillips, Nobel Prize for Physics, 1997) with evaporative cooling and magneto-optical traps provided experimental systems of cold atoms, which were primarily used to observe the long-awaited phenomenon of Bose-Einstein condensation (E. A. Cornell, W. Ketterle, E. Wieman, Nobel Prize for Physics 2001). The full tunability of magnetic and optical traps opens an extraordinary opportunity to study in practice not only 1D and 2D Bose systems, but also dimensional crossovers under the influence of the number of particles, size and shape of the system, interaction strength and temperature. These new developments have triggered a revival of theoretical interest in low-dimensional systems, when the old theoretical predictions are to be tested or carefully revised in order to address finite-size experimental systems, and a large field of new phenomena are to be explained.
While the first experimental indications of the BKT transition in weakly-interacting Bose system have been recently obtained Stock *et al.* (2005), many questions remain unanswered. One of the most interesting is, whether topological order survives under some conditions in the inhomogeneous trapped system, or is it dominated by the true LRO and Bose-Einstein condensation prevails? Can we control and directly observe the formation of vortex pairs in 2D quantum gases? These and other problems serve as the main motivation for this Colloquium.
In the next section we present a succinct overview of the history of work with dilute Bose systems, outlining some of the important theoretical problems relevant to weakly-interacting Bose gases.
### I.2 Historical overview
The condensation of conserved particles that obey the same statistics as photons was predicted by Einstein in 1924 even before the concept of Fermi statistics was introduced (1926). Einsteinโs prediction was preceded by an ingenious conjecture of Bose, who realized that black body radiation can be treated as a gas of indistinguishable photons. Einstein generalized ideas of Bose to material particles and published two famous papers, in which he developed what we now call Bose-Einstein statistics Einstein (1924, 1925).
The ideal gas of Bose particles is remarkably the only example of a non-interacting system in condensed matter physics that undergoes a phase transition upon decreasing the temperature. However, experimental realization of ideal Bose-Einstein condensates is extraordinarily difficult, since realistic systems always involve interactions. Largely for this reason Einsteinโs ideas did not receive a wide recognition in the scientific community for many years as being devoid of any practical significance. The condensation phenomenon did not even appear in the textbooks, until in 1938 F. London recognized the analogy between superfluidity of liquid <sup>4</sup>He, discovered by Kapitza (1938), and Allen and Misener (1938) and an ideal Bose gas and emphasized that Einsteinโs statement was โerroneously discreditedโ London (1938).
In support of Londonโs phenomenological ideas, the first microscopic theory of superfluidity in a system of weakly-interacting Bose particles was introduced in a brilliant paper by Bogoliubov (1947). Subsequent discussion about the connection between superfluidity and Bose Einstein condensation led Penrose and Onsager (1956) to formulate the generalized criterion for BE condensation. This line of research culminated in a paper of C. N. Yang, who in 1962 extended this criterion to superfluidity and superconductivity and proposed the concept of off-diagonal long-range order (ODLRO) Yang (1962). The condensed phase is characterized then by a non-vanishing asymptotic of a one-body density matrix at large distances.
During the decades which followed the work of Bogoliubov, successful field-theoretical approaches were developed and many important predictions about the thermodynamics of the interacting Bose system were made. However, apart from the successful observation of superfluidity in liquid Helium systems, the quest to create Bose-Einstein condensates (BEC) proved unrewarding for several decades. Finally, in 1995 Bose-Einstein condensates were realized in a fascinating series of experiments on rubidium and sodium vapours Ketterle *et al.* (1999); Ketterle (2001); Cornell and Wieman (2002)). The importance of this experimental achievement was recognized in the 2001 Nobel Prize for Physics, shared by E. A. Cornell, W. Ketterle, and E. Wieman.
The experimental realization of BEC has offered a unique opportunity to probe and control many interesting phenomena, not accessible or unstudied in the field of superfluidity, such as dimensional transitions, the crossover from Bose Einstein condensation to BCS pair condensation, interference effects, and disorder effects. Exotic links to cosmology Fedichev and Fischer (2003), quantum optics Recati *et al.* (2005) (two-state atomic quantum dots within a condensate), and even wetting phenomena Indekeu and Van Schaeybroeck (2004) have been recently proposed. The growing interest in Bose systems has resulted in more than 600 studies per year during last decade and the list of references related to BEC now exceeds 200 pages!
The actual observation of condensation was hindered by enormous technical difficulties, so that even 15 years ago researchers dared not to believe that nature would ever provide them with the โrightโ system. The main problem to overcome is the condensation of most systems into a solid or liquid upon cooling to low temperatures, which by-passes the BEC transition. In particular, the formation of clusters or molecules is driven by three-body collisions. The hard task for an experimentalist was therefore the creation of a gaseous system, in which three-body collisions occur much less frequently than two-body interactions.
The gas in which the two-body interactions prevail is called dilute. Diluteness implies a very low density of the gas, so that the characteristic range $`a_s`$ of the potential between the Bose particles is small compared to the mean particle distance, proportional to $`n^{1/3}`$ in three dimensions ($`n=N/V`$ being the density of the gas). The diluteness condition is therefore equivalent to the requirement that the gas parameter $`n^{1/3}a_s`$ be small
$$n^{1/3}a_s1.$$
(1)
Ultra low density of the system leads to extremely low condensation temperatures (in the nanokelvin range), realization of which was another technical obstacle for the experimentalists. At low temperature the thermal velocity of the particles $`v_T`$, which is proportional to the inverse De Broglie wave length
$$\lambda _T=\sqrt{\frac{2\pi \mathrm{}^2}{mk_BT}},$$
(2)
becomes very small ($`v_T=\mathrm{}/m\lambda _T`$ 1 mm/sec) and at temperatures of the order of a few nK all the particles โjumpโ into a coherent ground state. Sufficient diluteness of the gas is therefore one of the crucial conditions for BEC to be observed in the experiment.
In order to reach the required regime of temperature and density, various cooling and trapping techniques have been developed Ketterle *et al.* (1999). Before being cooled atoms are confined in an external potential created by an applied magnetic field. The finite extent of the condensate cloud and its inherent inhomogeneity introduce a number of important differences between BEC in a trap and uniform gas. For example, a trapped gas of Bose atoms exhibits a BEC transition not only in momentum space, but in coordinate space as well Dalfovo *et al.* (1999). In practice however, condensates are so small that the literal observation of their size and shape is limited by the resolution of existing experimental equipment. Nevertheless real space Bose condensates provide a novel resource for exploring many interesting phenomena, such as quantum interference effects and frequency dependent collective excitations.
The effect of a magnetic trap becomes more dramatic for lower dimensionality of the system. For example, in 2D a noninteracting trapped gas undergoes a BEC phase transition at finite temperature Widom (1968); Bagnato and Kleppner (1991); Li *et al.* (1999) in contrast to the 2D uniform case, where condensation is possible only at zero temperature. This difference arises due to modification of the density of states of the gas in the presence of a trap.
The description of an interacting system in a 2D harmonic potential is not trivial. In the case of a uniform gas, long range order does not develop because of the preponderance of long wavelength phase fluctuations, inherent to low-dimensional systems. This can be also seen as an infrared divergence of the integral $`N(๐ฉ)\frac{d^2p}{(2\pi \mathrm{})^2}`$, where $`N(๐ฉ)`$ is the number of particles out of the condensate with momentum $`๐ฉ`$. This divergence, on the other hand, is a consequence of the fact, that the energy of the system depends only on the phase gradient, and not on the phase itself, because the latter is not a well-defined quantity Lifshitz and Pitaevskii (2004). The absence of long-range order in 2D systems with a continuous symmetry is often referred to as the Bogoliubov $`k^2`$, or Hohenberg-Mermin-Wagner (BHMW) theorem (see works by Bogoliubov (1961, 1991), Wagner (1966), Mermin and Wagner (1966), and Hohenberg (1967)), and we discuss this issue in more detail in Chapter IV.2. Fisher and Hohenberg (1988) pointed out that a consequence of the long-wavelength phase fluctuations is a drastic modification of the diluteness condition, so that the conventional low-density requirement for weakly-interacting 2D Bose gas, $`n^{1/2}a_s1`$ is replaced by an inequality
$$\mathrm{ln}\mathrm{ln}\frac{1}{na_s^2}1.$$
(3)
Taken literally, condition (3) rules out the possibility of experimental realization of a 2D dilute Bose system. However, this condition does not work away from the transition. One can show from the analysis of quantum fluctuations (see Petrov *et al.* (2004) for review) that in this case the diluteness criterion amounts to $`1/\mathrm{ln}(1/na^2)1`$, previously derived by Schick (1971).
It is also intuitively clear that the trapping potential introduces a lower bound for the momentum of excitations and thus prevents the establishment of the long-range thermal fluctuations which destroy the condensate. Based on these arguments, Petrov *et al.* (2000) showed the existence of a true condensate in a quasi-2D system in a wide parameter range.
More generally, the BHMW approach is not really suitable for a proper analysis of an inhomogeneous system, such as trapped atomic vapour, as pointed out by Fischer (2002, 2005). In his work Fischer (2002, 2005) obtained a geometrical equivalent of the BHMW theorem, independently of the Hamiltonian of the system and showed that in the marginal $`d=2`$ case true condensation is still possible in an appropriately defined thermodynamic limit.
In support of theoretical estimations, the first experimental confirmations of macroscopic occupation of the harmonic oscillator ground state Gรถrlitz *et al.* (2001); Rychtarik *et al.* (2004) became known in sodium atom vapours, confined to optical and magnetic traps. A rapid progress in experimental techniques made it possible to increase the aspect ratio (anisotropy) of the trap from 79 Gรถrlitz *et al.* (2001) to 700 Smith *et al.* (2005). This large anisotropy of the new traps is sufficient to confine condensates with $`10^5`$ atoms in a quasi-2D regime Smith *et al.* (2005). Signs of local coherence were also observed in a two-dimensional gas of hydrogen atoms, absorbed on liquid <sup>4</sup>He surface Safonov *et al.* (1998). Quasi-2D condensate have been also recently created by Stock *et al.* (2005) and interesting phase defects have been measured. The crossover from 3D condensates to two- and ultimately 1D can be observed by changing the aspect ratio of the trap.
As indicated in previous section, the recent progress in laser-based trapping techniques and creation of optical lattices has led to a new generation of remarkable experiments. With controllable interparticle interaction it is now possible to observe the transition from the superfluid state to a Mott insulator Bloch (2004). Optical lattices provide a way to investigate various intriguing aspects of low-dimensional systems as well. Interest in 2D configurations of Bose particles has arisen in the context of high-temperature superconductivity and the fractional quantum Hall effect. All in all, ultracold atomic gases have the potential to impact a very broad range of physics.
In this Colloquium we discuss a number of selected issues related to two-dimensional weakly-interacting neutral Bose gases. If necessary, 3D problems are mentioned. We attempt to cover many references and otherwise refer the reader to numerous resources, such as several excellent theoretical reviews Dalfovo *et al.* (1999); Castin (2001); Leggett (2001); Fetter (2002); Petrov *et al.* (2004); Yukalov (2004) and books Pines (1962); Griffin (1993); Pethick and Smith (2002); Pitaevskii and Stringari (2003), a Resource Letter for BEC Hall (2003) and BEC web-sites. Though a certain level of subjectivity is unavoidable, we aim to provide the necessary information about the field to those who feel lost after a preliminary contact with current literature but want to learn more about the main problems in the fascinating area of Bose-Einstein condensates in 2D.
## II Ideal Bose gas
Consider a macroscopic system of noninteracting Bose particles at finite temperature in the grand-canonical ensemble. The total number of particles in such a system is defined by the equation
$$N=\underset{k}{}n_B(ฯต_k)=\rho (ฯต)n_B(ฯต)๐ฯต,$$
(4)
where $`n_B(ฯต_k)=1/(\mathrm{exp}\beta (ฯต_k\mu )1)`$ is the Bose-Einstein distribution function, $`\beta =1/k_BT`$ and $`\rho (ฯต)`$ is the density of states.
The chemical potential $`\mu `$ of the Bose gas, being negative, increases as the temperature drops and vanishes at the critical temperature $`T_c`$, indicating the phase transition to a condensed state. The transition temperature is therefore defined by Eq.(4) with $`\mu =0`$. In $`d`$ dimensions the density of states $`\rho (ฯต)=dN_ฯต/dฯตฯต^{(d2)/2}`$, and the particle density is proportional to the integral
$$n\frac{N}{V}\frac{ฯต^{(d2)/2}dฯต}{exp(ฯต/T_c)1}.$$
(5)
In 3D this integral converges and the Bose-Einstein condensation temperature has a finite value, $`T_c^{3D}n^{2/3}`$. This result can be also understood as a temperature scale at which the thermal wave length becomes comparable with the average interparticle spacing $`\lambda _Tln^{1/3}`$. As $`\lambda `$ is proportional to $`T^{1/2}`$, $`T_c`$ is proportional to $`n^{2/3}`$.
One can also calculate the and the number of particles occupying the ground state
$$N_0=N\left(1\frac{T}{T_0}\right)^{3/2}.$$
(6)
It is readily seen that $`N_0`$ increases as the temperature decreases. This phenomenon of macroscopic occupation by particles of the state with minimal energy at low temperatures is referred to as Bose-Einstein condensation. Note that the actual condensation occurs in momentum space.
In 2D, a constant density of states leads to an infrared divergent integral in expression (5) and condensation is not possible at any finite temperature.
Let us discuss how this picture changes in the presence of a trap. The general treatment of this problem was considered by Bagnato and Kleppner (1991). They studied the possibility of the Bose-Einstein condensation of an ideal gas, confined by one- or two-dimensional power-law trap: $`V_{ext}x^\eta `$. Bagnato and Kleppner (1991) showed that a two-dimensional system undergoes BEC for any finite value of $`\eta `$, moreover, $`T_c^{2D}`$ has a broad maximum in the vicinity of $`\eta =2`$, i.e. for a trapping potential close to parabolic. (A one-dimensional system displays BEC only for $`\eta <2`$.)
Practically the confining trap is well approximated by a harmonic potential
$$V_{ext}(๐ซ)=\frac{m}{2}(\omega _x^2x^2+\omega _y^2y^2+\omega _z^2z^2).$$
(7)
For non-interacting particles we can write the many-body Hamiltonian as a sum of one-particle Hamiltonians $`H_{MB}=_{i=1}^NH_{SP}(i)`$, whose eigenvalues are
$$ฯต_{n_xn_yn_z}=(n_x+\frac{1}{2})\mathrm{}\omega _x+(n_y+\frac{1}{2})\mathrm{}\omega _y+(n_z+\frac{1}{2})\mathrm{}\omega _z.$$
(8)
The lowest energy of the system in the trap is $`ฯต_{000}=\frac{3}{2}\mathrm{}\overline{\omega }`$, where for the sake of simplicity we introduced the average frequency $`\overline{\omega }=(\omega _x+\omega _y+\omega _z)/3`$.
Note, that in the ground state all $`N`$ particles occupy the level $`ฯต_{000}`$ and the wave function of the โcloudโ of these particles is easy to find
$`\varphi (๐ซ_\mathrm{๐}..๐ซ_๐)={\displaystyle \underset{i}{}}\phi _0(๐ซ_๐ข)`$
$`\phi _0(๐ซ_๐ข)=\left({\displaystyle \frac{m\omega _{ho}}{\pi \mathrm{}}}\right)^{3/4}\mathrm{exp}\left({\displaystyle \frac{m}{\mathrm{}}}(\omega _xx^2+\omega _yy^2+\omega _zz^2)\right)`$
where
$$\omega _{h0}=(\omega _x\omega _y\omega _z)^{1/3}.$$
(10)
In this case the density distribution of the particles is position dependent
$$n(๐ซ)=N|\phi _0(๐ซ)|^2$$
(11)
and the first important length scale appearing in the problem is the size of the cloud
$$a_{h0}=\sqrt{\frac{\mathrm{}}{m\omega _{h0}}}$$
(12)
which is just the average width of the Gaussian distribution (II) (Fig.1). Experimentally $`a_{h0}`$ is usually of order of 1$`\mu m`$.
Since we are mostly interested in low-dimensional effects, it is instructive to mention the experimental realization of a two-dimensional atomic trap. An axially symmetric harmonic potential can be written in the form, $`V_{ext}(r)=\frac{1}{2}m\omega _{}^2r_{}^2+\frac{1}{2}m\omega _z^2z^2=\frac{1}{2}m\omega _{}^2(r_{}^2+\lambda ^2z^2)`$, where $`\lambda =\omega _z/\omega _{}`$ characterizes the degree of anisotropy. For $`k_BT\mathrm{}\omega _z`$ and $`k_BT>\mathrm{}\omega _{}`$ the motion of atoms along the $`z`$ direction is frozen (particles only undergo zero point oscillations), and kinematically the gas can be considered as two dimensional. Thus by making one dimension of the trap very narrow, oscillator states become widely separated, and an effective 2D system is realized.
At finite temperature only a fraction of the particles $`N_0`$ occupies the lowest energy level and the others are thermally distributed over higher energy levels. However, we still can treat $`N_0`$ as a macroscopic number. Thermal excitations will cause the size of the atomic cloud to grow with temperature. In the semiclassical approximation $`k_BT>>\mathrm{}\omega _{h0}`$, where the relevant excitation energies are much larger than the interlevel spacing, it can be shown that the size of the cloud increases as a square root of temperature $`R_T=a_{h0}\sqrt{\frac{k_BT}{\mathrm{}\omega _{h0}}}`$. The important conclusion of this short discussion is that in harmonic traps, Bose condensation manifests itself as sharp peak in the central region of the density distribution in real space. The appearance of such a peak in both coordinate and momentum space is a peculiar feature of the trapped condensates, with significant impact on both theory and experiment. This is very different from the uniform gas discussed above, where the condensation cannot be revealed in real space, for the condensate and uncondensed particles occupy the same volume.
The total number of particles in the trap is defined by
$$N=\underset{n_xn_yn_z}{}\frac{1}{\mathrm{exp}(\frac{ฯต_{n_xn_yn_z}\mu }{T})1}$$
(13)
which is derived from Eq.(4) with a discrete energy spectrum (8). Note, that in this case the chemical potential at the transition point acquires a non-zero value of the lowest energy level: $`\mu (TT_c)\mu _c=(3/2)\mathrm{}\overline{\omega }`$.
In the semiclassical approximation we can simplify (13) by replacing the summation with integration and a straightforward solution for $`\mu =\mu _c`$ gives the Bose-condensation temperatures for the trapped gas in three and two dimensions
$$T_c^{3D}=\frac{\mathrm{}}{(\zeta (3))^{1/3}}\omega _{ho}N^{1/3}$$
(14)
$$T_c^{2D}=\frac{\mathrm{}\sqrt{6}}{\pi }\omega _{ho}N^{1/2}$$
(15)
The two-dimensional condensation temperature is now finite (nonzero). This is related to the density-of-states effect of the gas in the trap. Indeed, in the semiclassical approximation we can introduce a coordinate system defined by the three variables $`ฯต_{x,y,z}=\mathrm{}n_{x,y,z}\omega _{x,y,z}`$, in terms of which the surface of constant energy (8) is the plane $`ฯต=ฯต_x+ฯต_y+ฯต_z`$. Then the number of states $`N(ฯต)`$ is proportional to the volume in the first octant bounded by this plane
$$N_ฯต=\frac{1}{\mathrm{}^3\omega _{h0}^3}_0^ฯต๐ฯต_x_0^{ฯตฯต_x}๐ฯต_y_0^{ฯตฯต_xฯต_y}๐ฯต_z=\frac{ฯต^3}{6\mathrm{}^3\omega _{h0}^3}$$
(16)
The density of states $`\rho =dN_ฯต/dฯต`$ is then quadratic in energy $`\rho ^{3D}ฯต^2`$ in three dimensions and linear in energy in two dimensions $`\rho ^{2D}ฯต`$, in contrast to the constant density of states of a uniform 2D gas, and the integral in the Eq.(13) for $`\mu =\mu _c`$ is not infra-red divergent until $`d=1`$.
It is now straightforward to calculate the condensate fraction (e.g. 3D)
$$\frac{N_0}{N}=1\left(\frac{T}{T_c^{3D}}\right)^3$$
(17)
and total energy of the system and correspondingly all the interesting thermodynamic quantities. In 2D the condensate fraction is Bagnato and Kleppner (1991); Petrov *et al.* (2004)
$$\frac{N_0}{N}1\left(\frac{T}{T_c^{2D}}\right)^2.$$
(18)
Sign โ$``$โ in the expression (18) is related to the fact, that at $`T=T_c^{2D}`$ the condensate fraction is not exactly zero, because there is a small correction to the result due to the finite number of particles in the system Petrov *et al.* (2004). One should be therefore careful with the word โphase transitionโ in the context of trapped gases, because they are finite size systems and the phase transition notion is strictly defined only in thermodynamic limit. It is better to say that at $`T_c`$ there is a sharp crossover to the BEC state in the system. Note also, that at $`T_c^{2D}`$ the de-Broglie wavelength $`\lambda _T`$ becomes comparable with the mean interparticle separation $`\sqrt{T_c/(Nm\omega _{ho}^2})`$.
We end the section by remarking on the proper definition of the thermodynamic limit in the trapped case. It is well known that the transition temperature should be well-defined in the thermodynamic limit. The usual definition when the ratio $`N/V`$ is kept constant while the number of particles $`N`$ and the volume $`V`$ tend to infinity is apparently not suitable for the inhomogeneous situation. The appropriately defined limit is then obtained by letting $`N\mathrm{}`$ and $`\omega _{h0}0`$, while keeping the product $`N\omega _{h0}^3`$ (or $`N\omega _{h0}^2`$ in 2D) constant. In this case the temperatures (14) and (15) are well-defined.
A comprehensive survey of various issues related to the behaviour of the ideal Bose gas in a harmonic potential can be found in the paper by Mullin (1997).
The ideal-gas results are summarized in the Table 1.
## III Ground state of a weakly interacting Bose gas
### III.1 Bogoliubov approximation
In his seminal paper โOn the theory of superfluidityโ Bogoliubov (1947), published in 1947, Bogoliubov introduced the microscopic description of the ground state of a uniform, weakly interacting Bose gas. The assumption about the uniformity of the unperturbed ground state is crucial to his results. To assure a uniform Bose gas, Bogoliubov considered the case of repulsive interactions and made use of periodic boundary conditions. The gas is also assumed to be dilute ($`na^31`$), which permits to simplify the many-body problem and account for interactions in a rather fundamental way. In contrast to the uniform case, the nonuniform ground state is very โsensitiveโ to the introduction of any interactions and makes the solution of the many body problem highly nontrivial.
The standard Hamiltonian of an interacting Bose gas is
$`H`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle }\mathrm{\Psi }^+(r)\mathrm{\Psi }(r)dr+`$ (19)
$`+`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \mathrm{\Psi }^+(r)\mathrm{\Psi }^+(r^{})U(rr^{})\mathrm{\Psi }(r^{})\mathrm{\Psi }(r)๐r๐r^{}}`$
where $`U(rr^{})`$ is the interaction between particles. In momentum space this Hamiltonian reads
$$H=\underset{p}{}ฯต_pa_p^+a_p+\frac{1}{2V}\underset{pp^{}q}{}U_qa_p^+a_p^{}^+a_{p^{}q}a_{p+q},$$
(20)
$`U_q=e^{iqr}U(r)๐r`$ is a Fourier component of the interaction, the bosonic field operator $`\mathrm{\Psi }(x)=1/\sqrt{V}_pe^{ipx}a_p`$ (here $`x`$ is a four vector), and the boson creation and annihilation operators satisfy the usual commutation relations $`[a_p,a_p^{}^+]=\delta _{pp^{}}`$.
Without interactions all $`N`$ particles of the system occupy the state with zero energy and zero momentum. The number of condensed particles $`N_0`$ in this case is equal to the total number of particles $`N`$. When we switch on the interaction, two particles can scatter out of the condensate and occupy one of the the many zero-total-momentum states with separate momenta $`๐ค`$ and $`๐ค`$ (in the lowest order perturbation theory) and $`N_0`$ naturally decreases.
For a dilute weakly interacting Bose gas one can assume that the total depletion of the condensate is small ($`\delta N/N_01`$) and most of the particles remain in the condensate $`N_01`$. The key observation of Bogoliubov is that in this case the second-quantized condensate operators can be simply replaced by the โcโ- number $`\sqrt{N}_0`$
$$\widehat{a}_0,\widehat{a}_0^+\sqrt{N_0}.$$
(21)
The drawback of this prescription is that it leads to a Hamiltonian which no longer conserves the number of particles. This problem can be partly resolved by working in the grand-canonical ensemble, in which additional terms $`\mu N_p`$ ($`N_p=_{p0}a_p^+a_p`$) are introduced into the Hamiltonian (75). This secures the conservation of particles on the average. It is also worth mentioning that the Bogoliubov approximation is equivalent to the neglect of any dynamics in the condensed state.
In the weak coupling limit the Hamiltonian (75) can be diagonalized by applying the Bogoliubov canonical transformation
$`a_k=u_k\alpha _kv_k\alpha _k^+`$
$`a_k^+=u_k\alpha _k^+v_k\alpha _k`$ (22)
and the resultant Hamiltonian describes the system of non-interacting quasiparticles with the spectrum
$$\xi _k=\sqrt{n_0U_0\frac{k^2}{m}+\frac{k^4}{4m^2}},$$
(23)
where $`n_0=N_0/V`$ is the density of condensed particles.
From this dispersion relation (23) it follows that in the long wavelength limit the Bogoliubov quasiparticles behave as โphononsโ with a sound velocity $`s=\sqrt{\frac{n_0U_0}{m}}`$, and all of the low temperature thermodynamics of a Bose-condensed system is governed by this phonon spectrum. In the opposite limit of short wavelength the quasiparticles behave as free particles with an energy $`\frac{k^2}{2m}`$. By equating the kinetic energy and the โHartreeโ interaction energy $`n_0U_0`$ one can straightforwardly find the โtransitionโ wave vector $`k_c=\sqrt{2mn_0U_0}\sqrt{2}ms`$, which separates the phonon-like behavior of elementary excitations from the free particle one. $`k_c`$ introduces an important length scale into the system (Fig.2)
$$\lambda _c=\frac{\mathrm{}}{k_c}=\mathrm{}/\sqrt{2mn_0U_0},$$
(24)
over which the coherence effects are important in the interaction between particles. It is usually called the healing length (as in the context of trapped condensates), or sometimes the correlation or coherence length, and it refers to correlations between excitations in the system. These correlations are distinct from the long-range correlations, which lead to condensation in the $`k=0`$ mode.
One should also mention that the Bogoliubov canonical transformation is equivalent to a summation over the most divergent terms in the perturbation-series expansion for the ground state energy. The summation of such series is also equivalent to making the random-phase approximation (RPA).
It was important in the theory of superfluidity, that the low-lying Bogoliubov quasiparticles follow a linear dispersion. This kind of behavior is fully consistent with the Landau criterion for superfluidity, i.e., that no excitation can be created in a liquid moving with a velocity $`v`$ less than that of a sound ($`v<s`$). In case of noninteracting particles the dispersion is quadratic for all $`k`$ and superfluidity is not possible.
### III.2 Field-theoretical approaches: $`t`$-matrix approximation
To go beyond the Bogoliubov approximation, one needs to take both multiple scattering diagrams and RPA contributions into account. That can be done for example by means of a pseudopotential method Lee *et al.* (1957), or by field-theoretical methods, first applied to the Bose gas of small density at $`T=0`$ by Beliaev (1958a, b) and by Hugenholtz and Pines (1959).
The presence of the many particle condensate in the ground state of the interacting Bose gas was the main obstacle to application of the usual technique of Feynman diagrams to this system. Consider for example the one-particle Greenโs function in the interaction representation
$$G(xx^{})=i\frac{<T\{\mathrm{\Psi }(x)\mathrm{\Psi }^+(x^{})S\}>}{<S>}$$
(25)
Here the average is taken over the ground state of N non-interacting Bose particles, which are all in the condensate ($`N_0=N`$). The $`S`$-matrix is expressed as usual
$`S`$ $`=`$ $`T\{exp({\displaystyle \frac{i}{2}}{\displaystyle }d^4x_1d^4x_2U(x_1x_2)\times .`$ (26)
$`\times `$ $`.\mathrm{\Psi }^+(x_1)\mathrm{\Psi }^+(x_2)\mathrm{\Psi }(x_2)\mathrm{\Psi }(x_1))\}`$
where $`x_1`$ and $`x_2`$ are the four-vectors, and the interaction is $`U(x_1x_2)=U(r_1r_2)\delta (t_1t_2)`$. In order to derive the diagram series for the Greenโs function, we need to expand the $`S`$\- matrix in powers of $`H_{int}`$. Usually the terms containing the odd number of annihilation operators vanish after averaging over the ground state, which unfortunately does not happen in the case of a Bose gas due to the above mentioned peculiarities of the ground state. The expectation value of the $`N`$ product containing $`a_0`$ apparently does not disappear and the standard method of constructing diagrams cannot be applied in the case of an interacting Bose gas.
This difficulty was successfully resolved by Beliaev in 1958. He noticed that for a large number of particles $`N`$ the diagrammatic approach can be applied to particles with momenta $`p0`$, while the condensed phase (which does not disappear when the interactions are turned on) can be described as a sort of external field. It is thus convenient to separate the operators $`a_0`$ and $`a_0^+`$ (which act only on the ground state) from $`\mathrm{\Psi }`$ and $`\mathrm{\Psi }^+`$
$$\mathrm{\Psi }=\mathrm{\Psi }^{}+a_0/\sqrt{V};\mathrm{\Psi }^+=\mathrm{\Psi }^{}_{}{}^{}++a_0^+/\sqrt{V}.$$
(27)
The Greenโs function (25) is then also divided into two parts, and the operations $`T`$ and $`<\mathrm{}>`$ are represented as two successive operations, the former acting only on $`\mathrm{\Psi }^{}`$ and $`\mathrm{\Psi }^{}_{}{}^{}+`$, and the latter acting only on $`a_0`$ and $`a_0^+`$. The operators $`a_0`$ and $`a_0^+`$, occurring in the $`S`$ matrix are treated as parameters, and the expectation values over $`\mathrm{\Psi }^{}`$, $`\mathrm{\Psi }^{}_{}{}^{}+`$ ground state can be now calculated, using standard techniques.
With these ideas in mind, Beliaev succeeded in deriving a general expression for the one-particle Greenโs function of the interacting system in terms of some effective self-energies $`\mathrm{\Sigma }_{ik}`$ and the chemical potential $`\mu `$. However, the exact calculation of the Greenโs functions proved to be very complicated, and approximate methods of summing the series of Feynman graphs were developed.
For simplicity, Beliaev considered a short-range, central interaction potential $`U_๐ฉ=U_0`$ for $`p<1/a`$ and $`U_๐ฉ=0`$ for $`p>1/a`$. In the low density limit $`n_0a^31`$, where $`n_0`$ is the density of the particles in the condensate, he obtained a crucial result, that the main contributions to the self-energies of the Greenโs function originate from ladder diagrams. In this case the real interaction $`U`$ is replaced by an effective two-particle interaction $`\mathrm{\Gamma }`$, representing the sum of contributions from all ladder type Feynman graphs (Fig.3). The integral equation for the vertex $`\mathrm{\Gamma }`$, called the Bethe-Salpeter equation is
$`\mathrm{\Gamma }(x_1,x_2;x_3,x_4)=U_{x_1x_2}\delta (x_1x_3)\delta (x_2x_4)`$
$`+i{\displaystyle d^4x_5d^4x_6U_{x_1x_2}G^0(x_1x_5)G^0(x_2x_6)\mathrm{\Gamma }(x_5,x_6;x_3,x_4)},`$
(28)
where $`x(๐ซ,t)`$. In momentum representation this equation reads
$`\mathrm{\Gamma }(p_1,p_2;p_3,p_4)=U_{p_1p_2}+`$
$`i{\displaystyle d^4p_5d^4p_6U_{p_1p_5}G^0(p_5)G^0(p_6)\mathrm{\Gamma }(p_5,p_6;p_3,p_4)},`$
(29)
where the momentum conservation condition $`p_1+p_2=p_3+p_4=p_5+p_6`$ is implied and $`p_i(๐ฉ_๐ข,p_0^i)`$.
It is convenient to introduce relative and total momenta according to
$`p_1+p_2=P^{};p_3+p_4=P`$
$`p_1p_2=2p^{};p_3p_4=2p.`$ (30)
This transformation leads to the following equation
$`t(p^{},p,P)=U(p^{}p)+`$
$`+i{\displaystyle \frac{d^4q}{2\pi ^4}U(p^{}q)G^0(P/2+q)G^0(P/2q)t(q,p,P)}`$
(31)
where we denote $`\mathrm{\Gamma }`$ in the center of mass representation by โ$`t`$โ, and the free particle Greenโs function is $`G^0(p)=(p^0p^2/2m+i\delta )^1`$.
A conventional $`t`$ matrix equation is obtained from (31) after carrying out the integration over $`q_0`$. In two dimensions this results in the following equation
$$t(๐ฉ^{},๐ฉ,P)=U_{๐ฉ^{}๐ฉ}\frac{d๐ช}{(2\pi )^2}U_{๐ฉ^{}๐ช}\frac{t(๐ช,๐ฉ,P)}{k_0^2q^2/m+i\delta },$$
(32)
where $`k_0^2=P^0๐^\mathrm{๐}/4m`$. In the scattering theory this equation is also known as the Lippmann-Schwinger equation. Physically the $`t`$-matrix corresponds to the renormalization of the interaction by multiple scattering of one particle off another.
The standard way to treat the dilute Bose gas is thus to replace the real potential, which is usually strongly singular, by the zero momentum $`t`$ matrix generated from multiple two-particle scattering, represented by the infinite summation of the ladder diagrams described above.
The $`t`$ matrix equation (32) cannot be solved explicitly, but in general its solution can be expressed in terms of the scattering amplitude of two particles in vacuum. The scattering amplitude $`f(๐ฉ^{},๐ฉ)`$ for a transition from the initial relative wave vector $`๐ฉ`$ to a finite relative vector $`๐ฉ^{}`$ is defined by an expression
$$f(๐ฉ^{},๐ฉ)=๐๐ชU(๐ฉ^{}๐ช)\mathrm{\Psi }_๐ฉ(๐ช),$$
(33)
where $`\mathrm{\Psi }_p`$ is a wave function of a scattering problem with potential $`U`$ that satisfies the following Schrรถdinger equation in momentum representation
$$(k^2p^2)\mathrm{\Psi }_k(๐ฉ)๐๐ชU(๐ฉ๐ช)\mathrm{\Psi }_k(๐ช)=0.$$
(34)
According to elementary scattering theory Dalfovo *et al.* (1999); Castin (2001); Leggett (2001); Fetter (2002), at low energies $`s`$-wave scattering becomes dominant, and the scattering amplitude $`f_0`$ is approximated to leading order by
$$f_0\frac{4\pi \mathrm{}^2a_s}{m},$$
(35)
(where the momentum dependence of the scattering amplitude can be ignored in the low energy limit). Thus at low energies, in vacuum the only remaining parameter characterizing the interaction is the $`s`$-wave scattering length $`a_s`$.
In general, the t-matrix (32) requires the knowledge of the scattering amplitude for $`k_0^2q^2/m`$, known as โoff-the-energy-shellโ t-matrix. For two-particle scattering in vacuum, discussed above, only on-shell t-matrix is physically relevant. In the situation when three-body collisions become important, the calculation of the off-shell t-matrix is necessary Fadeev (1960). In the context of the dilute Bose gases the off-shell t-matrix arises in connection with so-called many-body t-matrix approach Stoof and Bijlsma (1993); Bijlsma and Stoof (1997); Proukakis *et al.* (1998), which we discuss in the next Chapter. The many-body t-matrix takes into account the effect of the medium (mean field) in which the collisions occur. At the low energy limit the many-body t-matrix is approximated by the off-shell two-body t-matrix Morgan *et al.* (2002). The solution of the off-shell t-matrix was first proposed by Beliaev (1958b) and Galitskii (1958). The alternative approach based on the inhomogeneous Schrรถdinger equation, which allows to treat the hard-sphere central potentials in one, two and three dimensions, was considered by Morgan *et al.* (2002). Morgan *et al.* (2002) have shown for any dimension that for all potentials with a finite range, the long-wavelength limit of the off-shell t-matrix depends only on energy and not on the initial and final relative momenta of the scattered particles. This result means that low-energy collisions can be represented by a contact potential.
Consider now the quasiparticle spectrum within the first order Beliaev approach. It turns out one can reproduce the Bogoliubov result (23) with the only difference that instead of potential $`U_0`$ the momentum independent scattering amplitude $`f_0`$ appears, for in the first order $`U_0`$ is equal to $`f_0`$. The healing length (24) can then be related to a scattering length
$$\lambda _c=\frac{1}{\sqrt{8\pi a_sn_0}}.$$
(36)
The second order approximation does not modify the physical picture of the low temperature behaviour of the interacting Bose gas, but provides the corrections to the sound velocity, and a damping proportional to $`p^5`$ related to the process of decay of one phonon into two. The third order corrections involve the solution of a three particle problem, which to date has not been solved.
We now turn to the 2D system. Following the methods developed by Beliaev, Schick (1971) examined a two-dimensional system of hard-disk bosons of diameter $`a`$ at low densities and absolute zero (see also the recent study of Ovchinnikov (1993). The dimensionless expansion parameters are the interaction $`U_0`$ and the gaseous parameter $`na^2`$, which is small in the dilute limit. The application of Beliaevโs method to 2D systems is not as straightforward as it is for 3D systems. In the 3D case, the ladder diagrams are the only contributions which do not depend on the small parameter $`na^3`$ and therefore it is natural to take them into account while calculating the first term in the expansions of all quantities in terms of density. In 2D the contributions from the ladder diagrams depend logarithmically on the parameter $`na^2`$, in particular, the effective interaction, or $`t`$ matrix is proportional to $`1/\mathrm{ln}(1/na^2)`$
$$f_0^{2D}\frac{4\pi }{m\mathrm{ln}(1/na^2)}.$$
(37)
The key conclusion of Schick (1971) is that $`1/\mathrm{ln}(1/na^2)`$ plays the role of the small parameter in the 2D dilute system at zero temperature and the dominant contributions are derived from the diagrams of first order in this parameter. In this approximation he calculated the leading order correction to the chemical potential
$$\mu =\frac{4\pi \mathrm{}^2n}{m\mathrm{ln}(na^2)}\left\{1+O[1/\mathrm{ln}(na^2)]\right\}$$
(38)
and the quasi-particle excitation spectrum
$$\xi _k=\sqrt{\mu \frac{k^2}{m}+\frac{k^4}{4m}}=\sqrt{\frac{k^4}{4m}+\frac{4\pi n}{m\mathrm{ln}(1/na^2)}k^2}.$$
(39)
In the long-wavelength limit the quasi-particles behave as phonons with a speed of propagation $`s=\sqrt{4\pi n/(m\mathrm{ln}(na^2))}`$. The spectrum changes from phonon-like to free particle-like in the vicinity of the momentum $`k_c`$ defined as
$$kak_ca16\pi na^2(m\mathrm{ln}(na^2))^11.$$
The ground state energy per particle and the condensate fraction take the form Schick (1971)
$`E/N`$ $`=`$ $`{\displaystyle \frac{2\pi \mathrm{}^2n}{m\mathrm{ln}(na^2)}}\left\{1+O[1/\mathrm{ln}(na^2)]\right\}`$ (40)
$`{\displaystyle \frac{n_0}{n}}`$ $`=`$ $`1+{\displaystyle \frac{1}{\mathrm{ln}(na^2)}}+O\left[1/(\mathrm{ln}(na^2))^2\right].`$ (41)
### III.3 Gross-Pitaevskii mean-field theory
The ground state and thermodynamic properties of an interacting Bose system confined to an external potential $`V_{ext}=\frac{1}{2}\mathrm{}\omega _{h0}(r/a_{h0})^2`$ ($`a_{h0}`$ is the trap size (12)) can be directly calculated from the Hamiltonian
$`H={\displaystyle }dr\psi ^+(r)({\displaystyle \frac{\mathrm{}^2}{2m}}^2+V_{ext}(r))\psi (r)+`$
$`+{\displaystyle \frac{1}{2}}{\displaystyle ๐r๐r^{}\psi ^+(r)\psi ^+(r^{})U(rr^{})\psi (r^{})\psi (r)}`$ (42)
using numerical methods, such as quantum Monte Carlo. Nevertheless for most experimentally relevant situations (when the number of atoms is large) the mean-field description of the system proves to be sufficient. In this case the macroscopic low energy behaviour of the system can be explored under the assumption that the order parameter varies over distances larger than the mean interparticle spacing.
Such a mean-field approximation was first developed by Gross and Pitaevskii. Their approach, which is valid in the dilute limit, is a straightforward generalization of Bogoliubov theory for the gas in the trap. One should bear in mind that the diluteness condition $`n_{max}a_s^31`$ does not automatically secure the weakness of the interactions. The interaction strength is specified by an extra parameter (see in particular the review of Dalfovo *et al.* (1999) and the paper by Fetter (1999)). The interaction energy, which is of the order of $`gNn`$ is to be compared with the kinetic energy, proportional to $`Na_{h0}^2`$. Since the average density of atoms $`nN/a_{h0}^3`$, the interaction strength can be characterized by a dimensionless parameter $`N|a_s|/a_{h0}`$. When $`Na_sa_{h0}`$, it means that the coherence length $`\lambda _c`$ (24) is large in comparison with the size of the trap $`a_{h0}`$ and the system is assumed to be nearly ideal and is described by a Gaussian distribution (II). In the opposite limit, $`Na_sa_{h0}`$, the coherence length is small and the dilute gas exhibits important nonideal behaviour Dalfovo *et al.* (1999).
The mean-field Gross-Pitaevskii approximation is extensively presented in the literature (see for instance review by Dalfovo *et al.* (1999), and paper by Leggett (2003), and a review with an emphasis on experiment by Angilella*et al.* (2006)), therefore we only mention briefly the key concepts of its derivation. Gross and Pitaevskiiโs approach is based on the Bogoliubov prescription for the condensate (21), according to which the boson field operators $`\psi `$ are written as a sum of a classical field $`\varphi `$, having the meaning of the order parameter, and a small perturbation $`\psi ^{}`$
$$\psi (r,t)=\varphi (r,t)+\psi ^{}(r,t),$$
(43)
implying that the depletion of the condensate is small. As a side note, we mention that in principle the problem of the order parameter definition in a finite inhomogeneous system arises in this case, but it turns out that the wave function of the condensate has a clear meaning, if determined through the diagonalization of the one-body density matrix in analogy with liquid-helium drops Lewart *et al.* (1988). This issue is also discussed in detail in a review by Leggett (2001).
One can expand the theory in the parameter $`\psi ^{}`$ and derive the equation for $`\varphi `$ either from the standard Heisenberg equation, or alternatively by taking the variation of the classical action $`S`$ of the type $`S=๐t๐๐ซ\overline{\varphi }\left[i_t\frac{\mathrm{}^2}{2m}^2V_{ext}\frac{g}{2}\overline{\varphi }\varphi \right]\varphi `$ with respect to $`\overline{\varphi }`$ (saddle point approximation). The derivation of the Gross-Pitaevskii (GP) equation and the next order corrections within the bosonic field theory can be found in the paper by Stenholm (1998).
The resulting Gross-Pitaevskii equation is
$$i\mathrm{}\frac{}{t}\varphi (r,t)=(\frac{\mathrm{}^2}{2m}^2+V_{ext}+g|\varphi (r,t)|^2)\varphi (r,t),$$
(44)
where we have approximated the potential by a $`\delta `$-function, $`V(rr^{})=g\delta (rr^{})`$ (which we can do under the assumption that the interparticle spacing is much larger than the interaction range), and where $`g=\frac{4\pi \mathrm{}^2a_s}{m}`$ is the 3D coupling constant. This coupling constant is equal to the zero momentum limit of the scattering amplitude (35) discussed above.
In the limit $`N\mathrm{}`$ (Thomas-Fermi approximation) the kinetic energy contribution can be neglected and the Gross-Pitaevskii equation can be solved analytically. This classical Thomas Fermi approximation breaks down in the vicinity of the boundary of the condensate, where the gradient of the condensate density is no longer small.
We discuss now the coupling constant of the 2D Bose gas. It was first demonstrated by Y. Lozovik in 1971 (see the review of Petrov *et al.* (2004)) that to zero order in perturbation theory the coupling constant $`g^{2D}=\mathrm{}^2f_0^{2D}/m`$, where $`f_0^{2D}`$ is the scattering amplitude at energy of the relative motion $`E=2\mu `$.
This coupling constant can be treated as a parameter, as in the work by Bayindir and Tanatar (1998) (see also references therein), where the two-dimensional Bose gases described by the GP equation have been studied. For some range of interaction strength it was shown that interacting bosons behave similarly to the noninteracting case in a harmonic trap. For weak short-range interparticle interactions, a finite temperature BEC phase transition was found to occur.
On the other hand, the coupling constant in 2D is expected to display a logarithmic dependence on density (cf.(37)) in accordance with estimations by Schick (1971) for $`f_0`$ in case of a homogeneous gas. The precise choice of $`g^{2D}`$ has in fact been a controversial issue (see Lieb *et al.* (2001) and references therein). For example, Kim *et al.* (1999) suggested $`g^{2D}1/\mathrm{ln}(1/ka)`$, where $`0<ka1`$ and $`k`$ is the infrared cut off introduced by the trap at $`1/a_{h0}`$, so that $`g^{2D}1/\mathrm{ln}(a_{h0}/a)`$. This kind of approximation may be reasonable when the size of the trap is much larger than all other length scales in the problem.
Note, that for quasi-2D gas in a trap the coupling constant was derived by Petrov *et al.* (2000)
$$g^{Q2D}=\frac{2\sqrt{2\pi }\mathrm{}^2}{m}\frac{1}{a_{ho}/a+(1/\sqrt{2\pi })\mathrm{ln}(1/\pi k^2a_{h0}^2)}.$$
(45)
The rigorous derivation of the Gross-Pitaevskii functional for a two-dimensional interacting gas was provided by Lieb *et al.* (2001). Their analysis leads to the following expression for the coupling constant
$$g^{2D}=\frac{1}{|ln(\overline{n}a^2)|},$$
(46)
where $`\overline{n}`$ is the average density of the particles, proportional to $`\sqrt{N}`$. The mean density is defined as $`\overline{n}=1/Nn^{TF}(r)^2d^2r`$, with the Thomas-Fermi density being $`n^{TF}(r)=(\mu ^{TF}V_{ext}(r))/8\pi `$, and $`\mu ^{TF}`$ chosen so that the constraint $`n^{TF}=N`$ holds. The density expansion has been applied to the case of a 2D Bose gas at zero temperature by Cherny and Shanenko (2001) in order to derive the Gross-Pitaevskii equation.
The modification of the GP equation due to the many-body renormalization of the scattering, mentioned in the section III.2, has been provided by Lee *et al.* (2002). The effective interparticle interaction in 2D is modeled by the off-shell two-body t-matrix, that at low energies depends on the energy of the collision. The energy dependence of the effective interaction can be written in the density-dependent form and applied to trapped 2D gas. This leads to the GP equation, describing the condensate wave-function that no longer has a cubic non-linearity in $`\mathrm{\Psi }`$, but instead goes as $`(|\psi |^2/\mathrm{ln}|\psi |^2)\psi `$ Lee *et al.* (2002).
It is also interesting to analyze the deviations from the mean-field behaviour, since the experimental system is well controlled nowadays and different regimes can be realized. The corrections to the mean-field ground state solution stem from the quantum fluctuations, and their effect becomes more prominent with the growth of the gas parameter, as has been observed in Monte Carlo simulations. For the calculation of quantum corrections in a systematic way we refer the reader to the paper by Andersen and Haugerud (2002) and references therein. Many references on the GP approximation and beyond can be found elsewhere Angilella *et al.* (2004); Kolomeisky *et al.* (2000). For the effects of a third spatial dimension and the self-consistent calculation of the coupling constant see the paper of Cherny and Brand (2004) and references therein.
## IV Finite temperature problems
Zero-temperature techniques are not really suitable for controlling IR thermal fluctuations, and new methods have to be devised to describe the interacting system at finite $`T`$. At the beginning of the Chapter IV we present the generic properties of the 2D XY models and the concept of quasi-long-range order, which is the central concept in the phase transition theory in 2D. Why the true long-range order cannot form in 2D uniform system, we discuss in detail in section IV.2, and especially the way familiar concepts from 2D phase transition theory should be revised in the trapped case.
In Section IV.3 we present the theory of Popov, who pioneered the finite-temperature generalization of Beliaevโs field-theoretic approach, and described the low-temperature superfluid state of the 2D Bose gas. In Section IV.4 we show how the diluteness condition of Fisher and Hohenberg, discussed in the introduction, arises as an applicability limit of the Popovโs t-matrix approach. In Section IV.5 we describe methods, which generalized and/or improve the results of Popov, and also the Monte Carlo simulations, which are up to date most reliable numerical calculations of the superfluid phase of the 2D Bose system. Before concluding, we mention how unique the 2D system with a contact interaction is, for it possesses an inherent symmetry, which leads to the birth of the special breathing modes, which in principle can be checked experimentally.
### IV.1 Introduction: 2D XY models
For our further analysis it is important to recognize that a uniform, interacting Bose system belongs to the XY universality class, characterized by a vector order parameter (for a comprehensive analysis see the book by Chaikin and Lubensky (1995)). It means that the finite-temperature behaviour of the 2D Bose gas is determined by generic properties of the 2D XY model.
We know, that the 2D XY models are very special, for the long-range thermal fluctuations destroy the long-range order at finite temperatures (Bose-Einstein condensation in case of 2D Bose gas). The existence of these long-wavelength modes in a 2D Bose fluid was first pointed out by Bogoliubov in his โ$`k^2`$ theoremโ in 1961, and later confirmed by Hohenberg (1967) and by Mermin and Wagner (1966) (this issue is discussed in more detail in section IV.2).
However, a special type of order - topological order - which gives rise to superfluidity, can develop in a two-dimensional Bose fluid below the โKosterlitz-Thoulessโ temperature $`T_{KT}`$, as predicted by Kosterlitz and Thouless (1973) and Berezinskii (1970, 1971) using the renormalization group method (RG). Below $`T_{KT}`$ the continuous U(1) symmetry (rotations in a two-dimensional plane) is broken and the system acquires a finite rigidity, or phase stiffness $`\rho _s`$. The order parameter correlations decay algebraically (for any coupling of the XY model), and the average order parameter is zero. However locally, the order parameter can have a well-defined value. This unique situation is described in terms of quasi-long-range order (QLRO) Chaikin and Lubensky (1995). Important low-lying excitations of the QLRO phase are vortex pairs (two vortices with opposite winding numbers) whose fugacity decreases with distance, thus not destroying the connectivity of the state (therefore $`\rho _s0`$).
The phase transition to a disordered state (with $`\rho _s=0`$) is associated with a dissociation of the coupled vortex pairs. Above $`T_{KT}`$ the vortex fluid can be treated as a kind of vortex plasma, where vortices play the role of mobile โchargesโ, interacting via a Coulomb potential. In this language the state below $`T_{KT}`$ can be described as an โinsulatingโ state of bound โchargesโ. The mapping of the 2D XY model onto the two-dimensional Coulomb gas is considered in detail in the review by Minnhagen (1987).
The rigidity, or superfluid density $`\rho _s`$ does not go continuously to zero at the critical temperature, but experiences a universal jump
$$\frac{m^2k_BT_{KT}}{\mathrm{}^2\rho _s(T_{KT_{}})}=\frac{\pi }{2},$$
(47)
first predicted by Nelson and Kosterlitz (1977) and successfully verified in experiments on superfluid <sup>4</sup>He films, absorbed on a substrate Bishop and Reppy (1978).
An interesting interpretation of Kosterlitz-Thouless physics in the context of bosonic systems was put forward by Kagan *et al.* (1987) almost 20 years ago. They propose that below $`T_{KT}`$ the system forms a โquasi-condensateโ, a condensed state achieved in a local sense. The introduction of the quasi-condensate concept was motivated by a peculiar behaviour of the one-particle density matrix $`\rho (r)`$ at large distances in 2D (Fig.4).
There are two length scales associated with the behaviour of $`\rho (r)`$: the aforementioned correlation length $`\lambda _c`$ at which $`\rho (r)`$ relaxes from the value $`n`$ at $`r=0`$ to $`n_0`$, and the characteristic radius of the phase fluctuations $`R_\theta `$, which is rather large $`R_\theta \lambda _c`$. The appearance of large $`R_\theta `$ can be understood in the following way: at large distances $`\rho `$ falls off as a power law of $`r`$ (Kane and Kadanoff (1967)) $`\rho (r)n_0(r/r^{})^\alpha `$, where $`rr^{}`$ and the coefficient $`\alpha `$ is proportional to the temperature and the Schickโs parameter $`\alpha T/(T^{}ln(1/na^2))`$ ($`T^{}(2m\lambda _c^2)^1`$) and therefore is very small $`\alpha 1`$. As a result of this the density matrix $`\rho `$ decays over a large length scale $`R_\theta r^{}e^{1/\alpha }`$ Kagan *et al.* (1987).
Conceptually, the system can be divided into blocks of size $`L`$, which is smaller than $`R_\theta `$. In each block one can introduce the wave-function of the condensate with a well-defined phase. The whole system is then described in terms of an ensemble of wave-functions of the blocks. Condensate wave-functions within the ensemble corresponding to blocks separated by a distance greater than $`R_\theta `$ have uncorrelated phases, and it is impossible to define the condensate wave function for the entire system as a whole. The state of matter with a fluctuating phase is called a โquasi-condensateโ Kagan *et al.* (1987). See also the extention of Bogoliubov methods to quasi-condensates by Mora and Castin (2003).
What happens to the XY universality class concepts in an experimentally realizable system of cold atoms confined in a trap, remains a controversial issue. We will see in the next Section, that many issues should be crucially reformulated in order to address the physics of trapped cold gases.
### IV.2 Problems of the long-range order formation in 2D
The notion that the development of the long-range order (LRO) is not possible in 2D dates back to the work of Peierls (1935), who argued that the thermal motion of low energy phonons will ruin the LRO in a 2D solid. A rigorous proof of the Peierlsโ statement was provided later by Mermin (1968).
Subsequent work by Mermin and Wagner (1966) provided a proof that there is no spontaneous magnetization or sublattice magnetization in an isotropic Heisenberg model with finite range interactions. At the same time Hohenberg (1967) succeeded in ruling out the existence of a conventional superfluid or superconducting ordering in one and two dimensions. It was also shown by Coleman (1972), that โthere are no Goldstone bosons in 2Dโ, which is equivalent to saying that there is no LRO in 2D.
A rigorous proof of the Mermin-Wagner-Hohenberg results exploits the Bogoliubov and Schwartz inequalities (Appendix I) and leads to the following result for the average occupation number of $`๐ค`$ states
$$a_k^+a_kn_k\frac{1}{2}+\frac{mTn_0}{k^2n}.$$
(48)
Here $`n_0`$ is a condensate density and $`n`$ is a total density. It is clear now that the appearance of the condensate (macroscopic occupation of a single state) in 2D for finite temperatures fails due to the mathematical fact that the function $`k^2`$ is not integrable at small momenta in two-dimensional $`k`$ space. Physically, the long-range thermal fluctuations prevent the formation of a coherent condensate.
The same result can be obtained from the infrared asymptote of the one-particle Greenโs function at zero frequency
$$G(๐ค,0)\frac{n_0m}{n_sk^2},$$
(49)
and was first derived by Bogoliubov (1961). The derivation of the asymptotic behaviour (49) in a functional integrals approach can be found in Popov (1983). Since the Greenโs function defines the average number of particles with momentum $`k`$, it is readily seen we arrive to the same result (48). The statement, that the condensate does not appear in a 2D interacting Bose system at any finite temperature, is also known as the Bogoliubov $`k^2`$ theorem.
We have already mentioned that in the context of modern condensed matter theory the absence of the LRO in 2D is discussed in terms of general properties of the XY models. A respective direction or a phase of the $`d`$-dimensional XY order parameter is specified by an angle $`\theta `$. The variance in the fluctuation of the order parameter phase is given by the integral
$$\theta ^2(๐ซ)\frac{T}{\rho _s}\frac{d^dq}{(2\pi )^d}q^2=\frac{T\mathrm{\Lambda }^{d2}}{\rho _s(d2)},$$
(50)
where $`\mathrm{\Lambda }`$ is the wave number cutoff Chaikin and Lubensky (1995). It is readily seen that $`d=2`$ is the critical dimension of the XY universality class and the fluctuations destroy long-range order in the 2D $`XY`$ model in accordance with the conclusions of Bogoliubov, Mermin, Wagner and Hohenberg. Quasi-long-range order, discussed in the previous section is nevertheless possible in 2D.
In case of a trapped gas the Bogoliubov-Mermin-Wagner-Hohenberg (BMWH) theorem rules out BEC in 2D in the interacting system (see Mullin (1997)). However, the question arises if one can actually apply BMHW theorem to a system confined within a harmonic potential. Is it still possible to unambiguously rule out the condensate formation in 2D atom traps? The applicability of the BMWH theorem to the inhomogeneous case requires careful consideration, for the Bogoliubov-Hohenberg inequality was derived assuming an infinite uniform system. In this approximation, many special features of practically realized condensates, such as their formation in real space, are excluded.
An alternative version of the Hohenberg inequality, suitable for the experimentally realizable Bose systems, has been recently proposed by Fischer (2002, 2005). Taking the dimension of the trap to be an experimentally controlled parameter, Fischer addressed the issue of a spatially localized Bose condensate, with the question in mind of how far one could โstretchโ the 3D condensate cloud before the coherence will be destroyed. Fischer derived an inequality, which controls the size of the smallest possible condensate for a given condensate and density profile. In Appendix I we briefly sketch the underlying concepts of his derivation.
The resulting inequality reads
$$\frac{nn_0}{n_0}\frac{2\pi R_c^2}{n\lambda _{dB}^2}C(๐ค)\frac{1}{2n_0}(1|\psi _0(๐ค)|^2/V_0)),$$
(51)
where $`\psi _0(๐ค)`$ is the condensate wave-function, $`R_c`$ is the effective radius of the condensate wave function (effective radius of the curvature of the condensate)
$$R_c=\left(V_0/nd^dr\psi _0(r)[\mathrm{\Delta }_r\psi _0^{}(r)]n(r)\right)^{1/2}$$
(52)
and
$$C(๐ค)=\left|d^dr|\psi _0|^2exp(ikr)\psi _0(k)d^dr\psi _0^{}(r)|\psi _0(r)|^2\right|^2.$$
(53)
Note, that only requirement on the Hamiltonian of the system, that is needed to derive the inequality (51) is that it should not contain any explicit velocity dependence in the interaction and external potentials.
Since in 2D $`R_c^2`$ scales as $`n`$, this case can be considered as marginal and the condensate still can emerge even in an interacting system. This is because the usual log divergences inherent for 2D are cut off by a trap. The inequality (51) is a geometrical equivalent of the Bogoliubov-Hohenberg inequality, since it gives the lower bound for the ratio of the effective radius of the condensate to the de Broglie wavelength $`\lambda _{dB}`$. The second term of rhs (51) can be used to obtain an upper limit on the possible condensate fraction as a function of temperature. Concrete examples of the application of (51) to quasi-1D systems are given in the paper by Fischer (2002).
One can also approach the problem of the condensate formation by directly analyzing the phase fluctuations of the order parameter (for a review see Hellweg *et al.* (2001)). Phase fluctuations are caused by the thermal excitations and are always present at finite temperatures. Note, that at very low temperatures density fluctuations in equilibrium are suppressed due to their energetic cost and can therefore be ignored. This assumption is not valid in the vicinity of a vortex core, but at very low temperatures, the vortex formation is negligible.
As an aside, we mention that the concept of phase in quantum systems, introduced by Dirac as a canonical conjugate observable to the number operator $`\widehat{n}`$, remains a controversial issue in certain circles. Formally it is known, that if $`\widehat{n}`$ is an operator with a purely discrete spectrum (which is always true for the number operator), then there can exist no operator $`\widehat{\theta }`$ such as the commutator $`[\widehat{n},\widehat{\theta }]=i\widehat{1}`$ holds. Different versions of phase-related operators have been constructed in order to overcome this difficulty (see for example the review by Carruthers and Nieto (1968) and the textbook on Quantum optics by Mandel and Wolf (1995)). Alternatives to conventional symmetry breaking approaches have even been proposed (see the paper of Stenholm (2002) and references therein ). An intriguing suggestion that the interference patterns of two atomic condensates can be explained without ever evoking the notion of phase was put forward by Javanainen and Yoo (1996).
In the present article we adopt the โconventionalโ and certainly more convenient approach, according to which the bosonic field operator takes on the form
$$\psi (๐ซ)=\sqrt{n_0(๐ซ)}\mathrm{exp}[i\theta (๐ซ)]$$
(54)
for the large number of particles. Here $`\theta (๐ซ)`$ is the operator of the phase fluctuations and $`n_0(๐ซ)`$ is the condensate density at $`T=0`$.
To proceed with calculations it is convenient to expand the phase operator in terms of the creation and annihilation operators for Bogoliubov quasiparticles (see Shevchenko (1992))
$$\widehat{\theta }(๐ซ)=\frac{1}{2\sqrt{n_0(๐ซ)}}\underset{k}{}\left((u_k+v_k)\widehat{a}_k+(u_kv_k)\widehat{a}_k^+\right),$$
(55)
where $`a_k`$ is the annihilation operator for the Bogoliubov excitation with energy $`ฯต_k`$, and $`u_k`$, $`v_k`$ are excitation functions, determined by a bosonic equivalent of the Bogoliubov-de Gennes equations (for a general reference see the book of de Gennes (1966) ). Expression (55) can be obtained in the formalism of Bogoliubov transformation generalized to an inhomogeneous case.
Phase fluctuations in a quasi 2D system can be analyzed within the formalism of the one-particle density matrix (see the works by Petrov *et al.* (2000, 2001))
$$\psi ^+(๐ซ)\psi (0)=\sqrt{n_0(๐ซ)n_0(0)}\mathrm{exp}[(\mathrm{\Delta }\theta (๐ซ))^2/2)]$$
(56)
One should mention that the quasi two-dimensionality of the system implies that the scattering of particles acquires a 3D character, while the kinetic properties of the gas remain two-dimensional.
It is clear from (55) that the estimation of the phase fluctuations $`\mathrm{\Delta }\theta (๐ซ)^2`$ requires a knowledge of the Bogoliubov quasiparticle spectrum in the inhomogeneous systems (see papers of Stringari (1996) and รnberg *et al.* (1997) and references in papers by Petrov *et al.* (2000, 2001)). This spectrum is discrete for $`T\mu `$ and for $`T\mu `$ one can use the local density approximation. In the Thomas-Fermi regime for $`T\mu `$ one obtains the following approximation
$$\mathrm{\Delta }\theta (๐ซ)^2T\mathrm{ln}(R/\lambda _{dB}).$$
(57)
Note, that (57) does not depend on precise expression for the repulsive coupling constant.
From (57) one can estimate the characteristic radius $`R_\theta `$ of phase fluctuations (the characteristic length at which phase changes by $`2\pi `$) to be $`R_\theta \lambda _{dB}exp(T_\theta /T)`$ with $`k_BT_\theta =N(\mathrm{}\omega _{})^2/\mu `$. We thus arrive at the conclusion that at low temperatures $`TT_\theta `$ the characteristic radius of the phase fluctuations is larger than the size of the trap $`R_\theta R_{}`$, so a true condensate exists. The emergence of a true condensate is attributed to the weakening of phase fluctuations induced by a trap, which introduces a low momenta cut off into the excitations in the system. At higher temperatures, $`TT_\theta `$ the system is characterized as a quasicondensate ($`R_\theta R_{}`$).
The crucial effect of a trap for 2D Bose gases was also emphasized by Ho and Ma (1999). They pointed out that long wavelength quantum fluctuations will be partially suppressed due to the gapped spectrum of collective modes Stringari (1996) and off-diagonal order will survive in 2D.
Since there is an experimental evidence in support of BEC existence in 2D, the discussion is not yet closed. Quantum Monte Carlo simulations for bosons in a two-dimensional harmonic trap do indeed show that a significant fraction of the particles is still present in the lowest state at low energies Heinrichs and Mullin (1998).
### IV.3 Popovโs approach
In this section we consider the way Popov (1983) generalized the field-theoretical methods, developed by Beliaev to finite temperatures. It is curious that the method, suggested by Popov in 1965, is conceptually very similar to renormalization-group approach, successfully applied in the 1970th to phenomena, unaccessible to perturbative methods, such as Kondo effect Hewson (1993).
As usual one starts with the introduction of the temperature Greenโs function
$`G(x,\tau ;x^{},\tau ^{})`$ $`=`$ $`<\psi (x,\tau )\overline{\psi }(x^{},\tau ^{})>`$ (58)
$`=`$ $`{\displaystyle \frac{e^S\psi (x,\tau )\overline{\psi }(x^{},\tau ^{})๐\psi ๐\overline{\psi }}{e^S๐\psi ๐\overline{\psi }}},`$ (59)
where $`S`$ is the classical action of the Bose gas
$$S=_0^\beta ๐\tau d^3x\overline{\psi }(x,\tau )_\tau \psi (x,\tau )_0^\beta ๐\tau H(\tau )$$
(60)
and
$`H(\tau )={\displaystyle d^3x\overline{\psi }(x,\tau )(\frac{^2}{2m}\mu )\psi (x,\tau )}+`$ (61)
$`+{\displaystyle \frac{1}{2}}{\displaystyle d^3xd^3yU(xy)\overline{\psi }(x,\tau )\overline{\psi }(y,\tau )\psi (y,\tau )\psi (x,\tau )}.`$
Next step is the construction of the perturbation theory and corresponding diagrams arising from integrals of the type (59), by performing the usual โtrickโ of separating out the condensate operators (27). However, in case of Bose system the perturbation series converges very poorly for small momenta and frequencies. In other words, the infrared asymptote of the Greenโs function is singular. In order to avoid these difficulties, Popov suggested following modifications: the bosonic fields $`\psi `$
$$\psi (x,\tau )=\sqrt{\frac{1}{\beta V}}\underset{k,\omega }{}exp[i(kx\omega \tau )]a(k,\omega )$$
(63)
is divided into a short wavelength โfastโ component $`\psi _1`$ and a long wavelength โslowโ component $`\psi _0`$ ($`\psi (x,\tau )=\psi _0(x,\tau )+\psi _1(x,\tau )`$) (see Fig.5). The momentum $`k_0`$ which separates the slow modes from the rapidly oscillating modes depends on the particular Bose system and only its order of magnitude can be estimated. Introduction of $`k_0`$ removes the divergences at small momenta, regularizing the perturbation theory.
A method of successive integration, first over the โrapidโ and then over the โslowโ fields, is then applied, using different schemes of perturbation theory at different stages of the integration (see chapter 4 in Popov (1983)). The fast modes โseeโ the slow modes as an effective condensate (โbareโ condensate according to Popov) with a superfluid density $`\rho _0=|\psi _0|^2`$. In Appendix II we give a succinct derivation of main Popovโs results.
This method of subsequent integration, developed by Popov, allows to estimate the low temperature asymptotic behaviour of the one -particle Greenโs function, and to derive a power-law decay of $`G(x,y)|xy|^\alpha `$ for $`|xy|\mathrm{}`$ (in 1D and 2D) rather than the exponential decay that occurs at high temperatures. In 2D, as we have mentioned in IV.1 this signals the development of topological LRO at low temperatures.
The analysis is based on the t-matrix description of the effective interactions, and the key property of the 2D t-matrix is that at low energies it vanishes, and at high energy cut-off the t-matrix diverges (see Appendix II). This results in an extremely small critical temperature
$$T_c\frac{\mu \mathrm{ln}(ฯต_0/\mu )}{4\mathrm{ln}\mathrm{ln}(ฯต_0/\mu )},$$
(64)
where $`ฯต`$ is a high energy cut-off and $`\mu `$ is a chemical potential. Bear in mind, that this is a mean-field derivation and the condition for the superfluid transition was assumed to be $`\rho =\rho _n`$, because Popov (as well as Berezinskii (1970, 1971)) thought that at the critical temperature $`T_c`$, the superfluid density vanishes.
The applicability of Popovโs mean-field description is based on the assumption of a very small exponent $`\alpha `$. For large $`\alpha `$ the probability of creation of quantum vortices becomes big and even this modified perturbation theory is invalid (see also discussion in IV.4 and the โcorrectedโ many-body mean-field theory in IV.5). The applicability of the Popovโs t-matrix description and the diluteness condition, derived by Fisher and Hohenberg is the main subject of Section IV.4.
### IV.4 Diluteness condition and validity of $`t`$-matrix approximation
We have already discussed that the perturbative treatment of the dilute weakly interacting Bose gas amounts to replacing the real potential by an effective two-particle $`t`$-matrix, obtained by summing up all ladder diagrams. From this point of view the diluteness condition determines the range of validity of the $`t`$-matrix approximation.
An explicit form for the diluteness condition of 2D interacting Bose gas at finite temperatures was first introduced by Fisher and Hohenberg (1988). They pointed out that singularities inherent to 2D systems (vanishing of scattering $`t`$ matrix at zero temperature and classical divergence of phase fluctuations) might lead to drastic modifications of the usual dilute gas expansion.
As we have discussed (see section III.2) at zero temperature in 2D the diluteness condition $`na^21`$ is replaced by
$$\frac{1}{\mathrm{ln}(1/na^2)}1.$$
(65)
Popovโs theory can be used to demonstrate that at finite temperature, the above condition (65) is replaced by an even more stringent inequality Fisher and Hohenberg (1988)
$$\mathrm{ln}\mathrm{ln}\frac{1}{na^2}1.$$
(66)
Fisher and Hohenberg (FH) provided a heuristic derivation of this result, based on the Bogoliubov quasiparticle picture. Their analysis is based on the simple observation that the usual Landau quasiparticle formula for the superfluid density
$$\frac{\rho _s}{\rho }=1\frac{\beta }{\rho d}\frac{d^dk}{(2\pi )^2}k^2\frac{e^{\beta \xi _k}}{(e^{\beta \xi _k}1)^2},$$
(67)
where $`d`$ is the dimension, does not have any singularities for $`d=2`$, except in case when the chemical potential is small ($`\mu `$ is introduced in (67) via the Bogoliubov quasi-particle spectrum $`\xi _k^2=n_0U_0k^2/m+k^4/4m^2\mu k^2/m+k^4/4m^2`$. The validity of this approximation for $`\mu `$ is discussed in Beliaev (1958a, b)). By introducing the infrared cut off ($`k_0\sqrt{\mu }`$) via the ansatz
$$\mu \frac{n}{|ln(a^2\mu )|},$$
(68)
the regularization of the integral of Eq.(67) can be achieved, and one arrives at Popovโs equations for the โsuperfluidโ and โnormalโ densities (91)-(93).
The analysis of the temperature dependence of the superfluid density allows one to separate out three characteristic regimes: (i) the low temperature region, the physics of which is defined by phononic behaviour of the quasiparticles, leading to a superfluid density which depends on temperature as $`(1\alpha T^3)`$; (ii) a free particle region, where $`\rho _s`$ behaves linearly with temperature; (iii) and a critical region, determined by the fluctuations around the critical temperature $`T_c`$ ( $`\rho _s`$ vanishes at $`T_c`$) (see Fig.6.)
The diluteness criterion is determined by the condition that the critical region is small enough so that all three regimes can be well separated. The width of the first regime is in fact given by Schickโs small parameter (65), while the size of the critical region is characterized by the double log (66). The problem however, is that for all practically relevant situations, even for very small $`na^2`$, Fisher and Hohenbergโs small parameter $`1/lnln(1/na^2)`$ is still orders of magnitude greater than $`1/ln(1/na^2)`$. This means that in practice the critical region associated with Kosterlitz-Thouless transition is so large, that mean-field based approaches do not give any reliable results. Note, that the double log result was also reproduced by Fisher and Hohenberg in a more accurate way within a renormalization group treatment of the same problem. They have also estimated the superfluid transition temperature, which reads
$$T_c\frac{2\pi n}{m\mathrm{ln}\mathrm{ln}[1/(na_s)^2]}.$$
(69)
The results of FH work Fisher and Hohenberg (1988) have been confirmed in other approaches, see for example virial expansion of a dilute Bose gas by Ren (2004) or RG analysis by Kolomeisky and Straley (1992a, b) and by Crisan *et al.* (2001). Pieri *et al.* (2001a, b) demonstrated by analyzing the normal state by standard diagram technique that the transition temperature (69) appears as a lower bound for the validity of the $`t`$-matrix as a controlled approximation for the dilute Bose gas.
The FH diluteness condition (66) is extremely stringent, and if straightforwardly applied to experimentally relevant situations Gรถrlitz *et al.* (2001); Rychtarik *et al.* (2004), would mean that the systems observed to undergo a BEC phase transition in 2D are not actually dilute, and could never be so. This line of reasoning motivated Liu and Wen Liu and Wen (2002) to come up with an exotic alternative scenario involving a two-dimensional strongly-correlated spin liquid.
The extreme conclusions drawn from the FH diluteness criterion are nevertheless related to the general drawbacks of the Popov approximation. We will see in the next section, that in a more realistic model, which takes into account interactions in a self-consistent way, the diluteness condition becomes much weaker. Moreover, in view of our previous discussion about the inapplicability of arguments based on homogeneous systems in the thermodynamic limit to trapped gases, it would seem that the FH diluteness requirement is not really relevant for the experimental situation of the Bose gas in a magnetic trap.
### IV.5 Other approaches: RPA, Many-body $`t`$-matrix, Monte Carlo
In this section we review a range of diagrammatic approaches that have built upon the early RPA and $`t`$-matrix approximation in order to improve the description of the superfluid of BKT transition and also the numerical methods, which allow to directly probe the critical region of the 2D transition.
The first finite temperature generalization of Bogoliubov random phase approximation (III.1) was introduced by Tserkovnikov (1964) (english version Tserkovnikov (1965)). His concern was to calculate the finite-temperature correction to the condensate density in 3D dilute Bose gases with weak interactions. Tserkovnikov assumed that the average single particle kinetic energy is small, compared to the potential energy for all temperatures below $`T_c^{3D}`$. He also remarks that his approximation does not meet the Landau superfluidity criterion and that more precise equations should be sought in future work.
The RPA method was further developed in the papers of Szepfalusy and Kondor (1974), whose main interest was the investigation of dynamics of the second order phase transition. Around the same time a large-$`N`$ approach was applied to the Bose gas by Abe (1974) and Abe and Hikami (1974), who calculated the dynamical scaling for one-particle Greenโs function up to $`O(1/N)`$. Here the idea of the large-$`N`$ approach is to expand the number of independent components of the Bose field from unity to $`N/2`$ using $`1/N`$ as an expansion parameter. To produce a controlled large $`N`$ limit, the interaction strength is scaled to be of order $`1/N`$.
The RPA large-$`N`$ method has been applied to 2D Bose gas by Nogueira and Kleinert (2006). The interaction in their approach is approximated by a 2D coupling constant, derived in the $`t`$-matrix approximation, considered by Popov and Schick (see sections III.2, IV.3 and IV.4). It is however known that in a large-$`N`$ approach one can not simultaneously account for both particle-hole channel (RPA) and particle-particle channel in a well-controlled fashion. Nevertheless, the authors claim that the diluteness condition leads to the appearing of $`t`$-matrix diagrams first, while the next class of diagrams are those from particle-hole channel Nogueira and Kleinert (2006). This approximation results in the Bogoliubov quasiparticle dispersion containing a log correction due to low dimensionality $`\xi _k=\sqrt{ฯต_k^2+2g_{2D}nฯต_k\left[1\frac{Tm}{\pi n}ln(ka)\right]}`$, so that excitations in the system exhibit a roton-like minimum. Note that the excitation spectrum is calculated under the assumption that one-particle Greenโs function and density correlators share the same poles (this property was derived by Hohenberg and Martin (1965) in case of 3D condensed Bose system). Would be interesting to check if these RPA results are confirmed in other approaches.
We now proceed to a discussion of the various generalizations of the two-body $`t`$-matrix approach. Though simple and elegant, the perturbative two-body $`t`$-matrix approach does have its drawbacks. The main problem is related to its inability to properly describe the critical region in low dimensions. For example, the $`t`$-matrix method does not predict the Nelson-Kosterlitz universal jump in the superfluid density. In 3D the two-body $`t`$-matrix approach leads to a first order phase transition for the condensate density, which is the consequence of non-self-consistency of this first order perturbative approximation (see also Griffin (1988); Lee and Yang (1958); Reatto and Straley (1969)).
Many of these problems can be solved if the many-body corrections, arising due to the surrounding gaseous medium are taken into account. This is the key idea in the โmany-body $`t`$-matrix approximationโ (see the comprehensive review by Shi and Griffin (1998), low-dimensional systems within many-body $`t`$-matrix approach are analyzed in the papers by Stoof and Bijlsma (1993); Andersen *et al.* (2002); Al Khawaja *et al.* (2002), for Hartree-Fock -Bogoliubov study of a two-dimensional gas see recent works by Gies *et al.* (2004); Gies and Hutchinson (2004); Gies *et al.* (2005)).
Since the โmany-body $`t`$-matrixโ methods are extensively discussed in the literature, here we restrict ourself to a brief description of the method, providing all relevant references. The Bogoliubov-Hartree-Fock (BHF) approximation (see Griffin (1996) and the analysis of excitations in a trapped 3D gas paper by Hutchinson *et al.* (1997)) has a Heisenberg equation of motion for a Bose field operator of the kind (43) as a starting point
$`i\mathrm{}{\displaystyle \frac{}{t}}\psi (r,t)`$ $`=`$ $`({\displaystyle \frac{\mathrm{}^2}{2m}}^2+V_{ext}(๐ซ)\mu )\psi (r,t)+`$ (70)
$`+`$ $`g\psi ^+(r,t)\psi (r,t)\psi (r,t).`$
A short-range interaction is assumed among the atoms $`U(๐ซ๐ซ^{})=g\delta (๐ซ๐ซ^{})`$. Treating the interaction term in this equation in the self-consistent mean-field approximation, one arrives at the equation
$`({\displaystyle \frac{\mathrm{}^2}{2m}}^2+V_{ext}(๐ซ)\mu )\varphi (r)+`$
$`+g(n_0(r)+2n^{}(r))\varphi (r)+gm^{}(r)\varphi ^{}(r)=0,`$ (71)
where $`n_0`$ is the condensate density, $`n^{}(r)=\psi ^+(r)\psi ^{}(r)`$ and $`m^{}(r)=\psi ^{}(r)\psi ^{}(r)`$ (anomalous average). In order to describe excitations in the system one should also write down the equation of motion for $`\psi ^{}`$
$`i\mathrm{}{\displaystyle \frac{}{t}}\psi ^{}(r)`$ $`=`$ $`({\displaystyle \frac{\mathrm{}^2}{2m}}^2+V_{ext}(๐ซ)\mu )\psi ^{}(r,t)+`$ (72)
$`+2gn(r)\psi ^{}(r,t)+gm(r)\psi ^+(r,t)`$
(with $`n(r)=n_0+n^{}`$ and $`m(r)=\varphi ^2+m^{}`$). Equations (71) and (72) correspond to the Bogoliubov-Hartree-Fock approximation: a bosonic analog of the finite-temperature Bogoliubov-de Gennes equations. These equations can also be re-expressed in terms of a Greenโs function formalism. Note, that the appearance of anomalous averages in BHF formalism (which are not present in Popov approach), leads to a gap in the quasiparticle excitation spectrum.
Many body effects are also effectively treated in a variational method, applied to dilute Bose gases by Bijlsma and Stoof (1997), and many-body $`t`$-matrix methods which have been significantly improved in recent years (see Andersen *et al.* (2002); Al Khawaja *et al.* (2002)). A time dependent BHF approximation has recently been developed by Proukakis *et al.* (1998). Here, the authors Proukakis *et al.* (1998) claim that the pseudopotential approximation $`U(๐ซ๐ซ^{})=g\delta (๐ซ๐ซ^{})`$ should be imposed only after the effective interaction is expressed in terms of many-body $`t`$-matrix. Both approaches, the one used by Griffin et al. and another, developed by Stoof and collaborators, are qualitatively similar, in that they treat interactions in many-body $`t`$-matrix approach, but they differ in some details, for instance in selecting out the important diagrams.
Let us now discuss some of the results of the many-body $`t`$-matrix method for 2D systems. More than a decade ago Stoof and Bijlsma (1993) demonstrated that the infrared divergences, appearing in the two-body $`t`$-matrix treatment, can be elegantly eliminated when the effects of surrounding gas are taken into account. With this approach the universal jump in superfluid density, predicted by Nelson and Kosterlitz can be reproduced. A weaker diluteness condition, namely that of Schick (65), defines the applicability of this many-body $`t`$-matrix approximation, one that is satisfied experimentally in systems, such as spin polarized atomic hydrogen, absorbed on <sup>4</sup>He surface Stoof and Bijlsma (1993).
The conclusions of Petrov *et al.* (2000, 2001) have been confirmed in recent investigations by Andersen *et al.* (2002) and Gies *et al.* (2004); Gies and Hutchinson (2004); Gies *et al.* (2005). The theory of Andersen *et al.* (2002), free of infrared divergences in all dimensions, allows for calculation of the density profile of a (quasi)-condensate cloud of a gas for any aspect ratio of the trap (within local density and Thomas-Fermi approximation). At very low temperature, depending on the trapping geometry, the presence of a true condensate in the equilibrium state is found. Hutchinson and coworkers Gies *et al.* (2004); Gies and Hutchinson (2004); Gies *et al.* (2005) also โseeโ within their HFB approach a macroscopic occupation of the ground state at low temperatures, implying the presence of a condensate state.
To conclude, the presence of the trap appears to stabilize the condensate against the long-wavelength fluctuations and the BEC state can form at finite, though very low temperatures, when the discrete nature of the energy spectrum is taken into account.
Most reliable description of the 2D Bose gas to date is provided by Monte Carlo simulations Kagan *et al.* (2000); Prokofโev *et al.* (2001); Prokofโev and Svistunov (2002), because it allows to study the critical region of the BKT transition, which is effectively very large and therefore unaccessible to perturbative methods. The numerical analysis is simplified, by the fact that the critical properties of all XY models are the same (see IV.2). It suffices therefore to study the classical $`|\mathrm{\Psi }|^4`$ model on the lattice within a Monte Carlo algorithm.
Consider, for instance, the temperature-dependence of the particle density in the critical region of the BKT transition, which follows from perturbative analysis Popov (1983); Kagan *et al.* (1987); Fisher and Hohenberg (1988) of weakly-interacting system
$$n=\frac{mT}{2\pi }\mathrm{ln}\frac{C}{mU_{eff}},$$
(73)
where $`U_{eff}`$ is an effective interaction, proportional to $`f_0^{2D}`$ (37), and $`C`$ is the constant, which is not possible to evaluate within perturbative expansion in powers of $`U_{eff}`$ Prokofโev *et al.* (2001). Monte Carlo estimation gives $`C=380\pm 3`$; this large value of $`C`$ makes it virtually impossible to reach the limit of small $`U_{eff}`$ for weakly-interacting system.
At transition we obtain an an accurate microscopic expression for the critical temperature of BKT transition Prokofโev *et al.* (2001); Prokofโev and Svistunov (2002).
$$\frac{T_{KT}}{n}=\frac{2\pi }{m\mathrm{ln}(C/mU_{eff})}.$$
(74)
It is interesting to compare this density $`n`$ to the quasi-condensate density $`n_q`$ and superfluid density $`n_s`$ in the critical region. It turns out, that $`n_q/n`$ is of order of unity, unless $`mU_{eff}`$ is exponentially small, while the ratio $`n_q/n_s`$ is of order of 2, which means that superfluid density is substantially smaller than quasi-condensate density at the transition.
The temperature behaviour of various densities, obtained in Monte Carlo procedure, can be used for checking whether RG and perturbative approaches essentially overlap. Indeed, Monte Carlo simulations have been able to capture the crossover between the mean-field behaviour and the critical fluctuation region described by the KT transition Prokofโev and Svistunov (2002). Prokofโev and Svistunov (2002) show that this crossover is characterized by a universal ratio of the superfluid and quasi-condensate density. One can also see that the conventional mean-filed result $`n_s/n1T/T_{KT}`$ is not valid anywhere, while the modified mean-field theory introduced by Prokofโev and Svistunov (2002) can predict accurately the behaviour of the quasi-condensate density up to $`T_{KT}`$.
### IV.6 Breathing modes of 2D systems
At the end of this chapter we consider a universal property of a two-dimensional gas with a contact interaction, confined in a harmonic potential. Pitaevskii and Rosch (1997) predict that such a system develops oscillations or breathing modes, which can be probed experimentally or in simulations and thus can serve as a practical criterion of the two-dimensional nature of a system.
The appearance of breathing modes is related to a hidden โLorentzโ symmetry inherent to any two-dimensional Hamiltonian of the following general form
$$H=H_0+H_{ext}$$
(75)
where
$$H_0=\underset{i}{}\left(\frac{1}{2m}\mathrm{}_i\right)+\underset{i<j}{}U(๐ซ_i๐ซ_j),$$
(76)
and $`H_{ext}=_i\frac{1}{2}m\omega _0^2r_i^2`$ is a harmonic potential.
It is readily seen that $`H_0`$ is scale invariant in case of a local 2D interaction
$$U(๐ซ_i๐ซ_j)=\frac{g}{2}\delta ^2(๐ซ_i๐ซ_j),$$
(77)
(in fact it is scale invariant for any potential with a property $`U(l๐ซ)=U(๐ซ)/l^2`$). The presence of a trap breaks the scale invariance of $`H_0`$. Note, that in principle the scaling invariance of the Hamiltonian $`H_0`$ is broken in 2D, because then scattering phase shift is energy dependent due to the logarithmic dependences characteristic of two dimensions (phase shift is proportional to the coupling constant $`g^{2D}`$ or to $`1/ln(ka)`$). The energy-dependent phase shift signals the breaking of scale invariance at the quantum level Cabo *et al.* (1998). But this symmetry breaking is explicit and is not attributed to any phase transition physics.
In spite of the breaking of scale invariance, because of a special property of the harmonic oscillator, a powerful spectrum generating symmetry still exists. That can be seen from the commutator $`[H_{ext},H]=i\omega _0^2Q`$, where $`Q=1/2_i(๐ฉ_i๐ซ_i+๐ซ_i๐ฉ_i)`$ is the generator of scale transformations. One can check that $`[Q,H_0]=2iH_0`$ and $`[Q,H_{ext}]=2iH_{ext}`$. These results can be formulated within the well-known algebra of SU(1,1) or SO(2,1) symmetry group, i.e. the 2D Lorentz group.
Starting from the lowest energy state $`ฯต_0`$ one can produce higher order states with energies $`ฯต_0+2n\omega _0`$ ($`n=1,2,\mathrm{}`$) by applying one of the SO(2,1) group generators $`L^+=(L_1+iL_2)/\sqrt{2}`$ where $`L_1=(H_0H_{pot})/2\omega _0,L_2=Q/2`$ (the correspondent annihilation operator is $`L^{}=(L_1iL_2)/\sqrt{2}`$). The excitations with the energies $`2n\omega _0`$ are associated with the breathing, or pulsating modes of the system.
As an example, the authors Pitaevskii and Rosch (1997) considered the classical Gross-Pitaevskii equation and predicted the existence of undamped breathing modes in the condensate. The appearance of transverse breathing modes with a frequency equal to an integer multiple of the trap oscillation frequency was observed experimentally in an elongated condensate of $`{}_{}{}^{87}Rb`$ atoms Chevy *et al.* (2002). Numerical simulations (exact diagonalization) seem to indicate the existence of dipole or breathing modes in a 2D system even for relatively small number of atoms Haugset and Haugerud (1998).
BHF study of a 2D Bose gas by Gies et al. Gies *et al.* (2004); Gies and Hutchinson (2004); Gies *et al.* (2005) has also showed that at low temperatures the frequency of the lowest lying excitation ($`n=0`$ mode) is precisely $`2\omega _0`$, independent of the interaction strength. At high temperatures the frequency of this mode shifts to a lower frequency region, being modified by the addition of the potential from the static thermal cloud.
Nevertheless, it is important to note that a $`\delta `$-function is not well-defined in two dimensions due to logarithmic UV divergences Pitaevskii and Rosch (1997), that are cut off by the finite range of interaction, whether it is a small or large effect, should be investigated.
## V Conclusions and open questions
We have surveyed a number of theoretical issues arising in the field of a weakly-interacting uniform or confined in a trap dilute Bose system at low temperatures in 2D. The underlying physics of such a system depends on the size of the system, degree of its inhomogeneity, and temperature.
If the system is very large and uniform, one might expect realization of the BKT transition, characterized by the presence of a topological order below the critical temperature $`T_{KT}`$ down to zero temperature when the true long-range order (BEC) forms. Perturbative approaches, based on a low density approximation and point-like or short-range interactions, surveyed in this Colloquium, are not really suited to a description of the vortex excitations of the ordered phase of BKT transition, however, these methods provide a good description of many physical properties. For example, a modified version of the mean-field theory, the โmany-body $`t`$-matrix approachโ is able to capture the second order nature of the phase transition in 3D and the Nelson-Kosterlitz universal jump of superfluidity in 2D.
In experiment the low temperature regime of a quantum gas is achieved by confining the atomic system in an external potential. The proper description of a practically realized system requires therefore the inclusion of the inhomogeneities, introduced by a trap. In 3D the effect of the trap is not very pronounced, and it is possible to calculate the correction to the critical temperature due to interactions perturbatively Arnold and Tomasik (2001), while in the uniform case the RG study is required Ledowski *et al.* (2004) (for details see the book by Pitaevskii and Stringari (2003)).
In 2D, as we have seen in the case of a gas without interactions, the presence of the trap dramatically modifies properties of the system (density of states), so that BEC becomes possible at finite temperature. The inclusion of interactions into the picture is a complicated task. First of all, the system is inhomogeneous and all previously developed perturbative methods, such as the $`t`$-matrix approach are strictly speaking not suitable for its analysis. In principle, one should solve the many-body scattering problem in a trap, taking into account the discrete spectrum and the finite range of a potential, which is extremely difficult. One can of course consider a simplified problem of a quasi-homogeneous trap and adapt well-studied techniques for that case. As we have seen, the main effects of the trap are indeed captured at least qualitatively within such a scheme.
Intuitively, it is clear, that inhomogeneities would tend to suppress the universal jump of superfluidity, and would rather favor the โtrueโ BEC state at low temperatures in a system with finite number of particles (because the long-range thermal fluctuations will be quenched by a trap).
Very recently Holzmann *et al.* (2005) have analyzed the behaviour of the weakly-interacting trapped system in the thermodynamic limit within local density approximation (LDA). They have shown, that although the universal Nelson-Kosterlitz jump is indeed not present, the system does undergo BKT transition at the temperature, somewhat lower than $`T_c^{2D}`$ (15).
In a paper by Simula *et al.* (2005) the authors have predicted that both a BE condensed phase and a KT superfluid phase, separated by a first order transition, will be present in a 2D trapped gas. The authors arrive at their conclusions by comparing the Helmholtz free energy of the ground state, characterized by a condensate wave-function and an excited state, containing a vortex pair in the ordering field. Based on entropy considerations, they find a critical temperature $`T_c`$ above which a thermally activated transition to the state containing vortex-pair excitations becomes favourable. These conclusions require, however, solid confirmations from both theory and experiment.
A challenging problem is however the actual estimation of the BEC transition temperature in a 2D interacting system and the relationship of the quasi-condensate state to the KT vortex-pair plasma state, depending on the geometry of the system and number of particles. It is also important to develop numerical methods adapted for experiment. One of the promising developments in this direction has been made by Davis *et al.* (2001a, b, 2002), who have proposed a projected Gross-Pitaevskii equation formalism, which allows the efficient investigation of finite temperature properties of the equilibrium condensate state - even in the region where the Bogoliubov theory fails. An extension of this method to harmonically trapped condensates has been considered in the recent paper by Blakie and Davis (2005). Just this year Simula and Blakie (2006) have analyzed the 2D Bose gas by classical field methods, adopted from quantum optics. They have demonstrated that two distinct superfluid phases, separated by thermal vortex-antivortex pair creation, exist in experimentally producible quasi-2D Bose gas. Simula and Blakie (2006) provided a strong evidence that a strange (โzipperโ) interference pattern observed in a recent experiment by Stock *et al.* (2005) can be explained by a presence of a vortex excitations in an experimental system.
One can thus tentatively identify the following phases in a trapped 2D Bose gas (Fig.7):
* phase I: low -temperature true BEC phase;
* phase II: KT vortex-antivortex pair superfluid, or quasi-condensate, or condensate with a fluctuating phase, at the transition $`T_{KT}`$ superfluid universal jump is almost suppressed;
* phase III: critical region of the BKT transitiona; vortex pair dissolve, above $`T_{KT}`$ vortices are unbound and free;
* phase IV: above the mean-field temperature it is a normal fluid, no local order parameter or vortices exist.
The phase diagram of a real system will depend on many factors, as we stressed at the beginning of this Chapter. Many aspects of the diagram, depicted on Fig. 7 require careful investigation, and reliable confirmation from both theory and experiment.
In ending we attempt to summarize some of the open questions:
* Nature of superfluid phases of a 2D weakly-interacting Bose gas: what is the nature of the crossover to a superfluid phase? What is the explicit relation between superfluidity and the โquasi-condensateโ state? Is there a crossover to BEC state at low temperatures? and if yes, under which conditions?
* Can we help experimentalists to โseeโ the vortex excitations in the superfluid state, to really identify the KT state? There is a clear need in good vortex detection methods. Can, for example, disorder help us to pin the vortices?
* What are in general measurable physical properties, which delineate between the coherent condensed state and superfluidity? Initial progress in this direction has been already made, Polkovnikov *et al.* (2005) have suggested how to identify the KT transition from experimentally measured interference pattern.
* Can we solve the many-body scattering problem in the trap? Does diluteness of the gas simplify this problem?
* Is local density a good approximation for description of the experimental systems, and if yes, under which conditions?
* Could one justify a large $`N`$ approach which improves on the existing method by incorporating the $`t`$-matrix approximation?
We have also discussed the diluteness condition, derived by Fisher and Hohenberg (1988) under certain conditions of the transition to superfluid state. In a finite-sized experimental system this condition is not really applicable, and one should use the criterion of Schick (65), which can be seen from the analysis of quantum fluctuations of the 2D BEC at zero temperature Petrov *et al.* (2004).
We have considered only low density approximations. When the gas is dense, approaches such as that of Gross and Pitaevskii are not applicable. In such high density regimes, new โstrong couplingโ approaches are required. One of the possible solutions may be the slave-boson approximation, which is valid for hard-core bosons at any density Ziegler and Shukla (1997).
We have not discussed in our Colloquium the role of disorder in the continuum Bose system, though it could be a subject of a separate review and opens up a lot of interesting perspectives. Recent Monte Carlo studies predict, for example, that for strong disorder the system enters an unusual regime, where the superfluid fraction is smaller than the condensate fraction Astrakharchik *et al.* (2002). Weak disorder can be treated within Bogoliubov theory Huang and Meng (1992); Giorgini *et al.* (1994) and the striking result of this is that disorder is more active in reducing superfluidity than in depleting the condensate. These results suggest that the relation of superfluidity and Bose-Einstein condensation require further theoretical and experimental investigation.
One can not but mention the rapidly increasing interest to the cold gases with dipole-dipole interactions, which are responsible for a variety of novel phenomena in ultracold dipolar systems, see, for instance, Santos *et al.* (2003, 2000); Pedri and Santos (2005); Stuhler *et al.* (2005); Fischer (2006) and references therein.
Finally, the problem of measurable quantities is in fact one of the most important in the context of trapped Bose gases. Unlike an electron system, one cannot attach โleadsโ to the system and measure transport properties of a condensate cloud. One of the most pressing practical needs for theorists and experimentalists is therefore the development of controllable new methods to probe the trapped condensate.
## Appendix I (for section IV.2)
Sketch of the derivation of Mermin-Wagner-Hohenberg theorem
One should use the Bogoliubov inequality
$$\frac{1}{2}\{A,A^+\}[[C,H],C^+]T|[C,A]|^2,$$
(78)
where the average $`X=Tr\left(Xe^{\beta H}\right)/Tre^{\beta H}`$, and operators $`A`$ and $`C`$ are such, that the ensemble averages on the lhs of (78) exist. The inequality (78) follows quite straightforwardly from the Schwartz inequality
$$(A,A)(B,B)|(A,B)|^2$$
(79)
where a scalar product is defined by $`(A,B)=T\frac{d\omega }{\pi }\frac{1}{\omega }\chi _{AB}(k\omega )`$, where $`\chi _{AB}`$ is the Fourier transform of the response function $`\chi _{AB}(rt,r^{}t^{})=\frac{1}{2\mathrm{}}[A(rt),B(r^{}t^{})]`$ (see for example the textbook by Forster (1990)).
In the case of a Bose system the derivation of Hohenberg (1967) is based on Bogoliubov and Schwartz inequalities and the $`f`$ sum rule
$$T\frac{d\omega }{\pi }\frac{1}{\omega }\chi _{AA^+}(k\omega )=\frac{k^2n}{m}.$$
(80)
Derivation by Fischer of geometrical analog of the Hohenberg inequality
The bosonic field operator is as usual decomposed into condensate and noncondensate parts
$$\psi (๐ซ)=\psi _0(๐ซ)a_0+\delta \psi (๐ซ).$$
(81)
The key observation of Fischer (2002) is that the Bogoliubov prescription should be applied after implementing the commutation relation
$$[\delta \psi (๐ซ),\delta \psi ^+(๐ซ^{})]=\delta (๐ซ๐ซ^{})\psi _0(๐ซ)\psi _0^{}(๐ซ^{}),$$
(82)
for otherwise, the second term on the rhs of (82), which turns out to be crucial for calculating the condensate fraction correctly, would vanish.
The operators $`A`$ and $`C`$ in the Bogoliubov inequality (78) are chosen to be smeared excitation and total density operators
$$\widehat{A}=d^drf_A(๐ซ)\delta \psi (๐ซ);\widehat{C}=d^drf_C(๐ซ)\delta \rho (๐ซ),$$
(83)
where $`f_A`$ and $`f_C`$ are carefully chosen โsmearing functionsโ ($`f_C(๐ซ)\psi _0^{}(๐ซ)`$, $`f_A(๐ซ)e^{i\mathrm{๐ค๐ซ}}`$). Next, the $`f`$-sum rule analogous to (80), can be derived in coordinate space.
## Appendix II (for the section IV.3)
Popovโs approach
To derive the phase transition curve for a two-dimensional interacting Bose gas, one needs to explore the finite-temperature behaviour of the $`t`$-matrix (32). The Bethe-Salpeter equation (Fig.3) in the Matsubara representation reads
$`\mathrm{\Gamma }(p_1,p_2;p_3,p_4)=U_{๐ค_\mathrm{๐}๐ค_\mathrm{๐}}{\displaystyle \frac{1}{\beta V}}{\displaystyle \underset{q,i\omega _l}{}}U_๐ชG^0(๐ค_\mathrm{๐}๐ช,i\omega _1i\omega _l)`$
$`\times G^0(๐ค_\mathrm{๐}+๐ช,i\omega _2+i\omega _l)\mathrm{\Gamma }(p_1q,p_2+q;p_3,p_4),`$ (84)
where $`\omega _j=2\pi jT`$ is an even Matsubara frequency, and the four-dimensional vector $`p_j(๐ค_j,\omega _j)`$ represents the momentum $`๐ค_j`$ and frequency $`\omega _j`$ of the particle before scattering $`(j=1,2)`$ and after $`(j=3,4)`$. Energy and momentum conservation requires $`p_1+p_2=p_3+p_4`$.
The main contribution to the sum over internal momenta in (84) comes from $`ka^1`$ which is due to the potential, discussed above. Since $`a^1\sqrt{T}\sqrt{\mu }`$, the $`\mu `$ dependence in the Greenโs function can be safely neglected Popov (1983). Consequently, after integrating over frequencies Eq. (84) is reduced to a $`t`$-matrix equation
$`t(p_1,p_3,z)`$ $`+`$ $`{\displaystyle \frac{dp^{}}{(2\pi )^2}U(p_1p^{})\frac{1}{p^2/mz}t(p^{},p_3,z)}=`$ (85)
$`=`$ $`U(p_1p_3)`$
Schematically, the $`t`$-matrix equation can be expressed as $`t_z+UR_zt_z=U`$, with $`R_z=1/(p^2/mz)`$. The operator $`(1+UR_z)`$ can be inverted and we get $`U=t_z(1R_zt_z)^1`$. In this fashion the interaction is eliminated from the $`t`$-matrix equation and we arrive at the Hilbert identity $`t_zt_{z_0}=t_z(R_{z0}R_z)t_{z_0}`$. This last equation is readily integrable, since at low energies ($`p_1,p_21/a,|z|1/(ma^2)`$) we can neglect the momentum dependence of the $`t`$ matrix. The energy $`z_0`$ defines an arbitrary high-energy cut off of the order $`1/(ma^2)`$, so that $`t(z_0)t(z)`$, and we obtain the long-wavelength asymptotic of the $`t`$-matrix in 2D
$$t\frac{4\pi }{m\mathrm{ln}[ฯต_0/(z)]}$$
(86)
We see that in 2D, the $`t`$-matrix vanishes in the limit $`p_1,p_2,z0`$ and in fact diverges at the high-energy cut off $`ฯต_0=|z_0|`$.
Next we need to integrate out the high-energy modes with momenta $`k>k_0^{}`$ in our functional $`\mathrm{exp}S`$ (60). The cut-off $`k_0^{}`$ is defined as
$$max(|\mu |,T)<<\frac{(k_0^{})^2}{2m}<<ฯต_0$$
(87)
The result of this integration is the reduced action
$`S^{}={\displaystyle \underset{\omega ,k<k_0^{}}{}}\left(i\omega {\displaystyle \frac{k^2}{2m}}+\mu \right)a^+(p)a(p)`$
$`{\displaystyle \frac{1}{2\beta V}}{\displaystyle \underset{p_1+p_2=p_3+p_4}{}}t^{}a^+(p_1)a^+(p_2)a(p_3)a(p_4)`$ (88)
where all the summations are cut-off at $`k=k_0^{}`$ and the potential is replaced by a $`t`$-matrix with
$$t^{}=t^{}(\omega )=\frac{4\pi }{m\mathrm{ln}(ฯต_0/[(k_0^2/m)i\omega ])}$$
(89)
Now the functional $`\mathrm{exp}S^{}`$ is to be integrated over the variables $`a^+(p),a(p)`$ within the momentum shell $`k_0<k<k_0^{}`$, where $`k_0`$ is defined from the inequality $`k_0^2/m>>T/\mathrm{ln}|ฯต_0/\mu |`$ and serves to distinguish between slow and rapid particles. Variables with momenta smaller than $`k_0`$ are taken into account in the action by the transformation
$$a^+(p),a(p)(\rho _0(k_0)\beta V)^{1/2}\delta _{p0}$$
(90)
here $`\rho _0(k_0)`$ is the density of slow particles and one can introduce the density of fast particles $`\rho _1(k_0)`$.
After the transformation (90) one can make use of standard perturbation theory formalism and derive expressions for the densities $`\rho _0`$ and $`\rho _1`$. Their sum gives the total density $`\rho =\rho _0+\rho _1=\rho _n+\rho _s`$, which is independent of auxiliary momenta $`k_0`$ and $`k_0^{}`$
$$\rho =\frac{m\mu }{4\pi }(\mathrm{ln}ฯต_0/\mu 1)\frac{1}{(2\pi )^2}d^2k\frac{k^2}{2mฯต(k)}n_B(k)$$
(91)
where $`n_B(k)=(e^{\beta ฯต_k}1)^1`$, and the formulae for the normal and the superfluid component densities read
$`\rho _n={\displaystyle \frac{\beta }{(2\pi )^2}}{\displaystyle d^2k\frac{k^2}{2m}n_B(k)(1+n_B(k))}`$ (92)
$`\rho _s={\displaystyle \frac{m\mu }{4\pi }}(\mathrm{ln}ฯต_0/\mu 1)`$
$`{\displaystyle \frac{1}{(2\pi )^2}}{\displaystyle \frac{d^2k}{2m}k^2n_B(k)\left[\frac{1}{ฯต_k}\beta (n_B(k)+1)\right]}.`$ (93)
Eqs (91) and (93) define the thermodynamics of the system below the phase transition. Above the phase transition one can use the standard perturbation theory for calculation of the Greenโs functions and the thermodynamical quantities. The critical temperature is derived now from the condition $`\rho =\rho _n`$ at the transition.
###### Acknowledgements.
I would like to acknowledge discussions with P.B. Blakie, A. Chubukov, D. Efremov, M. Garst, M. Greiter, A. Rosch and P. Woelfle. I highly appreciate many insightful discussions with M. Eschrig, U. Fischer and F. Nogueira. I am indebted to P. Coleman for reading the manuscript and numerous critical comments. I acknowledge the Humboldt foundation for support and Kavli Institute for Theoretical Physics of Santa Barbara for hospitality, under NSF grant PHW 99/07949.
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# Statistical properties of an ensemble of vortices interacting with a turbulent field
## 1 Introduction
This work presents a study of the statistical properties of a system of vortices interacting with random waves. This is motivated by the necessity to describe quantitatively the statistics of a turbulent plasma in the regime where structures are generated at random, persist for a certain time and are destroyed by perturbations. This regime is supposed to be reached at stationarity in drift wave turbulence when there is a competition between different space scales. These scales range from the radially extended eddies of the ion temperature gradient (ITG) driven modes to the intermediate $`\sqrt{L_n\rho _s}`$ interval where the scalar nonlinearity is dominant over the vectorial one and with random fast decays to few $`\rho _s`$ scales where robust vortex-like structures are generated. In a physical model the description of these states requires to consider simultaneously the long time response and the much faster events of trapping of the energy in small vortices, condensed from transient Kelvin-Helmholtz instabilities of the locally ordered sheared flow.
The analytical model we propose here is essentially *nonperturbative* in the sense that the structures are explicitely represented and we take into account their particular analytical expression (or a reasonable approximation). It is in the spirit of the semiclassical methods, frequently used in field theory. In this sense this treatment presents substantial differences compared with the more usual approaches, aiming in general to obtain a renormalization of the linear response or to calculate the correlations of the fluctuating field using a closure method. When structures are present in the plasma any perturbative approach is inefficient, because it can only go not too far beyond Gaussian statistics, *i.e.* calculate few cumulants beyond the second one. The problem consists in that they take as the base-point for developing perturbative expansions the state of equilibrium of the plasma or a Gaussian ensemble of waves. Or, from this point it is impossible to reach the strongly correlated states of coherent structures. The idea of semiclassical methods is precisely to place the equilibrium state on the structures and then to explore a neighborhood in the space of systemโs configurations, to include random waves. This is the approach we will adopt in the present work. To have a tractable problem we assume the simplest case where the vortices are not created or destroyed dynamically. The final outcome from this approach is a list in which a particular dependence of a contribution to the correlation, *e.g.* the exponent $`\mu `$ in $`\varphi \varphi _๐คk^\mu `$, is associated to a particular contribution in the physical model: statistics of the gas of vortices, interaction energy, nonlinearity of the physical model, etc.
### 1.1 The formulation of the problem
In agreement with the experimental observation and numerical simulations it has been found theoretically that the nonlinear differential equations describing strongly nonlinear electrostatic drift waves in two-dimensional plasma can have (1) turbulent solutions, consisting of very irregular fluctuations which can only be described by statistical quantities (irreducible correlations = cumulants); and (2) solutions that are cuasi-coherent structures of vortex type, which are remarkably robust and for which we have in certain cases explicit analytical expression. Although there is no conceptual difference between these two aspects, we can say that we have two manifestations of the nonlinearity: one is nonlinear mode coupling and energy transfer between waves ; and the second is the generation of organized flow with vortical pattern with strong stability and coherence of form. A first approximation is to break up the field (the electric potential $`\varphi `$) in two distinct elements : vortices and random waves. We will suppose that there is a finite number $`N`$ of vortices in a two-dimensional plasma and that these vortices have random position. The vortices are individually affected by the turbulent background. The turbulent background, in turn, is affected by the presence of the structures at random positions. In addition the turbulent background has statistical properties generated by the nonlinear nature of the waves interactions, even at amplitudes below those necessary to condense vortices. Random growth and decay of modes at marginal stability is included as a drive with Gaussian statistics.
Consider the field $`\varphi ^{Vortices}\varphi _V`$ of the discrete set of $`N`$ vortices individually represented by the electric potential $`\varphi _s^{\left(a\right)}(x,y)`$
$$\varphi _V(x,y)=\underset{a=1}{\overset{N}{}}\varphi _s^{\left(a\right)}(x,y)$$
(1)
interacting with a turbulent wave field $`\varphi ^{waves}\varphi `$. The total field is
$$\phi (x,y)=\underset{a=1}{\overset{N}{}}\varphi _s^{\left(a\right)}(x,y)+\varphi (x,y)$$
(2)
and we want to determine the statistical properties of the field $`\phi (x,y)`$. We will construct an action functional and we will calculate the generating functional of the irreducible correlations (cumulants) of the fluctuating field $`\phi `$. The action functional is expressed in terms of the field $`\phi `$ which we see as composed from vortices and turbulence, Eq.(2). For example, the two-point correlation is composed of four terms
$`\phi \phi `$ $`=`$ $`\varphi ^{Vortices}\varphi ^{Vortices}+\varphi ^{Vortices}\varphi ^{waves}`$ (3)
$`+\varphi ^{waves}\varphi ^{Vortices}+\varphi ^{waves}\varphi ^{waves}`$
Our procedure consists in absorbing the two intermediate terms into the first and the last terms. This is done by calculating the auto-correlations of each component taking into account the presence of the other. In this operation the second and the third terms, although are identical in Eq.(3) are regarded differently: the second term is seen as the contribution of the turbulent background to the auto-correlation of a gas of vortices and the third term is seen as the contribution of a set of vortices to the auto-correlations of a turbulent field.
Therefore we consider that the action functional is composed of two distinct parts: $`S_{V\phi }`$ the action for the system of vortices interacting with the random field; and $`S_{\phi V}`$ the action for the random field interacting with the vortices. Then the statistical ensemble of realizations of the fluctuating field $`\phi (x,y)`$ will actually consists of the Cartesian product of two distinct parts. The generating functional will be
$$Z=Z_{V\phi }Z_{\phi V}$$
(4)
where each factor is calculated using the action defined above. Although the generating functional is factorized, which corresponds to splitting the statistical ensemble into two parts, these two parts are not independent since each will be calculated such as to include the effect of the field from the other subsystem.
The interaction with an external current must be included in each of the action in order to calculate the correlations using functional derivatives. In the full action the current is introduced by adding a term in the Lagrangian density, $`J\phi `$, where $`\phi `$ is the total field, vortices plus random waves. When we separate the two components we have $`\phi =\varphi ^{Vortices}+\varphi ^{waves}`$ and $`J\phi =J\varphi ^{Vortices}+J\varphi ^{waves}`$ . This expression must be inserted in the action functional for the full system, before the intermediate terms in Eq.(3) are absorbed into the first and the last. This means that the *same* current $`J`$ will appear in the final expressions of the two generating functionals Eq.(4)
$$Z\left[J\right]=Z_{V\phi }\left[J\right]Z_{\phi V}\left[J\right]$$
(5)
It will be shown below that the value of the field at a particular point can be obtained by functional derivation to the current at that point $`\phi =\frac{\delta }{\delta J}`$. We obtain the average of the fluctuating field as
$`\varphi ={\displaystyle \frac{1}{Z\left[J=0\right]}}{\displaystyle \frac{\delta }{\delta J}}Z\left[J\right]|_{J=0}`$ (6)
$`=`$ $`{\displaystyle \frac{1}{Z_{V\phi }\left[J=0\right]}}{\displaystyle \frac{\delta }{\delta J}}Z_{V\phi }\left[J\right]|_{J=0}+{\displaystyle \frac{1}{Z_{\phi V}\left[J=0\right]}}{\displaystyle \frac{\delta }{\delta J}}Z_{\phi V}\left[J\right]|_{J=0}`$
The two-point correlation is obtained from the generating functional by applying two times this functional derivative, with the current in two distinct points and taking finally the current to zero.
We now explain how this will be done effectively.
### 1.2 Outline of the present analytical approach
#### 1.2.1 Gas of vortices in the turbulent background
In the case of the vortices interacting with random waves, the basic reasoning takes into consideration two distinct elements.
The first is concerned with the statistical properties of a collection of discrete vortices with zero, or very short, range of interaction. At this stage of the problem the particular shape of the potential distribution in a vortex is not essential and will be simplified in those situations where only the positions of the centers are important. The statistical ensemble consists of the configurations of randomly placed vortices, seen as a dilute gas.
The second aspect is concerned with the interaction between one (generic) vortex and the random waves in the surrounding turbulence and calculates the effect on the form of the potential of the vortex. Instead of the exact solution we will have now a form resulting from the scattering due to the random perturbation produced by the turbulent background.
The statistical properties result from the combination of the two elements, which will now be explained briefly.
The first problem is very similar to the Coulomb gas in two dimensions or to any problem related to dilute gas of interacting particles. The partition function is calculated using the sum of the energy of the individual vortices and the energy of interaction. The first part can be calculated from the exact analytical solution of the differential equation of the model (Eq.(10) below). The second part contains the sum of the energies of pair interaction, a problem that is principially complicated in our case, *i.e.* in turbulent plasma. Our *real* vortices are neutral (they are simply a deformation, stable and regular, of the scalar potential of the velocity) and there should be no interaction of Coulomb or Ampere type. However there is an elastic medium between the vortices and any motion of one of them generates sound waves that may couple with the potential of other vortices, influencing their motion. This may suggest that the interaction is of the same nature as the Casimir effect but this is beyond our primary interest here. What we minimally need to include in our partition function is the effect of scattering of the vortices at close encounter, since we do not take into account in the present treatment the variation of the number of vortices due to merging (which would imply a chemical potential). There is an intrinsic spatial scale in all aspects related with drift-wave vortices and this is the sonic Larmor radius, altered by the combined effect of the diamagnetic and a uniform flow. We show below that the pair interaction contains a kernel with fast spatial decay, $`K_0\left(r/\rho _s\right)`$ (the modified Bessel function). There are several reasons to assume that this is a physically correct choice (see ).
The second aspect of the problem of vortices acted upon by turbulence is related to the direct modification of the exact vortex by the random waves around it. The case of a single vortex interacting with random waves has been treated in the papers , . The starting point is the exact vortex-type solution of the differential equation we investigate. This is of course a non-random object and when it is isolated the statistical ensemble is trivial. However its interaction with random waves of the background turbulence makes it also a fluctuating object. We calculate the statistical properties of the ensemble of states of the fluctuating potential corresponding to the shapes of the vortex-type solution.
The action functional is extremum at the vortex solution, whose explicit expression we will use. A purely turbulent field also realizes the extremum of the action, but our separation between structures and random waves amounts to an approximative representation of the turbulent field, as a complimentary part to the structures. All configurations consisting of structures with randomly deformed, fluctuating shapes, together with random waves are a very good approximation of reality and must be found in close proximity of the extremum of the action. Then we have to explore the functional space of the systemโs configuration in the neighborhood of the structure (the extremum) trying to include as many as possible nearby configurations in the calculation of the generating functional. This will include into the generating functional the turbulent field, besides the structure. In few words, the idea is that the structure and the random waves, although very different in geometry, share a common property: they obtain or are very close to the extremum of the action functional. The technical procedure will consists of expanding the action functional to second order around the structure and to integrate over the space of configurations. This will automatically exclude bad approximations, *i.e.* the states which are too far from realizing the extremum of the action since for them the Boltzmann weight is exponentially small.
We conclude by noting that this is the standard semiclassical treatment . In this sense it is similar to the treatment of vortex statistics of the Abelian-Higgs model of superfluids and of many other systems.
#### 1.2.2 The turbulent field influenced by random vortices
For the turbulent field interacting with vortices, our approach combines two distinct elements. The first is the inclusion of the vortices as a random perturbation in the generating functional of the turbulent waves. The second is a perturbative treatment of the intrinsic nonlinearity of the turbulent waves, already modified by the inclusion of the effect of the random vortices.
In the first part we follow the similar approaches as for the electron conduction in the presence of random impurities, or as for flexible polymers in porous media. For the case where the vortices have uncorrelated random positions and take at random positive or negative amplitudes of equal magnitude we find that the problem is mapped onto the *sine*-Gordon model (actually *sinh*). For our model equation this will amount to a renormalization of the coefficient which plays the role of a physical โmassโ of the turbulent field , , or, in other terms to a shift of the spatial scale from $`\rho _s\left(1v_d/u\right)^{1/2}`$ to higher values.
The second part consists of a systematic perturbative treatment of the nonlinearity in order to get, as much as possible, a correct representation of the nonlinear content of the field of the turbulent waves (this means the nonlinear interaction and nonlinear energy transfer between low amplitude random waves , \[Horton1\], , ). The nonlinearity is included in a functional perturbative treatment where the turbulent plasma is driven by random rise and decay of modes at marginal stability. This induces a diffusive behavior of the turbulent field, at lowest order. We develop the treatment to one-loop, which means of order two in the strength of the nonlinearity.
One may inquire if this is not in contradiction with the separation operated at the beginning, where the most characteristic aspect of the nonlinearity, the generation of structure, is treated separately. However, it is clear that there is no danger of overlapping and double counting of the nonlinear effects. Since a statistical perturbative treatment is essentially an expansion in a parameter representing the departure from Gaussian statistics we can only hope to include higher cumulants beyond the second one (which means Gaussianity) but only few of them are accessible to effective calculation. Or, any structure needs a very large number of cumulants since, by definition, is an almost coherent field. It is illusory to try to capture the structure using a perturbative treatment starting with a state of equilibrium (*i.e.* no perturbation or, alternatively, a Gaussian collection of linear waves). Since any term of the perturbation can be represented by a Feynman diagram, we face the well known problem that the proliferation of diagrams at high orders leads to an effectively intractable problem. This is actually one reason for the use of the semiclassical methods.
## 2 The physical model
### 2.1 The equation
The model of the ion drift instability in magnetically confined plasmas can be formulated using the fluid equations of continuity and momentum conservation for electrons and for ions. It has been shown by a multiple space time scale analysis that the dynamics is dominated by two nonlinearities: the Charney-Hasegawa-Mima type, or vectorial nonlinearity, generated by the ion polarization drift and the Korteweg-De Vries, or scalar nonlinearity, related to the space variation of the density gradient length. The former is of high differential degree and is dominant at small spatial scales, of the order of few sonic Larmor radii $`\rho _s`$. The latter is dominant at โmesoscopicโ spatial scales, of the order of $`\sqrt{\rho _sL_n}`$. The numerical studies \[Mikhailovskaya\], the fluid-tank experiments \[Nezlin\] and the multiple space-time scale analytical analysis show that the scalar nonlinearity becomes dominant at late regimes in the statistical stationarity of the drift wave turbulence. However the possible manifestation of the two types of nonlinearity rises the problem of *structural* stability of either regime where only one of these nonlinearity is considered dominant: inclusion of the other strongly changes the behavior of the system. It is then pertinent to consider that the turbulence generated in a realistic regime may include manifestations of both types. The turbulence is dominated by the larger scales sustained by the scalar nonlinearity (described by the Flierl-Petviashvili equation ) together with robust vortices generated at the scales of few $`\rho _s`$, typical the CHM (*i.e.* vectorial) nonlinearity \[Nycander2\].
When the scalar nonlinearity is prevailing , \[MH2\] the equation has the form
$$\left(1_{}^2\right)\frac{\phi }{t}+v_{}\frac{\phi }{y}v_T\phi \frac{\phi }{y}=0$$
(7)
In a moving frame and restricting to stationarity we obtain
$$_{}^2\phi \alpha \phi \beta \phi ^2=0$$
(8)
The physical parameters are , ,
$$\alpha =\frac{1}{\rho _s^2}\left(1\frac{v_d}{u}\right),\beta =\frac{c_s^2}{2u^2}\frac{}{x}\left(\frac{1}{L_n}\right)$$
(9)
where $`\rho _s=c_s/\mathrm{\Omega }_i`$, $`c_s=\left(T_e/m_i\right)^{1/2}`$ and the potential is scaled as $`\varphi e\varphi /T_e`$. Here $`L_n`$ and $`L_T`$ are respectively the gradient lengths of the density and temperature. The velocity is the diamagnetic velocity $`v_d=\rho _sc_s/L_n`$. The condition for the validity of this equation are: $`\left(k_x\rho _s\right)\left(k\rho _s\right)^2\eta _e\rho _s/L_n`$, where $`\eta _e=L_n/L_{T_e}`$. The coefficients $`\alpha `$ and $`\beta `$ have the dimension $`\left(length\right)^2`$. This form will be used below.
### 2.2 The structures
The exact solution of the equation is
$`\phi _s(y,t;y_0,u)`$ $`=`$ $`3\left({\displaystyle \frac{u}{v_d}}1\right)`$ (10)
$`\times \mathrm{sec}\mathrm{h}^2\left[{\displaystyle \frac{1}{2\rho _s}}\left(1{\displaystyle \frac{v_d}{u}}\right)^{1/2}\left(yy_0ut\right)\right]`$
where the velocity is restricted to the intervals $`u>v_d`$ or $`u<0`$. In Ref. the radial extension of the solution is estimated as: $`\left(\mathrm{\Delta }x\right)^2\rho _sL_n`$. In our work we shall assume that $`u`$ is close to $`v_d`$ , $`uv_d`$ (i.e. the structures have small amplitudes). The monopolar vortex in this regime is discussed in Ref. \[Nycander2\]. For asymptotic form of the CHM equation see Ref.. We will adopt the one-dimensional section of the solution (*i.e.* Eq.(10)) when we calculate the eigenmodes of the determinant of the second functional derivative of the action. In conclusion, we consider coherent structures which are monopolar vortices of both signs of vorticity, with equal magnitudes and with random positions in plane.
## 3 Statistical analysis of the physical system
### 3.1 General functional framework
The general method for constructing the action functional for a classical stochastic system is described by Martin, Siggia and Rose (MSR) , and in path integral formalism, by Jensen . Two reviews by Krommes are very useful references on this point , . The functional method has been applied in several concrete problems and references may be consulted for details , , , , . We here review few elementary procedures (see ).
Consider a differential equation $`F\left[\varphi \right]=0`$ whose solution is $`\varphi ^z`$. The unknown function belongs to a space of functions $`\varphi (x,y)`$. We want to select from this space of functions precisely the one that is the solution of the differential equation and for this we can use the functional Dirac $`\delta `$ : $`\delta \left(\varphi \varphi ^z\right)`$. This can be represented as a product of usual $`\delta `$ functions in every point of space
$$\delta \left(\varphi \varphi ^z\right)=\underset{k=1}{\overset{N}{}}\delta \left[\varphi \left(x_k\right)\varphi ^z\left(x_k\right)\right]$$
(11)
Any operation that will be done on a functional of $`\varphi `$ can be now particularized to the solution $`\varphi ^z`$ by simply inserting this Dirac functional. For example the calculation of a functional $`\mathrm{\Omega }\left(\varphi \right)`$ at the function $`\varphi ^z`$ can be done by a functional integration over the space of all functions, with insertion of this $`\delta `$ functional. Using the Fourier representation of the ordinary Dirac functions and going to the continuous limit we note the appearence of the dual function $`\chi `$
$`\mathrm{\Omega }\left(\varphi ^z\right)`$ $`=`$ $`{\displaystyle ๐\left[\varphi \right]\mathrm{\Omega }\left(\varphi \right)\left(\left|\frac{\delta F}{\delta \varphi }\right|_{\varphi ^z}\right)\delta \left[F\left(\varphi \right)\right]}`$ (12)
$`=`$ $`\left(\left|{\displaystyle \frac{\delta F}{\delta \varphi }}\right|_{\varphi ^z}\right){\displaystyle ๐\left[\varphi \right]๐\left[\chi \right]\mathrm{\Omega }\left(\varphi \right)}`$
$`\times \mathrm{exp}\left\{i{\displaystyle ๐x\chi \left(x\right)F\left[\varphi \left(x\right)\right]}\right\}`$
We will define the partition function as usual, by the functional integral of Boltzmann weights calculated on the base of the MSR action
$$Z\left[J\right]=๐\left[\varphi \right]๐\left[\chi \right]\mathrm{exp}\left\{๐x๐y\left(\chi F\left[\varphi \right]+J_\varphi \varphi +J_\chi \chi \right)\right\}$$
(13)
The functional integration takes into account the fluctuations of the physical field $`\varphi `$ and of its dual $`\chi `$. The โfree-energyโ functional is defined by $`\mathrm{exp}\left\{W\left[J\right]\right\}=Z\left[J\right]`$ from which the irreducible correlations (cumulants) are calculated by functional derivatives to $`J`$.
$$\varphi (x,y)=\mathrm{exp}\left\{W\left[J\right]\right\}\frac{\delta }{\delta J_\varphi (x,y)}\mathrm{exp}\left\{W\left[J\right]\right\}|_{J=0}$$
(14)
and similar for higher cumulants. The two-point irreducible correlation for the field is
$`\varphi (x,y)\varphi (x^{},y^{})`$ (15)
$`=`$ $`{\displaystyle ๐\left[\varphi \right]๐\left[\chi \right]\varphi (x,y)\varphi (x^{},y^{})}`$
$`\times \mathrm{exp}\left\{{\displaystyle ๐x๐y\left(\chi F\left[\varphi \right]+J_\varphi \varphi +J_\chi \chi \right)}\right\}|_{J=0}`$
$`=`$ $`{\displaystyle \frac{1}{Z\left[J=0\right]}}{\displaystyle \frac{\delta }{\delta J_\varphi (x,y)}}{\displaystyle \frac{\delta }{\delta J_\varphi (x^{},y^{})}}Z\left[J\right]|_{J=0}`$
This is the general analytical instrument that will be used in the following calculations. The calculations from this work are presented in greater detail in Ref. \[FlorinMadi6\].
## 4 Vortices with random positions
### 4.1 The discrete set of $`N`$ vortices in plane
The action of the discrete set of vortices is determined by the sum of the actions of the individual vortices plus a part that results from the interaction between them . The first part is simply the time integration of the energy, *i.e.* (since time factorizes) the space integration of the expression of the product of the static potentials associated with a single vortex, Eq.(10), and to its dual, which in the end simply means the square of the wave-form of the vortex potential. This quantity is repeated for each of the $`N`$ vortices.
The partition function is
$`Z_V`$ $`=`$ $`{\displaystyle \frac{1}{N!}}\left({\displaystyle \underset{j=1}{\overset{N}{}}}Z_V^{\left(0\right)}\right){\displaystyle \underset{\left\{\alpha \right\}}{}}{\displaystyle \left(\underset{a=1}{\overset{N}{}}\frac{1}{A}d๐^{\left(a\right)}\right)}`$ (16)
$`\times \mathrm{exp}[\pi \underset{a>b}{{\displaystyle \underset{a=1}{\overset{N}{}}}{\displaystyle \underset{b=1}{\overset{N}{}}}}{\displaystyle }dxdy{\displaystyle }dx^{}dy^{}\rho _\omega ^{\left(a\right)}(x,y)`$
$`\times K_0\left(\rho _s^1\right|๐^{\left(a\right)}๐^{\left(b\right)}\left|\right)\rho _\omega ^{\left(b\right)}(x^{},y^{})]`$
with the following meaning.
The first factor simply takes into account $`N`$ independent vortices with arbitrary positions in plane and expresses the fact that this part of the partition function results from a Cartesian product of the $`N`$ statistical ensembles, one for each vortex. The factor $`N!`$ takes into account the permutation symmetry. The generating functional for a static vortex with structure given by the interaction with random waves, is calculated in Eq.(65) below. We have $`Z_V^{\left(0\right)}=Z_{V\phi }\left[J\right]`$ where we have indicated by the index $`V\phi `$ that the partition function of the vortex includes the interaction vortex-turbulence, and that the expression depends on the external current $`J`$, and will contribute to any correlation that we will obtain by functional derivations at $`J`$.
The sum is over the set of configurations $`\left\{\alpha \right\}`$ characterized by random choices of positive and negative vortices.
The integrations over the positions of the centers $`๐^{\left(a\right)}`$ of the vortices, $`a=1,N`$, express the fact that we allow arbitrary positions in plane, with equal probability. Each integral is normalized with the area of the physically interesting two-dimensional region of the plane, $`A`$.
The exponent of the Boltzmann weight contains action resulting from the interaction between vortices. We expect that for a dilute gas of vortices, where the distance between the centers $`๐^{\left(a\right)}๐^{\left(b\right)}`$ is much larger than the core diameter $`d`$, $`๐^{\left(a\right)}๐^{\left(b\right)}d,a,b=1,N`$ , the interaction is very weak. In order to describe the interaction between vortices we start from the well-known alternative model of the drift waves sustained by the ion polarization drift nonlinearity , . In this model it is considered a set of $`N_\omega `$ point-like vortices of strength $`\omega _i`$ interacting in plane by a short range potential expressed as the function $`K_0`$ (modified Bessel function) of the relative distance between vortices. The potential $`\varphi ^p`$ in a point $`๐`$ is a sum of contributions from all the $`N_\omega `$ point-vortices $`\varphi ^p\left(๐\right)=_{i=1}^{N_\omega }\omega _iK_0\left(\rho _s^1\left|๐๐_i\right|\right)`$ and the equations of motion $`d๐_i/dt=\varphi ^p\times \widehat{๐}_z`$ (where $`\widehat{๐}_z`$ is the versor perpendicular on the plane). The distribution of vorticity of the physical system (in particular the quasi-coherent vortical structures) represents spatial variations of density of these point-like vortices. The interaction between the *physical* vortices will result from the interaction between the point-like vortices, taking into account the density of these objects. The energy of interaction is
$$H=\underset{i>j}{\underset{i=1}{\overset{N_\omega }{}}\underset{j=1}{\overset{N_\omega }{}}}\omega _i\omega _jK_0\left(\rho _s^1\left|๐_i๐_j\right|\right)$$
called the Kirchhoff function. The range of spatial decay of the interaction is the Larmor sonic radius $`\rho _s`$, which however may be modified to an *effective* Larmor radius, in the presence of gradients and flow. When we approach the continuum limit $`N_\omega \mathrm{}`$ the envelope of the density becomes the physical vorticity $`\omega (x,y)`$ which, for this stage of the problem is sufficient to be considered as highly concentrated in the cores of the physical vortices and almost vanishing in the rest. Taking the elementary point-vortices of equal strength $`\omega _j\omega _0`$ we have that each physical vortex is an integer multiple $`N^{\left(a\right)}`$ of this quantity. Now we will associate with each physical vortex a continuous function, *i.e.* its vorticity defined on the whole plane, $`\rho _\omega ^{\left(a\right)}(x,y)`$, which is, as said, concentrated in $`(x,y)`$
$$\rho _\omega ^{\left(a\right)}(x,y)=N^{\left(a\right)}\omega _0\delta \left(๐๐^{\left(a\right)}\right)$$
(17)
Then the energy is
$`H`$ $`=`$ $`\underset{a>b}{{\displaystyle \underset{a=1}{\overset{N}{}}}{\displaystyle \underset{b=1}{\overset{N}{}}}}{\displaystyle ๐x๐y๐x^{}๐y^{}}`$ (18)
$`\times \rho _\omega ^{\left(a\right)}(x,y)K_0\left({\displaystyle \frac{\left|๐^{\left(a\right)}๐^{\left(b\right)}\right|}{\rho _s}}\right)\rho _\omega ^{\left(b\right)}(x^{},y^{})`$
Due to Eq.(17) the interaction energy is only the interaction between the *centers* $`R^{\left(a\right)}`$ and $`R^{\left(b\right)}`$ of the vortices. The summation proceeds by grouping the point-vortices into physical vortices, then assuming that these (for only this stage of the problem) have $`\delta `$-function shape and finally formally replacing this $`\delta `$ with a continuous distribution $`\rho _\omega ^{\left(a\right)}(x,y)`$. In this operation a number of infinities arise from the energy of the interaction of the point-vortices which are grouped into one physical vortex, since the relative distances are zero for them. This singular part can be removed since it does not participate to the functional variations induced by the โexternal excitationโ current $`J`$.
For simplification of the computation we now only consider physical vortices of equal amplitude (positive or negative) and then the vorticity distribution $`\rho _\omega ^{\left(a\right)}(x,y)`$ has amplitude $`\omega _v=p\omega _0`$ ($`p`$ is an integer), multiplied by the integer $`n^{\left(a\right)}`$ which can take the values $`\pm 1`$ for positive or negative vorticity. We have $`N^{\left(a\right)}=pn^{\left(a\right)}`$. The sum over the physical vorticesโ positions suggests to define a formal unique function of vorticity $`\rho _\omega (x,y)`$ for all the $`N`$ physical vortices
$$\rho _\omega (x,y)\underset{a=1}{\overset{N}{}}\rho _\omega ^{\left(a\right)}(x,y)=\underset{a=1}{\overset{N}{}}\omega _vn^{\left(a\right)}\delta \left(๐๐^{\left(a\right)}\right)$$
Further, the energy is normalized with a constant dimensional factor. The interaction part can be rewritten
$`\mathrm{exp}[{\displaystyle \frac{\pi }{\rho _s^4\omega _v^2}}\underset{a>b}{{\displaystyle \underset{a=1}{\overset{N}{}}}{\displaystyle \underset{b=1}{\overset{N}{}}}}{\displaystyle }dxdy{\displaystyle }dx^{}dy^{}`$ (19)
$`\times \rho _\omega ^{\left(a\right)}(x,y)K_0\left(\rho _s^1\right|๐^{\left(a\right)}๐^{\left(b\right)}\left|\right)\rho _\omega ^{\left(b\right)}(x^{},y^{})]`$
$`=`$ $`\mathrm{exp}\left[{\displaystyle \frac{1}{2\rho _s^4\omega _v^2}}{\displaystyle ๐x๐y๐x^{}๐y^{}\rho _\omega (x,y)G\left(๐๐^{}\right)\rho _\omega (x^{},y^{})}\right]`$
where $`G`$ is the kernel of interaction,
$`G\left(๐๐^{}\right)`$ $``$ $`{\displaystyle \frac{1}{2\pi }}\underset{a>b}{{\displaystyle }}K_0\left(\rho _s\left|๐^{\left(a\right)}๐^{\left(b\right)}\right|\right)`$ (20)
$`\times \delta \left(๐๐^{\left(a\right)}\right)\delta \left(๐^{}๐^{\left(b\right)}\right)`$
The differential equation for $`K_0`$ is $`\left(\mathrm{\Delta }1/\rho _s^2\right)K_0\left(r/\rho _s\right)=2\pi \delta \left(๐ซ\right)`$ . This helps to replace the interaction part with a Gaussian functional integral, by introducing an auxiliary field $`\psi `$
$`\mathrm{exp}\left[{\displaystyle \frac{1}{2\rho _s^4\omega _v^2}}{\displaystyle ๐x๐y๐x^{}๐y^{}\rho _\omega \left(๐\right)G\left(๐๐^{}\right)\rho _\omega \left(๐^{}\right)}\right]`$ (21)
$`=`$ $`p_1^1{\displaystyle ๐\left[\psi \right]\mathrm{exp}\left\{\frac{1}{2}๐x๐y\left[\left(\psi \right)^2+\frac{1}{\rho _s^2}\psi ^2\right]\right\}}`$
$`\times \mathrm{exp}\left[i{\displaystyle \frac{1}{\rho _s^2\omega _v}}{\displaystyle ๐x๐y\rho _\omega \left(๐\right)\psi \left(๐\right)}\right]`$
with $`p_1`$ a normalization constant. We make a change of variable in the functional integration $`2\pi \psi \chi `$ (this also changes the normalization constant $`p_1p`$) and return to the partition function Eq.(16)
$`Z_V`$ $`=`$ $`p^1{\displaystyle ๐\left[\chi \right]\mathrm{exp}\left\{\frac{1}{8\pi ^2}๐x๐y\left[\left(\chi \right)^2+\frac{1}{\rho _s^2}\chi ^2\right]\right\}}`$ (22)
$`\times {\displaystyle \frac{1}{N!}}\left(Z_V^{\left(0\right)}\right)^N{\displaystyle \underset{\left\{\alpha \right\}}{}}{\displaystyle \left(\underset{a=1}{\overset{N}{}}\frac{1}{A}d๐^{\left(a\right)}\right)\mathrm{exp}\left[i\underset{a}{}n^{\left(a\right)}\chi \left(๐^{\left(a\right)}\right)\right]}`$
In the last factor we note that in the sum each term consists of two contributions, corresponding to positive and negative vorticity, $`n^{\left(a\right)}=\pm 1`$, and they are weighted with the same factor, $`1/2`$
$`{\displaystyle \frac{1}{N!}}\left(Z_V^{\left(0\right)}\right)^N{\displaystyle \underset{\left\{\alpha \right\}}{}}{\displaystyle \underset{a=1}{\overset{N}{}}\frac{1}{A}d๐^{\left(a\right)}\mathrm{exp}\left[i\underset{a}{}n^{\left(a\right)}\chi \left(๐^{\left(a\right)}\right)\right]}`$ (23)
$`=`$ $`{\displaystyle \frac{\left(Z_V^{\left(0\right)}\right)^N}{N!}}\left[{\displaystyle \underset{a=1}{\overset{N}{}}}{\displaystyle \frac{1}{A}๐๐\frac{1}{2}\left(\mathrm{exp}\left[i\chi \left(๐\right)\right]+\mathrm{exp}\left[i\chi \left(๐\right)\right]\right)}\right]`$
In the last line we have removed the upper index $`\left(a\right)`$ since all factors in the product are now identical. For a fixed number $`N`$ of vortices the partition function $`Z_V^{\left(0\right)}`$ (with nonvanishing contribution to the derivatives to $`J`$) is decoupled from the other factors and will provide $`N`$-times the same contribution. The other factors, *i.e.* the functional integral that contains the interaction between the vortices can only appear in the final answer as a constant, multiplying contributions coming from $`Z_V^{\left(0\right)}\left[J\right]`$. When $`N`$ is arbitrary the partition function must also include a sum over terms each corresponding to a number $`N`$ of vortices
$`{\displaystyle \underset{N=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\left(Z_V^{\left(0\right)}\right)^N}{N!}}\left[{\displaystyle \frac{1}{A}๐๐\frac{1}{2}\left(\mathrm{exp}\left[i\chi \left(๐\right)\right]+\mathrm{exp}\left[i\chi \left(๐\right)\right]\right)}\right]^N`$ (24)
$`=`$ $`\mathrm{exp}\left\{{\displaystyle \frac{Z_V^{\left(0\right)}}{A}}{\displaystyle ๐x๐y\mathrm{cos}\left[\chi (x,y)\right]}\right\}`$
Then the partition function becomes
$$Z_V=p^1๐\left[\chi \right]\mathrm{exp}\left\{\frac{1}{8\pi ^2}๐x๐y\left[\left(\chi \right)^2+\frac{1}{\rho _s^2}\chi ^2\frac{8\pi ^2}{A}Z_V^{\left(0\right)}\mathrm{cos}\left[\chi (x,y)\right]\right]\right\}$$
(25)
In this expression the quantity $`Z_V^{\left(0\right)}`$ is a functional integral over the space of fields $`\varphi _s(x,y)`$ representing a single vortex. In the absence of the background turbulence the field $`\varphi _s(x,y)`$ is a deterministic quantity \[Eq.(10)\] and the statistical ensemble is trivially composed of one element. The interaction with the background turbulence induce a fluctuating form and the statistical properties can be obtained from $`Z_V^{\left(0\right)}`$. In other words, $`Z_V`$ includes two sources of fluctuations: one is the fluctuation of the field of vorticity due to the random positions in plane of the vortices (the gas of vortices) and the other is the fluctuation of the shape of the generic vortex due to interaction with the turbulent background.
If we want to use the expression (25) as a generating functional for correlation we must be able to drop into the functional integral the field representing the vortices, *i.e.*
$$\varphi _V\left(๐\right)=\underset{a}{}\varphi _s^{\left(a\right)}(x,y)$$
(26)
as in Eq.(15). An external excitation by the current $`J`$ will produce a change in $`Z_V`$ from the change of the vorticity of the gas of vortices and from the change of $`Z_V^{\left(0\right)}`$. We have
$$\varphi _V(x,y)=\frac{1}{Z_V\left[j=0\right]}\left[\left(\frac{\delta Z_V}{\delta J(x,y)}\right)_{Vort}+\frac{\delta Z_V}{\delta Z_V^{\left(0\right)}}\frac{\delta Z_V^{\left(0\right)}}{\delta J(x,y)}\right]$$
(27)
$`\varphi _V(x,y)\varphi _V(x^{},y^{})`$ $`=`$ $`{\displaystyle \frac{1}{Z_V\left[j=0\right]}}[\left({\displaystyle \frac{\delta ^2Z_V}{\delta J(x,y)\delta J(x^{},y^{})}}\right)_{vort}`$
$`+{\displaystyle \frac{\delta ^2Z_V}{\delta \left(Z_V^{\left(0\right)}\right)^2}}{\displaystyle \frac{\delta Z_V^{\left(0\right)}}{\delta J(x,y)}}{\displaystyle \frac{\delta Z_V^{\left(0\right)}}{\delta J(x^{},y^{})}}`$
$`+{\displaystyle \frac{\delta Z_V}{\delta Z_V^{\left(0\right)}}}{\displaystyle \frac{\delta ^2Z_V^{\left(0\right)}}{\delta J(x,y)\delta J(x^{},y^{})}}]`$
The formulas are taken at $`J=0`$. The first terms in these equations are related with the fluctuations of the vorticity of the gas of vortices as a continuous version of the discrete set of physical vortices with arbitrary positions in plane. The other terms are related to the fluctuation of the shape of a vortex and in order to calculate these contribution we need the explicit expression of $`Z_V`$, as a functional of $`Z_V^{\left(0\right)}`$. Further, we will need the detailed expression of $`Z_V^{\left(0\right)}\left[J\right]`$ and this will be calculated in the next Section.
In order to obtain the contribution from the fluctuation of the gas of vortices (the first terms in Eqs.(27) and (4.1)), we introduce a new term in the action, consisting of the interaction between the vorticity distribution $`\rho _\omega \left(๐\right)`$ and an external current, $`J_\omega `$
$$i\frac{1}{\omega _v}๐x๐y\left[\rho _\omega (x,y)J_\omega (x,y)\right]$$
(29)
(the factor $`i`$ is introduced for compatibility with Eq.(21)). This current $`J_\omega `$ is an external excitation applied on the field of the vorticity and not on the field of potential $`\varphi _V`$ as we would need in the Eq.(5). We may assume that there is a connection between the current $`J_\omega `$ and the current $`J`$ (which acts on the field $`\varphi _V`$) but there is no need to specify this relation. Indeed, the Eq.(6) shows that at the end both currents should be taken zero.
The last line of Eq.(21) transforms as follows
$$\mathrm{exp}\left[i\frac{1}{\omega _v}๐x๐y\rho _\omega \left(๐\right)\psi \left(๐\right)\right]\mathrm{exp}\left\{i\frac{1}{\omega _v}๐x๐y\rho _\omega \left(๐\right)\left[\psi \left(๐\right)+J_\omega \left(๐\right)\right]\right\}$$
(30)
All the calculations following Eq.(21) are repeated without modification but in the last term, instead of the function $`\chi (x,y)`$ we will have
$$\mathrm{cos}\left[\chi (x,y)\right]\mathrm{cos}\left[\chi (x,y)+J_\omega (x,y)\right]$$
(31)
since this was the term which resulted from the presence of the function $`\rho _\omega \left(๐\right)`$ in the Eq.(21). Making the change of variable in the functional integration (of Jacobian $`1`$)
$$\chi \chi J_\omega $$
(32)
the integrand in the action is expressed as
$$\frac{1}{8\pi ^2}\left\{\left[\left(\chi J_\omega \right)\right]^2+\frac{1}{\rho _s^2}\left(\chi J_\omega \right)^2\frac{8\pi ^2}{A}Z_V^{\left(0\right)}\mathrm{cos}\left[\chi (x,y)\right]\right\}$$
(33)
The two-point correlation of the field of the vorticity fluctuations can be calculated from
$$\frac{1}{\omega _v^2}\rho _\omega (x,y)\rho _\omega (x^{},y^{})=\frac{1}{Z_V\left[J_\omega =0\right]}\frac{\delta Z_V\left[J\right]}{i\delta J_\omega (x,y)i\delta J_\omega (x^{},y^{})}|_{J_\omega =0}$$
(34)
The squares in the action Eq.(33) are expanded
$``$ $`{\displaystyle \frac{1}{8\pi ^2}}\{(\chi )^2+{\displaystyle \frac{1}{\rho _s^2}}\chi ^2{\displaystyle \frac{8\pi ^2}{A}}Z_V^{\left(0\right)}\mathrm{cos}\left[\chi (x,y)\right]`$
$`+\left(J_\omega \right)^2+{\displaystyle \frac{1}{\rho _s^2}}J_\omega ^2`$
$`2(\chi )(J_\omega ){\displaystyle \frac{2}{\rho _s^2}}\chi J_\omega \}`$
In the part of the action that depends on $`J_\omega `$ we make integrations by parts
$`\left(J_\omega \right)^2+{\displaystyle \frac{1}{\rho _s^2}}J_\omega ^22\left(\chi \right)\left(J_\omega \right){\displaystyle \frac{2}{\rho _s^2}}\chi J_\omega `$
$``$ $`J_\omega \left(\mathrm{\Delta }J_\omega \right)+{\displaystyle \frac{1}{\rho _s^2}}J_\omega ^2+2J_\omega \left(\mathrm{\Delta }\chi \right){\displaystyle \frac{2}{\rho _s^2}}\chi J_\omega `$
The first line at the exponent in Eq.(4.1) does not contain the current and in the following functional derivations to $`J_\omega `$ we will temporary omit it. Consider the application of the first operator of derivation to $`iJ_\omega (x,y)`$
$`{\displaystyle \frac{\delta }{i\delta J_\omega (x,y)}}\mathrm{exp}\{{\displaystyle \frac{1}{8\pi ^2}}{\displaystyle }dxdy[J_\omega \left(\mathrm{\Delta }J_\omega \right)+{\displaystyle \frac{1}{\rho _s^2}}J_\omega ^2`$
$`+2J_\omega \left(\mathrm{\Delta }\chi \right){\displaystyle \frac{2}{\rho _s^2}}\chi J_\omega ]\}`$
$`=`$ $`{\displaystyle \frac{1}{i}}\mathrm{exp}\left\{\mathrm{}\right\}`$
$`\times \rho _s^2{\displaystyle \frac{\left(1\right)}{8\pi ^2}}\left(2\mathrm{\Delta }J_\omega +{\displaystyle \frac{2}{\rho _s^2}}J_\omega +2\mathrm{\Delta }\chi {\displaystyle \frac{2}{\rho _s^2}}\chi \right)_{(x,y)}`$
Every derivation to the current $`J_\omega `$ suppresses a space integration and in consequence the result is multiplied with factors $`\rho _s`$ which render the space integral dimensionless. The subscript shows that the functions inside the bracket are calculated in the point $`(x,y)`$.
The second operator of derivation is now applied on Eq.(4.1)
$`{\displaystyle \frac{\delta }{i\delta J_\omega (x^{},y^{})}}\mathrm{exp}\{{\displaystyle \frac{1}{8\pi ^2}}{\displaystyle }dxdy[J_\omega \left(\mathrm{\Delta }J_\omega \right)+{\displaystyle \frac{1}{\rho _s^2}}J_\omega ^2`$
$`+2J_\omega \left(\mathrm{\Delta }\chi \right){\displaystyle \frac{2}{\rho _s^2}}\chi J_\omega ]\}`$
$`\times {\displaystyle \frac{1}{i}}\rho _s^2{\displaystyle \frac{\left(1\right)}{8\pi ^2}}\left(2\mathrm{\Delta }J_\omega +{\displaystyle \frac{2}{\rho _s^2}}J_\omega +2\mathrm{\Delta }\chi {\displaystyle \frac{2}{\rho _s^2}}\chi \right)_{(x,y)}`$
$`=`$ $`{\displaystyle \frac{1}{i}}\rho _s^2{\displaystyle \frac{\left(1\right)}{8\pi ^2}}\left(2\mathrm{\Delta }J_\omega +{\displaystyle \frac{2}{\rho _s^2}}J_\omega +2\mathrm{\Delta }\chi {\displaystyle \frac{2}{\rho _s^2}}\chi \right)_{(x^{},y^{})}`$
$`\times {\displaystyle \frac{1}{i}}\rho _s^2{\displaystyle \frac{\left(1\right)}{8\pi ^2}}\left(2\mathrm{\Delta }J_\omega +{\displaystyle \frac{2}{\rho _s^2}}J_\omega +2\mathrm{\Delta }\chi {\displaystyle \frac{2}{\rho _s^2}}\chi \right)_{(x,y)}`$
$`\times \mathrm{exp}\left\{\mathrm{}\right\}`$
$`+{\displaystyle \frac{1}{i}}{\displaystyle \frac{1}{i}}\rho _s^2{\displaystyle \frac{\left(1\right)}{8\pi ^2}}\left[2\mathrm{\Delta }\delta (xx^{},yy^{})+{\displaystyle \frac{2}{\rho _s^2}}\delta (xx^{},yy^{})\right]`$
$`\times \mathrm{exp}\left\{\mathrm{}\right\}`$
At this moment we can take $`J0`$. The exponentials in the Eq(4.1) are equal to $`1`$. The auto-correlation is
$`{\displaystyle \frac{1}{\omega _v^2}}\rho _\omega (x,y)\rho _\omega (x^{},y^{})`$
$`=`$ $`{\displaystyle \frac{1}{Z_V\left[J_\omega =0\right]}}p^1{\displaystyle ๐\left[\chi \right]\mathrm{exp}\left\{\frac{1}{8\pi ^2}๐x๐y\left[\left(\chi \right)^2+\frac{1}{\rho _s^2}\chi ^2\frac{8\pi ^2}{A}Z_V^{\left(0\right)}\mathrm{cos}\left[\chi (x,y)\right]\right]\right\}}`$
$`\times \{{\displaystyle \frac{\rho _s^2}{8\pi ^2}}[2\mathrm{\Delta }\delta (xx^{},yy^{})+{\displaystyle \frac{2}{\rho _s^2}}\delta (xx^{},yy^{})]`$
$`\left({\displaystyle \frac{\rho _s^2}{8\pi ^2}}\right)^2(2\mathrm{\Delta }\chi {\displaystyle \frac{2}{\rho _s^2}}\chi )_{(x,y)}(2\mathrm{\Delta }\chi {\displaystyle \frac{2}{\rho _s^2}}\chi )_{(x^{},y^{})}\}`$
The last two lines in Eq.(4.1) (the curly bracket) arise from derivation. The normalization gives by definition
$`Z_V\left[J_\omega =0\right]`$ $``$ $`1`$
$`=`$ $`p^1{\displaystyle ๐\left[\chi \right]\mathrm{exp}\left\{\frac{1}{8\pi ^2}๐x๐y\left[\left(\chi \right)^2+\frac{1}{\rho _s^2}\chi ^2\frac{8\pi ^2}{A}Z_V^{\left(0\right)}\mathrm{cos}\left[\chi (x,y)\right]\right]\right\}}`$
The functional integration does not affect the first two terms in the curly bracket in Eq.(4.1). The functional integral of the last term represents the average of the product of two functions $`\chi `$.
$`{\displaystyle \frac{1}{\omega _v^2}}\rho _\omega (x,y)\rho _\omega (x^{},y^{})`$
$`=`$ $`{\displaystyle \frac{\rho _s^2}{4\pi ^2}}\left[\mathrm{\Delta }\delta (xx^{},yy^{})+{\displaystyle \frac{1}{\rho _s^2}}\delta (xx^{},yy^{})\right]`$
$`4\left({\displaystyle \frac{\rho _s^2}{8\pi ^2}}\right)^2\left(\mathrm{\Delta }{\displaystyle \frac{1}{\rho _s^2}}\right)_{(x,y)}\left(\mathrm{\Delta }{\displaystyle \frac{1}{\rho _s^2}}\right)_{(x^{},y^{})}\chi (x,y)\chi (x^{},y^{})`$
It is now easier if we use the Fourier representation of the fields, which we denote by the symbol $`\stackrel{~}{}`$. Writting $`๐ฑ(x,y)`$,
$$\frac{1}{\omega _v^2}\rho _\omega (x,y)\rho _\omega (x^{},y^{})=\frac{1}{\omega _v^2}๐๐ค\mathrm{exp}\left(i๐ค๐ฑ\right)๐๐ค^{}\mathrm{exp}\left(i๐ค^{}๐ฑ^{}\right)\stackrel{~}{\rho }_\omega \left(๐ค\right)\stackrel{~}{\rho }_\omega \left(๐ค^{}\right)$$
We have
$`{\displaystyle \frac{1}{\omega _v^2}}{\displaystyle ๐๐ค\mathrm{exp}\left(i๐ค๐ฑ\right)๐๐ค^{}\mathrm{exp}\left(i๐ค^{}๐ฑ^{}\right)\stackrel{~}{\rho }_\omega \left(๐ค\right)\stackrel{~}{\rho }_\omega \left(๐ค^{}\right)}`$
$`=`$ $`{\displaystyle \frac{\rho _s^2}{4\pi ^2}}{\displaystyle ๐๐ค\mathrm{exp}\left[i๐ค\left(๐ฑ๐ฑ^{}\right)\right]\left(k^2+\frac{1}{\rho _s^2}\right)}`$
$`\left({\displaystyle \frac{\rho _s^2}{4\pi ^2}}\right)^2{\displaystyle ๐๐ค\mathrm{exp}\left(i๐ค๐ฑ\right)๐๐ค^{}\mathrm{exp}\left(i๐ค^{}๐ฑ^{}\right)}`$
$`\times \left(k^2+{\displaystyle \frac{1}{\rho _s^2}}\right)\left(k^2+{\displaystyle \frac{1}{\rho _s^2}}\right)\stackrel{~}{\chi }\left(๐ค\right)\stackrel{~}{\chi }\left(๐ค^{}\right)`$
We can take a fixed reference point
$`๐ฑ^{}`$ $`=`$ $`๐`$
$`๐ฑ`$ $`=`$ $`๐+๐ฑ๐ฑ^{}`$
and write
$`{\displaystyle \frac{1}{\omega _v^2}}{\displaystyle ๐๐ค\mathrm{exp}\left[i๐ค\left(๐ฑ๐ฑ^{}\right)\right]๐๐ค^{}\mathrm{exp}\left[i\left(๐ค^{}+๐ค\right)๐\right]\stackrel{~}{\rho }_\omega \left(๐ค\right)\stackrel{~}{\rho }_\omega \left(๐ค^{}\right)}`$
$`=`$ $`{\displaystyle \frac{\rho _s^2}{4\pi ^2}}{\displaystyle ๐๐ค\mathrm{exp}\left[i๐ค\left(๐ฑ๐ฑ^{}\right)\right]\left(k^2+\frac{1}{\rho _s^2}\right)}`$
$`\left({\displaystyle \frac{\rho _s^2}{4\pi ^2}}\right)^2{\displaystyle ๐๐ค\mathrm{exp}\left[i๐ค\left(๐ฑ๐ฑ^{}\right)\right]๐๐ค^{}\mathrm{exp}\left[i\left(๐ค^{}+๐ค\right)๐\right]}`$
$`\times \left(k^2+{\displaystyle \frac{1}{\rho _s^2}}\right)\left(k^2+{\displaystyle \frac{1}{\rho _s^2}}\right)\stackrel{~}{\chi }\left(๐ค\right)\stackrel{~}{\chi }\left(๐ค^{}\right)`$
The parameter $`๐`$ has no particular role : non of our assumption has imposed a nonuniformity of the statistical properties on the plane. Therefore we can integrate Eq.(4.1) over the position $`๐`$, *i.e.* on the plane
$$\frac{1}{A}๐๐\mathrm{}$$
Obviously, this will produce in the left hand side a function $`\delta `$
$$\delta \left(๐ค^{}+๐ค\right)$$
after which the integration over the second wavenumber, $`๐ค^{}`$ , will impose
$$๐ค^{}=๐ค$$
For the first term in the right hand side, the integration over $`๐`$ will have no effect. For the second term the effect is the same as in the left hand side, *i.e.* we have $`๐ค^{}=๐ค`$. we will now replace $`๐ฑ๐ฑ^{}`$ by $`๐ฑ`$ and obtain
$`{\displaystyle \frac{1}{\omega _v^2}}{\displaystyle ๐๐ค\mathrm{exp}\left(i๐ค๐ฑ\right)\stackrel{~}{\rho }_\omega \left(๐ค\right)\stackrel{~}{\rho }_\omega \left(๐ค\right)}`$
$`=`$ $`{\displaystyle \frac{\rho _s^2}{4\pi ^2}}{\displaystyle ๐๐ค\mathrm{exp}\left(i๐ค๐ฑ\right)\left(k^2+\frac{1}{\rho _s^2}\right)}`$
$`\left({\displaystyle \frac{\rho _s^2}{4\pi ^2}}\right)^2{\displaystyle ๐๐ค\mathrm{exp}\left(i๐ค๐ฑ\right)\left(k^2+\frac{1}{\rho _s^2}\right)^2\stackrel{~}{\chi }\left(๐ค\right)\stackrel{~}{\chi }\left(๐ค\right)}`$
In physical space the correlation also reflects the statistical uniformity,
$$\frac{1}{\omega _v^2}๐๐ค\mathrm{exp}\left(i๐ค๐ฑ\right)\stackrel{~}{\rho }_\omega \left(๐ค\right)\stackrel{~}{\rho }_\omega \left(๐ค\right)=\frac{1}{\omega _v^2}\rho _\omega \left(๐ฑ\right)\rho _\omega \left(\mathrm{๐}\right)$$
The equation is
$`{\displaystyle \frac{1}{\omega _v^2}}\stackrel{~}{\rho }_\omega \left(๐ค\right)\stackrel{~}{\rho }_\omega \left(๐ค\right)`$
$`=`$ $`{\displaystyle \frac{\rho _s^2}{4\pi ^2}}\left(k^2+{\displaystyle \frac{1}{\rho _s^2}}\right)\left[1{\displaystyle \frac{\rho _s^2}{4\pi ^2}}\left(k^2+{\displaystyle \frac{1}{\rho _s^2}}\right)\stackrel{~}{\chi }\left(๐ค\right)\stackrel{~}{\chi }\left(๐ค\right)\right]`$
The second term in the bracket of Eq.(4.1) contains the two-point correlation of the function $`\chi `$ and can be obtained by explicit calculation of the functional integration in Eq.(LABEL:chichi). The same analytical problem as for the explicit calculation of $`Z_V\left[J=0\right]`$ and $`\chi \chi _๐ค\stackrel{~}{\chi }\left(๐ค\right)\stackrel{~}{\chi }\left(๐ค\right)`$ will appear later (for the turbulence scattered by the random vortices) and there we will give some details of calculation. At this moment few explanations are sufficient. For small amplitude of the auxiliary field $`\chi `$ the function $`\mathrm{cos}`$ is approximated with its first two terms
$$\mathrm{cos}\chi 1\frac{\chi ^2}{2}$$
(42)
The constant $`1`$ is only a shift of the action. However it leads to a term that depends on $`Z_V^{\left(0\right)}`$, which is integrated in the exponential over all volume *i.e.* the area $`A`$ on the plane. In detail, replacing Eq.(42) in Eq.(LABEL:eq525)
$`Z_V\left[J_\omega =0\right]`$ $``$ $`p^1{\displaystyle ๐\left[\chi \right]}`$
$`\times \mathrm{exp}\left[{\displaystyle \frac{1}{8\pi ^2}}{\displaystyle ๐x๐y\left(\frac{8\pi ^2}{A}Z_V^{\left(0\right)}\right)}\right]`$
$`\times \mathrm{exp}\left\{{\displaystyle \frac{1}{8\pi ^2}}{\displaystyle ๐x๐y\left[\left(\chi \right)^2+\frac{1}{\rho _s^2}\chi ^2+\left(\frac{8\pi ^2}{A}Z_V^{\left(0\right)}\right)\frac{\chi ^2}{2}\right]}\right\}`$
The first factor can be taken outside the functional integration
$`\mathrm{exp}\left\{{\displaystyle \frac{1}{8\pi ^2}}{\displaystyle ๐x๐y\left(\frac{8\pi ^2}{A}Z_V^{\left(0\right)}\right)}\right\}`$
$`=`$ $`\mathrm{exp}\left[Z_V^{\left(0\right)}\right]`$
Since it is determined by the non-interacting vortices, it must exist even if we would neglect completely the interaction between the vortices taking $`\chi 0`$. The rest of the Eq.(4.1) is the Gaussian functional integral
$`p^1{\displaystyle ๐\left[\chi \right]}`$
$`\times \mathrm{exp}\left\{{\displaystyle \frac{1}{8\pi ^2}}{\displaystyle ๐x๐y\left[\left(\chi \right)^2+\frac{1}{\rho _s^2}\chi ^2+\left(\frac{8\pi ^2}{A}Z_V^{\left(0\right)}\right)\frac{\chi ^2}{2}\right]}\right\}`$
$`=`$ $`p^1{\displaystyle ๐\left[\chi \right]}`$
$`\times \mathrm{exp}\left\{{\displaystyle \frac{1}{8\pi ^2}}{\displaystyle ๐x๐y\chi (x,y)\left[\mathrm{\Delta }+\left(\frac{1}{\rho _s^2}+\frac{4\pi ^2}{A}Z_V^{\left(0\right)}\right)\right]\chi (x,y)}\right\}`$
$`=`$ $`q\left[det\left(\mathrm{\Delta }+\alpha ^2\right)\right]^{1/2}`$
where
$$\alpha ^2\frac{1}{\rho _s^2}+\frac{4\pi ^2}{A}Z_V^{\left(0\right)}$$
The determinant can be calculated explicitely, by the product of the eigenvalues of the operator. This product (besides an infinite factor that will disappear) is convergent. However we keep this formal expression
$$Z_V=Z_V\left[J_\omega =0\right]=q\mathrm{exp}\left[Z_V^{\left(0\right)}\right]\left[det\left(\mathrm{\Delta }+\alpha ^2\right)\right]^{1/2}$$
(47)
The second term in the bracket of Eq.(4.1) contains the two-point correlation of the function $`\chi `$ and can be obtained by explicit calculation of the functional integration in Eq.(4.1). The same analytical problem as for the explicit calculation of $`Z_V\left[J_\omega =0\right]`$ (Eq.(LABEL:eq4592)) and $`\chi \chi _๐ค\stackrel{~}{\chi }\left(๐ค\right)\stackrel{~}{\chi }\left(๐ค\right)`$ will appear later (for the turbulence scattered by the random vortices) and there we will give some details of calculation. At this moment few explanations are sufficient. For small amplitude of the auxiliary field $`\chi `$ the function $`\mathrm{cos}`$ is approximated with its first two terms
$$\mathrm{cos}\chi 1\frac{\chi ^2}{2}$$
(48)
The constant $`1`$ is only a shift of the action. However it leads to a term that depends on $`Z_V^{\left(0\right)}`$, which is integrated in the exponential over all volume *i.e.* the area $`A`$ on the plane. In detail, replacing Eq.(42) in Eq.(LABEL:eq525)
$`Z_V\left[J_\omega =0\right]`$ $``$ $`p^1{\displaystyle ๐\left[\chi \right]\mathrm{exp}\left[\frac{1}{8\pi ^2}๐x๐y\left(\frac{8\pi ^2}{A}Z_V^{\left(0\right)}\right)\right]}`$
$`\times \mathrm{exp}\left\{{\displaystyle \frac{1}{8\pi ^2}}{\displaystyle ๐x๐y\left[\left(\chi \right)^2+\frac{1}{\rho _s^2}\chi ^2+\left(\frac{8\pi ^2}{A}Z_V^{\left(0\right)}\right)\frac{\chi ^2}{2}\right]}\right\}`$
The first factor can be taken outside the functional integration and is $`\mathrm{exp}\left[Z_V^{\left(0\right)}\right]`$. Since it is determined by the non-interacting vortices, it must exist even if we would neglect completely the interaction between the vortices taking $`\chi 0`$. The rest of the Eq.(4.1) is the Gaussian functional integral. Introducing the notation $`\alpha ^21/\rho _s^2+4\pi ^2Z_V^{\left(0\right)}/A`$
$$Z_V\left[J_\omega =0\right]=q\mathrm{exp}\left[Z_V^{\left(0\right)}\right]\left[det\left(\mathrm{\Delta }+\alpha ^2\right)\right]^{1/2}$$
(49)
and $`q`$ is a constant. The determinant can be calculated explicitely, by the product of the eigenvalues of the operator. We need this explicit expression because we need the functional dependence of $`Z_V=Z_V\left[J_\omega =0\right]`$ on $`Z_V^{\left(0\right)}`$ as results from Eq.(4.1). As will become clear later, the factor with the determinant, which comes from the influence of the fluctuating shape of a vortex on the correlations of the vorticity of a gas of vortices in the plane, is affected by a factor $`\rho _s^2/A`$, which is small compared with the exponential in Eq.(47). Therefore we calculate in a one dimensional cartezian approximation the eigenvalues, instead of a cylindrical problem. We have to solve
$$\left(\frac{d^2}{dx^2}+\alpha ^2\right)\eta _n(x,y)=\lambda _n^\eta \eta _n(x,y)$$
on an interval $`L`$. The eigenvalues are $`\lambda _n^\eta =\left(2\pi n/L\right)^2+\alpha ^2`$ , where $`n`$ is an integer, and we obtain
$`det\left(\mathrm{\Delta }+\alpha ^2\right)`$ $`=`$ $`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\lambda _n^\eta ={\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\left[\left(2\pi n/L\right)^2+\alpha ^2\right]`$
$`=`$ $`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\left(2\pi n/L\right)^2{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\left[1+{\displaystyle \frac{\alpha ^2L^2/\left(2\pi \right)^2}{n^2}}\right]`$
$`=`$ $`{\displaystyle \frac{\mathrm{sinh}\left(\alpha L/2\right)}{\alpha L/2}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\left(2\pi n/L\right)^2`$
The infinite product is eliminated since we always use the ratios of $`Z_V\left[J_\omega \right]`$ and $`Z_V\left[J_\omega =0\right]`$. We also take $`L=\sqrt{A}`$ and obtain
$$Z_V=q\mathrm{exp}\left[Z_V^{\left(0\right)}\right]\left\{\frac{\mathrm{sinh}\left[\left(A\rho _s^2+4\pi ^2Z_V^{\left(0\right)}\right)^{1/2}/2\right]}{\left(A\rho _s^2+4\pi ^2Z_V^{\left(0\right)}\right)^{1/2}/2}\right\}^{1/2}$$
(51)
The quantity $`A\rho _s^2`$ is very large and an approximation is possible
$`{\displaystyle \frac{\mathrm{sinh}\left[\left(A\rho _s^2+4\pi ^2Z_V^{\left(0\right)}\right)^{1/2}/2\right]}{\left(A\rho _s^2+4\pi ^2Z_V^{\left(0\right)}\right)^{1/2}/2}}`$
$`=`$ $`{\displaystyle \frac{1}{\rho _s^1\sqrt{A}\left(1+\frac{1}{2}\frac{4\pi ^2Z_V^{\left(0\right)}}{A\rho _s^2}\right)}}\mathrm{exp}\left[{\displaystyle \frac{\rho _s^1\sqrt{A}}{2}}\left(1+{\displaystyle \frac{1}{2}}{\displaystyle \frac{4\pi ^2Z_V^{\left(0\right)}}{A\rho _s^2}}\right)\right]`$
$``$ $`{\displaystyle \frac{\mathrm{exp}\left(\rho _s^1\sqrt{A}/2\right)}{\rho _s^1\sqrt{A}}}\mathrm{exp}\left({\displaystyle \frac{\pi ^2Z_V^{\left(0\right)}}{\sqrt{A}\rho _s^1}}\right)\left(1{\displaystyle \frac{2\pi ^2Z_V^{\left(0\right)}}{A\rho _s^2}}\right)`$
The first factor is large but constant and can absorbed into the coefficient $`q`$. We get in Eq.(51)
$`Z_V`$ $`=`$ $`q\mathrm{exp}\left[Z_V^{\left(0\right)}\right]\mathrm{exp}\left({\displaystyle \frac{\pi ^2Z_V^{\left(0\right)}}{2\sqrt{A}\rho _s^1}}\right)\left(1{\displaystyle \frac{2\pi ^2Z_V^{\left(0\right)}}{A\rho _s^2}}\right)^{1/2}`$
$`=`$ $`q\mathrm{exp}\left[Z_V^{\left(0\right)}\left(1{\displaystyle \frac{\pi ^2}{2\sqrt{A}\rho _s^1}}\right)\right]\left(1+{\displaystyle \frac{\pi ^2Z_V^{\left(0\right)}}{A\rho _s^2}}\right)`$
We can neglect the second term in the first exponential
$$Z_V=Z_V\left[J_\omega =0\right]=q\left(1+\frac{\pi ^2Z_V^{\left(0\right)}}{A\rho _s^2}\right)\mathrm{exp}\left[Z_V^{\left(0\right)}\right]$$
(53)
The second term in Eq.(4.1), *i.e.* the auto-correlation of $`\chi `$ in $`๐ค`$-space, may be calculated starting from the real-space correlation
$`\chi (x,y)\chi (x^{},y^{})`$
$`=`$ $`{\displaystyle \frac{1}{Z_V\left[J=0\right]}}p^1\mathrm{exp}\left[Z_V^{\left(0\right)}\right]{\displaystyle ๐\left[\chi \right]\chi (x,y)\chi (x^{},y^{})}`$
$`\times \mathrm{exp}\left\{{\displaystyle \frac{1}{8\pi ^2}}{\displaystyle ๐x๐y\chi (x,y)\left(\mathrm{\Delta }+\alpha ^2\right)\chi (x,y)}\right\}`$
As usual we return to the form of Eq.(4.1), and only for this step, we insert an external current $`J_e(x,y)`$ interacting with $`\chi `$. The auxilliary functional is denoted $`Z_\chi \left[J_e\right]`$
$`Z_\chi \left[J_e\right]`$ $`=`$ $`{\displaystyle \frac{1}{Z_V\left[J=0\right]}}p^1\mathrm{exp}\left[Z_V^{\left(0\right)}\right]{\displaystyle ๐\left[\chi \right]}`$
$`\times \mathrm{exp}\left\{{\displaystyle \frac{1}{8\pi ^2}}{\displaystyle ๐x๐y\left[\chi (x,y)\left(\mathrm{\Delta }+\alpha ^2\right)\chi (x,y)+J_e\chi \right]}\right\}`$
$`=`$ $`{\displaystyle \frac{1}{Z_V\left[J=0\right]}}p^1\mathrm{exp}\left[Z_V^{\left(0\right)}\right]{\displaystyle ๐\left[\varphi \right]}`$
$`\times \mathrm{exp}\{{\displaystyle \frac{1}{8\pi ^2}}{\displaystyle }dxdy\varphi (x,y)(\mathrm{\Delta }+\alpha ^2)\varphi (x,y)`$
$`{\displaystyle \frac{1}{8\pi ^2}}{\displaystyle }dxdy{\displaystyle \frac{1}{4}}J_e(x,y)(\mathrm{\Delta }+\alpha ^2)^1J_e(x,y)\}`$
To obtain the above equation we have made a change of variables $`\chi \varphi =\chi +\frac{1}{2}\left(\mathrm{\Delta }+\alpha ^2\right)^1J_e`$ of Jacobian $`1`$. The functional integration over $`\varphi `$ can now be carried out and the rest of the expression at the exponent appears in a factor
$$\left[det\left(\mathrm{\Delta }+\alpha ^2\right)\right]^{1/2}\mathrm{exp}\left\{\frac{1}{8\pi ^2}๐x๐y\frac{1}{4}J_e(x,y)\left(\mathrm{\Delta }+\alpha ^2\right)^1J_e(x,y)\right\}$$
where the symbol $``$ means that there also result constant factors. But these are the same as those contained in the factor $`q`$ introduced in the Eq.(47). We then have
$`Z_\chi \left[J_e\right]`$ $`=`$ $`{\displaystyle \frac{1}{Z_V\left[J=0\right]}}p^1\mathrm{exp}\left[Z_V^{\left(0\right)}\right]{\displaystyle ๐\left[\chi \right]}`$
$`\times \mathrm{exp}\left\{{\displaystyle \frac{1}{8\pi ^2}}{\displaystyle ๐x๐y\left[\chi (x,y)\left(\mathrm{\Delta }+\alpha ^2\right)\chi (x,y)+J_e\chi \right]}\right\}`$
$`=`$ $`p^1\mathrm{exp}\left\{{\displaystyle \frac{1}{8\pi ^2}}{\displaystyle ๐x๐y\frac{1}{4}J_e(x,y)\left(\mathrm{\Delta }+\alpha ^2\right)^1J_e(x,y)}\right\}`$
where we have taken into account Eq.(47). The correlation is
$`\chi (x,y)\chi (x^{},y^{})`$
$`=`$ $`{\displaystyle \frac{1}{Z_\chi \left[J_e=0\right]}}{\displaystyle \frac{\delta ^2Z_\chi \left[J_e\right]}{\delta J_e(x,y)\delta J_e(x^{},y^{})}}|_{J_e=0}`$
$`=`$ $`{\displaystyle \frac{1}{Z_\chi \left[J_e=0\right]}}p^1{\displaystyle \frac{\delta }{\delta J_e(x^{},y^{})}}\left[\left({\displaystyle \frac{1}{8\pi ^2}}2\left(\mathrm{\Delta }+\alpha ^2\right)^1J_e(x,y)\right)\mathrm{exp}\left\{\mathrm{}\right\}\right]_{J_e=0}`$
$`=`$ $`{\displaystyle \frac{1}{Z_\chi \left[J_e=0\right]}}p^1[{\displaystyle \frac{1}{4\pi ^2}}(\mathrm{\Delta }+\alpha ^2)^1\delta (๐ฑ๐ฑ^{})\mathrm{exp}\left\{\mathrm{}\right\}`$
$`+({\displaystyle \frac{1}{8\pi ^2}}2(\mathrm{\Delta }+\alpha ^2)^1J_e(x,y))({\displaystyle \frac{1}{8\pi ^2}}2(\mathrm{\Delta }+\alpha ^2)^1J_e(x^{},y^{}))\mathrm{exp}\left\{\mathrm{}\right\}]_{J_e=0}`$
(the factor $`\mathrm{exp}\left[Z_V^{\left(0\right)}\right]`$ has not been written since it disappears). Finally we have
$$\chi (x,y)\chi (x^{},y^{})=\frac{1}{4\pi ^2}\left(\mathrm{\Delta }+\alpha ^2\right)^1\delta \left(๐ฑ๐ฑ^{}\right)$$
(54)
or
$$\chi \chi _๐ค=\left[\frac{1}{4\pi ^2}\rho _s^2\left(k^2+\frac{1}{\rho _s^2}+\frac{4\pi ^2}{A}Z_V^{\left(0\right)}\right)\right]^1$$
(55)
Introducing this in Eq.(4.1) the correlation of the field of the scalar potential will have from this a contribution (with normalization)
$`{\displaystyle \frac{1}{\varphi _0^2}}\varphi _V\varphi _V_๐ค^{vort}`$ $`=`$ $`{\displaystyle \frac{1}{k^2\rho _s^2}}\left(1+{\displaystyle \frac{1}{k^2\rho _s^2}}\right)`$
$`\times {\displaystyle \frac{1}{4\pi ^2}}\left[1{\displaystyle \frac{1+k^2\rho _s^2}{1+k^2\rho _s^2+Z_V^{\left(0\right)}4\pi ^2\rho _s^2/A}}\right]`$
which represents, as said, the first term of the Eq.(4.1).
The other terms in the correlation $`\varphi _V\varphi _V`$ (Eq.(4.1)) are due to the fluctuation of the form of the generic vortex from interaction with the background turbulence. We calculate, using Eq.(LABEL:eq5237), the derivatives
$$\frac{\delta Z_V}{\delta Z_V^{\left(0\right)}}=q\left(1+\frac{\pi ^2}{A\rho _s^2}+\frac{\pi ^2Z_V^{\left(0\right)}}{A\rho _s^2}\right)\mathrm{exp}\left[Z_V^{\left(0\right)}\right]$$
(57)
and
$$\frac{\delta ^2Z_V}{\delta \left[Z_V^{\left(0\right)}\right]^2}=q\left(1+\frac{2\pi ^2}{A\rho _s^2}+\frac{\pi ^2Z_V^{\left(0\right)}}{A\rho _s^2}\right)\mathrm{exp}\left[Z_V^{\left(0\right)}\right]$$
(58)
Now we have to calculate the derivatives of $`Z_V^{\left(0\right)}`$ to $`J`$ as shown by the last two terms Eq.(4.1). This part will be added after $`Z_V^{\left(0\right)}`$ is calculated, in the next subsection.
### 4.2 A single vortex interacting with a turbulent environment
The partition function for a single vortex in interaction with turbulence is defined as
$$Z_V^{\left(0\right)}=๐ฉ^1D\left[\chi \right]D\left[\varphi \right]\mathrm{exp}\left\{S_V[\chi ,\varphi ]\right\}$$
(59)
and the equation is the stationary form of the equation used in Ref.,
$$F\left[\varphi \right]_{}^2\phi \alpha \phi \beta \phi ^2$$
(60)
The density of Lagrangean
$$[\chi (x,y),\varphi (x,y)]=\chi \left(_{}^2\phi \alpha \phi \beta \phi ^2\right)$$
(61)
is obtained from the Martin-Siggia-Rose method for classical stochastic systems. Then the action is
$$S_V[\chi ,\varphi ]=๐x๐y[\chi (x,y),\varphi (x,y)]$$
(62)
As usual we introduce the interaction with the external current $`๐=(J_\varphi ,J_\chi )`$. However since there is no need to calculate functional derivatives to $`\chi `$, we can only keep $`JJ_\varphi `$. For uniformity of notation in this paper, we will not use the factor $`i`$ in front of the action in contrast with . We have then
$$Z_V^{\left(0\right)}\left[J\right]=๐ฉ^1D\left[\chi \right]D\left[\varphi \right]\mathrm{exp}\left\{๐x๐y\left[\chi \left(_{}^2\phi \alpha \phi \beta \phi ^2\right)+J\varphi \right]\right\}$$
(63)
We need the explicit expression of the functional integral Eq.(63) and this has been obtained in the references and . For convenience we will recall briefly the steps of the calculation restricting to the results we need in the present work.
To calculate the generating functional of the vortex in the background turbulence we proceed in two steps: we solve the Euler-Lagrange equation for the action (62) obtaining the configuration of the system which extremises this action; further, we expand the action to second order in the fluctuations around this extremum (which will include the turbulent field) and integrate. The Euler-Lagrange equations have the solutions
$`\phi _{Js}(x,y)`$ $``$ $`\phi _s(x,y)`$ (64)
$`\chi _{Js}(x,y)`$ $`=`$ $`\phi _s(x,y)+\stackrel{~}{\chi }_J(x,y)`$
The first is the static form of the solution Eq.(10) and does not depend on $`J`$. The dual function is $`\phi _s(x,y)`$ plus a term resulting from the excitation by $`J`$ in its equation. This additional term $`\stackrel{~}{\chi }_J(x,y)`$ is calculated by the perturbation of the KdV soliton solution according to the modification of the Inverse Scattering Transform when an inhomogeneous term (*i.e.* $`J`$) is included. The action functional is calculated for these two functions
$$S_{Vs}\left[J\right]S_V[\phi _s(x,y),\phi _s(x,y)+\stackrel{~}{\chi }_J(x,y)]$$
Then the first part of our calculation is $`Z_V^{\left(0\right)}\left[J\right]๐ฉ^1\mathrm{exp}\left\{S_{Vs}\left[J\right]\right\}`$. Expanding the action around this extremum
$$S_V[\chi ,\phi ;J]=S_V[\phi _{Js},\chi _{Js}]+\frac{1}{2}\left(\frac{\delta ^2S_V\left[J\right]}{\delta \phi \delta \chi }|_{\phi _{Js},\chi _{Js}}\right)\delta \phi \delta \chi $$
we calculate the Gaussian integral and obtain
$$Z_V^{\left(0\right)}\left[J\right]=๐ฉ^1\mathrm{exp}\left\{S_{Vs}\left[J\right]\right\}\left[det\left(\frac{\delta ^2S_V\left[J\right]}{\delta \phi \delta \chi }|_{\phi _{Js},\chi _{Js}}\right)\right]^{1/2}$$
If we can neglect the advection of vortices by large scale wave-like fluctuations, we can calculate the determinant since the product of the eigenvalues converges without the need for regularization. The result is
$$Z_V^{\left(0\right)}\left[J\right]=๐ฉ^1\mathrm{exp}\left\{S_{Vs}\left[J\right]\right\}AB$$
(65)
Where
$$A=A\left[J\right]\left[\frac{\beta /2}{\mathrm{sinh}\left(\beta /2\right)}\right]^{1/4}$$
(66)
$$B=B\left[J\right]\left[\frac{\sigma /2}{\mathrm{sin}\left(\sigma /2\right)}\right]^{1/2}$$
(67)
The eigenvalue problem depends functionally on $`\stackrel{~}{\chi }_J(x,y)`$ which implies that $`\beta `$ and $`\sigma `$ depend on $`J`$. Their expressions can be found in . With those detailed formulas we can consider that we have the necessary knowledge to proceed to the calculation of the functional derivatives of $`Z_V^{\left(0\right)}\left[J\right]`$ at $`J`$, using Eq.(65).
$$\frac{1}{Z_V^{\left(0\right)}\left[J=0\right]}\frac{\delta Z_V^{\left(0\right)}\left[J\right]}{\delta J}=\frac{\delta S_{Vs}\left[J\right]}{\delta J}+\frac{1}{A}\frac{\delta A}{\delta J}+\frac{1}{B}\frac{\delta B}{\delta J}$$
(68)
For the present problem it is sufficient to take the first term as Eq.(10)
$$\frac{\delta S_{Vs}\left[J\right]}{\delta J}|_{J=0}\varphi _s$$
(69)
The next two terms in Eq.(68) represent the averaged, systematic modification of the shape of the field around the vortex due to the mutual interaction. They will be equally neglected, assuming that the main effect is contained in the dispersion of the fluctuations of the shape of the vortex interacting with the random field (which actually is our main concern here). Then
$$\frac{1}{Z_V^{\left(0\right)}\left[J=0\right]}\frac{\delta Z_V^{\left(0\right)}\left[J\right]}{\delta J}\varphi _s$$
(70)
The second derivative to the excitations in two points $`y_1`$ and $`y_2`$ is
$`{\displaystyle \frac{1}{Z_V^{\left(0\right)}\left[J=0\right]}}{\displaystyle \frac{\delta ^2Z_V^{\left(0\right)}\left[J\right]}{\delta J\left(y_2\right)\delta J\left(y_1\right)}}`$
$`=`$ $`{\displaystyle \frac{\delta S_{Vs}\left[J\right]}{\delta J\left(y_2\right)}}{\displaystyle \frac{\delta S_{Vs}\left[J\right]}{\delta J\left(y_1\right)}}+{\displaystyle \frac{\delta ^2S_{Vs}\left[J\right]}{\delta J\left(y_2\right)\delta J\left(y_1\right)}}`$
$`+{\displaystyle \frac{1}{A}}{\displaystyle \frac{\delta A}{\delta J\left(y_2\right)}}{\displaystyle \frac{\delta S_{Vs}\left[J\right]}{\delta J\left(y_1\right)}}+{\displaystyle \frac{1}{B}}{\displaystyle \frac{\delta B}{\delta J\left(y_2\right)}}{\displaystyle \frac{\delta S_{Vs}\left[J\right]}{\delta J\left(y_1\right)}}`$
$`+{\displaystyle \frac{1}{A}}{\displaystyle \frac{\delta A}{\delta J\left(y_1\right)}}{\displaystyle \frac{\delta S_{Vs}\left[J\right]}{\delta J\left(y_2\right)}}+{\displaystyle \frac{1}{B}}{\displaystyle \frac{\delta B}{\delta J\left(y_1\right)}}{\displaystyle \frac{\delta S_{Vs}\left[J\right]}{\delta J\left(y_2\right)}}`$
$`+{\displaystyle \frac{1}{A}}{\displaystyle \frac{\delta A}{\delta J\left(y_1\right)}}{\displaystyle \frac{1}{B}}{\displaystyle \frac{\delta B}{\delta J\left(y_2\right)}}+{\displaystyle \frac{1}{A}}{\displaystyle \frac{\delta A}{\delta J\left(y_2\right)}}{\displaystyle \frac{1}{B}}{\displaystyle \frac{\delta B}{\delta J\left(y_1\right)}}`$
$`+{\displaystyle \frac{1}{A}}{\displaystyle \frac{\delta ^2A}{\delta J\left(y_2\right)\delta J\left(y_1\right)}}+{\displaystyle \frac{1}{B}}{\displaystyle \frac{\delta ^2B}{\delta J\left(y_2\right)\delta J\left(y_1\right)}}`$
Since we have assumed as an acceptable approximation to neglect the averaged change produced by the turbulence on the soliton shape the second term in the RHS is zero. For the first term we use Eq.(69). We have
$`{\displaystyle \frac{1}{Z_V^{\left(0\right)}\left[J=0\right]}}{\displaystyle \frac{\delta ^2Z_V^{\left(0\right)}\left[J\right]}{\delta J\left(y_2\right)\delta J\left(y_1\right)}}`$
$`=`$ $`\varphi _s\left(y_2\right)\varphi _s\left(y_1\right)+`$
$`+\varphi _s\left(y_1\right)\left[{\displaystyle \frac{\delta }{\delta J\left(y_2\right)}}\mathrm{ln}A+{\displaystyle \frac{\delta }{\delta J\left(y_2\right)}}\mathrm{ln}B\right]`$
$`+\varphi _s\left(y_2\right)\left[{\displaystyle \frac{\delta }{\delta J\left(y_1\right)}}\mathrm{ln}A+{\displaystyle \frac{\delta }{\delta J\left(y_1\right)}}\mathrm{ln}B\right]`$
$`+{\displaystyle \frac{\delta \mathrm{ln}A}{\delta J\left(y_1\right)}}{\displaystyle \frac{\delta \mathrm{ln}B}{\delta J\left(y_2\right)}}+{\displaystyle \frac{\delta \mathrm{ln}A}{\delta J\left(y_2\right)}}{\displaystyle \frac{\delta \mathrm{ln}B}{\delta J\left(y_1\right)}}`$
$`+{\displaystyle \frac{1}{A}}{\displaystyle \frac{\delta ^2A}{\delta J\left(y_2\right)\delta J\left(y_1\right)}}+{\displaystyle \frac{1}{B}}{\displaystyle \frac{\delta ^2B}{\delta J\left(y_2\right)\delta J\left(y_1\right)}}`$
The expressions are complicated (see the Appendix of Ref.) and some numerical calculation of these expression is unavoidable. For small amplitude of the turbulent field the expression can be rewritten
$`{\displaystyle \frac{1}{Z_V^{\left(0\right)}\left[J=0\right]}}{\displaystyle \frac{\delta ^2Z_V^{\left(0\right)}\left[J\right]}{\delta J\left(y_2\right)\delta J\left(y_1\right)}}|_{J=0}`$
$`=`$ $`\varphi _s\left(y_2\right)\varphi _s\left(y_1\right)\left[1+f\left(y\right)\right]`$
*i.e.* in a form that expresses the fact that the two-point correlation is basically the auto-correlation of the potential of the exact soliton modified by a function $`f`$ which collects the contributions from the interaction with the random field. In $`k`$-space we have
$$\frac{1}{Z_V^{\left(0\right)}\left[J=0\right]}\frac{\delta ^2Z_V^{\left(0\right)}\left[J\right]}{\delta J\left(y_2\right)\delta J\left(y_1\right)}|_{J=0,๐ค}=\varphi _s\left(๐ค\right)\varphi _s\left(๐ค\right)\left[1+f\left(๐ค\right)\right]$$
At the limit where we do not expand the action to include configurations resulting from the interaction vortex-turbulence, we have $`\sigma 0`$ and $`\beta 0`$ and it results $`A=B=1`$. In this case $`f0`$.
The detailed expressions of these terms are given in the paper , . For the purpose of comparisons we will express the spectrum as
$$\frac{1}{\varphi _0^2}\varphi _V\varphi _V_๐ค=S\left(๐ค\right)\left[1+f\left(๐ค\right)\right]$$
(72)
where $`\varphi _0`$ is amplitude of a vortex, $`S\left(๐ค\right)`$ has been derived by Meiss and Horton
$$S\left(๐ค\right)=\left\{12\sqrt{2}\pi ^{3/2}k\rho _s\frac{u}{v_{}}\mathrm{csc}h\left[\frac{\pi k\rho _s}{\left(1v_{}/u\right)^{1/2}}\right]\right\}^2$$
(73)
and $`f\left(๐ค\right)`$ is function that is the correction to the Fourier transform of the squared secant-hyperbolic, produced by the turbulent waves.
Before proceeding further with the calculations based on the Eq.(72) we need to discuss the formal term $`f\left(๐ค\right)`$. Since this term represents the difference from the simple isolated vortex to the vortex perturbed by turbulence, one would like to have a quantitative connection between the amplitude of this term and at least two elements characterizing the background turbulence: (1) the amplitude and (2) the spectrum.
In the way we have conducted the calculations of the generating functional Eqs.(65), (66), (67) the new terms in the expression of the auto-correlation due to the factors $`A`$ and $`B`$ are expressed in real space, not in Fourier space. They are obtained from the product of eigenvalues of the operator representing the second order functional derivative of the action, *i.e.* they are connected with the geometry of the function space around the exact, vortex, nonlinear solution. The determinant of the operator $`\delta ^2S_V/\left(\delta \phi \delta \chi \right)`$ may be seen as a volume in the function space, centered on the vortex solution. The inverse of any eigenvalue gives an idea of the extension along a particular direction (eigenfunction) in function space. When an eigenvalue is very small, the operator almost vanishes on functions along that direction. At the limit this is a *zero mode* and corresponds to a translational symmetry of the physical system along that direction. The correlations depend on the sensitivity of this volume (the product of the eigenvalues) on the excitation $`J`$ applied on the system. The excitation is first manifested in the appearence of $`\stackrel{~}{\chi }_J(x,y)`$. This one consists of a part that will modify the shape of the exact vortex plus the oscillating tail generated when a soliton is perturbed. The latter can be considered as a component of the background turbulence. The โpropagationโ of the influence from an excitation $`J`$ can be summarised symbolically in the chain : $`J\stackrel{~}{\chi }_J(x,y)`$ eigenvalues of the operator $`\delta ^2S_V/\left(\delta \phi \delta \chi \right)`$ $`A`$ and $`B`$ (or $`\sigma `$ and $`\beta `$). The expressions of $`\sigma `$ and $`\beta `$ can be found in .
Now we can return to the Eq.(4.1). Using Eq.(70) the second term is
$`{\displaystyle \frac{1}{Z_V\left[J=0\right]}}{\displaystyle \frac{\delta ^2Z_V}{\delta \left(Z_V^{\left(0\right)}\right)^2}}{\displaystyle \frac{\delta Z_V^{\left(0\right)}}{\delta J(x,y)}}{\displaystyle \frac{\delta Z_V^{\left(0\right)}}{\delta J(x^{},y^{})}}`$
$`=`$ $`{\displaystyle \frac{1}{\mathrm{exp}\left[Z_V^{\left(0\right)}\right]q\left(1+Z_V^{\left(0\right)}\rho _s^2\pi ^2/A\right)}}\mathrm{exp}\left[Z_V^{\left(0\right)}\right]q\left[1+\left(2+Z_V^{\left(0\right)}\right)\rho _s^2\pi ^2/A\right]`$
$`\times \varphi _s(x,y)\varphi _s(x^{},y^{})`$
$`=`$ $`\left(1+{\displaystyle \frac{2\rho _s^2\pi ^2/A}{1+Z_V^{\left(0\right)}\rho _s^2\pi ^2/A}}\right)\varphi _s(x,y)\varphi _s(x^{},y^{})`$
The third term in Eq.(4.1) is
$`{\displaystyle \frac{1}{Z_V\left[j=0\right]}}{\displaystyle \frac{\delta Z_V}{\delta Z_V^{\left(0\right)}}}{\displaystyle \frac{\delta ^2Z_V^{\left(0\right)}}{\delta J(x,y)\delta J(x^{},y^{})}}`$
$`=`$ $`{\displaystyle \frac{1}{\mathrm{exp}\left[Z_V^{\left(0\right)}\right]q\left(1+Z_V^{\left(0\right)}\rho _s^2\pi ^2/A\right)}}q\left(1+{\displaystyle \frac{\pi ^2}{A\rho _s^2}}+{\displaystyle \frac{\pi ^2Z_V^{\left(0\right)}}{A\rho _s^2}}\right)\mathrm{exp}\left[Z_V^{\left(0\right)}\right]`$
$`\times \varphi _s(x,y)\varphi _s(x^{},y^{})\left(1+f\right)`$
$`=`$ $`\left(1+{\displaystyle \frac{\rho _s^2\pi ^2/A}{1+Z_V^{\left(0\right)}\rho _s^2\pi ^2/A}}\right)\varphi _s(x,y)\varphi _s(x^{},y^{})\left(1+f\right)`$
In $`๐ค`$ space the two contributions reads
$`{\displaystyle \frac{1}{Z_V\left[J=0\right]}}{\displaystyle \frac{\delta ^2Z_V}{\delta \left(Z_V^{\left(0\right)}\right)^2}}{\displaystyle \frac{\delta Z_V^{\left(0\right)}}{\delta J(x,y)}}{\displaystyle \frac{\delta Z_V^{\left(0\right)}}{\delta J(x^{},y^{})}}+{\displaystyle \frac{1}{Z_V\left[j=0\right]}}{\displaystyle \frac{\delta Z_V}{\delta Z_V^{\left(0\right)}}}{\displaystyle \frac{\delta ^2Z_V^{\left(0\right)}}{\delta J(x,y)\delta J(x^{},y^{})}}`$
$`=`$ $`\varphi _0^2S\left(๐ค\right)\left[1+{\displaystyle \frac{2\rho _s^2\pi ^2/A}{1+Z_V^{\left(0\right)}\rho _s^2\pi ^2/A}}+\left(1+f\right)\left(1+{\displaystyle \frac{\rho _s^2\pi ^2/A}{1+Z_V^{\left(0\right)}\rho _s^2\pi ^2/A}}\right)\right]`$
$`=`$ $`\varphi _0^2S\left(๐ค\right)\left(2+f+{\displaystyle \frac{3+f}{A/\rho _s^2+Z_V^{\left(0\right)}}}\right)`$
The results for Eq.(4.1) can now be collected
$`{\displaystyle \frac{1}{\varphi _0^2}}\varphi _V\varphi _V_๐ค^{vort+cs}`$
$`=`$ $`{\displaystyle \frac{1}{k^2\rho _s^2}}\left(1+{\displaystyle \frac{1}{k^2\rho _s^2}}\right){\displaystyle \frac{1}{8\pi ^2}}\left(1{\displaystyle \frac{\rho _s^2k^2+1}{\rho _s^2k^2+1+Z_V^{\left(0\right)}4\pi ^2\rho _s^2/A}}\right)`$
$`+S\left(๐ค\right)\left(2+f+{\displaystyle \frac{3+f}{A/\rho _s^2+Z_V^{\left(0\right)}}}\right)`$
We can make few remarks here. If the arbitrary position in plane of the vortices and the interaction between physical vortices were neglected, the only term that would persist is $`\mathrm{exp}\left[Z_V^{\left(0\right)}\right]`$. The first $`๐ค`$-dependent factor in Eq.(LABEL:eq54) comes from assuming that a statistical ensemble of realizations of the vorticity field is generated from the random positions in plane of the vortices, even reduced at a $`\delta `$-type shape. In practical terms this may be represented as follows: in a plane, an ensemble of vortices can be placed at arbitrary positions. We construct the statistical ensemble of the realizations of this stochastic system. If we measure in one point the field, it will be zero for most of the realizations and it will be finite when it happens that a vortex is there. This is a random variable. Now, if we measure in two points and collect the results for all realizations, the statistical properties of this quantity (the two-point auto-correlation) has a Fourier transform that is given by the two factors multiplying the square bracket in Eq.(LABEL:eq54), divided to $`k^4`$ (since we have the auto-correlation of the vorticity). When the interaction is considered, the factor in the curly bracket appears.
## 5 Random field influenced by vortices with random positions
Consider the equation
$$F\left[\varphi \right]_{}^2\phi \alpha \phi \beta \phi ^2=0$$
(76)
and extract from the total function the part that is due to the vortices
$$\phi (x,y)=\underset{a=1}{\overset{N}{}}\varphi _s^{\left(a\right)}(x,y)+\varphi (x,y)$$
(77)
Replacing in the equation we have
$`_{}^2\left[{\displaystyle \underset{a=1}{\overset{N}{}}}\varphi _s^{\left(a\right)}(x,y)\right]\alpha {\displaystyle \underset{a=1}{\overset{N}{}}}\varphi _s^{\left(a\right)}(x,y)\beta \left[{\displaystyle \underset{a=1}{\overset{N}{}}}\varphi _s^{\left(a\right)}(x,y)\right]^2`$
$`2\beta \left[{\displaystyle \underset{a=1}{\overset{N}{}}}\varphi _s^{\left(a\right)}(x,y)\right]\varphi (x,y)`$
$`+_{}^2\varphi \alpha \varphi \beta \varphi ^2`$
$`=`$ $`0`$
The first line is zero and we have
$$_{}^2\varphi \left[\alpha +2\beta \underset{a=1}{\overset{N}{}}\varphi _s^{\left(a\right)}(x,y)\right]\varphi \beta \varphi ^2=0$$
(79)
We write a Lagrangean for the random field according to the MSR procedure
$$[\chi ,\varphi ]=\chi \left\{_{}^2\varphi \left[\alpha +2\beta \underset{a=1}{\overset{N}{}}\varphi _s^{\left(a\right)}(x,y)\right]\varphi \beta \varphi ^2\right\}$$
(80)
and the action functional is
$`S_{\phi V}[\chi ,\varphi ]`$ $`=`$ $`{\displaystyle ๐x๐y\chi \left\{_{}^2\varphi \left[\alpha +2\beta \underset{a=1}{\overset{N}{}}\varphi _s^{\left(a\right)}(x,y)\right]\varphi \beta \varphi ^2\right\}}`$
$`=`$ $`{\displaystyle ๐x๐y\left\{\left(\chi \right)\left(\varphi \right)\chi \left[\alpha +2\beta \underset{a=1}{\overset{N}{}}\varphi _s^{\left(a\right)}(x,y)\right]\varphi \beta \chi \varphi ^2\right\}}`$
The generating functional is defined from the functional integral
$$๐ฉ^1D\left[\chi \right]D\left[\varphi \right]\mathrm{exp}\left(S_{\phi V}[\chi ,\varphi ]\right)$$
(82)
Now we will modify the action by considering as usual the interaction with external currents,
$$\mathrm{\Xi }\left[J\right]๐ฉ^1D\left[\chi \right]D\left[\varphi \right]\mathrm{exp}\left(S_{\phi V}[\chi ,\varphi ]+J_\chi \chi +J_\varphi \varphi \right)$$
(83)
With this functional integral we will have to calculate the *free energy* functional. The functional $`\mathrm{\Xi }\left[J\right]`$ depends on the function representing the vortices. The vortices are assumed known but their position in plane is random therefore we have to average over them: $`W\left[J\right]=\mathrm{ln}\left(\mathrm{\Xi }\left[J\right]\right)`$.
### 5.1 The average over the positions
To perform the statistical average over the random positions of the vortices. The functional which we have to average is
$`\mathrm{\Xi }`$ $`=`$ $`๐ฉ^1{\displaystyle D\left[\chi \right]D\left[\varphi \right]}`$
$`\times \mathrm{exp}\{{\displaystyle }dxdy[(\chi )(\varphi )+\alpha \chi \varphi +\beta \chi \varphi ^2`$
$`+\left(2\beta {\displaystyle \underset{a=1}{\overset{N}{}}}\varphi _s^{\left(a\right)}(x,y)\right)\chi \varphi ]\}`$
The part that depends on the positions can be written
$`\mathrm{exp}\left\{{\displaystyle ๐x๐y\left(2\beta \underset{a=1}{\overset{N}{}}\varphi _s^{\left(a\right)}(x,y)\right)\chi \varphi }\right\}`$
$`=`$ $`\mathrm{exp}\left\{{\displaystyle ๐x๐y2\beta \varphi ^s\underset{a=1}{\overset{N}{}}\delta \left(๐ซ๐ซ_a\right)\chi \varphi }\right\}`$
$`=`$ $`\mathrm{exp}\left[2\beta \varphi ^s{\displaystyle \underset{a=1}{\overset{N}{}}}\chi \left(๐ซ_a\right)\varphi \left(๐ซ_a\right)\right]`$
where $`\varphi ^s`$ is now a simple amplitude. Consider the more general situation where we have to average in addition over the amplitudes $`\varphi ^s`$ of the vortices. If $`\varphi ^s`$ is a stochastic variable we have to perform an average of the type
$`\mathrm{exp}\left[2\beta \varphi ^s{\displaystyle \underset{a=1}{\overset{N}{}}}\chi \left(๐ซ_a\right)\varphi \left(๐ซ_a\right)\right]_{๐ซ_a,\varphi ^s}`$
$`=`$ $`{\displaystyle \underset{a=1}{\overset{N}{}}}\mathrm{exp}\left[2\beta \varphi ^s\chi \left(๐ซ_a\right)\varphi \left(๐ซ_a\right)\right]_{๐ซ_a,\varphi ^s}`$
Consider that the (now) random variable $`\varphi ^s`$ has the probability density
$$g\left(\varphi ^s\right)$$
(87)
Then, restricting for the moment to only the average over $`\varphi ^s`$,
$`\mathrm{exp}\left[2\beta \varphi ^s\chi \left(๐ซ_a\right)\varphi \left(๐ซ_a\right)\right]_{\varphi ^s}`$
$`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}๐\varphi ^sg\left(\varphi ^s\right)\mathrm{exp}\left[2\beta \varphi ^s\chi \left(๐ซ_a\right)\varphi \left(๐ซ_a\right)\right]`$
$`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}๐\varphi ^sg\left(\varphi ^s\right)\mathrm{exp}\left(i\lambda \varphi ^s\right)`$
where
$$i\lambda \left(๐ซ_a\right)2\beta \chi \left(๐ซ_a\right)\varphi \left(๐ซ_a\right)$$
(89)
Then
$$\mathrm{exp}\left[2\beta \varphi ^s\chi \left(๐ซ_a\right)\varphi \left(๐ซ_a\right)\right]_{\varphi ^s}=\stackrel{~}{g}\left[\lambda \left(๐ซ_a\right)\right]$$
(90)
where $`\stackrel{~}{g}`$ is the Fourier transform of the probability distribution function $`g`$. We use the notation
$$\stackrel{~}{g}\left(๐ซ_k\right)\stackrel{~}{g}\left[\lambda \left(๐ซ_a\right)\right]$$
(91)
and we have to calculate
$$\underset{a=1}{\overset{N}{}}\stackrel{~}{g}\left(i2\beta \chi \left(๐ซ_a\right)\varphi \left(๐ซ_a\right)\right)_{๐ซ_a}=\underset{a=1}{\overset{N}{}}\stackrel{~}{g}\left(๐ซ_a\right)$$
(92)
For this we introduce the function $`h`$
$$\stackrel{~}{g}\left(๐ซ_a\right)h\left(๐ซ_a\right)+1$$
(93)
and we rewrite the average as
$`{\displaystyle \underset{a=1}{\overset{N}{}}}\stackrel{~}{g}\left(๐ซ_a\right)`$
$`=`$ $`{\displaystyle \underset{a=1}{\overset{N}{}}}h\left(๐ซ_a\right)+1`$
$`=`$ $`{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{i_1<i_2<\mathrm{}<i_l}{}}h\left(๐ซ_{i_1}\right)h\left(๐ซ_{i_2}\right)\mathrm{}h\left(๐ซ_{i_l}\right)`$
$`=`$ $`\mathrm{exp}\left[{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{l!}}{\displaystyle ๐๐ซ_{i_1}๐๐ซ_{i_2}\mathrm{}๐๐ซ_{i_l}h\left(๐ซ_{i_1}\right)h\left(๐ซ_{i_2}\right)\mathrm{}h\left(๐ซ_{i_l}\right)C^{\left(l\right)}(๐ซ_{i_1},๐ซ_{i_2},\mathrm{},๐ซ_{i_l})}\right]`$
where we have introduced the cumulants of the distribution of points $`๐ซ_{i_1}`$ in the plane.
According to our assumption
$$C^{\left(l\right)}=\{\begin{array}{cc}1/A& l=1\\ 0& l>1\end{array}$$
(95)
where $`A`$ is the area in plane. Therefore we have the result of averaging
$`{\displaystyle \underset{a=1}{\overset{N}{}}}\stackrel{~}{g}\left(๐ซ_a\right)`$ $`=`$ $`\mathrm{exp}\left[{\displaystyle \frac{1}{A}}{\displaystyle ๐๐ซh\left(๐ซ\right)}\right]`$
$`=`$ $`\mathrm{exp}\left\{{\displaystyle \frac{1}{A}}{\displaystyle ๐๐ซ\left[\stackrel{~}{g}\left(๐ซ\right)1\right]}\right\}`$
The part of the generating functional which depended on the positions can now be written
$$\mathrm{exp}\left[2\beta \varphi ^s\underset{a=1}{\overset{N}{}}\chi \left(๐ซ_a\right)\varphi \left(๐ซ_a\right)\right]_{๐ซ_a,\varphi ^s}=\mathrm{exp}\left\{\frac{1}{A}๐๐ซ\left[\stackrel{~}{g}\left(๐ซ\right)1\right]\right\}$$
(97)
and the full partition function is
$`\mathrm{\Xi }_{๐ซ_a,\varphi ^s}`$ $`=`$ $`N^1{\displaystyle D\left[\chi \right]D\left[\varphi \right]}`$
$`\times \mathrm{exp}\{{\displaystyle }dxdy[(\chi )(\varphi )+\alpha \chi \varphi +\beta \chi \varphi ^2`$
$`+{\displaystyle \frac{1}{A}}{\displaystyle }d๐ซ[\stackrel{~}{g}\left(๐ซ\right)1]]\}`$
The action at the exponent is
$$๐x๐y\left\{\left(\chi \right)\left(\varphi \right)+\alpha \chi \varphi +\beta \chi \varphi ^2+\frac{1}{A}\left[\stackrel{~}{g}(x,y)1\right]\right\}$$
(99)
The function $`\stackrel{~}{g}\left(๐ซ\right)`$ has as argument the expression $`i2\beta \chi \left(๐ซ\right)\varphi \left(๐ซ\right)`$.
A reasonable assumption is that the vortices can only be positive or negative and with the same magnitude. Then we have
$$g\left(\varphi ^s\right)=\frac{1}{2}\left[\delta \left(\varphi ^s\varphi _0\right)+\delta \left(\varphi ^s+\varphi _0\right)\right]$$
(100)
and the Fourier transform is
$`\stackrel{~}{g}\left(q\right)`$ $`=`$ $`{\displaystyle ๐\varphi ^s\mathrm{exp}\left(iq\varphi ^s\right)\frac{1}{2}\left[\delta \left(\varphi ^s\varphi _0\right)+\delta \left(\varphi ^s+\varphi _0\right)\right]}`$
$`=`$ $`{\displaystyle \frac{1}{2}}\left[\mathrm{exp}\left(iq\varphi _0\right)+\mathrm{exp}\left(iq\varphi _0\right)\right]`$
$`=`$ $`\mathrm{cos}\left(q\varphi _0\right)`$
This must be calculated for the argument $`i2\beta \chi \left(๐ซ\right)\varphi \left(๐ซ\right)`$ and gives
$`\stackrel{~}{g}\left(๐ซ\right)`$ $`=`$ $`\mathrm{cos}\left[i2\beta \varphi _0\chi \left(๐ซ\right)\varphi \left(๐ซ\right)\right]`$
$`=`$ $`\mathrm{cosh}\left[2\beta \varphi _0\chi \left(๐ซ\right)\varphi \left(๐ซ\right)\right]`$
The functional integral that must be calculated becomes
$`\mathrm{\Xi }_{๐ซ_k,\varphi ^s}`$ $`=`$ $`๐ฉ^1{\displaystyle D\left[\chi \right]D\left[\varphi \right]}`$
$`\times \mathrm{exp}\left\{{\displaystyle ๐x๐y\left[\left(\chi \right)\left(\varphi \right)+\alpha \chi \varphi +\beta \chi \varphi ^2+\frac{1}{A}\mathrm{cosh}\left(2\beta \varphi _0\chi \varphi \right)\right]}\right\}`$
The system is perturbed with an external current $`J(x,y)`$ acting on the field $`\varphi (x,y)`$.
$`Z_J`$ $``$ $`\mathrm{\Xi }_{๐ซ_a,\varphi ^s}\left[J\right]`$
$`=`$ $`๐ฉ^1{\displaystyle D\left[\chi \right]D\left[\varphi \right]\mathrm{exp}\left(S_J\right)}`$
$$S_J๐x๐y\left[\left(\chi \right)\left(\varphi \right)+\alpha \chi \varphi +\beta \chi \varphi ^2+\frac{1}{A}\mathrm{cosh}\left(2\beta \varphi _0\chi \varphi \right)+J\varphi \right]$$
(105)
## 6 Approximation for small amplitude vortices
Consider that the amplitudes of the vortices are not high, $`\varphi _0`$. Then
$$\mathrm{cosh}\left(2\beta \varphi _0\chi \varphi \right)1+\frac{1}{2}\left(2\beta \varphi _0\chi \varphi \right)^2$$
(106)
and the action (removing some terms without significance)
$$S_J๐x๐y\left[\left(\chi \right)\left(\varphi \right)+\alpha \chi \varphi +\beta \chi \varphi ^2+\frac{\left(2\beta \varphi _0\right)^2}{A}\chi ^2\varphi ^2+J\varphi \right]$$
(107)
This action in principle can lead to a perturbative treatment but with two vertices, of order three and of order four, a very difficult and unusual problem. We remark however that
$$S_J๐x๐y\left[\chi \left(^2\varphi \right)+\alpha \chi \varphi +\beta \chi \varphi ^2+\frac{\left(2\beta \varphi _0\right)^2}{A}\chi ^2\varphi ^2+J\varphi \right]$$
(108)
may become quadratic in $`\chi `$ and in $`\varphi `$ (therefore the functional integral is Gaussian) if we succeed to separate the product $`\chi ^2\varphi ^2`$.
#### 6.0.1 Technical step
Consider the following formula to disentangle the two variables
$$\mathrm{exp}\left(\frac{1}{2}U^2\right)=\frac{1}{\sqrt{2\pi }}_{\mathrm{}}^{\mathrm{}}๐s\mathrm{exp}\left(\frac{1}{2}s^2Us\right)$$
(109)
then
$$\mathrm{exp}\left(\frac{1}{2}2Q\chi ^2\varphi ^2\right)=\frac{1}{\sqrt{2\pi }}_{\mathrm{}}^{\mathrm{}}๐s\mathrm{exp}\left(\frac{1}{2}s^2\sqrt{2Q}\chi \varphi s\right)$$
(110)
and we have the action
$`Z_J`$ $`=`$ $`๐ฉ^1{\displaystyle \frac{1}{\sqrt{2\pi }}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}๐s\mathrm{exp}\left({\displaystyle \frac{1}{2}}s^2\right)`$
$`\times {\displaystyle }D\left[\varphi \right]D\left[\chi \right]\mathrm{exp}\left\{{\displaystyle }dxdy[\chi \left(^2\varphi \right)+\alpha \chi \varphi +\beta \chi \varphi ^2\sqrt{2Q}\chi \varphi s+J\varphi ]\right\}`$
or
$$S_J=๐x๐y\left[J\varphi \right]+๐x๐y\left[\chi \left(^2\varphi \right)+\chi \left(\alpha s\sqrt{2Q}\right)\varphi +\beta \chi \varphi ^2\right]$$
(112)
We have used the notation
$$Q\frac{\left(2\beta \varphi _0\right)^2}{A}$$
(113)
The functional integration over $`\chi `$ can be done immediately and gives a functional $`\delta `$ with the argument the equation in the modified form. If we make the integration over $`\chi `$ we obtain
$`{\displaystyle D\left[\chi \right]\mathrm{exp}\left\{๐x๐y\left[\chi \left(^2\varphi \right)+\chi \left(\alpha s\sqrt{2Q}\right)\varphi +\beta \chi \varphi ^2\right]\right\}}`$
$`=`$ $`{\displaystyle D\left[\chi \right]\mathrm{exp}\left[๐x๐y\chi \overline{F}\left(\varphi \right)\right]}`$
$`=`$ $`\left[\left|i{\displaystyle \frac{\delta \overline{F}}{\delta \varphi }}\right|_{\varphi ^z}\right]^1\delta \left(\varphi \varphi ^z\right)`$
where we have introduced the notation
$$\overline{F}\left(\varphi \right)^2\varphi +\left(\alpha s\sqrt{2Q}\right)\varphi +\beta \varphi ^2$$
(115)
and the function $`\overline{\varphi }^z`$ is the solution of the differential equation $`\overline{F}\left(\varphi \right)=0`$ *i.e.*
$$\overline{F}\left(\overline{\varphi }^z\right)=0$$
(116)
Inserting this result in the integral over the functions $`\varphi `$, we have
$`๐ฉ^1{\displaystyle D\left[\varphi \right]\mathrm{exp}\left(๐x๐yJ\varphi \right)\left[\left|i\frac{\delta \overline{F}}{\delta \varphi }\right|_{\overline{\varphi }^z}\right]^1\delta \left(\varphi \overline{\varphi }^z\right)}`$
$`=`$ $`\left[\left|i{\displaystyle \frac{\delta \overline{F}}{\delta \varphi }}\right|_{\overline{\varphi }^z}\right]^1\mathrm{exp}\left({\displaystyle ๐x๐yJ\overline{\varphi }^z}\right)`$
Then
$$Z_J=๐ฉ^1\frac{1}{\sqrt{2\pi }}_{\mathrm{}}^{\mathrm{}}๐s\mathrm{exp}\left(\frac{1}{2}s^2\right)\left[\left|i\frac{\delta \overline{F}}{\delta \varphi }\right|_{\overline{\varphi }^z}\right]^1\mathrm{exp}\left(๐x๐yJ\overline{\varphi }^z\right)$$
(118)
The most important part is to calculate
$$i\frac{\delta \overline{F}}{\delta \varphi }|_{\overline{\varphi }^z}=i\frac{\delta }{\delta \varphi }\left[^2\varphi +\left(\alpha s\sqrt{2Q}\right)\varphi +\beta \varphi ^2\right]|_{\overline{\varphi }^z}$$
(119)
But this simplifies with the part of the factor
$$๐ฉ^1\left|\frac{\delta F}{\delta \varphi }\right|_{\varphi ^z}$$
(120)
The fact that the initial factor is written before averaging over the positions and the second factor is obtained after averaging will make a small difference which can only be of higher order. What remains is a factor of normalization that we call $`c^1`$.
The functional derivatives are
$$\frac{\delta Z_j}{\delta J(x,y)}=\frac{1}{c}\frac{1}{\sqrt{2\pi }}_{\mathrm{}}^{\mathrm{}}๐s\mathrm{exp}\left(\frac{1}{2}s^2\right)\overline{\varphi }^z(x,y)\mathrm{exp}\left(๐x๐yJ\overline{\varphi }^z\right)$$
(121)
and
$`{\displaystyle \frac{\delta ^2Z_J}{\delta J(x,y)\delta J(x^{},y^{})}}`$ $`=`$ $`{\displaystyle \frac{1}{c}}{\displaystyle \frac{1}{\sqrt{2\pi }}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}๐s\mathrm{exp}\left({\displaystyle \frac{1}{2}}s^2\right)`$
$`\times \overline{\varphi }^z(x,y)\overline{\varphi }^z(x^{},y^{})\mathrm{exp}\left({\displaystyle ๐x๐yJ\overline{\varphi }^z}\right)`$
and taking
$$J0$$
(123)
we can calculate the contribution to the two-point correlation arising from this part of the partition function.
We have to solve
$$^2\varphi +\left(\alpha s\sqrt{2Q}\right)\varphi +\beta \varphi ^2=0$$
(124)
for a small amplitude field $`\varphi `$. This can be done by approximations, starting from the solution of the equation without random centers.
We can see that the immediate effect of the scattering of the turbulent field by the vortices is a renormalization of the coefficient $`\alpha `$ of the equation.
The solution that we need for Eq.(124) must reflect the fact that we are examining a turbulent field, in the dynamics of which the randomly placed vortices have been included. We can see that the integration over the auxiliary variable $`s`$ must not necessarily be truncated to ensure that $`\alpha _s\alpha s\sqrt{2Q}`$ remains positive. This is because the field is turbulent and there is no restriction concerning the existence of the vortex solution which should be imposed to the turbulent component of the field. We also note that the approximation taking only the $`s=0`$ contribution in the integral (6.0.1) returns us to the problem of the turbulent field without interaction with the random vortices. Although the main contribution (as shown by the Gaussian integration) comes from the free turbulent field itself, the presence of vortices is contained in the rest of the integral, for $`\left|s\right|`$ not close to $`0`$.
These remarks suggest to study the statistical properties of the transformed field, where $`\alpha `$ is replaced with $`\alpha _s`$.
## 7 The background turbulence: perturbative treatment
The vortices represent the strongly nonlinear part of the system and they have been extracted and treated separately as shown in Section IV. However, there is still nonlinear interaction in the remaining random field. This is the nonlinear mode coupling which produces the stationary turbulent states of the system, even in absence of any definite structure formation. This interaction must be taken into account when we analyse the statistical properties of the turbulent field. From general consideration we know what can be expected from this analysis: we will obtain a small departure from pure Gaussian statistics, expressed in nonlinear renormalization of the propagator and of the vertex of the interaction. The presence of random interaction with elements of the system that are beyond our simple model is accounted, as usual, by a noise term acting like a drive for the system of random waves. The noise will be assumed with the simplest (white) statistics
$$\zeta (x,y)\zeta (x^{},y^{})=D\delta \left(xx^{}\right)\delta \left(yy^{}\right)$$
(125)
and should be considered as a random stirring force composed, for example, of random growths and decays of marginally stable modes, thus injecting at random places some energy into the system.
We will provide a standard perturbative treatment for the differential equation
$$^2\varphi +\left(\alpha s\sqrt{2A}\right)\varphi +\beta \varphi ^2=\zeta $$
(126)
with the objective to calculate correlation functions and other statistical properties.
As usual we start by defining the generating functional of the statistical correlation
$`\mathrm{exp}\left(W\left[๐\right]\right)`$
$`=`$ $`{\displaystyle ๐\left[\eta \right]๐\left[\varphi \right]\mathrm{exp}\left\{๐x๐y\left[\eta ^2\varphi +\eta \alpha _s\varphi +\eta \beta \varphi ^2+J_\eta \eta +J_\varphi \varphi \right]\right\}}`$
$`\times {\displaystyle ๐\left[\eta \right]๐\left[\varphi \right]\mathrm{exp}\left\{๐x๐y\eta \zeta \right\}}`$
where we have introduced the notation
$$\alpha _s\alpha s\sqrt{2Q}$$
(128)
The second factor can easily be calculated
$`{\displaystyle ๐\left[\eta \right]๐\left[\varphi \right]\mathrm{exp}\left(๐x๐y\eta \zeta \right)}`$
$`=`$ $`{\displaystyle ๐\left[\eta \right]๐\left[\varphi \right]\mathrm{exp}\left[๐x๐y\left(2D\eta ^2\right)\right]}`$
In the generating functional there are the two dual functions, $`\eta (x,y)`$ and $`\varphi (x,y)`$, the second being the physical field of the random waves. We have inserted the action related with interaction with โexternalโ currents, $`J_\eta `$ and $`J_\varphi `$ which will permit to express the averages as functional derivatives.
Defining the Lagrangean density (after an integration by parts)
$$=\left(\eta \right)\left(\varphi \right)+\eta \alpha _s\varphi +\eta \beta \varphi ^2+D\eta ^2+J_\eta \eta +J_\varphi \varphi $$
(130)
we will separate into the Gaussian the nonlinear interaction part
$$=_0+_I$$
(131)
$$_0=\left(\eta \right)\left(\varphi \right)+\eta \alpha _s\varphi +D\eta ^2+J_\eta \eta +J_\varphi \varphi $$
(132)
$$_I=\eta \beta \varphi ^2$$
(133)
The first part of the Lagrangean is not linear because of the square term, but it is Gaussian therefore it would not pose any problem to the functional integration. We note that we have an order three vertex.
The action functional is correspondingly divided into two parts and can be written, as usual
$`Z\left[๐\right]`$ $`=`$ $`{\displaystyle ๐\left[\eta \right]๐\left[\varphi \right]\mathrm{exp}\left(๐x๐y_0\right)\mathrm{exp}\left(๐x๐y_I\right)}`$
$`=`$ $`\mathrm{exp}\left(\beta {\displaystyle ๐x๐y\frac{\delta }{\delta J_\eta }\frac{\delta }{\delta J_\varphi }\frac{\delta }{\delta J_\varphi }}\right){\displaystyle ๐\left[\eta \right]๐\left[\varphi \right]\mathrm{exp}\left(๐x๐y_0\right)}`$
For the linear part we have the Euler-Lagrange equations
$$\frac{d}{dx}\frac{\delta _0}{\delta \left(\eta /x\right)}+\frac{d}{dy}\frac{\delta _0}{\delta \left(\eta /y\right)}\frac{\delta _0}{\delta \eta }=0$$
(135)
$$\frac{d}{dx}\frac{\delta _0}{\delta \left(\varphi /x\right)}+\frac{d}{dy}\frac{\delta _0}{\delta \left(\varphi /y\right)}\frac{\delta _0}{\delta \varphi }=0$$
(136)
The equations are
$$\mathrm{\Delta }\varphi \alpha _s\varphi 2D\eta =J_\eta $$
(137)
and
$$\mathrm{\Delta }\eta \alpha _s\eta =J_\varphi $$
(138)
The latter is an inhomogeneous Helmholtz equation and has the solution
$$\eta _0(x,y)=๐x^{}๐y^{}G_{\eta \varphi }(x,y;x^{},y^{})J_\varphi (x^{},y^{})$$
(139)
in terms of the Green function appropriate for the space domain of our analysis and taking into account the boundary conditions for $`\eta `$.
The first Euler-Lagrange equation gives
$`\mathrm{\Delta }\varphi \alpha _s\varphi `$ $`=`$ $`J_\eta +2D\eta _0`$
$`=`$ $`J_\eta +2D{\displaystyle ๐x^{}๐y^{}G_{\eta \varphi }(x,y;x^{},y^{})J_\varphi (x^{},y^{})}`$
with the solution
$`\varphi _0(x,y)`$ $`=`$ $`{\displaystyle ๐x^{}๐y^{}G_{\varphi \eta }(x,y;x^{},y^{})J_\eta (x^{},y^{})}`$
$`+2D{\displaystyle ๐x^{}๐y^{}๐x^{\prime \prime }๐y^{\prime \prime }G_{\varphi \eta }(x,y;x^{\prime \prime },y^{\prime \prime })G_{\eta \varphi }(x^{\prime \prime },y^{\prime \prime };x^{},y^{})J_\varphi (x^{},y^{})}`$
We can introduce the matrix of propagators
$$๐\left(G_{ij}\right)=\left(\begin{array}{cc}G_{\varphi \varphi }& G_{\varphi \eta }\\ G_{\eta \varphi }& 0\end{array}\right)$$
(142)
where
$$G_{\varphi \varphi }(x,y;x^{},y^{})=2D๐x^{\prime \prime }๐y^{\prime \prime }G_{\varphi \eta }(x,y;x^{\prime \prime },y^{\prime \prime })G_{\eta \varphi }(x^{\prime \prime },y^{\prime \prime };x^{},y^{})$$
(143)
Introducing the notation
$$\mathrm{\Psi }\left(\begin{array}{c}\varphi \\ \eta \end{array}\right)$$
(144)
we will write symbolically (summation over repeated indices is implied)
$$\mathrm{\Psi }_i=๐x^{}๐y^{}G_{ik}J_k$$
(145)
for
$$\mathrm{\Psi }_i=๐x^{}๐y^{}G(x,y;x^{},y^{})J_k(x^{},y^{})$$
(146)
We can now calculate the linear part of the action along the system โs configurations given by these solutions. The action will be, of course, a functional of the current.
$`S_{0J}`$ $`=`$ $`{\displaystyle ๐x๐y\left[\left(\eta _0\right)\left(\varphi _0\right)+\alpha _s\eta _0\varphi _0+D\eta _0^2+J_\eta \eta _0+J_\varphi \varphi _0\right]}`$
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle ๐x๐y\left[\eta _0\mathrm{\Delta }\varphi _0+\alpha _s\eta _0\varphi _0+2D\eta _0^2+J_\eta \eta _0+J_\varphi \varphi _0\right]}`$
$`+{\displaystyle \frac{1}{2}}{\displaystyle ๐x๐y\left[\varphi _0\mathrm{\Delta }\eta _0+\alpha _s\eta _0\varphi _0+J_\eta \eta _0+J_\varphi \varphi _0\right]}`$
or
$`S_{0J}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle ๐x๐y\left[\eta _0\left(\mathrm{\Delta }\varphi _0+\alpha _s\varphi _0+2D\eta _0+J_\eta \right)+J_\varphi \varphi _0\right]}`$
$`+{\displaystyle \frac{1}{2}}{\displaystyle ๐x๐y\left[\varphi _0\left(\mathrm{\Delta }\eta _0+\alpha _s\eta _0+J_\varphi \right)+J_\eta \eta _0\right]}`$
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle ๐x๐y\left(J_\varphi \varphi _0+J_\eta \eta _0\right)}`$
We now express the solution by the Green functions
$`S_{0J}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle }dxdy[J_\varphi {\displaystyle }dx^{}dy^{}G_{\varphi \eta }(x,y;x^{},y^{})J_\eta (x^{},y^{})`$
$`+J_\eta {\displaystyle }dx^{}dy^{}G_{\eta \varphi }(x,y;x^{},y^{})J_\varphi (x^{},y^{})]`$
which can be written using the convention introduced above
$$S_{0J}=\frac{1}{2}๐x๐y๐x^{}๐y^{}J_i(x,y)G_{ij}(x,y;x^{},y^{})J_k(x^{},y^{})$$
(150)
We return to the generating functional
$`Z\left[๐\right]`$ $`=`$ $`\mathrm{exp}\left(\beta {\displaystyle ๐x๐y\frac{\delta }{\delta J_\eta }\frac{\delta }{\delta J_\varphi }\frac{\delta }{\delta J_\varphi }}\right)\mathrm{exp}\left(S_{0J}\right)`$
$`=`$ $`\mathrm{exp}\left({\displaystyle ๐x๐yC_{ijk}\frac{\delta }{\delta J_i}\frac{\delta }{\delta J_j}\frac{\delta }{\delta J_k}}\right)`$
$`\times \mathrm{exp}\left[{\displaystyle \frac{1}{2}}{\displaystyle ๐x๐y๐x^{}๐y^{}J_i(x,y)G_{ij}(x,y;x^{},y^{})J_k(x^{},y^{})}\right]`$
where we have introduced, for uniformity of notation
$$C_{ijk}(x,y)\beta \delta _{i\eta }\delta _{j\varphi }\delta _{k\varphi }$$
(152)
This is a classical framework for diagrammatic expansion.
In general the functional approach is well suited for statistical problems where the statistical ensemble of the realization of the systemโs configurations is produced by the effect of a noise or by a random choice of initial condition, as explained in MSR and Jensen. In the most usual case, where there is a noise acting on the system, using the (known) statistical properties of the noise means implicitly that we make a change of variables, from the field $`\varphi `$ to the noise $`\zeta `$ and this introduces a Jacobian,
$$๐ฅ\left[\varphi \right]=๐\left[\zeta \right]/๐\left[\varphi \right]$$
(153)
The Jacobian cancels all the diagrams that correspond to the *vacuum* to *vacuum* transitions, or equal-coordinates correlations. For this reason it is not mentioned. This is explained in Appendix B.
When applying the operator to the second exponential we have to remind that we need at least two open ends with currents $`J_\varphi `$ since we intend to determine the spectrum. On the other hand, it is clear that any open end with isolated either $`J_\varphi `$ or $`J_\eta `$ would vanish after taking the currents to zero.
The first term is naturally the diffusion driven by the random rise and decay of marginally stable modes, represented here as a noise.
The next term that is useful consists of the product of two vertex-like operators applied on a product of four propagator-like terms
$$J_{i_1}G_{i_1i_2}C_{i_2i_3i_4}G_{i_4i_5}G_{i_3i_6}C_{i_5i_6i_7}G_{i_7i_1}J_{i_1}$$
(154)
Obviously, this is the *one-loop* diagram. We will derivate to the first and last factors (currents) so that, more explicitly, its structure is
$$G_{\varphi \eta }C_{\eta \varphi \varphi }G_{\varphi \varphi }G_{\varphi \varphi }C_{\varphi \varphi \eta }G_{\eta \varphi }$$
(155)
We can represent the Green function by its Fourier transform
$$\left(\mathrm{\Delta }\alpha _s\right)G(x,y;x^{},y^{})=\delta \left(xx^{}\right)\delta \left(yy^{}\right)$$
(156)
$`G(x,y;x^{},y^{})`$ $`=`$ $`{\displaystyle ๐k_x๐k_y\frac{1}{k_x^2+k_y^2+\alpha _s}}`$
$`\times \mathrm{exp}\left[ik_x\left(xx^{}\right)\right]\mathrm{exp}\left[ik_y\left(yy^{}\right)\right]`$
We obtain the propagators for the coupling between the field $`\varphi `$ and its dual $`\eta `$
$`G_{\varphi \eta }(x,y;x^{},y^{})`$ $`=`$ $`{\displaystyle ๐k_x๐k_y\frac{1}{k_x^2+k_y^2+\alpha _s}}`$
$`\times \mathrm{exp}\left[ik_x\left(xx^{}\right)\right]\mathrm{exp}\left[ik_y\left(yy^{}\right)\right]`$
$`=`$ $`G_{\eta \varphi }(x,y;x^{},y^{})`$
$`G_{\varphi \varphi }(x,y)`$ $`=`$ $`2D{\displaystyle ๐x^{\prime \prime }๐y^{\prime \prime }G_{\varphi \eta }(x,y;x^{\prime \prime },y^{\prime \prime })G_{\eta \varphi }(x^{\prime \prime },y^{\prime \prime };x^{},y^{})}`$
$`=`$ $`{\displaystyle ๐k_x๐k_y\frac{1}{\left(k_x^2+k_y^2+\alpha _s\right)^2}}`$
$`\times \mathrm{exp}\left[ik_x\left(xx^{}\right)\right]\mathrm{exp}\left[ik_y\left(yy^{}\right)\right]`$
The intermediate integration over space $`(x^{\prime \prime },y^{\prime \prime })`$ yields equality of the Fourier variables of the two Green functions. We then have
$`\stackrel{~}{G}_{\varphi \varphi }`$ $`=`$ $`2D\left(k_x^2+k_y^2+\alpha _s\right)^2`$ (160)
$`\stackrel{~}{G}_{\varphi \eta }`$ $`=`$ $`\stackrel{~}{G}_{\eta \varphi }=\left(k_x^2+k_y^2+\alpha _s\right)^1`$
The one-loop term is expressed as two integrals over the intermediate momenta at the two vertices where we have a product of three Green functions with the vertex. Conservation of momentum in the loop and overall conservation of the diagram (which means that the $`๐ค`$ at input line must be the same at the output line) lead to
$`\varphi (x,y)\varphi (x^{},y^{})`$
$`=`$ $`{\displaystyle ๐k_x๐k_y\varphi \varphi _๐ค\mathrm{exp}\left[ik_x\left(xx^{}\right)ik_y\left(yy^{}\right)\right]}`$
$`=`$ $`{\displaystyle \frac{2D}{\left(๐ค^2+\alpha _s\right)^2}}`$
$`+{\displaystyle ๐๐ค\mathrm{\Gamma }๐๐ฉ\frac{1}{๐ค^2+\alpha _s}\beta \frac{2D}{\left(๐ฉ^2+\alpha _s\right)^2}\frac{2D}{\left[\left(๐ค๐ฉ\right)^2+\alpha _s\right]^2}\beta \frac{1}{๐ค^2+\alpha _s}}`$
The coefficient $`\mathrm{\Gamma }`$ represents the multiplicity of this diagram and factors of normalization.
$`\varphi \varphi _๐ค`$
$`=`$ $`\mathrm{\Gamma }\left(2D\beta \right)^2\left({\displaystyle \frac{1}{๐ค^2+\alpha _s}}\right)^2{\displaystyle ๐๐ฉ\frac{1}{\left(๐ฉ^2+\alpha _s\right)^2}\frac{1}{\left[\left(๐ค๐ฉ\right)^2+\alpha _s\right]^2}}`$
$`=`$ $`\mathrm{\Gamma }\left(2D\beta \right)^2\left({\displaystyle \frac{1}{๐ค^2+\alpha _s}}\right)^2{\displaystyle ๐๐ฉ\frac{1}{\left[\left(๐ค/2๐ฉ\right)^2+\alpha _s\right]^2}\frac{1}{\left[\left(๐ค/\mathrm{๐}+๐ฉ\right)^2+\alpha _s\right]^2}}`$
We transform the integral
$`I`$ $``$ $`{\displaystyle p๐p๐\theta \frac{1}{\left[๐ฉ^2๐ค๐ฉ+๐ค^2/4+\alpha _s\right]^2}}`$
$`\times {\displaystyle \frac{1}{\left[๐ฉ^2+๐ค๐ฉ+๐ค^2/4+\alpha _s\right]^2}}`$
Now we introduce the variables that will replace the constant terms
$$b๐ค^2/4+\alpha _s$$
(164)
and
$`๐ค๐ฉ`$ $`=`$ $`kp\mathrm{cos}\theta `$
$``$ $`up`$
The integral becomes
$$I=p๐p๐\theta \frac{1}{\left[p^2up+b\right]^2}\frac{1}{\left[p^2+up+b\right]^2}$$
(166)
We write the denominator in the form
$`\left[p^2up+b\right]^2\left[p^2+up+b\right]^2`$
$`=`$ $`\left[\left(p^2+b\right)^2u^2p^2\right]^2`$
$`=`$ $`\left[p^4+2p^2b+b^2u^2p^2\right]^2`$
$`=`$ $`\left[p^4+\left(2bu^2\right)p^2+b^2\right]^2`$
$`=`$ $`\left(p^4+2\gamma p^2+\delta \right)^2`$
where
$`2\gamma `$ $``$ $`2bu^2`$ (168)
$`\delta `$ $``$ $`b^2`$
The integral becomes
$$I=๐\theta _0^{\mathrm{}}\frac{pdp}{\left(p^4+2\gamma p^2+\delta \right)^2}$$
(169)
we make the substitution
$$p^2x$$
(170)
$$I=๐\theta \frac{1}{2}_0^{\mathrm{}}\frac{dx}{\left(x^2+2\gamma x+\delta \right)^2}$$
(171)
The integration over the intermediate momentum is
$`{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dx}{\left(x^2+2\gamma x+\delta \right)^2}}`$
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{1}{\left(\delta \gamma ^2\right)^{3/2}}}\left[{\displaystyle \frac{\pi }{2}}\mathrm{arctan}{\displaystyle \frac{\gamma }{\sqrt{\delta \gamma ^2}}}\right]+\left({\displaystyle \frac{\gamma }{2\delta }}\right){\displaystyle \frac{1}{\delta \gamma ^2}}`$
Here
$`\gamma `$ $`=`$ $`bu^2/2`$
$`=`$ $`{\displaystyle \frac{1}{4}}\left[k^2\left(12\mathrm{cos}^2\theta \right)+4\alpha _s\right]`$
$`{\displaystyle \frac{\gamma }{\sqrt{\delta \gamma ^2}}}`$ $`=`$ $`{\displaystyle \frac{bu^2/2}{\left|u\right|\sqrt{bu^2/4}}}`$
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{k^2\left(12\mathrm{cos}^2\theta \right)+4\alpha _s}{k\left|\mathrm{cos}\theta \right|\sqrt{k^2\mathrm{sin}^2\theta +4\alpha _s}}}`$
$`\delta \gamma ^2`$ $`=`$ $`u^2\left(b{\displaystyle \frac{u^2}{4}}\right)`$
$`=`$ $`{\displaystyle \frac{1}{4}}k^2\mathrm{cos}^2\theta \left(k^2\mathrm{sin}^2\theta +\alpha _s\right)`$
$$\frac{\gamma }{\delta }=4\frac{k^2\left(12\mathrm{cos}^2\theta \right)+4\alpha _s}{\left(k^2+4\alpha _s\right)^2}$$
(176)
Then
$`{\displaystyle p๐p๐\theta \frac{1}{\left[๐ฉ^2๐ค๐ฉ+๐ค^2/4+\alpha _s\right]^2}\frac{1}{\left[๐ฉ^2+๐ค๐ฉ+๐ค^2/4+\alpha _s\right]^2}}`$
$`=`$ $`{\displaystyle }d\theta \{{\displaystyle \frac{8}{\left[k^2\mathrm{cos}^2\theta \left(k^2\mathrm{sin}^2\theta +\alpha _s\right)\right]^{3/2}}}`$
$`\times \left[{\displaystyle \frac{\pi }{2}}\mathrm{arctan}{\displaystyle \frac{k^2\left(12\mathrm{cos}^2\theta \right)+4\alpha _s}{2k\left|\mathrm{cos}\theta \right|\sqrt{k^2\mathrm{sin}^2\theta +4\alpha _s}}}\right]`$
$`{\displaystyle \frac{8}{k^2\mathrm{cos}^2\theta \left(k^2\mathrm{sin}^2\theta +\alpha _s\right)}}{\displaystyle \frac{k^2\left(12\mathrm{cos}^2\theta \right)+4\alpha _s}{\left(k^2+4\alpha _s\right)^2}}\}`$
We can show that the integral is not singular at the limit
$$\mathrm{cos}\theta 0$$
(178)
so that the integration over $`\theta `$ is safe.
As order of magnitude the integral is dominated by terms like
$$\mathrm{\Gamma }\left(2D\beta \right)^2\frac{1}{k^2\left(k^2+\alpha _s\right)^{3/2}}$$
(179)
while the diffusive part is
$$\frac{2D}{\left(k^2+\alpha _s\right)^2}$$
(180)
After this we obtain for the turbulent background field
$$\varphi \varphi _๐ค^{turbulence}=a\frac{2D}{\left(k^2+\alpha _s\right)^2}+b\left(2D\beta \right)^2\frac{1}{\left(k^2+\alpha _s\right)^2}\frac{1}{k^2\left(k^2+\alpha _s\right)^{3/2}}$$
(181)
where we have collected the factors in two numbers, $`a`$ and $`b`$.
## 8 Summary
It is usually assumed that a turbulent plasma (in prticular with embedded structures) should exhibit a spectrum of exponential or algebraic type. There is no universal theoretical basis for such an assumption except for cases where the scaling invariance is justified on physical ground. We find that it would be more adequate to extract a particular behavior (on different spectral intervals) from expressions like Eq.(4.2) and Eq.(181), via regression on simple numerical values. This can be done for a particular physical system. Here, however, we will return to the traditional exponents in order to exhibit some associations we consider to be general.
We can now collect the results of the analysis. In the left column we write the approximative behavior of the contributions and in the right column the physical origin, according to this theory.
$$\begin{array}{cc}k^2\hfill & \begin{array}{c}\text{gas of vortices}\hfill \\ \text{(from }N=1\text{ to closely packed)}\hfill \end{array}\hfill \\ k^4\hfill & \begin{array}{c}\text{background turbulence + vortices }\frac{2D}{\left(k^2+\alpha _s\right)^2}\hfill \\ \text{weak interaction of vortices }K_0\hfill \end{array}\hfill \\ S\left(๐ค\right)k^2\text{(weak }\rho _s^2/A\text{)}\hfill & \text{geometry of c.s., }\text{via}\text{ interaction}\hfill \\ S\left(๐ค\right)\left[1+f\left(๐ค\right)\right]\hfill & \text{perturbed c.s., }\text{via}\text{ rare events of large }N\hfill \\ b\left(2D\beta \right)^2\frac{1}{\left(k^2+\alpha _s\right)^2}\frac{1}{k^2\left(k^2+\alpha _s\right)^{3/2}}\hfill & \text{one-loop mode coupling}\hfill \end{array}$$
We see two ways in which these results can be useful.
First, the combination of various exponential dependences from the listed contributions will result in an overall dependence with exponents that may be compared with the experiments or numerical simulations , , , , . The weight of each contribution is dependent on factors like the shape of the vortices (including the spatial extension compared to the area $`A`$), amplitude of the vortices, strength of the random drive ($`D`$). Also in the problem enters as parameter the function of distribution on the plane (assumed here uniform). The energy of interaction between vortices may be reconsidered along the example of vortices in superfluids.
Second, the spectrum can be seen as dominated in different spectral domains by one or another of these contributions. This is compatible with the known fact that the spectrum has regions with different exponents $`\mu `$.
In both these ways we have to solve an inverse problem, starting from the experimental (or numerical) spectrum and using the above list to map the exponential form to a physical process which may have been at its origin.
Various extensions of this treatment are possible. One can introduce a chemical potential for the description of the statistical equilibrium consisting of generation and suppression of coherent vortices. There are treatments of this type for the conversion of the global rotation of superfluids into localised vortices, an example that may also be useful for the consideration of the zonal flow saturation in tokamak, besides the Kelvin-Helmholtz instability and the collisions.
The treatment by generating functional allows in principle determination of statistical correlations at any high order desired. However, while the functional derivatives can easily produce the $`n`$-th order cumulant, we must be sure that the generating functional has been calculated with the necessary precision, or, in other words, we must be sure that we have incorporated the physical origin of these correlations. This method should be accompagned by the more standard analysis of closure of the hierarchy of equations for the correlations and these two approaches must be seen as complementary.
As a final remark, the physical model adopted in the present treatment may be extended to cover more complex regimes, with, of course, a certain increase in the analytical work.
Aknowledgments. The authors gratefully aknowledge the discussions with Professor David Montgomery. This work has been partly supported by the Japan Society for the Promotion of Science. The authors are grateful for this support and for the hospitality of Professor S.-I. Itoh and Professor M. Yagi.
## 9 Appendix A : physics of the equation
It is of interest to study a realistic two-dimensional model, like Hasegawa-Wakatani or similar. However we need for the beginning a simpler equation and if possible with a vortex solution with known analytical expression.
A possibility is the equation for the ion drift instabilities but here we must specify the scales.
Consider the equations for the ITG model in two-dimensions with adiabatic electrons:
$`{\displaystyle \frac{n_i}{t}}+\left(๐ฏ_in_i\right)`$ $`=`$ $`0`$ (182)
$`{\displaystyle \frac{๐ฏ_i}{t}}+\left(๐ฏ_i\right)๐ฏ_i`$ $`=`$ $`{\displaystyle \frac{e}{m_i}}\left(\varphi \right)+{\displaystyle \frac{e}{m_i}}๐ฏ_i\times ๐`$
We assume the quasineutrality
$$n_in_e$$
(183)
and the Boltzmann distribution of the electrons along the magnetic field line
$$n_e=n_0\mathrm{exp}\left(\frac{\left|e\right|\varphi }{T_e}\right)$$
(184)
In general the electron temperature can be a function of the radial variable
$$T_eT_e\left(x\right)$$
(185)
The velocity of the ion fluid is perpendicular on the magnetic field and is composed of the diamagnetic, electric and polarization drift terms
$`๐ฏ_i`$ $`=`$ $`๐ฏ_i`$
$`=`$ $`๐ฏ_{dia,i}+๐ฏ_E+๐ฏ_{pol,i}`$
$`=`$ $`{\displaystyle \frac{T_i}{\left|e\right|B}}{\displaystyle \frac{1}{n_i}}{\displaystyle \frac{dn_i}{dr}}\widehat{๐}_y`$
$`+{\displaystyle \frac{\varphi \times \widehat{๐ง}}{B}}`$
$`{\displaystyle \frac{1}{B\mathrm{\Omega }_i}}\left({\displaystyle \frac{}{t}}+\left(๐ฏ_E_{}\right)\right)_{}\varphi `$
The diamagnetic velocity will be neglected. Introducing this velocity into the continuity equation, one obtains an equation for the electrostatic potential $`\varphi `$.
Before writing this equation we introduce new dimensional units for the variables.
$$\varphi ^{phys}\varphi ^{}=\frac{\left|e\right|\varphi ^{phys}}{T_e}$$
(187)
$$(x^{phys},y^{phys})(x^{},y^{})=(\frac{x^{phys}}{\rho _s},\frac{y^{phys}}{\rho _s})$$
(188)
$$t^{phys}t^{}=t^{phys}\mathrm{\Omega }_i$$
(189)
The new variables $`(t,x,y)`$ and the function $`\varphi `$ are non-dimensional. In the following the *primes* are not written.
### 9.1 The equation
With these variables the equation obtained is
$`{\displaystyle \frac{}{t}}\left(1_{}^2\right)\varphi `$
$`\left(_{}\varphi \times \widehat{๐ง}\right)v_d`$
$`+\left(_{}\varphi \times \widehat{๐ง}\right)๐ฏ_T\varphi `$
$`+\left[\left(_{}\varphi \times \widehat{๐ง}\right)_{}\right]\left(_{}^2\varphi \right)`$
$`=`$ $`0`$
where
$`v_d`$ $``$ $`_{}\mathrm{ln}n_0_{}\mathrm{ln}T_e`$ (191)
$`๐ฏ_T`$ $``$ $`_{}\mathrm{ln}T_e`$
(This is Eq.(8) from the paper ).
### 9.2 No temperature gradient
From the various versions of the nonlinear drift equation, (in particular the Hasegawa-Mima equation ) we choose the radially symmetric Flierl-Petviashvili soliton equation:
$$\left(1\rho _s^2_{}^2\right)\frac{\phi }{t}+v_d\frac{\phi }{y}v_d\phi \frac{\phi }{y}=0$$
(192)
where $`\rho _s=c_s/\mathrm{\Omega }_i`$, $`c_s=\left(T_e/m_i\right)^{1/2}`$ and the potential is scaled as $`\phi =\frac{L_n}{L_{T_e}}\frac{e\mathrm{\Phi }}{T_e}`$ . Here $`L_n`$ and $`L_T`$ are respectively the gradient lengths of the density and temperature. The velocity is the diamagnetic velocity $`v_d=\frac{\rho _sc_s}{L_n}`$. The condition for the validity of this equation are: $`\left(k_x\rho _s\right)\left(k\rho _s\right)^2\eta _e\frac{\rho _s}{L_n}`$, where $`\eta _e=\frac{L_n}{L_{T_e}}`$.
The exact solution of the equation is
$$\phi _s(y,t;y_0,u)=3\left(\frac{u}{v_d}1\right)\mathrm{sec}h^2\left[\frac{1}{2\rho _s}\left(1\frac{v_d}{u}\right)^{1/2}\left(yy_0ut\right)\right]$$
(193)
where the velocity is restricted to the intervals $`u>v_d`$ or $`u<0`$. In the Ref. the radial extension of the solution is estimated as: $`\left(\mathrm{\Delta }x\right)^2\rho _sL_n`$. In our work we shall assume that $`u`$ is very close to $`v_d`$ , $`uv_d`$ (i.e. the solitons have small amplitudes).
## 10 Appendix B : Connection between the MSR formalism and Onsager-Machlup
In our approach the most natural way of proceeding with a stochastic differential equation is to use the MSR type reasoning in the Jensen formulation. The equation is discretized in space and time and selected with $`\delta `$ functions in an ensemble of functions (actually in sets of arbitrary numbers at every point of discretization). The result is a functional integral. There is however a particular aspect that needs careful analysis, as mentioned in the previous Subsection. It is the problem of the Jacobian associated with the $`\delta `$ functions. This problem is discussed in Ref. and for consistency we include here the essential of their original treatment.
The equation they analyse is in the time domain and is presented in most general form as
$$\frac{\varphi _j\left(t\right)}{t}=\left(\mathrm{\Gamma }_0\right)_{jk}\frac{\delta H}{\delta \varphi _k\left(t\right)}+V_j\left[\varphi \left(t\right)\right]+\theta _j$$
where the number of stochastic equations is $`N`$ , $`H`$ is functional of the fields, $`V_j`$ is the streaming term which obeys a current-conserving type relation
$$\frac{\delta }{\delta \varphi _j}V_j\left[\varphi \right]\mathrm{exp}\left\{H\left[\varphi \right]\right\}=0$$
The noise is $`\theta _j`$.
The following generating functional can be written
$$Z_\theta =\text{D}\left[\varphi _j\left(t\right)\right]\mathrm{exp}๐t\left[l_j\varphi _j\left(t\right)\right]\underset{j,t}{}\delta \left(\frac{\varphi _j\left(t\right)}{t}+K_j\left[\varphi \left(t\right)\right]\theta _j\right)J\left[\varphi \right]$$
the functions $`l_j\left(t\right)`$ are currents,
$$K_j\left[\varphi \left(t\right)\right]\left(\mathrm{\Gamma }_0\right)_{jk}\frac{\delta H}{\delta \varphi _k\left(t\right)}+V_j\left[\varphi \left(t\right)\right]$$
and $`J\left[\varphi \right]`$ is the Jacobian associated to the Dirac $`\delta `$ functions in each point of discretization.
The Jacobian can be written
$$J=det\left[\left(\delta _{jk}\frac{}{t}+\frac{\delta K_j\left[\varphi \right]}{\delta \varphi _k}\right)\delta \left(tt^{}\right)\right]$$
Up to a multiplicative constant
$$J=\mathrm{exp}\left(Tr\mathrm{ln}\left[\left(\frac{}{t}+\frac{\delta K}{\delta \varphi }\right)\frac{\delta \left(tt^{}\right)}{\frac{}{t}\delta \left(tt^{}\right)}\right]\right)$$
or
$$J=\mathrm{exp}\left(Tr\mathrm{ln}\left[1+\left(\frac{}{t}\right)^1\frac{\delta K\left(t\right)}{\delta \varphi \left(t^{}\right)}\right]\right)$$
Since the operator $`\left(\frac{}{t}\right)^1`$ is retarded, only the lowest order term survives after taking the trace
$$J=\mathrm{exp}\left[\frac{1}{2}๐t\frac{\delta K_j\left[\varphi \left(t\right)\right]}{\delta \varphi _j\left(t\right)}\right]$$
The factor $`1/2`$ comes from value of the $`\mathrm{\Theta }`$ function at zero.
In the treatment which preserves the dual function $`\widehat{\varphi }`$ associated to $`\varphi `$ in the functional, there is a part of the action
$$\widehat{\varphi }K\left[\varphi \right]$$
Then a $`\widehat{\varphi }`$ and a $`\varphi `$ of the same coupling term from $`\widehat{\varphi }K\left[\varphi \right]`$ close onto a loop. Since $`G_{\widehat{\varphi }\varphi }`$ is retarded, all these contributions vanish except the one with a single propagator line. This cancels exactly, in all orders, the part coming from the Jacobian.
Then it is used to ignore all such loops and together with the Jacobian.
In conclusion we can compared the two starting points in a functional approach: The one that uses *dual functions* $`\varphi \left(t\right)`$ and $`\chi \left(t\right)`$, closer in spirit to MSR; And the approach based on Onsager-Machlup functional, traditionally employed for the determination of the probabilities. Either we keep $`\chi \left(t\right)`$ and ignore the Jacobian (the first approach) or integrate from the beginning over $`\chi \left(t\right)`$ and include the Jacobian. The approaches are equivalent but we have followed the first one.
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# UAHEP051 Growth of a Susy Bubble: Inhomogeneity Effects
## 1 Introduction
In the past couple of years there has been a burst of theoretical activity discussing possible transitions between string vacua of differing amounts of supersymmetry . These include a string theory study of a transition from a local minimum of positive vacuum energy to an exactly supersymmetric phase . It is reasonable, therefore, to consider possible phenomenological signals for such transitions. We live in a broken-susy phase with, apparently, a positive vacuum energy density
$`ฯต=3560MeV/m^3`$ (1.1)
leading to an acceleration in the expansion of the universe. The basic string theories, on the other hand, suggest a true vacuum of exact supersymmetry with massless ground state supermultiplets. In flat space these have zero vacuum energy. It is, therefore, interesting to consider whether we are living in a false vacuum which will ultimately decay to an exact susy ground state. False vacuum decays were treated in some generality many years ago . In a false vacuum with energy density $`ฯต`$ it is expected that bubbles of true vacuum with radius $`r`$ and surface tension $`S`$ will have an effective potential given by a sum of a volume term and a surface term,
$`V(vac)=4\pi (r^3ฯต/3r^2S).`$ (1.2)
For small $`ฯต`$ the surface tension $`S`$ can be treated as a constant. Bubbles will be constantly nucleated from the vacuum with a steeply falling distribution in initial radii. However, only those with initial radii greater than some critical radius
$`R_c(vac)={\displaystyle \frac{3S}{ฯต}}`$ (1.3)
will grow to effect a phase transition while smaller bubbles will be rapidly quenched. The probability per unit time per unit volume to produce a bubble of radius $`R_c`$ or greater and, therefore, to effect a phase transition to the true vacuum has been given in the form
$`{\displaystyle \frac{d^2P}{dtd^3r}}=Ae^{B(vac)}`$ (1.4)
where, assuming a thin wall between the phases,
$`B(vac)={\displaystyle \frac{27\pi ^2S^4}{2ฯต^3}}.`$ (1.5)
In this picture, the fact that our broken susy world has existed so long is due to the smallness of $`ฯต^3`$ relative to $`S^4`$. In the physical vacuum, once a bubble of critical radius is nucleated the surface will rapidly accelerate and engulf the entire universe. The fact that this has not happened as yet suggests that
$`R_c(vac)>R_{galaxy}4.710^{20}m`$ (1.6)
or
$`S>5.610^{23}MeV/m^2=210^{23}M_{}R_E^2.`$ (1.7)
In much of this article we use the solar mass, $`M_{}=1.210^{60}`$ MeV, and earth radius, $`R_E=6.3810^6`$ m as convenient units. We also use natural units $`\mathrm{}=c=1`$.
Reasonable expectations exist that vacuum decay will be accelerated in dense matter. Heuristically, this can be seen by noting that, in the vacuum case, $`ฯต`$ is the energy advantage per unit volume of making a transition to the exact susy phase. One might expect, therefore, that the above equations will be modified in dense matter by replacing $`ฯต`$ by the energy advantage per unit volume of trading the broken susy phase for the exact susy phase, i.e.
$`ฯตฯต+\mathrm{\Delta }\rho `$ (1.8)
where $`\mathrm{\Delta }\rho `$ is the ground state matter density in the broken susy phase minus the ground state matter density in the exact susy phase as shown in fig. 1.
The difference $`\mathrm{\Delta }\rho `$ is the fermionic excitation energy density. The parameter controlling the exponential factor in the transition rate would then be
$`B={\displaystyle \frac{27\pi ^2S^4}{2(ฯต+\mathrm{\Delta }\rho )^3}}.`$ (1.9)
The value of $`\mathrm{\Delta }\rho `$ in a white dwarf star is calculated as follows. In a degenerate electron gas of $`N`$ electrons in a volume $`V`$ the Fermi momentum is
$`p_F=\left({\displaystyle \frac{3\pi ^2N}{V}}\right)^{1/3}`$ (1.10)
with, assuming equal numbers of neutrons and protons,
$`N/V={\displaystyle \frac{\rho }{2M_N}}`$ (1.11)
$`M_N`$ being the nucleon mass. The average kinetic energy is
$`<E>m=m\left(1+{}_{2}{}^{}F_{1}^{}(1/2,3/2;5/2;p_F^2/m^2)\right).`$ (1.12)
In the limit of zero electron mass which we will use for simplicity, this is
$`<E>={\displaystyle \frac{3p_F}{4}}.`$ (1.13)
The kinetic energy density, which is equal to the difference in ground state energy densities between the broken susy state and the exact susy state, is then
$`\mathrm{\Delta }\rho ={\displaystyle \frac{3p_FN}{4V}}={\displaystyle \frac{1}{4\pi ^2}}\left({\displaystyle \frac{3\pi ^2\rho }{2M_N}}\right)^{4/3}.`$ (1.14)
We could double this estimate since a comparable contribution is expected from nuclear excitation energies but, for the present, we will neglect corrections by factors of order a few. From 1.9 we have
$`B=\left({\displaystyle \frac{\stackrel{~}{\rho }}{\rho }}\right)^4`$ (1.15)
with
$`\stackrel{~}{\rho }=\left({\displaystyle \frac{8}{3\pi ^6}}\right)^{1/4}SM_N.`$ (1.16)
The phase transition probability per unit time per unit volume, $`Ae^B`$, increases rapidly with the density of the medium until the density becomes of order $`\stackrel{~}{\rho }`$ at which point it saturates. For more dense media, the transition rate is proportional to the volume.
The longevity of the universe (eq.1.7) implies that
$`\stackrel{~}{\rho }>0.140M_{}R_E^3`$ (1.17)
not far from the average density of the prototype white dwarf
$`\rho _{WD}={\displaystyle \frac{3}{4\pi }}M_{}R_E^3.`$ (1.18)
This suggests that the susy phase transition rate per unit volume could be appreciable for white dwarf stars but negligible for less dense objects and for denser objects of much smaller volume. The relative transition rates in a variety of astrophysical and terrestrial objects treated as of constant density are tabulated in table 1. One can see here that if $`\stackrel{~}{\rho }`$ is close to its lower limit, the transition rate in white dwarfs is orders of magnitude greater than in the other considered bodies. In the future, it would be of interest to explore larger values of $`\stackrel{~}{\rho }`$ for which the transition rates for white dwarf stars and neutron stars are comparable. However, since the density distributions in neutron stars and the radiative processes after the neutron to sneutron pair conversion are greatly different from the current calculations, these investigations are beyond the scope of the present paper.
In two recent articles we have explored the possibility that such a transition in a dense star is the central engine of gamma ray bursts. The intense, collimated gamma radiation released in this way could provide the power to accelerate a macroscopic portion of the star to relativistic energies as in the cannonball model . For reviews of the susy star idea, see . In these articles we have presented zeroth order predictions based on a transition in a typical white dwarf star (e.g. Sirius B) neglecting density inhomogeneity effects. In the current paper, we proceed to incorporate these effects as well as consequences of existing variations in dense star mass and radii.
In dense matter one would expect the critical radius to be
$`R_c(r)={\displaystyle \frac{3S}{ฯต+\mathrm{\Delta }\rho (r)}}12\pi ^2S\left({\displaystyle \frac{2M_N}{3\pi ^2\rho (r)}}\right)^{4/3}.`$ (1.19)
In an inhomogeneous medium, a bubble of radius $`r`$ will grow as long as
$`r>R_c(r).`$ (1.20)
Even ignoring density inhomogeneity within the star, the critical radius could be quite small inside the star but jump at the surface to a value that is much larger than the radius of the star. This has the effect of efficiently confining the susy phase to the interior of the star.
In the exact susy phase the particle and sparticle have a common mass which we have assumed to be that of the particle in the broken phase. In order for the phase transition to proceed, it is a crucial assumption that the common mass is no greater than the particle mass. This assumption is, perhaps, supported by the fact that the ground state masses in the exactly supersymmetric string theories are zero. Also, susy models with radiative electroweak symmetry breaking (EWSB) relate the EWSB to the susy breakdown so that, in the absence of susy breaking, all particle masses vanish. There are no comparable physical models suggesting the opposite assumption, namely that the common mass in the exact susy phase is at higher energy. Nevertheless, this opposite assumption could also be considered but perhaps only briefly since, in this case, there is no exothermic phase transition unless, perhaps, the Fermi energy was greater than the mass difference.
If a SUSY bubble forms in a dense material medium, the interior of the bubble will find itself greatly out of equilibrium since many of the fermionic constituents occupy high energy levels whereas scalar particles could all occupy ground state energy levels. In such a situation particle pairs will rapidly convert to sparticle pairs. Electrons, for example, will pair convert into scalar electrons (selectrons) thus evading the energy storing property of the Pauli exclusion principle:
$`e^{}e^{}\stackrel{~}{e}^{}\stackrel{~}{e}^{}`$ (1.21)
Note that, although the final state is a selectron pair and not a selectron-antiselectron pair, this process does not require R parity violation.
Other types of phase transition models have also been proposed to explain gamma ray bursts. Among these are transitions to quark matter , transitions to a color superconducting state , and transitions to mirror fermions . The calculations of the current paper on the behavior of a susy bubble can also be applied to hypothetical bubbles of these phases if they begin in the high density regions.
The susy phase transition model applied to the typical white dwarf star neglecting density inhomogeneity predicts, correctly though roughly, in a relatively parameter-free way, the minimum duration of the burst (the light crossing time), the mean gamma ray energy (the mean electron kinetic energy), and the total burst energy (the total electron kinetic energy). In addition (and in distinction to the other phase transition models mentioned above), the fact that the final state of the susy transition consists of scalar constituents predicts a significant amount of jet collimation due to Bose enhancement.
In section II of this paper we review the density profile of white dwarf stars produced by the balance of inward gravitational pressure and outward electron degeneracy pressure. The electron momentum distribution differs greatly from that of a degenerate Fermi gas at uniform density. In section III we treat the classical free collapse time of an inhomogeneous susy star relieved of Pauli blocking, taking into account the full range of white dwarf stars. We also compute the variations in burst duration assuming that, in a dense star, the true vacuum bubble expands at the density dependent speed of sound which, in the high density limit approaches $`c/\sqrt{(}3)`$. Section IV is reserved for conclusions.
## 2 Density and momentum space gradients in a white dwarf star
Since Chandrasekhar , it has been axiomatic in astrophysics that isolated stars below a mass of about $`1.41`$ solar masses are stable against collapse due to electron exchange degeneracy. The density profile of such white dwarf stars is determined by hydrodynamic equilibrium between gravity and this outward degeneracy pressure initially augmented by thermal pressure.
The degeneracy pressure of the electrons in a white dwarf takes a simple form in either the extreme relativistic or non-relativistic limits. In intermediate regimes, it can be written in terms of the variable
$`x={\displaystyle \frac{\mathrm{}}{m_ec}}\left({\displaystyle \frac{3\rho (r)}{8\pi \mu _eM_N}}\right)^{1/3}=b\rho (r)^{1/3}`$ (2.1)
where $`\mu _e=A/Z=2`$ and $`M_N`$ is the nucleon mass. The degeneracy pressure is proportional to the function
$`f(x)={\displaystyle \frac{1}{8}}\left(x(2x^23)\sqrt{x^2+1}+3\mathrm{sinh}^1(x)\right);`$ (2.2)
Specifically,
$`P_d=af(x)`$ (2.3)
with
$`a={\displaystyle \frac{8\pi m_e^4c^5}{3\mathrm{}^3}}.`$ (2.4)
The degeneracy pressure gradient is
$`{\displaystyle \frac{dP_d}{dr}}={\displaystyle \frac{ab}{3}}\rho ^{2/3}f^{}(x){\displaystyle \frac{d\rho }{dr}}.`$ (2.5)
In low temperature equilibrium this must balance the gravitational pressure gradient
$`{\displaystyle \frac{dP_G}{dr}}={\displaystyle \frac{\rho (r)G_NM(r)}{r^2}}.`$ (2.6)
Here, $`M(r)`$ is the mass within radius $`r`$ and $`G_N`$ is the gravitational constant. The resulting integro-differential equation for $`\rho `$ can be solved by choosing an arbitrary starting value $`\rho _0`$ at the center of the star and integrating outward until the density falls to zero, recording at each step the value of $`M(r)`$. This defines the radius $`R`$ of the star and the corresponding mass $`M(R)`$ as a function of the peak (central) density. The resulting mass-radius relation is shown in fig. 2. In standard astrophysics, all isolated stars with a mass of less than $`1.41M_{}`$ will ultimately decrease in radius as they cool until they reach a point on the curve of fig. 2 at which point they become absolutely stable. In fig. 3 we plot the central and average density of white dwarfs at zero temperature as a function of the stellar mass. The central densities are more than an order of magnitude greater than the average densities.
The density distribution of these white dwarf stars as a function of distance from the center is a family of curves of which a representative seven are illustrated in fig. 4. Previous work on the susy phase transition in dense stars has ignored the strong density variation and relied on average densities only. The stellar density goes to zero at the surface of the star so, as can be seen from eq. 1.20, the susy bubble will stall at some distance from the surface creating a thin atmosphere of normal matter. The thickness of this shell is determined by the surface tension for which, at present, we know only the limit of eq.1.7. The above formulae imply, in this limit, a skin thickness of about $`100\mu m`$ for a star of solar mass and earth radius. The interior of the bubble constitutes a resonant cavity whose scalar constituents will continue to radiate gammas until the star radiates all its excitation energy or collapses under gravitational pressure.
In the inhomogeneous white dwarf, there is a local Fermi momentum given as in 1.10
$`p_F(r)=\left({\displaystyle \frac{3\pi ^2\rho (r)}{2m_N}}\right)^{1/3}.`$ (2.7)
The momentum distribution in the electron sea is no longer simply quadratic but takes the step function dependent form
$`{\displaystyle \frac{dN}{dp}}={\displaystyle \frac{4p^2}{\pi }}{\displaystyle _0^R}r^2๐r\theta (p_F(r)p).`$ (2.8)
The peak electron momentum and therefore the peak photon energy after the phase transition is given by eq. 2.7 evaluated at the stellar center, $`r=0`$. These peak energies, tabulated in table 2 are much greater than that of white dwarfs treated as of constant density as in ref. . Even higher peak energies, of course, can be found in stars closer to the Chandrasekhar limit since, there, the radius approaches zero and the density diverges. In this region there is a gradual approach to a neutron star structure. The gamma ray burst observations, for comparison, suggest a total burst energy of about $`510^{50}`$ ergs and mean photon energies between $`0.1`$ and $`1.0`$ MeV.
In table 2 we record various properties of these seven.
## 3 Factors affecting burst duration
In the susy phase transition model there are several physical effects influencing the burst duration. These are
1. The bubble growth time. This is the time it takes for a bubble nucleated in a high density region near the center to grow to the stellar surface. A lower limit to this time is the light crossing time but, more plausibly, in dense matter the bubble surface should expand at some rate comparable to the speed of sound in matter of that density.
2. The light crossing time. After the bubble has engulfed the star, there could be an additional time required for light emitted on the far side of the star to cross the stellar diameter. At high density, these escaping photons might undergo a random walk leading to a time proportional to the square of the radius divided by a mean free path. However, it has been observed that the Landau-Pomeranchuk-Migdal effect would greatly increase the transparency of dense matter to gamma rays or, equivalently, the mean free path.
3. The free collapse time. The conversion of fermions to bosons following a susy phase transition eliminates the degeneracy pressure and the star will undergo gravitational collapse. Classically, the star would collapse to a point in a matter of seconds but, once a radius a few times the Schwarzschild radius is achieved, general relativistic effects begin to dominate and will cause collapse to the Schwarzschild radius, as seen by a distant observer, to require an infinite amount of time during which photon emission will be progressively red-shifted leading to a certain amount of afterglow below the gamma ray spectrum. In ref. it was found that the electrons higher in the Pauli tower have a higher probability of pair conversion implying an earlier pair conversion than that of lower energy electrons. This implies a natural progression to smaller frequencies as the burst progresses. The growth of the susy bubble into lower density regions has similar consequences. These effects by themselves have a much shorter time scale than observed afterglows and therefore are relevant only if the gravitational collapse and bubble cooling are greatly slowed by the following consideration.
4. Radiation pressure. Radiation released by nucleons converting to scalar nucleons, which then drop into the nuclear ground state, is expected to rapidly expand the star and, thus, decrease the stellar density. This will increase the bubble growth time, the light crossing time, and the stellar collapse time. Density waves could be created in the wake of the initial blast and might lead to the rapid time variability observed in gamma ray bursts. Existing studies of susy bubble behavior, including the present article, do not include the effects of radiation pressure.
We will study first the bubble growth time assuming this is governed approximately by the speed of sound in dense matter. The speed of sound at radius $`r`$ depends on the local pressure and density and is given by
$`v_s(r)=\sqrt{{\displaystyle \frac{3P}{\rho }}}.`$ (3.1)
In the case of constant density this takes the simple form
$`v_s(r)=\sqrt{{\displaystyle \frac{2}{3}}\gamma \pi G_N\rho (R^2r^2)}.`$ (3.2)
Here $`\gamma `$ is the ratio of specific heats ($`5/3`$ for a monatomic gas), and $`R`$ is the stellar radius at which the pressure vanishes. eq. 3.2 may be derived by considering the downward force exerted on a column of matter from radius $`r`$ to the surface. The bubble growth time is then
$`\tau ={\displaystyle _0^R}๐r/v_s(r)={\displaystyle \frac{\pi R}{2v_s(0)}}.`$ (3.3)
This time is about $`2`$ s for our typical white dwarf but given the variations in density among the full sample of white dwarfs, the growth times based on the average densities as recorded in table 2 have a ratio of maximum to minimum of about $`148`$. This is close to the observed ratio of short bursts but far from the observed ratio of about $`10,000`$ including long bursts and short bursts together. In any case we know from figure 4 that the densities are rapidly varying especially for the higher mass dwarfs. In the case of non-constant density the pressure at radius $`r`$ remains governed by the differential equation of eq. 2.6. One begins with zero pressure at the stellar surface and integrates inward to find the pressure as a function of distance from the center. The local speed of sound is given by eq. 3.1 subject to the limit $`c/\sqrt{(}3)`$ in the high density limit. The bubble growth time is given by the left-most equality of eq. 3.3 . The stellar mass, radius, and bubble growth time are then given as a function of the central density. The probability per unit time of a susy phase transition in a star of radius R is
$`{\displaystyle \frac{1}{N}}{\displaystyle \frac{dN}{dt}}=4\pi A{\displaystyle _0^R}r^2๐re^{(\frac{\stackrel{~}{\rho }}{\rho (r)})^4}.`$ (3.4)
Because of the exponential suppression at low density, $`\stackrel{~}{\rho }`$ is a key parameter in identifying the gamma ray burst progenitor in the phase transition model. Although $`A`$ is a free parameter at this point, the dominant contributions to the bursts are likely to come from the largest bodies of density $`\stackrel{~}{\rho }`$ or greater. The assumption here is that $`A`$ is not strongly density dependent.
It is interesting to explore the possibility that $`\stackrel{~}{\rho }`$ is close to its lower limit from eq. 1.17. This would imply, on the cosmological time scale, an imminent end to our current world of broken symmetry. If, on the other hand, $`\stackrel{~}{\rho }`$ is very much greater than its lower limit, the bursts would come primarily from denser objects than most isolated white dwarfs, i.e. either neutron stars or larger bodies in the process of gravitational collapse. In this case, density enhancement through accretion such as in the collapsar model , could play a role. However, since the Fermi momentum of the electron sea is proportional to $`\rho ^{1/3}`$ as in eq. 1.10, the peak photon energy potentially gives, in the susy star model, an upper limit on $`\stackrel{~}{\rho }`$.
In the current paper we will explore the possibilities that $`\stackrel{~}{\rho }`$ is related to its lower limit by a factor of 1,5, or 25. The rate of bursts as a function of total progenitor mass is given by multiplying 3.4 by the number of white dwarf stars of given radius or, equivalently, of given mass. We use the Sloan survey mass distribution given by assuming that the relative numbers do not change greatly when extrapolated to zero temperature. The distribution is strongly peaked at $`0.56`$ solar masses as shown in fig.5.
The observed burst rate is about
$`{\displaystyle \frac{dN}{dt}}510^7yr^1gal^1.`$ (3.5)
but the local rate is much lower suggesting that the fuel of gamma ray bursts has been largely exhausted and we are now primarily seeing bursts that happened long ago in more distant galaxies. The observed sample of bursts with identified redshifts seems to cluster around redshift $`1`$. This interpretation is, perhaps, supported by the observed shortage of low luminosity dwarfs in our galaxy and the, at first glance contradictory, apparent excess of dark objects of white dwarf mass in the galactic halo .
From the total energy release, table 3 might suggest a $`\stackrel{~}{\rho }`$ value near or slightly above $`3.5M_{}R_E^3`$. However, the investigation of beaming and other contributions to the energy release is at too early a stage to rule out larger values of $`\stackrel{~}{\rho }`$. Since the bubble growth time is a monotonic function of dwarf mass, we can use the mass distribution of fig. 5 and the transition probability of eq. 3.4 to calculate the shape of the growth time distribution. This is shown in table. 3 for the three chosen values of $`\stackrel{~}{\rho }`$. It is encouraging that the growth times are in the range of the observed short burst durations but a critical test of the model must await the incorporation of the other factors influencing the burst duration as enumerated above.
One of these factors is the collapse time of a star relieved of Pauli blocking. In we have noted that the classical collapse time of a susy white dwarf is
$`\tau _c={\displaystyle \frac{\pi }{2}}\left({\displaystyle \frac{8\pi G_N\rho }{3}}\right)^{1/2}.`$ (3.6)
Although that paper considered only stars of constant density, the collapse time remains the same for inhomogeneous stars as long as one uses the average density. The classical collapse times of the representative sample of white dwarfs is given in table 3. However once the star approaches the Schwarzschild radius, general relativistic effects dilate the collapse time as seen by a distant observer. The approach to the Schwarzschild radius, $`r^{}=\frac{2G_NM}{c^2}`$ is given by
$`rr^{}=(r_0r^{})e^{(tt_0)c/r^{}}.`$ (3.7)
During this time the star can still radiate although the photon energies will be redshifted from their emission energies according to the relation
$`E=E_{em}{\displaystyle \frac{e^{(tt_0)c/2r^{}}}{\sqrt{(}1+e^{(tt_0)c/r^{}})}}.`$ (3.8)
The time dependence of this component of the afterglow is independent of frequency. The time constant, $`c/r^{}`$ for this contribution to the afterglow will be of order of ten microseconds for a star of near solar mass. With currently available techniques, it will be impossible to observe such a short afterglow. Another source of afterglow will be the emission from circumstellar material irradiated by the burst. This second source may be absent if the burst comes from the decay of an isolated star. Since, at present, afterglows have only been definitely observed for some of the long bursts, it might be interesting to consider a proposal where all of the short bursts and many of the long bursts originate in isolated stars while the bursts with extensive afterglows originate in stars with significant circumstellar material as in the collapsar model. A very massive star in the process of gravitational collapse will necessarily pass through stages of high fermion degeneracy where the possibility of a susy phase transition might become significant even though the time spent in these stages is not long.
## 4 Conclusions
We have considered the effects of mass variations and density inhomogeneities in white dwarf stars under the assumption that such stars experience a phase transition to the exact susy ground state. It is clear that much work remains to be done. Nevertheless, it is encouraging that the basic assumption with few free parameters produces a gamma ray emission on the right time scale, with mean gamma energy and total energy release close to observations. The density inhomogeneities increase the expected peak photon energy by an order of magnitude over that of constant density stars of the same average density.
The assumption that the surface tension of a susy bubble is independent of density needs to be examined and probably relaxed. Similarly, the strong density dependence of eqs. 1.4,1.9 might be somewhat softened by corrections to the thin wall approximation of the vacuum decay studies . Effects due to the radiation pressure at non-zero temperature need to be incorporated and may significantly affect the duration distribution. Radiation from the collapse of the Pauli tower in nuclei as well as the enhanced energy release from snuclear fusion and beta decay need to be studied. The resulting radiation pressure is expected to slow the final collapse of a star relieved of Pauli blocking. In addition, we hope to explore the possibility that the profound nuclear explosion at the stellar center due to the phase transition sets up density standing waves. This could be related to the spikey behavior of the observed gamma ray bursts.
Acknowledgements
This work was supported in part by the US Department of Energy under grant DE-FG02-96ER-40967. We gratefully acknowledge discussions with Paul Cox, Ben Harms, Irina Perevalova, George Karatheodoris, Phil Hardee, Bill Keel, and Sanjoy Sarkar.
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# Coherent ฮโบ and ฮโข(1520) photoproduction off the deuteron
## I Introduction
The first evidence for the pentaquark hadron $`\mathrm{\Theta }^+`$ discovered by the LEPS collaboration at SPring-8 Nakano03 was subsequently confirmed in other experiments OtherPenta . However, some other experiments failed to find the $`\mathrm{\Theta }^+`$ signal (for a review see Hicks ; Kabana05 ). Since then the situation concerning the existence of the pentaquarks remained controversial. Independent studies of the manifestation of a $`\mathrm{\Theta }^+`$ state in different processes are, therefore, urgently desired.
$`\mathrm{\Theta }^+`$ photoproduction in the reaction $`\gamma DnpK^+K^{}`$ seems to be very interesting and important NakanoP04 ; TedeschiP04 . First, it allows to study simultaneously the $`\gamma p\mathrm{\Lambda }(1520)K^+`$ and $`\gamma n\mathrm{\Theta }^+K^{}`$ subreactions characterized by the similarity in the production mechanisms, i.e. both processes are described by the same set of the tree level Feynman diagrams Hosaka0503 ; TEHN04 ; Hosaka0505 . Therefore, one hopes to define the ratio of $`\mathrm{\Theta }^+`$ to $`\mathrm{\Lambda }(1520)`$ photoproduction with minimal uncertainty of the production mechanisms, which is important for understanding the nature of $`\mathrm{\Theta }^+`$. Second, in case of the $`\gamma D`$ interaction one can study qualitatively a new basic process - the coherent $`\mathrm{\Theta }^+\mathrm{\Lambda }(1520)`$ photoproduction. This reaction has its own physics interest and unambiguously will shed new light to pentaquark properties and the mechanism of the $`\mathrm{\Theta }^+`$ photoproduction.
It is commonly supposed now that the total width of the $`\mathrm{\Theta }^+`$ is as small as $`\mathrm{\Gamma }_\mathrm{\Theta }1`$ MeV SmallWidth , being much smaller than the total $`\mathrm{\Lambda }^{}`$ decay width, $`\mathrm{\Gamma }_\mathrm{\Lambda }^{}15.6`$ MeV PDG . (Throughout this paper, for simplicity, we use notation $`\mathrm{\Lambda }^{}\mathrm{\Lambda }(1520)`$.) This means that the most promising way for studying the coherent $`\mathrm{\Lambda }^{}\mathrm{\Theta }^+`$ production is to analyze the invariant $`nK^+`$ mass, $`M_{nK^+}`$, distribution at fixed invariant mass of the $`pK^{}`$ pair, $`M_{pK^{}}`$. The enhancement of the $`\mathrm{\Theta }^+`$ photoproduction, when $`M_{pK^{}}`$ is in the vicinity of the $`\mathrm{\Lambda }^{}`$ mass, would indicate the manifestation of the coherent $`\mathrm{\Lambda }^{}\mathrm{\Theta }^+`$ photoproduction. This particular channel will appear in strong competition with the resonant and non-resonant background processes. By the notation โresonant processโ we mean, for example, the $`\mathrm{\Lambda }^{}`$ photoproduction from the proton inside the deuteron, when the neutron is a spectator, and similarly the $`\mathrm{\Theta }^+`$ photoproduction from a neutron, when the deuteronโs proton is a spectator. The notation โnon-resonantโ process denotes $`K^+K^{}`$ photoproduction from a nucleon without excitation of $`\mathrm{\Lambda }^{}`$ or $`\mathrm{\Theta }^+`$. It is clear that the coherent photoproduction and the background processes must be analyzed together using the same theoretical approaches. This allows to define the kinematical conditions where the coherent channel manifests itself clearly above strong background processes.
The aim of the present paper is to discuss these important aspects. Our model includes the elementary subprocesses of $`\gamma N\mathrm{\Lambda }^{}K`$ and $`\gamma N\mathrm{\Theta }^+\overline{K}`$ reactions. For the latter ones we use a model based on the effective Lagrangian approach of Ref. TEHN04 which is, generally speaking, similar to the models developed by other authors in Refs. Hosaka03 ; OKL031 ; NT03 ; LiuKo031 ; Zhao03 ; ZhaoKhal04 ; Oh-2 ; CloseZhao ; Roberts04 ; Mart:2004at ; Oh:2004wp . All these approaches predict the approximate equality of the cross sections of the $`\gamma n\mathrm{\Theta }^+K^{}`$ and $`\gamma p\mathrm{\Theta }^+\overline{K}^0`$ reactions. This equality may be changed into a suppression of the $`\gamma p\mathrm{\Theta }^+\overline{K}^0`$ transition Hosaka0503 ; Vita . However, we are going to demonstrate that the amplitude of the coherent $`\mathrm{\Lambda }^{}\mathrm{\Theta }^+`$ photoproduction, when $`\mathrm{\Lambda }^{}`$ is produced in the forward hemisphere in the $`\gamma D`$ center of mass system, is defined by the product of the $`\mathrm{\Lambda }^{}`$ photoproduction amplitude in $`\gamma N`$ interaction and the amplitude of the $`\mathrm{\Theta }^+NK`$ transition. In other words, the coherence effect of the $`\mathrm{\Lambda }^{}\mathrm{\Theta }^+`$ photoproduction in the forward hemisphere does not depend on the $`\mathrm{\Theta }^+`$ photoproduction amplitude and remains finite even if the cross section of the $`\gamma p\mathrm{\Theta }^+\overline{K}^0`$ reaction is vanishing. The coherence effect in the backward hemisphere is sensitive to the $`\mathrm{\Theta }^+`$ photoproduction amplitude, and it is suppressed in parallel with the suppression of the $`\gamma p\mathrm{\Theta }^+\overline{K}^0`$ reaction.
Our paper is organized as follows. In Sec. II, we discuss the resonant $`\mathrm{\Theta }^+`$ and $`\mathrm{\Lambda }^{}`$ photoproduction from a nucleon. In Sec. III, we consider the coherent $`\gamma D\mathrm{\Lambda }^{}\mathrm{\Theta }^+`$ reaction. Our model is similar to the approach of Ref. gammaD , developed for coherent $`\mathrm{\Theta }^+\mathrm{\Lambda }(\mathrm{\Sigma }^0)`$ photoproduction from a deuteron. In Sec. IV, we discuss the background processes. We start thereby from an analysis of the non-resonant background in โelementaryโ $`\gamma N\mathrm{\Theta }^+\overline{K}`$ and $`\gamma N\mathrm{\Lambda }^{}K`$ reactions. Then we apply these subprocesses to an analysis of the background spectator channels. Finally, we estimate the contribution of the coherent semi-resonant processes, which differ from the coherent photoproduction by the replacement of one hyperon by $`NK`$ or $`N\overline{K}`$ pairs. The results of our numerical calculations are presented in Sec. V. The summary is given in Sec. VI. In Appendix A, we show an explicit form of the transition operators for the resonance amplitude.
## II photoproduction from a nucleon
### II.1 $`๐ฏ^\mathbf{+}`$ photoproduction
The main diagrams for the amplitude of the resonance $`\mathrm{\Theta }^+`$ photoproduction in the reaction $`\gamma NNK\overline{K}`$ are shown in Fig. 1. We neglect here the contribution resulting from the photon interacting with the final decay vertex NT03 . In view of the chosen kinematics, where the invariant mass of the final $`KN`$ pair is near the resonance position, this is a good approximation since in the neglected graphs the $`\mathrm{\Theta }^+`$ is far off-shell and the graphs of Figs. 1 a - d dominate the resonance contribution. From a formal point of view gauge invariance is lost without contributions arising from the electromagnetic interaction in the decay vertex.
However, following Ref. hhgauge for the initial photoproduction process, we construct an overall conserved current by an appropriate choice of the contact term of Fig. 1d.
In this section $`k`$, $`p`$, $`q`$, $`\overline{q}`$, and $`p^{}`$ denote the four-momenta of the incoming photon, the initial nucleon, the outgoing $`K`$ and $`\overline{K}`$ mesons, and the recoil nucleon, respectively. The standard Mandelstam variables for the virtual $`\mathrm{\Theta }^+`$ photoproduction are defined by $`t=(\overline{q}k)^2`$, $`sW^2=(p+k)^2`$. The $`\overline{K}`$ meson production angle $`\theta `$ in the center-of-mass system (c.m.s.) is given by $`\mathrm{cos}\theta =๐ค\overline{๐ช}/(|๐ค||\overline{๐ช}|)`$, and the corresponding solid angle is $`\mathrm{\Omega }`$. We consider the integrated $`\mathrm{\Theta }^+`$ decay distribution. The differential cross section $`\gamma N\mathrm{\Theta }^+\overline{K}NK\overline{K}`$ as a function the $`\overline{K}`$ meson production angle and $`NK`$ invariant mass, $`M_{nK^+}`$, at the resonance position with $`M_{nK^+}=M_\mathrm{\Theta }=1.54`$ GeV is related to the cross section of the $`\mathrm{\Theta }^+`$ photoproduction in the $`\gamma N\mathrm{\Theta }^+\overline{K}`$ reaction as
$`{\displaystyle \frac{d\sigma _{fi}^R}{d\mathrm{\Omega }dM_{nK^+}}}|_{M_{nK^+}=M_\mathrm{\Theta }}={\displaystyle \frac{1}{\pi \mathrm{\Gamma }_\mathrm{\Theta }}}{\displaystyle \frac{d\sigma _{fi}^{\mathrm{\Theta }^+}}{d\mathrm{\Omega }}}`$ (1)
with $`\mathrm{\Gamma }_\mathrm{\Theta }`$ as $`\mathrm{\Theta }^+`$ decay width and
$`{\displaystyle \frac{d\sigma _{fi}^{\mathrm{\Theta }^+}}{d\mathrm{\Omega }}}={\displaystyle \frac{1}{64\pi ^2s}}{\displaystyle \frac{p_{\mathrm{out}}}{p_{\mathrm{in}}}}{\displaystyle \frac{1}{4}}{\displaystyle \underset{m_i,m_f,\lambda _\gamma }{}}|A_{m_f;m_i,\lambda _\gamma }^{\mathrm{\Theta }^+}|^2.`$ (2)
Here, $`A^{\mathrm{\Theta }^+}`$ is the $`\mathrm{\Theta }^+`$ photoproduction amplitude in the $`\gamma N\mathrm{\Theta }^+\overline{K}`$ reaction, $`m_i`$ and $`m_f`$ are the nucleon and $`\mathrm{\Theta }^+`$ spin projections, respectively, and $`\lambda _\gamma `$ denotes the incoming photon helicity; $`p_{\mathrm{in}}`$ and $`p_{\mathrm{out}}`$ are the relative momenta in the initial and the final states in c.m.s., respectively. Further on we will concentrate on the calculation of $`A^{\mathrm{\Theta }^+}`$. For simplicity, in this analysis we limit our consideration to the isoscalar, spin-1/2 $`\mathrm{\Theta }^+`$. Generalization for higher spin Hosaka0505 may be done in a straightforward manner.
The effective Lagrangians which define the Born terms for the diagrams shown in Fig. 1a - d are discussed in many papers (for references see Ref. TEHN04 ). Note that different phase conventions are often employed. Therefore, for the sake of definiteness, we list here the effective Lagrangians used in the present work<sup>1</sup><sup>1</sup>1Throughout this paper, isospin operators will be suppressed in all Lagrangians and matrix elements for simplicity. They can be easily accounted for in the corresponding coupling constants.:
$`_{\gamma KK}`$ $`=`$ $`ie(K^{}^\mu K^+K^+^\mu K^{})A_\mu ,`$ (3a)
$`_{\gamma \mathrm{\Theta }\mathrm{\Theta }}`$ $`=`$ $`e\overline{\mathrm{\Theta }}\left(\gamma _\mu {\displaystyle \frac{\kappa _\mathrm{\Theta }}{2M_\mathrm{\Theta }}}\sigma _{\mu \nu }^\nu \right)A^\mu \mathrm{\Theta },`$ (3b)
$`_{\gamma NN}`$ $`=`$ $`e\overline{N}\left(e_N\gamma _\mu {\displaystyle \frac{\kappa _N}{2M_N}}\sigma _{\mu \nu }^\nu \right)A^\mu N,`$ (3c)
$`_{\mathrm{\Theta }NK}^{\pm [\text{pv}]}`$ $`=`$ $`{\displaystyle \frac{g_{\mathrm{\Theta }NK}}{M_\mathrm{\Theta }\pm M_N}}\overline{\mathrm{\Theta }}\mathrm{\Gamma }_\mu ^\pm (^\mu K)N+\text{h.c.},`$ (3d)
$`_{\gamma \mathrm{\Theta }NK}^{[\mathrm{pv}]}`$ $`=`$ $`i{\displaystyle \frac{eg_{\mathrm{\Theta }NK}}{M_\mathrm{\Theta }\pm M_N}}\overline{\mathrm{\Theta }}\mathrm{\Gamma }_\mu ^\pm A^\mu KN+\text{h.c.},`$ (3e)
$`_{\mathrm{\Theta }NK}^{\pm [\text{ps}]}`$ $`=`$ $`ig_{\mathrm{\Theta }NK}\overline{\mathrm{\Theta }}\mathrm{\Gamma }^\pm KN+\text{h.c.},`$ (3f)
$`_{\gamma KK^{}}`$ $`=`$ $`{\displaystyle \frac{eg_{\gamma KK^{}}}{M_K^{}}}ฯต^{\alpha \beta \mu \nu }_\alpha A_\beta _\mu \overline{K}_\nu ^{}K+\text{h.c.},`$ (3g)
$`_{\mathrm{\Theta }NK^{}}^\pm `$ $`=`$ $`g_{\mathrm{\Theta }NK^{}}\overline{\mathrm{\Theta }}\mathrm{\Gamma }^{}\left(\gamma _\mu {\displaystyle \frac{\kappa ^{}}{M_\mathrm{\Theta }+M_N}}\sigma _{\mu \nu }^\nu \right)\overline{K}^\mu N+\text{h.c.},`$ (3h)
where $`A^\mu ,\mathrm{\Theta }`$, $`K`$, and $`N`$ are the photon, $`\mathrm{\Theta }^+`$, kaon, and the nucleon fields, respectively, $`K^{}`$ stands for the vector kaon field; $`\mathrm{\Gamma }_\mu ^\pm \mathrm{\Gamma }^\pm \gamma _\mu `$ (with $`\mathrm{\Gamma }^+=\gamma _5`$ and $`\mathrm{\Gamma }^{}=1`$ for positive and negative parity, respectively), $`e_p=1`$, $`e_n=0`$, and $`\kappa _N`$ denotes the nucleon anomalous magnetic moment ($`\kappa _p=1.79`$ and $`\kappa _n=1.91`$), $`\kappa _\mathrm{\Theta }`$ stands for the anomalous magnetic moment of $`\mathrm{\Theta }^+`$ and $`\kappa ^{}`$ denotes the tensor coupling of nucleon and strange vector mesons. The superscripts โPSโ and โPVโ correspond to the pseudo-scalar and pseudo-vector $`\mathrm{\Theta }^+NK`$ coupling schemes. Equation (3e) describes the contact (Kroll-Ruderman) interaction in the pseudo-vector coupling scheme (see Fig. 1d), which does not appear in case of the pseudo-scalar coupling (cf. Eq. (3f)).
In calculating the invariant amplitudes we dress the vertices by form factors. In the present tree-level approach and within our chosen kinematics, only the lines connecting the electromagnetic vertex with the initial $`\mathrm{\Theta }^+KN`$ vertex correspond to off-shell hadrons. We describe the product of both the electromagnetic and the hadronic form-factor contributions along these off-shell lines by the covariant phenomenological function
$`F(M,p^2)={\displaystyle \frac{\mathrm{\Lambda }^4}{\mathrm{\Lambda }^4+(p^2M^2)^2}},`$ (4)
where $`p`$ is the corresponding off-shell four-momentum of the virtual particle, $`M`$ denotes its mass, and $`\mathrm{\Lambda }`$ stands for the cut-off parameter. The electromagnetic current of the complete amplitude is conserved by making the initial photoproduction process gauge invariant. To this end, we apply the gauge invariance prescription by Haberzettl hhgauge with the modification by Davidson and Workman DavWork to construct a contact term for the initial process $`\gamma N\mathrm{\Theta }^+\overline{K}`$ free of kinematical singularities. We emphasize that contributions of the latter type are necessary even for pure pseudo-scalar coupling.
Since the coupling scheme and the $`\mathrm{\Theta }^+`$ parity are unknown one has to define the corresponding parameters in such a way to get the corresponding cross sections independently of $`\mathrm{\Theta }^+`$ parity and coupling scheme. We follow Ref. TEHN04 , where parameters of the model are fixed by a comparison of the resonant $`\mathrm{\Theta }^+`$ photoproduction cross section and non-resonant background with experiment and it is shown that one can find such a parameter set which parallels the prediction for PS and PV couplings and for positive and negative $`\mathrm{\Theta }^+`$ parity states as well, at least for the unpolarized, single and double polarization spin observables. Therefore, we can limit the present analysis to the PS coupling and a positive $`\mathrm{\Theta }^+`$ parity.
The resonance amplitudes obtained for the $`\gamma n`$ and $`\gamma p`$ reactions read
$`A_{fi}^{\mathrm{\Theta }^+}(\gamma n)=\overline{u}_\mathrm{\Theta }(p_\mathrm{\Theta })\left[_{}^{s}{}_{\mu }{}^{}+_{}^{t}{}_{\mu }{}^{}+_{}^{u}{}_{\mu }{}^{}+_{}^{c}{}_{\mu }{}^{}+_{}^{t}{}_{\mu }{}^{}(K^{})\right]u_n(p)\epsilon ^\mu ,`$ (5a)
$`A_{fi}^{\mathrm{\Theta }^+}(\gamma p)=\overline{u}_\mathrm{\Theta }(p_\mathrm{\Theta })\left[_{}^{s}{}_{\mu }{}^{}+_{}^{u}{}_{\mu }{}^{}+_{}^{c}{}_{\mu }{}^{}+_{}^{t}{}_{\mu }{}^{}(K^{})\right]u_p(p)\epsilon ^\mu .`$ (5b)
The explicit forms of the transition operators $`_\mu ^i`$ for the $`\gamma n\mathrm{\Theta }^+K^{}`$ and $`\gamma p\mathrm{\Theta }^+\overline{K}^0`$ reactions are exhibited in Appendix A.
For a positive $`\mathrm{\Theta }^+`$ parity the coupling constant $`g_{\mathrm{\Theta }NK}`$ is found from the $`\mathrm{\Theta }^+`$ decay width as
$`\mathrm{\Gamma }_\mathrm{\Theta }={\displaystyle \frac{[g_{\mathrm{\Theta }NK}]^2p_F}{2\pi M_\mathrm{\Theta }}}(\sqrt{M_N^2+p_F^2}M_N).`$ (6)
We choose a small width, $`\mathrm{\Gamma }_\mathrm{\Theta }=1`$ MeV SmallWidth , assuming that the observed width in the invariant mass distribution is determined by the experimental resolution. The magnitude of the coupling constant $`g_{\gamma KK^{}}`$ is extracted from the width of the $`K^{}\gamma K`$ decay PDG . Its sign is fixed by SU(3) symmetry. This delivers $`eg_{\gamma K^0K^0}=0.35`$ and $`eg_{\gamma K^\pm K^\pm }=0.23`$. The contribution of the $`s`$-channel (Fig. 1b) is small causing to a rather weak dependence of the total amplitude on the tensor coupling $`\kappa _\mathrm{\Theta }`$ in the $`\gamma \mathrm{\Theta }\mathrm{\Theta }`$ vertex within a โreasonableโ range of $`0|\kappa _\mathrm{\Theta }|0.5`$ MagMom . Therefore, we can choose $`\kappa _\mathrm{\Theta }=0`$. The coupling constant $`g_{\mathrm{\Theta }NK^{}}`$ is written as $`g_{\mathrm{\Theta }NK^{}}=\alpha _\mathrm{\Theta }g_{\mathrm{\Theta }NK}`$, where the parameter $`\alpha _\mathrm{\Theta }`$ depends on the choice of the tensor coupling $`\kappa ^{}`$ in Eq. (3h) and cut-off parameters $`\mathrm{\Lambda }_K^{}`$ in the form factors of the $`K^{}`$ exchange amplitude. Increasing value of $`\mathrm{\Lambda }_K^{}`$ leads to a decreasing $`\alpha _\mathrm{\Theta }`$. Following Ref. TEHN04 we use $`\mathrm{\Lambda }_K^{}=1.5`$ GeV and $`\alpha _\mathrm{\Theta }=1.875`$ at $`\kappa ^{}=0`$. This value of $`\alpha _\mathrm{\Theta }`$ is close to the quark model estimates $`\alpha _\mathrm{\Theta }=\sqrt{3}`$ QMKK\* .
Another cut-off parameter, $`\mathrm{\Lambda }_B`$, defines the Born terms of the $`s`$-, $`u`$-, and $`t`$-channels and the current-conserving contact terms. Note that the inclusion of the $`\mathrm{\Sigma }`$ and $`\mathrm{\Lambda }`$ photoproduction processes LambdaSigma results in a larger ambiguity in the choice of $`\mathrm{\Lambda }_B`$ which varies from 0.5 to 2 GeV depending on the coupling scheme and the method of conserving the electromagnetic current etc. The analysis of the vector meson photoproduction TitovLee and $`\gamma n\mathrm{\Theta }^+K^{}`$ favor a small value of the cut-off, $`\mathrm{\Lambda }_B0.5`$ GeV. For the $`\gamma p\mathrm{\Theta }^+\overline{K}^0`$ reaction the $`K^{}`$ exchange channel remains to be dominant at $`\mathrm{\Lambda }_B1.5`$ GeV and, therefore, in this paper we use a โuniversalโ value, $`\mathrm{\Lambda }_B0.5`$ GeV, for all Born terms.
In Fig. 2 we exhibit the differential cross sections of the reactions $`\gamma n\mathrm{\Theta }^+K^{}`$ (a) and $`\gamma p\mathrm{\Theta }^+\overline{K}^0`$ (b) in the c.m.s. at $`E_\gamma =2`$ GeV. One can see that the $`t`$-channel $`K^{}`$ exchange depicted in Fig. 1e gives the dominant contribution compared to the Born terms shown in Fig. 1a - d in both reactions.
### II.2 $`๐ฒ\mathbf{(}\mathrm{๐๐๐๐}\mathbf{)}`$ photoproduction
The main diagrams for the amplitudes of the excitation of the $`\mathrm{\Lambda }`$ hyperon in the $`\gamma NNK\overline{K}`$ reaction at low energies are shown in Fig. 3. Similarly to the $`\mathrm{\Theta }^+`$ photoproduction we neglect the photon interaction within the decay vertex and restore the gauge invariance by a proper choice of the contact terms. The Mandelstam variables for the virtual $`\mathrm{\Lambda }^{}`$ photoproduction are defined by $`t=(qk)^2`$, $`sW^2=(p+k)^2`$. The $`K`$ meson production angle $`\theta `$ (in $`\gamma p`$ c.m.s.) is given by $`\mathrm{cos}\theta =๐ค๐ช/(|๐ค||๐ช|)`$.
For the description of the $`\mathrm{\Lambda }^{}`$ excitation with $`J^P=\frac{3}{2}^{}`$ we use the following effective Lagrangians 3-2R ; TitovLee
$`_{\mathrm{\Lambda }^{}NK}`$ $`=`$ $`{\displaystyle \frac{g_{\mathrm{\Lambda }^{}NK}}{M_\mathrm{\Lambda }^{}}}\overline{\mathrm{\Lambda }}_{}^{}{}_{\mu }{}^{}\theta ^{\mu \nu }(Z)(_\nu \overline{K})\gamma _5N+\text{h.c.},`$ (7a)
$`_{\gamma \mathrm{\Lambda }^{}NK}`$ $`=`$ $`i{\displaystyle \frac{eg_{\mathrm{\Lambda }^{}NK}}{M_\mathrm{\Lambda }^{}}}\overline{\mathrm{\Lambda }}_{}^{}{}_{\mu }{}^{}\gamma _5A^\mu \overline{K}N+\text{h.c.},`$ (7b)
$`_{\mathrm{\Lambda }^{}NK^{}}^\pm `$ $`=`$ $`i{\displaystyle \frac{g_{\mathrm{\Lambda }^{}NK^{}}}{M_\mathrm{\Lambda }^{}}}\overline{\mathrm{\Lambda }}_{}^{}{}_{\mu }{}^{}\theta ^{\mu \nu }(Y)\gamma ^\lambda F_{\overline{K}}^{}{}_{\lambda \nu }{}^{}N+\text{h.c.},`$ (7c)
where $`\mathrm{\Lambda }^{}`$ is the $`\mathrm{\Lambda }`$(1520) field, $`M_\mathrm{\Lambda }^{}`$ denotes the $`\mathrm{\Lambda }^{}`$ mass, $`F_K^{\mu \nu }`$ is related to the vector $`K^{}`$ meson field as $`F_K^{\mu \nu }=^\nu K_{}^{}{}_{}{}^{\mu }^\mu K_{}^{}{}_{}{}^{\nu }`$. The operator $`\theta _{\mu \nu }(X)`$ is a function of the โoff-shellโ parameter $`X`$: $`\theta _{\mu \nu }(X)=g_{\mu \nu }(\frac{1}{2}+X)\gamma _\mu \gamma _\nu `$. In this paper we consider such a kinematics where the invariant mass of the outgoing $`N\overline{K}`$ pair is close to $`M_\mathrm{\Lambda }^{}`$, $`\mathrm{\Lambda }^{}`$ is almost on-shell, and therefore, the contribution from terms proportional to $`\gamma _\mu \gamma _\nu `$ in $`\theta _{\mu \nu }(X)`$ disappears. This means that $`\theta _{\mu \nu }(X)`$ may be replaced by $`g_{\mu \nu }`$. We assume a vanishing value of the anomalous magnetic moment of $`\mathrm{\Lambda }^{}`$ and, therefore, neglect the $`\mathrm{\Lambda }^{}\gamma `$ interaction, and, correspondingly, the contribution of the $`u`$-channel shown in Fig. 3c. All vertices are dressed by the form factors similarly to the case of the $`\mathrm{\Theta }^+`$ photoproduction with the same cut-off parameters. The amplitudes for the $`\gamma p\mathrm{\Lambda }^{}K^+`$ and $`\gamma n\mathrm{\Lambda }^{}K^0`$ reactions read
$`A_{}^{\mathrm{\Lambda }^{}}{}_{fi}{}^{}(\gamma p)=\overline{u}_\mathrm{\Lambda }^{}^\sigma (p_\mathrm{\Lambda }^{})\left[_{}^{s}{}_{\sigma \mu }{}^{}+_{}^{t}{}_{\sigma \mu }{}^{}+_{}^{c}{}_{\sigma \mu }{}^{}+_{}^{t}{}_{\sigma \mu }{}^{}(K^{})\right]u_p(p)\epsilon ^\mu ,`$ (8a)
$`A_{}^{\mathrm{\Lambda }^{}}{}_{fi}{}^{}(\gamma n)=\overline{u}_\mathrm{\Lambda }^{}^\sigma (p_\mathrm{\Lambda }^{})\left[_{}^{s}{}_{\sigma \mu }{}^{}+_{}^{t}{}_{\sigma \mu }{}^{}(K^{})\right]u_n(p)\epsilon ^\mu .`$ (8b)
The explicit transition operators $`_{\sigma \mu }^i`$ for these reactions are listed in Appendix A.
The coupling constant $`g_{\mathrm{\Lambda }^{}NK}`$ is found from the $`\mathrm{\Lambda }^{}`$ decay width,
$`\mathrm{\Gamma }_{\mathrm{\Lambda }^{}N\overline{K}}={\displaystyle \frac{[g_{\mathrm{\Lambda }^{}NK}]^2p_F^3}{6\pi M_\mathrm{\Lambda }^{}^3}}(\sqrt{M_N^2+p_F^2}M_N),`$ (9)
where $`p_F`$ is $`\mathrm{\Lambda }^{}N\overline{K}`$-decay momentum. Taking $`\mathrm{\Gamma }_{\mathrm{\Lambda }^{}N\overline{K}}0.45\times 15.6`$ MeV PDG , one finds $`|g_{\mathrm{\Lambda }^{}NK}|=32.6`$.
Analog to the above considered $`\mathrm{\Theta }^+`$ photoproduction we denote $`g_{\mathrm{\Lambda }^{}NK^{}}=\alpha _\mathrm{\Lambda }^{}g_{\mathrm{\Lambda }^{}NK}`$. The parameter $`\alpha _\mathrm{\Lambda }^{}`$ must be defined by a comparison of calculated cross sections with experimental data at $`E_\gamma 2`$ GeV. However, the available experimental data for the $`\gamma p\mathrm{\Lambda }^{}K^+`$ reaction cover the energy range $`E_\gamma =2.84.8`$ (GeV) Barber1980 , beyond the applicability of the effective Lagrangian formalism. Thus, in this region the total cross section decreases with energy as $`E_\gamma ^{2.1}`$, whereas the amplitudes of Eq. (8) predict a strong increase. The energy dependence at high energy is reasonably well described by the Regge phenomenology. Since the $`\mathrm{\Lambda }^{}`$ decay angular distribution supports the dominance of the $`t`$-channel natural parity exchange processes, one can assume that the dominant contribution to the $`\mathrm{\Lambda }^{}`$ photoproduction at high energy comes from the leading $`K^{}`$ trajectory Collins . The corresponding amplitude is obtained from the $`t`$-channel $`K^{}`$ meson exchange in Eq. (8) by the Reggezation of the $`K^{}`$ meson exchange propagator, i.e.
$`{\displaystyle \frac{1}{tM_K^{}^2}}\gamma (t)\left({\displaystyle \frac{s}{s_0}}\right)^{\alpha (t)},`$ (10)
where $`\alpha (t)=\alpha (0)+\alpha ^{}t`$ is the Regge trajectory and $`\gamma (t)`$ denotes the normalization function
$`\gamma (t)=C_R(\mathrm{Tr}[RR^{}])^1,`$
$`R=\overline{u}_\mathrm{\Lambda }^{}^\sigma (p_\mathrm{\Lambda }^{})\left[\epsilon ^{\nu \mu \alpha \beta }k^\nu q_{}^{}{}_{}{}^{\alpha }(q_{}^{}{}_{\sigma }{}^{}\gamma ^\beta q^{}/g_\sigma ^\beta )\right]u_n(p)\epsilon ^\mu `$ (11)
with $`q^{}=p_\mathrm{\Lambda }^{}p`$. In the following we assume that at energies near the threshold, the production amplitude is defined by the effective Lagrangian model of Eq. (8), $`A_{eff.L.}^\mathrm{\Lambda }^{}`$, whereas at high energies it is described by the Regge phenomenology, $`A_R^\mathrm{\Lambda }^{}`$, as
$`A^\mathrm{\Lambda }^{}=A_{eff.L.}^\mathrm{\Lambda }^{}\theta (E_0E_\gamma )+A_R^\mathrm{\Lambda }^{}\theta (E_\gamma E_0).`$ (12)
We take $`E_0=2.3`$ GeV as matching point between the two regimes. The choice of parameters in Eq. (11) as $`s_0=1`$ GeV, $`\alpha (t)=0.1+0.9t`$ and $`C_R=29.6`$ gives a satisfactory description of the high energy data, as exhibited in Fig. 4 for the differential cross section at $`E_\gamma =3.7`$ GeV.
In Fig. 5 we show the energy dependence of the total cross section.
The dot-dashed curve is the fit of the data $`\sigma 6.55(E_\gamma /\mathrm{GeV})^{2.1}`$ ($`\mu `$b) from Barber1980 . For illustration we also show the cross section calculated with a constant amplitude where the energy dependence is defined by the phase space volume alone. The strength parameter $`\alpha _\mathrm{\Lambda }^{}`$ is adjusted by fitting the calculated cross section to the experimental extrapolation (dot-dashed curve) at the normalization point. Two solutions $`\alpha _\mathrm{\Lambda }^{}=+0.372`$ and $`0.657`$ result in two different energy dependencies of the cross section at low energy. Both solutions exceed the experimental data above the normalization point. The solution with positive $`\alpha _\mathrm{\Lambda }^{}`$ at low energies is close to the pure phase space dependence shown by the long-dashed curve.
In Fig. 6 we show the differential cross sections of the $`\mathrm{\Lambda }^{}`$ photoproduction at $`E_\gamma =2`$ GeV. The differential cross sections of the $`\gamma p\mathrm{\Lambda }^{}K^+`$ reaction for positive $`\alpha _\mathrm{\Lambda }^{}`$ together with the separate contributions of the Born and $`K^{}`$ exchange channels are shown in Fig. 6a. In case of the $`\gamma n\mathrm{\Lambda }^{}K^0`$ reaction, shown in Fig. 6b by the solid curve, the Born term ($`s`$-channel exchange) is negligible. In the $`\gamma p`$ reaction, the interplay of the Born terms and the $`K^{}`$ exchange amplitude is important at forward angles that leads to a dependence of the total cross section on the sign of $`\alpha _\mathrm{\Lambda }^{}`$ (see Fig. 6b). However, as we will see later, in the coherent $`\gamma D\mathrm{\Lambda }^{}\mathrm{\Theta }^+`$ reaction the region of backward angles of the $`K^+`$ photoproduction gives the main contribution and, therefore, the final result is not sensitive to the choice of the solution. Nevertheless, for further consideration we chose the solution with positive $`\alpha _\mathrm{\Lambda }^{}`$ because it describes better the total $`K^+K^{}`$ production in $`\gamma p`$ interaction at low energies.
Finally we note that a similar approach for the $`\mathrm{\Lambda }^{}`$ photoproduction based on the effective Lagrangian formalism was developed in the recent paper Hosaka0503 . Differences consist in a different choice of the form factors and parameters, which results in slightly different predictions for the differential and total cross sections. This difference may be resolved experimentally.
## III Reaction $`๐ธ๐ซ\mathbf{}๐ฒ^{\mathbf{}}๐ฏ^\mathbf{+}`$
The tree level diagrams for the coherent $`\gamma D\mathrm{\Lambda }^{}\mathrm{\Theta }^+`$ photoproduction are shown in Fig. 7.
First of all note that the amplitudes from the charge and neutral meson exchange shown in Figs. 7a and c and/or b and d give a constructive interference in the total cross section. That is because in the elementary amplitudes of $`\gamma N\mathrm{\Lambda }^{}K`$ and $`\gamma N\mathrm{\Theta }^+\overline{K}`$ reactions the dominant contribution comes from the $`K^{}`$ exchange. The different signs in $`\gamma K^0\overline{K}^0`$ and $`\gamma K^+K^{}`$ vertices are compensated by the different signs in $`n\mathrm{\Theta }^+K^{}`$ and $`p\mathrm{\Theta }^+\overline{K}^0`$ interactions. The latter is a consequence of the assumed isospin $`I=0`$ of the pentaquark.
The amplitudes of the coherent $`\mathrm{\Lambda }^{}\mathrm{\Theta }^+`$ photoproduction are expressed through the transition operators of the โelementaryโ processes $`\gamma N\mathrm{\Lambda }^{}K`$ and $`\gamma N\mathrm{\Theta }^+\overline{K}`$ shown in Fig. 7a,c and b,d, respectively, as
$`A_{(a,c)}`$ $`=`$ $`g_{\mathrm{\Theta }NK}{\displaystyle \frac{d^4p}{(2\pi )^4}\overline{u}_\mathrm{\Theta }\gamma _5\frac{1}{q^2M_K^2}\overline{u}_\mathrm{\Lambda }^{}^\sigma _{\sigma \mu }^\mathrm{\Lambda }^{}\frac{p/+M}{p^2M^2}\mathrm{\Gamma }_D\frac{p/^{}+M}{p_{}^{}{}_{}{}^{2}M^2}U_Dฯต^\mu },`$ (13a)
$`A_{(b,d)}`$ $`=`$ $`{\displaystyle \frac{g_{\mathrm{\Lambda }^{}NK}}{M_\mathrm{\Lambda }^{}}}{\displaystyle \frac{d^4p}{(2\pi )^4}\overline{u}_\mathrm{\Theta }_\mu ^\mathrm{\Theta }\frac{1}{q^2M_K^2}\overline{u}_\mathrm{\Lambda }^{}^\sigma q_\sigma \gamma _5\frac{p/+M}{p^2M^2}\mathrm{\Gamma }_D\frac{p/^{}+M}{p_{}^{}{}_{}{}^{2}M^2}U_Dฯต^\mu },`$ (13b)
where the transition operators $``$ are described in the previous section, $`\mathrm{\Gamma }_D`$ and $`U_D`$ stand for the deuteron $`np`$ coupling vertex and the deuteron spinor, respectively, $`p^{}=p_Dp`$ and $`q`$ is the momentum of the exchanged kaon.
Following Ref. gammaD we assume that the dominant contribution to the loop integrals comes from their imaginary parts which may be evaluated by summing all possible cuttings of the loops, as shown in Fig. 8.
Calculating the imaginary parts we use the following substitutions for the propagators of the on-shell particles (shown by crosses)
$`{\displaystyle \frac{1}{q^2M_K^2}}2\pi \delta (q^2M_K^2),`$
$`{\displaystyle \frac{p/+M}{p^2M^2}}2\pi (p/+M)\delta (p^2M^2)`$ (14)
and the identity
$`{\displaystyle d^4p\delta (p^2M^2)}={\displaystyle \frac{d^3๐ฉ}{2E}}`$ (15)
with $`E^2=๐ฉ^2+M^2`$. We also use the standard representation of the product of the deuteron vertex function and the attached nucleon propagator through the non-relativistic deuteron function
$`\mathrm{\Gamma }_D{\displaystyle \frac{\overline{u}_1(p)\overline{u}_2(p_Dp)}{U_D}}=\sqrt{2M_D}\psi _{m_D,m_1m_2},`$ (16)
where $`\psi _{m_D,m_1m_2}`$ is the deuteron wave function with the spin projection $`m_D`$ and the nucleons spin projections $`m_1`$ and $`m_2`$. By using Eqs. (14) - (16), one can express the principal parts of the invariant amplitudes in Eq. (13) as
$`A_{(a,c)}^P`$ $`=`$ $`g_{\mathrm{\Theta }NK}{\displaystyle \underset{m_1m_2}{}}[\overline{u}_\mathrm{\Theta }(p_\mathrm{\Theta })\gamma _5u_{m_1}(r)][\overline{u}_\mathrm{\Lambda }^{}^\sigma (p_\mathrm{\Lambda }^{})_{\sigma \mu }^\mathrm{\Lambda }^{}ฯต^\mu u_{m_2}(r)]S_{m_1m_2}^\mathrm{\Lambda }^{},`$ (17a)
$`A_{(b,d)}^P`$ $`=`$ $`{\displaystyle \frac{g_{\mathrm{\Lambda }^{}NK}}{M_\mathrm{\Lambda }^{}}}{\displaystyle \underset{m_1m_2}{}}[\overline{u}_\mathrm{\Theta }_\mu ^\mathrm{\Theta }u_{m_1}(r)ฯต^\mu ][\overline{u}_\mathrm{\Lambda }^{}^\sigma q_\sigma \gamma _5u_{m_2}(r)]S_{m_1m_2}^{\mathrm{\Theta }^+},`$ (17b)
where $`r=p_D/2`$, and
$`S_{m_1m_2}^\mathrm{\Lambda }^{}`$ $`=`$ $`I_{m_1m_2}^i(p_\mathrm{\Theta })+I_{m_1m_2}^j(kp_\mathrm{\Lambda }^{}),S_{m_1m_2}^{\mathrm{\Theta }^+}=I_{m_1m_2}^i(p_\mathrm{\Lambda }^{})+I_{m_1m_2}^j(kp_\mathrm{\Theta }),`$
$`I_{m_1m_2}^{i,j}(p_X)`$ $`=i`$ $`{\displaystyle \frac{\sqrt{2M_D}}{16\pi }}{\displaystyle \frac{pdp}{Ep_X}\theta (1|a_{i,j}(p,p_X)|)\varphi _{m_D,m_1m_2}(p,a(p,p_X))},`$
$`a_i(p,p_X)`$ $`=`$ $`{\displaystyle \frac{2EE_X+M_K^2M_X^2M^2}{2pp_X}},`$
$`a_j(p,p_X)`$ $`=`$ $`{\displaystyle \frac{2EE_XM_K^2+M_X^2+M^2}{2pp_X}},`$
$`\varphi _{m_D,m_1m_2}(p,a)`$ $`=`$ $`\sqrt{4\pi }{\displaystyle \frac{1}{2}}m_1{\displaystyle \frac{1}{2}}m_2|1m_D\left(u_0(p)+{\displaystyle \frac{1}{\sqrt{8}}}(3a^21)(13\delta _{m_D0})u_2(p)\right),`$ (18)
where $`M_X^2=E_X^2p_X^2`$ and $`u_l`$ with $`l=0,2`$ is the radial deuteron wave function in the momentum space, normalized as
$`{\displaystyle \frac{d๐ฉ}{(2\pi )^3}\mathrm{\Phi }(p)}=1,`$
where
$`\mathrm{\Phi }(p)=4\pi \left(u_0^2(p)+u_2^2(p)\right).`$ (19)
In deriving Eqs. (17) we neglect the weak dependence of the โelementaryโ amplitudes of $`\gamma N\mathrm{\Lambda }^{}K`$ and $`\gamma N\mathrm{\Theta }^+\overline{K}`$ on $`p`$ (see Figs. 2 and 4), compared to the sharp $`p`$ dependence of $`\mathrm{\Phi }(p)`$. In our calculation we use the deuteron wave function for the โrealisticโ Paris potential Paris . We checked that the final result does not depend on the fine structure of the deuteron wave function and practically does not depend on the choice of the potential.
The differential cross section of the coherent $`\mathrm{\Lambda }^{}\mathrm{\Theta }^+`$ photoproduction reads
$`{\displaystyle \frac{d\sigma ^{\gamma D\mathrm{\Lambda }^{}\mathrm{\Theta }^+}}{d\mathrm{\Omega }}}={\displaystyle \frac{1}{64\pi ^2}}{\displaystyle \frac{1}{S}}{\displaystyle \frac{P_{\mathrm{out}}}{P_{\mathrm{in}}}}|A_{a,c}+A_{b,d}|^2,`$ (20)
where $`S,P_{\mathrm{in}}`$ and $`P_{\mathrm{out}}`$ are the square of the total energy, momenta in initial and the final states in $`\gamma D`$ c.m.s., respectively; averaging and summing over the spin projections in the initial and the final states are assumed. Note that the interference between amplitudes $`A_{a,c}`$ and $`A_{b,d}`$ is negligible and they can be summed incoherently.
In Fig. 9 we show the differential cross section of the reaction $`\gamma D\mathrm{\Lambda }^{}\mathrm{\Theta }^+`$ at $`E_\gamma =2`$ GeV as a function of the angle between the beam direction and direction of flight of $`\mathrm{\Lambda }^{}`$ in the $`\gamma D`$ c.m.s. The non-monotonous behaviour of the cross section is completely defined by the spectral functions $`S^\mathrm{\Lambda }^{}`$ and $`S^{\mathrm{\Theta }^+}`$ in Eqs. (17a) and (17b), respectively. The spectral functions $`S^\mathrm{\Lambda }^{}`$ and $`S^{\mathrm{\Theta }^+}`$ have sharp peaks in forward ($`\theta _{\gamma \mathrm{\Lambda }^{}}27.5^o`$) and backward ($`\theta _{\gamma \mathrm{\Lambda }^{}}152.5^o`$ ) hemispheres, respectively.
## IV Background contribution
Since the $`\mathrm{\Lambda }^{}`$ and $`\mathrm{\Theta }^+`$ are unstable baryons, the typical experiment for studying the coherent $`\gamma D\mathrm{\Lambda }^{}\mathrm{\Theta }^+`$ process must include a simultaneous measurement of the $`pK^{}`$ and $`nK^+`$ invariant masses. Therefore, the question is whether the predicted cross section of the coherent $`\mathrm{\Lambda }^{}\mathrm{\Theta }^+`$ photoproduction is large enough to be seen above the background of competing resonance and non-resonance processes in the $`\gamma DnpK^+K^{}`$ reaction.
We consider three types of background processes. One is the photoproduction of a $`K^+K^{}`$ pair in a $`\gamma p`$ interaction when the neutron is a spectator. This process includes the resonant $`\gamma p\mathrm{\Lambda }^{}K^+pK^+K^{}`$ photoproduction and the non-resonant $`\gamma ppK^+K^{}`$ reaction shown in Fig. 10a and b, respectively.
Similarly, a $`K^+K^{}`$ pair can be produced in a $`\gamma n`$ interaction, when the proton is a spectator. The corresponding processes are depicted in Fig. 10c and d.
The third process is the coherent โbackgroundโ when the $`K^+K^{}`$ pair is produced in a $`\gamma N`$ interaction and one of the kaons together with the second nucleon forms the outgoing $`\mathrm{\Theta }^+`$ or $`\mathrm{\Lambda }^{}`$, as shown in Fig. 10 e and f, respectively. We denote it as a coherent semi-resonant background.
### IV.1 Spectator channels
First, let us consider the $`K^+K^{}`$ photoproduction in a $`\gamma D`$ interaction where the neutron or proton are merely spectators. As an input, we have to describe the elementary processes $`\gamma ppK^+K^{}`$ and $`\gamma nnK^+K^{}`$ which consist of the resonant and non-resonant parts.
#### IV.1.1 $`๐ธ๐\mathbf{}๐๐ฒ^\mathbf{+}๐ฒ^{\mathbf{}}`$
The dominant contribution to the non-resonant part in $`\gamma p`$ reactions comes from the virtual vector meson decay and $`\mathrm{\Lambda }(1405)`$ excitation TEHN04 ; Oh:2004wp as depicted in Fig. 11a and b. The contribution from excitations of other hyperons is strongly suppressed since they are far off-shell.
The vector meson channel $`\gamma pVppK^+K^{}`$, where $`V=\varphi ,\rho ,\omega `$ has been analyzed in detail in Ref. TEHN04 . In the present study we use this model where the vector mesons are produced through the Pomeron and meson ($`\pi ,\eta ,\sigma `$) exchanges with the same parameters. The only difference with Ref. TEHN04 is that now we do not use a cut on the invariant mass of the $`K^+K^{}`$ pair around the $`\varphi `$ meson mass.
We parameterize the amplitude of the virtual $`\mathrm{\Lambda }(1405)`$ excitation through the $`K^{}`$ exchange process. This assumption is supported by the $`K^{}`$ exchange dominance in $`\mathrm{\Lambda }^{}`$ and $`\mathrm{\Theta }^+`$ photoproduction and allows to reduce the number of unknown parameters. The amplitude of this channel reads
$`A_{fi}^\mathrm{\Lambda }^{}`$ $`=`$ $`\overline{u}(p^{})_\mu ^\mathrm{\Lambda }^{}u(p)\epsilon ^\mu ,`$
$`_\mu ^\mathrm{\Lambda }^{}`$ $`=`$ $`i{\displaystyle \frac{eg_{\gamma KK^{}}g^{}}{M_K^{}(tM_K^{}^2)}}\epsilon ^{\mu \nu \alpha \beta }k_\nu q_\alpha {\displaystyle \frac{(p/_{\mathrm{\Lambda }^{}+M_\mathrm{\Lambda }^{}})\gamma _5\gamma _\beta }{p_\mathrm{\Lambda }^{}^2M_\mathrm{\Lambda }^{}^2+i\mathrm{\Gamma }_\mathrm{\Lambda }^{}M_\mathrm{\Lambda }^{}}}F_K^{}(t),`$ (21)
where $`\mathrm{\Lambda }^{}\mathrm{\Lambda }(1405)`$, $`\mathrm{\Gamma }_\mathrm{\Lambda }^{}=50`$ MeV is the total decay width of $`\mathrm{\Lambda }^{}`$ PDG , $`F_K^{}(t)`$ is the $`K^{}`$ exchange form factor, the constant $`g^{}`$ is a product of two coupling constants $`g_{\mathrm{\Lambda }^{}NK}`$ and $`g_{\mathrm{\Lambda }^{}NK^{}}`$. The choice $`g^{}7.8`$ gives the correct value of the total yield of $`K^+K^{}`$ mesons at $`E_\gamma 2`$ GeV. Note that the interference between the resonance and non-resonance channels in the total cross section is rather weak and, therefore, they can be added incoherently. Thus, the total cross section of the $`\gamma ppK^+K^{}`$ reaction reads
$`{\displaystyle \frac{d\sigma }{d\mathrm{\Omega }dM_{pK^{}}}}`$ $`=`$ $`\left({\displaystyle \frac{d\sigma }{d\mathrm{\Omega }}}\right)^{\gamma p\mathrm{\Lambda }^{}K^+}F^\mathrm{\Lambda }^{}(M_{pK^{}})`$ (22)
$`+`$ $`{\displaystyle \frac{1}{64\pi ^2}}{\displaystyle \frac{1}{s}}{\displaystyle \frac{p_{\mathrm{out}}}{p_{\mathrm{in}}}}{\displaystyle \frac{\overline{q}_F}{16\pi ^3}}{\displaystyle \left(|A_{fi}^V(\gamma p)|^2+|A_{fi}^\mathrm{\Lambda }^{}|^2\right)๐\mathrm{\Omega }_F},`$
where $`\mathrm{\Omega }`$ is the solid angle of the $`K^{}`$ meson photoproduction in the $`\gamma p`$ c.m.s., $`\overline{q}_F`$ is the momentum of the $`K^{}`$ meson in the c.m.s. of the $`pK^{}`$-pair, $`\mathrm{\Omega }_F`$ is the $`K^{}`$ meson solid angle in this system. Summing and averaging over the spin projection in the initial and the final states is to be included. $`F^\mathrm{\Lambda }^{}(M_{pK^{}})`$ stands for the $`\mathrm{\Lambda }^{}`$ decay distribution which is obtained straightforwardly from the general expression of the $`\gamma ppK^+K^{}`$ amplitude with the virtual excitation of a $`\mathrm{\Lambda }^{}`$ hyperon,
$`F^\mathrm{\Lambda }^{}(M_x)={\displaystyle \frac{\mathrm{\Gamma }_{\mathrm{\Lambda }^{}pK^{}}}{\pi \mathrm{\Gamma }_{\mathrm{tot}}}}{\displaystyle \frac{2M_xM_\mathrm{\Lambda }^{}\mathrm{\Gamma }_{\mathrm{tot}}}{(M_x^2M_\mathrm{\Lambda }^{}^2)^2+(\mathrm{\Gamma }_{\mathrm{tot}}M_\mathrm{\Lambda }^{})^2}},`$ (23)
where $`\mathrm{\Gamma }_{\mathrm{tot}}=15.6`$ MeV and $`\mathrm{\Gamma }_{\mathrm{\Lambda }^{}pK^{}}=(0.45/2)\times \mathrm{\Gamma }_{\mathrm{tot}}`$ PDG .
The $`pK^{}`$ invariant mass distribution at $`E_\gamma =2`$ GeV integrated over $`\mathrm{\Omega }`$ is shown in Fig. 12. One can see that the $`\mathrm{\Lambda }(1405)`$ excitation contributes at $`M_{pK^{}}`$ below the $`\mathrm{\Lambda }^{}`$ resonance position, and the vector meson channels contribute mainly at large $`M_{pK^{}}`$, above $`M_\mathrm{\Lambda }^{}`$. The partial contributions to the total $`\gamma ppK^+K^{}`$ cross section are the following: $`\sigma (\mathrm{\Lambda }^{})0.19\mu `$b, $`\sigma (V)0.17\mu `$b and $`\sigma (\mathrm{\Lambda }(1405))0.07\mu `$b. The total cross section $`\sigma _{\mathrm{tot}}0.43\mu `$b is in agreement with the experimental data of Ref. exp-gp-KK : $`\sigma _{\mathrm{tot}}^{\mathrm{exp}}=(0.47\pm 0.12)\mu `$b at $`E_\gamma =22.5`$ (GeV).
#### IV.1.2 $`๐ธ๐\mathbf{}๐๐ฒ^\mathbf{+}๐ฒ^{\mathbf{}}`$
In this case the non-resonance part is dominated by the vector meson excitation and, therefore, the $`nK^+`$ invariant mass distribution may be written in obvious notation as
$`{\displaystyle \frac{d\sigma }{d\mathrm{\Omega }dM_{nK^+}}}`$ $`=`$ $`\left({\displaystyle \frac{d\sigma }{d\mathrm{\Omega }}}\right)^{\gamma n\mathrm{\Theta }^+K^{}}F^{\mathrm{\Theta }^+}(M_{nK^+})+{\displaystyle \frac{1}{64\pi ^2}}{\displaystyle \frac{1}{s}}{\displaystyle \frac{p_{\mathrm{out}}}{p_{\mathrm{in}}}}{\displaystyle \frac{q_F}{16\pi ^3}}{\displaystyle |A_{fi}^V(\gamma n)|^2๐\mathrm{\Omega }_F}`$ (24)
with
$`F^\mathrm{\Theta }(M_x)={\displaystyle \frac{1}{2\pi }}{\displaystyle \frac{2M_xM_\mathrm{\Theta }\mathrm{\Gamma }_\mathrm{\Theta }}{(M_x^2M_\mathrm{\Theta }^2)^2+(\mathrm{\Gamma }_{\mathrm{tot}}M_\mathrm{\Theta })^2}}.`$ (25)
We will also use the Gaussian distribution taking into account the small $`\mathrm{\Theta }^+`$ decay width and the finite experimental resolution
$`F_G^\mathrm{\Theta }(M_x)={\displaystyle \frac{1}{2}}{\displaystyle \frac{1}{\sigma \sqrt{2\pi }}}e^{\frac{(M_xM_\mathrm{\Theta })^2}{2\sigma ^2}}.`$ (26)
The $`nK^+`$ invariant mass distribution at $`E_\gamma =2`$ GeV integrated over $`\mathrm{\Omega }`$ is shown in Fig. 13. One can see the sharp peak of $`\mathrm{\Theta }^+`$ excitation. In case of a Gaussian $`\mathrm{\Theta }^+`$ decay distribution the peak is modified. The height of the peak is reduced by the factor $`\sigma /\mathrm{\Gamma }_\mathrm{\Theta }`$ and the width becomes proportional to $`\sigma `$.
#### IV.1.3 Spectator reactions $`๐ธ๐ซ\mathbf{}๐๐ฒ^\mathbf{+}๐ฒ^{\mathbf{}}\mathbf{(}๐\mathbf{)}`$ and $`๐ธ๐ซ\mathbf{}๐๐ฒ^\mathbf{+}๐ฒ^{\mathbf{}}\mathbf{(}๐\mathbf{)}`$
The differential cross section of the $`\gamma DpK^+K^{}(n)`$ reaction, where the neutron is a spectator, reads
$`{\displaystyle \frac{d\sigma ^{\mathrm{sp}.(\mathrm{n})}}{d\mathrm{\Omega }dM_{pK^{}}dM_{nK+}}}=\left({\displaystyle \frac{d\sigma }{d\mathrm{\Omega }dM_{pK^{}}}}\right)^{\gamma ppK^+K^{}}W_{nK}(M_{nK^+}),`$
$`W_{nK}(M_{nK^+})=2M_{nK^+}{\displaystyle \frac{d๐ฉ_n}{(2\pi )^3\sqrt{1+๐ฉ_n^2/M_N^2}}\delta (M_{nK^+}^2(p_n+q)^2)\mathrm{\Phi }(๐ฉ_n)},`$ (27)
where we neglect the smooth dependence of $`d\sigma ^{\gamma ppK^+K^{}}`$ on $`๐ฉ_n`$ in comparison to the sharp $`๐ฉ_n`$ dependence of the momentum distribution in the deuteron, $`\mathrm{\Phi }(๐ฉ_n)`$, defined in Eq. (19).
If the invariant mass of the $`nK^+`$ pair is not fixed then the integration over $`M_{nK+}`$ leads to the obvious result
$`{\displaystyle ๐M_{nK^+}\frac{d\sigma ^{\mathrm{sp}.(\mathrm{n})}}{d\mathrm{\Omega }dM_{pK^{}}dM_{nK+}}}\left({\displaystyle \frac{d\sigma }{d\mathrm{\Omega }dM_{pK^{}}}}\right)^{\gamma ppK^+K^{}}.`$ (28)
When the invariant mass is fixed then the function $`W_{nK}(M_{nK^+})`$ becomes important and, moreover, it mainly defines the dependence of the cross section on $`M_{nK^+}`$. Indeed, let us assume that the momentum distribution in a deuteron behaves like a delta function, i.e. $`\mathrm{\Phi }(๐ฉ)(2\pi )^3\delta (๐ฉ)`$. Then one gets
$`W_{nK}(M_{nK^+})\mathrm{\hspace{0.17em}2}M_{nK^+}\delta (M_{nK^+}^2(M_N^2+M_K^2+2E_{K^+}M_N)).`$ (29)
That is, the distribution $`W_{nK}(M_{nK^+})`$ has a peak around the point $`M_{nK^+}^{}{}_{0}{}^{}\sqrt{M_N^2+M_K^2+2E_{K^+}M_N}`$ which is determined by the energy of the $`K^+`$ meson in the laboratory system. On the other hand, this energy depends on the invariant mass of the $`pK^{}`$ pair and the angle of the $`K^+`$ production in the $`\gamma p`$ c.m.s. In reality, the distribution function reads
$`W_{nK}(M_{nK^+})`$ $`=`$ $`2M_{nK^+}{\displaystyle \frac{pdp}{8\pi ^2q_L\sqrt{1+p^2/M_N^2}}\mathrm{\Phi }(p)\theta (1|a|)},`$
$`a`$ $`=`$ $`{\displaystyle \frac{2\sqrt{(q_L^2+M_K^2)(p^2+M_N^2)}+M_N^2+M_K^2M_{nK^+}^2}{2pq_L}},`$ (30)
where $`q_L`$ is the momentum of $`K^+`$ meson in laboratory system. The distribution function $`W_{nK}`$ is shown in Fig. 14a as a function of $`M_{nK^+}`$ at fixed angle of $`pK^{}`$ pair photoproduction, $`\theta _{\gamma (pK^{})}`$ (in $`\gamma D`$ c.m.s.) for three different invariant masses of the $`pK^{}`$ pair: $`M_{pK^{}}=`$1.52, 1.57 and 1.47 GeV. The choice of $`\theta _{\gamma (pK^{})}=27.5^o`$ corresponds to the position of the maximum of the coherent $`\gamma D\mathrm{\Lambda }^{}\mathrm{\Theta }^+`$ photoproduction cross section at forward angles (see Fig. 9). This angle corresponds to the backward $`K^+`$ photoproduction in $`\gamma p\mathrm{\Lambda }^{}K^+`$: $`\theta _{\gamma K^+}119^o`$ in the $`\gamma p`$ c.m.s.
The differential cross section of the $`\gamma DnK^+K^{}(p)`$ reaction, where the proton is spectator, may be obtained from Eq. (27), using the substitution $`np`$, $`K^+K^{}`$ and $`M_{nK^+}M_{pK^{}}`$,
$`{\displaystyle \frac{d\sigma ^{\mathrm{sp}.(\mathrm{p})}}{d\mathrm{\Omega }dM_{pK^{}}dM_{nK+}}}=\left({\displaystyle \frac{d\sigma }{d\mathrm{\Omega }dM_{nK^+}}}\right)^{\gamma nnK^+K^{}}W_{pK}(M_{pK^{}}).`$ (31)
The essential difference is that now we analyze the dependence of the distribution function $`W_{pK}`$ not on $`M_{pK^{}}`$ but on the invariant mass $`M_{nK^+}`$. This dependence is included in $`W_{pK}`$ implicitly through the dependence of the momentum of $`K^{}`$ on $`M_{nK^+}`$ and therefore, in general, we have no narrow peak structure of $`W_{pK}`$ as a function of $`M_{nK^+}`$. As an example, in Fig. 14b we show the distribution $`W_{pK}`$ as a function of $`M_{nK^+}`$ at fixed values of $`M_{pK^{}}=`$1.52, 1.57 and 1.47 GeV and $`\theta _{\gamma (pK^{})}=152.5^o`$. One can see a broad maximum at $`M_{pK^{}}=`$1.52 GeV and an almost monotonic behaviour at 1.47 and 1.57 GeV.
### IV.2 Coherent semi-resonant background
The amplitude of the process shown in Fig. 10e is calculated similarly to the amplitude of the coherent $`\mathrm{\Lambda }^{}\mathrm{\Theta }^+`$ photoproduction described by Eq. (17a). The corresponding cross section reads
$`{\displaystyle \frac{d\sigma ^e}{d\mathrm{\Omega }dM_{pK^{}}dM_{nK+}}}={\displaystyle \frac{1}{64\pi ^2}}{\displaystyle \frac{1}{s}}{\displaystyle \frac{p_{\mathrm{out}}}{p_{\mathrm{in}}}}{\displaystyle \frac{\overline{q}_F}{16\pi ^3}}{\displaystyle \frac{1}{2}}\left|{\displaystyle \frac{d\mathrm{\Omega }^{}}{d\mathrm{\Omega }}}\right|{\displaystyle ๐\mathrm{\Omega }_F|A_e|^2F^\mathrm{\Theta }(M_{nK^+})},`$
$`A_e=g_{\mathrm{\Theta }NK}{\displaystyle \underset{m_1m_2}{}}[\overline{u}_\mathrm{\Theta }(p_\mathrm{\Theta })\gamma _5u_{m_1}(r)][\overline{u}_\mathrm{\Lambda }^{}^\sigma (p_\mathrm{\Lambda }^{})_{\sigma \mu }^{\gamma ppK^+K^{}}\epsilon ^\mu u_{m_2}(r)]S_{m_1m_2}^\mathrm{\Lambda }^{},`$ (32)
where $`p_{\mathrm{in}}`$, $`p_{\mathrm{out}}`$ are the momenta of the proton and $`pK^{}`$ pair in $`\gamma p`$ c.m.s., $`\mathrm{\Omega }`$ and $`\mathrm{\Omega }^{}`$ are the solid angles of the $`pK^{}`$ pair in $`\gamma D`$ and $`\gamma p`$ reactions, respectively, $`\overline{q}_F`$ is the momentum of $`K^{}`$ meson in the rest frame of the $`pK^{}`$ pair, $`\mathrm{\Omega }_F`$ is the solid angle of $`K^{}`$ in this frame. The additional factor $`1/2`$ assumes renormalization of the flux in the $`\gamma D`$ system compared to the $`\gamma p`$ interaction. The function $`F^\mathrm{\Theta }(M_{nK^+})`$ is defined in Eq. (25). Averaging and summing over the spin projections in initial and the final states, respectively, have to be performed. Actually, here we have a sum of two cross sections. One is the contribution of the virtual vector meson and another one is the contribution of the virtual $`\mathrm{\Lambda }(1405)`$ excitation.
Similarly, one can write the cross section of the process shown in Fig. 10f as
$`{\displaystyle \frac{d\sigma ^f}{d\mathrm{\Omega }dM_{pK^{}}dM_{nK+}}}={\displaystyle \frac{1}{64\pi ^2}}{\displaystyle \frac{1}{s}}{\displaystyle \frac{p_{\mathrm{out}}}{p_{\mathrm{in}}}}{\displaystyle \frac{q_F}{16\pi ^3}}{\displaystyle \frac{1}{2}}\left|{\displaystyle \frac{d\mathrm{\Omega }^{}}{d\mathrm{\Omega }}}\right|{\displaystyle ๐\mathrm{\Omega }_F|A_f|^2F^\mathrm{\Lambda }^{}(M_{pK^{}})},`$
$`A_f={\displaystyle \frac{g_{\mathrm{\Lambda }^{}NK}}{M_\mathrm{\Lambda }^{}}}{\displaystyle \underset{m_1m_2}{}}[\overline{u}_\mathrm{\Theta }_\mu ^{\gamma nnK^+K^{}}u_{m_1}(r)\epsilon ^\mu ][\overline{u}_\mathrm{\Lambda }^{}^\sigma q_\sigma \gamma _5u_{m_2}(r)]S_{m_1m_2}^{\mathrm{\Theta }^+},`$ (33)
where the function $`F^\mathrm{\Lambda }^{}(M_{pK^{}})`$ is defined in Eq. (23) and other notations are similar to the previous case.
Let us now compare the contribution of the coherent $`\mathrm{\Lambda }^{}\mathrm{\Theta }^+`$ photoproduction and the coherent semi-resonant background described by Eqs. (32) and (33) in the vicinity of the $`\mathrm{\Theta }^+`$ and $`\mathrm{\Lambda }^{}`$ resonance position
$`{\displaystyle \frac{d\stackrel{~}{\sigma }^{ch.}}{d\mathrm{\Omega }}}={\displaystyle \underset{M_\mathrm{\Lambda }^{}\mathrm{\Delta }}{\overset{M_\mathrm{\Lambda }^{}+\mathrm{\Delta }}{}}}{\displaystyle \underset{M_\mathrm{\Theta }\mathrm{\Delta }}{\overset{M_\mathrm{\Theta }+\mathrm{\Delta }}{}}}๐M_{pK^{}}๐M_{nK^+}{\displaystyle \frac{d\sigma ^{\gamma D\mathrm{\Lambda }^{}\mathrm{\Theta }^+}}{d\mathrm{\Omega }}}F_\mathrm{\Lambda }^{}(M_{pK^{}})F_{\mathrm{\Theta }^+}(M_{nK^+}),`$
$`{\displaystyle \frac{d\stackrel{~}{\sigma }^{ch.bg.}}{d\mathrm{\Omega }}}={\displaystyle \underset{M_\mathrm{\Lambda }^{}\mathrm{\Delta }}{\overset{M_\mathrm{\Lambda }^{}+\mathrm{\Delta }}{}}}{\displaystyle \underset{M_\mathrm{\Theta }\mathrm{\Delta }}{\overset{M_\mathrm{\Theta }+\mathrm{\Delta }}{}}}๐M_{pK^{}}๐M_{nK^+}\left({\displaystyle \frac{d\sigma ^e}{d\mathrm{\Omega }dM_{pK^{}}dM_{nK+}}}+{\displaystyle \frac{d\sigma ^f}{d\mathrm{\Omega }dM_{pK^{}}dM_{nK+}}}\right),`$ (34)
where $`\mathrm{\Delta }=20`$ MeV. In Fig. 15 we show result of such a comparison. One can see that the coherent background contribution has local maxima caused by the spectral functions $`S`$, but the values of these contributions at the peak positions are much smaller compared to the coherent process. Therefore, the dominant background contribution comes from the spectator processes.
## V Results and discussion
As pointed out above, the coherent $`\mathrm{\Lambda }^{}\mathrm{\Theta }^+`$ photoproduction seems to be accessible most effectively by a search for a sharp $`\mathrm{\Theta }^+`$ peak in the invariant $`nK^+`$ mass distribution at fixed invariant masses of the $`pK^{}`$ pair
$`{\displaystyle \frac{d\sigma ^{\gamma DnpK^+K^{}}(M_0)}{d\mathrm{\Omega }dM_{nK^+}}}={\displaystyle \underset{M_0\mathrm{\Delta }}{\overset{M_0+\mathrm{\Delta }}{}}}๐M_{pK^{}}{\displaystyle \frac{d\sigma ^{\gamma DnpK^+K^{}}}{d\mathrm{\Omega }dM_{nK^+}dM_{pK^{}}}}.`$ (35)
In our further analysis we choose $`M_0=1.52,\mathrm{\hspace{0.17em}1.57}`$, and $`1.47`$ GeV and $`\mathrm{\Delta }=20`$ MeV. One can expect that the coherent photoproduction appears at $`M_0=M_\mathrm{\Lambda }^{}=1.52`$ GeV and it is suppressed relative to the strong background when we go above or below this point. Since the cross section of the coherent photoproduction at $`E_\gamma =2`$ GeV has bumps at $`\theta _{\gamma \mathrm{\Lambda }^{}}27.5^o`$ and $`152.5^o`$ in $`\gamma D`$ c.m.s. (see Fig. 9), then it is natural to expect that the regions around these angles are more favored for a manifestation of the coherence effect.
Note that at forward and backward angles of the $`pK^{}`$ pair photoproduction, $`\theta _{\gamma (pK^{})}`$, some of spectator processes shown in Fig. 10 are suppressed dynamically. To illustrate this point let us consider the dependence of $`\mathrm{cos}\theta _{\gamma K^{}}`$ in $`\gamma p\mathrm{\Lambda }^{}K^+`$ photoproduction and $`\mathrm{cos}\theta _{\gamma K^+}`$ in $`\gamma n\mathrm{\Theta }^+K^{}`$ photoproduction as a function of $`\mathrm{cos}\theta _{\gamma \mathrm{\Lambda }^{}}`$. Here we assume that $`\theta _{\gamma K}`$ is the $`K`$ meson photoproduction angle in the $`\gamma N`$ c.m.s. and $`\theta _{\gamma \mathrm{\Lambda }^{}}`$ is the $`\mathrm{\Lambda }^{}`$ photoproduction angle in $`\gamma D`$ c.m.s. One can see that the region of $`0\theta _{\gamma \mathrm{\Lambda }^{}}76^o`$ is forbidden kinematically for $`\mathrm{\Theta }^+K^{}`$ photoproduction from the resting neutron. Similarly, the region of $`107^o\theta _{\gamma \mathrm{\Lambda }^{}}\pi `$ is forbidden for $`\mathrm{\Lambda }^{}K^+`$ photoproduction from the resting proton. In the kinematically forbidden regions the corresponding processes can be proceeded only through the high-momentum component in the deuteron wave function and, therefore, are exponentially small.
Consider first $`\gamma DnpK^+K^{}`$ photoproduction at a forward angle of the $`pK^{}`$ pair at $`\theta _{\gamma (pK^{})}27.5^o`$ and $`E_\gamma =2`$ GeV. The corresponding invariant mass distributions for $`M_0=1.52`$, 1.57 and 1.47 are shown in Fig. 17a, b and c, respectively.
At $`M_0=M_\mathrm{\Lambda }^{}`$, the background is dominated by the resonant $`\mathrm{\Lambda }^{}`$ photoproduction in the spectator mechanism shown in Fig. 10a. The next important contribution comes from the non-resonant spectator channel (Fig. 10b). The shape of the background spectrum has a resonance like behavior with the center close to the mass of $`\mathrm{\Theta }^+`$ and a width of about 15 MeV. This behaviour is defined by the spectral distribution function $`W_{nK}`$ (or the deuteron momentum distribution) in Eq. (27) and the kinematics (see Fig. 14a). At $`M_{pK^{}}=1.52`$ GeV, $`W_{nK}`$ has a sharp peak at $`M_{nK^+}1.54`$ GeV. For $`M_{pK^{}}=`$1.57 and 1.47 GeV the peak position is shifted to lower or higher masses, respectively. Similarly, one can see the corresponding shift in the background contribution at $`M_0=1.57`$ and 1.47 GeV, shown in Figs. 17b and c. Here, the background is dominated by the non-resonant spectator channels. Its value is almost similar for all considered values of $`M_0`$ being much smaller than the total background at $`M_0=1.52`$ GeV.
At $`M_0=1.52`$ GeV, the height of the peak of the coherent $`\mathrm{\Lambda }^{}\mathrm{\Theta }^+`$ channel is about one third of the total background contribution. This ratio decreases for $`M_0=M_\mathrm{\Lambda }^{}\pm 70`$ MeV. Thus, a summary plot of the total $`nK^+`$ invariant mass distribution for three fixed intervals of the $`pK^{}`$ invariant mass is shown in Fig. 18.
One can conclude that, since the width of the coherent photoproduction is much smaller than the effective width of the background, this contribution can be extracted experimentally under the condition of a high resolution measurement of the $`nK^+`$ invariant mass.
In case of a energy resolution comparable to the width of the background peak one has to smear this peak.
The simplest way to do it is integrating the $`nK^+`$ invariant mass distribution over $`\mathrm{\Omega }`$ in the forward hemisphere of the $`pK^{}`$ pair photoproduction. The corresponding predictions for $`M_0=1.52`$ GeV and a summary plot for three values of $`M_0`$ are shown in Figs. 19a and b, respectively. One can see that again at $`M_0=M_\mathrm{\Lambda }^{}`$ the background is dominated by the resonance $`\mathrm{\Lambda }^{}`$ photoproduction where the neutron is a spectator. But the shape of the background is quite different from the previous case. Instead of the narrow peak one observes a monotonous increase of the background contribution. This behavior allows to extract the sharp $`\mathrm{\Theta }^+`$ peak of the coherent $`\mathrm{\Lambda }^{}\mathrm{\Theta }^+`$ photoproduction. The peak becomes negligible at $`M_0=M_\mathrm{\Lambda }^{}\pm 70`$ MeV, as shown in Fig. 19b. Here one can also see the prediction for a Gaussian smearing of the $`\mathrm{\Theta }^+`$ peak with $`\sigma =5`$ MeV.
Consider now the backward hemisphere of the $`pK^{}`$ pair photoproduction in the reaction $`\gamma DnpK^+K^{}`$, say for $`\theta _{\gamma (pK^{})}152.5^o`$. The corresponding invariant mass distributions at different $`M_0`$ are exhibited in Fig. 20.
Now, the dominant contribution to the background comes from the spectator resonant $`\mathrm{\Theta }^+`$ photoproduction, depicted in Fig. 10c. The other channels are rather weak. At $`M_0=1.52`$ GeV the background contribution is enhanced by the distribution function $`W_{pK}`$ which at $`M_{nK^+}1.54`$ GeV is much greater for $`M_0M_\mathrm{\Lambda }^{}`$ (see Fig. 14b). The coherent contribution of the $`\mathrm{\Lambda }^{}\mathrm{\Theta }^+`$ photoproduction is a factor of four smaller than the background contribution.
The summary plot of the total invariant mass distribution of the $`nK^+`$ for three fixed intervals of the $`pK^{}`$ invariant mass is displayed in Fig. 21. One can see a strong increase of the invariant mass distribution at $`M_0=1.52`$ GeV. But this increase is caused mainly by the properties of the distribution function $`W_{pK}`$. Here, we have no striking qualitative effect of the coherent $`\mathrm{\Lambda }^{}\mathrm{\Theta }^+`$ photoproduction. Therefore, studying the coherent $`\mathrm{\Lambda }\mathrm{\Theta }^+`$ photoproduction seems to be difficult in this kinematical region.
Now we would like to make three comments. First, since in the forward hemisphere of the $`\mathrm{\Lambda }^{}`$ photoproduction the dominant contribution comes from the backward angles of the $`K^+`$ photoproduction in the elementary $`\gamma p\mathrm{\Lambda }^{}\overline{K}^+`$ subprocess, our predictions are not sensitive to the choice of the solution for the coupling strength $`\alpha _\mathrm{\Lambda }^{}`$ discussed in Sec. II (see Fig. 6 b).
Second, in our analysis we have assumed that the $`\mathrm{\Theta }^+`$ photoproduction from the nucleon is dominated by the $`t`$-channel $`K^{}`$ exchange process. This assumption leads to a similarity of the $`\mathrm{\Theta }^+`$ photoproduction from the neutron and proton. A violation of this similarity (or a suppression of the photoproduction from the proton, with keeping the cross section of the $`\gamma n\mathrm{\Theta }^+K^{}`$ on the same level) discussed recently Hosaka0505 ; Vita would result in a suppression of the process shown in Fig. 7d. As a consequence, the coherent cross section of $`\mathrm{\Lambda }^{}\mathrm{\Theta }^+`$ photoproduction would be suppressed around the second peak at backward angles of the $`pK^{}`$ pair photoproduction, shown in Figs. 9 and 15 leaving the first peak at forward angles photoproduction without change. The corresponding calculation of the differential cross section with and without contribution of the $`\gamma p\mathrm{\Theta }^+\overline{K}^0`$ subprocess is presented in Fig. 22. Since the coherent $`\mathrm{\Lambda }^{}\mathrm{\Theta }^+`$ photoproduction is determined by the first peak, our main result shown in Fig. 19 remains unchanged.
Third, the โbump-likeโ structure of the differential cross section of the coherent $`\gamma D\mathrm{\Lambda }^{}\mathrm{\Theta }^+`$ reaction is caused mainly by the spectral functions $`S`$ in Eqs. (17). Thus in Eq. (17a), the amplitude of the $`\mathrm{\Theta }^+nK^+`$ transition is a smooth function compared to the spectral function $`S^\mathrm{\Lambda }^{}`$ independently on the properties of $`\mathrm{\Theta }^+`$. Therefore, our predictions remain to be valid for the $`J^P=\frac{3}{2}^\pm `$ of $`\mathrm{\Theta }^+`$, considered in recent Ref. Hosaka0505 .
When our prediction is to be compared with experiments, one should pay attention, at least, the following two points. First, an energy spread in the beam photon may change the shape of the background, which is mainly determined by the quasi-free $`\mathrm{\Lambda }^{}`$ production. However, our conclusion indicated by Fig. 19 is not changed qualitatively. Second, the shape of the background is sensitive to the acceptance of the measurement. In particular, the effect of coherent $`\mathrm{\Lambda }^{}\mathrm{\Theta }^+`$ production may be significantly suppressed when the detector does not have acceptance to detect $`pK^{}`$ pair in the forward angles. In contrast, the acceptance to the forward $`pK^{}`$ like one in the case of LEPS of SPring-8 Nakano03 may make the effect more pronounced.
## VI Summary
In summary we analyzed the coherent $`\mathrm{\Lambda }^{}\mathrm{\Theta }^+`$ photoproduction in $`\gamma D`$ interaction with taking into account different background processes. We found that the behavior and the strength of the background processes depend strongly on the kinematics where the momentum distribution in the deuteron plays a key role. Thus, at fixed angle of the $`pK^{}`$ photoproduction the $`nK^+`$ invariant mass distribution of the background processes looks like a narrow peak with maximum around the $`\mathrm{\Theta }^+`$ mass. This behaviour hampers the extraction of the coherent process at finite invariant mass resolution. Most promising is an experimental analysis of the distributions integrated over the $`pK^{}`$ production angles in the forward hemisphere of c.m.s. In this case the background processes increase monotonously with $`M_{nK^+}`$ in the vicinity of $`M_{\mathrm{\Theta }^+}`$, which allows to extract the coherent $`\gamma D\mathrm{\Lambda }^{}\mathrm{\Theta }^+`$ channel even with finite invariant mass resolution. We demonstrated that the coherent $`\mathrm{\Theta }^+\mathrm{\Lambda }(1520)`$ photoproduction does not depend on the $`\mathrm{\Theta }^+`$ photoproduction amplitude, but rather it is defined by the probabilities of the $`\mathrm{\Lambda }(1520)`$ photoproduction and the $`\mathrm{\Theta }^+NK`$ transition. Therefore, this effect may be used as an independent method for studying the mechanism of $`\mathrm{\Theta }^+`$ production and $`\mathrm{\Theta }^+`$ properties.
Our model estimates for the $`\gamma D`$ reaction may be considered as an example why the $`\mathrm{\Theta }^+`$ peak is seen under certain experimental conditions and why it does not appear above the strong background in other ones.
Finally, we note that the predicted process of the coherent $`\mathrm{\Lambda }^{}\mathrm{\Theta }^+`$ photoproduction may be studied experimentally at the electron and photon facilities at LEPS of SPring-8, JLab, Crystal-Barrel of ELSA, and GRAAL of ESFR.
###### Acknowledgements.
We appreciate fruitful discussions with T. Nakano who initiated this study, and we thank H. Ejiri, M. Fujiwara, K. Hicks and A. Hosaka for useful comments and suggestions. One of authors (A.I.T.) thanks E. Grosse for offering the hospitality at FZR. This work was supported by BMBF grant 06DR121, GSI-FE.
## Appendix A Transition operators for the resonance amplitudes
### A.1 The $`๐ฏ^\mathbf{+}`$ photoproduction amplitude
We show here the explicit expressions for the transition operators $`_\mu `$ in Eq. (5) for a positive $`\mathrm{\Theta }^+`$ parity and the PS coupling scheme.
The specific parameters for the form factor in Eq. (4) are defined by
$$F_s=F(M_N,s),F_u=F(M_\mathrm{\Theta },u),\text{and}F_t=F(M_{K^+},t).$$
(36)
In addition, we need the form factor combinations
$$\stackrel{~}{F}_{tu}=F_t+F_uF_tF_u\text{and}\stackrel{~}{F}_{su}=F_s+F_uF_sF_u$$
(37)
to construct the contact terms $`_\mu ^c`$ given below that make the initial photoproduction amplitude gauge invariant hhgauge ; DavWork . The four-momenta in the following equations are defined according to the arguments given in the reaction equation
$$\gamma (k)+N(p)\mathrm{\Theta }^+(p_\mathrm{\Theta })+\overline{K}(\overline{q}).$$
(38)
#### A.1.1 $`๐ธ๐\mathbf{}๐ฏ^\mathbf{+}๐ฒ^{\mathbf{}}`$
$`_\mu ^t`$ $`=i{\displaystyle \frac{eg_{\mathrm{\Theta }NK}(k_\mu 2\overline{q}_\mu )\gamma _5}{tM_{K^+}^2}}F_t,`$ (39a)
$`_\mu ^s`$ $`=ieg_{\mathrm{\Theta }NK}\gamma _5{\displaystyle \frac{p/+k/+M_N}{sM_N^2}}\left(i{\displaystyle \frac{\kappa _p}{2M_N}}\sigma _{\mu \nu }k^\nu \right)F_s,`$ (39b)
$`_\mu ^u`$ $`=ieg_{\mathrm{\Theta }NK}\left(\gamma _\mu +i{\displaystyle \frac{\kappa _\mathrm{\Theta }}{2M_\mathrm{\Theta }}}\sigma _{\mu \nu }k^\nu \right){\displaystyle \frac{p/_\mathrm{\Theta }k/+M_\mathrm{\Theta }}{uM_\mathrm{\Theta }^2}}\gamma _5F_u,`$ (39c)
$`_\mu ^c`$ $`=ieg_{\mathrm{\Theta }NK}\gamma _5\left[{\displaystyle \frac{(k2\overline{q})_\mu }{tM_{K^+}^2}}(\stackrel{~}{F}_{tu}F_t)+{\displaystyle \frac{(2p_\mathrm{\Theta }k)_\mu }{uM_\mathrm{\Theta }^2}}(\stackrel{~}{F}_{tu}F_u)\right].`$ (39d)
The transition operator of $`t`$-channel $`K^{}`$ exchange amplitude is given by
$`_\mu ^t(K^{})={\displaystyle \frac{eg_{\gamma KK^{}}g_{\mathrm{\Theta }NK^{}}}{M_K^{}}}{\displaystyle \frac{\epsilon _{\mu \nu \alpha \beta }k^\alpha \overline{q}^\beta }{tM_K^{}^2}}\left[\gamma ^\nu i{\displaystyle \frac{\sigma ^{\nu \lambda }(pp_\mathrm{\Theta })_\lambda }{M_\mathrm{\Theta }+M_N}}\kappa ^{}\right]F(M_K^{},t).`$ (40)
#### A.1.2 $`๐ธ๐\mathbf{}๐ฏ^\mathbf{+}\overline{๐ฒ}^\mathrm{๐}`$
$`_\mu ^s`$ $`=i{\displaystyle \frac{eg_{\mathrm{\Theta }NK}}{M_\mathrm{\Theta }+M_N}}\gamma _5\overline{q}/{\displaystyle \frac{p/+k/+M_N}{sM_N^2}}\left(\gamma _\mu +i{\displaystyle \frac{\kappa _p}{2M_N}}\sigma _{\mu \nu }k^\nu \right)F_s,`$ (41a)
$`_\mu ^u`$ $`=i{\displaystyle \frac{eg_{\mathrm{\Theta }NK}}{M_\mathrm{\Theta }+M_N}}\left(\gamma _\mu +i{\displaystyle \frac{\kappa _\mathrm{\Theta }}{2M_\mathrm{\Theta }}}\sigma _{\mu \nu }k^\nu \right){\displaystyle \frac{p/_\mathrm{\Theta }k/+M_\mathrm{\Theta }}{uM_\mathrm{\Theta }^2}}\gamma _5\overline{q}/F_u,`$ (41b)
$`_\mu ^c`$ $`=i{\displaystyle \frac{eg_{\mathrm{\Theta }NK}}{M_\mathrm{\Theta }+M_N}}\gamma _5\overline{q}/\left[{\displaystyle \frac{(2p+k)_\mu }{sM_N}}(\stackrel{~}{F}_{su}F_s)+{\displaystyle \frac{(2p_\mathrm{\Theta }k)_\mu }{uM_\mathrm{\Theta }^2}}(\stackrel{~}{F}_{su}F_u)\right].`$ (41c)
### A.2 $`๐ฒ^{\mathbf{}}`$ photoproduction amplitude
We show here the explicit expressions for the transition operators $`_{\sigma \mu }`$ in Eq. (8) for the reactions $`\gamma p\mathrm{\Lambda }^{}K^+`$ and $`\gamma n\mathrm{\Lambda }^{}K^0`$.
#### A.2.1 $`๐ธ๐\mathbf{}๐ฒ^{\mathbf{}}๐ฒ^\mathbf{+}`$
$`_{\sigma \mu }^t`$ $`=i{\displaystyle \frac{eg_{\mathrm{\Lambda }^{}NK}M_\mathrm{\Lambda }^{}(2q_\mu k_\mu )(k_\sigma q_\sigma )\gamma _5}{tM_{K^+}^2}}F_t,`$ (42a)
$`_{\sigma \mu }^s`$ $`=i{\displaystyle \frac{eg_{\mathrm{\Lambda }^{}NK}}{M_\mathrm{\Lambda }^{}}}q^\sigma \gamma _5{\displaystyle \frac{p/+k/+M_N}{sM_N^2}}\left(\gamma _\mu +i{\displaystyle \frac{\kappa _p}{2M_N}}\sigma _{\mu \nu }k^\nu \right)F_s,`$ (42b)
$`_{\sigma \mu }^c`$ $`=i{\displaystyle \frac{eg_{\mathrm{\Lambda }^{}NK}}{M_\mathrm{\Lambda }^{}}}\gamma _5[{\displaystyle \frac{(2qk)_\mu (kq)_\sigma }{tM_{K^+}^2}}(\stackrel{~}{F}_{ts}F_t){\displaystyle \frac{(2p+k)_\mu }{sM_N}}(\stackrel{~}{F}_{ts}F_s)`$
$`+g_{\sigma \mu }\stackrel{~}{F}_{ts}].`$ (42c)
The corresponding form factors are defined by
$$F_s=F(M_N,s),F_t=F(M_{K^+},t),\text{and}\stackrel{~}{F}_{ts}=F_t+F_sF_tF_s.$$
(43)
The transition operator of $`t`$-channel $`K^{}`$ meson exchange amplitude is given by
$`_{\sigma \mu }^t(K^{})={\displaystyle \frac{eg_{\gamma KK^{}}g_{\mathrm{\Lambda }^{}NK^{}}}{M_K^{}M_\mathrm{\Lambda }^{}}}{\displaystyle \frac{\epsilon _{\nu \mu \alpha \beta }k^\nu q^\alpha }{tM_K^{}^2}}\left[q_\sigma ^{}\gamma _\sigma q/^{}g_{\sigma \beta }\right]F(M_K^{},t)`$ (44)
with $`q^{}=p_\mathrm{\Lambda }^{}p`$.
#### A.2.2 $`๐ธ๐\mathbf{}๐ฒ^{\mathbf{}}๐ฒ^\mathrm{๐}`$
$`_{\sigma \mu }^s=i{\displaystyle \frac{eg_{\mathrm{\Theta }NK}}{M_\mathrm{\Lambda }^{}}}q^\sigma \gamma _5{\displaystyle \frac{p/+k/+M_N}{sM_N^2}}\left(i{\displaystyle \frac{\kappa _p}{2M_N}}\sigma _{\mu \nu }k^\nu \right)F_s.`$ (45)
The $`t`$-channel $`K^{}`$-exchange operator is defined by Eq. (44) with appropriate coupling constant $`g_{\gamma KK^{}}`$.
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# 1 Introduction
## 1 Introduction
While there may be good reason to believe that planet-disk interactions may excite the eccentricities of at least isolated planets (Goldreich and Sari 2003; Ogilvie and Lubow 2003) there are currently no studies to guide our understanding of likely outcomes in the case of multiple planet (or satellite) systems. Therefore, one may consider circular, coplanar giant planets as a starting condition. Tsiganis et al. (2005) have recently argued that evolution through the 1:2 Jupiter-Saturn mean motion resonance (MMR) of a quasi-circular, coplanar and compact solar system that is allowed to evolve by planetesimal scattering is consistent with the observed semi-major axes, eccentricities and mutual inclinations of Jupiter, Saturn, Uranus and Neptune. In a companion publication, Gomes et al. (2005) argue that this resonance crossing may be linked to the Late Heavy Bombardment of the terrestrial planets taking place $`700`$ Myr after their formation (Hartmann et al. 2000). For such a scenario to work the bulk of the divergent migration between Jupiter and Saturn and the resonance passage itself must take place following gas dissipation, and also later than simulations with an evenly spread disk of planetesimals would indicate. Whether or not the identification with the Late Heavy Bombardment is correct, in this viewpoint the solar system must have been more compact and regular prior to gas dissipation (during its first $`10^610^7`$ years) than it is today. The authors also briefly consider the potentially disruptive consequences of such a scenario for the regular and irregular satellites of the giant planets, but conclude that at least the regular satellites might have โsurvivedโ unscathed.<sup>1</sup><sup>1</sup>1Though the issue is not spelled out in detail (presumably due to space limitations), in the case of the Saturnian satellite system Iapetusโ relatively low eccentricity $`e0.03`$ may be the main constraint for such a scenario. For this reason, it would appear unlikely that Titan could owe its eccentricity to a close encounter between Saturn and another giant planet.
Jupiterโs low obliquity $`3^{}`$ may be indicative of its formation by hydrodynamic gas accretion. Yet, it has been noted that secular spin-orbit resonances can complicate this straightforward interpretation. In particular, adiabatic passage through a resonance matching the spin axis precession rate to the $`\nu _{16}`$ precession frequency of the orbit plane due to the gravitational perturbation of Saturn (Hamilton and Ward 2002; Canup and Ward 2002; Ward 2003) and $`\nu _{17}`$ due to the gravitational perturbation of Uranus (Hamilton, pers. comm.) may result in obliquities significantly larger than observed. For Jupiter the amplitude of the $`\nu _{16}`$ term (with a period of $`50,000`$ years) is $`0{}_{}{}^{}.36`$ and the $`\nu _{17}`$ term (with a period of $`P_{16}450,000`$ yrs) is $`0{}_{}{}^{}.055`$, which could have resulted in obliquities of up to $`26^{}`$ and $`14^{}`$, respectively, had these resonances been crossed adiabatically as the circumplanetary gas disk was dissipated (Ward 2003). This leads Canup and Ward (2002) to argue that the circumplanetary disk must have viscously evolved in a timescale sufficiently short ($`O(10^5)`$ yr) as to preclude adiabatic passage, resulting in a gas-starved (or gas-poor \[Mosqueira et al. 2000; Estrada and Mosqueira 2005\]) satellite disk.
The issue we tackle here is whether our decaying turbulence<sup>2</sup><sup>2</sup>2Consistent with numerical simulations that show turbulence decay in the absence of a โstirringโ mechanism (Hawley et al. 1999). SEMM model (Mosqueira and Estrada 2003a,b; hereafter MEa,b) is especially susceptible to secular spin-orbit resonances, and inconsistent with Jupiterโs low obliquity. In particular, we focus on the $`\nu _{16}`$ term. The reasons for this are: First, the period of the orbital precession $`\nu _{17}`$ ($`P_{17}450,000`$ yrs) is already very close to the precession period of Jupiterโs spin axis due to the solar torque on the Galilean satellites (which are locked to the Jupiterโs equator plane by this planetโs oblateness; Goldreich 1965). A slight adjustment of satellite or planet positions might be enough to place Jupiter spin axis precession period in one side or the other of this resonance, so that any formation model is apt to be affected. Second, even if the resonance is crossed the limiting obliquity for adiabatic passage is significantly smaller than that for the $`\nu _{16}`$ term. Furthermore, the timescale for adiabatic passage in the case of the $`\nu _{17}`$ secular spin-orbit resonance may be longer ($`O(10^7)`$ yrs) than the timescale for gas dissipation by photoevaporation<sup>3</sup><sup>3</sup>3Shu et al. (1993) conclude that EUV from the central star may photoevaporate a T Tauri disk in $`10^7`$ yrs outside of a gravitational radius $`r_g10`$ AU. But there is considerable uncertainty in this estimate. For instance, Adams et al. (2004) argue that the disk is likely to be photoevaporated by FUV radiation from neighboring stars. Furthermore, these authors state that photoevaporation is effective to a significantly smaller radius $`0.10.2r_g`$.. Third, it is likely that Uranus and Neptune formed after Jupiter and Saturn, once most of gas in the planetary and circumplanetary disks had already dissipated. Thus, from here on we focus on $`\nu _{16}`$ and consider Jupiter and Saturn only.
## 2 Secular Perturbation Theory
In a solar system consisting of Jupiter and Saturn it is straightforward to construct a Laplace-Lagrange secular solution (e.g., Brouwer and Clemence 1961; Murray and Dermott 1999). For $`I<<1`$, the orbital inclination $`I`$ and the longitude of the ascending node with respect to the invariant plane $`\mathrm{\Omega }`$ are given by
$$I\mathrm{sin}\mathrm{\Omega }=I_1\mathrm{sin}\gamma _1+I_2\mathrm{sin}(f_2t+\gamma _2),$$
(1)
and
$$I\mathrm{cos}\mathrm{\Omega }=I_1\mathrm{cos}\gamma _1+I_2\mathrm{cos}(f_2t+\gamma _2),$$
(2)
where $`f_2`$ is the eigenfrequency, $`I_1`$ and $`I_2`$ are eigenvector components, and $`\gamma _1`$ and $`\gamma _2`$ are phases. If we use parameters for Jupiter and Saturn as observed, we obtain $`f_2=7.06\times 10^3^{}`$ yr<sup>-1</sup> and $`I_2=6.30\times 10^3`$ (in radians). This yields a secular oscillation with period of $`P_{16}51,000`$ yr. These values are not too dissimilar from the secular solution of the planetary system (e.g., Murray and Dermott 1999).
Let us now consider a time early on before Jupiter and Saturn had passed through the 1:2 MMR but after most of the planetary gas disk had dissipated by photoevaporation in a timescale $`10^610^7`$ yr. Since the planetary nebula may shield it to some degree, it may be appropriate to assume that the much denser subnebula takes longer to dissipate. If so, at this time both the precession of Jupiterโs spin axis and orbital plane may be significantly faster than they are today. Because of the scattering on nearby planetesimals (Gomes et al. 2005) and the dissipation of the nebula, Saturn would dominate the precession of Jupiterโs orbital plane.<sup>4</sup><sup>4</sup>4Ward (1981) treats the two-orbit/nebula problem as dispersal takes place. Here we consider the case when the nebula (but not the subnebula) has already been dissipated by some means. We take as nominal values for the orbital parameters of Jupiter and Saturn the starting conditions in Tsiganis et al. (2005), namely, $`a_J=5.45`$ AU, $`a_S=8.50`$ AU (a few tenths of an AU inside the 1:2 MMR) and $`\mathrm{sin}I_{JS}10^3`$, where $`I_{JS}`$ is the relative inclination of the orbital planes of Jupiter and Saturn. With these parameters<sup>5</sup><sup>5</sup>5 An estimate for the gap opening timescale may be obtained using the tidal torque formula (Lin and Papaloizou 1993) $`\tau _{gap}(\mathrm{\Delta }/a)^5P/\mu ^2`$, where $`\mathrm{\Delta }`$ is the gapโs half-width, $`\mu `$ is the mass ratio of the secondary to the primary, and $`P`$ is the orbital period of the secondary. Using $`\mathrm{\Delta }(a_Sa_J)/21.5`$ AU, we find that Jupiter and Saturn would have cleared the gas disk in between in a short timescale $`t10^4`$ yrs. In particular, this timescale is shorter than satellite formation timescale $`10^410^6`$ (depending on location; Mosqueira et al. 2001) in MEa,b. we obtain $`f_2=1.60\times 10^2^{}`$ yr<sup>-1</sup> and $`I_2=2.72\times 10^4`$ (in radians). As shown in Fig. 1, in this case the period of Jupiterโs orbital plane precession is $`P_{16}23,000`$ yrs.
## 3 Obliquity Variation
Jupiterโs current spin axis precession period is $`4.5\times 10^5`$ yrs due mostly to the solar torque exerted on the Galilean satellites (e.g., Ward 1975)<sup>6</sup><sup>6</sup>6We obtain a precession period of $`4.8\times 10^5`$ yrs using a moment of inertia for Jupiter of $`K=0.26`$, which yields a spin angular momentum $`J_p=4.4\times 10^{45}`$ g cm<sup>2</sup>/s. There is some uncertainty in the moment of inertia $`5\%`$ (Fortney J., pers. comm.) but this doesnโt affect the argument. What is important to note, however, is that about $`33\%`$ of the precession constant $`\alpha `$ is contributed by the torque of the Sun directly on the planet because of its oblateness. Adding a circumplanetary gas disk decreases this fraction, which we then ignore.. However, a massive circumplanetary gas disk would result in a much shorter precession period. The inner parts of the disk ($`r<r_t=\left(2M_p/M_{}J_2R_p^2a_p^3\right)^{1/5}40R_J`$, where $`M_p=M_J`$, $`a_p=5.45`$ AU, $`R_p2R_J`$, and $`J_21/R_p0.008`$ are the planetary mass, semi-major axis, radius and quadrupole gravitational harmonic following envelope collapse \[which may take place in a fast $`10^410^5`$ yrs timescale, Hubickyj O., pers. comm.\], and $`M_{}`$ and $`R_J`$ are the Sunโs mass and Jupiterโs present radius) would orbit in the plane of the planetโs equator and precess as a unit with the planet, whereas the outer parts of the disk ($`r>r_t`$) would not (Goldreich 1966). The spin axis $`\widehat{๐ฌ}`$ then precesses around the orbit normal $`\widehat{๐ง}`$ at a rate given by (e.g., Tremaine 1991)
$$\frac{d\widehat{๐ฌ}}{dt}=\alpha (\widehat{๐ฌ}\widehat{๐ง})(\widehat{๐ฌ}\times \widehat{๐ง}),$$
(3)
where $`\widehat{๐ฌ}`$ and $`\widehat{๐ง}`$ are unit vectors, and the precession constant<sup>7</sup><sup>7</sup>7In actuality, there should be another term added to this precession rate due to the torque of the extended part of the disk $`r>r_t`$ on the inner part of disk. However, the contribution of this extended region drops rapidly with distance, and in our model the gas surface density drops-off at a radial location $`r_Dr_t`$. is given by
$$\alpha =\frac{3\pi \mathrm{\Omega }_p^2}{2H}_{R_p}^{r_D}\mathrm{\Sigma }(a)a^3๐a,$$
(4)
where the surface density of circumplanetary disk $`\mathrm{\Sigma }(a)`$ is assumed to drop-off sharply at $`r_D2r_c`$, where $`r_c=R_H/4815R_J`$ is the centrifugal radius (MEa), $`\mathrm{\Omega }_p`$ is the planetโs orbital frequency and the total angular momentum of the precessing system is given by
$$H=J_p+2\pi (GM_p)^{1/2}_{R_p}^{r_D}\mathrm{\Sigma }(a)a^{3/2}๐a,$$
(5)
where $`G`$ is the gravitational constant and $`J_p`$ is the spin angular momentum of the planet. For $`\mathrm{\Sigma }1/a`$, we can write $`\alpha =M_D\mathrm{\Omega }_p^2r_D^2/(4H)`$ and $`H=J_p+2/3M_D\sqrt{GM_pr_D}`$, where $`M_D`$ is the disk mass. The precession period is given by $`T=2\pi /(\alpha \mathrm{cos}\theta )`$, where $`\mathrm{cos}\theta =\widehat{๐ฌ}\widehat{๐ง}`$ is the obliquity. Given that in our SEMM model $`M_D10M_{sats}`$, where $`M_{sats}4\times 10^{26}`$ g is the mass of the Galilean satellites<sup>8</sup><sup>8</sup>8There is some ambiguity here because in the model of MEa,b a significant fraction of the mass of Callisto is derived from the extended part of the disk, whereas $`M_D`$ is the mass of the inner disk out to about Callisto. On the other hand, Io and Europa should be reconstituted for unaccreted volatiles., the spin of the planet provides $`90\%`$ of the angular momentum of the system $`H5\times 10^{45}`$ g cm<sup>2</sup>/s (see MEa Table 3). Using an obliquity of $`\theta 3^{}`$ and $`r_D=40R_J`$ (which implies a surface density $`\mathrm{\Sigma }2\times 10^4`$ g/cm<sup>2</sup> at $`15R_J`$ consistent with the value obtained by applying the inviscid gap-opening criterion to Ganymede in a disk with aspect ratio $`0.1`$ \[MEb\]), we obtain a precession period of $`T4\times 10^4`$ yrs, which is slightly shorter than Jupiterโs current orbital plane secular precession period $`P_{16}5\times 10^4`$ yrs, but it is longer than Jupiterโs orbital plane precession period when Saturn was inside the 1:2 MMR, i.e., $`P_{16}2\times 10^4`$ yrs. Hence, in our decaying turbulence SEMM model it is possible for Jupiter either to have crossed the secular spin-orbit resonance before the Keplerian disk reached its quiescent phase (and satellites formed) at a time when the viscous evolution of the disk was likely driven by Roche-lobe gas inflow and the resonance passage was non-adiabatic, or not to have crossed this resonance at all. However, it may be possible to alter this conclusion by choosing different parameters, such as a more massive and extended circumplanetary disk. Thus, next we consider the case of resonance passage.
The limiting obliquity $`\theta _{max}`$ that could be generated by the obliquity โkickโ incurred during adiabatic resonance passage in the non-capture direction ($`\dot{\alpha }<0`$) (which implies resonance crossing as the circumplanetary gas disk is dissipated and the spin-axis precession period increases) is given by (Henrard and Murigande 1987; Ward and Hamilton 2004)
$$\mathrm{cos}\theta _{max}=\frac{2}{\left(1+\mathrm{tan}^{2/3}|I_2|\right)^{3/2}}1,$$
(6)
which yields $`\theta _{max}26^{}`$ using the current value for $`|I_2|0{}_{}{}^{}.36`$, but a significantly smaller value of $`\theta _{max}9.1^{}`$ using $`|I_2|0{}_{}{}^{}.015`$ obtained from our nominal, low mutual inclination case. Furthermore, the minimum time for adiabatic crossing is (Ward and Hamilton 2004)
$$\tau _{min}P_{16}\left(\frac{\theta }{2\pi |I_2|}\right)^2.$$
(7)
Taking $`\theta 9^{}`$, we find $`2\times 10^8`$ yrs for our nominal case, which is much longer than the gas dissipation timescale. This means that at least for the nominal case the crossing must be non-adiabatic and the final obliquity is rate dependent. Taking $`\theta 3^{}`$, we calculate $`\tau _{min}9\times 10^4`$ yrs for current solar-system parameters, but it is $`\tau _{min}2\times 10^7`$ yrs (here $`P_{16}23,000`$ yrs is about a factor of two shorter than the present value, but this is more than compensated by the much smaller value of $`|I_2|`$) for our nominal case, which is comparable to or longer than the gas dissipation timescale. We may estimate the resulting obliquity to be $`\theta \mathrm{tan}^{1/3}(|I_2|)4^{}`$ for our nominal case, which is consistent with Jupiterโs observed value. At any rate, this would all be irrelevant if the resonance were not crossed.
## 4 Conclusions
We have briefly investigated the consequences for Jupiterโs obliquity of a satellite formation model in which a massive (compared to the Galilean satellites yet enhanced in solids by a factor of $`10`$ compared to the solar composition minimum mass model) subnebula is allowed to dissipate in a long timescale ($`10^610^7`$ yrs) after the dispersal of the nebula itself. For the sake of specificity, we have adopted the model of Tsiganis et al. (2005) in which the solar system was more compact and regular before Jupiter and Saturn crossed the 1:2 MMR than it is today. We find that such a combined scenario (by no means unique) does not imply the likelihood of a larger obliquity for Jupiter than is observed. This is both because the secular $`\nu _{16}`$ spin-orbit resonance may not be crossed, and because the resulting obliquity may be consistent with Jupiterโs value even if it is crossed. We conclude that Jupiterโs low obliquity is compatible with our SEMM satellite formation model (MEa,b) provided one allows for solar-system conditions early-on unlike those presently observed.
Acknowledgements
This work is supported by grants from PGG and the National Research Council.
References
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# A tail-matching method for the linear stability of multi-vector-soliton bound states
## 1 Introduction
Nonlinear optics and fiber communication systems are advancing very rapidly these days. In this process, a widely-used mathematical model is the coupled nonlinear Schrรถdinger (NLS) equations which govern pulse propagation in birefringent fibers . Similar equations with a saturable nonlinearity also govern the interaction of two incoherently-coupled laser beams . These equations arise in water-wave interactions as well . Solution properties of the coupled NLS equations have been examined extensively in the past ten years. It is known that these equations admit single-hump vector solitons due to a nonlinear mutual trapping effect . When these solitons are perturbed, they undergo long-lived internal oscillations . When they collide with each other, a fractal structure can arise in the parameter space .
Multi-vector-soliton bound states also exist in the coupled NLS equations . These states are pieced together by several single-hump vector solitons. They received much attention because of several reasons. First, in fiber communication systems, pulse-pulse interference impairs the system performance. If multi-soliton bound states exist, these solutions would have implications to system designs. Second, the existence of such bound states is noteworthy because they can not exist in the NLS equation . Thirdly, these states are closely related to similar states in incoherent laser beams, which have been observed experimentally .
After the numerical discovery of these multi-vector-soliton bound states in , their analytical construction was made in by a tail-matching technique. Their linear-stability problem was studied later in by an extension of the Karpman-Solevโev-Gorshkov-Ostrovsky (KSGO) method , and small eigenvalues bifurcated from the zero eigenvalue of single vector solitons were calculated. These calculations show that multi-soliton bound states are always linearly unstable. But this KSGO method is quite involved, thus a simpler technique for the calculation of eigenvalue bifurcations is called upon.
In this paper, we use a new tail-matching method to analyze the linear stability of two-vector-soliton bound states in the coupled NLS equations. Under the condition that individual vector solitons in these bound states are wave-and-daughter-waves (i.e., one component is much larger than the other one), and are widely separated, small eigenvalues of these bound states that bifurcate from the zero eigenvalue of single solitons are calculated. These small eigenvalues are all the non-zero discrete eigenvalues of the two-soliton bound states. We found that unstable eigenvalues from phase-mode bifurcations always exist, thus the bound states are always linearly unstable. The present technique is much simpler, but it gives identical results as the KSGO method .
## 2 Two-vector-soliton bound states: a review
The coupled NLS equations
$$iA_t+A_{xx}+(|A|^2+\beta |B|^2)A=0,$$
(2.1)
$$iB_t+B_{xx}+(|B|^2+\beta |A|^2)B=0,$$
(2.2)
govern optical pulse propagation in birefringent fibers . Here $`\beta `$ is the cross-phase-modulational coefficient. When $`\beta =0`$ or 1, these equations are integrable by the inverse scattering method .
These equations admit solitary-wave solutions of the following form:
$$A(x,t)=r(x)e^{i\omega ^2t},B(x,t)=R(x)e^{i\mathrm{\Omega }^2t},$$
(2.3)
where $`\omega `$ and $`\mathrm{\Omega }`$ are frequencies, and the real-valued amplitude functions $`r(x)`$ and $`R(x)`$ satisfy the ordinary differential equations (ODEs):
$$r_{xx}\omega ^2r+(r^2+\beta R^2)r=0,$$
(2.4)
$$R_{xx}\mathrm{\Omega }^2R+(R^2+\beta r^2)R=0.$$
(2.5)
After a simple rescaling of variables, we normalize $`\omega =1`$. Since Eqs. (2.1) and (2.2) are Galilean-invariant, moving solitary waves can be readily deduced from the stationary ones (2.3) (see ).
Solitary waves in Eqs. (2.4) and (2.5) have been studied extensively before (see and the references therein). It has been shown that for any frequency $`\mathrm{\Omega }[\mathrm{\Omega }_c,1/\mathrm{\Omega }_c]`$ where
$$\mathrm{\Omega }_c\frac{\sqrt{1+8\beta }1}{2},$$
(2.6)
this ODE system admits a unique, single-hump, and positive vector-soliton solution which is symmetric in both $`r`$ and $`R`$ components. We call this solution the fundamental vector soliton, and denote it as $`[r_0(x),R_0(x)]`$. The asymptotic behavior of this fundamental soliton at infinity is
$$r_0(x)ce^{|x|},R_0(x)Ce^{\mathrm{\Omega }|x|},|x|\mathrm{},$$
(2.7)
where $`c`$ and $`C`$ are tail coefficients. When $`\mathrm{\Omega }`$ is close to the boundary value $`\mathrm{\Omega }_c`$, $`R_0(x)r_0(x)`$, $`c2\sqrt{2}`$, and $`C1`$; if $`\mathrm{\Omega }`$ is close to $`1/\mathrm{\Omega }_c`$, $`r_0(x)R_0(x)`$, $`c1`$, and $`C2\sqrt{2}/\mathrm{\Omega }_c`$. These vector solitons with $`R_0r_0`$ or $`r_0R_0`$ are the so-called wave-and-daughter-waves.
Two-vector-soliton bound states also exist in the ODE system (2.4) and (2.5) . These bound states look like two single-humped vector solitons glued together, while the two solitons are in-phase in one component, and out-of-phase in the other component. The in-phase component of the bound states are symmetric around the bound-state center, and the out-of-phase component are anti-symmetric around the bound-state center. In the limit of large soliton separation, these bound states approach a superposition of two fundamental solitons (to the leading order):
$$r(x)r_0(x)r_0(x\mathrm{\Delta }),$$
(2.8)
$$R(x)R_0(x)\pm R_0(x\mathrm{\Delta }),$$
(2.9)
where the separation $`\mathrm{\Delta }1`$. These widely-separated bound states exist in two parameter-regions : (i) $`\mathrm{\Omega }\mathrm{\Omega }_c`$ or $`1/\mathrm{\Omega }_c`$, and $`0<\beta <1`$; (ii) $`\mathrm{\Omega }1`$, and $`\beta >0`$. In the first region, the bound states look like two wave-and-daughter-waves glued together; while in the second region, the bound states look like two nearly-equal-amplitude vector solitons glued together. In this article, we only consider the bound states in the first region. In this region, the spacing between the two wave-and-daughter-waves in the bound state is given by the formula
$$\mathrm{\Delta }=\frac{\mathrm{ln}c^2\mathrm{ln}\mathrm{\Omega }^2C^2}{1\mathrm{\Omega }}.$$
(2.10)
Below, the bound states with the minus sign in (2.8) and the plus sign in (2.9) will be termed type-I bound states, while the ones with the plus sign in (2.8) and the minus sign in (2.9) termed type-II bound states (as we have done in ). These bound states at parameter values $`\beta =2/3`$ and $`\mathrm{\Omega }=0.85`$ are displayed in Fig. 1. A comparison between the analytical spacing formula (2.10) and numerical values has been made in , and excellent agreement has been obtained.
Analytical construction of multi-soliton bound states in a general nonlinear wave system was made in using a tail-matching method, and the spacing formula (2.10) for the coupled NLS equations was derived only as a special case (see for an application of this method for the construction of other types of multi-pulse bound states). Below, we use a simplified version of โs method to construct multi-vector-soliton bound states in the coupled NLS equations \[i.e., (2.4) and (2.5)\], and reproduce the spacing formula (2.10). There are two reasons for our doing this: (i) to highlight the key ideas in the tail-matching method for the construction of multi-pulse bound states; (ii) to motivate a similar tail-matching idea for the linear-stability analysis of multi-pulse bound states (see Sec. 3).
Suppose we have a bound state of two vector solitons located at $`x=0`$ and $`\mathrm{\Delta }(1)`$. As $`\mathrm{\Delta }\mathrm{}`$, the leading-order asymptotics of this bound state is
$$r(x)r_0(x)+s_1r_0(x\mathrm{\Delta }),R(x)R_0(x)+s_2R_0(x\mathrm{\Delta }),$$
(2.11)
where $`s_1`$ and $`s_2`$ are sign-constants and are either 1 or $`1`$. Our task is to determine values of $`s_1,s_2`$ and the spacing $`\mathrm{\Delta }`$. For this purpose, we consider the bound state in two $`x`$-regions: $`\mathrm{}<x\mathrm{\Delta }`$, and $`0x<\mathrm{}`$. Since the treatments for these two regions are the same, we only look at the first region $`\mathrm{}<x\mathrm{\Delta }`$. In this region, the bound state is a slightly perturbed fundamental vector soliton, i.e.,
$$r(x)=r_0(x)+\stackrel{~}{r}(x),R(x)=R_0(x)+\stackrel{~}{R}(x),$$
(2.12)
where $`\stackrel{~}{r},\stackrel{~}{R}1`$. The actual forms and sizes of $`\stackrel{~}{r}`$ and $`\stackrel{~}{R}`$ are not important in this analysis, but their asymptotics in the region $`x\frac{1}{2}\mathrm{\Delta }1`$ is crucial. This asymptotics can be obtained by matching $`[r(x),R(x)]`$โs expressions (2.12) with their asymptotics (2.11). This is the key idea of the method. This matching gives the leading-order asymptotics of $`(\stackrel{~}{r},\stackrel{~}{R})`$ as
$$\left[\begin{array}{c}\stackrel{~}{r}(x)\\ \stackrel{~}{R}(x)\end{array}\right]\left[\begin{array}{c}s_1r_0(x\mathrm{\Delta })\\ s_2R_0(x\mathrm{\Delta })\end{array}\right]\left[\begin{array}{c}s_1ce^{x\mathrm{\Delta }}\\ s_2Ce^{\mathrm{\Omega }(x\mathrm{\Delta })}\end{array}\right],x\frac{1}{2}\mathrm{\Delta }1.$$
(2.13)
As $`x\mathrm{}`$, $`(\stackrel{~}{r},\stackrel{~}{R})0`$.
When Eq. (2.12) is substituted into (2.4) and (2.5), and higher-order terms dropped, the linearized equations for perturbations $`(\stackrel{~}{r},\stackrel{~}{R})`$ are found to be
$$L\left[\begin{array}{c}\stackrel{~}{r}\\ \stackrel{~}{R}\end{array}\right]=0,$$
(2.14)
where the linearization operator $`L`$ is
$$L=\left[\begin{array}{cc}_{xx}1+3r_0^2(x)+\beta R_0^2(x)& 2r_0(x)R_0(x)\\ 2r_0(x)R_0(x)& _{xx}\mathrm{\Omega }^2+3R_0^2(x)+\beta r_0^2(x)\end{array}\right],$$
(2.15)
which is self-adjoint. Eq. (2.14) has one localized homogeneous solution $`[r_0^{}(x),R_0^{}(x)]^\text{T}`$ when $`\beta 0`$ and 1. Here the prime is the derivative, and the superscript โTโ is the transpose of a vector. Thus, in order for the solution $`(\stackrel{~}{r},\stackrel{~}{R})`$ of Eq. (2.14) to exist, a solvability condition
$$_{\mathrm{}}^{x_m}[r_0^{}(x),R_0^{}(x)]L\left[\begin{array}{c}\stackrel{~}{r}\\ \stackrel{~}{R}\end{array}\right]๐x=0$$
(2.16)
must be satisfied. Here $`x_m\frac{1}{2}\mathrm{\Delta }`$. This condition can be simplified through integration by parts as
$$\left[\stackrel{~}{r}^{}(x)r_0^{}(x)\stackrel{~}{r}(x)r_0^{\prime \prime }(x)+\stackrel{~}{R}^{}(x)R_0^{}(x)\stackrel{~}{R}(x)R_0^{\prime \prime }(x)\right]_{\mathrm{}}^{x_m}=0.$$
(2.17)
Finally, substitution of the asymptotics (2.7) and (2.13) into the above condition gives
$$e^{(1\mathrm{\Omega })\mathrm{\Delta }}=\frac{s_1s_2C^2\mathrm{\Omega }^2}{c^2}.$$
(2.18)
This equation readily shows that in order for the bound state to exist, we must have $`s_1s_2=1`$. In that case, the spacing formula (2.18) becomes exactly the same as (2.10). Hence the above simplified tail-matching method reproduces the results from the more general analysis in .
The relative errors of the leading-order bound states (2.8), (2.9), and the spacing formula (2.10) have been discussed in , and these errors are $`O(e^\mathrm{\Delta },e^{\mathrm{\Omega }\mathrm{\Delta }})`$. For the bound states, we have
$$r(x)=\left[r_0(x)r_0(x\mathrm{\Delta })\right]\left[1+O(e^\mathrm{\Delta },e^{\mathrm{\Omega }\mathrm{\Delta }})\right],$$
(2.19)
$$R(x)=\left[R_0(x)\pm R_0(x\mathrm{\Delta })\right]\left[1+O(e^\mathrm{\Delta },e^{\mathrm{\Omega }\mathrm{\Delta }})\right].$$
(2.20)
These error estimates can also be obtained as follows. Let us consider the type-I bound state. Write
$$r(x)=r_0(x)r_0(x\mathrm{\Delta })+\widehat{r}(x),$$
(2.21)
$$R(x)=R_0(x)+R_0(x\mathrm{\Delta })+\widehat{R}(x),$$
(2.22)
where $`\widehat{r},\widehat{R}1`$. Substituting these equations into (2.4) and (2.5), we find that in the region $`\mathrm{}<x\mathrm{\Delta }`$,
$$L\left[\begin{array}{c}\widehat{r}\\ \widehat{R}\end{array}\right]=\left[\begin{array}{c}\left[3r_0^2(x)+\beta R_0^2(x)\right]r_0(x\mathrm{\Delta })+2\beta r_0(x)R_0(x)R_0(x\mathrm{\Delta })\\ \left[3R_0^2(x)+\beta r_0^2(x)\right]R_0(x\mathrm{\Delta })2\beta r_0(x)R_0(x)r_0(x\mathrm{\Delta }),\end{array}\right].$$
(2.23)
Here terms which are higher-order than the ones kept in (2.23) have been dropped. A similar equation can be obtained in the region $`0x<\mathrm{}`$. Inspection of these equations shows that for all values of $`x`$, $`(\widehat{r},\widehat{R})`$ is exponentially small compared to the leading-order terms in (2.8) and (2.9), and the relative errors are $`O(e^\mathrm{\Delta },e^{\mathrm{\Omega }\mathrm{\Delta }})`$ as shown in (2.19) and (2.20).
The spacing formula (2.10) together with the results in reveals that, in order for $`\mathrm{\Delta }1`$, we must have $`0<\beta <1`$. In addition, we must have either $`C/c1`$, $`\mathrm{\Omega }<1`$, or $`C/c1`$, $`\mathrm{\Omega }>1`$. In the former limit, $`R_0(x)r_0(x)`$, and $`\mathrm{\Omega }\mathrm{\Omega }_c`$; while in the latter limit, $`r_0(x)R_0(x)`$, and $`\mathrm{\Omega }1/\mathrm{\Omega }_c`$. In both limits, the bound states look like two wave-and-daughter-waves glued together. These two limits are actually equivalent through a scaling of variables (the so-called reciprocal relation in ). Thus in the rest of this article, we will just consider the former limit where $`\mathrm{\Omega }\mathrm{\Omega }_c`$. Fundamental solitons in this limit have been perturbatively determined in . Utilizing those results, we readily find that the asymptotic formula for the spacing $`\mathrm{\Delta }`$ is
$$\mathrm{\Delta }\frac{1}{1\mathrm{\Omega }_c}\left\{\mathrm{ln}(\mathrm{\Omega }\mathrm{\Omega }_c)+K+o(1)\right\},\mathrm{\Omega }\mathrm{\Omega }_c,$$
(2.24)
where the constant $`K`$ is
$$K=(32\mathrm{\Omega }_c)\mathrm{ln}22\mathrm{ln}\mathrm{\Omega }_c+\mathrm{ln}\gamma ,$$
and
$$\gamma =\frac{(1\mathrm{\Omega }_c^3)_{\mathrm{}}^{\mathrm{}}\text{sech}^{4\mathrm{\Omega }_c}x๐x}{2\mathrm{\Omega }_c_{\mathrm{}}^{\mathrm{}}\text{sech}^{2\mathrm{\Omega }_c}x๐x}.$$
(2.25)
## 3 Linear-stability analysis of two-vector-soliton <br>bound states
To study the linear stability of the above two-vector-soliton bound states, we perturb these states as
$$A=e^{it}\left\{r(x)+\psi _1(x)e^{i\lambda t}+\psi _2^{}(x)e^{i\lambda ^{}t}\right\},$$
(3.1)
$$B=e^{i\mathrm{\Omega }^2t}\left\{R(x)+\psi _3(x)e^{i\lambda t}+\psi _4^{}(x)e^{i\lambda ^{}t}\right\},$$
(3.2)
where $`[r(x),R(x)]`$ is a two-vector-soliton bound state, $`\psi _k(1k4)`$ are infinitesimal disturbances, $`\lambda `$ is the eigenvalue, and $`\psi ^{}`$ is the complex conjugate of $`\psi `$. Substituting these equations into (2.1) and (2.2) and dropping higher-order terms, we get the following eigenvalue relation
$$\mathrm{\Psi }=\lambda \mathrm{\Psi },$$
(3.3)
where the linearization operator $``$ is
$$=\left(\begin{array}{cccc}_{xx}1+2r^2+\beta R^2& r^2& \beta rR& \beta rR\\ r^2& _{xx}+12r^2\beta R^2& \beta rR& \beta rR\\ \beta rR& \beta rR& _{xx}\mathrm{\Omega }^2+2R^2+\beta r^2& R^2\\ \beta rR& \beta rR& R^2& _{xx}+\mathrm{\Omega }^22R^2\beta r^2\end{array}\right),$$
(3.4)
and $`\mathrm{\Psi }=(\psi _1,\psi _2,\psi _3,\psi _4)^\text{T}`$. The spectrum of $``$ contains all information on the linear stability of the two-vector-soliton bound state. If $``$ has eigenvalues with Im($`\lambda )<0`$, then the bound state is linearly unstable; otherwise, it is linearly stable. Obviously, the continuous spectrum of $``$ always lies on the real axis, thus we only need to look at the discrete eigenvalues of $``$. Notice that $`^2`$ is self-adjoint, thus $``$โs eigenvalues are either purely real, or purely imaginary. In addition, if $`\lambda `$ is an eigenvalue of $``$, so are $`\lambda `$. Hence $``$โs eigenvalues always come in pairs on the real or imaginary axis.
It is insightful to view $``$โs discrete eigenvalues for the bound states $`(r,R)`$ in the perspective of eigenvalue bifurcations from a similar operator $`_0`$ for fundamental vector solitons $`(r_0,R_0)`$. Here
$$_0=|_{(r=r_0,R=R_0)}.$$
(3.5)
The spectrum structure of $`_0`$ has been determined completely in . This operator has 6 or 8 discrete eigenvalues (multiplicity counted), depending on whether an internal mode exists or not. The zero eigenvalue always has multiplicity 6 โ three from position and phase invariances (Goldstein modes), and the other three from velocity and frequency (or equivalently, amplitude) variations. When an internal mode exists, a pair of real eigenvalues of opposite sign are present as well. If two vector solitons form a widely-separated stationary bound state $`(\mathrm{\Delta }1)`$, both solitons must be either wave-and-daughter-waves (with $`\mathrm{\Omega }\mathrm{\Omega }_c`$ and $`0<\beta <1`$), or having nearly equal amplitudes in the two components (with $`\mathrm{\Omega }1`$ and $`\beta >0`$) (only the former type of bound states is studied in this paper). It has been shown in that wave-and-daughter-waves do not admit internal modes. For instance, single vector solitons in Fig. 1โs bound states do not have internal modes . Thus $`_0`$ has only a single discrete eigenvalue zero of multiplicity 6. In a bound state of two wave-and-daughter-waves, $``$ will then have 12 discrete eigenvalues (multiplicities counted) โ double that for a single wave-and-daughter-wave. Here, no new discrete eigenvalues of $``$ can be generated from edge bifurcations of $`_0`$ because the edge points of $`_0`$โs continuous spectrum are not in the continuous spectrum. Now the zero eigenvalue of $``$ still has multiplicity 6. Another six eigenvalues of $``$ must bifurcate from zero due to tail interactions between solitons. Tracking of these six bifurcated eigenvalues will then provide a complete characterization of linear stability for these two-soliton bound states.
In this section, we propose a new and general tail-matching method to derive explicit analytical formulas for the six eigenvalues of $``$ that bifurcate from zero, under the condition that the individual vector solitons in the bound state are widely separated. The idea of this method is to perturbatively determine the bifurcated eigenstate around each vector soliton. By matching their tail asymptotics in the center region of the eigenstate, their asymptotics in that region can then be explicitly obtained. Finally by utilizing the solvability conditions for the eigenstate, analytical formulas for the bifurcated eigenvalues will be derived. These formulas turn out to be exactly the same as those obtained by the KSGO method , but the present method is much simpler.
The bifurcated eigenvalues are of two types. One type consists of a pair of eigenvalues which bifurcate from the position-related zero eigenvalue. At infinite soliton separation, these eigenstates are equal to the sum of two position-induced Goldstein modes separated infinitely apart. The other type consists of two pairs of eigenvalues which bifurcate from the phase-related zero eigenvalue. At infinite soliton separation, these eigenstates are equal to the sum of two phase-induced Goldstein modes separated infinitely apart. These two types of eigenvalues have their counterparts in the linear stability analysis by the KSGO method .
Below, we consider eigenvalue bifurcations in type-I and II bound states. It turns out that eigenvalues of type-II bound states are simply equal to those of type-I states multiplied by $`i`$ (see also ). Thus calculations for only type-I states will be presented. These states are anti-symmetric in $`r`$, and symmetric in $`R`$, around the center of the bound states \[see Fig. 1(a)\]. When $`\mathrm{\Delta }\mathrm{}`$,
$$r(x)r_0(x)r_0(x\mathrm{\Delta }),R(x)R_0(x)+R_0(x\mathrm{\Delta }).$$
(3.6)
Bifurcations of position- and phase-related eigenvalues are studied separately next.
### 3.1 Bifurcation of position-related eigenvalues
When $`\mathrm{\Delta }\mathrm{}`$, the eigenstate bifurcated from the position-related zero eigenvalue is a sum of two position-induced Goldstein modes of fundamental vector solitons:
$$\mathrm{\Psi }(x)\mathrm{\Psi }_0(x)+\widehat{\mathrm{\Psi }}_0(x),\mathrm{\Delta }\mathrm{},$$
(3.7)
where
$$\mathrm{\Psi }_0(x)=\left[\begin{array}{c}r_0^{}(x)\\ r_0^{}(x)\\ R_0^{}(x)\\ R_0^{}(x)\end{array}\right],\widehat{\mathrm{\Psi }}_0(x)=\left[\begin{array}{c}r_0^{}(x\mathrm{\Delta })\\ r_0^{}(x\mathrm{\Delta })\\ R_0^{}(x\mathrm{\Delta })\\ R_0^{}(x\mathrm{\Delta })\end{array}\right].$$
(3.8)
The first two components of $`\mathrm{\Psi }`$ are anti-symmetric around $`x=\frac{1}{2}\mathrm{\Delta }`$, and the last two components are symmetric around $`x=\frac{1}{2}\mathrm{\Delta }`$. Note that in the same limit, the eigenstate $`\mathrm{\Psi }_0(x)\widehat{\mathrm{\Psi }}_0(x)`$ is simply $`[r^{}(x),r^{}(x),R^{}(x),R^{}(x)]^\text{T}`$, which is the un-bifurcated Goldstein eigenmode with eigenvalue zero and is thus not considered.
In the limit $`\mathrm{\Delta }1`$, we consider the bifurcated eigenstate in the region $`\mathrm{}<x\mathrm{\Delta }`$, and expand it as a perturbation series in powers of the small eigenvalue $`\lambda `$:
$$\mathrm{\Psi }(x)=\mathrm{\Psi }_0(x)+\lambda \mathrm{\Psi }_1(x)+\lambda ^2\mathrm{\Psi }_2(x)+o(\lambda ^2),\mathrm{}<x\mathrm{\Delta }.$$
(3.9)
In this region, we also expand
$$=_0+\lambda ^2_2+o(\lambda ^2).$$
(3.10)
It is noted that when $`\mathrm{\Delta }1`$, single solitons in the bound state considered are wave-and-daughter-waves whose two components have different orders (one component asymptotically much smaller than the other). This fact certainly has implications in the stability analysis. In particular, different components of $`\mathrm{\Psi }`$ and $``$ in the region $`\mathrm{}<x\mathrm{\Delta }`$ may have slightly different orders of magnitude. Thus, a perturbation expansion with a more-detailed ordering of different components than that in (3.9) and (3.10) might be needed. But as we will see next, (3.9) and (3.10) still work.
Now we substitute the perturbation expansions (3.9) and (3.10) into the eigenvalue equation (3.3). At $`O(1)`$, we get $`_0\mathrm{\Psi }_0(x)=0`$ which is satisfied automatically. At $`O(\lambda )`$, we get
$$_0\mathrm{\Psi }_1=\mathrm{\Psi }_0,$$
(3.11)
whose solution is
$$\mathrm{\Psi }_1(x)=[\frac{1}{2}xr_0(x),\frac{1}{2}xr_0(x),\frac{1}{2}xR_0(x),\frac{1}{2}xR_0(x)]^\text{T}.$$
(3.12)
Note that this function is an inhomogeneous solution of Eq. (3.11). In general, $`\mathrm{\Psi }_1`$ should also include the homogeneous solutions which are a linear combination of the three Goldstein modes: $`\mathrm{\Psi }_0(x)`$ in (3.8), and \[$`\mathrm{\Phi }_0^{(1)}(x)`$, $`\mathrm{\Phi }_0^{(2)}(x)`$\] in (3.26). But the $`\mathrm{\Psi }_0(x)`$ term in $`\mathrm{\Psi }_1`$ can be combined with the $`O(1)`$ term in the expansion (3.9), and the other two Goldstein modes \[$`\mathrm{\Phi }_0^{(1)}(x)`$, $`\mathrm{\Phi }_0^{(2)}(x)`$\] are phase-related and do not arise here. Thus the solution of Eq. (3.11) can be taken as (3.12) without loss of generality.
At $`O(\lambda ^2)`$, we get
$$_0\mathrm{\Psi }_2=\mathrm{\Psi }_1_2\mathrm{\Psi }_0.$$
(3.13)
It is more convenient to express $`_2\mathrm{\Psi }_0`$ in a different form. Recall that the un-bifurcated position-related Goldstein mode $`\mathrm{\Psi }_g(x)`$ of $``$ has the leading-order asymptotics $`\mathrm{\Psi }_0(x)\widehat{\mathrm{\Psi }}_0(x)`$ for $`\mathrm{\Delta }1`$. When this asymptotics and $``$โs expansion (3.10) are substituted into the Goldstein-mode relation $`\mathrm{\Psi }_g(x)=0`$, we find that asymptotically,
$$\lambda ^2_2\mathrm{\Psi }_0=_0\widehat{\mathrm{\Psi }}_0.$$
(3.14)
Thus Eq. (3.13) can be rewritten as
$$_0\left[\mathrm{\Psi }_2+\lambda ^2\widehat{\mathrm{\Psi }}_0\right]=\mathrm{\Psi }_1.$$
(3.15)
Note that $`_0\widehat{\mathrm{\Psi }}_0`$ is $`O(\lambda ^2)`$, thus the above equation is not dis-ordered.
The solvability condition of Eq. (3.15) will produce formulas for the eigenvalue $`\lambda `$. To do this, we need the asymptotics of function $`\mathrm{\Psi }_2`$ in the region $`x\frac{1}{2}\mathrm{\Delta }1`$, which we derive using the tail-matching idea. We have known that $`\mathrm{\Psi }`$โs leading-order asymptotics at $`\mathrm{\Delta }1`$ is given by (3.7) for all values of $`x`$. Combining this asymptotics with the perturbation expansion (3.9) and solutions (3.8) and (3.12), we see that
$$\mathrm{\Psi }_2(x)\frac{1}{\lambda ^2}\widehat{\mathrm{\Psi }}_0(x),x\frac{1}{2}\mathrm{\Delta }1.$$
(3.16)
Of course, $`\mathrm{\Psi }_2(x)0`$ when $`x\mathrm{}`$.
The homogeneous equation of (3.15) has three linearly independent solutions which are the Goldstein modes $`\mathrm{\Psi }_0(x)`$ in (3.8) and \[$`\mathrm{\Phi }_0^{(1)}(x),\mathrm{\Phi }_0^{(2)}(x)`$\] in (3.26). Thus Eq. (3.15) has three solvability conditions. These conditions can be readily derived by noting that diag$`(1,1,1,1)_0`$ is self-adjoint. It turns out that two of the solvability conditions induced by the Goldstein modes \[$`\mathrm{\Phi }_0^{(1)},\mathrm{\Phi }_0^{(2)}`$\] are satisfied automatically. The remaining solvability condition reads,
$`{\displaystyle _{\mathrm{}}^{x_m}}\mathrm{\Psi }_0|\text{diag}(1,1,1,1)|_0(\mathrm{\Psi }_2+\lambda ^2\widehat{\mathrm{\Psi }}_0)๐x`$
$`={\displaystyle _{\mathrm{}}^{x_m}}\mathrm{\Psi }_0|\text{diag}(1,1,1,1)|\mathrm{\Psi }_1๐x,`$ (3.17)
where $`x_m\frac{1}{2}\mathrm{\Delta }`$, $`|`$ and $`|`$ are the Dirac ket and bra notations . Integrating by parts to its left-hand-side and simplifying its right-hand-side, we get
$$\left[\mathrm{\Psi }_0|(\mathrm{\Psi }_{2x}+\lambda ^2\widehat{\mathrm{\Psi }}_{0x})\mathrm{\Psi }_{0x}|(\mathrm{\Psi }_2+\lambda ^2\widehat{\mathrm{\Psi }}_0)\right]_{\mathrm{}}^{x_m}=\frac{1}{2}_{\mathrm{}}^{x_m}(r_0^2+R_0^2)๐x.$$
(3.18)
The left-hand-side of (3.18) can be calculated using the asymptotics (2.7), (3.16) and equation (3.8), while the integral on the right-hand-side of (3.18) is asymptotically equal to a similar integral but with the upper limit $`x_m`$ replaced by $`\mathrm{}`$. After these calculations, the eigenvalue $`\lambda `$ is finally found to be
$$\lambda ^2=\frac{16(1\mathrm{\Omega })c^2e^\mathrm{\Delta }}{M+N},$$
(3.19)
where $`M`$ and $`N`$ are the masses of the $`r`$ and $`R`$ components in a fundamental vector soliton:
$$M_{\mathrm{}}^{\mathrm{}}r_0^2๐x,N_{\mathrm{}}^{\mathrm{}}R_0^2๐x.$$
(3.20)
We immediately see that formula (3.19) is identical to the one derived in using the KSGO method. Thus, the present tail-matching method has the same accuracy as the KSGO method, but is only simpler. Recall that we only consider the $`\mathrm{\Omega }\mathrm{\Omega }_c(<1)`$ limit. In this limit, the type-I state flips sign in its larger component, and does not flip sign in its smaller component \[see Fig. 1(a)\]. According to formula (3.19), this position-related eigenvalue $`\lambda `$ is real, thus it does not create instability. A comparison between the analytical formula (3.19) and numerical values at $`\beta =\frac{2}{3}`$ has been made in , and excellent agreement has been obtained.
Formula (3.19) shows that $`\lambda =O(e^{\frac{1}{2}\mathrm{\Delta }})`$. When $`\mathrm{\Omega }\mathrm{\Omega }_c`$ and $`0<\beta <1`$, $`\mathrm{\Delta }\mathrm{}`$ (see Sec. 2). In this limit, the asymptotic formula for $`\lambda `$ can be obtained more explicitly from (2.24) and (3.19) as
$$\lambda ^2\alpha (\mathrm{\Omega }\mathrm{\Omega }_c)^{\frac{1}{1\mathrm{\Omega }_c}},\mathrm{\Omega }\mathrm{\Omega }_c,$$
(3.21)
where the constant $`\alpha `$ is
$$\alpha =32(1\mathrm{\Omega }_c)\left(\frac{\mathrm{\Omega }_c^24^{\mathrm{\Omega }_c}}{8\gamma }\right)^{\frac{1}{1\mathrm{\Omega }_c}}.$$
(3.22)
Here $`\mathrm{\Omega }_c`$ and $`\gamma `$ are defined in Eqs. (2.6) and (2.25).
### 3.2 Bifurcation of phase-related eigenvalues
Calculations for the bifurcation of phase-related eigenvalues are quite similar to that done above. In this case, the eigenmode has the following asymptotics:
$$\mathrm{\Psi }(x)\mathrm{\Phi }_0(x)+\widehat{\mathrm{\Phi }}_0(x),\mathrm{\Delta }\mathrm{},$$
(3.23)
where
$$\mathrm{\Phi }_0(x)=\mathrm{\Phi }_0^{(1)}(x)+\delta \mathrm{\Phi }_0^{(2)}(x),$$
(3.24)
$$\widehat{\mathrm{\Phi }}_0(x)=\mathrm{\Phi }_0^{(1)}(x\mathrm{\Delta })\delta \mathrm{\Phi }_0^{(2)}(x\mathrm{\Delta }),$$
(3.25)
$$\mathrm{\Phi }_0^{(1)}(x)=\left[\begin{array}{c}r_0(x)\\ r_0(x)\\ 0\\ 0\end{array}\right],\mathrm{\Phi }_0^{(2)}(x)=\left[\begin{array}{c}0\\ 0\\ R_0(x)\\ R_0(x)\end{array}\right],$$
(3.26)
and $`\delta `$ is some constant (to be determined). The first two components of this mode are symmetric around $`x=\frac{1}{2}\mathrm{\Delta }`$, and the last two components are anti-symmetric around $`x=\frac{1}{2}\mathrm{\Delta }`$. Note that the eigenstate with asymptotics $`\mathrm{\Phi }_0(x)\widehat{\mathrm{\Phi }}_0(x)`$ is simply $`[r(x),r(x),\delta R(x),\delta R(x)]^\text{T}`$, which is the un-bifurcated Goldstein eigenmode with eigenvalue zero and is not a concern.
Next, we construct a perturbation-series solution for $`\mathrm{\Psi }(x)`$ in the region $`\mathrm{}<x\mathrm{\Delta }`$. The perturbation series is similar to (3.9):
$$\mathrm{\Psi }(x)=\mathrm{\Phi }_0(x)+\lambda \mathrm{\Phi }_1(x)+\lambda ^2\mathrm{\Phi }_2(x)+o(\lambda ^2),\mathrm{}<x\mathrm{\Delta },$$
(3.27)
where $`\mathrm{\Phi }_0(x)`$ is given by (3.24). Substituting this expansion and (3.10) into the eigenvalue relation (3.3), the $`O(1)`$ equation is satisfied automatically. At $`O(\lambda )`$, we get
$$_0\mathrm{\Phi }_1=\mathrm{\Phi }_0,$$
(3.28)
whose solution is
$$\mathrm{\Phi }_1(x)=\frac{1}{2}\frac{}{\omega }\left[\begin{array}{c}r_0(x;\omega ,\mathrm{\Omega })\\ r_0(x;\omega ,\mathrm{\Omega })\\ R_0(x;\omega ,\mathrm{\Omega })\\ R_0(x;\omega ,\mathrm{\Omega })\end{array}\right]_{\omega =1}+\frac{\delta }{2\mathrm{\Omega }}\frac{}{\mathrm{\Omega }}\left[\begin{array}{c}r_0(x;\omega ,\mathrm{\Omega })\\ r_0(x;\omega ,\mathrm{\Omega })\\ R_0(x;\omega ,\mathrm{\Omega })\\ R_0(x;\omega ,\mathrm{\Omega })\end{array}\right]_{\omega =1}.$$
(3.29)
Here $`(r_0,R_0)`$ is the fundamental vector soliton obtained from ODEs (2.4) and (2.5) without rescaling of $`\omega =1`$. Again, the homogeneous solution of Eq. (3.28), which is a linear combination of Goldstein modes $`\mathrm{\Psi }_0(x)`$ in (3.8) and \[$`\mathrm{\Phi }_0^{(1)}(x),\mathrm{\Phi }_0^{(2)}(x)`$\] in (3.26), is not included because the latter modes can be absorbed into the $`O(1)`$ term in the perturbation expansion (3.27), and the former mode does not arise.
At $`O(\lambda ^2)`$, we get
$$_0\mathrm{\Phi }_2=\mathrm{\Phi }_1_2\mathrm{\Phi }_0.$$
(3.30)
Again, utilizing the un-bifurcated phase-related Goldstein modes of $``$ with asymptotics $`\mathrm{\Phi }_0(x)\widehat{\mathrm{\Phi }}_0(x)`$ โ similar to what we have done in Sec. 3.1, we can rewrite $`_2\mathrm{\Phi }_0`$ so that Eq. (3.30) becomes
$$_0\left[\mathrm{\Phi }_2+\lambda ^2\widehat{\mathrm{\Phi }}_0\right]=\mathrm{\Phi }_1.$$
(3.31)
This equation has three solvability conditions induced by the three Goldstein modes in the homogeneous solution. The condition induced by the mode $`\mathrm{\Psi }_0(x)`$ in (3.8) is satisfied automatically. The other two conditions are
$`{\displaystyle _{\mathrm{}}^{x_m}}\mathrm{\Phi }_0^{(k)}|\text{diag}(1,1,1,1)|_0(\mathrm{\Phi }_2+\lambda ^2\widehat{\mathrm{\Phi }}_0)๐x`$
$`={\displaystyle _{\mathrm{}}^{x_m}}\mathrm{\Phi }_0^{(k)}|\text{diag}(1,1,1,1)|\mathrm{\Phi }_1๐x,`$ (3.32)
where $`x_m\frac{1}{2}\mathrm{\Delta }`$, and $`k=1,2`$. Integration by parts simplifies these conditions as
$$\left[\mathrm{\Phi }_0^{(1)}|(\mathrm{\Phi }_{2x}+\lambda ^2\widehat{\mathrm{\Phi }}_{0x})\mathrm{\Phi }_{0x}^{(1)}|(\mathrm{\Phi }_2+\lambda ^2\widehat{\mathrm{\Phi }}_0)\right]_{\mathrm{}}^{x_m}=\frac{1}{2}M_\omega +\frac{\delta }{2\mathrm{\Omega }}M_\mathrm{\Omega },$$
(3.33)
and
$$\left[\mathrm{\Phi }_0^{(2)}|(\mathrm{\Phi }_{2x}+\lambda ^2\widehat{\mathrm{\Phi }}_{0x})\mathrm{\Phi }_{0x}^{(2)}|(\mathrm{\Phi }_2+\lambda ^2\widehat{\mathrm{\Phi }}_0)\right]_{\mathrm{}}^{x_m}=\frac{1}{2}N_\omega +\frac{\delta }{2\mathrm{\Omega }}N_\mathrm{\Omega }.$$
(3.34)
To calculate the left-hand-sides of the above two equations, we need the asymptotics of function $`\mathrm{\Phi }_2(x)`$ in the region $`x\frac{1}{2}\mathrm{\Delta }`$. This asymptotics can be obtained by comparing $`\mathrm{\Psi }(x)`$โs asymptotics (3.23) with its perturbation expansion (3.27). This comparison shows that $`\mathrm{\Phi }_2(x)`$ must have the asymptotics
$$\mathrm{\Phi }_2(x)\frac{1}{\lambda ^2}\widehat{\mathrm{\Phi }}_0(x),x\frac{1}{2}\mathrm{\Delta }1.$$
(3.35)
When this asymptotics as well as (2.7) is substituted into Eqs. (3.33) and (3.34) and parameter $`\delta `$ eliminated, the eigenvalue $`\lambda `$ is then given by the quartic equation
$$\lambda ^4\frac{16c^2(N_\mathrm{\Omega }M_\omega )e^\mathrm{\Delta }}{M_\omega N_\mathrm{\Omega }M_\mathrm{\Omega }N_\omega }\lambda ^2\frac{16^2c^4e^{2\mathrm{\Delta }}}{M_\omega N_\mathrm{\Omega }M_\mathrm{\Omega }N_\omega }=0.$$
(3.36)
Again, this formula for phase-related eigenvalues is identical to that obtained in by the KSGO method. As pointed in , this formula shows that a two-vector-soliton bound state always has one phase-related unstable eigenvalue which bifurcates from zero, thus is always linearly unstable. Comparison between this formula and numerical values has also been made in , and excellent agreement has been observed.
Formula (3.36) shows that phase-related eigenvalues $`\lambda =O(e^{\frac{1}{2}\mathrm{\Delta }})`$, the same as position-related eigenvalues \[see Eq. (3.19)\]. The asymptotics of these eigenvalues in the limit $`\mathrm{\Omega }\mathrm{\Omega }_c`$ can also be obtained from (2.24) and (3.36), but is not pursued here.
Next, we briefly discuss the linear stability of type-II vector-soliton bound states \[see Fig. 1(b)\]. In the limit of large separation, these solitons have the asymptotics
$$r(x)r_0(x)+r_0(x\mathrm{\Delta }),R(x)R_0(x)R_0(x\mathrm{\Delta }).$$
(3.37)
Repeating the above analytical calculations, we find that the eigenvalues for type-II states are equal to those of type-I states multiplied by $`i`$. Thus type-II states are always linearly unstable as well. But different from type-I states, the instability of type-II states in the limit $`\mathrm{\Omega }\mathrm{\Omega }_c(<1)`$ is caused by two unstable eigenvalues: one related to position-mode bifurcations \[see (3.19)\], and the other one related to phase-mode bifurcations \[see (3.36)\]. This result agrees with that by the KSGO method as well as the numerical result .
## 4 Discussion
In this paper, we have analytically studied the linear stability of two-vector-soliton bound states in the coupled NLS equations by a new tail-matching method. Under the condition that the two vector solitons are wave-and-daughter-waves and are widely separated, we have calculated small eigenvalues of these bound states that bifurcate from the zero eigenvalue of single vector solitons. These small eigenvalues calculated are all the discrete non-zero eigenvalues of the bound states. We have shown that these bound states are always linearly unstable due to the existence of one unstable phase-induced eigenvalue. The analytical formulas for eigenvalues derived from this tail-matching method turn out to be exactly the same as those from the KSGO method , but the present method is much simpler. Even though our calculations were performed for two-vector-soliton bound states, these calculations can apparently be generalized to $`n`$-vector-soliton bound states.
This tail-matching method for the linear stability analysis of multi-soliton bound states is apparently quite general, and it can be applied to a wide range of other wave systems where similar multi-pulse bound states have been reported . In addition, this method should be applicable to the stability analysis of other types of multi-pulses from nonlocal bifurcations in coupled-NLS-type equations . Preliminary linear stability results through numerical studies shows that such multi-pulses are also linearly unstable , consistent with the results of the present article. This tail-matching method can also be used to calculate eigenvalue bifurcations of multi-pulses from internal modes (non-zero eigenvalues) of single pulses โ a bifurcation which incidentally does not arise in the present problem because internal modes do not exist in single wave-and-daughter-waves .
Under what conditions can this tail-matching method give useful results? The main condition is that the individual pulses in the multi-pulse bound state are widely separated. This condition will dictate in what parameter regions such multi-pulses can exist, and thus tail-matching can proceed. For instance, for the coupled NLS equations (2.1) and (2.2), the spacing formulas (2.10) and (2.24) dictate that widely-separated multi-vector-soliton bound states ($`\mathrm{\Delta }1`$) exist when $`\mathrm{\Omega }`$ is near the local bifurcation boundary $`\mathrm{\Omega }=\mathrm{\Omega }_c`$. This is precisely the parameter region where our tail-matching linear stability analysis is performed.
The two-vector-soliton bound states studied in this paper reside outside the continuous spectrum of the corresponding linear-wave system. It is known that multi-pulse embedded solitons residing inside the continuous spectrum exist in various wave systems as well . An interesting open issue is whether the tail-matching method can also be applied to the linear stability analysis of such multi-pulse embedded solitons. In the third-order NLS equation, it has been shown numerically that such embedded solitons are all linearly stable, but nonlinearly semi-stable .
Lastly, we relate this tail-matching method to other existing techniques for the linear stability analysis of multi-pulse bound states. Currently, the following techniques exist: the KSGO method , the dynamical-system method , and the effective-interaction-potential method . The present method is asymptotically accurate. It gives the same results as the KSGO method, but is much simpler. The dynamical-system method can count the number of unstable eigenvalues, or express the eigenvalues as the zeros of the Evans function. But it generally does not produce explicit formulas for eigenvalues. The effective-potential method can only capture position-related eigenvalue bifurcations, not phase-related eigenvalue bifurcations. (For the coupled NLS equations, phase-related eigenvalues are more important). In view of this comparison, we feel that the tail-matching method for the linear stability of multi-soliton bound states is very promising.
Department of Mathematics and Statistics
University of Vermont
Burlington, VT 05401, USA
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# Untitled Document
A Class of Cosmological Matrix Models
Miao Li
Institute of Theoretical Physics
Academia Sinica, P.O. Box 2735
Beijing 100080, P. R. China
and
Interdisciplinary Center of Theoretical Studies
Academia Sinica
mli@itp.ac.cn
We discuss a class of matrix models describing cosmology with a light-like singularity, generalizing the model proposed by Craps et al. in hep-th/0506180.
June 2005
Craps et al. recently proposed to study a simple cosmological background with a null singularity , this cosmology admits a matrix model description, thus lends itself to a rigorous study. An earlier example of cosmology with a null singularity is proposed in \[2,,3\] and is later studied by many authors \[4--6\], they find that this singularity is highly unstable due to gravitational back-reaction. However, the model of Craps et al. seems to avoid this problem.
The construction of the model in is simple. One starts in type IIA string theory with a flat string metric and a linear dilaton background $`\varphi =Qx^+`$, where $`x^+`$ is a light-like coordinate. The linear dilaton does not require modifying the critical dimension, since its linearity is along a null direction. Although the string metric is flat, the Einstein metric is nontrivial:
$$ds^2=e^{Qx^+/2}(2dx^+dx^{}+(dx^i)^2).$$
For a positive $`Q`$, the metric contracts to a singularity at $`x^+=\mathrm{}`$, this is actually a curvature singularity. The corresponding 11 dimensional M theory metric also has a singularity at $`x^+=\mathrm{}`$. The authors of show that the singularity at $`x^+=\mathrm{}`$ lies in a finite geodesic distance, while the other singularity is at infinite distance. This background preserves half of total 32 supersymmetries, one expects that there is a control over the null singularity.
Since the string metric is flat, this time-dependent background admits a matrix string description, in which the Yang-Mills coupling constant is time-dependent. In fact, the authors of show that this Yang-Mills theory can be regarded as one with a constant coupling on a world sheet with a time-dependent metric.
It goes without saying that this is an important observation, and it may provide the first example in which we can study a time-dependent background in a controlled fashion. Thus, it is interesting to ask whether this model is unique, or it has many cousins. In this note we will see that indeed there is a large class of such models.
We shall limit ourselves in M theory in this note. Consider metric
$$ds^2=e^{2\alpha x^+}(2dx^+dx^{}+(dx^i)^2)+e^{2\beta x^+}(dx^a)^2,$$
where $`i=2,\mathrm{}10d`$, $`a=11d,\mathrm{}10`$, namely, there are $`d`$ coordinates $`x^a`$, the total dimensions of spacetime is 11. The metric of is given by taking $`d=1`$, $`\alpha =Q/3`$ and $`\beta =2Q/3`$. $`X^{10}`$ is taken as the M theory circle, thus the metric reduces to the IIA flat string metric with a linear dilaton background.
To find a solution, we use an orthonormal basis
$$e^\pm =e^{\alpha x^+}dx^\pm ,e^i=e^{\alpha x^+}dx^1,e^a=e^{\beta x^+}dx^a.$$
The non-vanishing components of spin connection are
$$\omega _+=\alpha e^{\alpha x^+}e^+,\omega _{i+}=\alpha e^{\alpha x^+}e^i,\omega _{a+}=\beta e^{\alpha x^+}e^a.$$
The non-vanishing curvature 2-forms are
$$\begin{array}{cc}\hfill R_{i+}& =\alpha ^2e^{2\alpha x^+}e^ie^+,\hfill \\ \hfill R_{a+}& =\beta (2\alpha \beta )e^{2\alpha x^+}e^ae^+.\hfill \end{array}$$
The only non-vanishing component of the Ricci tensor is $`R_{++}`$ given by
$$R_{++}=[(9d)\alpha ^2+d\beta (2\alpha \beta )]e^{2\alpha x^+},$$
and the Einstein equation $`R_{++}=0`$ has two solutions
$$\beta =(1\pm \frac{3}{\sqrt{d}})\alpha .$$
For $`d=1`$, choosing the minus sign in (1) reproduces the background considered in .
We now show that the background given by eqs.(1) and (1) not only is a solution, but also preserves half of supersymmetries. The only interesting SUSY transformation is that for the gravitino $`\delta \mathrm{\Psi }_\mu =D_\mu ฯต`$. The components of spin connection of interest are $`\omega _+`$, $`\omega _i`$ and $`\omega _a`$:
$$\begin{array}{cc}\hfill \omega _+& =2\alpha 2\alpha \gamma ^{}\gamma ^+,\omega _i=2\alpha \gamma ^i\gamma ^+,\hfill \\ \hfill \omega _a& =2\beta e^{(\beta \alpha )x^+}\gamma ^a\gamma ^+.\hfill \end{array}$$
Since
$$[D_+,D_a]=\frac{1}{2}\beta (\beta 2\alpha )e^{(\beta \alpha )x^+}\gamma ^a\gamma ^+,$$
The compatibility condition for the constraints $`D_\mu ฯต=0`$ is $`\gamma ^+ฯต=0`$. This condition eliminates half of the components in $`ฯต`$. The constraints $`D_iฯต=D_aฯต=0`$ are solved if $`ฯต`$ is independent of $`x^i`$ and $`x^a`$. Finally, $`D_+ฯต=0`$ is solved if $`ฯต=\mathrm{exp}(\frac{1}{2}\alpha x^+)\eta `$ for a constant $`\eta `$. We conclude that the unbroken supersymmetry is parametrized by $`\eta `$ with the constraint $`\gamma ^+\eta =0`$. Thus, just like the original background of , our more general metric also preserves 16 supersymmtries, this should be enough to guarantee that our proposed matrix model to be described shortly is a valid description of dynamics over this background.
In the case $`d=9`$, $`\beta =0,2\alpha `$, although $`[D_+,D_a]=0`$, but $`[D_+,D_i]=\frac{1}{2}\alpha ^2\gamma ^i\gamma ^+`$, again we obtain the same unbroken supersymmetries.
In fact, (1) is only a special case of a more general class of solutions preserving half of supersymmetries (to be discussed in the end of this note). Before we study the matrix model for (1), let us discuss the geometric properties of (1). Apparently, the sign of $`\alpha `$ and its absolute value can be changed by changing $`x^+`$ and $`x^{}`$, so we always assume $`\alpha >0`$. As $`x^+\mathrm{}`$, the factor $`e^{2\alpha x^+}`$ approaches zero, the transverse dimensions $`x^i`$ shrink to zero size, this is the big bang point with regard to these coordinates. This singularity locates a finite distance away in view of the affine parameter, indeed, let $`dX^+=e^{2\alpha x^+}dx^+`$, the first term in the metric (1) becomes
$$2dX^+dx^{},$$
thus $`X^+`$ is the affine parameter for the null geodesic $`x^{}=`$const. The big bang singularity occurs at $`X^+=0`$.
The geometry of other transverse dimensions $`x^a`$ depends on the choice $`\beta `$ in (1). We name the choice $`\beta =(1+3/\sqrt{d})\alpha `$ case 1, and the choice $`\beta =(13/\sqrt{d})\alpha `$ case 2. For case 1, $`\beta >0`$, so all $`x^a`$ shrink at $`x^+=\mathrm{}`$. For case 2, $`\beta <0`$ except for the $`d=9`$ case, thus all $`x^a`$ get to infinity at $`x^+=\mathrm{}`$ and shrink to zero size at $`x^+=\mathrm{}`$. $`d=9`$ is special, in this case, all 9 dimensions $`x^a`$ do not evolve in time, and one can redefine $`x^+`$ such that the metric is Minkowski. Supersymmetry is also enhanced, since there is no longer constraint $`\gamma ^+ฯต=0`$ following from (1).
The model studied in is a special example of case 2 when $`d=1`$. As in , we can compactify $`x^{10}`$ on a circle to obtain type IIA string theory. The string metric and the dilaton $`\varphi `$ are related to the M theory metric through
$$ds^2=e^{2\varphi /3}ds_{st}^2+e^{4\varphi /3}(dx^{10})^2,$$
we obtain
$$\varphi =\frac{3}{2}\beta x^+,ds_{st}^2=e^{(2\alpha +\beta )x^+}[2dx^+dx^{}+(dx^i)^2]+e^{3\beta x^+}(dx^a)^2.$$
The 10D Einstein metric reads
$$ds_E^2=e^{(2\alpha +\beta /4)x^+}[2dx^+dx^{}+(dx^i)^2]+e^{(9\beta /4)x^+}(dx^a)^2.$$
Since the string coupling โconstantโ $`g_s=\mathrm{exp}(\varphi )=\mathrm{exp}(\frac{3}{2}\beta x^+)`$, for case 1, strings are weakly coupled at big bang, while all transverse dimensions are zero-sized, so it is not clear whether we can trust the perturbative string theory. As we shall see shortly in the matrix model, spacetime breaks down and we need to employ a full non-abelian description. At later times, strings become strongly coupled, while all dimensions expand, we shall see that a rather simple theory emerges, that is, only abelian degrees of freedom survive.
For case 2, strings are strongly coupled at big bang. transverse dimensions $`x^i`$ start with a zero size, while transverse dimensions $`x^a`$ have infinite size and contract as time evolves. This picture in valid both in terms of the string metric as well as the 10D Einstein metric. We will study in more detail the spacetime properties in the matrix model later.
The string frame metric (1) is in general non-flat. $`\beta `$ never vanishes, so to get a flat metric, $`d=1`$ is necessary, and in this case one chooses $`\beta =2\alpha `$, this is the special case considered in . More generally, one obtains a world-sheet theory with an action explicitly depending on time. One can attempt to quantize the string in the light-cone gauge, to do so, introduce the new light-cone coordinate $`dy^+=\mathrm{exp}((2\alpha +\beta )x^+)dx^+`$, the first term in the string frame metric (1) becomes $`2dy^+dx^{}`$. The lignt-cone momentum $`p_{}=p^+`$ conjugate to $`x^{}`$ is conserved. In the light-cone quantization, we can set the light-cone gauge $`y^+=\tau `$, the bosonic part of the world-sheet action reads
$$S=\frac{1}{4\pi \alpha ^{}}๐\tau ๐\sigma [e^{(2\alpha +\beta )x^+}^\alpha X^i_\alpha X^i+e^{3\beta x^+}^\alpha X^a_\alpha X^a],$$
where the period of $`\sigma `$ is $`2\pi \alpha ^{}p^+`$, and $`\mathrm{exp}((2\alpha +\beta )x^+)=(2\alpha +\beta )\tau `$. As $`x^+\mathrm{}`$, $`\tau 0`$. Since $`\tau `$ starts at a finite point, it is more useful to use the old light-cone coordinate $`x^+=t`$, the action becomes
$$S=\frac{1}{4\pi \alpha ^{}}๐t๐\sigma [e^{(4\alpha +2\beta )t}^\alpha X^i_\alpha X^i+e^{(2\alpha +4\beta )t}^\alpha X^a_\alpha X^a].$$
The time-dependent coefficients may be interpreted as the effective tensions. For transverse coordinates $`X^i`$, the effective tension is $`T_i=\frac{1}{2\pi \alpha ^{}}e^{(4\alpha +2\beta )t}`$, for coordinates $`X^a`$, the effective tension is $`T_a=\frac{1}{2\pi \alpha ^{}}e^{(2\alpha +4\beta )t}`$. In case 1, both effective tensions get to zero at big bang, while the string coupling also gets to zero. Since the string spectrum becomes very dense, the usual free string picture does not apply. In later times, effective string tensions get large, and the string coupling constant also becomes strong, we will have to deal with a strongly coupled massless sector. In case 2, $`4\alpha +2\beta >0`$ except for $`d=1`$, our above analysis still applies to $`T_i`$. $`2\alpha +4\beta =6\alpha (12/\sqrt{d})`$, for $`d>4`$, $`T_a`$ becomes small in earlier time, the same as case 1. For $`d<4`$, $`T_a`$ becomes large in earlier time, thus these transeverse dimensions are effectively frozen. $`d=4`$ is a special case, $`T_a`$ is independent of time.
The vertex operator of a string state in the background (1) in general is quite complicated. For instance, consider a massless scalar satisfying equation of motion
$$_\mu \left(e^{2\varphi }\sqrt{g}g^{\mu \nu }_\nu \mathrm{\Phi }\right)=0.$$
As function of $`x^+`$, the scaling of components $`g^{ab}`$ is different from the scaling of components $`g^+`$ and $`g^{ij}`$, so there is no simple plane wave solution in general, which implies that the vertex operator of this massless scalar field is not simple. However, we can consider special cases when $`\mathrm{\Phi }`$ is independent of $`x^a`$, in this case, a vertex operator $`V=\mathrm{exp}(ik_+x^++ik_{}x^{}+ik_ix^i)`$ must satisfy the on-shell condition:
$$k_{}(2k_+i\gamma )k_i^2=0,\gamma =(9d)\alpha +d\beta ,$$
which require a imaginary part of $`k_+`$: $`\mathrm{}k_{}=\gamma /2`$, thus in the vertex operator, there is an exponential factor
$$e^{\frac{1}{2}\gamma x^+}.$$
Since $`\gamma =(9\pm 3\sqrt{d})\alpha 0`$, this exponential factor always blows up at $`x^+=\mathrm{}`$ for $`d<9`$. The effective string coupling constant for these states is $`g_{eff}=g_s\mathrm{exp}(\frac{1}{2}\gamma x^+)`$. Dimensions $`x^a`$ for these states are effectively compactified, thus when we discuss interactions among these states, the space-time dimensionality is $`11d`$, and there is an effective string tension in action (1), or $`\alpha _{eff}^{}=\alpha ^{}\mathrm{exp}((4\alpha +2\beta )x^+)`$. The effective Newton constant is $`G_{eff}=g_{eff}^2\alpha ^{(9d)/2}`$, as a function of $`x^+`$, it scales
$$G_{eff}e^{\gamma x^+(9d)(2\alpha +\beta )x^+}.$$
Now, $`\gamma `$ is positive and $`(9d)(2\alpha +\beta )`$ is non-negative if $`d<9`$, the effective Newton constant blows up at $`x^+=\mathrm{}`$, the perturbative sstring picture is not valid at least for these states.
Next, we study the matrix model. It is not necessary to compactify any dimension of (1). In the usual flat background, the Matrix Theory action reads
$$S=๐t\mathrm{Tr}\left(\frac{1}{2R}(D_tX^i)^2+\frac{R}{4}[X^i,X^j]^2+i\theta ^TD_t\theta R\theta ^T\gamma _i[X^i,\theta ]\right),$$
where $`R`$ is the longitudinal cut-off, or it may be viewed as the radius of $`x^{}`$ in the DLCQ M theory . For simplicity, we set the M theory Planck length $`l_p=1`$. A derivation of this matrix action is given in (see also ). With our metric (1), it is straightforward to write the corresponding matrix action. To do this, we need to use the light-cone coordinate in (1), and identify $`X^+`$ with $`\tau `$, the world-line time of D0-branes. In the matrix model, we apply the M theory metric (1) directly. Since the action is rather lengthy, we separate the action into the bosonic part and the fermionic part. The bosonic part reads
$$\begin{array}{cc}\hfill S_B& =d\tau \mathrm{Tr}\{\frac{1}{2R}e^{2\alpha x^+}(D_\tau X^i)^2+\frac{1}{2R}e^{2\beta x^+}(D_\tau X^a)^2+\frac{R}{4}e^{4\alpha x^+}[X^i,X^j]^4\hfill \\ & +\frac{R}{4}e^{4\beta x^+}[X^a,X^b]^4+\frac{R}{2}e^{(2\alpha +2\beta )x^+}[X^i,X^a]^4\}.\hfill \end{array}$$
Note that $`x^+`$ appearing in the above action is the old light-cone coordinate. The fermionic part reads
$$S_F=๐\tau \{i\theta ^TD_\tau \theta Re^{\alpha x^+}\theta ^T\gamma _i[X^i,\theta ]Re^{\beta x^+}\theta ^T\gamma _a[X^a,\theta ]\}.$$
It is rather awkward to use $`\tau `$ as time, since it has a finite beginning. Let us switch back to the coordinate $`x^+`$ and on the world-line identify $`t=x^+`$, thus $`d\tau =\mathrm{exp}(2\alpha t)dt`$, we have
$$\begin{array}{cc}\hfill S_B& =dt\mathrm{Tr}\{\frac{1}{2R}(D_tX^i)^2+\frac{1}{2R}e^{2(\beta \alpha )t}(D_tX^a)^2+\frac{R}{4}e^{6\alpha t}[X^i,X^j]^4\hfill \\ & +\frac{R}{4}e^{(2\alpha +4\beta )t}[X^a,X^b]^4+\frac{R}{2}e^{(4\alpha +2\beta )t}[X^i,X^a]^4\},\hfill \end{array}$$
and
$$S_F=๐t\{i\theta ^TD_t\theta Re^{3\alpha t}\theta ^T\gamma _i[X^i,\theta ]Re^{(2\alpha +\beta )t}\theta ^T\gamma _a[X^a,\theta ]\}.$$
Before study the simplest properties of this matrix model, let us show that we can recover the matrix string action of . In this case, there is only one $`X^a`$, call it $`X^{10}`$. The compactification scheme is given in . Up a dimensionful parameter, we replace the trace in the matrix action by $`๐\sigma \mathrm{Tr}`$, the commutator $`R[X^{10},X^i]`$ is replaced by the covariant derivative $`iD_\sigma X^i`$. Finally, use $`\beta =2\alpha `$, we find
$$S_B=๐t๐\sigma \mathrm{Tr}\{\frac{1}{2R}(D_\alpha X^i)^2+\frac{1}{2R^3}g_s^2F_{t\sigma }^2+\frac{R}{4}g_s^2[X^i,X^j]^2\},$$
and
$$S_F=dtd\sigma \mathrm{Tr}[\theta ^T\sigma ^\alpha D_\alpha \theta Rg_s^1\theta ^T\gamma _i[X^i,\theta ],$$
where $`g_s`$ is the time-dependent string coupling $`g_s=\mathrm{exp}(3\alpha t)`$. Our action is identical to that in up to an identification of a dimensionful parameter $`l_s`$. This matrix string theory can be regarded as a 2D Yang-Mills theory with a time-dependent coupling or a 2D Yang-Mills theory with a constant coupling in the world-sheet metric
$$ds^2=g_s^2(dt^2+d\sigma ^2).$$
To study the properties of matrix model defined by (1) and (1), we consider cases 1 and 2 separately.
$``$ Case 1, $`\beta =(1+\frac{3}{\sqrt{d}})\alpha `$.
The kinetic term of $`X^i`$ is always simple, as is the kinetic term of $`\theta `$. The coefficients of the remaining terms all vanish in the limit $`t\mathrm{}`$, so there is no constraint arising from these terms, this implies that all matrices are fully non-abelian. On the other hand, as $`t\mathrm{}`$, these coefficients blow up, thus, all matrices must commute with one another, the only surviving degrees of freedom are diagonal elements. Moreover, $`X^a`$ must be independent of time in this limit, so $`X^a`$ become frozen abelian moduli. Recall that in the string picture, we found previously that the effective tension $`T_a`$ becomes infinitely heavy, even heavier than $`T_i`$ for large $`t`$, this is related to the fact that $`X^a`$ become moduli in the matrix model, while $`X^i`$ still have dynamics. If we compactify some of the transverse dimensions, the story becomes slightly more involved. For instance, compactifying $`X^{10}`$ on a circle, all other matrices become function on a circle $`\sigma `$. In the limit $`t\mathrm{}`$, although they have to be diagonal, they are not always periodic functions of $`\sigma `$, the eigen-values can get permuted after circling along $`\sigma `$, these twisted sectors describe strings of various lengths.
$``$ Case 2, $`\beta =(1\frac{3}{\sqrt{d}})\alpha `$.
$`6\alpha `$ and $`4\alpha +2\beta `$ are always positive, at big bang, there is no constraint on commutators $`[X^i,X^j]`$, $`[X^i,X^a]`$, $`[X^i,\theta ]`$ and $`[X^a,\theta ]`$. On the other hand, as $`t\mathrm{}`$, these commutators are forced to vanish. $`2(\beta \alpha )`$ is always negative, so at bing bang, $`X^a`$ must be independent of time, they become non-abelian moduli in the model. $`2\alpha 4\beta =6\alpha (12/\sqrt{d})`$, $`d=4`$ becomes the critical dimension. For a larger $`d`$, there is no constraint on the commutators $`[X^a,X^b]`$ at big bang, and they have to vanish as $`t\mathrm{}`$, thus, for $`d>4`$, all matrices commute in this limit and only abelian degrees of freedom survive. For $`d<4`$, $`[X^a,X^b]`$ have to vanish at big bang, together with the fact that $`X^a`$ are independent of time, these matrices become abelian moduli at big bang, this fact is also reflected in our previous analysis in the string picture, where we found $`T_a`$ become infnitely heavy at big bang. When $`d=4`$, $`X^a`$ remain nonabelian all the time. In the string picture, recall that when $`d=4`$, the effective tension $`T_a`$ is constant.
It is also of interest to study compactification. Just as in the flat background, we do not know how to write down matrix model action if we compactify more than 5 dimensions. We may choose to compactify all $`X^a`$, or some of $`X^a`$, or some of $`X^a`$ and some of $`X^i`$. For simplicity, let us compactify all $`X^a`$ on a torus $`T^d`$. Replacing $`R[X^a,X^i]`$ by $`iD_aX^i`$, $`R^2[X^a,X^d]`$ by $`F_{ab}`$ etc, the bosonic part of the matrix model action reads
$$\begin{array}{cc}\hfill S_B& =d^{d+1}\sigma \{\frac{1}{2R}(D_tX^i)^2\frac{1}{2R}e^{(4\alpha +2\beta )t}(D_aX^i)^2+\frac{1}{2R^3}e^{2(\beta 2\alpha )t}F_{ta}^2\hfill \\ & +\frac{1}{4R}e^{(2\alpha +4\beta )t}F_{ab}^2+\frac{R}{4}e^{6\alpha t}[X^i,X^j]^2\}.\hfill \end{array}$$
The analysis of this action is similar to our previous analysis of the action without compactification, a statement there can be simply translated to one for the action (1). For instance, demanding $`[X^a,X^b]=0`$ is translated to $`F_{ab}=0`$, that is, the spatial connection must be flat.
It is interesting that the previous noticed โcritical dimensionโ $`d=4`$ corresponds to the situation that the action (1) is no longer complete. Compactification on a torus $`T^4`$, the complete matrix model is to be given by the world volume theory of coincident M5-branes . Of course, that theory will also be time-dependent, and hopefully one can figure out more details of the time dependence by studying field theory behavior of (1). Likewise, action (1) is again incomplete for compactification on $`T^5`$, in this case the complete theory is given by the little string theory .
Unlike the case when $`d=1`$, the Yang-Mills theory (1) can not be simply interpreted as a theory with a constant coupling on a world volume of a nontrivial metric. To interpret action (1) as Yang-Mills theory, we need to introduce both a nontrivial metric on the world volume as well as a time-dependent Yang-Mills coupling. Since a factor $`g_{YM}^2`$ appears in the coefficient of $`F_{\mu \nu }^2`$ and a factor $`g_{YM}^2`$ appears in the coefficient of $`[X^i,X^j]^2`$, we find for $`d>1`$
$$\begin{array}{cc}\hfill ds_{WV}^2& =e^{\frac{d}{d1}(4\alpha +2\beta )t}dt^2+e^{\frac{1}{d1}(4\alpha +2\beta )t}(d\sigma ^a)^2,\hfill \\ \hfill g_{YM}^2& =e^{\frac{1}{d1}(4\alpha +2\beta )t+2(\alpha \beta )t},\hfill \end{array}$$
where the subscript $`WV`$ stands for world volume. the world volume metric components $`g_{aa}`$ must be all same, there are only 3 independent functions in the geometry and the Yang-Mills coupling, however, we need to match 5 different kinds of coefficients in (1), so it is nontrivial that there is a solution as presented in (1). This is quite important especially for $`d>3`$, since we need to complete the action (1) by introducing more degrees of freedom, this is possible only when (1) admits Yang-Mills theory interpretation. One can check that the identification (1) also works for the fermionic part of the matrix action.
Finally, we discuss generalizations of background (1). There are two directions to generalize (1). The first is to consider
$$ds^2=2e^{2\alpha x^+}dx^+dx^{}+\underset{i}{}e^{2\beta _ix^+}(dx^i)^2.$$
The Einstein equations are solved provided
$$\alpha ^2+\underset{i}{}\beta _i(2\alpha \beta _i)=0.$$
This background again preserves 16 supersymmetries of the form $`ฯต=\mathrm{exp}(\frac{1}{2}\alpha x^+)\eta `$, with a constant $`\eta `$ satisfying $`\gamma ^+\eta =0`$.
The second direction to generalize (1) is to consider metric of the form
$$ds^2=e^{2f(x^+)}(2dx^+dx^{}+(dx^i)^2)+e^{2g(x^+)}(dx^a)^2.$$
The non-vanishing curvature 2 forms are
$$\begin{array}{cc}\hfill R_{i+}& =(f^2f^{\prime \prime })e^{2f}e^ie^+,\hfill \\ \hfill R_{a+}& =(2f^{}g^{}g^2g^{\prime \prime })e^{2f}e^ae^+.\hfill \end{array}$$
The Einstein equations are solved if
$$(9d)(f^2f^{\prime \prime })+d(2f^{}g^{}g^2g^{\prime \prime })=0.$$
Once again, there are 16 unbroken supersymmetries, parameterized by $`ฯต=\mathrm{exp}(\frac{1}{2}f)\eta `$, $`\gamma ^+\eta =0`$.
Needless to say that equation (1) admits infinite many solutions. Two special solutions deserve attention. One solution is when $`g=0`$, in this case $`f^2f^{\prime \prime }=0`$. All the curvature 2 forms vanish thus it may appear that we obtain a flat space solution. This is almost case except for a simple twist. The solution of $`f^2f^{\prime \prime }=0`$ is $`\mathrm{exp}(2f)=1/(x^+)^2`$, by changing coordinates $`dy^+=dx^+/(x^+)^2`$, $`y^{}=x^{}`$ we obtain a metric
$$ds^2=2dy^+dy^{}+(y^+)^2(dx^i)^2+(dx^a)^2.$$
If one of $`x^i`$ is periodic, we find a nontrivial background similar to the orbifold discussed in \[2,,3\]. Of course if all $`x^i`$ are noncompact, (1) is equivalent to the Minkowski space. Similarly, one may set $`f=0`$ and find $`\mathrm{exp}(2g)=(x^+)^2`$, a nontrivial orbifold results provided one of $`x^a`$ is compactified. One may attempt to construct a matrix model for the metric (1) too. However, time $`\tau =y^+`$ seems to have to terminate at $`\tau =0`$. One may try to avoid this problem by going back to the original light-cone coordinate $`x^+`$ and take $`t=x^+`$ in the matrix model, although this helps to eliminate the coefficient $`\tau ^2`$ in the kinetic term of $`X^i`$ but it also introduces a coefficient $`t^2`$ in the kinetic term of $`X^a`$, again $`t`$ must be terminated at $`t=0`$. Also the supersymmetry parameter $`ฯต=\sqrt{x^+}\eta `$ does not exist beyond $`x^+=0`$ since $`ฯต`$ has to be real. This problem may cause disease in the definition of the matrix model, and may be a reflection of the gravitational instability caused by a test particle.
Acknowledgments.
I am grateful to B. Chen, Q. G. Huang and W. Song for discussions. This work was supported by a grant of CNSF.
References
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# Parametric Autoresonance in Faraday waves
## I Introduction
There are two ways of driving a classical nonlinear oscillator by a small oscillating force: via either external or parametric resonance. In both cases, the initial growth of the amplitude of the oscillator is arrested, even without dissipation, when nonlinear effects come into play. This is due to the fact that the natural frequency of a nonlinear oscillator is amplitude-dependent, so a mismatch between the (invariable) driving frequency and the natural frequency appears Landau1 .
To overcome the nonlinear mismatch and maintain phase locking between the driving force and the oscillator, one can slowly vary the driving frequency with time so as to achieve a persistent growth of the oscillations. This simple and versatile method is called autoresonance: either external, or parametric. To emphasize the difference between the two, let us watch a child on a swing. When a parent pushes the swing (once in each cycle), he gradually increases the time interval between the pushes as the swing amplitude grows. Here he employs external autoresonance. On the contrary, when the child swings himself (he achieves it by moving the position of his center of mass up and down twice in each cycle), he gradually increases the period of these modulations as the swing amplitude grows. This is parametric autoresonance.
The simple model of a nonlinear oscillator, excited via external autoresonance, has found numerous applications in physics, see e.g. Ref. Friedland1 for a brief review. The external autoresonance scheme has been also generalized to systems with an infinite number of degrees of freedom, such as nonlinear waves Meerson5 ; Meerson3 ; Lazar\_3wave and vortices vortices . In contrast to the external autoresonance, parametric autoresonance has received much less attention Khain . In this work we generalize to nonlinear waves the theory, developed in Ref. Khain for a nonlinear oscillator. Specifically, we show that the parametric autoresonance mechanism can be used for driving nonlinear standing gravity waves with a steadily growing amplitude on a free surface of a fluid.
Michael Faraday Faraday was the first to observe that, when a tank containing a fluid is periodically vibrated in the vertical direction, a standing wave pattern forms at the free surface of the fluid when the vibration frequency is twice the frequency of the surface vibrations. This phenomenon is a classic example of parametric resonance, because the vertical acceleration of the tank - an intrinsic parameter of the system - depends on time via the periodic vibration. Lord Rayleigh Rayleigh carried out a further series of experiments, which supported Faradayโs observations, and also developed a linear theory for these waves in terms of a linear Mathieu equation. Benjamin and Ursell Ursell advanced the linear theory further. Subsequently, Miles Miles1 ; Miles2 ; Miles3 , Douady douady , Milner milner , and Decent and Craik decent formulated a weakly nonlinear theory based on amplitude expansion, while some of these and indeed numerous other works dealt with experimental studies of Faraday waves.
Being interested in parametric autoresonance, we add a new dimension to the problem of Faraday waves and investigate weakly nonlinear standing gravity waves formed when the vibration frequency is slowly decreased (chirped downwards) in time. We show that the negative frequency chirp causes a persistent growth of the wave amplitude. Like in other instances of autoresonance, the exact form of the frequency chirp is unimportant once the chirp sign is correct, and the chirp rate is not too high. The autoresonance excitation is expected to terminate at large amplitudes, when an underlying dynamical system, corresponding to the case of a constant frequency, ceases to exhibit a non-trivial stable fixed point.
Here is the layout of the rest of the paper. Section II presents a brief overview of theory of weakly nonlinear Faraday waves with a constant driving frequency. Sections III and IV deal with theory of chirped Faraday waves, in inviscid (Sec. III) and low-viscosity (Sec. IV) fluids. Section V presents a brief discussion of our results.
## II Weakly nonlinear Faraday waves
### II.1 Inviscid fluid
To set the stage for a theory of chirped Faraday waves, we need to briefly review the theory of weakly nonlinear Faraday waves with a constant driving frequency. Consider a quasi-two-dimensional rectangular tank with a fluid of length $`l`$, width $`w`$ and depth $`h`$, so that $`lw`$. We assume that the elevation of the fluid, caused by the wave, depends only on the longitudinal coordinate $`x`$ and time $`t`$, so that we have a quasi-two-dimensional flow in the $`xz`$ plane ($`z`$ is the vertical coordinate). The unperturbed level of the fluid is at $`z=0`$. The vertical displacement of the vibrating tank is described by the equation
$$\zeta (t)=a_0\mathrm{cos}(2\omega t).$$
(1)
We assume weak forcing, that is the vibration acceleration is much less than the gravity acceleration $`g`$, and introduce a small dimensionless parameter $`\epsilon `$:
$$\epsilon =\frac{\omega ^2a_0}{g}\mathrm{\hspace{0.17em}1}.$$
(2)
In the limit of inviscid fluid the flow remains potential once it is potential at $`t=0`$, and the external forces are potential Lamb , as is the case here. The assumption of a potential flow is also approximately valid in a low-viscosity fluid Ursell . We also assume that the wavelength of the standing wave is much larger than the capillary length of the fluid and neglect the capillary effects throughout the paper. The linear dispersion relation for the wave is $`\omega _n^2=gk_n\mathrm{tanh}(k_nh)`$, where $`\omega _n`$ is the natural frequency of the $`n`$-th mode, $`k_n=2\pi /\lambda _n=n\pi /l`$ is the wave number of the $`n`$-th mode, and $`n=1,2,\mathrm{}`$.
The governing equations for the velocity potential $`\phi (x,z,t)`$ and the wave profile $`\eta (x,t)`$ are Ursell ; Lamb ; Landau2 ; Currie :
$`^2\phi `$ $`=`$ $`0,`$ (3)
$`\left[\phi _t+{\displaystyle \frac{1}{2}}(\phi _x^2+\phi _z^2)+(g+\ddot{\zeta })\eta \right]|_{_{z=\eta }}`$ $`=`$ $`0,`$ (4)
$`\left(\eta _t+\phi _x\eta _x\phi _z\right)|_{_{z=\eta }}`$ $`=`$ $`0,`$ (5)
$`\phi _z|_{_{z=h}}`$ $`=`$ $`0,`$ (6)
where indices denote partial derivatives. The Laplaceโs equation (3) describes a potential flow of an incompressible fluid. Equations (4) and (5) are the Navier-Stokes equation and the kinematic boundary condition, respectively, evaluated at the free surface. Finally, Eq. (6) is the boundary condition for the vertical velocity component at the bottom of the tank.
Let the vibration frequency be close to twice the natural frequency of the primary mode $`n=1`$, i.e. $`\omega \omega _1`$, so that this mode is excited via parametric resonance. In a weakly nonlinear regime it suffices to account for the excitation of only one higher order mode: the secondary mode $`n=2`$, which is enslaved to the primary mode Miles2 . Therefore, one should look for $`\phi (x,z,t)`$ and $`\eta (x,t)`$ in Eqs. (3)-(6) in the following form Miles1 ; Miles2 :
$`\phi (x,z,t)`$ $`=`$ $`\phi _0(t)+\phi _1(t)\psi _1(x){\displaystyle \frac{\mathrm{cosh}\left[k_1(z+h)\right]}{\mathrm{cosh}(k_1h)}}+`$
$`+`$ $`\phi _2(t)\psi _2(x){\displaystyle \frac{\mathrm{cosh}\left[k_2(z+h)\right]}{\mathrm{cosh}(k_2h)}}+\mathrm{},`$
$`\eta (x,t)`$ $`=`$ $`\eta _1(t)\psi _1(x)+\eta _2(t)\psi _2(x)+\mathrm{},`$ (7)
where the eigenfunctions $`\psi _n(x)=\sqrt{2}\mathrm{cos}(k_nx)`$. The higher order terms will be neglected in the following. In the deep-water limit $`k_nh1`$, the linear dispersion relation for the wave becomes $`\omega _n^2gk_n`$. For this approximation to hold with an error less than 0.5%, it suffices to demand that $`h>l`$.
The perturbation theory we are using employs the smallness of $`\epsilon `$. As will be seen later, this smallness implies a smallness of the wave amplitude compared with the wavelength, so that the dimensionless parameter $`\kappa =k_1\eta `$ is small. Expanding $`\phi (x,z=\eta ,t)`$ in the vicinity of the unperturbed surface $`z=0`$ in a power series in $`\kappa `$, and substituting it and the second of Eqs. (II.1) into Eqs. (4) and (5), we obtain in the leading and sub-leading orders of $`\kappa `$:
$$\eta _2(t)\frac{k_1\eta _1^2(t)}{\sqrt{2}},$$
(8)
$$\phi _0(t)\dot{\eta }_1\eta _1,\phi _1(t)\frac{\dot{\eta }_1}{k_1},\phi _2(t)0.$$
(9)
As we see, the next-order corrections $`\phi _0`$ and $`\eta _2`$ are enslaved to the primary mode, and their magnitudes are $`๐ช(\kappa \eta _1)`$. In addition, we obtain a nonlinear differential equation of the second-order for the time-dependent amplitude of the primary mode $`\eta _1(t)`$:
$$\ddot{\eta }_1+\frac{1}{2}k_1^2(5\dot{\eta }_1^2\eta _13\omega _1^2\eta _1^3)+\omega _1^2\left[1+4\epsilon \mathrm{cos}(2\omega t)\right]\eta _1=0,$$
(10)
where we have used Eq. (1) and kept terms up to $`๐ช(\kappa ^3)`$. Equation (10) is a generalization of the linear Mathieu equation Bogoliubov .
Now we employ the method of averaging Bogoliubov . We make an Ansatz $`\eta _1(t)=A_1(t)\mathrm{cos}\left[\omega _1t+\varphi _1(t)\right]`$ and $`\dot{\eta }_1(t)=\omega _1A_1(t)\mathrm{sin}\left[\omega _1t+\varphi _1(t)\right]`$ in Eq. (10) and treat the amplitude $`A_1(t)`$ and phase $`\varphi _1(t)`$ as slow functions of time. (Being interested in the first-order equations with respect to $`\epsilon `$, one can omit higher temporal harmonics in $`\eta _1(t)`$ Bogoliubov ; perturbative .) Let $`\delta =\omega _1\omega `$ be the (small) detuning from the exact linear resonance. Then, for $`\epsilon \mathrm{\hspace{0.17em}1}`$ and $`|\delta |\omega _1`$, we can perform averaging over the fast time $`\omega _1^1`$ Bogoliubov . Introducing a new phase variable $`\varphi =\delta t+\varphi _1`$, we obtain:
$`\dot{A_1}`$ $`=`$ $`\epsilon \omega _1A_1\mathrm{sin}(2\varphi ),`$
$`\dot{\varphi }`$ $`=`$ $`\epsilon \omega _1\mathrm{cos}(2\varphi ){\displaystyle \frac{k_1^2A_1^2\omega _1}{4}}+\delta .`$ (11)
The second term in the right side of the equation for $`\dot{\varphi }`$ describes the nonlinear frequency shift of the standing wave. One can see that, as the wave amplitude grows, its frequency goes down only\_standing . This fact is important in the autoresonance excitation scheme introduced below. Rescaling time, amplitude and detuning,
$$\tau =\epsilon \omega _1t,B=\frac{k_1}{2\sqrt{\epsilon }}A_1,\mathrm{\Delta }=\frac{\delta }{\epsilon \omega _1},$$
(12)
we rewrite Eqs. (11) in a scaled form:
$`\dot{B}`$ $`=`$ $`B\mathrm{sin}(2\varphi ),`$
$`\dot{\varphi }`$ $`=`$ $`\mathrm{cos}(2\varphi )B^2+\mathrm{\Delta },`$ (13)
where the time derivatives are taken with respect to the slow time $`\tau `$. When $`\mathrm{\Delta }๐ช(1)`$, the typical value of $`B`$ (for example the stable fixed point, see below) is $`๐ช(1)`$. Going back to Eq. (12), we see that, in the dimensional units, the parameter $`k_1A_1\epsilon ^{1/2}1`$. As in the leading order $`\kappa k_1A_1`$, this validates our assumption that $`\kappa \mathrm{\hspace{0.17em}1}`$.
Equations (13) describe weakly-nonlinear constant-frequency Faraday waves in the leading order in $`\epsilon `$. In the context of Faraday waves, Eqns. (13) were first obtained by Miles Miles1 ; Miles2 , though he derived them in a different way, working with the Lagrangian of the fluid. In the sub-leading order in $`\epsilon `$, additional nonlinear terms appear Miles3 ; douady ; milner ; decent , which will not be considered here.
Equations (13) can be rewritten in a Hamiltonian form if we introduce the action and angle variables $`I=B^2/2\text{and}\varphi `$:
$`\dot{I}`$ $`=`$ $`{\displaystyle \frac{H}{\varphi }}=2I\mathrm{sin}(2\varphi ),`$
$`\dot{\varphi }`$ $`=`$ $`{\displaystyle \frac{H}{I}}=\mathrm{cos}(2\varphi )2I+\mathrm{\Delta },`$ (14)
where the Hamiltonโs function is
$$H(I,\varphi )=I\mathrm{cos}(2\varphi )I^2+\mathrm{\Delta }I,$$
(15)
The fixed points of this dynamical system are determined by the value of the scaled detuning $`\mathrm{\Delta }`$ Bogoliubov ; Struble :
* $`\mathrm{\Delta }<1`$. No fixed points.
* $`1<\mathrm{\Delta }<1`$. Three fixed points: an elliptic point $`[I_{},\varphi _{}]=[(1+\mathrm{\Delta })/2,0]`$ and two saddle points $`[I_{},\varphi _{}]=[0,\pm \mathrm{arccos}(\mathrm{\Delta })/2]`$.
* $`\mathrm{\Delta }>1`$. Two fixed points: the same elliptic point $`[I_{},\varphi _{}]=[(1+\mathrm{\Delta })/2,0]`$ as in case (b), and a saddle point $`[I_{},\varphi _{}]=[(\mathrm{\Delta }1)/2,\pi /2]`$.
Figure 1 shows the phase plane ($`\varphi `$,$`I`$) in the cases of $`0<\mathrm{\Delta }<1`$ and $`\mathrm{\Delta }>1`$. The phase portrait is periodic in $`\varphi `$ with period $`\pi `$. In the case of $`0<\mathrm{\Delta }<1`$, the separatrix is formed by the curve $`I=\mathrm{\Delta }+\mathrm{cos}(2\varphi )`$ and the straight line $`I=0`$. In the case of $`\mathrm{\Delta }>1`$, the separatrix is formed by the curves $`I=[\mathrm{cos}(2\varphi )+\mathrm{\Delta }\sqrt{\alpha }]/2`$ and $`I=[\mathrm{cos}(2\varphi )+\mathrm{\Delta }+\sqrt{\alpha }]/2`$, where $`\alpha =\mathrm{cos}^2(2\varphi )+2\mathrm{\Delta }\mathrm{cos}(2\varphi )+2\mathrm{\Delta }1`$. Notice that the maximum possible amplitude of phase-locked oscillations is achieved at a nonzero detuning from the exact linear resonance, like in many other instances of nonlinear resonance.
As the Hamiltonโs function (15) is a constant of motion, the system is integrable. In particular, one can find the โnonlinear periodโ: the period of motion along a closed trajectory in the phase plane. Denoting the constant Hamiltonโs function as $`H_0`$, we obtain
$$T_{nl}=2_\varphi _{}^{\varphi _+}\frac{d\varphi }{\left\{\left[\mathrm{\Delta }+\mathrm{cos}(2\varphi )\right]^24H_0\right\}^{1/2}},$$
(16)
where $`\varphi _\pm =\pm \mathrm{arccos}(2\sqrt{H_0}\mathrm{\Delta })/2`$. For a zero detuning, and initial conditions very close to the fixed point $`I_{}=1/2`$ and $`\varphi _{}=0`$ (so that $`H_0\mathrm{\hspace{0.17em}1}/4`$), we obtain $`T\pi `$. This corresponds to small harmonic oscillations around the elliptic fixed point. In the physical units the period of small oscillations is $`T_{nl}^{ph}\pi /(\epsilon \omega _1)`$, that is much longer than the wave period.
### II.2 Low-viscosity fluid
Taking into account a weak damping of the wave amounts to adding a linear damping term $`2\gamma \dot{\eta }_1`$ to the left side of Eq. (10), where $`\gamma `$ is defined in terms of the rate of loss of mechanical energy due to dissipation, Landau2 . The incorporation of only a linear damping term requires that $`\gamma /\omega _1\mathrm{\hspace{0.17em}1}`$, so that the damping is treated perturbatively. The specific damping mechanisms which contribute to the value of damping rate $`\gamma `$ are the bulk viscosity Landau2 , dissipation in the vicinity of the fixed walls Miles4 , dissipation at the free surface (especially if contaminated) Miles4 and contact line damping christiansen (see, e.g. Ref. christiansen for a review). In practice, one can interpret the damping rate as a phenomenological term, and determine it from a comparison with experiment.
Including the linear damping term in the first of Eqs. (13), we obtain:
$`\dot{B}`$ $`=`$ $`B\mathrm{sin}(2\varphi )\mathrm{\Gamma }B,`$
$`\dot{\varphi }`$ $`=`$ $`\mathrm{cos}(2\varphi )B^2+\mathrm{\Delta },`$ (17)
where $`\mathrm{\Gamma }=\gamma /(\epsilon \omega _1)>0`$ is a dimensionless damping rate. It follows from the first of Eqs. (17), that non-trivial fixed points $`B_{}0`$ can exist only when $`\mathrm{\Gamma }<1`$, that is for a small enough viscosity. In this case, any trajectory on the phase plane of the system (except for trajectories with a zero measure) converges to a stable focus, see Fig. 2. This is in contrast to the inviscid case, where starting from initial conditions outside the separatrix leaves the trajectory phase-unlocked. The fixed points $`[B_{},\varphi _{}]`$ of Eqs. (17) are determined by the values of $`\mathrm{\Delta }`$ and $`\mathrm{\Gamma }`$. Since Eqs. (17) \[and Eqs. (13)\] are invariant under the transformation $`BB`$, one needs to consider only fixed points with positive amplitudes. Let $`\sigma =\sqrt{1\mathrm{\Gamma }^2}`$ and $`\xi =\sqrt{1\mathrm{\Delta }^2}`$. Let us also denote the critical damping rate
$$\mathrm{\Gamma }_{cr}=\left[\frac{4}{5}\frac{8\mathrm{\Delta }}{25}\left(\mathrm{\Delta }\sqrt{\mathrm{\Delta }^2+\frac{5}{4}}\right)\right]^{1/2},$$
(18)
which will appear shortly. There are four different cases:
* $`\mathrm{\Delta }<1`$. No fixed points.
* $`1\mathrm{\Delta }<0`$. For $`\mathrm{\Gamma }>\xi `$, $`[0,\mathrm{arccos}(\mathrm{\Delta })/2]`$ is a stable node, and $`[0,\mathrm{arccos}(\mathrm{\Delta })/2]`$ is a saddle point. For $`0\mathrm{\Gamma }\xi `$, $`[0,\pm \mathrm{arccos}(\mathrm{\Delta })/2]`$ are two saddle points, and $`[(\sigma +\mathrm{\Delta })^{1/2},\mathrm{arcsin}(\mathrm{\Gamma })/2]`$ is a stable fixed point. For $`0\mathrm{\Gamma }<\mathrm{\Gamma }_{cr}`$ it is a stable focus, while for $`\mathrm{\Gamma }_{cr}<\mathrm{\Gamma }\xi `$ it is a stable node.
* $`0\mathrm{\Delta }<1`$. For $`\mathrm{\Gamma }\mathrm{\hspace{0.17em}1}`$, $`[0,\mathrm{arccos}(\mathrm{\Delta })/2]`$ is a stable node, and $`[0,\mathrm{arccos}(\mathrm{\Delta })/2]`$ is a saddle point. For $`\xi <\mathrm{\Gamma }<\mathrm{\hspace{0.17em}1}`$, $`[0,\mathrm{arccos}(\mathrm{\Delta })/2]`$ is a stable node, $`[0,\mathrm{arccos}(\mathrm{\Delta })/2]`$ and $`[(\mathrm{\Delta }\sigma )^{1/2},\pi /2\mathrm{arcsin}(\mathrm{\Gamma })/2]`$ are two saddle points, and $`[(\sigma +\mathrm{\Delta })^{1/2},\mathrm{arcsin}(\mathrm{\Gamma })/2]`$ is a stable fixed point. For $`\xi <\mathrm{\Gamma }<\mathrm{\Gamma }_{cr}`$ it is a stable focus, while for $`\mathrm{\Gamma }_{cr}<\mathrm{\Gamma }<1`$ it is a stable node. For $`0\mathrm{\Gamma }\xi `$, $`[0,\pm \mathrm{arccos}(\mathrm{\Delta })/2]`$ are two saddle points, and $`[(\sigma +\mathrm{\Delta })^{1/2},\mathrm{arcsin}(\mathrm{\Gamma })/2]`$ is a stable focus.
* $`\mathrm{\Delta }1`$. For $`0\mathrm{\Gamma }\mathrm{\hspace{0.17em}1}`$, $`[(\sigma +\mathrm{\Delta })^{1/2},\mathrm{arcsin}(\mathrm{\Gamma })/2]`$ is a stable fixed point. For $`0\mathrm{\Gamma }<\mathrm{\Gamma }_{cr}`$ it is a stable focus, while for $`\mathrm{\Gamma }_{cr}<\mathrm{\Gamma }\mathrm{\hspace{0.17em}1}`$ it is a stable node, and $`[(\mathrm{\Delta }\sigma )^{1/2},\pi /2\mathrm{arcsin}(\mathrm{\Gamma })/2]`$ is a saddle point.
Figure 3 shows two characteristic values of the scaled damping $`\mathrm{\Gamma }`$ as functions of $`\mathrm{\Delta }`$. The first one is $`\mathrm{\Gamma }_{cr}`$ from Eq. (18). The second one is the maximum value of $`\mathrm{\Gamma }`$ for which a nontrivial stable fixed point still exists. For $`1<\mathrm{\Delta }<0`$ this maximum value is equal to $`\sqrt{1\mathrm{\Delta }^2}`$, while for $`\mathrm{\Delta }\mathrm{\hspace{0.17em}0}`$ it is equal to $`1`$. To conclude the brief review of the constant-frequency theory, we notice that the dependence of $`B_{}`$ on $`\mathrm{\Gamma }`$ exhibits a pitchfork bifurcation. Figure 4 shows the bifurcation diagram in case (c) for $`\mathrm{\Delta }=0.5`$.
## III Chirped Faraday waves in an inviscid fluid
### III.1 Governing equations, phase portrait, criteria and numerical examples
Now let the vibration frequency be time-dependent (chirped): $`\omega =\omega (t)`$. In general, the dimensionless parameter $`\epsilon =\omega ^2a_0/g`$ will also become time-dependent. For simplicity, we shall assume that $`a_0`$ also varies in time so that $`\epsilon =const`$ constamp . Our objective is to keep a Faraday wave close to resonance in spite of its nonlinear frequency shift, so as to achieve a persistent growth of the wave amplitude. Like in other autoresonance schemes, the exact form of the function $`\omega (t)`$ is unimportant if this function satisfies three criteria:
1. The chirp sign coincides with the sign of the nonlinear frequency shift of the wave. For the standing Faraday waves $`\omega (t)`$ should decrease for the wave amplitude to increase.
2. The frequency chirp rate must be sufficiently small, so that the phase portrait of the system evolves adiabatically: $`|\dot{\omega (t)}|T_{nl}\omega (t)`$, where $`T_{nl}`$ is the characteristic nonlinear period, see Eq. (16).
3. The dynamic frequency mismatch, which we define as the absolute value of the increment of the vibration frequency during one nonlinear period, $`|\omega (t+T_{nl})\omega (t)|`$, should be small compared with the inverse nonlinear period, $`T_{nl}^1`$. In physical units, this yields $`\mu /(\epsilon \omega _1)^2\mathrm{\hspace{0.17em}1}`$.
Criteria 1 and 2 have appeared in previous works on autoresonance Friedland1 , while criterion 3 is new. We shall see shortly that, in the problem of parametric autoresonance, criterion 3 is more restrictive than criterion 2.
The derivation of the equation of motion for the primary mode amplitude, for a slowly time-dependent driving frequency, goes along the same lines as in the case of a constant driving frequency. The resulting equation is (compare with Eq. (10)\]:
$$\ddot{\eta }_1+\frac{1}{2}k_1^2(5\dot{\eta }_1^2\eta _13\omega _1^2\eta _1^3)+\omega _1^2\left[1+4\epsilon \mathrm{cos}(2\mathrm{\Phi }(t))\right]\eta _1=0,$$
(19)
where $`\mathrm{\Phi }(t)=_0^t\omega (t^{})๐t^{}`$. If $`\epsilon `$ is small, and the chirp rate is slow on the time scale of the wave period, one can again use, close to the parametric resonance, the method of averaging Khain . For concreteness, we assume in this work a constant chirp rate $`\mu `$:
$$\omega (t)=\omega _1\mu t,$$
(20)
so that $`\mathrm{\Phi }(t)=\omega _1t\mu t^2/2`$. Introducing a scaled chirp rate $`m=\mu /(\omega _1\epsilon )^2`$, and the same scaled time $`\tau `$ and amplitude $`B`$ as before \[see Eq. (12)\], we arrive at the following scaled equations:
$`\dot{B}`$ $`=`$ $`B\mathrm{sin}(2\varphi ),`$
$`\dot{\varphi }`$ $`=`$ $`\mathrm{cos}(2\varphi )B^2+m\tau ,`$ (21)
where now $`\varphi (t)=\mu t^2/2+\varphi _1(t)`$, and the differentiation is done with respect to $`\tau `$. In the action-angle variables we obtain:
$`\dot{I}`$ $`=`$ $`2I\mathrm{sin}(2\varphi ),`$
$`\dot{\varphi }`$ $`=`$ $`\mathrm{cos}(2\varphi )2I+m\tau .`$ (22)
Once $`B(t)`$ and $`\varphi (t)`$ are found, one can immediately reconstruct the standing way profile $`\eta (x,t)`$ by using the Ansatz $`\eta _1(t)=A_1(t)\mathrm{cos}\left[\omega _1t+\varphi _1(t)\right]`$ and the โenslaving relationโ (8) in the second of equations (II.1). Therefore, in the rest of the paper we shall focus on Eqs. (22) which coincide, up to notation, with those obtained by Khain and Meerson Khain , who investigated parametric autoresonance in a nonlinear oscillator. Equations similar to (22) also appear in the problem of the second-harmonic autoresonance in an externally driven oscillator Lazar7 .
The Hamiltonโs function of the system (22) is time dependent:
$$H(I,\varphi )=I\mathrm{cos}(2\varphi )I^2+m\tau I,$$
(23)
so $`H`$ is not a constant of motion anymore. In the following we shall use $`t`$ for the slow time $`\tau `$.
Figure 5 shows an example of parametric autoresonance: a persistent phase locking and a systematic growth of $`I`$ with time, with some oscillations on top of the systematic growth. Here the scaled chirp rate $`m`$ is less than some critical value $`m_{cr}`$ for these initial conditions. Figure 6 illustrates breakdown of autoresonance, observed when $`m>m_{cr}`$. As in other instances of autoresonance, a theory of parametric autoresonance appeals to the constant-frequency theory. Comparing Eqs. (22) and (14), one can see that the term $`mt`$ plays the role of an effective (time-dependent) detuning. Therefore, when the chirp rate is small, $`m\mathrm{\hspace{0.17em}1}`$, the phase portrait of the system almost coincides with that of the autonomous equations (14) (see Fig. 1), except that now it changes with time according to the current value of the effective detuning. The change of the phase portrait is adiabatically slow, except at times $`t1/m`$ and $`1/m`$ (corresponding to $`\mathrm{\Delta }=1`$ and $`1`$, respectively), when bifurcations occur. One consequence of the adiabatic evolution is that Eqs. (22) have โquasi-fixedโ points. The most important stable quasi-fixed point $`[I_{}(t),\varphi _{}(t)]`$ can be found by assuming that $`\varphi _{}(t)1`$, and that it varies with time slowly. Then the second of Eqs. (22) yields, in the leading order,
$$I_{}\frac{1}{2}\left(1+mt\right).$$
(24)
Substituting this into the first of Eqs. (22), we obtain
$$\varphi _{}\frac{m}{4(1+mt)}.$$
(25)
The stable quasi-fixed point, or trends (24) and (25), previously found by Khain and Meerson Khain (see also Ref. Lazar7 ), are the essence of parametric autoresonance. Shown in Fig. 5 are $`I(t)`$ and $`\varphi (t)`$ found numerically, and the trends (24) and (25). The trend (24) corresponds to a steady growth of the wave amplitude: $`B_{}(t)=[2I_{}(t)]^{1/2}(mt+1)^{1/2}`$. The important phase trend (25) was overlooked in Ref. Lazar7 . Notice that, at scaled time $`t1`$, the phase trend $`\varphi _{}1/(4t)`$ becomes independent of the chirp rate $`m`$. Importantly, for the expressions (24) and (25) to be valid, one can either demand $`m1`$, or go to long times: $`t1`$. Therefore, the stable quasi-fixed point keeps its meaning, at long times, even at finite (non-small) $`m`$.
We found that, surprisingly, unstable quasi-fixed points of the chirped system also play an important role in the dynamics. The unstable points are analogs of the constant-frequency saddle points discussed in the previous section \[see the text following Eq. (15)\]. To find the locations of the unstable quasi-fixed points in the leading order, one can simply replace the detuning $`\mathrm{\Delta }`$ by $`mt`$. Therefore, on the time interval $`0<t<1/m`$, there are two saddle quasi-fixed points $`[I_{},\varphi _{}][0,\pm (1/2)\mathrm{arccos}(mt)]`$. These points disappear at $`t1/m`$, and a new saddle points appears: $`[I_{},\varphi _{}][(mt1)/2,\pi /2)]`$. These expressions (including the boundaries of the corresponding time intervals) are valid in the leading order in $`m1`$. Higher-order corrections can be also calculated.
Now we are in a position to discuss criteria 2 and 3 for parametric autoresonance in this system. For a constant chirp rate $`\mu `$ \[see Eq. (20)\], criterion 2 can be written, in the physical units, as $`\mu (\epsilon \omega )^1<\omega `$, or $`\mu <\epsilon \omega ^2`$. Now, the dynamic frequency mismatch, acquired by the chirped system during time $`T_{nl}`$, can be estimated as $`\mu T_{nl}\mu /(\epsilon \omega )`$. Criterion 3 demands that this quantity be small compared to $`T_{nl}^1\epsilon \omega `$, which yields $`\mu <\epsilon ^2\omega ^2`$. As $`\epsilon `$ is small, criterion 3 is more restrictive than criterion 2. In the scaled variables, criterion 3 has the form of $`m<1`$, as can be expected from the form of scaled Eqs. (22). The inequalities here are written up to numerical factors which depend on the initial conditions, see below.
One more convenient description of the chirped system can be achieved if we rewrite Eqs. (22) as a second order equation for the phase:
$$\ddot{\varphi }+2mt\mathrm{sin}(2\varphi )+\mathrm{sin}(4\varphi )m=0,$$
(26)
or
$$\ddot{\varphi }+\frac{V(\varphi ,t)}{\varphi }=0,$$
(27)
where we have introduced a time-dependent potential
$$V(\varphi ,t)=\frac{1}{4}\mathrm{cos}(4\varphi )mt\mathrm{cos}(2\varphi )m\varphi .$$
(28)
This suggests new canonical variables $`\varphi `$ and $`u=\mathrm{cos}(2\varphi )2I+mt`$, so that in the new time-dependent Hamiltonian, $`H(\varphi ,u,t)=u^2/2+V(\varphi ,t)`$, there is a clear separation between the potential energy and the kinetic energy. The new Hamiltonian describes a โparticleโ of a unit mass and velocity $`u=\dot{\varphi }`$, moving in a time-dependent potential $`V`$. This picture is useful for a qualitative analysis of the dynamics of the โparticleโ when $`m`$ is small, so the potential slowly varies in time, see Fig. 7. In the variables $`u,\varphi `$ the stable quasi-fixed point becomes (approximately) $`[0,(m/4)(1+mt)^1]`$, while the saddle points are $`[0,\pm (1/2)\mathrm{arccos}(mt)]`$ at $`0<t<1/m`$, and $`[0,\pm \pi /2]`$ at $`t>1/m`$. We shall see shortly that each of the unstable points $`[0,(1/2)\mathrm{arccos}(mt)]`$ and $`[0,\pi /2]`$ plays an important role in this system.
Let us consider two typical cases of parametric autoresonant excitation of a Faraday wave. In the first case one first excites the wave at a constant frequency, so that the initial values of the action and phase are in the vicinity of the stable fixed point. Then, upon slowly reducing the driving frequency, one keeps the phase locked, as our โparticleโ oscillates in a potential well which slowly deepens with time, see Fig. 7. In the second case one starts the autoresonant driving from an almost zero wave amplitude. Here the saddle point $`[I_{},\varphi _{}][0,(1/2)\mathrm{arccos}(mt)]`$ plays an important role. Indeed, the stable manifold of this quasi-fixed point is along the $`\varphi `$ axis. Therefore, no matter what the initial phase is, the phase approaches, on a time scale $`๐ช(1)`$, the saddle point. The unstable manifold of this saddle point is along the $`I`$ axis, so $`I(t)`$ will grow with time. Still, if $`I(t=0)`$ is small enough, $`I(t)`$ remains small during this time scale $`๐ช(1)`$. Therefore, phase locking is always achieved at this stage, so the time interval $`0<t<1/m`$ can be called the โtrapping stageโ. Later on $`I(t)`$ grows significantly but, as we found numerically, the โparticlesโ remain inside the (slowly expanding) separatrix $`I=\mathrm{cos}(2\varphi )+mt`$. As a result, the phase starts to perform large-amplitude oscillation around the stable quasi-fixed point, and phase locking persists.
One more alternative description of the system of equations (21) is in terms of the complex amplitude $`\psi (t)=B(t)\mathrm{exp}[i\varphi (t)]`$:
$$i\psi _t+\psi ^{}(|\psi |^2mt)\psi =0,$$
(29)
where the subscript $`t`$ denotes differentiation with respect to the slow time. The long-time behavior of $`I_{}`$ can be obtained by looking at the asymptotic solutions of Eq. (29) at $`t\mathrm{}`$. For a solution such that $`|\psi |`$ grows with time like a power law, the leading terms are those in the parentheses. This immediately yields
$$|\psi (t)|(mt)^{1/2},$$
(30)
which corresponds to the leading term (when $`mt\mathrm{\hspace{0.17em}1}`$) of Eq. (24), and describes a phase-locked wave \[here $`\varphi `$ stays close to zero, see Eq. (25)\]. On the contrary, if $`|\psi (t)|`$ remains bounded and small, the first term of Eq. (29) is balanced by the last one, and we obtain
$$\psi (t)=\psi _0\mathrm{exp}\left(imt^2/2\right),$$
(31)
where $`\psi _0|\psi _0|\mathrm{exp}(i\varphi _0)=const.`$ This solution corresponds to an unlocked phase $`\varphi (t)=\varphi _0+mt^2/2`$ and a constant amplitude $`|\psi _0|`$. Of course, the phase of the wave $`\varphi _1(t)=\varphi (t)mt^2/2`$, which is defined by the Ansatz $`\eta _1(t)=A_1(t)\mathrm{cos}\left[\omega _1t+\varphi _1(t)\right]`$ (where $`t`$ is the physical time), stays constant in this regime, and is equal to $`\varphi _0`$.
We determined numerically, for several typical classes of initial conditions, the critical value of $`m`$, $`m=m_{cr}`$, which separates the phase locking regime from the phase unlocking regime alternative . At a fixed initial phase $`\varphi (t=0)=0`$, $`m_{cr}`$ grows with the initial amplitude $`B(t=0)`$, at least until the scaled amplitude becomes of order unity, see Fig. 8. This result agrees with those obtained by Fajans et. al. (see Fig. 3 in Ref. Lazar7 ) who presented them in terms of $`\epsilon _{cr}`$ versus $`\mu `$. When starting from the stable quasi-fixed point: $`I(t=0)=1/2`$ and $`\varphi (t=0)=0`$, we found that $`m_{cr}\mathrm{\hspace{0.17em}4.963}`$.
We also found $`m_{cr}`$ as a function of the initial phase, for a given, and very small, initial amplitude, see Fig. 9. This dependence is relatively weak. The largest $`m_{cr}`$ is obtained for $`\varphi =\pi /4`$, the smallest one for $`\varphi =\pi /4`$.
What is the signature of the special case $`m=m_{cr}`$? At $`m<m_{cr}`$ the phase oscillates in time. As $`m`$ approaches $`m_{cr}`$ from below, the onset of the phase oscillations is delayed more and more, see Fig. 10. Now, at $`m>m_{cr}`$ the phase initially changes slowly and then rapidly escapes to infinity. As $`m`$ approaches $`m_{cr}`$ from above, the point of rapid departure of the phase is delayed more and more, as shown in Fig. 10. This suggests that, in the special case $`m=m_{cr}`$ the phase neither oscillates around a trend, nor departs from it. Instead, the phase monotonically approaches $`\pi /2`$. Looking at the effective potential, shown in Fig. 7, one realizes that, in the special case $`m=m_{cr}`$, our โparticleโ neither oscillates in the potential well (which would correspond to phase locking), nor escapes from the well (which would correspond to phase unlocking). Instead, the โparticleโ lands, at $`t=\mathrm{}`$, on the peak of the potential at $`\varphi =\pi /2`$. In the following we shall find the asymptotic form of this special trajectory.
### III.2 Perturbative solutions
In this subsection we present three perturbative analytic solutions which illustrate the basic features of parametric autoresonance: persistent resonant growth, capture into resonance and the limiting trajectory which separates between phase locking and unlocking. In each of the three cases, a local analysis around one of the quasi-fixed points of the system is required.
#### III.2.1 In the vicinity of the stable quasi-fixed point
Let us linearize Eq. (26) in the vicinity of $`\varphi =0`$. As discussed above, this requires either a small chirp rate, $`m1`$, or a long time, $`t1`$. We obtain:
$$\ddot{\varphi }+4(1+mt)\varphi =m.$$
(32)
We look for the solution in the form
$$\varphi (t)=\frac{m}{4(1+mt)}+\delta \varphi (t),$$
(33)
where the first term is the phase trend (25). Neglecting higher-order terms, we arrive at the Airy equation
$$\ddot{\delta \varphi }+4(1+mt)\delta \varphi =0,$$
(34)
whose general solution is Abramowitz :
$`\delta \varphi (\tau )`$ $`=`$ $`C_1Ai\left[\left({\displaystyle \frac{2}{m}}\right)^{2/3}(1+mt)\right]+`$ (35)
$`C_2Bi\left[\left({\displaystyle \frac{2}{m}}\right)^{2/3}(1+mt)\right].`$
Here $`Ai(\tau )`$ and $`Bi(\tau )`$ are the Airy functions of the first and second kind, and $`C_1`$ and $`C_2`$ are constants depending on the initial conditions. Using the large-argument expansion of $`Ai(z)`$ and $`Bi(z)`$ Abramowitz , we obtain:
$$\varphi (t)\frac{m}{4(1+mt)}+\frac{A}{(1+mt)^{1/4}}\mathrm{sin}\left[\frac{4(1+mt)^{3/2}}{3m}+\xi \right],$$
(36)
where $`A`$ and $`\xi `$ are new constants. The action $`I(t)`$ can be found using the second of Eqs. (22):
$$I(t)\frac{mt+1}{2}A(1+mt)^{1/4}\mathrm{cos}\left[\frac{4(1+mt)^{3/2}}{3m}+\xi \right].$$
(37)
The first terms in Eqs. (36) and (37) are the systematic trends (25) and (24). The solutions (36) and (37) coincide, up to notation, with the WKB-solutions obtained in Ref. Khain . Fig. 11 shows excellent agreement between the analytical solutions (36) and (37) and numerical solutions.
#### III.2.2 Capture into resonance
Now we consider the case of driving a Faraday wave starting from a very small amplitude and a large initial detuning, that is far from resonance. Here one is interested in the capture into resonance. In the case of external autoresonance this phenomenon has been extensively studied by Friedland Friedland1 . In the case of parametric autoresonance this phenomenon has not been addressed, although equations similar to (22) were analyzed in Ref. Lazar7 in the context of the second-harmonic autoresonance in an externally driven oscillator.
We start, in Eqs. (21), at a large negative time $`t_0<0`$, and assume that the initial detuning $`mt_0`$ is very large in the absolute value, while the initial scaled amplitude $`B(t=t_0)`$ is much less than unity. As long as $`B(t)1`$, one can neglect the $`B^2`$-term in the second of Eqs. (21):
$$\dot{\varphi }=\mathrm{cos}(2\varphi )+mt.$$
(38)
This equation describes two distinct stages of the dynamics. In the first stage $`t`$ is large and negative, and $`|mt|1`$. Therefore, the term $`mt`$ is dominant, so Lazar7
$$\varphi \varphi (t_0)+\frac{m}{2}(tt_0)^2.$$
(39)
During this pre-locking stage, the phase varies rapidly. The second stage occurs roughly at $`1/m<t<1/m`$. Here the two terms on the right hand side of Eq. (38) are comparable. As $`m1`$, the duration of this stage is long compared with unity, and the phase approaches, on a time scale $`๐ช(1)`$, the unstable quasi-fixed point
$$\varphi _{}(t)\frac{\mathrm{arccos}(mt)}{2}\frac{m}{4[1(mt)^2]^{1/2}},$$
(40)
where we have included the next-order correction in $`m`$. As noted previously, this quasi-fixed saddle point ceases to exist at $`t>1/m`$; the correction term in Eq. (40) is invalid too close to $`t=1/m`$. We call this regime โlinear phase lockingโ stage.
Let us obtain the dependence of $`B`$ on time in the pre-locking and linear phase locking stages. In the pre-locking stage $`\varphi (t)`$ is given by Eq. (39). Then the first of Eqs. (21) yields
$$B(t)B(t_0)\mathrm{exp}\left\{_{t_0}^t\mathrm{sin}\left[m(st_0)^2+2\varphi (t_0)\right]๐s\right\}$$
(41)
(the integral can be expressed through the Fresnel integral). At this stage $`B(t)`$ oscillates rapidly, because of the rapid and monotonic change of the phase. At the linear phase locking stage, $`\varphi (t)\varphi _{}(t)`$ as given by Eq. (40). Then, neglecting higher-order terms in $`m`$, we arrive at the equation
$$\dot{B}B\sqrt{1(mt)^2}$$
(42)
which yields
$$BB_2\mathrm{exp}\left\{\frac{1}{2}\left[t\sqrt{1(mt)^2}+\frac{\mathrm{arcsin}(mt)}{m}\right]\right\},$$
(43)
where $`B_2`$ is a constant determined by the initial conditions.
Equation (43) breaks down when the earliest of the two events occurs: $`t`$ becomes larger than $`1/m`$, or $`B(t)`$ becomes comparable to unity, so that one cannot neglect the $`B^2`$-term in the second of Eqs. (21). The phase starts to oscillate, while the amplitude both oscillates and grows, see Fig. 12, and the system enters the autoresonance regime.
Similar results are observed when starting from a small amplitude and a zero initial detuning (that is, in exact linear resonance). Therefore, when the initial amplitude is very small, and $`m\mathrm{\hspace{0.17em}1}`$, phase locking is very robust.
#### III.2.3 Critical chirp rate and the limiting trajectory
As we have seen numerically, when $`m`$ approaches the critical value $`m_{cr}`$, the phase $`\varphi `$ approaches $`\pi /2`$ at $`t1`$. Let us assume that, at a time $`t_0`$, $`\varphi `$ is already in the vicinity of $`\pi /2`$: $`\varphi =\pi /2\delta \varphi `$, where $`\delta \varphi 1`$. Linearizing Eq. (26), we obtain
$$\ddot{\delta \varphi }4\delta \varphi (mt1)=m.$$
(44)
To find the trend, we neglect the term $`\ddot{\delta \varphi }`$. Therefore, at $`t1`$, we obtain $`\delta \varphi 1/(4t)`$, so
$$\varphi _{}(t)\frac{\pi }{2}\frac{1}{4t}.$$
(45)
Using the second of Eqs. (22), we obtain the respective trend of $`I(t)`$:
$$I_{}(t)\frac{mt}{2}\frac{1}{2}.$$
(46)
At this stage we notice that Eqs. (45) and (46) describe, at $`t1`$, one of the unstable (saddle) quasi-fixed points of the system: the one that appears close to $`t=1/m`$. This again shows that adiabaticity holds not only at $`m1`$, but also at $`t1`$.
Now consider small deviations from the unstable trends (45) and (46). Putting $`\varphi =\pi /21/(4t)\delta \varphi (t)`$, we again arrive at the Airy equation for $`\delta \varphi (t)`$. Its general solution is $`\delta \varphi (t)=C_1Ai[(4m)^{1/3}t]+C_2Bi[(4m)^{1/3}t]`$, where $`C_1`$ and $`C_2`$ are determined by the initial conditions. Assuming $`m^{1/3}t1`$, we employ the asymptotics of the Airy functions Abramowitz :
$$Ai(z)\frac{\mathrm{exp}\left[(2/3)z^{3/2}\right]}{2\sqrt{\pi }z^{1/4}},Bi(z)\frac{\mathrm{exp}\left[(2/3)z^{3/2}\right]}{\sqrt{\pi }z^{1/4}}$$
(47)
at $`z\mathrm{\hspace{0.17em}1}`$. The exponentially decaying term is negligible. The exponentially growing term describes instability of the special trajectory with respect to small perturbations. This instability occurs both at $`m<m_{cr}`$ (when the phase leaves the vicinities of the saddle point and goes to the vicinity of the stable quasi-fixed point), and at $`m>m_{cr}`$ (when the phase locking terminates). For the special trajectory, obtained for $`m=m_{cr}`$, the coefficient $`C_2`$ must vanish, which brings us back to Eqs. (45) and (46), where we must put $`m=m_{cr}`$. Now it is clear that the asymptotic behavior of the $`\varphi _{}(t)`$ at $`m=m_{cr}`$ is independent of $`m_{cr}`$ and, therefore, on the initial conditions. On the contrary, the asymptotic behavior of $`I_{}(t)`$ at $`m=m_{cr}`$ does depend on $`m_{cr}`$ and, therefore, on the initial conditions. Unfortunately, the local analysis does not enable one to find the value of $`m_{cr}`$ analytically.
Figures 13 and 14 show the behavior of the system when $`m`$ is very close to (just below) $`m_{cr}`$. Figure 13 shows that, after a time $`๐ช(1)`$, the phase follows the trend (45), independently of the initial conditions (and of the value of $`m_{cr}`$). Figure 14 shows that, after a time of $`๐ช(1)`$, good agreement between $`I(t)`$, found numerically, and the trend (46) holds until the time when the โparticleโ leaves the vicinity of the unstable point and transfers to the vicinity of the stable point.
## IV Chirped Faraday waves in a low-viscosity fluid
We now account for a small viscosity and briefly describe the dynamics of the phase and amplitude of the primary mode in the case of a slowly time-dependent vibration frequency. The weakly nonlinear governing equations are obtained by analogy with Eqs. (17):
$`\dot{B}`$ $`=`$ $`B\mathrm{sin}(2\varphi )\mathrm{\Gamma }B,`$
$`\dot{\varphi }`$ $`=`$ $`\mathrm{cos}(2\varphi )B^2+mt,`$ (48)
where $`t`$ is the slow time as before. For $`\mathrm{\Gamma }<1`$, there exists a stable quasi-fixed point which describes autoresonant excitation of the wave:
$`B_{}`$ $``$ $`\left[(1\mathrm{\Gamma }^2)^{1/2}+mt\right]^{1/2},`$
$`\varphi _{}`$ $``$ $`{\displaystyle \frac{1}{2}}\mathrm{arcsin}(\mathrm{\Gamma })+{\displaystyle \frac{m}{4[1\mathrm{\Gamma }^2+mt(1\mathrm{\Gamma }^2)^{1/2}]}}.`$ (49)
When $`\mathrm{\Gamma }1`$, Eqs. (49) become
$$B_{}\sqrt{1+mt},\varphi _{}\frac{\mathrm{\Gamma }}{2}+\frac{m}{4(1+mt)}.$$
(50)
Figure 15 shows a projection on the $`(\varphi ,B)`$ plane of a three-dimensional trajectory in the space of $`\varphi `$, $`B`$ and $`t`$. One can see phase locking and a steady growth of the wave amplitude with time.
Figures 16 and 17 show two different regimes of autoresonant growth in the dissipative system. Shown in the two figures are the primary mode amplitude versus time for the scaled damping rates $`\mathrm{\Gamma }<\mathrm{\Gamma }_{cr}`$ and $`\mathrm{\Gamma }>\mathrm{\Gamma }_{cr}`$, respectively. As the initial detuning is zero, $`\mathrm{\Gamma }_{cr}=2/\sqrt{5}=0.89\mathrm{}`$ here, see Eq. (18). Figure 16 shows decaying oscillations on top of the amplitude growth, given by the first of Eqs. (49). Figure 17 shows a non-oscillatory regime of the amplitude growth. Figure 18 shows breakdown of autoresonance when the chirp rate $`m`$ exceeds a critical value.
Similarly to the inviscid case, one can rewrite the governing equations (48) as a single equation for the complex amplitude, $`\psi =B\mathrm{exp}(i\varphi )`$:
$$i\psi _t+\psi ^{}(|\psi |^2mti\mathrm{\Gamma })\psi =0.$$
(51)
Assuming a solution growing in time (that is, phase locked solution), we obtain, at $`t1`$, $`B=|\psi |(mt)^{1/2}`$, as before.
Equations (48) can be also rewritten as a single second-order equation for $`\varphi (t)`$:
$`\ddot{\varphi }+2\mathrm{\Gamma }\dot{\varphi }`$ $`+`$ $`\mathrm{sin}(4\varphi )+2mt\mathrm{sin}(2\varphi )`$ (52)
$``$ $`2\mathrm{\Gamma }\mathrm{cos}(2\varphi )=2\mathrm{\Gamma }mt+m`$
\[compare it with Eq. (26)\]. This equation is convenient for a perturbative treatment in the vicinity of the stable quasi-fixed point. For $`\mathrm{\Gamma }1`$, we can linearize Eq. (52) around $`\varphi =0`$:
$$\ddot{\varphi }+2\mathrm{\Gamma }\dot{\varphi }+4(1+mt)\varphi =2\mathrm{\Gamma }(1+mt)+m.$$
(53)
Now we substitute $`\varphi (t)=\varphi _{}(t)+\delta \varphi (t)`$, where $`\varphi _{}`$ is given by the second of Eqs. (50), and obtain a linear equation
$$\ddot{\delta \varphi }+2\mathrm{\Gamma }\dot{\delta \varphi }+4(1+mt)\varphi =0.$$
(54)
Its approximate solution, at $`m1`$ and $`\mathrm{\Gamma }1`$, can be written as
$$\delta \varphi \frac{Ae^{\mathrm{\Gamma }t}}{(1+mt)^{1/4}}\mathrm{sin}\left[\frac{4(1+mt)^{3/2}}{3m}+\xi \right],$$
(55)
where $`A`$ and $`\xi `$ are constants depending on the initial conditions. The respective solution for $`B(t)`$ is
$$B(1+mt)^{1/2}\frac{Ae^{\mathrm{\Gamma }t}}{(1+mt)^{1/4}}\mathrm{cos}\left[\frac{4(1+mt)^{3/2}}{3m}+\xi \right],$$
(56)
These solutions are simple extensions of the non-viscous perturbative solutions. A comparison with numerical solutions is shown in Fig. 19, and excellent agreement is observed.
In the viscous case the critical value $`m_{cr}`$ of the scaled chirp rate $`m=\mu /(\epsilon \omega _1)^2`$ depends on the scaled damping rate of the wave $`\mathrm{\Gamma }=\gamma /(\epsilon \omega _1)`$. Figure 20 shows this dependence, which we found numerically when starting at $`t=0`$ from the stable fixed point of Eqs. (17) with $`\mathrm{\Delta }=0`$, that is from $`B_{}=(1\mathrm{\Gamma }^2)^{1/4}`$ and $`\varphi _{}=(1/2)\mathrm{arcsin}(\mathrm{\Gamma })`$. One can see that, as $`\mathrm{\Gamma }`$ increases, the critical chirp rate goes down monotonically. As $`\mathrm{\Gamma }`$ approaches $`1`$, $`m_{cr}`$ goes to $`0`$. Notice that, once we return to the physical (dimensional) critical chirp rate $`\mu _{cr}`$ and the wave damping rate $`\gamma `$, the dependence of $`\mu _{cr}`$ on $`\epsilon `$, at fixed $`\gamma `$, is not a power law.
Finally, we tested the accuracy of our reduced equations (48). We compared numerical solutions of Eqs. (48) with numerical solutions of the unreduced equation of motion (19) for the primary mode amplitude $`\eta _1(t)`$, with a linear damping term added. Rescaling the time $`\tau =\omega _1t`$, and the amplitude $`\widehat{\eta }_1=(k_1\eta _1)/(2\sqrt{\epsilon })`$, one can rewrite the unreduced equation as
$`\ddot{\widehat{\eta }}_1`$ $`+`$ $`2\epsilon \mathrm{\Gamma }\dot{\widehat{\eta }}_1+10\epsilon \dot{\widehat{\eta }}_1^2\widehat{\eta }_16\epsilon \widehat{\eta }_1^3+`$ (57)
$`+`$ $`\widehat{\eta }_1\left[1+4\epsilon \mathrm{cos}\left(2\tau m\epsilon ^2\tau ^2\right)\right]=0.`$
A typical example of this comparison is shown in Fig. 21, and a fairly good agreement between the envelope of $`\widehat{\eta }_1(t)`$ and the amplitude $`B(t)`$ is observed.
## V Discussion
This paper presents a theory of weakly nonlinear standing gravity waves, parametrically excited by weak vertical vibrations with a down-chirped vibration frequency. We have shown that autoresonance phase locking and a steadily growing wave amplitude can be achieved despite the nonlinear frequency shift of the wave. For typical initial conditions we have found the critical chirp rate, above which autoresonance breaks down. When starting from a very small wave amplitude, and slowly passing through resonance, phase locking always occurs. We have obtained approximate analytical expressions for the time-dependent wave profile in different regimes. We have demonstrated that each of the three quasi-fixed points of the reduced dynamic equation, describing the primary mode, plays an important role in the dynamics of the system and/or in determining the critical chirp rate.
Parametric autoresonance in Faraday waves is a robust phenomenon, and its experimental observation should not be difficult. To test our theory, one should use a quasi-two-dimensional tank with a low-viscosity liquid, and perform measurements of the standing wave elevation as a function of time. The long-term wave amplitude trend \[see the first of Eqs. (49)\], and the critical chirp rate (see Fig. 20), give examples of quantitative predictions of the theory that can be tested in experiment. One such experiment is presently under way Oded .
The applicability of the weakly nonlinear theory presented in this work is limited to weak forcing and weak damping. Our analysis neglected higher order terms in $`\epsilon `$, which include a nonlinear forcing Miles3 , a cubic damping Miles3 ; douady ; milner , a small correction to the linear detuning/frequency shift, a quintic conservative term decent , etc. These terms can be included in the theory of autoresonance in order to achieve a better accuracy. Importantly, once these higher-order terms are added to the constant-frequency model \[Eq. (17)\], the non-trivial stable fixed point will exist only up to a certain value of the frequency detuning Miles3 ; douady ; milner ; decent . Therefore, in experiment, the autoresonant growth of the wave is expected to terminate when the (time-dependent) frequency detuning comes close to the maximum value, for which the non-trivial stable fixed point in the underlying constant-frequency model still exists. The breakdown of autoresonance is expected to occur at an amplitude smaller than the wave-breaking amplitude, which is comparable to the wavelength Schultz .
## ACKNOWLEDGEMENTS
We thank Oded Ben-David and Jay Fineberg for useful discussions. This research was supported by the Israel Science Foundation.
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# Enhanced Symmetries in Multiparameter Flux Vacua
## 1 Introduction
Flux compactifications of string theory provide an interesting and increasingly well-studied avenue towards making contact with phenomenology. Compactifications with flux can stabilize moduli, potentially eliminating one of the chief difficulties of traditional string theory backgrounds. Although interesting progress in moduli stabilization via fluxes has been made in heterotic , M-theory and most recently type IIA compactifications , the most well-studied case remains that of compactifications of type IIB string theory . In type IIB string theory, fluxes can completely stabilize the complex structure moduli and the dilaton, while the Kรคhler moduli are unfixed; it has been argued that nonperturbative effects can then in some cases stabilze these remaining moduli .
As there are apparently a large number of possible flux compactifications, one is naturally led to ask about the properties that distinguish them. A particularly interesting characteristic to understand is the presence of enhanced symmetries. Other than the natural interest in classifying the different types of low-energy theories making up the โlandscapeโ of vacua , enhanced symmetries are interesting for a number of reasons. Discrete symmetries are useful for model-building in both phenomenological and cosmological contexts. In addition, it has been argued that the fraction of vacua with vanishing superpotential is associated to the prevalence of vacua with realistically low supersymmetry breaking scale ; see also . These vacua also have vanishing cosmological constant at tree level and hence fit onto the โthird branchโ of , which was argued to have different statistics than other branches; this was further studied in . Hence understanding how common such vacua are is of considerable interest.
In , supersymmetric<sup>1</sup><sup>1</sup>1Here we mean vacua solving the F-flatness conditions associated to the complex structure moduli and dilaton; when the superpotential vanishes these are supersymmetric at tree level. flux vacua with vanishing (tree level) superpotential and/or discrete symmetries were studied in compactifications with no more than one complex structure modulus. There it was found that enhanced symmetry vacua occured in some models but not others, though the associated geometries were naively similar. Given that one knows these kinds of vacua are indeed present in some cases, it is clear that it would be desirable to understand how common they are in more general cases.
Moreover, the presence of such vacua seemed in to be associated with a particular mathematical structure appearing in the periods characterizing the geometry: the periods took values in a finite-dimensional algebraic extension of the rational numbers, and this fact was important for obtaining a solution to the vacuum equations. Especially considering the relevance of such arithmetic features of geometries to issues such as attractor vacua , one is naturally led to wonder whether this sort of mathematical structure persists for more complicated models. Other work on the arithmetic properties of Calabi-Yau geometries can be found in .
In this paper, we address the question of the incidence and properties of enhanced symmetry vacua associated to geometries with more than one complex structure modulus. We focus on a particular class of spaces, the hypersurfaces in weighted projective space, though we see no reason why our results should not generalize. We look for constructions of the type of vacuum that was most interesting in : vacua with vanishing superpotential enforced by a preserved discrete R-symmetry. A complementary study, also searching for R-symmetric vacua in weighted projective spaces but using different methods, appears in .
Although naively it seems the vacuum equations associated to the additional moduli will greatly complicate the analysis, we obtain a construction under which these equations become trivial. We show that if certain constraints on the fluxes can be imposed, all vacuum equations follow and the $`W=0`$ vacua exist. A consequence of this is that not all complex structure moduli are fixed at tree level. Whether these constraints on the fluxes can hold is determined by the properties of the modular group transformations for a particular space, and for the simplest (Fermat) hypersurfaces, this can be determined entirely from the integer weights characterizing the projective space. We find that, as for the one-parameter case, for a fixed number of moduli some but not all spaces can support $`W=0`$ vacua. Estimating how many do brings in simple aspects of number theory, and one finds that in the most favorable cases, the number of $`W=0`$ vacua approaches the same scaling as the set of all vacua as the number of moduli becomes large.
Furthermore, we find that although the mathematical structure of the periods continues to play a vital role, in multiparameter models they generally do not live in an algebraic field extension of the rational numbers. Instead, they generate a vector space with basis elements generated by in general transcendental numbers, without any natural product. The properties of this vector space are critical to the existence of vacua.
In section 2, we review flux compactifications of type IIB string theory and the results of on enhanced symmetry vacua. In section 3 we review the relevant properties of the periods of hypersurfaces in weighted projective space. Section 4 describes our construction of enhanced symmetry vacua, while section 5 applies the method to a number of examples, including one outside our class of spaces. In section 6 we conclude.
## 2 Review of enhanced symmetry vacua
### 2.1 Type IIB flux compactifications
We use the same conventions as . We work with a Calabi-Yau 3-fold $``$ and choose a standard symplectic basis for the $`b_3=2(h_{2,1}+1)`$ 3-cycles $`\{A^a,B_b\}`$ with dual cohomology elements $`\alpha _a,\beta ^a`$ obeying
$`{\displaystyle _{A^a}}\alpha _b=\delta _b^a,{\displaystyle _{B_a}}\beta ^b=\delta _a^b,{\displaystyle \alpha _a}\beta ^b=\delta _a^b.`$ (1)
The relevant moduli are the complex combination of the RR axion and the dilaton $`\varphi C_0+ie^\phi `$, as well as the complex structure moduli, which are encoded in the periods of the holomorphic 3-form $`\mathrm{\Omega }`$,
$`z^a={\displaystyle _{A^a}}\mathrm{\Omega },๐ข_b={\displaystyle _{B_a}}\mathrm{\Omega },`$ (2)
where the $`z^a`$ may be used as projective coordinates on the $`h_{2,1}`$-dimensional complex structure moduli space, with the $`๐ข_b=_b๐ข(z)`$ taken as functions of the $`z^a`$. Defining the $`b_3`$-vector of periods $`\mathrm{\Pi }(z)(๐ข_b,z^a)`$, we may write the Kรคhler potential for the dilaton and the complex structure moduli as
$`K`$ $`=`$ $`K^\varphi +K^{cs},K^\varphi =\mathrm{log}(i(\varphi \overline{\varphi })),`$ (3)
$`K^{cs}`$ $`=`$ $`\mathrm{log}(i{\displaystyle _{}}\mathrm{\Omega }\overline{\mathrm{\Omega }})=\mathrm{log}(i\mathrm{\Pi }^{}\mathrm{\Sigma }\mathrm{\Pi }),`$ (4)
where $`\mathrm{\Sigma }`$ is the symplectic matrix $`\mathrm{\Sigma }\left(\begin{array}{c}\mathrm{\hspace{0.33em}0\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}1}\\ \mathrm{1\hspace{0.33em}\hspace{0.33em}0}\end{array}\right)`$.
We turn on the RR and NSNS 3-form field strengths,
$`F_3=(2\pi )^2\alpha ^{}(f_a\alpha ^a+f_{a+h_{2,1}+1}\beta _a),H_3=(2\pi )^2\alpha ^{}(h_a\alpha ^a+h_{a+h_{2,1}+1}\beta _a),`$ (5)
where flux quantization requires the $`b_3`$-vectors $`f`$, $`h`$ to have integer entries. The fluxes generate a tadpole for the charge associated to the $`C_4`$ field, with value
$`N_{\mathrm{flux}}={\displaystyle \frac{1}{(2\pi )^4(\alpha ^{})^2}}{\displaystyle _{}}F_3H_3=f\mathrm{\Sigma }h,`$ (6)
which must be canceled against negative-charge sources such as O3-planes and $`(p,q)`$ 7-branes; in counting flux vacua one typically leaves this charge sink as an arbitrary integer $`L`$ and requires
$`N_{\mathrm{flux}}L,`$ (7)
where the case of an inequality can be made up by mobile D3-branes. Most importantly, the fluxes (5) induce a superpotential for the moduli $`\varphi ,z^a`$ given by
$$W=_{}G_3\mathrm{\Omega }=(2\pi )^2\alpha ^{}(f\varphi h)\mathrm{\Pi },G_3F_3\varphi H_3.$$
(8)
In what follows we shall set $`(2\pi )^2\alpha ^{}=1`$. The Kรคhler moduli $`\rho _i`$ do not participate in $`W`$, and as a result one may derive the no-scale relation
$`{\displaystyle \underset{i}{}}|D_{\rho _i}W|^2=3|W|^2,`$ (9)
leading to the positive-definite potential
$`V=e^K\left(|D_\varphi W|^2+{\displaystyle \underset{a}{}}|D_{z^a}W|^2\right).`$ (10)
We shall neglect the Kรคhler moduli hereafter.
### 2.2 $`W=0`$ vacua
We are interested in vacua that satisfy the F-flatness conditions associated to the superpotential (8),
$`D_{z^a}W_{z^a}W+W_{z^a}K=0,D_\varphi W_\varphi W+W_\varphi K=0,`$ (11)
as well as the satisfying the vanishing of the superpotential,
$`W=0.`$ (12)
For solutions of (11), the tadpole $`N_{\mathrm{flux}}`$ (6) becomes positive definite, and a finite number of vacua exist satisfying (7). As the charge sink $`L`$ becomes large, the number of such vacua scales as
$`N_{\mathrm{vac}}\sqrt{L}^{2b_3},`$ (13)
where $`2b_3`$ is the total number of fluxes. One question is then how many vacua exist additionally satisfying (12). Vacua with $`W=0`$ are supersymmetric at tree level thanks to (9).
It is straightforward to show that (12) combined with the dilaton equation from (11) is equivalent to
$`{\displaystyle F_3}\mathrm{\Omega }={\displaystyle H_3}\mathrm{\Omega }=0f\mathrm{\Pi }(z^a)=h\mathrm{\Pi }(z^a)=0,`$ (14)
where all dependence on the dilaton has dropped out. Hence $`W=0`$ vacua may only exist at points on the complex structure moduli space where the period vector $`\mathrm{\Pi }`$ is orthogonal to integral flux vectors $`f`$, $`h`$; moreover to avoid a vanishing tadpole (6), the vectors $`f`$ and $`h`$ must not be aligned. These periods must furthermore satisfy the remaining equations (11) for the complex structure moduli.
In , examples of hypersurfaces in weighted projective space with a single complex structure modulus $`\psi `$ were considered. The equations (14) were studied at the Landau-Ginzburg point $`\psi =0`$, where the period vector $`\mathrm{\Pi }`$ acquired a particular arithmetic structure: up to an overall constant, the components of $`\mathrm{\Pi }`$ took values in the cyclotomic field $`_d`$, the extension of the rationals by $`d^{th}`$ roots of unity, with $`d`$ depending on the geometry.
This field can be viewed as a vector space over the rationals; hence $`f\mathrm{\Pi }`$ and $`h\mathrm{\Pi }`$ can be expanded in a basis for this vector space, where the basis vectors of $`_d`$ are generally irrational numbers, and their rational coefficients are determined by $`f`$ and $`h`$. The existence of $`W=0`$ vacua is then conditional on whether integer fluxes $`f`$, $`h`$ can be found to cancel independently every basis vector in this space. If this is done, one may always use the final $`D_\psi W=0`$ equation to solve for the dilaton $`\varphi `$, as $`\varphi `$ enters in no other equations.
The ability to find $`W=0`$ vacua thus depends on the dimension of this โvector space of the periodsโ, as a single equation like $`f\mathrm{\Pi }=0`$ will decompose into a number of relations on the fluxes equal to this dimension. The hypersurface called $`M_6`$ in had dimension two, and solutions were possible, while for the other three cases, the dimension was four and not enough freedom was present in the fluxes to find solutions. Thus the four one-parameter models, naively quite similar, had dramatically different spectra of $`W=0`$ vacua.
### 2.3 R-symmetries
The $`W=0`$ vacua in the $`M_6`$ model are accompanied by a discrete R-symmetry, a transformation of the moduli under which $`W`$ changes by a phase, โenforcingโ the vanishing of the superpotential at its fixed point $`\psi =0`$. It is worth explaining what constitutes a discrete symmetry in these models.
The string theory compactifications we consider each possess a modular group of symmetries acting both on the moduli $`z^a,\varphi `$ and on the fluxes $`f`$, $`h`$. The $`SL(2,)`$ S-duality of type IIB string theory is one example, as are geometric transformations of the complex structure moduli space. These symmetries are most usefully thought of as discrete gauge symmetries, since they correspond to a redundancy of the description: two apparent vacua related by such a transformation actually constitute only one genuine vacuum, and this redundancy must be accounted for to properly count vacua.
The symmetries we are interested in are not modular transformations, though they are related. We want to consider global symmetries of the low-energy effective field theory for the moduli. In this effective theory, the fluxes appear not as fields, but as coupling constants. Consequently our R-symmetries will transform the moduli fields, but in order to stay within a given effective field theory description, they must not act on the couplings.
In fact, the R-symmetries we find act on the moduli in the same fashion as the modular group, but leave out the transformation on the fluxes. Unlike the modular transformations, which are always symmetries, such a moduli-only transformation need not be a good symmetry: they are symmetries only for special values of the fluxes. When present, these are genuine global symmetries relating distinct vacua, not redunancies of the description.<sup>2</sup><sup>2</sup>2Some models have modular transformations that happen to leave the fluxes invariant; these should be considered gauge symmetries as well, and represent an identification on the moduli space in the low-energy theory. Such a symmetry descending from a modular transformation can arise for values of the fluxes that place the vacuum on the fixed locus of that modular symmetry , and we will use this princple in constructing $`W=0`$ vacua with R-symmetries.
### 2.4 Extension to multiparameter cases
The goal in this paper is to extend the construction of $`W=0`$ vacua with R-symmetries for one-parameter cases of to hypersurfaces with additional complex structure parameters $`\phi ^A`$. We shall continue to look at the Landau-Ginzburg point $`\psi =0`$ and will look for R-symmetries descending from modular transformations rotating around this point. The additional parameters complicate the analysis in two ways. First, the periods now depend on the additional moduli $`\phi ^A`$, affecting the solution to the โflux orthogonalityโ relations (14). Second, additional equations $`D_{\phi ^A}W=0`$ also need to be solved.
We shall describe how given a certain kind of solution to (14) the structure of the periods $`\mathrm{\Pi }`$ are such that the $`D_{\phi ^A}W=0`$ equations follow trivially. The $`\phi ^A`$ moduli are then unconstrained by the fluxes, and the $`D_\psi W=0`$ equation can be used to determine the dilaton as in the one-parameter case.
Finding $`W=0`$ vacua thus will come down to finding these particular solutions of (14). As we shall show, at $`\phi ^A=0`$ the periods again take values in a cyclotomic field. Away from the locus $`\phi ^A=0`$, while they in general lack arithmetic properties, the periods still take values in a finite-dimensional vector space over the rationals. The particular solutions in question are just choices of fluxes that solve (14) for this (in general larger) vector space at $`\phi ^A0`$. As in the one-parameter case, where one out of four hypersurfaces supported $`W=0`$ vacua, we shall find that while not all spaces support such vacua, a substantial number do.
The construction we employ to obtain these vacua utilizes the matrix that generates monodromy around the $`\psi =0`$ locus on moduli space, which will make manifest that all such $`W=0`$ vacua are accompanied by R-symmetries. Before we can describe this construction, we must describe the geometries and their periods in more detail.
## 3 Geometric considerations
In this section we review relevant properties of hypersurfaces in weighted projective space, their mirrors and their periods. Nothing in this section is new, though our emphasis on the โperiod indicesโ $`n_I`$ and the Euler totient function $`\varphi (d)`$ may be unfamiliar.
### 3.1 Hypersurfaces in weighted projective space
The period vector $`\mathrm{\Pi }(z)`$ depends on the particular geometry considered. The class of Calabi-Yau threefolds we work with are all hypersurfaces in weighted projective space. The weighted projective space $`๐_{k_1k_2k_3k_4k_5}^4`$ has weighted homogeneous coordinates $`x_i\lambda ^{k_i}x_i`$ and degree $`d_{i=1}^5k_i`$. We then define the Calabi-Yau manifold $``$ as the vanishing locus (with singularities resolved) of the polynomial
$`P(x_i)P_0(x_i)d\psi x_1x_2x_3x_4x_5+{\displaystyle \underset{\stackrel{~}{A}}{}}\phi ^{\stackrel{~}{A}}M_{\stackrel{~}{A}}(x_i).`$ (15)
Here $`P_0(x_i)`$ is a suitable โdefining polynomialโ
$`P_0(x_i)={\displaystyle \underset{j=1}{\overset{5}{}}}{\displaystyle \underset{i=1}{\overset{5}{}}}x_i^{a_{ij}},{\displaystyle \underset{i=1}{\overset{5}{}}}k_ia_{ij}=dj.`$ (16)
The $`M_{\stackrel{~}{A}}(x_i)`$ are a set of monomials associated with complex structure variations $`\phi ^{\stackrel{~}{A}}`$:
$`M_{\stackrel{~}{A}}(x_i)={\displaystyle \underset{i=1}{\overset{5}{}}}x_i^{q_{\stackrel{~}{A}}^i},{\displaystyle \underset{i=1}{\overset{5}{}}}k_iq_{\stackrel{~}{A}}^i=d\stackrel{~}{A},`$ (17)
and in (15) we have separated out the โfundamental monomialโ $`M_0=dx_1x_2x_3x_4x_5`$ and denoted its modulus $`\psi `$. In general not all complex structure deformations can be expressed in terms of a monomial $`M_{\stackrel{~}{A}}`$; these non-polynomial deformations are mirror to Kรคhler moduli that are not toric. We shall restrict to the case of spaces with only monomial deformations here.
The simplest class of such hypersurfaces occurs when each $`k_i`$ divides $`d`$. In this case, $`a_{ij}`$ can be chosen diagonal and $`P_0`$ is then a Fermat polynomial:
$`P_0^{\mathrm{Fermat}}(x_i)=x_1^{d/k_1}+x_2^{d/k_2}+x_3^{d/k_3}+x_4^{d/k_4}+x_5^{d/k_5}.`$ (18)
We call these Fermat cases. In the non-Fermat examples, all non-degenerate models have either a $`0`$ or a $`1`$ for the off-diagonal elements of $`a_{ij}`$, and take a block-diagonal form with Fermat blocks as well as either โtadpoleโ $`x_1^{a_{11}}x_2+x_2^{a_{22}}x_3+\mathrm{}+x_n^{a_{nn}}`$ or โloopโ $`x_1^{a_{11}}x_2+x_2^{a_{22}}x_3+\mathrm{}+x_n^{a_{nn}}x_1`$ blocks; see for example .
The defining polynomial $`P_0`$ (16) has a set of discrete phase rotation symmetries $`๐ฏ_j`$: each Fermat-type monomial $`x_j^{a_{jj}}`$ gives rise to a $`_{a_{jj}}`$, while tadpole and loop blocks give rise to $`_{a_{11}a_{22}\mathrm{}a_{nn}}`$ and $`_{a_{11}a_{22}\mathrm{}a_{nn}+(1)^{n1}}`$ symmetries, respectively. One linear combination of the $`๐ฏ_j`$, the โquantum symmetryโ $`Q_{}_j๐ฏ_j_d`$, is a trivial identification of the homogeneous coordinates; the remaining symmetries of $`P_0`$ form the geometric symmetries $`G_{}`$.
The mirror of $``$ is formed as follows<sup>3</sup><sup>3</sup>3We assume $`P`$ is transverse; more general methods involving toric geometry are discussed for example in . . Define the dual polynomial $`\widehat{P}_0`$ with transposed exponents, $`\widehat{a}_{ij}a_{ji}`$:
$`\widehat{P}_0(y_i)={\displaystyle \underset{j=1}{\overset{5}{}}}{\displaystyle \underset{i=1}{\overset{5}{}}}y_i^{\widehat{a}_{ij}},`$ (19)
which lives in a projective space $`๐_{\widehat{k}_1\widehat{k}_2\widehat{k}_3\widehat{k}_4\widehat{k}_5}^4`$ with coordinates $`y_i\lambda ^{\widehat{k}_i}y_i`$ and degree $`\widehat{d}_{i=1}^5\widehat{k}_i`$. The vanishing of this polynomial defines a space $`\widehat{๐ฒ}`$, and has phase symmetries $`\widehat{๐ฏ}_j`$ divided into $`Q_{\widehat{๐ฒ}}=_{\widehat{d}}`$ and $`G_{\widehat{๐ฒ}}`$ as before. The mirror $`๐ฒ`$ of $``$ is then the quotient $`\widehat{๐ฒ}/H`$, where $`H`$ is the subgroup of $`G_{\widehat{๐ฒ}}`$ under which the fundamental monomal $`y_1y_2y_3y_4y_5`$ is invariant. One finds that $`G_{\widehat{๐ฒ}}=H\times _d`$, such that after the quotient one has the symmetries
$`Q_๐ฒ=G_{}=H\times _{\widehat{d}},G_๐ฒ=Q_{}=_d.`$ (20)
The polynomial $`\widehat{P}`$ defining the mirror is then
$`\widehat{P}\widehat{P}_0(y_i)\widehat{d}\psi y_0y_1y_2y_3y_4+{\displaystyle \underset{\alpha =1}{\overset{m}{}}}\phi ^AM_A(y_i),`$ (21)
where the $`M_A(y_i)`$ are monomials of degree $`\widehat{d}`$ invariant under $`H`$.
In what follows we will work with the mirrors $`๐ฒ`$, for which there are a set of well-studied simple cases with a small number of complex structure moduli due to the restriction to $`H`$-invariant monomials; the corresponding $``$ have a small number of Kรคhler moduli but in general many complex structure moduli, and are thus less useful as examples. Our results should, however, generalize to the $``$ models as well.
For Fermat cases, where the $`k_i`$ divide $`d`$ and $`a_{ij}`$ is diagonal, we have $`\widehat{d}=d`$ and $`๐ฒ`$ is just a quotient of $`=\widehat{๐ฒ}`$; $`\widehat{P}`$ is then simply a truncation of $`P`$ to monomials invariant under $`H`$, and one may use $`x_i`$ in place of $`y_i`$ since they are defined in the same space. We may then think of merely restricting to $`H`$-invariant monomials in $``$ instead of actually taking the mirror; as discussed in , turning on $`H`$-invariant fluxes consistently leads to solutions in this subspace of the moduli space of $``$.
### 3.2 Periods
The periods for $`๐ฒ`$ may then be defined as follows . One begins with the fundamental period $`\varpi _0`$, defined by a canonical choice of holomorphic three-form and of a $`B_0`$-cycle. The integral over the cycle is evaluated and one obtains a series in $`1/\psi `$. For $`\phi ^A=0`$ this can be written as a generalized hypergeometric function:
$`\varpi _0=_qF_{q1}(\widehat{n}_I;\widehat{m}_J,1,1,1;({\displaystyle \underset{i=1}{\overset{5}{}}}k_i^{k_i}\psi )^d),`$ (22)
where the $`\{\widehat{n}_I\}`$ and $`\{\widehat{m}_J\}`$ are determined as follows: compare the set $`\{1/d,2/d,\mathrm{}(d1)/d\}`$ with the set $`\{1/k_1,\mathrm{}(k_11)/k_1,\mathrm{}1/k_5,\mathrm{}(k_51)/k_5\}`$, and remove from both sets any number appearing in both (but only remove it once in each for each match); what remains are the $`\widehat{n}_I`$ and the $`\widehat{m}_J`$, respectively. The integer $`qd1`$ is the number of the $`\widehat{n}_I`$. Note that (22) is a period of the mirror $`๐ฒ`$, but is given in terms of the weights $`k_i`$ and degree $`d`$ of the original space $``$.
The fundamental period can be analytically continued to the region of small $`\psi `$ via an integral of Barnes type; this is done in general in . The result is
$`\varpi _0={\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\mathrm{\Gamma }(\frac{n}{d})\alpha ^{n(d1)/2}(k_1^{k_1/d}d\psi )^{n1}}{\mathrm{\Gamma }(n)_{r=1}^{k_11}\mathrm{\Gamma }(\frac{r}{k_1}\frac{n}{d})_{i=2}^5\mathrm{\Gamma }(1\frac{k_in}{d})}},`$ (23)
where $`k_1`$ is the smallest of the $`k_i`$, $`\alpha `$ is a $`d^{th}`$ root of unity:
$`\alpha ^d=1,`$ (24)
and we have renormalized by an extra factor $`(2\pi )^{(1k_1)/2}/\psi `$ relative to (22) and to ; this Kรคhler transformation assures the periods and thus the Kรคhler potential (3) are regular at $`\psi =0`$, as $`W=0`$ is not predictive if $`e^K`$ diverges.
In what follows we specialize for simplicity to $`k_1=1`$, which is obeyed by all our examples. The inclusion of the neglected parameters $`\phi ^A`$ then gives the fundamental period
$`\varpi _0={\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\mathrm{\Gamma }(\frac{n}{d})\alpha ^{n(d1)/2}(d\psi )^{n1}}{\mathrm{\Gamma }(n)_{i=2}^5\mathrm{\Gamma }(1\frac{k_in}{d})}}U_n(\phi ^A){\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}c_nU_n(\phi ^A)\psi ^{n1},`$ (25)
with $`U_n(\phi ^A=0)=1`$; this will in general only be well-defined in some neighborhood of $`\phi ^A=0`$. The functional form of the $`U_n(\phi ^A)`$ will not matter for us.
The remaining periods are generated from the fundamental period $`\varpi _0`$ by a monodromy action associated to the discrete symmetries $`G_๐ฒ_d`$. The full polynomial (15) with nonzero $`\psi ,\phi ^A`$ is not preserved by $`G_๐ฒ`$ transformations; instead, the $`G_๐ฒ`$ action can be canceled by also transforming the moduli $`\psi ,\phi ^A`$. This generates an identification on the moduli space of $`\psi ,\phi ^A`$, since transforming the moduli can be absorbed into a change of coordinates. This is a subgroup $`๐`$ of the full monodromy group of the moduli space, with $`๐G_๐ฒ_d`$.
The monomials $`M_A`$ invariant under $`H`$ transform nontrivially under $`G_๐ฒ`$. In particular, the fundamental monomial always faithfully represents $`_d`$. This induces the monodromy action on the moduli space:
$`๐:\psi \alpha \psi ,\phi ^A\alpha ^{Q_A}\phi ^A.`$ (26)
The periodicities of the $`\phi ^A`$ determined by $`Q_A`$ will be important for us. For cases with $`k_1=1`$ one can take $`_d`$ to be generated by $`x_1\alpha ^1x_1`$, in which case we have $`Q_A=q_A^1`$, with $`q_A^1`$ as in (17).
From the fundamental period one may then use the monodromy $`๐`$ to define $`d1`$ additional periods:
$`\varpi _J(\psi ,\phi ^A)\varpi _0(\alpha ^J\psi ,\alpha ^{Q_AJ}\phi ^A)={\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}c_n\alpha ^{nJ}U_n(\alpha ^{Q_AJ}\phi ^A)\psi ^{n1}.`$ (27)
In general not all $`d`$ of the periods so constructed are independent, instead $`q`$ are independent where $`q`$ is the index of the hypergeometric (22) . This generates all the periods associated to complex structure moduli that can be written as monomial deformations of the fundamental polynomial; since we focus on examples where all moduli are of this sort, we have generated all periods from (27) and write $`b_3=q`$. We can choose a basis $`\varpi _0,\mathrm{}\varpi _{b_31}`$ and arrange them into a vector $`\stackrel{}{\varpi }`$:
$`\stackrel{}{\varpi }(\psi ,\phi ^A)={\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}c_n\stackrel{}{p}_n(\phi ^A)\psi ^{n1},c_n={\displaystyle \frac{\mathrm{\Gamma }(\frac{n}{d})\alpha ^{n(d1)/2}d^{n1}}{\mathrm{\Gamma }(n)_{i=2}^5\mathrm{\Gamma }(1\frac{k_in}{d})}}.`$ (28)
with $`(p_n)_J(\phi ^A)\alpha ^{nJ}U_n(\alpha ^{Q_AJ}\phi ^A)`$. We shall have special use for the term at $`\psi =0`$, and denote the corresponding functions at $`n=1`$ as $`p_J(\phi ^A)=\alpha ^JU(\alpha ^{Q_AJ}\phi ^A)`$.
Finally, the symplectic basis for the periods $`\mathrm{\Pi }`$ (which enters into the equations (11, 12) for vacua) is obtained from the Picard-Fuchs basis $`\varpi `$ via a transformation
$`\mathrm{\Pi }_I(\psi ,\phi ^A)=m_{IJ}\varpi _J(\psi ,\phi ^A),`$ (29)
where the matrix $`m_{IJ}`$ has rational entries . As we shall discuss, the precise form of $`m_{IJ}`$ will not be important for us. The expression (28), (29) for the periods is the result we will use in the following sections.
Note that not all values of $`n`$ in the expansion (28) have nonzero coefficients $`c_n`$. The $`\mathrm{\Gamma }`$-functions in the denominator have poles, killing the coefficient, whenever
$`n={\displaystyle \frac{d\mathrm{}}{k_i}},`$ (30)
for any positive integer $`\mathrm{}`$ and any $`k_i`$ with $`i=2,3,4,5`$. It is easy to see that the values of $`n`$ having nonzero $`c_n`$ are periodic mod $`d`$; and moreover those values $`1nd`$ with nonzero $`c_n`$, call them the $`n_I`$, are precisely the integers
$`n_I=d\widehat{n}_I,`$ (31)
with the $`\widehat{n}_I`$ as given below (22). We call the $`n_I`$, which are $`b_3`$ in number, the period indices.
The period indices will play an essential role in our analysis, so we mention a few of their properties. They are determined solely by the $`k_i`$ defining the weighted projective space. The set contains at minimum those integers $`1n_I<d`$ that share no common factors with $`d`$. The number of such integers is the definition of the Euler totient function $`\varphi (d)`$. Hence we have in general
$`\varphi (d)b_3<d.`$ (32)
As we shall see, whether $`b_3>\varphi (d)`$ or $`b_3=\varphi (d)`$ will determine whether there may be or may not be R-symmetric $`W=0`$ vacua at $`\psi =0`$ in a given model.
## 4 Construction of $`W=0`$ vacua
With this description of the periods in hand, we go on to construct vacua satisfying $`DW=W=0`$. The equations $`D_\varphi W=W=0`$ (14) at $`\psi =0`$ become using (28)
$`\stackrel{~}{f}p(\phi ^A)=\stackrel{~}{h}p(\phi ^A)=0,`$ (33)
where $`p(\phi ^A)`$ is (up to an overall coefficient) the $`\varpi `$ period vector at $`\psi =0`$, and $`\stackrel{~}{f}_Jf_Im_{IJ}`$ and $`\stackrel{~}{h}_Jh_Im_{IJ}`$ are vectors of rational numbers. If a solution to these โflux orthogonalityโ equations exists for rational $`\stackrel{~}{f}`$, $`\stackrel{~}{h}`$, we can always scale up to find a set of integer solutions for $`f`$, $`h`$; hence merely to demonstrate existence of vacua, we do not need to know the values of $`m_{IJ}`$, just that they are rational. We first discuss the algebraic properties of these equations, before turning to the remaining $`DW=0`$ equations and a construction for a solution.
### 4.1 The vector space of periods
For given $`\phi ^A`$, the elements of the vector $`p_J(\phi ^A)\alpha ^JU(\alpha ^{Q_AJ}\phi ^A)`$ are in general not rational, but instead generate a vector space $`๐ฑ`$ over the rationals, the โvector space of periodsโ, with basis
$`\mathrm{basis}๐ฑ=\{U(\phi ^A),\alpha U(\alpha ^{Q_A}\phi ^A),\mathrm{}\alpha ^{b_31}U(\alpha ^{Q_A(b_31)}\phi ^A)\}.`$ (34)
Satisfying (33) by tuning the rational rescaled fluxes $`\stackrel{~}{f}`$, $`\stackrel{~}{h}`$ can only be done by separately setting to zero the rational coefficient of each basis element of this vector space. It is thus vital for the possibility of a solution that the $`p_J`$ (34) are not linearly independent over the rationals; otherwise the only solution would be that all fluxes vanish. This lack of independence will arise because of a combination of two effects.
First, the $`d`$ basis elements of the cyclotomic field $`_d`$, $`\{1,\alpha ,\alpha ^2,\mathrm{}\alpha ^{d1}\}`$, are in general not independent when $`d`$ is not prime; only $`\varphi (d)`$ are independent, where $`\varphi (d)`$ is the Euler totient function defined in the last section. In one-parameter models, the functions $`U`$ are absent and the vector space of periods $`๐ฑ`$ is simply the cyclotomic field $`_d`$ extending the rationals by the root of unity $`\alpha `$ , and one has $`dim๐ฑ=\varphi (d)`$.
Secondly, in the multiparameter cases one also has the functions $`U(\phi ^A)`$, which in general take transcendental values. However, the $`U(\alpha ^{Q_AJ}\phi ^A)`$ need not be distinct for all $`J`$, as $`Q_A`$ may share a common factor with $`d`$, leading to $`U(\alpha ^{Q_AJ}\phi ^A)=U(\alpha ^{Q_A(J+D_A)}\phi ^A)`$ for some $`D_A`$. As a result, the function $`U(\phi ^A)`$ has an overall periodicity set by $`D=\mathrm{lcm}(\{D_A\})`$. Thus we find that $`p_{J+D}(\phi ^A)=\alpha ^Dp_J(\phi ^A)`$, and such elements may be linearly dependent regardless of the value of the $`\phi ^A`$.
Hence the $`p_J(\phi ^A)`$ will in general obey linear relationships for arbitrary values of $`\phi ^A`$, defining a vector space $`๐ฑ`$ over the rationals with $`\phi ^A`$-dependent basis elements, satisfying
$`\varphi (d)dim๐ฑb_3<d.`$ (35)
By tuning the $`f`$ and $`h`$ fluxes to obtain a zero coefficient for each basis element of $`๐ฑ`$, one finds a solution for (33) for any $`\phi ^A`$. If the dimension of $`๐ฑ`$ is equal to $`b_3`$, this is only possible for $`f=h=0`$. On general grounds one then expects that the dimension of the space of orthogonal flux vectors solving either equation in (33) is given by
$`b_3dim๐ฑ.`$ (36)
For fixed $`d`$ and $`b_3`$ this is largest when $`dim๐ฑ=\varphi (d)`$, which can only occur when the overall periodicity $`D`$ is a factor of $`\varphi (d)`$.
As an example, consider the two-parameter model $`(k_1,k_2,k_3,k_4,k_5)=(1,1,2,2,6)`$; this was studied in and will be further analyzed in section 5. This model has $`d=12`$, and the cyclotomic field $`_{12}`$ has $`\varphi (12)=4`$ linearly independent elements, which can be taken to be $`1,\alpha ,\alpha ^2,\alpha ^3`$ with the rest related via
$`\alpha ^4=\alpha ^21,\mathrm{for}\alpha ^{12}=1.`$ (37)
The monomial associated to the parameter $`\phi `$ is $`M=x_1^6x_2^6`$ with $`Q=q^1=6`$, meaning $`U(\alpha ^Q\phi )=U(\phi )`$, hence with overall periodicity $`D=2`$. For a two-parameter model there are $`b_3=6`$ elements of $`p_J(\phi )`$ generating the vector space of periods $`๐ฑ`$,
$`\mathrm{basis}๐ฑ_{11126}=\{U(\phi ),\alpha U(\phi ),\alpha ^2U(\phi ),\alpha ^3U(\phi ),\alpha ^4U(\phi ),\alpha ^5U(\phi )\},`$ (38)
but using (37), only the first four are independent; hence in this case $`dim๐ฑ=\varphi (d)=4`$, thanks to the compatibility between the overall periodicity $`D`$ and the dimension of the cyclotomic field $`\varphi (d)`$.
For a more complicated case where $`dim๐ฑ>\varphi (d)`$, consider another $`d=12`$ space, $`(1,2,3,3,3)`$. The cyclotomic field is the same as the previous example, and the relation (37) still holds, but this time there are two additional parameters $`\phi ^1`$, $`\phi ^2`$ with monomials $`M_1=x_1^4x_2^4`$ and $`M_2=x_1^8x_2^2`$, with $`Q_1=4`$ and $`Q_2=8`$, giving rise to overall periodicity $`D=3`$. This three-parameter model has the vector space of periods generated by the $`b_3=8`$ elements of $`p_J(\phi ^A)`$,
$`\mathrm{basis}๐ฑ_{12333}`$ $`=`$ $`\{U(\phi ^1,\phi ^2),\alpha U(\alpha ^4\phi ^1,\alpha ^8\phi ^2),\alpha ^2U(\alpha ^8\phi ^1,\alpha ^4\phi ^2),\alpha ^3U(\phi ^1,\phi ^2),`$
$`\alpha ^4U(\alpha ^4\phi ^1,\alpha ^8\phi ^2),\alpha ^5U(\alpha ^8\phi ^1,\alpha ^4\phi ^2),U(\phi ^1,\phi ^2),\alpha U(\alpha ^4\phi ^1,\alpha ^8\phi ^2)\},`$
where thanks to $`\alpha ^6=1`$ we see the last two elements are not independent; however despite the relation (37) we cannot write the $`\alpha ^4`$ and $`\alpha ^5`$ elements in terms of the others, because the periodicity of the $`U`$โs is not compatible, a reflection of $`\varphi (d)=4`$ not being a multiple of $`D=3`$. Hence in this example, $`dim๐ฑ=6>\varphi (d)=4`$.
At a specific value of $`\phi ^A`$ the vector space in which the periods live may become smaller; for example at $`\phi ^A=0`$ it reduces simply to $`_d`$. A more complicated example is described for the $`(1,1,2,2,6)`$ two-parameter model in sec. 8.3 of , where it is shown that solving the conditions for โattractor pointsโ at $`\psi =0`$ leads to a set of discrete values for $`\phi `$ where the periods take values in the rank two vector space (and field) $`[i]`$. However, as we now discuss, solving the additional $`D_{\phi ^A}W=0`$ equations in the fashion we describe necessitates a solution valid for a continuous range of $`\phi ^A`$.
### 4.2 Solving remaining complex structure equations
We will find that solutions to the flux orthogonality equations (33) for arbitrary $`\phi ^A`$ are not uncommon. Consider such a solution, and try to solve the equations $`D_{\phi ^A}W=0`$. Then since $`W=0`$, each equation reduces to
$`_AW=(\stackrel{~}{f}\varphi \stackrel{~}{h})_Ap(\phi ^A)=0.`$ (40)
One has $`_Ap_J(\phi ^A)=\alpha ^J_AU(\alpha ^{Q_AJ}\phi ^A)`$. The essential point is that for each $`A`$, the $`_Ap_J(\phi ^A)`$ will generate a vector space $`๐ฑ^{(A)}`$ with basis elements obeying the same linear relationships as the basis elements of $`๐ฑ`$ generated by the $`p_J(\phi ^A)`$.
This occurs because the $`\alpha ^J`$ coefficient is does not change, and the periodicities of the $`U`$โs โ which are all that go into determining the linear relations on the basis vectors โ are not disturbed by the derivative. So although the elements of the basis of the space $`๐ฑ^{(A)}`$ will in general be different transcendental numbers than those of $`๐ฑ`$, their coefficients in (40) will be zero if the coefficients in (33) are.
Moreover, the $`f`$ and $`h`$ parts are separately zero, meaning that again the dilaton $`\varphi `$ drops out of the equation. Since the dilaton is completely unconstrained thus far, we may simply use the final equation $`D_\psi W=0`$ to solve for it:
$`\varphi ={\displaystyle \frac{f_\psi \mathrm{\Pi }}{h_\psi \mathrm{\Pi }}}.`$ (41)
Hence we find that once a solution to (33) with the periods orthogonal to the flux vectors is found for arbitrary $`\phi ^A`$ , the remaining $`DW=0`$ equations can always be satisfied.
One can consider finding a solution to (33) valid only at a single point $`\phi _0^A`$; however if this solution does not extend to a continuous family of solutions over the space of $`\phi ^A`$, the $`D_{\phi ^A}W=0`$ solutions will not in general be satisfied. It is not impossible that the ranks of $`๐ฑ`$ and of the $`๐ฑ^{(A)}`$ may all reduce at the same point, or $`๐ฑ`$ and the $`๐ฑ^{(A)}`$ may at a certain point become isomorphic. Such circumstances would lead to isolated $`W=0`$ vacua; whether they exist is an interesting open questions.
### 4.3 Ansatz for vacua
Consequently, finding $`W=0`$ vacua reduces to solving the flux orthogonality equations (33) over a continuous range of $`\phi ^A`$. To do so, we are interested finding integer $`b_3`$-vectors $`g_J`$ obeying
$`g\mathrm{\Pi }(\psi =0,\phi ^A)=0.`$ (42)
A solution to (33) will then follow setting $`f_J=g_J^{(1)}`$ and $`h_J=g_J^{(2)}`$ for two such vectors $`g^{(1)}`$ and $`g^{(2)}`$, which must be non-parallel in order to produce nonzero tadpole (6). Our method for finding such a vector is essentially to require the existence of an R-symmetry in the low-energy effective theory. We do this as follows.
Consider the monodromy action $`๐`$ (26) generating phase rotations around the LG point $`\psi =0`$. It will be realized on the periods $`\mathrm{\Pi }_J`$ as a $`b_3\times b_3`$ matrix $`A_{IJ}`$ representing $`_d`$, hence obeying $`A^d=1`$. It has the general form
$`A=m\left(\begin{array}{c}\mathrm{\hspace{0.33em}0\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}1\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}0}\mathrm{}\mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}0}\\ \mathrm{\hspace{0.33em}0\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}0\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}1}\mathrm{}\mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}0}\\ \mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}\\ \mathrm{\hspace{0.33em}0\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}0\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}0}\mathrm{}\mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}1}\\ \mathrm{}\end{array}\right)m^1,`$ (43)
where the rational entries on the final line are determined by the expansion of $`\varpi _{b_3}`$ in the other $`\varpi _J`$ for each particular model. The monodromy action of this matrix on the periods $`\mathrm{\Pi }`$ translates into
$`A\mathrm{\Pi }(\psi ,\phi ^A)=\alpha \mathrm{\Pi }(\alpha \psi ,\alpha ^{Q_A}\phi ^A),`$ (44)
where the overall factor of $`\alpha `$ comes from our normalization of the $`\varpi _J`$.
As described in section 2, an R-symmetry must be a transformation on the moduli that leaves the fluxes invariant, such that the superpotential rotates by a nontrivial phase. Consider the action of $`๐`$ (26) some number $`N`$ times on the moduli, while holding the fluxes invariant. One has
$`W(\varphi ,\psi ,\phi ^A)=(f\varphi h)\mathrm{\Pi }(\psi ,\phi ^A)\alpha ^N(f\varphi h)A^N\mathrm{\Pi }(\psi ,\phi ^A).`$ (45)
This can become an overall phase rotation of $`W`$ if $`f`$ and $`h`$ are left eigenvectors of $`A^N`$.
As $`A^d=1`$, the eigenvalues of $`A`$ and its powers are $`d^{th}`$ roots of unity. Since $`A`$ and $`f`$, $`h`$ all contain rational numbers, the only possible eigenvalues are $`\pm 1`$. Correspondingly, our ansatz consists of looking for flux vectors $`g`$ satisfying the โR-symmetry constraintโ
$`g=gA^N,`$ (46)
for some integer $`1N<d`$. An R-symmetry implies the vanishing of the superpotential when the vacuum lies on the fixed locus of the symmetry, and $`g\mathrm{\Pi }`$ vanishes for analogous reasons. Given (44) and (46), we have
$`g\mathrm{\Pi }(\psi =0,\phi ^A)=gA^N\mathrm{\Pi }(\psi =0,\phi ^A)=\alpha ^Ng\mathrm{\Pi }(\psi =0,\alpha ^{NQ_A}\phi ^A).`$ (47)
We then find that flux orthogonality (42) follows if the R-symmetry constraint (46) is accompanied by an additional periodicity condition
$`\alpha ^{NQ_A}=1Q_A.`$ (48)
As discussed, the $`\phi ^A`$ have an overall periodicity $`D`$; (48) then requires $`N=\mathrm{}D`$ for integer $`\mathrm{}`$. Taking $`f`$, $`h`$ proportional to such $`g`$โs then implies the existence of a $`_{d/N}`$ R-symmetry
$`\psi \alpha ^N\psi ,\phi ^A\phi ^AW\alpha ^NW.`$ (49)
Since powers of $`A`$ may also have eigenvalue $`1`$, it is possible that the square root of the R-symmetry constraint may also hold:
$`g=gA^{N/2}.`$ (50)
Hence if $`N/2`$ is even, the R-symmetry may be promoted to $`_{2d/N}`$:
$`\psi \alpha ^{N/2}\psi ,\phi ^A\alpha ^{NQ_A/2}\phi ^AW\alpha ^{N/2}W.`$ (51)
Note that $`N/2`$ need not necessarily be compatible with the periodicity condition.
### 4.4 The R-symmetry constraint and eigenvalues of $`A`$
Obtaining $`W=0`$ vacua thus comes down to imposing the R-symmetry constraint (46) while satisfying the periodicity condition (48). To understand when we can impose the R-symmetry constraint, we consider the eigenvalues of $`A`$.
Since $`A^d=1`$, each eigenvalue is a $`d^{th}`$ root of unity, $`\alpha ^l`$ for some $`l`$. In principle we can calculate these by determining the relation between $`\varpi _{b_3}`$ and the other $`\varpi _J`$ to fill in the last line of the matrix (43), and diagonalize it. However, the eigenvalues may be determined more simply either by using the explicit formula for the periods $`\varpi `$ (27) expanded in a power series in $`\psi `$, or simply from the monodromy action (44). Using these one can show that
$`\mathrm{\Pi }^{(n)}{\displaystyle \underset{J=0}{\overset{b_31}{}}}_\psi ^n\mathrm{\Pi }(\psi ,\alpha ^{Q_AJ}\phi ^A)|_{\psi =0}c_{n+1}{\displaystyle \underset{J=0}{\overset{b_31}{}}}p_{n+1}(\alpha ^{Q_AJ}\phi ^A),`$ (52)
is an eigenvector of $`A`$, with eigenvalue $`\alpha ^{n+1}`$.
However, not all $`\mathrm{\Pi }^{(n)}`$ are nonzero as some of the coefficients $`c_n`$ vanish. The coefficients are nonzero precisely for the values $`n`$ mod $`d=n_J`$, with the period indices $`n_J`$ given in (31). Hence the eigenvalues $`a_J`$ of $`A`$ are just
$`a_J=\alpha ^{n_J}.`$ (53)
We have established previously that there are exactly $`b_3`$ nonzero $`n_J`$, giving the correct number of eigenvalues for $`A`$.
These eigenvalues represent $`_d`$ and hence always obey $`a_J^d=1`$. However, it is not necessarily the case that a given eigenvalue faithfully represents $`_d`$. When $`d`$ is not prime, there exist powers $`\alpha ^l`$ of $`\alpha `$ with $`l`$ sharing a factor with $`d`$; one then has $`(\alpha ^l)^{d/m}=1`$ where $`m=\mathrm{gcd}(l,d)`$, and the eigenvalue represents $`_{d/m}`$.
Eigenvalues that faithfully represent $`_d`$ will not help in imposing (46), as there is no power $`N<d`$ for which they become unity. Instead, unfaithful eigenvalues are what we are looking for: by definition there exists a power
$`N_J{\displaystyle \frac{d}{\mathrm{gcd}(a_J,d)}}<d`$ (54)
such that $`a_J^N=+1`$. Thus we find the mechanism we need: for each unfaithful eigenvalue $`a_J`$ of $`A`$, there exists an $`N_J`$ such that $`A^{N_J}`$ has one nontrivial eigenvector $`g`$ and the R-symmetry constraint (46) can be imposed. Let us try to count these unfaithful eigenvalues.
As remarked previously, the set $`\{n_J\}`$ of period indices always contains at least the integers $`1n_J<d`$ that share no common factors with $`d`$. These are precisely the powers leading to faithful eigenvalues $`\alpha ^{n_J}`$: hence the (unhelpful) faithful eigenvalues are always present. Since the number of faithful $`n_J`$ is the definition of the Euler totient function $`\varphi (d)`$, and there are $`b_3`$ of the $`n_J`$, we have
$`\varphi (d)b_3<d.`$ (55)
The remaining $`b_3\varphi (d)`$ eigenvalues never faithfully represent $`_d`$, since the exponents share a common factor with $`d`$. Hence, before taking into account the periodicity condition (48), we expect $`b_3\varphi (d)`$ independent solutions to (46).
In fact, one may quickly determine whether there are unfaithful eigenvalues of the monodromy matrix simply by looking at the weights $`k_i`$, as follows. Consider the prime factors $`p_\alpha `$ of the degree $`d`$. If any of the integers $`d/p_\alpha `$ is absent from the set $`\{k_i\}`$, there is an unfaithful eigenvalue $`\alpha ^{p_\alpha }`$ (and in general others of the form $`\alpha ^{mp_\alpha }`$ for some $`m`$) that can be used to impose the R-symmetry constraint (46) with $`N=d/p_\alpha `$. If all integers $`d/p_\alpha `$ are found among the $`k_i`$, all eigenvalues are faithful.
The conclusion of the study of the one-parameter cases was that the dimension of the space of fluxes solving $`DW=W=0`$ was given by the difference of the number of fluxes, and the number of constraints placed on them by expanding the equations in the basis of the extension over the rationals in which the periods lived โ there just the cyclotomic field $`_d`$. For each $`f`$ and $`h`$ flux, this difference was just $`b_3\varphi (d)`$, consistent with the results found here.
In the multiparameter case, however, one must also confront the periodicity condition (48), as we now discuss.
### 4.5 The periodicity condition
For Fermat models, the periodicity $`D`$ of the moduli $`\phi ^A`$ can be calculated from the $`k_i`$ as follows. Consider all sets of 2 or 3 of the $`k_i`$ and find the greatest common divisor. When this is not unity, call it $`D_A`$. There is then a corresponding monomal $`M_A`$, composed only of the coordinates $`x_i`$ whose indices do not contain $`D_A`$, of the form (17) with
$`q_A^i={\displaystyle \frac{d}{D_A}}+{\displaystyle \frac{db_A^i}{\mathrm{lcm}(D_A,k_i)}},`$ (56)
for any integers $`b_A^i>\mathrm{lcm}(D_A,k_i)/D_A`$, with total degree $`d`$ if $`_i(k_i+D_Ak_ib_A^i/\mathrm{lcm}(D_A,k_i))=D_A`$. (In general there can be more than one admissible monomial if several values of $`b_A^i`$ are permitted.) It is then easy to see that $`D_A`$ is the periodicity of this monomial, and hence the modulus $`\phi ^A`$ as well. Consequently the total periodicity is simply<sup>4</sup><sup>4</sup>4For Fermat examples this coincides with the weight of the hyperplane class of $``$ .
$`D\mathrm{lcm}(\{D_A\}),`$ (57)
and the periodicity requirement for an eigenvalue $`\alpha ^{n_J}`$ is that $`N_J=\mathrm{}_JD`$ for some integer $`\mathrm{}_J`$.
Hence for Fermat cases, one may deduce whether there are $`W=0`$ vacua entirely from the $`k_i`$. As discussed in the last subsection, any integer $`d/p_\alpha `$ with $`p_\alpha `$ a prime factor of $`d`$ that is absent from the $`k_i`$ becomes a value of $`N`$ for a set of unfaithful eigenvalues; the periodicity condition is satisfied if this $`d/p_\alpha `$ is a multiple of all common factors of the set of $`k_i`$. For non-Fermat cases, the monomials live in $`๐_{\widehat{k}_1,\widehat{k}_2,\widehat{k}_3,\widehat{k}_4,\widehat{k}_5}^4`$ and cannot be analyzed so simply; we comment on a few non-Fermat examples in the next section.
The requirement of the periodicity condition may remove dimensions from the space of possible fluxes. It corresponds precisely to the fact that the vector space of periods $`๐ฑ`$ can have dimension larger than the cyclotomic field $`_d`$, and hence can place additional constraints on the fluxes. Instead of $`b_3\varphi (d)`$, the more general formula is $`b_3dim๐ฑ`$, corresponding to the number of eigenvalues that unfaithfully represent $`_d`$ as well as complying with the periodicity condition (48) on $`U(\phi ^A)`$. This is our result for the incidence of R-symmetric $`W=0`$ vacua in hypersurfaces in weighted projective space, whose agreement with the counting (36) indicates that the construction via the R-symmetry constraint (46) produces all the relevant vacua.
### 4.6 Enhanced symmetries for special values of the dilaton
The R-symmetries we have constructed all involve a transformation of the complex structure moduli, with the dilaton inert. Imposing additional constraints on the fluxes, one may enhance these symmetries by allowing the dilaton to transform.
These transformations are associated with vacua where the dilaton sits at one of the fixed points of the $`SL(2,)`$ modular group, either the $`_2`$ fixed point $`\varphi =i`$ or the $`_3`$ fixed point $`\varphi =\mathrm{exp}(\pi i/3)`$. Consider a transformation including a $`_2`$ action on the dilaton,
$`\psi \alpha ^M\psi ,\phi ^A=\alpha ^{Q_AM}\phi ^A,\varphi {\displaystyle \frac{1}{\varphi }}.`$ (58)
If the fluxes satisfy the constraint
$`hA^M=\pm f,fA^M=h,`$ (59)
the transformation of the superpotential becomes
$`W(\psi ,\phi ^A,\varphi )\pm {\displaystyle \frac{\alpha ^M}{\varphi }}W(\psi ,\phi ^A,\varphi ).`$ (60)
Because the dilaton transformation acts as Kรคhler transformation on the Kรคhler potential, the superpotential must transform with the factor $`1/\varphi `$ for the total to be a symmetry.
The constraints (59) require that $`A^{2M}`$ has eigenvalue $`1`$. Hence we see that if an R-symmetry associated to a relation $`g=gA^N`$ obeys $`N=4M`$ for integer $`M`$, we can extend the $`_{d/N}`$ R-symmetry to its โfourth rootโ, the larger $`_{d/M}=_{4d/N}`$ transformation (58).
Using the flux constraints (59) and the formula (41) for the dilaton, one may show that $`\varphi =i`$ in the vacuum; hence as with the complex structure moduli space, the symmetries are present at fixed points of the modular group action.
An analogous situation holds for the $`_3`$ transformation on the dilaton
$`\psi \alpha ^M\psi ,\phi ^A=\alpha ^{Q_AM}\phi ^A,\varphi {\displaystyle \frac{1}{1\varphi }}.`$ (61)
In this case the flux constraint is
$`hA^M=\pm (hf),fA^M=\pm h,`$ (62)
and the transformation of the superpotential becomes
$`W(\psi ,\phi ^A,\varphi )\pm {\displaystyle \frac{\alpha ^M}{1\varphi }}W(\psi ,\phi ^A,\varphi ).`$ (63)
In this case the requirement is that $`A^{3M}`$ has eigenvalue $`1`$. Hence for an R-symmetry $`_{d/N}`$ associated to $`g=gA^N`$ with $`N=3M`$, one can promote it to a symmetry $`_{3d/N}`$ (61) that cubes to the original R-symmetry. Similarly, a symmetry $`_{2d/N}`$ corresponding to $`g=gA^{N/2}`$ can be promoted to $`_{6d/N}`$. The flux constraints (62) require the vacuum value of the dilaton to lie at the $`_3`$ point, $`\varphi =\mathrm{exp}(\pi i/3)`$.
The possibility of enlarging the R-symmetry for special values of the fluxes will in fact be realized in all of our explicit Fermat models, as we shall tabulate in section 5.
### 4.7 Counting $`W=0`$ vacua
As described in section 2, the set of all IIB flux vacua satisfying $`DW=0`$ with tadpole no greater than $`L`$ scales like $`\sqrt{L}^{2b_3}`$, where $`2b_3`$ is the number of fluxes. One is interested in the counting of vacua obeying $`W=0`$ as well.
In , an estimate was made for the number of $`W=0`$ vacua,
$`N_{\mathrm{vacua}}{\displaystyle \underset{H=1}{\overset{H_{\mathrm{max}}}{}}}H^{2๐\eta }L^{b_3(\eta +๐)/2},`$ (64)
with in our notation $`๐dim๐ฑ`$, $`\eta =b_3dim๐ฑ/2`$. Here the height $`H(\varphi )`$ takes into account the value of the dilaton appearing in the $`D_{\phi ^A}W=0`$ equations and its effect on the ability to solve the equations by varying the fluxes.
However, although the formula (64) correctly accounts for the dilaton dropping out of the equations $`D_\varphi W=W=0`$, it does not take into account the resulting triviality and dilaton-independence of the $`D_{\phi ^A}W=0`$ equations. Hence for the class of vacua we study, (64) will not give a proper counting of vacua; the dilaton-independence means the height $`H(\varphi )`$ does not contribute, and the triviality of the $`D_{\phi ^A}W=0`$ equations means fewer constraints are imposed.
Instead, the counting of vacua is much simpler. The dimension of the space of solutions for each $`f`$ and $`h`$ is simply the number of eigenvalues of $`A`$ that unfaithfully represent $`_d`$, as well as being compatible with the periodicity $`D`$. This is equivalent to imposing $`dim๐ฑ`$ constraints on the $`b_3`$ available fluxes. Taking into account both $`f`$ and $`h`$ flux, we find the counting
$`N_{\mathrm{vacua}}(W=0)L^{b_3dim๐ฑ},`$ (65)
with $`b_3dim๐ฑ`$ the number of acceptable eigenvalues. Recalling that all flux vacua scale as $`L^{b_3}`$ (13), we find that the suppression is given precisely by the dimension of the vector space of periods,
$`{\displaystyle \frac{N_{\mathrm{vacua}}(W=0)}{N_{\mathrm{vacua}}}}L^{dim๐ฑ}.`$ (66)
One would like to estimate (65) as $`b_3`$ grows. In particular, do $`W=0`$ vacua make up a sizable fraction of all vacua in this limit, or become a set of measure zero?
We recall from (35) that both $`b_3`$ and $`dim๐ฑ`$ lie between $`\varphi (d)`$ and $`d`$. Hence it is useful to understand the behavior of the Euler totient function $`\varphi (d)`$ as $`d`$ grows; the following results may be found, for example, in . This function does not behave smoothly, as one has $`\varphi (d)=d1`$ for $`d`$ prime, while it may be considerably smaller han $`d`$ when the integer is highly composite. One always has $`\varphi (d)\sqrt{d}`$ except for $`d=2,6`$, and for $`d`$ not prime there is an upper bound:
$`\sqrt{d}\varphi (d)d\sqrt{d},d\mathrm{composite}.`$ (67)
Models where the suppression factor (66) is least for fixed $`b_3`$ are hence those where $`b_3d`$ and $`dim๐ฑ\varphi (d)\sqrt{d}`$, giving
$`{\displaystyle \frac{N_{\mathrm{vacua}}(W=0)}{N_{\mathrm{vacua}}}}L^{\sqrt{b_3}}.`$ (68)
Hence, as $`b_3`$ grows large, $`W=0`$ vacua are invariably more and more suppressed; however, for these least suppressed cases the suppression is a strictly smaller power than the total number of vacua, indicating that the number of $`W=0`$ vacua to leading order in $`b_3`$ scales the same as the total number of vacua, going as $`L^{b_3}`$.
Furthermore, for $`d\mathrm{}`$ the totient function tends towards
$`\varphi (d\mathrm{})e^\gamma {\displaystyle \frac{d}{\mathrm{log}\mathrm{log}d}},`$ (69)
where $`\gamma `$ is the Euler-Mascheroni constant with $`e^\gamma 0.561`$. Hence for all $`b_dd`$, $`dim๐ฑ\varphi (d)`$ models the suppression factor (66) remains strictly smaller than the total number of vacua, and again to leading order in $`b_3`$ the number of $`W=0`$ vacua is $`L^{b_3}`$.
Making more precise statements requires being able to determine more precisely how the dimension of the vector space of periods $`dim๐ฑ`$ differs from $`\varphi (d)`$ in the limit of large $`b_3`$. In particular, in making the above estimates we have assumed that taking $`dim๐ฑ\varphi (d)`$ is valid for some models in this limit. We do not analyze this further here, but leave it as an interesting open question.
## 5 Examples
In this section we apply the analysis of the previous sections to some simple hypersurfaces in weighted projective space: the Fermat models with one, two and three parameters enumerated in . The one-parameter cases were already described in and our results agree with that discussion. We see that the case for the two- and three-parameter models are similar: not every space is compatible with $`W=0`$ vacua, but a significant fraction are.
We also describe a non-Fermat model possessing $`W=0`$ vacua. At the end of the section we discuss generalizations of these results, and present an example of a Calabi-Yau defined as the intersection of multiple polynomials.
### 5.1 Fermat models
In one-parameter models the $`U(\phi ^A)`$ are absent and the periodicity condition is therefore absent as well. Hence the only question is whether the R-symmetry constraint can be imposed, or equivalently, whether $`b_3\varphi (d)>0`$, which is the language in which the issue was studied in . Of the four models listed in table 1, only the $`d=6`$ model $`(1,1,1,1,2)`$ satisfies this relation.
Besides the values of $`d`$ and $`\varphi (d)`$ for each model, we have listed the period indices $`\{n_J\}`$, which give the eigenvalues $`\alpha ^{n_J}`$ of the monodromy matrix $`A`$, with $`\alpha ^d=1`$. We have denoted in bold those $`b_3\varphi (d)`$ eigenvalues that share a common factor with $`d`$ and hence do not faithfully represent $`_d`$. The last column displays the R-symmetry for the model that does possess $`W=0`$ vacua, and in parentheses the enlarged symmetry present when combined with a transformation of the dilaton; those with no R-symmetry do not have $`W=0`$ vacua. This analysis is completely in agreement with .
We now turn to the two- and three-parameter hypersurfaces as listed in (we leave out a three-parameter model that has a non-monomial deformation). Tables 2 and 3 contain additional columns listing the monomials corresponding to the $`\phi ^A`$ moduli, as well as the total periodicity $`D`$ of those monomials. Again those eigenvalues of $`A`$ that are not faithful are bold, and when $`W=0`$ vacua are present, the R-symmetry is listed. The $`W=0`$ vacua for the model $`(1,1,1,2,6)`$ were previously found in .
In most of these cases, all the eigenvalues of $`A`$ are compatible with the periodicity condition (48) set by $`D`$, namely that $`N_Jd/\mathrm{gcd}(n_J,d)=\mathrm{}_JD`$ for some integer $`\mathrm{}_J`$. The exception is the $`(1,2,3,3,3)`$ model, which has four unfaithful eigenvalues, $`\alpha ^2,\alpha ^3,\alpha ^9,\alpha ^{10}`$ with $`\alpha ^{12}=1`$. For $`\alpha ^2`$ and $`\alpha ^{10}`$, we can satisfy the condition $`g=gA^6`$, compatble with the periodicity $`D=3`$. For the other two eigenvalues, however, the corresponding condition would be $`g=gA^4`$, which is incompatbible with $`D=3`$; hence only the first two eigenvalues generate acceptable flux vector solutions. This corresponds precisely to the fact that in this example, $`dim๐ฑ>\varphi (d)`$ as discussed in section 4; we have $`dim๐ฑ=6`$ while $`\varphi (d)=4`$. The dimension of available fluxes is then $`b_3dim๐ฑ=2`$, rather than $`b_3\varphi (d)=4`$.
Hence we see for these simple examples that $`W=0`$ vacua, while not present in every model, are not uncommon and arise in an order one fraction of the cases considered.
### 5.2 Non-Fermat models
Most hypersurfaces in weighted projective space are not Fermat, so it is useful to identify non-Fermat models for which this construction applies. Of the sixteen non-Fermat models with two or three parameters listed in we find two cases with $`W=0`$ vacua. The simpler is the two-parameter model with weights $`(1,1,1,2,3)`$ and $`d=8`$. Here $`b_3=6`$, $`\varphi (8)=2`$ and the $`b_3\varphi (8)=2`$ unfaithful eigenvalues are $`\alpha ^2,\alpha ^6`$ with $`\alpha ^8=1`$, allowing us to impose the R-symmetry constraint
$`g=gA^4.`$ (70)
To calculate the periodicity $`D`$, we proceed as follows. A choice for the defining polynomial is
$`P_0(x_i)=x_1^8+x_2^8+x_3^8+x_4^4+x_5^2x_4,`$ (71)
leading to the mirror polynomial
$`\widehat{P}_0(y_i)=y_1^8+y_2^8+y_3^8+y_4^4y_5+y_5^2,`$ (72)
defined in $`๐_{1,1,1,1,4}^4`$ with $`\widehat{d}=8`$. The orbifold group defining the mirror is $`H=_8\times _8`$, and the monomial other than $`y_1y_2y_3y_4y_5`$ invariant under $`H`$ is $`(y_1y_2y_3y_4)^2`$. Hence $`Q=2`$ for this monomial, and the periodicity is $`D=d/Q=4`$. Since $`N=4=D`$ in (70), the periodicity condition is satisfied and we indeed find $`W=0`$ vacua in this example; the R-symmetry is
$`_4:\psi \alpha ^2\psi ,\phi \phi .`$ (73)
The three-parameter model $`(1,1,2,3,5)`$ with $`d=12`$ has $`W=0`$ vacua as well, as one can show in an analogous fashion; there only two of the four unfaithful eigenvalues are compatible with the periodicity condition.
Thus we see that $`W=0`$ vacua are not restricted to the Fermat cases. In this small sample the Fermat cases seem to have a higher incidence of $`W=0`$ vacua, but one should bear in mind that the non-Fermat class always includes those models with prime $`d`$ (except the quintic), which never support $`W=0`$ vacua by this construction. It would be interesting to quantify more precisely the statement of the periodicity condition for non-Fermat examples, and thus understand whether $`W=0`$ vacua are of greater (or lesser) incidence in Fermat cases than in non-prime non-Fermat models.
### 5.3 Complete intersection Calabi-Yau models
Hypersurfaces in weighted projective space represent a subclass of the more general class of spaces defined as the intersections of multiple hypersurfaces in products of projective spaces. We have not discussed this more general class at all, but one naturally wonders whether our methods may be applied in that more general context.
In this section we give an example of an analogous construction of $`W=0`$ vacua in a simple space defined via the vanishing of two cubic hypersurfaces in $`^5`$:
$`P_1=x_1x_2x_33\psi x_4x_5x_6,P_2=x_4x_5x_63\psi x_1x_2x_3,`$ (74)
where the single parameter $`\psi `$ is the only one surviving the $`H`$-projection to the mirror . There is a $`_6`$ monodromy around $`\psi =0`$, represented on a Picard-Fuchs basis of periods as
$`A=\left(\begin{array}{c}43\mathrm{2\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}1}\\ \mathrm{\hspace{0.33em}\hspace{0.33em}1\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}0\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}0\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}0}\\ \mathrm{\hspace{0.33em}\hspace{0.33em}0\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}1\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}0\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}0}\\ 86\mathrm{5\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}2}\end{array}\right).`$ (75)
This transformation has the interesting property that $`A^61`$ (in fact there is no power of $`A`$ that is the identity); this is related to a logarithmic singularity in the periods at $`\psi =0`$, absent for simple hypersurfaces. Nonetheless, the eigenvalues of $`A`$ do represent $`_6`$; in fact they represent $`_6`$ unfaithfully, as they are all cube roots of unity. Hence we can impose the R-symmetry constraint
$`f=fA^3,h=hA^3,`$ (76)
to obtain fluxes that lead to $`W=0`$ vacua.<sup>5</sup><sup>5</sup>5$`W=0`$ vacua in this model were independently found by A. Giryavets. The associated R-symmetry is
$`_2:\psi \psi .`$ (77)
Since all four eigenvalues of $`A^3`$ are $`+1`$, naively one might think that all fluxes produce $`W=0`$ vacua; however, two of the associated eigenvectors of $`A^3`$ actually vanish, a novel occurance for the CICY models that we did not encounter previously. Hence the $`W=0`$ vacua occur at codimension two in the space of fluxes. As this is a one-parameter model, there is no periodicity condition; the logarithmic singularity of the periods at $`\psi =0`$ is removed by a Kรคhler transformation.
For the other one-parameter CICY models presented in , there is always a $`_d`$ monodromy action around $`\psi =0`$, with $`d`$ the sum of the degrees of the polynomials. A similar analysis then holds for them. We have not considered more complicated CICY spaces in any further detail, but we see no reason in principle why this construction of $`W=0`$ vacua does not apply over the entire class.
## 6 Conclusions
We have demonstrated a general construction for obtaining type IIB flux vacua with vanishing superpotential and discrete R-symmetries. The discrete R-symmetry appears at fixed loci of the corresponding modular transformation, in particular at the Landau-Ginzburg point on the $`\psi `$-plane, when the monodromy action on the periods possesses eigenvalues unfaithfully representing the monodromy group. The remaining condition that these eigenvalues are compatible with the periodicity of the additional complex structure parameters can be thought of as a requirement that these moduli are at fixed points of the modular transformation as well. We formulated our construction for hypersurfaces in weighted projective space, but we see no reason why it should not generalize to the broader class of complete intersections in products of projective spaces, and we gave a simple example in section 5.
We have also showed that the counting (as usual at large charge tadpole $`L`$) of the $`W=0`$ vacua in the most favorable models approaches the same power law as the set of all vacua as the number of moduli grows large; the ratio is subleading in $`b_3`$. There are also many spaces with no $`W=0`$ vacua. Understanding this distribution better involves an improved understanding of the asymptotic behavior of the vector space of periods.
The nature of the construction, that it is valid over a range of parameters, means that not all complex structure moduli are fixed by the fluxes; this is in addition to the Kรคhler moduli, which are never stabilized by fluxes in a IIB vacuum. The nonperturbative effects that can stabilize Kรคhler moduli in general also depend on the complex structure moduli, so it is conceivable that all moduli could still be stabilized once all corrections are taken into account.
We make no claim that our construction exhausts all $`W=0`$ vacua. It is possible that vacua exist at isolated points, rather than only over a continuous range as we show here. The rank of the vector space of periods $`๐ฑ`$ can become smaller at subloci on moduli space, as can the ranks of the vector spaces $`๐ฑ^{(n)}`$ associated with derivatives of the periods; if this happens all at the same point isolated vacua with vanishing superpotential could arise.
This study was motivated partially by the findings of that $`W=0`$ vacua in one-parameter models were associated with particular arithmetic structure, with the periods living in field extensions of the rational numbers of small degree, in particular cyclotomic fields for the hypersurfaces in weighted projective space. We find here, however, that in general $`W=0`$ vacua are not associated with an arithmetic structure. The place of the field extension is taken by the โvector space of periodsโ, which as the name implies maintains the vector space properties of the extensions without preserving an obvious product. Special points in moduli space exist at least in some cases where the vector space reduces in dimension, and may be promoted to a field, associated with attractor points and complex multiplication . It seems likely that further study of these connections could shed more light on the mathematical structure of the periods and the associated vacua.
## Acknowledgments
I benefited from discussions and correspondence with M. Dine, A. Giryavets, S. Gukov, S. Kachru, S. Katz, D. Morrison, Z. Sun and W. Taylor. This work was supported by NSF grant PHY-0243680. Any opinions, findings, and conclusions or recommendations expressed in this material are those of the author and do not necessarily reflect the views of the National Science Foundation.
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# Next-to-Leading Order QCD Correction to ๐โบโข๐โปโ๐ฑ/๐+๐ผ_๐ at โ๐ =10.6 GeV
## Abstract
One of the most challenging open problems in heavy quarkonium physics is the double charm production in $`e^+e^{}`$ annihilation at B factories. The measured cross section of $`e^+e^{}J/\psi +\eta _c`$ is much larger than leading order (LO) theoretical predictions. With the nonrelativistic QCD factorization formalism, we calculate the next-to-leading order (NLO) QCD correction to this process. Taking all one-loop self-energy, triangle, box, and pentagon diagrams into account, and factoring the Coulomb-singular term into the $`c\overline{c}`$ bound state wave function, we get an ultraviolet and infrared finite correction to the cross section of $`e^+e^{}J/\psi +\eta _c`$ at $`\sqrt{s}=10.6`$ GeV. We find that the NLO QCD correction can substantially enhance the cross section with a K factor (the ratio of NLO to LO ) of about 1.8-2.1; hence it greatly reduces the large discrepancy between theory and experiment. With $`m_c=1.4\mathrm{GeV}`$ and $`\mu =2m_c`$, the NLO cross section is estimated to be $`18.9`$ fb, which reaches to the lower bound of experiment.
One of the most challenging open problems in heavy quarkonium physics and nonrelativistic QCD (NRQCD) is the double charm production in $`e^+e^{}`$ annihilation at B factories. The inclusive production cross section of $`J/\psi `$ via double $`c\overline{c}`$ in $`e^+e^{}J/\psi c\overline{c}`$ at $`\sqrt{s}=10.6`$GeV measured by Belle Collaboration Abe:2002rb is about a factor of 5 higher than theoretical predictions including both the color-singletcs and color-octetliu04 $`c\overline{c}`$ contributions in the leading order (LO) NRQCD BBL . Even more seriously, the exclusive production cross section of double charmonium in $`e^+e^{}J/\psi \eta _c`$ measured by Belle Abe:2002rb ; Pakhlov
$`\sigma [J/\psi +\eta _c]\times B^{\eta _c}[2]`$ $`=`$ $`\left(25.6\pm 2.8\pm 3.4\right)\mathrm{fb},`$ (1)
and BaBarBaBar:2005
$`\sigma [J/\psi +\eta _c]\times B^{\eta _c}[2]`$ $`=`$ $`\left(17.6\pm 2.8_{2.1}^{+1.5}\right)\mathrm{fb},`$ (2)
could be larger than theoretical predictions by an order of magnitude or at least a factor of 5. Here $`B^{\eta _c}[2]`$ is the branching fraction for the $`\eta _c`$ to decay into at least 2 charged tracks, so Eqs. (1) and (2) give the lower bound for this cross section. Theoretically, treating charmonium as a nonrelativistic $`c\overline{c}`$ bound state, two independent studies by Braaten and Lee Braaten:2002fi and by Liu, He, and Chao Liu:2002wq showed that at LO in the QCD coupling constant $`\alpha _s`$ and the charm quark relative velocity $`v`$ the cross-section of $`e^+e^{}J/\psi \eta _c`$ at $`\sqrt{s}=10.6`$GeV is about $`3.85.5`$fb (depending on the used parameters, e.g., the long-distance matrix element, $`m_c`$ and $`\alpha _s`$). In comparison with Eq. (1) or Eq. (2), such a large discrepancy between theory and data may present a challenge to our current understanding of charmonium production based on NRQCD and perturbative QCD.
Some theoretical studies have been suggested in order to resolve this large discrepancy problem. In particular, Bodwin, Braaten, and Lee proposed Bodwin:2002fk ; Bodwin:2002kk that processes proceeding via two virtual photons may be important, and Belle data for $`J/\psi +\eta _c`$ might essentially include the $`J/\psi +J/\psi `$ events which were produced via two photons. Brodsky, Goldhaber, and Lee suggested that since the dominant mechanism for charmonium production in $`e^+e^{}`$ annihilation is expected to be the color-singlet process $`e^+e^{}c\overline{c}gg`$, the final states observed by Belle might contain $`J/\psi `$ and a $`M3\mathrm{G}\mathrm{e}\mathrm{V}`$ spin-$`J`$ glueball $`๐ข_J`$ ($`J=0,2`$Brodsky:2003hv . Motivated by these proposals, Belle presented an updated analysis Abe:2004ww , and ruled out the $`J/\psi +J/\psi `$ and spin-0 glueball scenarios. Ma and Si studied this process by treating the charm quark as a light quark and using light-cone distribution amplitudes to parameterize nonperturbative effects related to the inner structure of charmonium MaandSi:2004 . Similar approaches were also considered by Bondar and Chernyak bondar . But the enhanced cross section is sensitive to the specific form of quark distributions. Hagiwara, Kou and Qiao obtained a result consistent with Ref.Braaten:2002fi and Ref.Liu:2002wq , and conjectured that higher-order corrections in $`\alpha _s`$ may be huge Hagiwara:2003cw . There are also other suggestions to resolve the double charmonium problem, and a comprehensive review on related topics and recent developments in quarkonium physics can be found in Ref. Brambilla:2004wf .
In order to further clarify this problem, in this paper we present a result for the next to leading order (NLO) QCD correction to the process of $`e^++e^{}J/\psi +\eta _c`$. As is known, the NLO QCD corrections are important for quarkonium production in inelastic $`J/\psi `$ photoproductionkra , in $`J/\psi `$ plus jet and plus prompt photon associated production in two photon collisionskni , and in gluon fragmentation functions for heavy quarkoniumbra .
At LO in $`\alpha _s`$, $`J/\psi +\eta _c`$ can be produced at order $`\alpha ^2\alpha _s^2`$, for which we refer to e.g. Ref Liu:2002wq . There are four Feynman diagrams, two of which are shown in Fig. 1, and the other two can be obtained by reversing the arrows on the quark lines. Momenta for the involved particles are assigned as $`e^{}(k_1)e^+(k_2)J/\psi (2p_1)+\eta _c(2p_2)`$. Using the NRQCD factorization formalism, we can write down the scattering amplitude in the nonrelativistic limit to describe the creation of two color-singlet $`c\overline{c}`$ pairs at short distances, which subsequently hadronize into $`J/\psi +\eta _c`$ at long distances in the $`e^+e^{}`$ annihilation process. (Note that here the color-octet $`c\overline{c}`$ contribution is of higher order in $`v`$ and therefore negligible). Choosing the Feynman gauge, we get the amplitude of Born diagrams
$`i_{Born}`$ $`=`$ $`{\displaystyle \frac{4096\pi e_c\alpha \alpha _sm|R_S(0)|^2}{3s^3}}\times `$ (3)
$`ฯต_{\alpha \beta \nu \rho }p_1^\alpha p_2^\beta \epsilon ^\nu \overline{v}_e(k_2)\gamma ^\rho u_e(k_1),`$
where $`s=(k_1+k_2)^2`$, $`e_c=\frac{2}{3}`$ is the electric charge of the charm quark, $`\rho `$ is the Lorentz indices of the virtual photon, $`\epsilon `$ is the polarization vector of $`J/\psi `$. $`2p_1`$ and $`2p_2`$ are the momenta of $`J/\psi `$ and $`\eta _c`$ respectively. $`R_S(0)`$ is the radial wave function at the origin of the ground state charmonium $`J/\psi `$ and $`\eta _c`$.
At NLO in $`\alpha _s`$, the cross section is
$`\mathrm{d}\sigma `$ $``$ $`|_{Born}+_{NLO}|^2`$ (4)
$`=`$ $`|_{Born}|^2+2\mathrm{R}\mathrm{e}(_{Born}_{NLO}^{})+๐ช(\alpha ^2\alpha _s^4).`$
The self-energy and triangle diagrams all correspond to propagators and vertexes of Born diagrams. There remain twenty-four box and pentagon diagrams. Twelve diagrams of them are shown in Fig. 2. The upper $`c\overline{c}`$ hadronize to $`J/\psi `$, and the lower to $`\eta _c`$. The other twelve diagrams are obtained by reversing the arrows on the quark lines. Specially, the associated diagram with Pentagon N12 exists only by reversing the arrows on the lower quark lines which hadronize to $`\eta _c`$.
The self-energy and triangle diagrams are in general ultraviolet (UV) divergent; while the triangle, box, and pentagon diagrams are in general infrared (IR) divergent. Box N5 and N8 and Pentagon N10, which have a virtual gluon line connected with the $`c\overline{c}`$ in a meson, also contain the Coulomb singularities due to the exchange of longitudinal gluons between $`c`$ and $`\overline{c}`$. In the practical calculation, the IR and UV singularities are regularized with $`D=42ฯต`$ space-time dimension, and the Coulomb singularities are regularized by a small relative velocity $`v`$ between $`c`$ and $`\overline{c}`$ kra , $`v=|\stackrel{}{p_{1c}}\stackrel{}{p_{1\overline{c}}}|/m`$ , defined in the meson rest frame. For the Coulomb-singular part of the virtual cross section, we find
$`\sigma `$ $`=`$ $`|R_S(0)|^4\widehat{\sigma }^{(0)}\left(1+{\displaystyle \frac{2\pi \alpha _sC_F}{v}}+{\displaystyle \frac{\alpha _s\widehat{C}}{\pi }}+๐ช(\alpha _s^2)\right)`$ (5)
$``$ $`|R_S(0)|^4\widehat{\sigma }^{(0)}\left[1+{\displaystyle \frac{\alpha _s}{\pi }}\widehat{C}+๐ช(\alpha _s^2)\right].`$
In the second step, the Coulomb-singularity term has to be factored out and mapped into the wave functions of $`J/\psi `$ and $`\eta _c`$. For the LO expressions of operators $`๐ช^{J/\psi }\left[{}_{}{}^{3}S_{1}^{(1)}\right]`$ and $`๐ช^{\eta _c}\left[{}_{}{}^{1}S_{0}^{(1)}\right]`$ are associated with $`R_S(0)`$, and the NLO are proportional to $`\pi \alpha _sC_F/v`$ BBL . And the two operators give a factor of $`2`$ at $`๐ช(\alpha _s)`$ , resulting in just the Coulomb-singular term in Eq. (5).
The self-energy and triangle diagrams contain UV singularities, which are removed by the renormalization of the QCD coupling constant $`g_s`$, the charm-quark mass $`m`$ and field $`\psi `$, and the gluon field $`A_\mu `$. Similar to the renormalization scheme in Ref.kni (see also kra ), we define
$$g_s^0=Z_gg_s,m^0=Z_mm,\psi ^0=\sqrt{Z_2}\psi ,A_\mu ^0=\sqrt{Z_3}A_\mu ,$$
(6)
where the superscript 0 labels bare quantities and $`Z_i=1+\delta Z_i`$, with $`i=g,m,2,3`$, are renormalization constants. The quantities $`\delta Z_i`$ are of $`๐ช(\alpha _s)`$ and they contain UV singularities and finite pieces which depend on the choice of renormalization scheme. We define $`Z_2`$ and $`Z_m`$ in the on-mass-shell (OS) scheme, and $`Z_3`$ and $`Z_g`$ in the modified minimal-subtraction ($`\overline{\mathrm{MS}}`$) scheme
$`\delta Z_2^{\mathrm{OS}}`$ $`=`$ $`C_F{\displaystyle \frac{\alpha _s}{4\pi }}\left[{\displaystyle \frac{1}{ฯต_{\mathrm{UV}}}}+{\displaystyle \frac{2}{ฯต_{\mathrm{IR}}}}3\gamma _E+3\mathrm{ln}{\displaystyle \frac{4\pi \mu ^2}{m^2}}+4\right],`$
$`\delta Z_m^{\mathrm{OS}}`$ $`=`$ $`3C_F{\displaystyle \frac{\alpha _s}{4\pi }}\left[{\displaystyle \frac{1}{ฯต_{\mathrm{UV}}}}\gamma _E+\mathrm{ln}{\displaystyle \frac{4\pi \mu ^2}{m^2}}+{\displaystyle \frac{4}{3}}\right],`$
$`\delta Z_3^{\overline{\mathrm{MS}}}`$ $`=`$ $`{\displaystyle \frac{\alpha _s}{4\pi }}(\beta _02C_A)\left[{\displaystyle \frac{1}{ฯต_{\mathrm{UV}}}}\gamma _E+\mathrm{ln}(4\pi )\right],`$
$`\delta Z_g^{\overline{\mathrm{MS}}}`$ $`=`$ $`{\displaystyle \frac{\beta _0}{2}}{\displaystyle \frac{\alpha _s}{4\pi }}\left[{\displaystyle \frac{1}{ฯต_{\mathrm{UV}}}}\gamma _E+\mathrm{ln}(4\pi )\right],`$ (7)
where $`\mu `$ is the renormalization scale, $`\gamma _E`$ is the Eulerโs constant and $`\beta _0=(11/3)C_A(4/3)T_Fn_f`$ is the one-loop coefficient of the QCD beta function, and $`n_f`$ is the number of active quark flavors. There are three massless light quarks $`u,d,s`$ and one heavy quark $`c`$, so $`n_f=4`$. Color factors are given by $`T_F=1/2,C_F=4/3,C_A=3`$ in $`SU(3)_c`$. Differing from Ref.kni , we take the $`\overline{\mathrm{MS}}`$ scheme for $`Z_3`$ with no external gluon legs and set $`n_f=4`$. In this scheme, we do not need to calculate the self-energy on external quark legs. It turned out that the difference for the calculated cross section in different schemes is of order of next to next to leading order and can therefore be neglected in the NLO result. In the NLO corrections we should use the two-loop formula for $`\alpha _s(\mu )`$,
$$\frac{\alpha _s(\mu )}{4\pi }=\frac{1}{\beta _0L}\frac{\beta _1\mathrm{ln}L}{\beta _0^3L^2},$$
(8)
where $`L=\mathrm{ln}\left(\mu ^2/\mathrm{\Lambda }_{\mathrm{QCD}}^2\right)`$, and $`\beta _1=(34/3)C_{A}^{}{}_{}{}^{2}4C_FT_Fn_f(20/3)C_AT_Fn_f`$ is the two-loop coefficient of the QCD beta function.
Pentagon diagrams N11 and N12 can be reduced to integrals with a lower number of external legs directly, since there are only two independent momenta. Then they can be calculated the same way as box diagrams. To treat Pentagon N10 in Fig. 2, we need to calculate the five-point function $`E_0[p_1,2p_1,p_2,2p_2,m,0,m,0,m]`$, and the finite term $`E_0^{fin}`$, where
$`E_0`$ $`=`$ $`E_0^{fin}+{\displaystyle \frac{2}{s}}D_0[p_1,p_1p_2,p_1,0,m,0,m]+{\displaystyle \frac{2}{s}}D_0[p_2,p_1+p_2,p_2,0,m,0,m],`$ (9)
$`E_0^{fin}`$ $`=`$ $`{\displaystyle \frac{4}{s}}D_0[p_1+p_2,p_1+2p_2,p_1,0,0,m,m]+{\displaystyle \frac{\mathrm{d}^Dq}{(2\pi )^D}\frac{2/s(s/24qp_1+4qp_28m^2)}{(q^2m^2)(q+p_1)^2((q+2p_1)^2m^2)(qp_2)^2((q2p_2)^2m^2)}}`$
$`=`$ $`{\displaystyle \frac{2\sqrt{4m^2s}\mathrm{tan}^1\frac{\sqrt{s}}{\sqrt{4m^2s}}\sqrt{s}\mathrm{ln}(\frac{s}{m^2})}{i\pi ^2m^2s^{5/2}}}+{\displaystyle \frac{2(4m^2s)^{3/2}\mathrm{tan}^1\frac{\sqrt{s}}{\sqrt{4m^2s}}+\sqrt{s}\left(i\pi (3m^2s)+(s4m^2)\mathrm{ln}(\frac{s}{m^2})\right)}{8im^4\pi ^2(4m^2s)s^{5/2}(16m^2s)^1}},`$
where the IR- and Coulomb-finite term $`E_0^{fin}`$ is calculated with dimension $`D=4`$ and velocity $`v=0`$, and $`\mathrm{ln}(s/m^2)=\mathrm{ln}((s+i0)/m^2)=\mathrm{ln}(s/m^2)i\pi `$. In Eq. (9) the $`D_0[p_1,p_1p_2,p_1,0,m,0,m]`$ term is,
$`D_0`$ $`=`$ $`{\displaystyle \frac{4}{s}}C_0[p_1,p_1,0,m,m]+{\displaystyle \frac{i}{(4\pi )^2}}{\displaystyle \frac{2i\pi 2\mathrm{l}\mathrm{n}4}{m^2s}}.`$ (11)
This term will appear in $`\mathrm{Box}\mathrm{N5},\mathrm{N8}`$. The other IR-divergence terms can be calculated like that. Then all the IR-divergence terms become $`C_0[p_1,p_2,0,m,m]`$ and $`C_0[p_{1c},p_{1\overline{c}},0,m,m]`$. We find that $`\mathrm{Box}\mathrm{N3},\mathrm{N6},\mathrm{N7}`$ and $`\mathrm{PentagonN12}`$ are IR-finite respectively, and sum of $`\mathrm{Box}\mathrm{N1}+\mathrm{N2}+\mathrm{N4}+\mathrm{N9}`$ is IR-finite, and IR-divergence term of $`\mathrm{Pentagon}\mathrm{N11}`$ is canceled by vertex diagrams. IR-divergence and Coulomb-singular terms of $`\mathrm{Box}\mathrm{N5}+\mathrm{N8}`$ and $`\mathrm{Pentagon}\mathrm{N10}`$ are all related to the $`C_0[p_{1c},p_{1\overline{c}},0,m,m]`$ term. With $`v=|\stackrel{}{p_{1c}}\stackrel{}{p_{1\overline{c}}}|/m0`$,
$`C_0={\displaystyle \frac{i}{2m^2(4\pi )^2}}\left({\displaystyle \frac{4\pi \mu ^2}{m^2}}\right)^ฯต\mathrm{\Gamma }(1+ฯต)\left[{\displaystyle \frac{1}{ฯต}}+{\displaystyle \frac{\pi ^2}{v}}2\right].`$ (12)
IR-divergence terms of $`\mathrm{Box}\mathrm{N5}+\mathrm{N8}+\mathrm{Pentagon}\mathrm{N10}`$ are canceled by counter terms, and the Coulomb singularity is mapped into $`R_S(0)`$. UV term is canceled by counter terms. Then the final NLO result for the cross section is UV-, IR-, and Coulomb-finite. Details of the calculation can be found in a forthcoming paper.
We now turn into numerical calculations for the cross section of $`e^++e^{}J/\psi +\eta _c`$. To be consistent with the NLO result the value of the wave function squared at the origin should be extracted from the leptonic width at NLO of $`\alpha _s`$ (see e.g. BBL ): $`|R_S(0)|^2=[(9m_{J/\psi }^2)/(16\alpha ^2(14C_F\alpha _s/\pi ))]\mathrm{\Gamma }(J/\psi e^+e^{})`$. Using the experimental value $`5.40\pm 0.15\pm 0.07`$ KeV PDG , we obtain $`|R_S(0)|^2=0.978\text{GeV}^3`$, which is a factor of 1.21 larger than $`0.810\text{GeV}^3`$ that was used in Refs.Braaten:2002fi ; Liu:2002wq from potential model calculations. Taking $`m_{J/\psi }=m_{\eta _c}=2m`$ (in the nonrelativistic limit), $`m=1.5GeV`$, $`\mathrm{\Lambda }_{\overline{\mathrm{MS}}}^{(4)}=338\mathrm{M}\mathrm{e}\mathrm{V}`$, with Eq.(8) we find $`\alpha _s(\mu )=0.259`$ for $`\mu =2m`$ (these are the same as in Ref.Liu:2002wq , except here a larger $`|R_S(0)|^2`$ is used), and get the cross section in NLO
$$\sigma (e^++e^{}J/\psi +\eta _c)=15.7\mathrm{fb},$$
(13)
which is a factor of $`1.96`$ larger than the LO cross section $`8.0`$ fb. If we set $`\mu =m`$ and $`\mu =\sqrt{s}/2`$, then $`\alpha _s=0.369`$ and $`0.211`$, which result in the cross section $`27.5`$ and $`11.2`$ fb respectively. If we set $`m=1.4\mathrm{GeV}`$ and $`\mu =2m`$, the cross section is $`18.9`$ fb and $`9.2`$ fb at NLO and LO respectively. (Our LO result is also consistent with Ref.Braaten:2002fi if we take their smaller value for $`|R_S(0)|^2`$ and $`\mu =\sqrt{s}/2`$.) In Fig. 3 we show the calculated $`e^++e^{}J/\psi +\eta _c`$ cross sections at LO and NLO as functions of the renormalization scale $`\mu `$ with two mass values $`m_c=1.4`$GeV and $`1.5`$GeV, as compared with the Belle and BaBar data. We see the NLO QCD correction enhances the cross section by about a factor of 2, despite of existing theoretical uncertainties.
The relativistic corrections may further significantly enhance the cross sectionHFC (see also Braaten:2002fi ). The reason for the enhancement is quite obvious that in Fig. 1 the virtuality of the gluon takes its maximum value of $`Q^2=s/4`$ in the nonrelativistic limit, and taking account of the relative momentum between the charm quarks in the charmonium will lower the value of the gluon virtuality.
In conclusion, we find that by taking all one-loop self energy, triangle, box, and pentagon diagrams into account, and factoring the Coulomb singular term associated with the exchange of longitudinal gluons between $`c`$ and $`\overline{c}`$ into the $`c\overline{c}`$ bound state wave function, we get an ultraviolet (UV) and infrared (IR) finite correction to the cross section of $`e^+e^{}J/\psi +\eta _c`$ at $`\sqrt{s}=10.6`$ GeV, and that the NLO QCD correction can substantially enhance the cross section with a K factor (the ratio of NLO to LO ) of about 1.8-2.1; and hence it crucially reduces the large discrepancy between theory and experiment. With $`m=1.4\mathrm{GeV}`$ and $`\mu =2m`$, the NLO cross section is estimated to be $`18.9`$ fb, which reaches to the lower bound of experiment.
###### Acknowledgements.
We would like to thank G.T. Bodwin, B.A. Kniehl and J. Lee for helpful comments and discussions. We also thank K.Y. Liu and C. Meng for valuable discussions. This work was supported in part by the National Natural Science Foundation of China (No 10421503), the Key Grant Project of Chinese Ministry of Education (No 305001), and the Research Found for Doctorial Program of Higher Education of China.
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# Piecewise harmonic subharmonic functions and positive Cauchy transforms
## 1. Introduction
One of the most frequently used constructions in complex analysis and geometry is to consider the maximum of a finite number of pairwise distinct harmonic functions. As is well known, the result is a subharmonic function which is also piecewise harmonic. A quite natural problem is to investigate the converse direction, namely study the class of functions generated by this basic albeit fundamental procedure. Its classical flavor and some important applications โ some of which are listed below โ further motivate a deeper study of this question on which surprisingly little seems to be known. In this paper we answer this question by giving a local characterization of the aforementioned class of functions in generic cases and in the process we establish several remarkable properties for this class. In particular, we show that any subharmonic piecewise harmonic function may essentially be realized as the maximum of finitely many harmonic functions.
### 1.1. Piecewise Harmonic and Piecewise Analytic Functions
Let us first define a fairly general notion.
###### Definition 1.
Let $`X`$ be a real or complex subspace of the space of smooth functions in a domain (open connected set) $`U`$ in $`^2`$ or $``$. We say that a function $`\phi `$ is piecewise in $`X`$ if one can find finitely many pairwise disjoint open sets $`M_i`$, $`1ir`$, in $`U`$ and pairwise distinct functions $`\phi _iX`$, $`1ir`$, such that
* $`\phi =\phi _i`$ in $`M_i`$, $`1ir`$;
* $`U_{i=1}^rM_i`$ is of Lebesgue measure $`0`$.
The set of all functions that are piecewise in $`X`$ is denoted by $`PX`$.
###### Remark 1.
It is not difficult to see that $`PX`$ is actually a (real or complex) vector space. This as well as further properties of $`PX`$ functions and related concepts are discussed in the Appendix.
Note that since $`PX`$ functions are locally integrable they define distributions and their derivatives are therefore defined in the distribution sense (and functions are identified if they define the same distributions). In particular, if $`\phi PX`$ one can form $`\mathrm{\Delta }\phi ๐^{}(U)`$ and also $`_z\phi ,_{\overline{z}}\phi ๐^{}(U)`$ if $`X`$ is complex.
We now specialize Definition 1 to obtain the main objects of our study, namely the spaces of piecewise harmonic and piecewise analytic functions, respectively.
###### Notation 1.
Fix a domain $`U`$, let $`H=H(U)`$ be the real space of (real-valued) harmonic functions in $`U`$ and $`A=A(U)`$ be the complex space of analytic functions in $`U`$. By Definition 1 the following holds:
* Given a piecewise harmonic function $`\phi PH`$ there exists a finite family of pairwise disjoint open sets $`\{M_i\}_{i=1}^r`$ in $`U`$ covering $`U`$ up to a set of Lebesgue measure $`0`$ and a corresponding family of pairwise distinct harmonic functions $`\{H_i(z)\}_{i=1}^r`$ in $`U`$ such that
$$\phi (z)=\underset{i=1}{\overset{r}{}}H_i(z)\chi _i(z)$$
(1.1)
a.e. in $`U`$, where $`\chi _i`$ is the characteristic function of the set $`M_i`$, $`1ir`$;
* Similarly, any piecewise analytic function $`\mathrm{\Phi }PA`$ may be represented as
$$\mathrm{\Phi }(z)=\underset{i=1}{\overset{r}{}}A_i(z)\chi _i(z)$$
(1.2)
a.e. in $`U`$, where $`M_i`$ and $`\chi _i`$, $`1ir`$, are as in (a) and $`\{A_i(z)\}_{i=1}^r`$ is a family of pairwise distinct analytic functions in $`U`$. Given this data and a point $`pU`$ we set
$$H_i(z)=\mathrm{}\left[_p^zA_i(w)๐w\right],zU,\mathrm{\hspace{0.17em}1}ir.$$
(1.3)
These are well-defined harmonic functions in $`U`$ provided that $`U`$ is simply connected, which we tacitly assume throughout unless otherwise stated.
We stress the fact that in the above definitions no regularity ($`C^1`$) conditions are assumed on the negligible set $`U_{i=1}^rM_i`$. Note also that Definition 1 and Notation 1 are merely a convenient way of saying that a $`PH`$ function $`\varphi `$ equals one of finitely many harmonic functions in certain prescribed sets. Therefore $`PH`$ functions need not be continuous nor subharmonic and one can hardly expect any interesting statements in this kind of generality. The same philosophy applies to $`PA`$ functions: as defined above, a function $`\mathrm{\Phi }`$ is $`PA`$ if it is equal to one of finitely many analytic functions in certain open sets. Thus $`PA`$ functions need not be continuous and this will not be case either in our situation.
### 1.2. Canonical Piecewise Decompositions
Note that conditions (i)โ(ii) in Definition 1 remain valid if non-empty Lebesgue negligible sets are subtracted from the sets $`M_i`$, so it is in general impossible to say something about the boundaries of these sets. However, the inclusions $`M_iU\text{supp}(\phi \phi _i)`$, $`1ir`$, always hold, where the supports are defined in the distribution sense (recall from ยง1.1 that $`PX`$ functions are locally integrable and $`L_{loc}^1(U)`$ is viewed as a subspace of $`๐^{}(U)`$). Now both $`X=H(U)`$ and $`X=A(U)`$ are examples of function spaces satisfying the unique continuation property, i.e., $`f0`$ in $`U`$ if $`fX`$ vanishes in some open non-empty subset of $`U`$. In view of the above inclusions, for spaces with this property one can reformulate Definition 1 in a more canonical way as follows.
###### Definition 2.
Let $`X`$ be a real or complex subspace of the space of smooth functions in a domain $`U`$ in $`^2`$ or $``$. Assume that $`X`$ satisfies the unique continuation property and let $`\phi L_{loc}^1(U)`$. Then $`\phi PX`$ ($`\phi `$ is piecewise in $`X`$) if one can find pairwise distinct elements $`\phi _iX`$, $`1ir`$, such that the set $`\mathrm{\Gamma }:=_{1ir}\text{supp}(\phi \phi _i)`$ is of Lebesgue measure $`0`$.
Setting $`M_i=U\text{supp}(\phi \phi _i)`$, $`1ir`$, in Definition 2 we see that $`M_i`$ is the largest open set in which $`\phi \phi _i`$ vanishes (as a distribution or almost everywhere). Further useful properties of the canonical piecewise decomposition of the $`PX`$ function $`\phi `$ given in Definition 2 are gathered in the next lemma. Henceforth by a โcontinuous functionโ we mean a function in $`L_{loc}^1(U)`$ which agrees almost everywhere with a continuous function in $`U`$.
###### Lemma 1.
In the above notation the following holds:
* $`_{1ir}M_i=U\mathrm{\Gamma }`$;
* $`\overline{M}_iM_j=\mathrm{}`$, $`1ijr`$;
* $`M_i=\stackrel{ฬ}{\overline{M}_i}`$ (i.e., $`M_i`$ is the interior of $`\overline{M}_i`$), $`1ir`$;
* $`\mathrm{\Gamma }=_{1i<jr}\overline{M}_i\overline{M}_j`$.
* If $`\phi `$ is continuous then $`\mathrm{\Gamma }g^1(0)`$, where $`g:=_{1i<jr}(\phi _i\phi _j)`$.
###### Proof.
The first statement is obviously true by the (canonical) definition of the sets $`M_i`$, $`1ir`$. To prove (ii) suppose that $`ij`$ and $`p\overline{M}_iM_j`$. Then one can find $`qM_i`$ arbitrarily close to $`p`$ with $`qM_iM_j`$. Since $`q\text{supp}(\phi \phi _i)`$ and $`q\text{supp}(\phi \phi _j)`$ one gets $`q\text{supp}(\phi _i\phi _j)`$ and the unique continuation property implies that $`\phi _i=\phi _j`$, which contradicts the fact that $`\phi _i\phi _j`$.
To show (iii) assume that $`p\stackrel{ฬ}{\overline{M}_i}`$. Then there exists an (open) neighborhood $`N`$ of $`p`$ which is contained in $`\overline{M}_i`$. Since $`\overline{M}_iM_j=\mathrm{}`$ if $`ji`$ (cf. (ii)) it follows that $`NM_i\mathrm{\Gamma }`$. Hence $`\phi =\phi _i`$ in $`N`$ and $`NM_i`$, so that in particular $`pM_i`$.
Clearly, $`_{1ir}\overline{M}_i=U`$. Therefore, if $`p\mathrm{\Gamma }`$ then $`p\overline{M}_i`$ for some $`i`$ and $`p`$ must then be a boundary point of $`M_i`$. Assume that $`p\overline{M}_j`$ whenever $`ji`$. Then there is a neighborhood $`N`$ of $`p`$ such that $`NM_j=\mathrm{}`$ for $`ji`$. Hence $`N\overline{M}_i`$ and it follows from (iii) that $`p\stackrel{ฬ}{M_i}`$. This gives a contradiction (since $`p`$ is a boundary point of $`M_i`$) and shows that $`p\overline{M}_i\overline{M}_j`$ for some $`ji`$, which proves (iv).
Finally, if $`\phi `$ is continuous then $`\phi =\phi _i`$ in $`\overline{M}_i`$ and $`\phi =\phi _j`$ in $`\overline{M}_j`$ hence $`\phi _i=\phi _j`$ in $`\overline{M}_i\overline{M}_j`$ and thus $`g0`$ in $`\overline{M}_i\overline{M}_j`$ for $`ij`$, so that by (iv) $`g0`$ in $`\mathrm{\Gamma }`$. โ
The familiar โmaximum constructionโ that we alluded to at the beginning of this introduction yields natural examples of $`PH`$ and $`PA`$ functions. We recall briefly the interplay between the classes of functions obtained in this case:
###### Example 1.
Let $`\{H_i(z)\}_{i=1}^r`$ be a finite family of pairwise distinct harmonic functions in a domain $`U`$. Then $`\phi (z):=\mathrm{max}_{1ir}H_i(z)`$ is a (subharmonic) $`PH`$ function. Indeed, set $`\mathrm{\Omega }:=\{zUH_k(z)H_l(z),1klr\}`$, let $`M_i`$ be the (open) set consisting of those $`z\mathrm{\Omega }`$ for which $`\phi (z)=H_i(z)`$ and denote by $`\chi _i`$ the characteristic function of $`M_i`$, $`1ir`$. It is clear that $`U\mathrm{\Omega }`$ is Lebesgue negligible, so that $`\{M_i\}_{i=1}^r`$ forms a covering of $`U`$ up to a set of Lebesgue measure $`0`$ and
$$\phi (z)=\underset{i=1}{\overset{r}{}}H_i(z)\chi _i(z)$$
a.e. in $`U`$. Moreover, the subharmonicity of $`\phi `$ implies that $`\nu :=^2\phi /\overline{z}z0`$ in the sense of distributions. In fact $`\nu `$ is a positive measure supported on the (finite) union of level curves $`\{zUH_i(z)H_j(z)=0\}`$, $`1ijr`$. One can show that in this case the support actually determines the measure (Theorem 2 in ยง2).
Now the derivative of $`\phi `$, again in the distribution sense, inherits a similar property only this time with respect to analytic functions. Classical results yield namely
$$\phi (z)/z=\underset{i=1}{\overset{r}{}}A_i(z)\chi _i(z)$$
a.e. in $`U`$, where $`A_i:=H_i/z`$, $`1ir`$, are analytic functions in $`U`$ (cf. Proposition 2 in ยง2). Hence $`\phi (z)/z`$ is a $`PA`$ function. Note that the above relation may be reformulated as saying that $`\phi `$ satisfies a.e. in $`U`$ the differential equation $`P(\phi (z)/z,z)=0`$, where $`P(y,z):=_{i=1}^r(yA_i(z))`$ is a polynomial in $`y`$ with coefficients that are holomorphic in $`U`$.
### 1.3. Main Problem and Results
$`PA`$ functions occur naturally โ and this was our original motivation โ in various contexts, such as the study of the asymptotic behavior of polynomial solutions to ordinary differential equations , the theory of Stokes lines and orthogonal polynomials . In the aforementioned contexts $`PA`$ functions are mostly constructed as limits and thus one has no control on the differentiable structure of the resulting sets $`M_i`$. It is therefore important to describe the local and global structure of $`PA`$ functions both with and without additional regularity assumptions โ such as piecewise $`C^1`$-boundary conditions on the sets $`M_i`$, see ยง2 โ and this is the primary objective of this paper. To state our main problem it is convenient to use:
###### Notation 2.
Given a domain $`U`$ let $`\mathrm{\Sigma }(U)=\{f๐^{}(U)_{\overline{z}}f0\}`$.
Clearly, $`_z\phi \mathrm{\Sigma }(U)`$ if $`\phi `$ is subharmonic in $`U`$, which holds e.g. for the maximum of finitely many harmonic functions. For a (known) converse see the Appendix.
###### Main Problem.
Let $`\mathrm{\Phi }\mathrm{\Sigma }(U)`$ be a $`PA`$ function in a given domain $`U`$. Find conditions that guarantee that $`\mathrm{\Phi }`$ is locally (or globally) of the form $`_z\phi `$, where $`\phi `$ is the maximum of a finite number of harmonic functions in $`U`$.
The necessity of assuming $`_{\overline{z}}\mathrm{\Phi }0`$ in the Main Problem will soon become quite clear and is further illustrated in Example 2, see also Lemma 12 in the Appendix. We give four answers to the above problem which may be summarized (in terms of the mutual implications among them) as follows:
$$\text{Theorem }\text{1}\text{Corollary }\text{5}\text{Corollary }\text{1}\text{Theorem }\text{3}.$$
(1.4)
We formulate here just the first (Theorem 1) and third (Corollary 1) main results of this paper. The fourth one (Theorem 3) is an alternative approach to the Main Problem suggested by our referee, as were several ideas used in this paper.
###### Theorem 1.
Let $`\mathrm{\Phi }\mathrm{\Sigma }(U)`$ be a $`PA`$ function as in (1.2) and assume that $`pU`$ satisfies the following conditions:
* $`p\overline{M}_i`$, $`1ir`$;
* $`A_i(p)A_k(p)(A_j(p)A_k(p))`$ for any triple of distinct indices $`(i,j,k)`$ in $`\{1,\mathrm{},r\}`$.
* $`A_i(p)A_k(p)`$ for any pair of distinct indices $`(i,k)`$ in $`\{1,\mathrm{},r\}`$.
There exists a neighborhood $`\stackrel{~}{N}(p)`$ of $`p`$ such that $`\mathrm{\Phi }=2\phi /z`$ a.e. in $`\stackrel{~}{N}(p)`$, where $`\phi (z)=\mathrm{max}_{1ir}H_i(z)`$ and the $`H_i`$โs are the harmonic functions defined in (1.3).
A word about each of the three conditions imposed in Theorem 1 is in order:
* suggests defining the following index set for any $`pU`$:
$$I(p)=\{j\{1,\mathrm{},r\}p\overline{M}_j\}$$
(1.5)
and $`i(p)=|I(p)|`$. Condition (i) then requires that $`i(p)=r`$, i.e., every set $`M_i`$ is โactiveโ. This will be tacitly assumed throughout;
* is the most important assumption and amounts to the requirement that for all distinct indices $`i,j,k\{1,\mathrm{},r\}`$ the level curves $`H_i=H_k`$ and $`H_j=H_k`$ should meet transversally at $`p`$ (i.e., the critical sets $`\mathrm{\Gamma }_{i,j,k}`$ defined in (3.1) below do not contain $`p`$). For an illustration of the necessity of this assumption see Example 3 and Figure 1 in ยง7;
* means that locally the ($`0`$-)level curves of $`H_iH_j`$, $`ij`$, form a foliation by $`1`$-dimensional smooth curves of a small enough neighborhood of $`p`$. As (ii) above, this assumption will also be used in an essential way.
Let $`K`$ be the convex hull of the points $`A_i(p)`$, $`iI(p)`$, and denote by $`K`$ its boundary, which is clearly an $`i(p)`$-gon. From Theorem 1 and its proof sketched in ยง3 and completed in ยง4 โ ยง5 (see, in particular, Lemma 6 in ยง4.1 and Corollary 5 in ยง4.4) we deduce the following:
###### Corollary 1.
Assume all the hypotheses of Theorem 1 except conditions (i)โ(ii) and set $`S(p)=\{iI(p)A_i(p)\text{ is an extreme point of }K\}`$. If $`A_k(p)K`$ for $`kI(p)S(p)`$ then the conclusion of Theorem 1 holds.
###### Remark 2.
In particular, Corollary 1 holds if $`S(p)=I(p)`$, i.e., all points $`A_i(p)`$, $`iI(p)`$, are extreme in $`K`$.
We emphasize the fact that results similar to those above cannot hold for arbitrary $`PA`$ functions. Indeed, as we already noted, the requirement that $`\mathrm{\Phi }/\overline{z}0`$ is crucial. In particular, it implies that the open sets $`\{M_i\}_{i=1}^r`$ and the analytic functions $`\{A_i(z)\}_{i=1}^r`$ associated with $`\mathrm{\Phi }`$ have to be intimately related to each other. The latter statement is illustrated (and further reinforced) in the next example.
###### Example 2.
Let $`r=2`$, $`A_1(z)1`$ and $`A_2(z)i`$. Then the subharmonic function $`\phi `$ defined in Theorem 1 becomes $`\phi (x,y)=\mathrm{max}(x,y)`$, that is, $`\phi (x,y)=x`$ if $`x+y0`$ and $`\phi (x,y)=y`$ for $`x+y0`$. Hence its derivative $`\frac{2\phi }{z}`$ equals $`1`$ if $`x+y0`$ and $`i`$ for $`xy0`$, respectively. Theorem 1 says (loosely) that among all $`PA`$ functions $`\mathrm{\Phi }`$ of the form $`1\chi _{M_1}+i\chi _{M_2}`$ for varying sets $`M_1`$ and $`M_2`$ (covering some neighborhood of the origin up to a Lebesgue negligible set) $`\frac{2\phi }{z}`$ is the only one that has a positive $`\overline{z}`$-derivative in the sense of distributions. To see why this is the case consider the following simple example: let $`l`$ be a line through the origin with unit normal $`n=n_1+in_2`$, so that $`l`$ consists of two half-planes. Let $`M_1`$ be the one with $`n`$ as interior normal to its boundary and $`M_2`$ the other half-plane. Set $`\mathrm{\Phi }=1\chi _{M_1}+i\chi _{M_2}`$. Then
$$\frac{\mathrm{\Phi }}{\overline{z}}=\frac{1}{2}(1i)(n_1+in_2)ds,$$
where $`ds`$ is Euclidean length measure along the common boundary $`l`$ to $`M_1`$ and $`M_2`$ (see Corollary 4). Clearly, $`\mathrm{\Phi }/\overline{z}0`$ only if $`n_1+in_2=\frac{1}{\sqrt{2}}(1+i)`$, i.e., if the line $`l`$ is given by $`x+y=0`$. In other words one must indeed have $`\mathrm{\Phi }=\frac{2\phi }{z}`$, where $`\phi `$ is the subharmonic function defined in Theorem 1. Note that in this particular example we used the fact that the boundaries of the $`M_i`$โs are $`C^1`$ in order to explicitly calculate the derivative of $`\mathrm{\Phi }`$. Our theorems show that the corresponding result is true in a much more general situation with no assumptions on the boundaries.
The local characterization of subharmonic functions with $`PA`$ derivatives is almost an immediate consequence of Theorem 1 and shows that at generic points such functions are indeed maxima of a finite set of harmonic functions:
###### Corollary 2.
Let $`\mathrm{\Psi }`$ be a subharmonic function such that $`\mathrm{\Psi }/z`$ is a $`PA`$ function with decomposition given by (1.2) and satisfying conditions (i)โ(iii) of Theorem 1. Then there exists a neighborhood $`\stackrel{~}{N}(p)`$ of $`p`$ and harmonic functions $`H_i`$, $`1ir`$, defined in $`\stackrel{~}{N}(p)`$ such that $`\mathrm{\Psi }(z)=\mathrm{max}_{1ir}H_i(z)`$ a.e. in $`\stackrel{~}{N}(p)`$.
Let $`\mathrm{\Phi }\mathrm{\Sigma }(U)`$, so that by Notation 2 and \[8, Theorem 2.1.7\] the measure $`\nu :=\mathrm{\Phi }/\overline{z}`$ is positive. Let further $`pU`$ and $`N(p)`$ be a neighborhood of $`p`$ such that $`\overline{N(p)}U`$. Then the (positive) measure $`\stackrel{~}{\nu }:=\chi _{_{\overline{N(p)}}}\nu `$ extends to $``$ and there exists some analytic function $`A`$ such that $`\mathrm{\Phi }=C_{\stackrel{~}{\nu }}+A`$ (as distributions) in $`N(p)`$, where $`C_{\stackrel{~}{\nu }}`$ is the Cauchy transform of $`\stackrel{~}{\nu }`$ defined by
$$C_{\stackrel{~}{\nu }}:=\frac{1}{\pi z}\stackrel{~}{\nu }.$$
The above decomposition for $`\mathrm{\Phi }`$ is a consequence of formula (4.4.2) in op. cit. asserting that $`\mathrm{\Phi }`$ and $`C_{\stackrel{~}{\nu }}`$ have the same derivative with respect to $`/\overline{z}`$, so that by \[8, Theorem 4.4.1\] they must differ by an analytic function. Hence we also have the following corollary to Theorem 1.
###### Corollary 3.
Let $`\mathrm{\Phi }\mathrm{\Sigma }(U)`$ be a $`PA`$ function with decomposition given by (1.2) and set $`\nu =\mathrm{\Phi }/\overline{z}`$. Assume that $`pU`$ satisfies conditions (i)โ(iii) of Theorem 1 and let $`N(p)`$ and $`\stackrel{~}{\nu }`$ be as above. Then $`\mathrm{\Phi }=C_{\stackrel{~}{\nu }}+A`$ in $`N(p)`$, where $`A`$ is an analytic function and the positive measure $`\stackrel{~}{\nu }`$ is supported in a union of segments of level sets for the functions $`H_iH_j`$, $`1ijr`$. Moreover, $`\nu `$ may be locally described by means of its support in the sense of formula (2.2) (see Theorem 2 (3) in ยง2).
Note that the above results hold in a surprisingly great generality as they assume no ร priori knowledge of the differentiable structure of $`\text{supp}\nu `$. We construct an example showing that the picture is even more complex in non-generic cases and in particular that Corollary 2 is not true if $`p`$ is special enough, see Example 3 in ยง7.
The special case when the $`A_i`$ in Theorem 1 are constant functions was treated in . Our crucial Lemma 3 is mutatis mutandis generalized from that paper. In the simpler situation of loc. cit. some additional global results were obtained. These show essentially that any (locally) $`PH`$ subharmonic function is globally (in $`U`$) a maximum of finitely many harmonic functions. Example 3 in ยง7 again shows that this is not true in general. However, it is not difficult to get complete results in the case when only two functions are involved, see ยง2. It would be interesting to establish when a subharmonic function with a $`PA`$ derivative is globally a maximum of finitely many harmonic functions (cf. Problem 2 in ยง7).
## 2. Derivatives of Sums
Recall the canonical piecewise decomposition of a $`PH`$ function from ยง1.2 (cf. Definition 2 with $`X=H(U)`$). If $`\mathrm{\Psi }(z)`$ is a $`PH`$ subharmonic function of the form (1.1) then the support of the associated Riesz measure $`\mathrm{\Delta }\mathrm{\Psi }`$ equals $`\mathrm{\Gamma }:=U_{i=1}^rM_i`$. Indeed, it is clear that $`\text{supp}(\mathrm{\Delta }\mathrm{\Psi })\mathrm{\Gamma }`$. For the reverse inclusion note that $`\mathrm{\Psi }`$ is harmonic in a neighborhood of any point $`p\mathrm{\Gamma }\text{supp}(\mathrm{\Delta }\mathrm{\Psi })`$. If such a point exists one can find $`ij`$ so that any neighborhood of $`p`$ intersects $`M_i`$ and $`M_j`$, and then $`H_i`$ and $`H_j`$ both agree with $`\mathrm{\Psi }`$ in some neighborhood of $`p`$ hence $`H_i=H_j`$ (by the unique continuation property), which is a contradiction.
In this section we first discuss the case of a $`PA`$ function $`\mathrm{\Phi }`$ with canonical piecewise decomposition as in Definition 2 such that the corresponding set $`\mathrm{\Gamma }=U_{i=1}^rM_i`$ is a locally finite union of piecewise $`C^1`$-curves. We show that if the distribution derivative $`\mathrm{\Phi }/\overline{z}`$ is positive then this measure is determined in a simple way by its support, see Theorem 2 (3) below. Note that in view of Lemma 1 (v) a situation where $`\mathrm{\Gamma }`$ is piecewise smooth occurs if one considers a $`PA`$ function of the form $`\mathrm{\Phi }=_{1ir}(H_i/z)\chi _i`$, where $`\mathrm{\Psi }=_{1ir}H_i\chi _i`$ is a continuous $`PH`$ function (for instance, $`\mathrm{\Psi }`$ could be the maximum of finitely many harmonic functions). In this case we show that the continuity assumption implies that $`\mathrm{\Phi }`$ is actually the distribution derivative of $`\mathrm{\Psi }`$ (without any $`C^1`$-assumptions on $`\mathrm{\Gamma }`$).
We start with the case when only two functions are involved. Assume that $`\mathrm{\Phi }(z)`$ is defined in a domain $`U`$ and that there exists a smooth curve $`\mathrm{\Gamma }U`$ dividing $`U`$ into two open connected components $`U=M_1\mathrm{\Gamma }M_2`$ such that $`\mathrm{\Phi }(z)=A_i(z)`$ in $`M_i`$, $`i=1,2`$, where $`A_i(z)`$ is a function analytic in some neighborhood of $`M_i`$. In particular, $`\mathrm{\Phi }(z)`$ is a $`PA`$ function.
###### Lemma 2.
If $`\nu :=\mathrm{\Phi }(z)/\overline{z}0`$ in the sense of distribution theory (i.e., $`\nu `$ is a positive measure) then at each point $`\stackrel{~}{z}`$ of $`\mathrm{\Gamma }`$ the tangent line $`l(\stackrel{~}{z})`$ to $`\mathrm{\Gamma }`$ is orthogonal to $`\overline{A_1(\stackrel{~}{z})}\overline{A_2(\stackrel{~}{z})}`$ and the measure $`\nu `$ at $`\stackrel{~}{z}`$ equals
$$\frac{|A_1(\stackrel{~}{z})A_2(\stackrel{~}{z})|ds}{2},$$
where $`ds`$ denotes length measure along $`\mathrm{\Gamma }`$.
Lemma 2 is an immediate consequence of the following well-known result, see e.g. \[8, Theorem 3.1.9\].
###### Proposition 1.
Let $`YX`$ be open subsets of $`^k`$ such that $`Y`$ has a $`C^1`$-boundary $`Y`$ in $`X`$ and let $`uC^1(X)`$. If $`\chi __Y`$ denotes the characteristic function of $`Y`$, $`dS`$ the Euclidean surface measure on $`Y`$ and $`n`$ the interior unit normal to $`Y`$ then
$$_j(u\chi __Y)=(_ju)\chi __Y+un_jdS,$$
where $`_j`$ and $`n_j`$ are the partial derivative with respect to the $`j`$-th coordinate and the $`j`$-th component of $`n`$, respectively.
###### Corollary 4.
In the notation of Proposition 1 one has
$$\begin{array}{cc}& \frac{(u\chi __Y)}{\overline{z}}=(\frac{u}{\overline{z}})\chi __Y+\frac{1}{2}u(n_1+in_2)ds,\hfill \\ & \frac{(u\chi __Y)}{z}=(\frac{u}{z})\chi __Y+\frac{1}{2}u(n_1in_2)ds.\hfill \end{array}$$
(2.1)
###### Proof of Lemma 2.
Suppose that the function $`\mathrm{\Phi }(z)=A_1(z)\chi _1(z)+A_2(z)\chi _2(z)`$ satisfies the conditions of Lemma 2, where $`\chi _i`$ is the characteristic function of $`M_i`$, $`i=1,2`$. Corollary 4 implies in particular that $`\nu `$ is supported on the smooth separation curve $`\mathrm{\Gamma }`$ and that with an appropriate choice of co-orientation one has $`\nu =\frac{(A_1A_2)nds}{2}`$, which proves the lemma. โ
Proposition 1 remains true if the boundary of $`Y`$ is assumed to be only piecewise $`C^1`$ or just Lipschitz continuous (cf. op. cit.). We may therefore apply it to functions of the form
$$\underset{1ir}{\mathrm{max}}H_i(z)=\underset{i=1}{\overset{r}{}}H_i(z)\chi _i(z)$$
in $`U`$ and get the description of their derivatives given in the introduction. In this case the normal $`n`$ is defined a.e. with respect to length measure on the boundary and the equality in Corollary 2 is interpreted in this sense.
###### Notation 3.
Given a $`PH`$ function $`\mathrm{\Psi }(z)=_{i=1}^rH_i(z)\chi _i(z)`$ as in (1.1) let $`\mathrm{\Gamma }_\mathrm{\Psi }=U_{i=1}^rM_i`$ and denote by $`\mathrm{\Gamma }_\mathrm{\Psi }^d`$ the set of points where the normal to $`\mathrm{\Gamma }_\mathrm{\Psi }`$ is not defined. In similar fashion, for a $`PA`$ function $`\mathrm{\Phi }(z)=_{i=1}^rA_i(z)\chi _i(z)`$ as in (1.2) we set $`\mathrm{\Gamma }_\mathrm{\Phi }=U_{i=1}^rM_i`$ and let $`\mathrm{\Gamma }_\mathrm{\Phi }^d`$ be the set of points where the normal to $`\mathrm{\Gamma }_\mathrm{\Phi }`$ is not defined.
Essentially the same arguments yield the following generalization of Lemma 2.
###### Theorem 2.
Let
$$\mathrm{\Phi }(z)=\underset{i=1}{\overset{r}{}}A_i(z)\chi _i(z)$$
be a $`PA`$ function in a simply connected domain $`U`$ such that
* $`\mathrm{\Gamma }_\mathrm{\Phi }`$ is a locally finite union of piecewise $`C^1`$-curves;
* $`\mathrm{\Phi }/\overline{z}0`$.
Let $`H_i`$, $`1ir`$, be real-valued harmonic functions as in (1.3). Then for any $`z\mathrm{\Gamma }_\mathrm{\Phi }\mathrm{\Gamma }_\mathrm{\Phi }^d`$ there is a neighborhood $`N(z)`$ such that
1. $`N(z)\mathrm{\Gamma }_\mathrm{\Phi }`$ consists of two components $`N(z)_i`$, $`N(z)_j`$ such that $`\mathrm{\Phi }(z)=A_k(z)`$ in $`N(z)_k`$ for $`k=i,j`$;
2. $`N(z)\mathrm{\Gamma }_\mathrm{\Phi }`$ is contained in a level curve of $`H_iH_j`$ for some $`i,j`$;
3. In $`N(z)`$ one has
$$\mathrm{\Phi }(z)/\overline{z}=\frac{|A_i(z)A_j(z)|ds}{2}.$$
(2.2)
The restriction of $`\mathrm{\Phi }(z)/\overline{z}`$ to $`U\mathrm{\Gamma }_\mathrm{\Phi }^d`$, determined locally by (2.2), extends to a measure $`\mu `$ on $`U`$ which is absolutely continuous with respect to length measure on $`\mathrm{\Gamma }_\mathrm{\Phi }`$. Furthermore $`\mathrm{\Phi }(z)/\overline{z}=\mu `$ in $`U`$. Moreover, if any two level curves $`\mathrm{\Gamma }_{ij}`$, $`\mathrm{\Gamma }_{kl}`$ with $`i<j`$, $`k<l`$, $`(i,j)(k,l)`$ intersect in at most a finite number of points, then the measure $`\mu `$ hence also $`\mathrm{\Phi }(z)/\overline{z}`$ is determined by its support $`\mathrm{\Gamma }_\mathrm{\Phi }`$.
###### Proof.
Assertions (1), (2) and identity (2.2) are direct consequences of Lemma 2. Since by (i) $`\mathrm{\Gamma }_\mathrm{\Phi }`$ is a locally finite union of piecewise $`C^1`$-curves the set $`\mathrm{\Gamma }_\mathrm{\Phi }^d`$ has measure $`0`$ with respect to length measure $`ds`$ on $`\mathrm{\Gamma }_\mathrm{\Phi }`$ and thus the measure $`\mu `$ extending the right-hand side of (2.2) to $`\mathrm{\Gamma }_\mathrm{\Phi }`$ exists. It remains to show that
$$\mathrm{\Phi }/\overline{z}=\mu .$$
(2.3)
Note that $`\mathrm{\Phi }/\overline{z}=\mu +G`$, where $`G`$ is a sum of Dirac measures supported at (singular) points in $`\mathrm{\Gamma }_\mathrm{\Phi }^d`$. Consider now a singular point $`p\mathrm{\Gamma }_\mathrm{\Phi }^d`$, a small neighborhood $`N`$ of $`p`$, and the Cauchy transform $`C_{\stackrel{~}{\mu }}`$ of (the extension to $``$ of) the measure $`\stackrel{~}{\mu }:=\chi _{_{\overline{N}}}\mu `$. Suppose that locally at $`p`$ the measure $`G`$ is given by $`c\delta _p`$ for some $`c0`$. Then the function
$$\mathrm{\Phi }C_{\stackrel{~}{\mu }}\frac{c}{zp}$$
is analytic at $`p`$. On the other hand, $`\mathrm{\Phi }`$ is bounded and by the classical Plemelj-Sokhotski formulas (cf., e.g., \[1, ยง3.6\]) the Cauchy transform $`C_{\stackrel{~}{\mu }}`$ has at most a logarithmic singularity at $`p`$. It follows that $`c=0`$, which proves (2.3). For the last statement in part (3) of the theorem note that the assumption on the level curves made there guarantees that each regular point of $`\mathrm{\Gamma }_\mathrm{\Phi }`$ belongs to a unique $`\mathrm{\Gamma }_{ij}`$, hence in view of (2.2) the measure $`\mathrm{\Phi }/\overline{z}`$ is locally determined by $`\mathrm{\Gamma }_{ij}`$. โ
In the remainder of this paper we will see that results similar to Theorem 2 actually hold without local regularity assumptions as in Theorem 2 (i).
Obviously, a $`PH`$ function $`\mathrm{\Psi }`$ has a $`PA`$ derivative almost everywhere. However, this is not necessarily the same as the distribution derivative of $`\mathrm{\Psi }`$. The next result shows that this is true for continuous $`PH`$ functions.
###### Proposition 2.
If the canonically decomposed PH function
$$\mathrm{\Psi }(z)=\underset{i=1}{\overset{r}{}}H_i(z)\chi _i(z)$$
is continuous in $`U`$ (cf. ยง1.2) then
$$\mathrm{\Psi }(z)/z=\underset{i=1}{\overset{r}{}}A_i(z)\chi _i(z)$$
(2.4)
in the sense of distributions, where $`A_i:=H_i/z`$, $`1ir`$.
###### Proof.
Let $`\mathrm{\Gamma }_\mathrm{\Psi }`$ be as in Notation 3. By Lemma 1 (5) $`\mathrm{\Gamma }_\mathrm{\Psi }`$ is contained in the zero set of the function $`g=_{1i<jr}(H_iH_j)`$. Let $`p\mathrm{\Gamma }_\mathrm{\Psi }\mathrm{\Gamma }_\mathrm{\Psi }^d`$ be a regular point of $`\mathrm{\Gamma }_\mathrm{\Psi }`$ and $`N`$ be a small (open) neighborhood of $`p`$. Let further $`N_\pm `$ be $`N`$ intersected with the two sides of $`\mathrm{\Gamma }_\mathrm{\Psi }`$. It follows that $`N_+M_i`$ and $`N_{}M_j`$ for some $`ij`$ if $`N`$ is small enough, and the restriction of $`\mathrm{\Psi }`$ to $`N`$ is a smooth function plus $`f\chi _i`$, where $`f0`$ in $`\mathrm{\Gamma }_\mathrm{\Psi }`$. Then $`(f\chi _i)/z`$ is a function in $`N`$ and we conclude that $`\mathrm{\Psi }/z=_{i=1}^rA_i\chi _i+G`$, where $`G`$ is a distribution supported at the singular points $`\mathrm{\Gamma }_\mathrm{\Psi }^d\mathrm{\Gamma }_\mathrm{\Psi }`$. Since $`\mathrm{\Gamma }_\mathrm{\Psi }^d`$ is a discrete set, by choosing a continuous solution $`h`$ to $`h/z=_{i=1}^rA_i\chi _i`$ we get a continuous solution $`\mathrm{\Psi }h`$ to $`(\mathrm{\Psi }h)/z=G`$ and it follows that $`G0`$, which proves the proposition. โ
## 3. Local Characterization in Generic Cases: Sketch of Proof
In this section we give an equivalent formulation of Theorem 1 and sketch its proof. Under some mild non-degeneracy assumptions, this provides a local description of functions with positive (distributional) $`\overline{z}`$-derivative which is equal a.e. to one of a finite number of given analytic functions.
Let us first fix notations and assumptions.
###### Notation 4.
Let $`\{M_i\}_{i=1}^r`$, $`r2`$, be a finite family of disjoint open subsets of a simply connected domain $`U`$ covering $`U`$ up to a set of zero Lebesgue measure and denote by $`\chi _i`$ the characteristic function of $`M_i`$. Given a family $`\{A_i(z)\}_{i=1}^r`$ of pairwise distinct analytic functions in $`U`$ define the (measurable) function
$$\mathrm{\Psi }(z)=\underset{i=1}{\overset{r}{}}A_i(z)\chi _i(z).$$
Fix a point $`pU`$. As in (1.3) we let
$$H_i(z)=\mathrm{}\left[_p^zA_i(w)๐w\right],1ir.$$
Note that each $`H_i`$ is a well-defined harmonic function in $`U`$ satisfying $`H_i/z=\frac{1}{2}A_i(z)`$. If $`r3`$ we associate to each triple $`(i,j,k)`$ of distinct indices in $`\{1,\mathrm{},r\}`$ the following โcritical setโ
$$\mathrm{\Gamma }_{i,j,k}=\{zUA_i(z)\text{}A_j(z)\text{}A_k(z)\text{ are collinear}\}.$$
(3.1)
Alternatively, $`\mathrm{\Gamma }_{i,j,k}`$ consists of the set of $`zU`$ such that $`A_i(z)A_k(z)`$ and $`A_j(z)A_k(z)`$ are linearily dependent over the reals. This is the set where the gradients of $`H_iH_k`$ and $`H_jH_k`$ are parallel, or equivalently, the level curves through $`z`$ to these functions are parallel. Clearly, $`\mathrm{\Gamma }_{i,j,k}`$ is either a real analytic curve or else there exists $`c`$ such that $`A_i(z)A_k(z)=c(A_j(z)A_k(z))`$ for all $`zU`$.
In this notation Theorem 1 may then be restated as follows. Suppose โ using the labeling in the theorem โ that $`i(p)=r`$ (cf. (1.5)), assume that $`\mathrm{\Psi }/\overline{z}0`$ as a distribution supported in $`U`$ and let $`pU`$ be such that
* $`p\overline{M}_i`$, $`1ir`$;
* There is no critical set $`\mathrm{\Gamma }_{i,j,k}`$ that contains $`p`$;
* $`A_i(p)A_j(p)`$ for $`1ijr`$, i.e., $`p`$ is a non-singular point of $`H_iH_j`$.
Then there exists a neighborhood $`\stackrel{~}{N}(p)`$ of $`p`$ such that
$$\mathrm{\Psi }=2\phi /z\text{ a.e. in }\stackrel{~}{N}(p),$$
where $`\phi `$ is the subharmonic function defined by
$$\phi (z)=\underset{1ir}{\mathrm{max}}H_i(z).$$
###### Remark 3.
Generically, the sets $`\mathrm{\Gamma }_{i,j,k}`$ are curves and so conditions (ii) and (iii) above hold outside some real analytic set.
Strategy of the proof and two fundamental lemmas. The proof of Theorem 1 is rather technical and the main parts of the argument are contained in Lemma 3 and Lemma 4 below, which to some extent hold independently of condition (ii) in Theorem 1. We will now show that Theorem 1 follows in fact from these two lemmas. First, a convenient reformulation of the conclusion of Theorem 1 is that for $`1ir`$ one has $`\chi _i=1`$ a.e. in the set where $`\phi (z)=H_i(z)`$, and this is what we will actually show. Clearly, it is enough to prove this statement for $`i=1`$.
Assumption I. By considering the function $`\mathrm{\Psi }A_1`$ and using the fact that $`A_1`$ is analytic in $`U`$ (hence $`A_1/\overline{z}=0`$) we may assume without loss of generality that
$$A_1(z)=H_1(z)=0\text{ }\text{for}\text{ }zU,$$
(I)
which we do, except when otherwise stated, throughout the remainder of this section as well as in ยง4 and ยง5.
Define now
$$\begin{array}{cc}& W=W_1(p):=\{zU\phi (z)=0\},\hfill \\ & W_i(p):=\{zU\phi (z)=H_i(z)\},\mathrm{\hspace{0.17em}2}ir.\hfill \end{array}$$
(3.2)
We have to prove that $`\mathrm{\Psi }=0`$ a.e. in $`NW`$, or equivalently $`\mathrm{\Psi }=0`$ a.e. in $`N\stackrel{ฬ}{W}`$ for some small enough neighborhood $`N`$ of $`p`$, where $`\stackrel{ฬ}{W}`$ denotes the interior of $`W`$.
The first lemma asserts that $`\chi _1`$ is increasing along every path along which all functions $`H_i`$, $`2ir`$, are decreasing.
###### Lemma 3.
Let $`pU`$ satisfy all the assumptions of Theorem 1 except condition (ii). If $`\gamma `$ is a piecewise $`C^1`$-path from $`z_1=\gamma (0)`$ to $`z_2=\gamma (1)`$ such that that each of the functions $`[0,1]tH_i(\gamma (t))`$, $`2ir`$, is decreasing then
$$(\chi _1\varphi )(z_1)(\chi _1\varphi )(z_2)$$
(3.3)
for any positive test function $`\varphi `$ with $`\text{supp}\varphi `$ small enough.
The second lemma guarantees that enough many points may be reached by paths of the form given in Lemma 3. To make a precise statement we need the following definition: to each $`zU`$ we associate the set
$$V(z)=\{\zeta U\text{ piecewise }C^1\text{-path from }z\text{ to }\zeta \text{ along which all }H_i\text{ decrease}\}.$$
###### Definition 3.
Given $`pU`$ and two subsets $`M,XU`$ with $`p\overline{M}`$ we say that $`V(z)`$ tends to $`X`$ through $`M`$ as $`zp`$, which we denote by $`lim_{Mzp}V(z)=X,`$ if for each $`\alpha X`$ and any sequence $`\{z_n\}_nM`$ converging to $`p`$ one has $`\alpha V(z_n)`$ for all but finitely many indices $`n`$.
###### Lemma 4.
Let $`pU`$ satisfy all the assumptions of Theorem 1, in particular $`p\mathrm{\Gamma }_{i,j,1}`$ for any $`i,j`$. Then there is a neighborhood $`N`$ of $`p`$ with
$$\underset{Uzp}{lim}V(z)=N\stackrel{ฬ}{W}.$$
###### Remark 4.
Note that there are actually no sets $`\mathrm{\Gamma }_{i,j,k}`$ at all if $`r=2`$ in Lemma 4.
###### Theorem 1: outline of the proof.
As noted in the paragraph preceding Lemma 3, we have to show that there exists a sufficiently small neighborhood $`N`$ of $`p`$ such that $`\mathrm{\Psi }=0`$ a.e. in $`N\stackrel{ฬ}{W}`$. This is trivially true if $`W`$ has no interior points (i.e., if $`\stackrel{ฬ}{W}`$ has zero Lebesgue measure) and so we may assume that $`\stackrel{ฬ}{W}`$ has positive Lebesgue measure.
Let now $`\{\varphi _s\}_s`$ be a sequence of test functions satisfying $`\text{supp}\varphi _s\{0\}`$ as $`s\mathrm{}`$ and $`\varphi _s๐\lambda =1`$, $`s`$, where $`\lambda `$ denotes Lebesgue measure. Note that $`\{\varphi _s\chi _1\}_s`$ converges in $`L_{loc}^1`$ to $`\chi _1`$. In particular, this implies that for all $`ฯต>0`$, $`\delta >0`$ there exist a sufficiently large $`s(ฯต,\delta )`$ such that if $`s`$, $`ss(ฯต,\delta )`$, there is a point $`z_1=z_1(ฯต,\delta ,s)U`$ satisfying
$$|z_1p|<\delta \text{ and }(\varphi _s\chi _1)(z_1)>1ฯต.$$
(3.4)
To see this let $`N_\delta =\{zU|zp|<\delta \}`$ and suppose that $`(\varphi _{s_k}\chi _1)(z)1ฯต`$ for some infinite sequence $`\{s_k\}_k`$ and almost all $`zN_\delta `$. Then
$$_{N_\delta }\left|(\varphi _{s_k}\chi _1)(z)\chi _1(z)\right|๐\lambda (z)>ฯต\lambda (M_1N_\delta )$$
and since by assumption $`\lambda (M_1N_\delta )>0`$ this contradicts the fact that $`\{\varphi _{s_k}\chi _1\}_s`$ converges to $`\chi _1`$ in $`L_{loc}^1`$ as $`k\mathrm{}`$, so that (3.4) must hold.
From (3.3) and (3.4) it follows that $`(\varphi _s\chi _1)(z)>1ฯต`$ for $`zV(z_1)`$, which together with the identity $`\varphi _s1=1`$ yields $`\left(\varphi _s_{i=2}^r\chi _i\right)(z)<ฯต`$ and therefore
$$\begin{array}{cc}\hfill \left|(\varphi _s\mathrm{\Psi })(z)\right|& =\left|\varphi _s(z\zeta )\mathrm{\Psi }(\zeta )๐\lambda (\zeta )\right|\hfill \\ & ฯต\underset{2dr}{\mathrm{max}}\underset{\zeta z\text{supp}\varphi _s}{sup}|A_d(\zeta )|=:ฯตC_s(z),zV(z_1).\hfill \end{array}$$
(3.5)
Now we assume in addition that all the conditions of Theorem 1 and Lemma 4 are true. Fix $`ฯต>0`$. The arguments above show that one can construct a sequence $`\{z_n\}_nU`$ such that
$$|z_np|<\frac{1}{n}\text{ and }(\varphi _{s_n}\chi _1)(z_n)>1ฯต$$
(3.6)
for some strictly increasing sequence of positive integers $`\{s_n\}_n`$. By Lemma 4 there exists a neighborhood $`N`$ of $`p`$ such that each $`zN\stackrel{ฬ}{W}`$ belongs to all but finitely many sets $`V(z_n)`$, $`n`$. Combined with (3.5) this shows that for every $`zN\stackrel{ฬ}{W}`$ there exists $`n_z`$ such that
$$\left|(\varphi _{s_n}\mathrm{\Psi })(z)\right|C_{s_n}(z)ฯต\text{ for }nn_z.$$
(3.7)
Since $`A_d`$, $`2dr`$, are analytic functions and $`\text{supp}\varphi _{s_n}\{0\}`$, $`n\mathrm{}`$, it follows from (3.5) that by shrinking the neighborhood $`N`$ (if necessary) one can find $`C>0`$ such that $`C_{s_n}(z)C`$ for $`n`$ and $`zN\stackrel{ฬ}{W}`$. Together with (3.7) and the fact that $`lim_n\mathrm{}\varphi _{s_n}\mathrm{\Psi }=\mathrm{\Psi }`$ in $`L_{loc}^1`$ this clearly implies that $`\mathrm{\Psi }=0`$ a.e. in $`N\stackrel{ฬ}{W}`$, which proves Theorem 1. โ
## 4. Proof of Lemma 4
To complete the proof of Theorem 1 it remains to show Lemma 3 and Lemma 4. We start with the latter, which we prove in this section.
### 4.1. Preliminaries
Let $`A(z)`$ be an analytic function defined in a neighborhood of some point $`z_0`$ and set $`H(z):=\mathrm{}\left[_{z_0}^zA(w)๐w\right]`$, so that $`H(z)/z=\frac{1}{2}A(z).`$ The directional derivative of $`H`$ with respect to a complex number $`v=\alpha +\beta i`$ is given by
$$D_vH(z)=\alpha H(z)/x+\beta H(z)/y=\mathrm{}\left[vA(z)\right]$$
(4.1)
and the gradient of $`H(x,y)`$ considered as a vector in $``$ is just
$$H(x,y)=2H(z)/\overline{z}=\overline{A(z)}.$$
(4.2)
If $`A(z_0)0`$ then $`z_0`$ is a non-critical point for $`H(z)`$ and locally the $`0`$-level curves of $`H`$ form a foliation by $`1`$-dimensional smooth curves of a small enough neighborhood $`N`$ of $`z_0`$ (\[13, Theorem 5.7\]). In particular, the ($`0`$-)level curve $`C_H`$ of $`H`$ through $`z_0`$ divides $`N`$ into two components
$$N_H^+=\{zNH(z)>0\},N_H=\{zNH(z)<0\}.$$
Correspondingly, the tangent to $`C_H`$ at $`z_0`$ divides the plane into two opposite half-planes
$$\begin{array}{cc}& \tau (z_0)^+=\{v+z_0vH(z_0)0\}=\{v+z_0\mathrm{}\left[vA(z_0)\right]0\},\hfill \\ & \tau (z_0)=\{v+z_0vH(z_0)0\}=\{v+z_0\mathrm{}\left[vA(z_0)\right]0\}.\hfill \end{array}$$
We now return to the functions $`A_i`$, $`1ir`$, suspending for the moment Assumption I in ยง3 stating that $`A_1=0`$. As before, we suppose that $`A_i(p)A_j(p)`$ if $`ij`$. Consider the convex hull $`K`$ of the points $`A_i(p)`$, $`1ir`$. For each $`i`$ define the dual cone (with vertex at $`p`$) to the sector consisting of all rays from $`H_i(p)=\overline{A_i(p)}`$ to points in the complex dual $`\overline{K}`$ by
$$\begin{array}{cc}\hfill \sigma _i(p):& =\underset{kK}{}\{v+pv(\overline{k}H_i(p))0\}\hfill \\ & =\underset{ji}{\overset{r}{}}\{v+pv(H_j(p)H_i(p))0\}\hfill \\ & =\underset{ji}{\overset{r}{}}\{v+p\mathrm{}\left[v(A_j(p)A_i(p))\right]0\}.\hfill \end{array}$$
(4.3)
Clearly, this cone is the infinitesimal analogue of the set $`W_i(p)`$ defined in (3.2). The interior of $`\sigma _i(p)`$ contains the directions in which $`H_i`$ grows faster (up to the first order) than any other $`H_k`$, $`ki`$.
There are several possibilities for the cone $`\sigma _i(p)`$: (a) it may have a top angle strictly between $`0`$ and $`\pi `$, in which case we say that it is a pointed cone (b) it consists just of the point $`p`$ or (c) it is either a line, a half-line or a half-plane.
The next lemma is a direct consequence of basic convex geometry.
###### Lemma 5.
With the above notations and assumptions the following holds:
* If $`A_i(p)`$ lies in the interior of $`K`$ then $`\sigma _i(p)=\{p\}`$;
* If $`K`$ is not a segment then $`A_i(p)`$ is an extreme point of $`K`$ if and only if $`\sigma _i(p)`$ is a pointed cone.
Now consider condition (ii) in Theorem 1, which is also part of the assumptions of Lemma 4. By Lemma 5 (ii) this condition is strictly stronger than the hypothesis in the following lemma.
###### Lemma 6.
Assume that the only points $`A_i(p)`$ contained in the boundary $`K`$ of $`K`$ are extreme points. If $`S(p)=\{i\{1,\mathrm{},r\}A_i(p)\text{ is an extreme point of }K\}`$ then:
* $`\mathrm{max}_{1ir}H_i(z)=\mathrm{max}_{iS(p)}H_i(z)`$ in a neighborhood of $`p`$;
* There is a neighborhood $`N`$ of $`p`$ such that $`_{iS(p)}NW_i=N`$.
###### Proof.
Clearly, (ii) follows from (i). Let now $`jS(p)`$, so that by Lemma 5 and the assumption of Lemma 6 one has $`\sigma _j(p)=\{p\}`$. This means that for each ray from $`p`$ in the unit vector direction $`vS^1`$ there is at least one $`H_i`$, $`iS(p)`$, such that
$$v\{u+pu(H_i(p)H_j(p))0\}.$$
Thus, for each $`vS^1`$ there is a product neighborhood $`I(v)\times J(v,p)S^1\times U`$ of $`\{v\}\times \{p\}`$ such that there exists $`i=i(v)S(p)`$ so that the continuous function $`u(H_i(z)H_j(z))`$ is positive if $`(u,z)I(v)\times J(v,p)`$. By the compactness of $`S^1\times \{p\}S^1\times U`$, a finite number of neighborhoods $`I(v_l)\times J(v_l,p)`$, $`1ls`$, cover $`S^1\times \{p\}`$. Hence the neighborhood $`J(p):=_{1ls}J(v_l,p)`$ of $`p`$ has the property that along each ray from $`p`$ with direction $`vS^1`$ there is some $`iS(p)`$ such that $`H_i(z)>H_j(z)`$ if $`zJ(p)\{p\}`$, which proves (i). โ
For the rest of this section we may again (and do) assume that $`W=W_1`$, $`A_1=H_1=0`$ (see Assumption I in ยง3), and furthermore that $`p=0`$. By condition (ii) in Theorem 1 (which, as we already pointed out, is also assumed in Lemma 4) and Lemma 6 it is then enough to prove Lemma 4 in the case when the index $`1`$ belongs to the set $`S(p)`$ defined above, which we now proceed to do.
### 4.2. Changing Coordinates
To prove Lemma 4 in the above situation we will further simplify the picture by making suitable coordinate changes as follows. Let $`G`$ be a $`C^1`$-homeomorphism from a domain $`U^{}`$ to $`U`$ that takes a neighborhood $`N^{}U^{}`$ of $`p^{}=G^1(p)`$ one-to-one onto $`N`$. Then $`W(p)N`$ is the homeomorphic image under $`G`$ of the set
$$W^{}(p)=\{wN^{}H_i(G(w))0,\mathrm{\hspace{0.17em}2}ir\}$$
(note that we do not need to assume that $`G`$ is analytic since we are not concerned with preserving subharmonicity in the present situation). Furthermore, if $`zU`$ and $`z^{}=G^1(z)`$ then $`V(z)`$ is the homeomorphic image under $`G`$ of the set
$$\begin{array}{cc}& V^{}(z^{})\hfill \\ & =\{\zeta ^{}U^{}\text{ piecewise }C^1\text{-path from }z^{}\text{ to }\zeta ^{}\text{ along which all }H_iG\text{ decrease}\}.\hfill \end{array}$$
Clearly, since $`G`$ is one-to-one it suffices for the proof of Lemma 4 to show that there exists a neighborhood $`N^{}`$ of $`p^{}`$ such that $`V^{}(z^{})`$ tends to $`\stackrel{ฬ}{W^{}}`$ through an appropriate set as $`z^{}p^{}`$ (cf. Definition 3).
As an immediate application of this observation we may prove Lemma 4 in the case when $`K`$ is a line segment. Indeed, suppose that $`A_1(0)=0`$ and $`A_2(0)`$ are the (only) two extreme points of $`K`$. By Lemma 6 the functions $`A_1(z)0`$ and $`A_2(z)`$ are the only active ones at $`p`$ and it suffices to show that $`V(z)`$ tends to $`W`$ through $`W`$ as $`zp`$ in a suitable neighborhood. We may change coordinates as above in order to reduce this case to the situation when $`H_2(x,y)=y`$. Then just consider the harmonic conjugate $`Q`$ of $`H_2`$ and note that $`Nz(Q(z),H_2(z))`$ is a local homeomorphism for a sufficiently small neighborhood $`N`$ of $`p=0`$. It follows that
$$V(z)N=\{wN\mathrm{}w\mathrm{}z\}\text{ and }\stackrel{ฬ}{W}N=\{wN\mathrm{}w<\mathrm{}p=0\},$$
so the conclusion of Lemma 4 is immediate in this case.
### 4.3. The General Case $`r3`$
From the discussion at the beginning of this section it follows that if $`W`$ is as in (3.2) and as before $`\stackrel{ฬ}{W}`$ is its interior we get that the open set
$$\mathrm{\Omega }(p):=\underset{i=2}{\overset{r}{}}N_{H_i}=\stackrel{ฬ}{W}N$$
is bounded by parts of some of the ($`0`$-)level curves through $`p=0`$ of $`H_i`$, $`2ir`$, and part of the boundary of $`N`$. Furthermore, $`\sigma _1(p)`$ is a pointed cone subtending an angle $`\alpha (0,\pi )`$ at its vertex (which is the origin), and it is bounded in a small neighborhood of $`p`$ by tangents to some level curves, say $`H_2=0`$ and $`H_3=0`$, that meet transversally at $`p`$. Since two non-identical real analytic curves can intersect each other only in a discrete set it follows that for a small enough neighborhood $`N`$ of $`p`$ the boundary of $`\mathrm{\Omega }(p)`$ will consist of at most part of two level curves (and part of the boundary of $`N`$).
By the inverse function theorem the map
$$(x,y)R(x,y):=(H_2(x,y),H_3(x,y))$$
is a homeomorphism from a neighborhood (also called $`N`$) of $`p`$ to a neighborhood of $`p`$. This map takes $`WN`$ to an open subset of the third quadrant and $`p`$ is an interior point in the induced topology of the third quadrant. Clearly, the homeomorphism $`G(x,y)=R^1(x,y)`$ satisfies $`H_3(G(x,y))=x`$ and $`H_2(G(x,y))=y`$ so that by ยง4.2 we may assume without loss of generality throughout the rest of this section that
> $`H_2(x,y)=y`$, $`H_3(x,y)=x`$, $`\sigma _1(p)`$ is the third quadrant and $`WN`$ is the corresponding quadrant of a disk.
The assumption on the boundary of the convex hull of the $`A_i(p)`$โs (cf. Lemma 6 and the discussion following it) implies that there are no other level curves through $`p`$ that are parallel to either of the level curves of $`H_2`$ or $`H_3`$ through $`p`$ except the latter curves themselves.
Now by viewing gradients as complex numbers for each $`zN`$ we may write
$$H_k(z)=|H_k(z)|e^{\sqrt{1}\theta _k(z)},\text{ where }\theta _k(z)[0,2\pi ),\mathrm{\hspace{0.17em}2}kr.$$
(4.4)
Our assumptions imply that $`0<\theta _k(p)<\pi /2`$ for $`2kr`$. Let us further shrink $`N`$ โ if necessary โ so that
$$0<\theta _k(z)<\pi /2\text{ for }k\{2,\mathrm{},r\}\{2,3\}\text{ if }zN.$$
(4.5)
###### Claim 1.
For any $`z\stackrel{ฬ}{W}N`$ there exists a neighborhood $`\stackrel{~}{N}_z`$ of $`0`$ such that every point in $`\stackrel{~}{N}_z`$ may be reached by a path from $`z`$ along which each of the functions $`H_k`$, $`2kr`$, increases.
###### Proof.
Let $`z\stackrel{ฬ}{W}N`$. Then clearly both coordinates $`x`$ and $`y`$ are increasing along the straight segment from $`z`$ to $`p=0`$ given by $`\{(1t)zt[0,1]\}`$. Moreover, there is a disk $`N_z`$ centered at $`p`$ such that $`wN_z`$ implies that both $`x`$ and $`y`$ increase along the path $`\gamma _w(t)=(1t)z+tw`$, $`t[0,1]`$, from $`z`$ to $`w`$. (Note that $`N_z`$ is the largest disk contained in $`N\{w\mathrm{}w\mathrm{}z,\mathrm{}w\mathrm{}z\}`$.) Thus both functions $`[0,1]tH_k(\gamma _w(t))`$, $`k\{2,3\}`$, are increasing. Let us show that this is true as well for each of the remaining functions $`[0,1]tH_k(\gamma _w(t))`$, $`k\{2,\mathrm{},r\}\{2,3\}`$. By (4.5) one has $`H_k(z)=(\alpha (z),\beta (z))`$, where $`\alpha (z),\beta (z)>0`$ if $`k\{2,3\}`$ and $`zN`$, so that the derivative
$$\frac{d}{dt}H_k(\gamma _w(t))=\alpha (\gamma _w(t))\mathrm{}(wz)+\beta (\gamma _w(t))\mathrm{}(wz)$$
(4.6)
is positive for $`w=0`$, $`2kr`$, and $`t[0,1]`$. Hence there is a neighborhood $`\stackrel{~}{N}_z`$ of $`0`$ such that the expression in (4.6) is positive for all $`w\stackrel{~}{N}_z`$ and $`t[0,1]`$. This means that each point in $`\stackrel{~}{N}_z`$ may be reached by a path from $`z`$ along which each of the functions $`H_k`$, $`2kr`$, increases. โ
The proof of Lemma 4 is now immediate: if $`\{z_n\}_n`$ is a sequence converging to $`p`$ there is $`n_0`$ such that $`nn_0`$ implies $`z_n\stackrel{~}{N}_z`$ and by Claim 1 there is a path from $`z`$ to $`z_n`$ along which all $`H_k`$, $`2kr`$, increase. Going in the other direction there is a path from $`z_n`$ to $`z`$ along which all $`H_k`$, $`2kr`$, decrease hence $`zV(z_n)`$ for $`nn_0`$. By the above remarks this completes the proof of Lemma 4.
### 4.4. A More Precise Version of Theorem 1
Revisiting the proof of Theorem 1 sketched in ยง3 we see that we can actually formulate a more precise result by using the terminology and arguments given in ยง4.1โยง4.3 above.
###### Corollary 5.
Assume that all hypotheses of Theorem 1 are satisfied except condition (ii). Let $`A_i(p)`$ be an extremal point in $`K`$ and consider the part $`K_i`$ of its boundary (i.e., the union of the two edges of $`K`$) connecting $`A_i(p)`$ to its two neighbouring extremal points. If $`A_k(p)K_i`$, $`ki`$, there exists a neighborhood $`N`$ of $`p`$ such that $`\mathrm{\Psi }=2\phi /z`$ a.e. in $`W_i(p)N`$.
## 5. Proof of Lemma 3
In this section we prove the remaining lemma, namely Lemma 3 that generalizes a corresponding result obtained in in the (simpler) case when the $`A_i`$ are constant functions. Recall Notation 4, the renormalization argument in Assumption I of ยง3 allowing $`A_10`$, and the assumptions of Lemma 3 and Theorem 1 for our given $`PA`$ function
$$\mathrm{\Psi }(z)=\underset{i=1}{\overset{r}{}}A_i(z)\chi _i(z)=0\chi _1(z)+\underset{i=2}{\overset{r}{}}A_i(z)\chi _i(z)$$
(5.1)
and for the path $`\gamma `$. In particular, we assume that condition (iii) in Theorem 1 is fulfilled at all points on $`\gamma `$, that is, $`\gamma `$ does not pass through singular points for the differences $`H_iH_j`$ with $`ij`$. We may reparametrize $`\gamma `$ by arc-length using the parameter interval $`[0,L]`$ and so we may assume that $`|\dot{\gamma }(t)|=1`$, $`t[0,L]`$. Note first that it is enough to prove the following modified form of Lemma 3: for each $`t_1[0,L]`$ there exists $`\eta >0`$ such that for any positive test function $`\varphi `$ with $`\text{supp}\varphi `$ small enough one has
$$\begin{array}{cc}& (\chi _1\varphi )(z_1)(\chi _1\varphi )(z_2),\hfill \\ & \text{where }z_1=\gamma (t_1)\text{ and }z_2=\gamma (t_2)\text{ with }0<t_2t_1<\eta .\hfill \end{array}$$
(5.2)
Indeed, the fact that (5.2) implies Lemma 3 follows easily by a compactness argument: fix $`t_1`$ and let $`s_2`$ be maximal such that (3.3) holds for $`t_2<s_2`$. If $`s_2L`$ then (5.2) gives a contradiction to the maximality of $`s_2`$. For simplicity we make a translation so that $`z_1=0`$. Clearly, we may also assume that $`\gamma `$ is $`C^1`$.
The idea of the proof of inequality (5.2) is to use the asymptotic properties of the logarithm of $`\mathrm{\Psi }`$. For this we need to take the logarithm of the $`A_i`$ and we must therefore make sure that it is possible to choose a suitable branch. To this end we first prove the following assertion.
###### Claim 2.
There exists a neighborhood $`M`$ of $`z_1=0`$ such that
$$A_i(z)\{t\overline{v}t(0,\mathrm{})\},zM,\mathrm{\hspace{0.17em}1}ir,$$
whenever $`v`$ is a unimodular complex number satisfying $`v\sigma (z_1)`$, where (cf. (4.3))
$$\sigma (z_1)=\underset{i=2}{\overset{r}{}}\{u\mathrm{}\left[uA_i(z_1)\right]0\}.$$
###### Proof.
Since $`A_10`$ this is immediate for $`i=1`$. By condition (iii) in Theorem 1 there exists $`c^{}>0`$ such that $`|A_i(z_1)|c^{}`$ for $`i\{2,\mathrm{},r\}`$, so that there is $`c(0,c^{}]`$ and a neighborhood $`M`$ of $`z_1`$ such that $`|A_i(z)|c`$ for $`i\{2,\mathrm{},r\}`$ and $`zM`$. It follows that for all unit vectors $`v\sigma (z_1)`$ we may assume up to shrinking $`M`$ that $`\mathrm{}\left[vA_i(z)\right]c/2`$ for $`zM`$. Thus the angle $`\rho `$ between $`A_i(z)`$ and $`\overline{v}`$ satisfies $`\rho (\pi /3,5\pi /3)`$ since $`\mathrm{cos}\rho =|A_i(z)|^1\mathrm{}\left[vA_i(z)\right]<1/2`$, which proves the claim. โ
We use the result that we have just established in order to simplify the situation. For this we choose $`\eta >0`$ such that $`\gamma (t)M`$, $`t[0,\eta ]`$, where the neighborhood $`M`$ of $`z_1=0`$ is as in Claim 2, and we let $`v=\dot{\gamma }(0)`$. Note that since by the assumption in Lemma 3 all functions $`[0,\eta ]tH_i(\gamma (t))`$, $`2ir`$, are decreasing we have $`v\sigma (z_1)`$ by (4.1). Up to replacing $`\mathrm{\Psi }`$ by the function $`e^{i\theta }\mathrm{\Psi }(e^{i\theta }z)`$, where $`v=e^{i\theta }`$, we may also assume that $`v=1`$. In particular, we deduce that $`\mathrm{}\left[\dot{\gamma }(0)\right]=1>0`$ so that by further shrinking $`M`$ and the corresponding $`\eta >0`$ we get the key property
$$\mathrm{}\left[\dot{\gamma }(t)\right]>0,t[0,\eta ].$$
(5.3)
Let $`\stackrel{~}{\mathrm{\Psi }}_ฯต=\mathrm{log}(\mathrm{\Psi }ฯต)`$, where $`ฯต>0`$ is arbitrary and we have chosen a branch of the logarithm that is defined in the complex plane cut along the positive real axis. The composite distribution $`\stackrel{~}{\mathrm{\Psi }}_ฯต`$ is then defined by the above rotation of the complex plane, since $`v=1\sigma (z_1)`$. We now study its derivative along the path $`\gamma `$.
Give $`\zeta M`$ define as above (cf. (4.3))
$$\sigma (\zeta )=\underset{i=2}{\overset{r}{}}\{u\mathrm{}\left[uA_i(\zeta )\right]0\}.$$
Then for any fixed $`ฯต>0`$ and $`u\sigma (\zeta )`$ with $`\mathrm{}u>0`$ one has
$$\mathrm{}\left[u(A_i(w)ฯต)\right]<0,1ir,$$
(5.4)
for all $`w`$ in a (sufficiently small) neighborhood of $`\zeta `$. In particular, inequality (5.4) holds for all vectors of the form $`u=\dot{\gamma }(t)`$ in view of (5.3) and the fact that all functions $`[0,\eta ]tH_i(\gamma (t))`$, $`2ir`$, are decreasing (and thus $`u\sigma (\zeta )`$ by (4.1)). It follows that if $`\varphi `$ is a positive test function with $`\varphi ๐\lambda =1`$ and $`\text{supp}\varphi `$ is small enough then
$$\mathrm{}\left[u(\varphi \mathrm{\Psi }ฯต)\right]<0$$
(5.5)
and therefore
$$\mathrm{}\left[\frac{\overline{u}}{\varphi \mathrm{\Psi }ฯต}\right]0$$
in a neighborhood of $`\zeta `$. Since $`(\varphi \mathrm{\Psi })/\overline{z}0`$ we get
$$\mathrm{}\left[\overline{u}\frac{}{\overline{z}}\mathrm{log}(\varphi \mathrm{\Psi }ฯต)\right]=\mathrm{}\left[\frac{\overline{u}}{\varphi \mathrm{\Psi }ฯต}\frac{(\varphi \mathrm{\Psi })}{\overline{z}}\right]0.$$
Letting $`\text{supp}\varphi 0`$ with $`\varphi ๐\lambda =1`$ we see that $`\mathrm{log}(\varphi \mathrm{\Psi }ฯต)\stackrel{~}{\mathrm{\Psi }}_ฯต`$ in $`L_{loc}^1`$ (hence as a distribution) and by passing to the limit we get
$$\mathrm{}\left[\overline{u}\frac{\stackrel{~}{\mathrm{\Psi }}_ฯต}{\overline{z}}\right]0.$$
Write now $`\stackrel{~}{\mathrm{\Psi }}_ฯต=\sigma _ฯต+i\tau _ฯต`$, where $`\sigma _ฯต`$ and $`\tau _ฯต`$ are real-valued distributions. Then the latter inequality yields
$$\mathrm{}\left[\overline{u}\frac{\sigma _ฯต}{\overline{z}}\right]\mathrm{}\left[\overline{u}\frac{\tau _ฯต}{\overline{z}}\right],$$
(5.6)
where (5.6) is interpreted as being valid for the restrictions of the corresponding distributions to a neighborhood of $`\zeta `$. Note that up to further shrinking $`M`$ (and the corresponding $`\eta >0`$) by our choice of the branch of the logarithm used in the definition of $`\stackrel{~}{\mathrm{\Psi }}_ฯต`$ we have
$$\tau _ฯต(z)(\frac{\pi }{2},\frac{3\pi }{2}),z2M=\{a+ba,bM\}.$$
(5.7)
Let us show that relations (5.6)โ(5.7) produce the desired result. Recall that for a real-valued function $`\omega (z)`$ one has
$$\frac{\omega (z)}{\overline{z}}=\overline{\frac{\omega (z)}{z}}$$
(5.8)
in the sense of distributions. We consider the derivative of $`\stackrel{~}{\mathrm{\Psi }}_ฯต`$ along the path $`\gamma `$: if $`\varphi `$ is a positive test function then since $`\sigma _ฯต`$ is a real-valued distribution we deduce from (5.8) and (5.6) that the following holds in the interval $`(0,\eta )`$:
$$\begin{array}{cc}\hfill \frac{d}{dt}\left[(\varphi \sigma _ฯต)(\gamma (t))\right]& =2\mathrm{}\left[\dot{\gamma }(t)\frac{\varphi \sigma _ฯต}{z}(\gamma (t))\right]=2\mathrm{}\left[\overline{\dot{\gamma }(t)}\frac{\varphi \sigma _ฯต}{\overline{z}}(\gamma (t))\right]\hfill \\ & =2\mathrm{}\left[\overline{\dot{\gamma }(t)}\frac{\varphi }{\overline{z}}(\gamma (t)w)\sigma _ฯต(w)\right]๐\lambda (w)\hfill \\ & 2\mathrm{}\left[\overline{\dot{\gamma }(t)}\frac{\varphi }{\overline{z}}(\gamma (t)w)\tau _ฯต(w)\right]๐\lambda (w).\hfill \end{array}$$
(5.9)
Now if $`\text{supp}\varphi `$ is small enough, say $`\text{supp}\varphi M`$, then from (5.7) and the fact that $`|\dot{\gamma }(t)|=1`$ for $`t[0,\eta ]`$ (cf. the reparametrization argument at the beginning of this section) we get
$$2|\mathrm{}\left[\overline{\dot{\gamma }(t)}\frac{\varphi }{\overline{z}}(\gamma (t)w)\tau _ฯต(w)\right]d\lambda (w)|2\frac{3\pi }{2}\frac{1}{2}\left(\right|\left|\frac{\varphi }{x}\right||_1+\left|\right|\frac{\varphi }{y}\left||_1\right)=:\kappa (\varphi ),$$
(5.10)
where $`||||_1`$ denotes the $`L^1`$-norm. Note that the (positive) constant $`\kappa (\varphi )`$ defined above does not depend on $`ฯต`$. Combining (5.9) and (5.10) we obtain
$$(\varphi \sigma _ฯต)(z_2)(\varphi \sigma _ฯต)(z_1)\kappa (\varphi )\eta .$$
(5.11)
On the other hand by (5.1) we have
$$\stackrel{~}{\mathrm{\Psi }}_ฯต(z)=\mathrm{log}\left[ฯต\chi _1(z)+\underset{i=2}{\overset{r}{}}(A_i(z)ฯต)\chi _i(z)\right]$$
hence
$$\sigma _ฯต(z)=(\mathrm{log}ฯต)\chi _1(z)+f_ฯต(z),\text{ where }f_ฯต(z)=\underset{i=2}{\overset{r}{}}\mathrm{log}|A_i(z)ฯต|\chi _i(z),$$
and therefore $`(\varphi \sigma _ฯต)(z)=(\mathrm{log}ฯต)(\varphi \chi _1)(z)+(\varphi f_ฯต)(z)`$. By condition (iii) in Theorem 1 there exists $`c>0`$ such that $`|A_i(z)|c`$ for $`i\{2,\mathrm{},r\}`$ and $`zM`$ (cf. the proof of Claim 2). We deduce that there exists $`c^{}>0`$ (independent of $`ฯต`$ and $`\varphi `$) such that $`|(\varphi f_ฯต)(z)|c^{}\varphi _{\mathrm{}}`$ for $`zM`$, where $`||||_{\mathrm{}}`$ denotes the $`L^{\mathrm{}}`$-norm. It follows that
$$(\varphi \sigma _ฯต)(z)=(\mathrm{log}ฯต)(\varphi \chi _1)(z)+O(1).$$
(5.12)
Substituting (5.12) in (5.11) and letting $`ฯต0`$ we conclude that (5.2) holds, which by the preliminary remarks at the beginning of this section completely settles Lemma 3.
## 6. An Alternative Approach Under Extra Conditions
In the previous sections we formulated and proved three results answering the Main Problem stated in ยง1 under fairly mild assumptions, namely Theorem 1 and its consequences Corollary 5 and Corollary 1 (cf. (1.4)). We will now prove Theorem 3 below that provides a fourth answer to the Main Problem under some extra (yet still mild) conditions. Although this result may be obtained directly from Corollary 1, the point in what follows is to present a different approach<sup>1</sup><sup>1</sup>1This approach and the subsequent proofs were suggested by the referee whom we would like to thank for generously sharing his ideas with us. from the one used in ยง3โยง5 that does not rely on Lemma 3 and Lemma 4.
###### Notation 5.
Let $`\mathrm{\Phi }PA`$ be as in (1.2), which we assume to be the canonical piecewise decomposition of $`\mathrm{\Phi }`$ in the sense of Definition 2. We may write
$$U\underset{i=1}{\overset{r}{}}M_i=Z,$$
where $`M_i`$, $`1ir`$, are pairwise disjoint open sets and $`Z`$ is Lebesgue negligible. Note that each $`M_i`$ is also Lebesgue negligible since $`M_iZ`$, $`1ir`$. As before we let $`\chi _i`$ be the characteristic function of $`M_i`$. Recall from (1.5) the set $`I(p)`$ and its cardinality $`i(p)`$ defined for any $`pU`$. To simplify some discussions, assume that $`U`$ is simply connected and choose $`f_iA(U)`$ such that $`f_i^{}(z)=A_i(z)`$, $`1ir`$, where the $`A_i`$ are the given (analytic) functions appearing in the decomposition (1.2) of $`\mathrm{\Phi }`$. Hence
$$\mathrm{\Phi }(z)=\underset{i=1}{\overset{r}{}}f_i^{}(z)\chi _i(z).$$
For arbitrarily fixed $`pU`$ we let
$$\varphi (z)(=\varphi _p(z))=\underset{jI(p)}{\mathrm{max}}\mathrm{}(f_j(z)f_j(p))=\underset{jI(p)}{\mathrm{max}}H_j(z),$$
where the $`H_i`$ are the harmonic functions defined in (1.3) (cf. ยง1 in the case when $`i(p)=r`$). Clearly, $`\varphi `$ is a continuous subharmonic function in $`U`$ which vanishes at $`p`$. Finally, if $`kI(p)`$ and $`i(p)>1`$ set
$$V_k(p)=\left\{\underset{jI(p)\{k\}}{}\theta _j\left(f_k^{}(p)f_j^{}(p)\right)\right|\theta _j0,jI(p)\{k\},\underset{jI(p)\{k\}}{}\theta _j>0\}.$$
Recall the definition of $`\mathrm{\Sigma }(U)`$ from Notation 2.
###### Theorem 3.
In the above notations assume that $`\mathrm{\Phi }\mathrm{\Sigma }(U)`$ and that the following conditions hold:
* The one-dimensional Hausdorff measure of $`M_jM_kM_l`$ is $`0`$ whenever $`j<k<l`$;
* If $`i(p)>1`$ and $`kI(p)`$ then $`0V_k(p)`$.
Then a.e. in a neighborhood of every $`pU`$ one has $`\mathrm{\Phi }=2_z\varphi (=2_z\varphi _p)`$.
###### Remark 5.
Recall the assumptions involving the (extremal points of the) convex hull $`K`$ of the points $`A_i(p)=f_i^{}(p)`$, $`1ii(p)`$, that we used in Corollary 1 and Corollary 5. Although still mild (since it is generically true), requirement (ii) in Theorem 3 is actually stronger than the aforementioned assumptions.
The remainder of this section is devoted to the proof of Theorem 3, which uses induction on $`i(p)`$.
Consider first the case $`i(p)=1`$. By relabeling the indices we may assume that $`I(p)=\{1\}`$, that is, $`p\overline{M}_j`$ for $`j>1`$. Hence $`p`$ is either an interior point of $`M_1`$ or $`pZ`$ and every neighborhood of $`p`$ intersects $`M_1`$. If the former occurs then $`\mathrm{\Phi }(z)=2_z\mathrm{}f_1(z)`$ in an open neighborhood of $`p`$ and thus $`\mathrm{\Phi }=2\varphi /z`$ in that neighborhood. If $`pZ`$ then there is a small open neighborhood $`\mathrm{\Omega }`$ of $`p`$ contained in $`M_1Z`$ and we conclude that $`\mathrm{\Phi }(z)=2_z\mathrm{}f_1(z)`$ a.e. in $`\mathrm{\Omega }`$, hence equality holds in $`\mathrm{\Omega }`$ in the distribution sense. This settles the case when $`i(p)=1`$.
Assume next that $`i(p)=2`$ and (without loss of generality) $`I(p)=\{1,2\}`$. Since the $`M_i`$ are pairwise disjoint it follows that $`pZ`$ and $`p\overline{M}_k`$ for $`k>2`$. Therefore, there is an open neighborhood $`\mathrm{\Omega }`$ of $`p`$ such that
$$\mathrm{\Phi }(z)=f_1^{}(z)\chi _1(z)+f_2^{}(z)\chi _2(z),z\mathrm{\Omega }.$$
Let $`\chi =\chi _2|_\mathrm{\Omega }`$, $`f=f_2f_1|_\mathrm{\Omega }`$, and define
$$\mathrm{\Psi }(z)=f^{}(z)\chi (z)=\mathrm{\Phi }(z)f_1^{}(z).$$
Note that $`_{\overline{z}}\mathrm{\Psi }(z)0`$ in $`\mathrm{\Omega }`$. Condition (ii) in Theorem 3 implies that $`f^{}(p)0`$ and we may assume (after shrinking $`\mathrm{\Omega }`$, if necessary) that $`f`$ is a diffeomorphism from $`\mathrm{\Omega }`$ onto some open disk $`D`$. We may then write $`\chi (z)=\eta (f(z))`$, where $`\eta =\eta (w)=\eta (u+iv)`$ is the characteristic function of some open subset $`\omega `$ of $`D`$, and we get
$$0_{\overline{z}}\mathrm{\Phi }(z)=_{\overline{z}}f^{}(z)\eta (f(z))=|f^{}(z)|^2(_{\overline{w}}\eta )(f(z)),$$
so that $`_{\overline{w}}\eta 0`$ in $`D`$. Since $`\eta `$ is real-valued this means that $`\eta `$ is an increasing function of $`u`$. Hence the open set $`\omega `$ is defined by an inequality of the form $`\mathrm{}w>a`$, and then $`M_2\mathrm{\Omega }`$ is defined by the inequality $`\mathrm{}(f_2(z)f_1(z))>a`$. Moreover, since $`p`$ is in the closure of the set where $`\chi =1`$ we must have $`a=\mathrm{}(f_2(p)f_1(p))`$. Clearly, we may assume that $`f_1(p)=f_2(p)=0`$. Then $`\mathrm{\Phi }(z)=f_1^{}(z)`$ when $`z\mathrm{\Omega }`$ and $`\mathrm{}f_1(z)>\mathrm{}f_2(z)`$ while $`\mathrm{\Phi }(z)=f_2^{}(z)`$ when $`z\mathrm{\Omega }`$ and $`\mathrm{}f_1(z)<\mathrm{}f_2(z)`$. This shows that $`\mathrm{\Phi }=2\varphi /z`$ in a neighborhood of $`p`$, which completes the proof in the case when $`i(p)=2`$.
The above observations also give us a result that will be used later on:
###### Lemma 7.
Assume that $`I(p)=\{j,k\}`$, where $`j<k`$, and that $`\gamma (t)`$ is a $`C^1`$-curve escaping from $`M_j`$ into $`M_k`$ when $`t=\tau `$ in the sense that $`\gamma (t)M_j`$ for $`t<\tau `$ and there is a sequence $`\{\tau _\nu \}_1^{\mathrm{}}`$ with $`\tau _\nu >\tau `$ and $`\tau _\nu \tau `$ as $`\nu \mathrm{}`$ such that $`\gamma (\tau _\nu )M_k`$. Then $`_t\mathrm{}\left(f_j(\gamma (t))f_k(\gamma (t))\right)|_{t=\tau }0`$.
Let us now pass to the case when $`i(p)3`$. Then $`pZ`$ and there is an open neighborhood of $`p`$ that does not intersect $`ri(p)`$ of the $`\overline{M}_j`$. By deleting these sets from $`U`$ we may assume that $`i(p)=r3`$ (cf. the comments after (1.5) in ยง1). We then know that $`p_{i=1}^rM_j`$. It is no restriction to further assume that the $`f_j`$ are normalized so that $`f_j(p)=0`$ for every $`j`$. Then $`\varphi (z)(=\varphi _p(z))=\mathrm{max}_j\mathrm{}f_j(z)`$ and we have to prove that
$$\mathrm{}f_k=\varphi \text{ in }M_kN,$$
(6.1)
where $`NU`$ is a sufficiently small open neighborhood of $`p`$. Let
$$N_k=\{zN\mathrm{}f_k(z)>\mathrm{}f_j(z)\text{ when }jk\}.$$
Suppose now that we can show the following:
$$N_k\overline{M}_k\text{ for every }k\text{ if }N\text{ is sufficiently small}.$$
(6.2)
Since the $`\mathrm{}f_j`$ must be pairwise distinct harmonic functions in $`U`$ (as a consequence of condition (ii) in Theorem 3), the set where $`\mathrm{}f_j=\mathrm{}f_k`$ for some $`j,k`$ with $`jk`$ is of Lebesgue measure $`0`$. It follows that $`N`$ is the disjoint union of the sets $`N_k`$ together with a set of measure $`0`$. Since the $`M_j`$ are pairwise disjoint and $`M_j`$ is of Lebesgue measure $`0`$ for every $`j`$ (since $`M_jZ`$, cf. Notation 5) we deduce that $`(M_kN)N_k`$ is Lebesgue negligible. From this we conclude that $`\mathrm{}f_k=\varphi `$ in $`M_kN`$ hence $`\mathrm{\Phi }=2_z\varphi `$ in $`N`$, which proves Theorem 3.
Thus the main issue is to show that (6.2) holds. When doing this we may assume that $`k=r`$ and consider the harmonic functions $`h_j=\mathrm{}(f_rf_j)`$, $`1jr1`$. We know that $`h_j(p)=0`$. Let $`qN_r`$, i.e., $`qN`$ and $`h_j(q)>0`$ for $`j<r`$. We want to show that $`q\overline{M}_r`$. For this we define
$$\mathrm{\Lambda }=\underset{j<k<l}{}\left(M_jM_kM_l\right).$$
By assumption (i) in Theorem 3, $`\mathrm{\Lambda }`$ has vanishing one-dimensional Hausdorff measure. We need the following lemma.
###### Lemma 8.
There is an open set $`NU`$ containing $`p`$ such that the following holds: if $`wN`$ and $`h_k(w):=\mathrm{}(f_r(w)f_k(w))>0`$ when $`k<r`$, then there exist an open neighborhood $`=_wU`$ of $`p`$ and for every $`z`$ a real analytic mapping $`\gamma =\gamma (s,t)`$ from a neighborhood of $`[0,1]\times [0,1]`$ into $`U`$ such that
* The restriction of $`\gamma `$ to any set where $`t<t_0<1`$ is a diffeomorphism onto its image;
* $`\gamma (1/2,0)=z`$ and $`\gamma (s,1)=w`$ for all $`s`$;
* $`_th_k(\gamma (s,t))>0`$ for all $`(s,t)`$ when $`k<r`$.
Assertion (6.2) โ and thus, as explained above, Theorem 3 as well โ is now a consequence of Lemma 8. Indeed, let $`N`$ be a small neighborhood of $`p`$ satisfying the assumptions of Lemma 8 and $`wN`$ be such that $`h_k(w)>0`$ for $`k<r`$. We need to prove that $`p\overline{M}_r`$. For this let $`=_w`$ be as in the conclusion of Lemma 8. Since $`p\overline{M}_r`$ we know that $``$ contains a point $`zM_r`$. Let $`\gamma `$ be the mapping corresponding to $`z`$ and $`w`$. By shrinking the domain in which the variable $`s`$ ranges we may assume that $`\gamma (s,0)M_r`$ when $`s[0,1]`$. Set
$$๐_\nu =\{(s,t)0s1,\mathrm{\hspace{0.17em}0}t1\nu ^1\}$$
for each integer $`\nu 2`$. Since the one-dimensional Hausdorff measure of $`\mathrm{\Lambda }`$ vanishes this is also true for the one-dimensional Hausdorff measure of
$$K_\nu :=\{(s,t)๐_\nu \gamma (s,t)\mathrm{\Lambda }\}.$$
It follows that
$$J_\nu :=\{s[0,1](s,t)K_\nu \text{ for some }t\}$$
is a closed set of Lebesgue measure $`0`$. In fact, $`J_\nu `$ is the projection of a set with vanishing one-dimensional Hausdorff measure, see, e.g., \[10, Theorem 7.5\]. Therefore, the set $`J_\nu `$ is of the first category, which implies that $`_\nu J_\nu `$ is also of the first category. This gives us an $`s[0,1]`$ such that $`\gamma (s,t)\mathrm{\Lambda }`$ when $`0t<1`$. From condition (c) in Lemma 8 and Lemma 7 it follows that the curve $`t\gamma (s,t)`$, which starts at $`\gamma (s,0)M_r`$, can not leave $`\overline{M}_r`$ until $`t=1`$. Hence $`w\overline{M}_r`$, which proves (6.2) and we are done.
It remains to prove Lemma 8. In doing so we will use the fact that the functions $`h_j=\mathrm{}(f_rf_j)`$, $`1jr1`$, introduced above are real-valued and real analytic, but we will make no use of their harmonicity. Condition (ii) in Theorem 3 implies that the set of all linear combinations $`_{j=1}^{r1}\theta _jdh_j(p)`$, where $`\theta _j0`$ for all $`j`$ and $`dh`$ denotes differential, is contained in a convex cone $`\mathrm{\Gamma }`$ with positive opening angle less than $`\pi `$. We make an affine change of coordinates, only keeping the affine space structure of $``$. This change of coordinates will allow us to replace $`\mathrm{\Gamma }`$ with any other cone with positive opening angle, and without loss of generality we may further assume that $`p`$ is the origin. Then we are in the situation where a set of $`m=r1`$ real analytic and real-valued functions $`h_1,\mathrm{},h_m`$ are defined in a neighborhood $`V`$ of the origin in $`^2`$ and satisfy the conditions
* $`h_j(0)=0`$ and $`dh_j(0)0`$ when $`1jm`$;
* The closed convex cone generated by the gradients $`h_j(0)`$, $`1jm`$, is contained in the cone $`\mathrm{\Gamma }:=\{(x,y)^2|x|y\}.`$
To complete the proof of Lemma 8 we only have to establish the following result.
###### Lemma 9.
Assume conditions (I)โ(II) above. Then there is an open set $`0NV`$ such that the following holds: if
$$w\mathrm{\Omega }_N:=\{z=(x,y)Nh_j(z)>0,\mathrm{\hspace{0.17em}1}jm\}$$
one can find an open neighborhood $`=_w`$ of the origin and for each $`z`$ a $`C^1`$-mapping $`\gamma (s,t)`$ from a neighborhood of $`[0,1]\times [0,1]`$ into $`V`$ such that
* The restriction of $`\gamma `$ to any set where $`t<t_0<1`$ is a diffeomorphism onto its image;
* $`\gamma (1/2,0)=z`$ and $`\gamma (s,1)=w`$ for all $`s`$;
* $`_th_k(\gamma (s,t))>0`$ for all $`(s,t)`$ when $`km`$.
###### Proof.
Define
$$\mathrm{\Omega }_N^\pm =\mathrm{\Omega }_N\{(x,y)^2\pm x0\}$$
whenever $`NV`$. It suffices to prove that there exist an open set $`0N=N_+V`$ such that the conclusion of the lemma holds when $`w\mathrm{\Omega }_N^+`$. Indeed, by replacing $`h_k(x,y)`$ with $`h_k(x,y)`$ we would obtain $`N=N_{}`$ for which the conclusion of the lemma would then be true when $`w\mathrm{\Omega }_N^{}`$ and thus the assertions in the lemma would follow for the open set $`N=N_+N_{}`$.
It is no restriction to assume that $`dh_j(0)`$ is proportional to $`dx+dy`$ for some $`j`$. By shrinking $`V`$ if necessary and applying the implicit function theorem we may also assume that every $`h_j`$ is of the form
$$h_j(x,y)=\beta _j(x,y)(yg_j(x)),$$
where $`\beta _j,g_j`$ are real analytic functions and $`\beta _j>0`$. Then by using the real analyticity of the functions $`g_j`$ we may further assume โ after shrinking $`V`$ and relabeling the indices, if necessary โ that $`V=(b,b)\times (b,b)`$ for some positive real number $`b`$ and that $`g_1(x)g_2(x)\mathrm{}g_m(x)`$ when $`0<x<b`$. With these normalizations it follows that
$$1g_1^{}(0)g_2^{}(0)\mathrm{}g_m^{}(0)=1$$
and finally, after making a non-linear change of the $`x`$-coordinate, we may additionally assume that $`g_m(x)=x`$.
Below we let $`a<b`$ and $`\delta `$ be small positive numbers and we make generic use of the letter $`C`$ to denote constants that are independent of $`a`$ and $`\delta `$ when these stay small. Define
$$\begin{array}{cc}\hfill N(a)& =\{z|z|<a\},\hfill \\ \hfill \mathrm{\Omega }^+(a)& =\{z=(x,y)N(a)x0\text{ and }h_k(z)>0\text{ for all }k\},\hfill \end{array}$$
so that $`\mathrm{\Omega }^+(a)=\{z=(x,y)0x<y,|z|<a\}.`$
Now, we clearly have the estimates
$$C^1\beta _j(z)\text{ and }|\beta _j(z)|C,zN(a).$$
(6.3)
Let $`w=(u,v)\mathrm{\Omega }^+(a)`$ and set $`\rho =vu`$. Then $`\rho `$ is a positive real number that depends on $`w`$ and we define
$$=_w=\{z|z|<\delta \rho \}.$$
Take $`z`$ and let $`\alpha ^2`$ be linearly independent from $`wz`$ and such that $`|\alpha |\delta \rho `$. Introduce the mapping
$$\gamma (s,t)=(x(s,t),y(s,t))=z+(s1/2)(1t)\alpha +t(wz)$$
(6.4)
defined for all $`(s,t)`$ in a small open neighborhood of $`[0,1]\times [0,1]`$. It is then immediate that assertions (a) and (b) in the lemma are satisfied.
In order to verify (c) we compute the $`t`$-derivative of $`h_j(\gamma (s,t))`$:
$$\begin{array}{cc}\hfill _t(h_j(\gamma (s,t)))=(y(s,t)& g_j(x(s,t)))_t(\beta _j(\gamma (s,t)))\hfill \\ & +\beta _j(\gamma (s,t))(_ty(s,t)g_j^{}(x(s,t))_tx(s,t)).\hfill \end{array}$$
(6.5)
We see that
$$\left|_t(\beta _j(\gamma (s,t)))\right|Ca.$$
(6.6)
Since $`g_j(x)g_m(x)=x`$ when $`0<x<a`$ we may write
$$g_j(x)=xp_j(x),$$
where $`p_j(x)0`$. If $`p_j(x)0`$ then $`p_j(x)=x^{\mu _j}q_j(x)`$, where $`\mu _j`$ is a positive integer and $`q_j(0)>0`$. By taking $`a`$ sufficiently small we may then assume that
$$p_j^{}(x)=\mu _jx^{\mu _j1}q_j(x)+x^{\mu _j}q_j^{}(x)C^1p_j(x)/x,0<x<a.$$
(6.7)
Moreover, since $`x(s,t)=(1t)x(s,0)+tx(s,1)(1t)x(s,0)`$ it follows that $`|x(s,t)|C\delta \rho `$ if $`x(s,t)0`$. Hence there is a constant $`C`$ such that
$$\left|p_j^{}(x(s,t))p_j^{}(|x(s,t)|)\right|C\delta \rho ,0s,t1.$$
(6.8)
Next, one has
$$\begin{array}{cc}\hfill y(s,t)& g_j(x(s,t))=(1t)y(s,0)+ty(s,1)x(s,t)+p_j(x(s,t))\hfill \\ & =(1t)y(s,0)+ty(s,1)(1t)x(s,0)tx(s,1)+p_j(x(s,t))\hfill \\ & =(1t)(y(s,0)x(s,0))+t(y(s,1)x(s,1))+p_j(x(s,t))\hfill \\ & =(1t)(y(s,0)x(s,0))+t\rho +p_j(x(s,t)).\hfill \end{array}$$
(6.9)
Recall that $`w\mathrm{\Omega }^+(a)`$, so that in particular $`|w|<a`$. Since $`|z|<\delta \rho `$ and $`|\alpha |\delta \rho `$ it follows from (6.4) that $`|x(s,t)|<a`$ if $`\delta `$ is small enough. We then deduce from (6.8) and (6.9) that
$$|y(s,t)g_j(x(s,t))|C\rho +p_j(|x(s,t)|).$$
(6.10)
Using (6.7) and (6.8) we find that
$$\begin{array}{cc}\hfill _ty(s,t)& (_tx(s,t))g_j^{}(x(s,t))=\rho (y(s,0)x(s,0))+(_tx(s,t))p_j^{}(x(s,t))\hfill \\ & =\rho (y(s,0)x(s,0))+(x(s,1)x(s,0))p_j^{}(x(s,t))\hfill \\ & =\rho (y(s,0)x(s,0))x(s,0)p_j^{}(x(s,t))+x(s,1)p_j^{}(x(s,t))\hfill \\ & (1C\delta )\rho +x(s,1)p_j^{}(x(s,t))(12C\delta )\rho +x(s,1)p_j^{}(|x(s,t)|)\hfill \\ & (12C\delta )\rho +C^1p_j(|x(s,t)|).\hfill \end{array}$$
We now choose $`\delta `$ small enough so that e.g. $`2C\delta <1/2`$. This gives the inequality
$$_ty(s,t)(_tx(s,t))g_j^{}(x(s,t))C^1(\rho +p_j(|x(s,t)|)).$$
(6.11)
Combining (6.11) with (6.3), (6.5), (6.6) and (6.10) we get
$$\begin{array}{cc}\hfill _th_j(\gamma (s,t))& \beta _j(\gamma (s,t))(_ty(s,t)(_tx(s,t))g_j^{}(x(s,t)))\hfill \\ & |(y(s,t)g_j(x(s,t)))_t\beta _j(\gamma (s,t))|\hfill \\ & C^2(\rho +p_j(|x(s,t)|))C^2a(\rho +p_j(|x(s,t)|))\hfill \\ & =(C^2C^2a)(\rho +p_j(|x(s,t)|)).\hfill \end{array}$$
Taking $`a<C^4/2`$ we obtain a positive bound from below for the right-hand side in the last expression, which completes the proof of the lemma. โ
## 7. Examples and Further Problems
### 7.1. The Necessity of Non-degeneracy Assumptions
If one of the cones $`\sigma _i(p)`$ in (4.3) is a line it may happen that $`W(p)\{p\}`$ is the union of two components $`W(p)_l`$ and $`W(p)_r`$, each bounded by level curves as above. In this case there might be several different subharmonic $`PH`$ functions that satisfy condition (i) in Theorem 1, as shown by Example 3 below. Hence something like condition (ii) is indeed necessary in order to obtain the conclusion of the aforementioned theorem.
###### Example 3.
Set $`H_1(x,y)=0`$, $`H_2(x,y)=4x+x^2y^2`$, and $`H_3(x,y)=x`$. There are three level curves through $`(0,0)`$ to functions of the form $`H_iH_j`$ with $`ij`$. These are depicted in Figure 1. Let $`\phi =\mathrm{max}\{H_10,H_2,H_3\}`$. The functions in the figure closest to the origin in each sector are the restriction of $`\phi `$ to that sector.
If one instead defines $`\mathrm{\Psi }(x,y)`$ by changing the value in the two upper sectors from $`0`$ to $`H_3`$ respectively $`H_2`$ then one obtains a different continuous $`PH`$ function that is again subharmonic. Clearly, every neighborhood of the origin still has the property that $`\mathrm{\Psi }`$ is equal to each of the three harmonic functions in some subset of positive Lebesgue measure. So $`\mathrm{\Psi }`$ is a maximum of harmonic functions along the curves, hence trivially subharmonic away from the origin. Letting $`0\chi C_0^{\mathrm{}}()`$ be equal to $`1`$ near the origin and $`\chi _ฯต(z):=\chi (z/ฯต),`$ this implies that $`(1\chi _ฯต)\mathrm{\Delta }\mathrm{\Psi }0`$ in $`๐^{}`$. But clearly $`\chi _ฯต\mathrm{\Delta }\mathrm{\Psi }0`$ in $`๐^{}`$ as $`ฯต0`$ since $`\mathrm{\Psi }=O(|z|)`$. Hence $`\mathrm{\Psi }`$ is subharmonic.
### 7.2. On Global Descriptions
In this paper we have only considered the problem of locally characterizing the maximum of a finite number of harmonic functions. A natural question is to study various situations when a subharmonic $`PH`$ function is globally the maximum of a finite number of harmonic functions. Such a situation occurs for instance in , where the given harmonic functions are linear. The same conclusion holds when the number of given harmonic functions is two as well as in certain other cases. We discuss some of these cases in the following examples, which were inspired by .
###### Example 4.
Let $`A_1`$ and $`A_2`$ be entire functions such that $`A_1(z)A_2(z)`$, $`z`$ and assume that $`\mathrm{\Phi }:=\chi _1A_1+\chi _2A_2`$ satisfies $`\mathrm{\Phi }/\overline{z}0`$, where $`\chi _1`$ and $`\chi _2`$ are the characteristic functions of the sets $`M_1`$ and $`M_2`$, respectively (cf. Notation 1). The first assumption implies that $`H_i(z)=\mathrm{}\left[_0^zA_i(w)๐w\right]`$, $`i=1,2`$ are well-defined functions in $``$ and that there are no singular points for $`H_1H_2`$. For simplicity assume further that level curves to $`H_1H_2`$ as well as the support $`\mathrm{\Phi }/\overline{z}`$ are connected. If $`p\overline{M}_1\overline{M}_2`$, it follows from Theorem 1 (condition (ii) there being vacuous in this case) that there exists a neighborhood $`N`$ of $`p`$ and constants $`c_1(p)`$, $`c_2(p)`$ such that
$$\mathrm{\Phi }=2\frac{}{z}\mathrm{max}(H_1+c_1(p),H_2+c_2(p))=2\frac{}{z}\mathrm{max}(H_1,H_2+c_2(p)c_1(p))$$
In particular, the common boundary of $`M_1`$ and $`M_2`$ in $`N`$ is the level curve $`H_1H_2=c_2(p)c_1(p)`$ and this is also the support of $`\mathrm{\Phi }/\overline{z}`$ in $`N`$. The local information implies, by the connectedness assumptions, that globally $`c_2(p)c_1(p)`$ is a constant $`c`$ independent of $`p`$, and that the support actually consists of the level curve $`H_1H_2=c`$, and finally that
$$\mathrm{\Phi }=2\frac{}{z}\mathrm{max}(H_1,H_2+c).$$
###### Example 5.
This example is essentially one-dimensional. Assume that
$$=\underset{j=1}{\overset{r}{}}\overline{I}_j,$$
where the $`I_j`$ are open pairwise disjoint intervals. Set $`M_j=I_j\times `$, $`1jr`$, and let $`\chi _j(x)`$ be the characteristic function of $`I_j`$, which we also view as the characteristic function of $`M_j`$. Let $`h_j(x+\sqrt{1}y)=a_jx+b_j`$, $`1jr`$, be linear functions on $``$ and assume as usual that
$$\chi :=\frac{}{\overline{z}}\left[\underset{j=1}{\overset{r}{}}\frac{h_j(z)}{z}\chi _j\right]=\underset{j=1}{\overset{r}{}}\frac{h_j(z)}{z}\frac{\chi _j}{\overline{z}}=\underset{j=1}{\overset{r}{}}\frac{a_j}{2}\frac{\chi _j}{\overline{z}}0.$$
Since $`\frac{\chi _j}{\overline{z}}=\frac{1}{2}\frac{\chi _j}{x}`$ we deduce that $`_{j=1}^ra_j\chi _j`$ is an increasing function of $`x`$ and thus $`h(x)=_0^x_{j=1}^ra_j\chi _j`$ is a convex function. Set
$$H(x,y)=h(y)+h^{}(y+0)(xy).$$
By convexity we have
$$h(x)H(x,y),x,y,$$
(7.1)
with equality when $`y=x`$. The functions $`H(x,y)`$ viewed as linear functions of $`x`$ are independent of $`y`$ when $`yI_j`$. We denote their common value for $`yI_j`$ by $`\stackrel{~}{h}_j(x)`$ and notice that $`\stackrel{~}{h}_jh_j=C_j`$, where $`C_j`$ is a constant. It follows from (7.1) that
$$h(x)=\underset{1kr}{\mathrm{max}}\stackrel{~}{h}_k(x)\text{ in }M_j$$
and then differentiation implies that
$$h^{}(x)=\frac{}{x}\underset{1kr}{\mathrm{max}}\stackrel{~}{h}_k(x)=\frac{}{x}\underset{1kr}{\mathrm{max}}\left(h_k(x)+C_k\right).$$
This means precisely that the $`PA`$ function $`\chi `$ satisfies
$$\chi =2\frac{}{z}\underset{1jr}{\mathrm{max}}\left(h_j(z)+C_j\right)$$
and is therefore globally the maximum of a finite number of harmonic functions.
### 7.3. Related Questions
Let us finally formulate and discuss some interesting related problems.
###### Problem 1.
At the moment we do not know although we strongly suspect that locally there are in fact only a finite number of possibilities for $`\mathrm{\Psi }`$ even when conditions (i)โ(iii) are weakened in Theorem 1. This holds e.g. for the function constructed in Example 3. In particular, it seems likely that there always exists a sufficiently small neighborhood of $`p`$ that can be dissected into sectors bounded by level curves to $`H_iH_j`$ such that $`\mathrm{\Psi }`$ is constant in each such sector. Example 3 suggests that the local behavior of a $`PH`$ subharmonic function is determined by the geometry of the level curves $`\mathrm{\Gamma }_{i,j,k}`$ whose study is essentially a problem of a combinatorial and topological nature. It would be interesting to give a description of this local behavior in terms of Morse theory (the study of level curves was Morseโs original motivation for his theory, see ).
###### Problem 2.
Another problem is to understand the global behavior of a $`PH`$ subharmonic function and in particular to give criteria saying precisely when $`\frac{\mathrm{\Psi }}{z}`$ is the derivative of the maximum of a finite number of harmonic functions as in the last two examples. This would have interesting applications to uniqueness theorems for Cauchy transforms that are algebraic functions as in .
###### Problem 3.
There are also several connections between the questions studied in the present paper and the theory of asymptotic solutions to differential equations. For instance, sets like those that occur as the support of the measures in Theorem 2 play a remarkable role in the latter theory (). Moreover, many similar techniques are used, e.g. the admissible sets in are closely related to (though not exactly the same as) the sets $`V(z)`$ in Lemma 3 above. These connections are quite close in the cases studied in (as well as other cases) and certainly deserve further investigation in view of their important applications.
###### Problem 4.
Let $`U`$ be a domain in $`C^n`$, where $`n1`$. By analogy with Definition 1 and Notation 1 one can define the notions of $`PH_n`$ and $`PA_n`$ functions in $`U`$ as natural higher-dimensional generalizations of the concepts of $`PH`$ and $`PA`$ functions, respectively. It seems reasonable to conjecture that appropriate higher-dimensional analogue of Theorem 1 hold for the class $`PA_n`$ and that as a consequence one would get a natural extension of e.g. Corollary 2 to the class $`PH_n`$.
## Appendix. Comments on Some Properties and Definitions
As before, $`\chi _\mathrm{\Omega }`$ denotes the characteristic function of a set $`\mathrm{\Omega }`$ (or $`^2`$). Let us introduce the following additional condition: an open set $`\mathrm{\Omega }^2`$ is said to have property (\*) if $`\mathrm{\Omega }`$ is of Lebesgue measure $`0`$ and $`_z\chi _\mathrm{\Omega },_y\chi _\mathrm{\Omega }`$ are measures.
###### Lemma 10.
If $`\mathrm{\Omega }_1,\mathrm{\Omega }_2^2`$ have property (\*) then so does $`\mathrm{\Omega }_1\mathrm{\Omega }_2`$.
###### Proof.
It is clear that $`(\mathrm{\Omega }_1\mathrm{\Omega }_2)`$ is Lebesgue negligible. Let $`K^2`$ be any compact set, choose $`\eta C_0^{\mathrm{}}(^2)`$ with $`\eta (x,y)๐x๐y=1`$, define $`\eta _ฯต=ฯต^2\eta (x/ฯต,y/ฯต)`$ for $`ฯต(0,1)`$ and set $`\chi _{j,ฯต}=\chi _j\eta _ฯต`$, where $`\chi _j=\chi _{_{\mathrm{\Omega }_j}}`$, $`j=1,2`$. Then $`0\chi _{j,ฯต}1`$, $`\chi _{j,ฯต}\chi _j`$ a.e. as $`ฯต0`$ and $`_x\chi _{j,ฯต}_{L^1(K)}=\eta _ฯต_x\chi _j_{L^1(K)}C_K,`$ where $`C_K`$ is independent of $`ฯต`$. Since $`_x(\chi _{1,ฯต}\chi _{2,ฯต})=\chi _{1,ฯต}_x\chi _{2,ฯต}+\chi _{2,ฯต}_x\chi _{1,ฯต}`$ it follows that if $`\varphi C_0^{\mathrm{}}(^2)`$ then
$$\begin{array}{c}\left|\chi _{1,ฯต}(x,y)\chi _{2,ฯต}(x,y)_x\varphi (x,y)dxdy\right|\hfill \\ \hfill |\varphi (x,y)|(|_x\chi _{1,ฯต}(x,y)|+|_x\chi _{2,ฯต}(x,y)|)๐x๐y2C_K\varphi _L^{\mathrm{}}.\end{array}$$
When $`ฯต0`$ this shows that
$$\left|\chi _1(x,y)\chi _2(x,y)_x\varphi (x,y)dxdy\right|2C_K\varphi _L^{\mathrm{}}$$
and thus $`_x(\chi _1\chi _2)`$ is a distribution of order $`0`$ (which extends to a measure). This finishes the proof since $`_y(\chi _1\chi _2)`$ can be dealt with in the same way. โ
Lemma 10 shows that if we define sets $`P^{}X`$ of functions โpiecewise\* in $`X`$โ as in Definition 1 by demanding in addition that all sets $`M_i`$ have property (\*) then $`P^{}X`$ are again vector spaces.
###### Lemma 11.
If $`uP^{}X`$ is continuous then $`_xu,_yuP^{}X`$, where derivatives are taken in the distribution sense.
###### Proof.
Let us write $`u=_{i=1}^ru_i\chi _i`$, where $`\chi _i`$ is the characteristic function of the (open) set $`M_i`$, $`_{i=1}^r\chi _i=1`$ a.e. and $`_x\chi _i,_y\chi _i`$ are measures, $`1ir`$. Since $`u`$ is continuous we can find $`u_ฯตC^{\mathrm{}}(U)`$ tending uniformly to $`u`$ on every compact set as $`ฯต0`$. Now
$$\begin{array}{cc}\hfill _xu_ฯต& =\underset{i=1}{\overset{r}{}}(_xu_i)\chi _i+\underset{i=1}{\overset{r}{}}\chi _i_x(u_ฯตu_i)\hfill \\ & =\underset{i=1}{\overset{r}{}}(_xu_i)\chi _i+_x\left(\underset{i=1}{\overset{r}{}}(u_ฯตu_i)\chi _i\right)\underset{i=1}{\overset{r}{}}(u_ฯตu_i)_x\chi _i.\hfill \end{array}$$
(A1)
For every $`i`$ one has $`u_ฯตu_i=u_ฯตu`$ in a dense subset of $`M_i`$. It follows that $`u_ฯตu_i`$ uniformly on every compact subset of $`\overline{M}_i`$, hence also on every compact subset of the support of the measure $`_x\chi _i`$. Therefore, $`(u_ฯตu_i)_x\chi _i0`$ in $`๐^{}(^2)`$ as $`ฯต0`$. This is true for $`(u_ฯตu_i)\chi _i`$ as well and so by letting $`ฯต0`$ in (A1) we conclude that $`_xu_ฯต=_{i=1}^r(_xu_i)\chi _i`$. The same argument applies to $`_yu`$. โ
Given a domain $`U`$ let $`S(U)`$ be the class of subharmonic functions in $`U`$. Recall Notation 2 from ยง1, where we already noted the (well-known) fact that $`_z\varphi \mathrm{\Sigma }(U)`$ whenever $`\varphi S(U)`$. For completeness we give here a proof of a (also well-known) partial converse to this statement.
###### Lemma 12.
If $`U`$ is simply connected and $`f\mathrm{\Sigma }(U)`$ then $`f=_z\varphi `$ for some $`\varphi S(U)`$ which is uniquely determined modulo an additive constant.
###### Proof.
Since the operator $`_z`$ is elliptic we may write $`f=_zw`$, where $`w=u+iv๐^{}(U)`$ (cf., e.g., ). We get $`\mathrm{\Delta }u+i\mathrm{\Delta }v=\mathrm{\Delta }w=4_{\overline{z}}_zw=4_{\overline{z}}f0,`$ which implies that $`uS(U)`$, $`vH(U)`$, and thus $`f=_zu+g`$, where $`g=i_zvA(U)`$. Let $`GA(U)`$ be such that $`G^{}(z)=g(z)`$ and define $`\varphi =u+G+\overline{G}`$. Then $`\varphi S(U)`$ and $`_z\varphi =_zu+_zG=_zu+g=f.`$ The last assertion in the lemma follows from the fact that a function $`h`$ in $`U`$ is constant whenever $`h=\overline{h}`$ and $`_zh=0`$. โ
## Acknowledgements
We would like to thank Jan-Erik Bjรถrk and Anders Melin for stimulating discussions and useful comments. We are especially grateful to the anonymous referee for his detailed reports (articles in their own right!) with numerous insightful suggestions and an alternative approach for deriving results similar to Theorem 1 under some mild extra assumptions (see Theorem 3 in ยง6). With his kind permission we reproduced large parts of his reports in ยง1, ยง6 and the Appendix.
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# Non-Gaussianity of Large-Scale Cosmic Microwave Background Anisotropies beyond Perturbation Theory
## I Introduction
Cosmological inflation lrreview is the dominant paradigm to understand the initial conditions for CMB anisotropies and structure formation. In the inflationary picture, the primordial cosmological perturbations are created from quantum fluctuations โredshiftedโ out of the horizon during an early period of superluminal expansion of the universe, where they remain โfrozenโ. They are observable as temperature anisotropies in the CMB at the last scattering surface. They were first detected by the Cosmic Background Explorer (COBE) satellite smoot92 ; bennett96 ; gorski96 . The last and most impressive confirmation of the inflationary paradigm has been recently provided by the data of the Wilkinson Microwave Anisotropy Probe (WMAP) mission wmap1 . Since the observed cosmological perturbations are of the order of $`10^5`$, one might think that first-order perturbation theory will be adequate for all comparison with observations. That may not be the case however, as the Planck satellite planck and its successors may be sensitive to the non-Gaussianity of the cosmological perturbations at the level of second- or higher-order perturbation theory review . Statistics like the bispectrum and the trispectrum of the CMB can be used to assess the level of primordial non-Gaussianity on various cosmological scales and to discriminate it from the one induced by secondary anisotropies and systematic effects review ; hu ; dt ; jul . Therefore, it is of fundamental importance to provide accurate theoretical predictions for the statistics of the large-angle CMB anisotropies as left imprinted by the primordial seeds originated during or immediately after inflation. Steps towards this goal have been taken in Refs. pc ; mm ; prl ; tomita at the level of second order perturbation theory.
In this paper we derive an expression for the anisotropies of the CMB on scales larger than the horizon at last scattering which is valid at any order in perturbation theory, providing a fully non-linear generalization to the Sachs-Wolfe effect at first- sw and second-order prl . In particular, for the standard single-field models of inflation, we provide the exact non-perturbative expression for the bispectrum and the trispecturm in the so-called โsqueezedโ limit in which some of the wavenumbers are much smaller than the others. Furthermore, we compute the generic expressions for the non-linearity parameters $`f_{\mathrm{NL}}`$ and $`g_{\mathrm{NL}}`$ characterizing respectively the quadratic and cubic non-linearity in the large-angle CMB anisotropies.
The paper is organized as follows. In Section II we provide the non-linear generalization of the Sachs-Wolfe effect which is expressed in terms of the comoving curvature perturbation from inflation in section III. In Section IV we show how to compute the connected $`n`$-point correlation functions of the CMB anisotropies, leaving the details to the Appendix. Finally, we draw our conclusions in Section V.
## II The non-linear Sachs-Wolfe effect
Our starting point is the Arnowitt-Deser-Misner (ADM) formalism which is particularly useful to deal with the non-linear evolution of cosmological perturbations. The line element is
$$ds^2=N^2dt^2+N_idtdx^i+\gamma _{ij}dx^idx^j,$$
where the three-metric $`\gamma _{ij}`$, the lapse $`N`$ and the shift $`N_i`$ functions describe the evolution of timelike hypersurfaces. In the ADM formalism the equations simplify considerably if we set $`N^i=0`$. Moreover we are interested only in scalar perturbations in a flat Universe and therefore we find it convenient to recast the metric as
$$ds^2=e^{2\mathrm{\Phi }}dt^2+a^2(t)e^{2\mathrm{\Psi }}\delta _{ij}dx^idx^j,$$
(1)
where $`a(t)`$ is the scale factor describing the evolution of the homogeneous and isotropic Universe and we have introduced two gravitational potentials $`\mathrm{\Phi }`$ and $`\mathrm{\Psi }`$. The expression (1) holds at any order in perturbation theory. To make contact with the usual perturbative approach, one may expand the gravitational potentials at first- and second-order, e.g, $`\mathrm{\Phi }=\mathrm{\Phi }_1+\mathrm{\Phi }_2/2`$. From Eq. (1) one recovers at linear order the well-known longitudinal gauge, $`N^2=(1+2\mathrm{\Phi }_1)`$ and $`\gamma _{ij}=a^2(12\mathrm{\Psi }_1)\delta _{ij}`$. At second-order, one finds $`\mathrm{\Phi }_2=\varphi _22\varphi _1^2`$ and $`\mathrm{\Psi }_2=\psi _2+2\psi _1^2`$ where $`\varphi _1`$, $`\psi _1`$ and $`\varphi _2`$, $`\psi _2`$ (with $`\varphi _1=\mathrm{\Phi }_1`$ and $`\psi _1=\mathrm{\Psi }_1`$) are the first and second-order gravitational potentials in the longitudinal (Poisson) gauge adopted in Refs. MMB ; review , $`N^2=(1+2\varphi _1+\varphi _2)`$ and $`\gamma _{ij}=a^2(12\psi _1\psi _2)\delta _{ij}`$ as far as scalar perturbations are concerned. In writing Eq. (1) we have neglected vector and tensor perturbation modes. For the vector perturbations the reason is that we are interested in long-wavelength perturbations, i.e. on scales larger than the horizon at last scattering, while vector modes will contain gradient terms being produced as non-linear combination of scalar-modes and thus they will be more important on small scales (linear vector modes are not generated in standard mechanisms for cosmological perturbations, as inflation). For example the results of Ref. MMB show clearly this for second-order perturbations. The tensor contribution can be neglected for two reasons. First, the tensor perturbations produced from inflation on large scales give a negligible contribution to the higher-order statistics of the Sachs-Wolfe effect being of the order of (powers of) the slow-roll parameters during inflation (this holds for linear tensor modes as well as for tensor modes generated by the non-linear evolution of scalar perturbations during inflation, for example see the results of Ref. maldacena for second-order perturbations). Moreover, while on large scales the tensor modes have been proven to remain constant in time Salopek1 , when they approach the horizon they have a wavelike contribution which oscillates with decreasing amplitude.
Since we are interested in the cosmological perturbations on large scales, that is in perturbations whose wavelength is larger than the Hubble radius at last scattering, a local observer would see them in the form of a classical โ possibly time-dependent โ (nearly zero-momentum) homogeneous and isotropic background. Therefore, it should be possible to perform a change of coordinates in such a way as to absorb the super-Hubble modes and work with a metric of an homogeneous and isotropic Universe (plus, of course, cosmological perturbations on scale smaller than the horizon). We split the gravitational potential $`\mathrm{\Phi }`$ as
$$\mathrm{\Phi }=\mathrm{\Phi }_{\mathrm{}}+\mathrm{\Phi }_s,$$
(2)
where $`\mathrm{\Phi }_{\mathrm{}}`$ stands for the part of the gravitational potential receiving contributions only from the super-Hubble modes; $`\mathrm{\Phi }_s`$ receives contributions only from the sub-horizon modes
$`\mathrm{\Phi }_{\mathrm{}}`$ $`=`$ $`{\displaystyle \frac{d^3k}{(2\pi )^3}\theta \left(aHk\right)\mathrm{\Phi }_\stackrel{}{k}e^{i\stackrel{}{k}\stackrel{}{x}}},`$
$`\mathrm{\Phi }_s`$ $`=`$ $`{\displaystyle \frac{d^3k}{(2\pi )^3}\theta \left(kaH\right)\mathrm{\Phi }_\stackrel{}{k}e^{i\stackrel{}{k}\stackrel{}{x}}},`$ (3)
where $`H`$ is the Hubble rate computed with respect to the cosmic time, $`H=\dot{a}/a`$, and $`\theta (x)`$ is the step function. Analogous definitions hold for the other gravitational potential $`\mathrm{\Psi }`$.
By construction $`\mathrm{\Phi }_{\mathrm{}}`$ and $`\mathrm{\Psi }_{\mathrm{}}`$ are a collection of Fourier modes whose wavelengths are larger than the horizon length and we may safely neglect their spatial gradients. Therefore $`\mathrm{\Phi }_{\mathrm{}}`$ and $`\mathrm{\Psi }_{\mathrm{}}`$ are only functions of time. This amounts to saying that we can absorb the large-scale perturbations in the metric (1) by the following redefinitions
$`d\overline{t}`$ $`=`$ $`e^\mathrm{\Phi }_{\mathrm{}}dt,`$ (4)
$`\overline{a}`$ $`=`$ $`ae^\mathrm{\Psi }_{\mathrm{}}.`$ (5)
The new metric describes a homogeneous and isotropic Universe
$$ds^2=d\overline{t}^2+\overline{a}^2\delta _{ij}dx^idx^j,$$
(6)
where for simplicity we have not included the sub-horizon modes. On super-horizon scales one can regard the Universe as a collection of regions of size of the Hubble radius evolving like unperturbed patches with metric (6) Salopek1 .
Let us now go back to the quantity we are interested in, namely the anisotropies of the CMB as measured today by an observer $`๐ช`$. If she/he is interested in the CMB anisotropies at large scales, the effect of super-Hubble modes is encoded in the metric (6). During their travel from the last scattering surface โ to be considered as the emitter point $``$ โ to the observer, the CMB photons suffer a redshift determined by the ratio of the emitted frequency $`\overline{\omega }_{}`$ to the observed one $`\overline{\omega }_๐ช`$
$$\overline{T}_๐ช=\overline{T}_{}\frac{\overline{\omega }_๐ช}{\overline{\omega }_{}},$$
(7)
where $`\overline{T}_๐ช`$ and $`\overline{T}_{}`$ are the temperatures at the observer point and at the last scattering surface, respectively.
What is then the temperature anisotropy measured by the observer? The expression (7) shows that the measured large-scale anisotropies are made of two contributions: the intrinsic inhomogeneities in the temperature at the last scattering surface and the inhomogeneities in the scaling factor provided by the ratio of the frequencies of the photons at the departure and arrival points. Let us first consider the second contribution. As the frequency of the photon is the inverse of a time period, we get immediately the fully non-linear relation
$$\frac{\overline{\omega }_{}}{\overline{\omega }_๐ช}=\frac{\omega _{}}{\omega _๐ช}e^{\mathrm{\Phi }_{\mathrm{}}+\mathrm{\Phi }_\mathrm{}๐ช}.$$
(8)
As for the temperature anisotropies coming from the intrinsic temperature fluctuation at the emission point, it maybe worth to recall how to obtain this quantity in the longitudinal gauge at first order. By expanding the photon energy density $`\rho _\gamma T_\gamma ^4`$, the intrinsic temperature anisotropies at last scattering are given by $`\delta _1T_{}/T_{}=(1/4)\delta _1\rho _\gamma /\rho _\gamma `$. One relates the photon energy density fluctuation to the gravitational perturbation first by implementing the adiabaticity condition $`\delta _1\rho _\gamma /\rho _\gamma =(4/3)\delta _1\rho _m/\rho _m`$, where $`\delta _1\rho _m/\rho _m`$ is the relative fluctuation in the matter component, and then using the energy constraint of Einstein equations $`\mathrm{\Phi }_1=(1/2)\delta _1\rho _m/\rho _m`$. The result is $`\delta _1T_{}/T_{}=2\mathrm{\Phi }_1/3`$. Summing this contribution to the anisotropies coming from the redshift factor (8) expanded at first order provides the standard (linear) Sachs-Wolfe effect $`\delta _1T_๐ช/T_๐ช=\mathrm{\Phi }_1/3`$. Following the same steps, we may easily obtain its full non-linear generalization.
Let us first relate the photon energy density $`\overline{\rho }_\gamma `$ to the energy density of the non-relativistic matter $`\overline{\rho }_m`$ by using the adiabaticity conditon. Again here a bar indicates that we are considering quantities in the locally homogeneous Universe described by the metric (6). Using the energy continuity equation on large scales $`\overline{\rho }/\overline{t}=3\overline{H}(\overline{\rho }+\overline{P})`$, where $`\overline{H}=d\mathrm{ln}\overline{a}/d\overline{t}`$ and $`\overline{P}`$ is the pressure of the fluid, one can easily show that there exists a conserved quantity in time at any order in perturbation theory KMNR
$$\mathrm{ln}\overline{a}+\frac{1}{3}^{\overline{\rho }}\frac{d\overline{\rho ^{}}}{\left(\overline{\rho ^{}}+\overline{P^{}}\right)}.$$
(9)
The perturbation $`\delta `$ is a gauge-invariant quantity representing the non-linear extension of the curvature perturbation $`\zeta `$ on uniform energy density hypersurfaces on superhorizon scales for adiabatic fluids KMNR . Indeed, expanding it at first and second order one gets the corresponding definition $`\zeta _1=\psi _1\delta _1\rho /\dot{\rho }`$ and the quantity $`\zeta _2`$ introduced in Ref. MW . At first order the adiabaticity condition corresponds to set $`\zeta _{1\gamma }=\zeta _{1m}`$ for the curvature perturbations relative to each component. At the non-linear level the adiabaticity condition generalizes to
$$\frac{1}{3}\frac{d\overline{\rho }_m}{\overline{\rho }_m}=\frac{1}{4}\frac{d\overline{\rho }_\gamma }{\overline{\rho }_\gamma },$$
(10)
or
$$\mathrm{ln}\overline{\rho }_m=\mathrm{ln}\overline{\rho }_\gamma ^{3/4}.$$
(11)
To make contact with the standard second-order result, we may expand in Eq. (11) the photon energy density perturbations as $`\delta \overline{\rho }_\gamma /\rho _\gamma =\delta _1\rho _\gamma /\rho _\gamma +\frac{1}{2}\delta _2\rho _\gamma /\rho _\gamma `$, and similarly for the matter component. We immediately recover the adiabaticity condition
$$\frac{\delta _2\rho _\gamma }{\rho _\gamma }=\frac{4}{3}\frac{\delta _2\rho _m}{\rho _m}+\frac{4}{9}\left(\frac{\delta _1\rho _m}{\rho _m}\right)^2$$
(12)
given in Ref. review .
Next we need to relate the photon energy density to the gravitational potentials at the non-linear level. The energy constraint inferred from the (0-0) component of Einstein equations in the matter-dominated era with the โbarredโ metric (6) is
$$\overline{H}^2=\frac{8\pi G_N}{3}\overline{\rho }_m.$$
(13)
Using Eqs. (4) and (5) the Hubble parameter $`\overline{H}`$ reads
$$\overline{H}=\frac{1}{\overline{a}}\frac{d\overline{a}}{d\overline{t}}=e^\mathrm{\Phi }_{\mathrm{}}(H\dot{\mathrm{\Psi }}_{\mathrm{}}),$$
(14)
where $`H=d\mathrm{ln}a/dt`$ is the Hubble parameter in the โunbarredโ metric. Eq. (13) thus yields an expression for the energy density of the non-relativistic matter which is fully nonlinear, being expressed in terms of the gravitational potential $`\mathrm{\Phi }_{\mathrm{}}`$
$$\overline{\rho }_m=\rho _me^{2\mathrm{\Phi }_{\mathrm{}}},$$
(15)
where we have dropped $`\dot{\mathrm{\Psi }}_{\mathrm{}}`$ which is negligible on large scales. By perturbing the expression (15) we are able to recover in a straightforward way the solutions of the (0-0) component of Einstein equations for a matter-dominated Universe in the large-scale limit obtained at second-order in perturbation theory. Indeed, recalling that $`\mathrm{\Phi }`$ is perturbatively related to the quantity $`\varphi =\varphi _1+\varphi _2/2`$ used in Ref. review by $`\mathrm{\Phi }_1=\varphi _1`$ and $`\mathrm{\Phi }_2=\varphi _22(\varphi _1)^2`$, one immediately obtains review ; BMR4
$`{\displaystyle \frac{\delta _1\rho _m}{\rho _m}}`$ $`=`$ $`2\varphi _1,`$
$`{\displaystyle \frac{1}{2}}{\displaystyle \frac{\delta _2\rho _m}{\rho _m}}`$ $`=`$ $`\varphi _2+4(\varphi _1)^2.`$ (16)
The expression for the intrinsic temperature of the photons at the last scattering surface $`\overline{T}_{}\overline{\rho }_\gamma ^{1/4}`$ follows from Eqs. (11) and (15)
$$\overline{T}_{}=T_{}e^{2\mathrm{\Phi }_{\mathrm{}}/3}.$$
(17)
Plugging Eqs. (8) and (17) into the expression (7) we are finally able to provide the expression for the CMB temperature which is fully nonlinear and takes into account both the gravitational redshift of the photons due to the metric perturbations at last scattering and the intrinsic temperature anisotropies
$$\overline{T}_๐ช=\left(\frac{\omega _๐ช}{\omega _{}}\right)T_{}e^{\mathrm{\Phi }_{\mathrm{}}/3}.$$
(18)
From Eq. (18) we read the non-perturbative anisotropy corresponding to the Sachs-Wolfe effect
$$\frac{\delta _{np}\overline{T}_๐ช}{T_๐ช}=e^{\mathrm{\Phi }_{\mathrm{}}/3}1.$$
(19)
Eq. (19) is one of the main results of this paper and represents at any order in perturbation theory the extension of the linear Sachs-Wolfe effect. At first order one gets
$$\frac{\delta _1T_๐ช}{T_๐ช}=\frac{1}{3}\mathrm{\Phi }_1,$$
(20)
and at second order
$$\frac{1}{2}\frac{\delta _2T_๐ช}{T_๐ช}=\frac{1}{6}\mathrm{\Phi }_2+\frac{1}{18}\left(\mathrm{\Phi }_1\right)^2,$$
(21)
which exactly reproduces the generalization of the Sachs-Wolfe effect at second-order in the perturbations found in Ref. review ; BMR4 (where $`\mathrm{\Phi }_1=\varphi _1`$ and $`\mathrm{\Phi }_2=\varphi _22(\varphi _1)^2`$).
## III Relating the CMB anisotropies to the inflationary comoving curvature perturbation
In this section we relate the gravitational potentials $`\mathrm{\Phi }_{\mathrm{}}`$ to $`\mathrm{\Psi }_{\mathrm{}}`$ to the curvature perturbation $`\zeta =\delta `$ at any order in perturbation theory (for notational simplicity we drop the subscrip โ$`\mathrm{}`$โ from now on). This will allow to express the non-linear temperature fluctuations in terms of the initial conditions provided by inflation. We use the evolution equation in the ADM formalism
$$N_{|k}^{|i}+\frac{1}{3}N_{|l}^{|l}\delta _k^i+N^{(3)}\overline{R}_k^i=N\mathrm{\hspace{0.17em}\hspace{0.17em}8}\pi G\overline{S}_k^i,$$
(22)
where a vertical bar denotes a covariant derivative with respect to $`\gamma _{ij}`$ and we have used that fact that the traceless part of the extrinsic curvature $`\overline{K}_{ij}=K_{ij}K\gamma _{ij}/3`$ vanishes in our metric (6). The extrinsic curvature is defined as $`K_{ij}=\dot{\gamma _{ij}}/(2N)`$ and $`K=\gamma ^{ij}K_{ij}`$. Here $`{}_{}{}^{(3)}\overline{R}_{k}^{i}=\gamma ^{ij}{}_{}{}^{(3)}\overline{R}_{jk}^{}=^{(3)}R_k^i^{(3)}R\delta _k^i/3`$ where $`{}_{}{}^{(3)}R_{jk}^{}`$ is the Ricci tensor of constant time hypersurfaces associated with the metric $`\gamma _{ij}`$. Analogous definitions hold for $`\overline{S}_k^i`$ constructed from the matter stress-energy three-tensor $`S_{ij}=T_{ij}`$.
From Eq. (22) we want to obtain a constraint between the gravitational potentials $`\mathrm{\Phi }`$ and $`\mathrm{\Psi }`$ keeping track of non-local terms. Thus we cannot neglect the gradient terms appearing in this equation. Since $`N=e^\mathrm{\Phi }`$ we find
$`N_{|k}^{|i}`$ $`=`$ $`{\displaystyle \frac{e^{2\mathrm{\Psi }}e^\mathrm{\Phi }}{a^2(t)}}\left(\mathrm{\Phi }_{,k}\mathrm{\Phi }^{,i}+\mathrm{\Phi }_{,k}^{,i}+\mathrm{\Psi }_{,k}\mathrm{\Phi }^{,i}+\mathrm{\Psi }^{,i}\mathrm{\Phi }_{,k}\mathrm{\Psi }_{,j}\mathrm{\Phi }^{,j}\delta _k^i\right),`$
$`{}_{}{}^{(3)}R_{k}^{i}`$ $`=`$ $`{\displaystyle \frac{e^{2\mathrm{\Psi }}}{a^2(t)}}\left(\mathrm{\Psi }_{,k}^{,i}+^2\mathrm{\Psi }\delta _k^i+\mathrm{\Psi }^{,i}\mathrm{\Psi }_{,k}\mathrm{\Psi }_{,m}\mathrm{\Psi }^{,m}\delta _k^i\right),`$
$`{}_{}{}^{(3)}R`$ $`=`$ $`2{\displaystyle \frac{e^{2\mathrm{\Psi }}}{a^2(t)}}\left(2^2\mathrm{\Psi }\mathrm{\Psi }_{,i}\mathrm{\Psi }^{,i}\right).`$ (23)
Hence Eq. (22) reads
$``$ $`\left(\mathrm{\Phi }_{,k}\mathrm{\Phi }^{,i}+\mathrm{\Phi }_{,k}^{,i}+\mathrm{\Psi }_{,k}\mathrm{\Phi }^{,i}+\mathrm{\Psi }^{,i}\mathrm{\Phi }_{,k}\mathrm{\Psi }_{,j}\mathrm{\Phi }^{,j}\delta _k^i\right)+{\displaystyle \frac{1}{3}}\left(^2\mathrm{\Phi }+\mathrm{\Phi }^{,l}\mathrm{\Phi }_{,l}\mathrm{\Psi }_{,l}\mathrm{\Phi }^{,l}\right)\delta _k^i`$
$`+`$ $`\left(\mathrm{\Psi }_{,k}^{,i}+^2\mathrm{\Psi }\delta _k^i+\mathrm{\Psi }^{,i}\mathrm{\Psi }_{,k}\mathrm{\Psi }_{,m}\mathrm{\Psi }^{,m}\delta _k^i\right){\displaystyle \frac{2}{3}}\left(2^2\mathrm{\Psi }\mathrm{\Psi }_{,l}\mathrm{\Psi }^{,l}\right)\delta _k^i=8\pi Ga^2(t)e^{2\mathrm{\Psi }}\overline{S}_k^i,`$
where the indices of the partial derivatives are raised by $`\delta ^{ij}`$. Notice that at first order this equation gives the usual constraint $`\mathrm{\Phi }_1=\mathrm{\Psi }_1`$. As far as the matter content is concerned we can consider the perfect fluid energy momentum tensor $`T_{\mu \nu }=\left(\overline{\rho }+\overline{P}\right)u_\mu u_\nu +\overline{P}g_{\mu \nu }`$ and we find
$`S_k^i`$ $`=`$ $`\gamma ^{ij}T_{jk}=a^2(t)e^{2\mathrm{\Psi }}\left(\overline{\rho }+\overline{P}\right)u^iu_k+\overline{P}\delta _k^i,`$
$`\overline{S}_k^i`$ $`=`$ $`a^2(t)e^{2\mathrm{\Psi }}\left(\overline{\rho }+\overline{P}\right)(u^iu_ku^iu_i\delta _k^i/3),`$ (25)
where $`u^i=\delta ^{ij}u_j`$. We need an expression for the spatial velocities. Thus we use the momentum constraint of the ADM equations which reads Bardeen1
$$8\pi GJ_i=\frac{2}{3}K_{|i},$$
(26)
where $`J_i=NT_i^0=e^\mathrm{\Phi }T_i^0=e^\mathrm{\Phi }T_{0i}`$ is the momentum density and $`K=3e^\mathrm{\Phi }(H(t)\dot{\mathrm{\Psi }})`$. One can use the normalization $`g^{\mu \nu }u_\mu u_\nu =1`$ to express $`u_0`$, finding $`u_0=(1+a^2(t)e^{2\mathrm{\Psi }}u^iu_i)^{1/2}`$. Let us just consider scalar velocities $`u^i=_iu`$ (on sufficiently large scales vector modes may be safely neglected). From Eq. (26) we obtain
$$(\overline{\rho }+\overline{P})u_i(1+a^2(t)e^{2\mathrm{\Psi }}u^iu_i)^{1/2}=\left[\frac{e^\mathrm{\Phi }}{4\pi G}(H(t)\dot{\mathrm{\Psi }})\right]_{|i}$$
(27)
In fact we are interested in the case of non-relativistic matter $`\overline{P}_m=0`$ and $`\overline{\rho }_m=\rho _me^{2\mathrm{\Phi }}`$. As Eq. (III) contains at least two gradients (also $`\overline{S}_k^i`$ contains at least two spatial gradients), using the gradient expansion we may restrict ourselves to the solution
$$u_i=\frac{e^{2\mathrm{\Phi }}}{4\pi G\rho _m}\left[e^\mathrm{\Phi }(H(t)\dot{\mathrm{\Psi }})\right]_{|i}=\frac{e^\mathrm{\Phi }H(t)}{4\pi G\rho _m}\mathrm{\Phi }_{,i},$$
(28)
where we have neglected $`\dot{\mathrm{\Psi }}`$. We thus find
$$\overline{S}_k^i=\frac{e^{2\mathrm{\Psi }}}{6\pi Ga^2}\left(\mathrm{\Phi }^{,i}\mathrm{\Phi }_{,k}\frac{1}{3}\mathrm{\Phi }^{,j}\mathrm{\Phi }_{,j}\delta _k^i\right).$$
(29)
Inserting Eq. (29) into Eq (III) and applying the operator $`_i^k`$ we find
$`^4\left(\mathrm{\Psi }\mathrm{\Phi }\right)`$ $`=`$ $`{\displaystyle \frac{3}{2}}_i^k\left(\mathrm{\Psi }^{,i}\mathrm{\Psi }_{,k}\right)+{\displaystyle \frac{1}{2}}^2\left(\mathrm{\Psi }^{,i}\mathrm{\Psi }_{,i}\right)`$ (30)
$`+`$ $`{\displaystyle \frac{7}{2}}_i^k\left(\mathrm{\Phi }^{,i}\mathrm{\Phi }_{,k}\right){\displaystyle \frac{7}{6}}^2\left(\mathrm{\Phi }^{,i}\mathrm{\Phi }_{,i}\right)`$
$`+`$ $`3_i^k\left(\mathrm{\Phi }^{,i}\mathrm{\Psi }_{,k}\right)^2\left(\mathrm{\Phi }^{,i}\mathrm{\Psi }_{,i}\right).`$
Notice that Eq. (30) is the non-linear generalization of the constraint between the gravitational potentials $`\mathrm{\Phi }`$ and $`\mathrm{\Psi }`$ in the longitudinal (Poisson) gauge and is valid at any order in perturbation theory. At linear order one recovers the well-known result $`\mathrm{\Phi }_1=\mathrm{\Psi }_1`$, while at second order the relation first found in Refs. enhanc ; review follows $`\psi _2\varphi _2=4\psi _1^2+10^4_i^k(\psi _1^{,i}\psi _{1,k})\frac{10}{3}^2(\psi _1^{,i}\psi _{1,i})`$ (where $`\mathrm{\Phi }_2=\varphi _22\varphi _1^2`$ and $`\mathrm{\Psi }_2=\psi _2+2\psi _1^2`$ and $`\varphi _1=\mathrm{\Phi }_1`$, $`\psi _1=\mathrm{\Psi }_1`$).
In the following we find it convenient to write the relation (30) as $`\mathrm{\Psi }=\mathrm{\Phi }+๐ฆ[\mathrm{\Phi },\mathrm{\Psi }]`$ where the kernel $`๐ฆ`$ is obtained by acting through the operator $`^4`$ onto the r.h.s. of Eq. (30) and takes into account non-local terms. Going to momentum space, one easily realizes that in the specific โsqueezedโ limit of Ref. maldacena , where one of the wavenumbers is much smaller than the other two, e.g. $`k_1k_{2,3}`$, the kernel $`๐ฆ0`$. From Eq. (9) the curvature perturbation is given by
$$\zeta \delta =\mathrm{\Psi }+\frac{1}{3}\mathrm{ln}\frac{\overline{\rho }}{\rho }=\mathrm{\Psi }\frac{2}{3}\mathrm{\Phi },$$
(31)
where we have used $`\overline{\rho }_m=\rho _me^{2\mathrm{\Phi }}`$. Hence we finally find $`5\mathrm{\Phi }/3=\zeta +๐ฆ`$ and we can recast the non-linear temperature anisotropies (19) as
$$\frac{\delta _{np}\overline{T}_๐ช}{T_๐ช}=e^{\zeta /5๐ฆ/5}1,$$
(32)
which is the starting point to evaluate the $`n`$point correlation function for CMB temperature anisotropies. The comoving curvature perturbation $`\zeta `$ is conserved at any order in perturbation theory on large scales and therefore one can fix its properties right at the end of inflation. Since for standard inflation the curvature perturbation can be considered as a Gaussian distributed quantity acquaviva ; maldacena (deviations from non-Gaussianity are proportional to deviations of the spectral index from unity and are, therefore, tiny), we will adopt $`\zeta `$ as our Gaussian seed in the case of standard single-field models of inflation.
## IV Correlation functions of large-scale CMB anisotropies
As a warm-up excercise, let us first consider a simpler case and evaluate the $`n`$-point correlation function of $`e^{\phi (๐ฑ)}`$ where $`\phi (๐ฑ)`$ is a Gaussian random field. Applying the well-known techniques of quantum field theory, it turns out that
$$e^{\phi (๐ฑ_1)}\mathrm{}..e^{\phi (๐ฑ_N)}=e^{\frac{1}{2}{\scriptscriptstyle ๐๐ฑ๐๐ฒJ(๐ฑ)\phi (๐ฑ)\phi (๐ฒ)J(๐ฒ)}},$$
(33)
where $`J(๐ฑ)=_{i=1}^N\delta (๐ฑ๐ฑ_i)`$ corresponds in fact to the source term appearing in the path integral formulation. The calculation of the integral in Eq. (33) brings
$$e^{\phi (๐ฑ_1)}\mathrm{}..e^{\phi (๐ฑ_N)}=e^{\frac{1}{2}_{i,j}\phi (๐ฑ_i)\phi (๐ฑ_j)}.$$
(34)
Notice, for example, that for the 2-point correlation function we find the usual result
$$e^{\phi (๐ฑ_1)}e^{\phi (๐ฑ_2)}=e^{\phi ^2}e^{\phi (๐ฑ_1)\phi (๐ฑ_2)},$$
(35)
which for $`๐ฑ_1=๐ฑ_2`$ gives $`e^{2\phi (๐ฑ)}=e^{2\phi ^2}`$, where $`\phi ^2=\phi ^2(๐ฑ)`$.
For the 3-point function one finds
$$e^{\phi (๐ฑ_1)}e^{\phi (๐ฑ_2)}e^{\phi (๐ฑ_3)}=e^{\frac{3}{2}\phi ^2}e^{\phi _1\phi _2+\phi _1\phi _3+\phi _2\phi _3},$$
(36)
where for simplicity we have used the notation $`\phi _i\phi _j=\phi (๐ฑ_i)\phi (๐ฑ_j)`$. If now we expand the exponential in the limit in which the two point function is small we obtain
$`e^{\phi (๐ฑ_1)}e^{\phi (๐ฑ_2)}e^{\phi (๐ฑ_3)}`$ $``$ $`e^{\frac{3}{2}\phi ^2}[1+(\phi _1\phi _2+\mathrm{cycl}.)`$ (37)
$`+`$ $`{\displaystyle \frac{1}{2}}(\phi _1\phi _2^2+\mathrm{cycl}.)`$
$`+`$ $`(\phi _1\phi _2\phi _1\phi _3+\mathrm{cycl}.)].`$
It is in fact a term analogous to the last contribution in Eq. (37) that enters in the three-point function of the CMB anisotropies (19) if we are interested in non-linearities up to second-order terms only. The complication arises from the fact that the gravitational potential $`\mathrm{\Phi }`$ appearing in Eq. (19) is not a Gaussian variable, since it can already contain quadratic (and higher order) terms in the Gaussian variable $`\zeta `$ in the case of single-field models of inflation. These will add non-Gaussian contributions which are contained in the kernel $`๐ฆ`$. In other scenarios for the generation of the cosmological perturbations $`\zeta `$ is not a Gaussian quantity and we will properly take into account the primordial non-Gaussian contributions.
Let us see how to generalize the previous procedure to this case. First of all we apply an iterative procedure to express the kernel $`๐ฆ`$ in terms of powers of $`\zeta `$. We will use Eqs. (30) and (31). For the zeroth and first-order terms of the iteration we find (the suffix does not refer to the order of the expansion in the perturbations, but to the order of the approximation given by the iteration procedure: each $`r`$-th term contains up to $`(r+1)`$ powers of $`\zeta `$)
$`\mathrm{\Phi }^{(0)}`$ $`=`$ $`{\displaystyle \frac{3}{5}}\zeta `$ (38)
$`\mathrm{\Phi }^{(1)}`$ $`=`$ $`{\displaystyle \frac{3}{5}}\zeta \left({\displaystyle \frac{3}{5}}\right)^3๐ฆ[\zeta ^2],`$ (39)
and for the next terms ($`n=1,2,..`$)
$`\mathrm{\Phi }^{(2n)}`$ $`=`$ $`\mathrm{\Phi }^{(2n1)}+๐ฆ_1[\mathrm{\Phi }^{(0)},\mathrm{\Phi }^{(2n2)}\mathrm{\Phi }^{(2n1)}]`$
$`+`$ $`{\displaystyle \underset{m=0}{\overset{n2}{}}}๐ฆ_1[\mathrm{\Phi }^{(m)}\mathrm{\Phi }^{(m+1)},\mathrm{\Phi }^{(2nm3)}\mathrm{\Phi }^{(2nm2)}],`$
$`\mathrm{\Phi }^{(2n+1)}`$ $`=`$ $`\mathrm{\Phi }^{(2n)}+๐ฆ_1[\mathrm{\Phi }^{(0)},\mathrm{\Phi }^{(2n1)}\mathrm{\Phi }^{(2n)}]`$ (41)
$`+`$ $`{\displaystyle \underset{m=0}{\overset{n2}{}}}๐ฆ_1[\mathrm{\Phi }^{(m)}\mathrm{\Phi }^{(m+1)},\mathrm{\Phi }^{(2nm2)}\mathrm{\Phi }^{(2nm1)}]`$
$`+`$ $`๐ฆ_2[(\mathrm{\Phi }^{(n1)}\mathrm{\Phi }^{(n)})^2],`$
where we have introduced the bilinear operators
$`๐ฆ_1[(),()]`$ $``$ $`^4(6_i^k[()^{,i}()_{,k}]2^2[()^{,i}()_{,i}]),`$
$`๐ฆ_2[(),()]`$ $``$ $`^4\left({\displaystyle \frac{1}{2}}_i^k\left[()^{,i}()_{,k}\right]{\displaystyle \frac{1}{2}}^2[()^{,i}()_{,i}]\right).`$ (42)
Notice that for equal entries $`๐ฆ_1=6๐ฆ/5`$. If the upper limit of the sums appearing in these expressions turns out to be negative the sum must be taken to be vanishing.
Thus we can use Eq. (38) and (39) to find the next order approximations for the expression of $`\mathrm{\Phi }`$ in terms of the Gaussian curvature perturbation $`\zeta `$. For example for $`\mathrm{\Phi }^{(2)}`$ up to $`๐ช(\zeta ^3)`$ contributions we have
$$\mathrm{\Phi }^{(2)}=\frac{3}{5}\zeta \left(\frac{3}{5}\right)^3๐ฆ[\zeta ^2]+๐ฆ_1[\frac{3}{5}\zeta ,\left(\frac{3}{5}\right)^3๐ฆ[\zeta ^2]].$$
(43)
Using this iterative procedure we express the kernel $`๐ฆ`$ in Eq. (32) as a function of $`\zeta `$ and we can determine the $`n`$-point connected correlation functions for the temperature anisotropies (32) accounting for the information contained in $`๐ฆ[\zeta ]`$. In the following we outline the procedure and we give the results, while more details on the computation can be found in the Appendix.
In order to obtain the $`n`$-point (connected) correlation functions we will borrow again some techniques of functional-integral analysis from quantum field theory Ramond ; Zinn ; Mosel (for different applications of the path-integral approach in cosmology, see Refs. PW ; GW ; MLB ; Bertschinger ; MVJ ; MHV ). In particular we need the so called generating functional of the correlation functions $`W`$, which in our case can be written as
$$Z[J]=๐[\zeta ]๐ซ[\zeta ]e^{i{\scriptscriptstyle ๐๐ฑJ(๐ฑ)(e^{\zeta /5๐ฆ\left(\zeta \right)/5}1)}}.$$
(44)
Here $`J(๐ฑ)`$ is an external source perturbing the underlying statistics, $`๐ซ[\zeta ]`$ is the probability density functional which in our case is a Gaussian one (see Eq. (A.3) of the Appendix), and $`๐[\zeta ]`$ is a suitable measure such that the total probability is normalized to unity, $`๐[\zeta ]๐ซ[\zeta ]=1`$. The kernel $`๐ฆ[\zeta ]`$ will be given by the iterative procedure. From the generating functional one can obtain the correlation functions by taking the functional derivatives of $`Z[J]`$ with respect to the source $`J`$ evaluated at $`J=0`$ Zinn
$`(e^{\zeta _1/5๐ฆ(\zeta _1)/5}1)\mathrm{}(e^{\zeta _n/5๐ฆ(\zeta _n)/5}1)=`$
$`i^n{\displaystyle \frac{\delta ^nZ[J]}{\delta J(๐ฑ_1)\mathrm{}\delta J(๐ฑ_n)}}|_{(J=0)}Z^{(n)}(๐ฑ_1,\mathrm{},๐ฑ_n),`$ (45)
where $`\zeta _i\zeta (๐ฑ_i)`$.
By defining the new functional $`W[J]=\mathrm{ln}Z[J]`$ one can obtain the connected correlation functions by the same operations
$`(e^{\zeta _1/5๐ฆ(\zeta _1)/5}1)\mathrm{}(e^{\zeta _n/5๐ฆ(\zeta _n)/5}1)_{\mathrm{conn}.}=`$
$`i^n{\displaystyle \frac{\delta ^nW[J]}{\delta J(๐ฑ_1)\mathrm{}\delta J(๐ฑ_n)}}|_{(J=0)}W^{(n)}(๐ฑ_1,\mathrm{},๐ฑ_n).`$ (46)
It is clear that the complete knowledge of the statistical properties of the perturbations, i.e. the complete knowledge of the correlation functions at all orders, can be achieved if one knows the generating functionals. In turn from the definitions (IV) and (IV) the correlation functions will appear in a power series expansion of the generating functionals, as
$`Z[J]`$ $`=`$ $`1+{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{i^n}{n!}}{\displaystyle ๐๐ฑ_1\mathrm{}๐๐ฑ_nZ^{(n)}(๐ฑ_1,\mathrm{},๐ฑ_n)J(๐ฑ_1)\mathrm{}J(๐ฑ_n)},`$
$`W[J]`$ $`=`$ $`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{i^n}{n!}}{\displaystyle ๐๐ฑ_1\mathrm{}๐๐ฑ_nW^{(n)}(๐ฑ_1,\mathrm{},๐ฑ_n)J(๐ฑ_1)\mathrm{}J(๐ฑ_n)}.`$ (47)
To compute the functional derivatives of $`W[J]`$ one can follow the standard procedure to evaluate the connected correlation functions used in field theory Ramond . The computation makes use of a perturbative expansion around some known solution, which corresponds to the connected correlation functions of the free scalar field. In our case, the known solution corresponds to the correlation functions when the kernel $`๐ฆ`$ vanishes, which we have computed previously in Eq. (34). As explained in detail in the Appendix, we perform the expansion around $`๐ฆ=0`$ using as expansion parameter the r.m.s amplitude of the cosmological perturbations themselves, that is $`\zeta ^2^{1/2}10^5`$.
### IV.1 The Bispectrum
From the general results presented in the Appendix we provide now the expression for the 3-point connected correlation function. It suffices to expand the kernel $`๐ฆ`$ up to second order. The total kernel can be written as a convolution in configuration space
$$๐ฆ(\zeta )=๐๐ฑ_1๐๐ฑ_2K_2(๐ฑ๐ฑ_1,๐ฑ๐ฑ_2)\zeta (๐ฑ_1)\zeta (๐ฑ_2),$$
(48)
where $`K(๐ฑ๐ฑ_1,๐ฑ๐ฑ_2)`$ is the double inverse Fourier transform of the expression
$$\stackrel{~}{K}_2(๐ค_1,๐ค_2)=\left(a_{\mathrm{NL}}1\right)+\frac{9}{5}\left[(๐ค_1๐ค_3)(๐ค_2๐ค_3)/k^4(1/3)(๐ค_1๐ค_2)/k^2\right],$$
(49)
where $`k=\left|๐ค_3\right|`$ and $`๐ค_3=(๐ค_1+๐ค_2)`$. In Eq. (48) we have added the constant $`a_{\mathrm{NL}}`$ whose role is to parametrize the primordial non-Gaussianity generated during or after inflation in the various scenarios for the generation of the cosmological perturbations, $`\zeta =\zeta _\mathrm{L}+(a_{\mathrm{NL}}1)\zeta _\mathrm{L}^2`$ (from now on we will remove the subscript โLโ). For instance, in single field models of inflation $`a_{\mathrm{NL}}=1`$ (plus tiny contributions proportional to the deviation from scale invariance), in the curvaton scenarion $`a_{\mathrm{NL}}=(3/4r)r/2`$, where $`r(\rho _\sigma /\rho )_D`$ represents the relative curvaton contribution to the total energy density at curvaton decay review .
We find
$`W^{(3)}(๐ฑ_1,๐ฑ_2,๐ฑ_3)`$ $`=`$ $`(e^{\zeta (๐ฑ_1)/5}1)(e^{\zeta (๐ฑ_2)/5}1)(e^{\zeta (๐ฑ_3)/5}1)_{\mathrm{connected}}`$
$``$ $`5{\displaystyle \underset{p}{}}{\displaystyle }d๐ฒ_1d๐ฒ_2K_2(๐ฑ_{p_2}๐ฒ_1,๐ฑ_{p_2}๐ฒ_2)[\stackrel{~}{w}_2(๐ฑ_{p_1},๐ฒ_1)`$
$`\times `$ $`\stackrel{~}{w}_2(๐ฑ_{p_3},๐ฒ_2)+{\displaystyle \frac{1}{2}}\stackrel{~}{w}_4(๐ฑ_{p_1},๐ฑ_{p_3},๐ฒ_1,๐ฒ_2)],`$
where the sum is over all permutations $`p_1,p_2,p_3`$ taking the values $`(1,2,3)`$ and
$`\stackrel{~}{w}_2(๐ฑ_1,๐ฑ_2)`$ $``$ $`{\displaystyle \frac{1}{5}}(e^{\zeta _1/5}1)\zeta _2)_{\mathrm{conn}.}={\displaystyle \frac{1}{5^2}}e^{\zeta ^2/50}\zeta _1\zeta _2,`$ (51)
$`\stackrel{~}{w}_4(๐ฑ_1,๐ฑ_2,๐ฑ_3,๐ฑ_4)`$ $``$ $`{\displaystyle \frac{1}{25}}(e^{\zeta _1/5}1)(e^{\zeta _2/5}1)\zeta _3\zeta _4_{\mathrm{conn}.}`$
$`=`$ $`{\displaystyle \frac{e^{\zeta ^2/25}}{3\times 5^4}}\left(e^{\zeta _1\zeta _2/25}1\right)\left(\zeta _1\zeta _4+\zeta _2\zeta _4\right)\left(\zeta _1\zeta _3+\zeta _2\zeta _3\right)+\mathrm{cyclic}.`$
Despite the fact that we expanded the kernel up to second-order, the expression (IV.1) becomes exact at any order in perturbation theory in the squezeed limit for which $`๐ฆ`$ tends to zero and for single field models of inflation for which $`a_{\mathrm{NL}}=1`$. In such a case, the exact three-point correlation function for the temperature anisotropies on large-scales is
$`W^{(3)}(๐ฑ_1,๐ฑ_2,๐ฑ_3)`$ $`=`$ $`W_0^{(2)}(๐ฑ_1,๐ฑ_2)W_0^{(2)}(๐ฑ_1,๐ฑ_3)+W_0^{(2)}(๐ฑ_1,๐ฑ_3)W_0^{(2)}(๐ฑ_2,๐ฑ_3)`$
$`+`$ $`W_0^{(2)}(๐ฑ_1,๐ฑ_2)W_0^{(2)}(๐ฑ_2,๐ฑ_3)+W_0^{(2)}(๐ฑ_1,๐ฑ_2)W_0^{(2)}(๐ฑ_2,๐ฑ_3)W_0^{(2)}(๐ฑ_3,๐ฑ_1),`$
where
$`W_0^{(2)}(๐ฑ_i,๐ฑ_j)`$ $``$ $`e^{\zeta ^2/50}\left(e^{\zeta _i\zeta _j/50}1\right){\displaystyle \frac{d^3k}{(2\pi )^3}e^{i๐ค(๐ฑ_i๐ฑ_j)}P(k)}`$ (54)
$``$ $`{\displaystyle \frac{1}{50}}{\displaystyle \frac{d\mathrm{ln}k}{2\pi ^2}j_0\left(k\left|๐ฑ_i๐ฑ_j\right|\right)๐ซ_\zeta (k)},`$
and $`๐ซ_\zeta =A(k_0)^2(k/k_0)^{n_S1}`$ is the primordial power spectrum of the comoving curvature perturbation with amplitude $`A`$ and spectral index $`n_S`$. The expression for the exact bispectrum of temperature anisotropies valid at any order in perturbation theory is
$$<\frac{\delta _{np}T(๐ค_1)}{T}\frac{\delta _{np}T(๐ค_2)}{T}\frac{\delta _{np}T(๐ค_3)}{T}>=(2\pi )^3\delta ^{(3)}\left(๐ค_1+๐ค_2+๐ค_3\right)B(๐ค_1,๐ค_2,๐ค_3),$$
(55)
where
$`B(๐ค_1,๐ค_2,๐ค_3)`$ $`=`$ $`P(k_1)P(k_2)+P(k_1)P(k_3)+P(k_2)P(k_3)`$
$`+`$ $`{\displaystyle \frac{d๐ช}{(2\pi )^3}P\left(\left|๐ช๐ค_1\right|\right)P\left(\left|๐ช๐ค_2\right|\right)P\left(\left|๐ช๐ค_3\right|\right)}`$
$``$ $`2P(k_1)P(k_2)+{\displaystyle \frac{d๐ช}{(2\pi )^3}P\left(\left|๐ช\right|\right)P\left(\left|๐ช๐ค_2\right|\right)P\left(\left|๐ช+๐ค_2\right|\right)},(k_1k_2,k_3).`$
For generic momenta configurations and for models for which $`a_{\mathrm{NL}}`$ is sizeable, the exponentials present in the first line of Eq. (IV.1) has to be expanded at second-order to consistently match the order of the kernel $`K`$
$$e^{\zeta /5}1\frac{1}{5}\zeta +\frac{1}{2\times 5^2}\zeta ^2+๐ช\left(\zeta ^3\right).$$
(57)
We immediately recover the expression obtained at second-order in perturbation theory for the non-linearity parameter $`f_{\mathrm{NL}}`$ defined as the coefficient of the gravitational potential $`\mathrm{\Phi }`$ expanded at second-order in terms of the linear Gaussian gravitational potential $`\mathrm{\Phi }_\mathrm{L}=\varphi _1`$, $`\mathrm{\Phi }=\mathrm{\Phi }_\mathrm{L}+f_{\mathrm{NL}}\left(\mathrm{\Phi }_\mathrm{L}\right)^2`$, with the convential Sachs-Wolfe effect expressed as $`\delta T/T=(\mathrm{\Phi }/3)`$ prl
$$f_{\mathrm{NL}}=\left[\frac{5}{3}(1a_{\mathrm{NL}})+\frac{1}{6}\right]+\left[3(๐ค_1๐ค_3)(๐ค_2๐ค_3)/k^4(๐ค_1๐ค_2)/k^2\right].$$
(58)
It may be worth noticing that the coefficient 1/6 in the last expression is simply the result of the expansion of the exponential (57) expressed in terms of the gravitational potential $`\varphi _1`$: $`e^{\zeta /5}=e^{\mathrm{\Phi }_\mathrm{L}/3}1\frac{1}{3}\mathrm{\Phi }_\mathrm{L}+\frac{1}{18}\mathrm{\Phi }_\mathrm{L}^2`$, which gives the contribution $`3\times (1/18)=1/6`$ to $`f_{\mathrm{NL}}`$.
### IV.2 The Trispectrum
Let us now follow a similar procedure to obtain the 4-point connected correlation function. In this case the kernel $`๐ฆ`$ appearing in Eq. (32) must be expanded up to third order. In configuration space it can be written as a convolution
$`๐ฆ(\zeta )`$ $`=`$ $`{\displaystyle ๐๐ฑ_1๐๐ฑ_2K_2(๐ฑ๐ฑ_1,๐ฑ๐ฑ_2)\zeta (๐ฑ_1)\zeta (๐ฑ_2)}`$ (59)
$`+`$ $`{\displaystyle ๐๐ฑ_1๐๐ฑ_2๐๐ฑ_3K_3(๐ฑ๐ฑ_1,๐ฑ๐ฑ_2,๐ฑ๐ฑ_3)\zeta (๐ฑ_1)\zeta (๐ฑ_2)\zeta (๐ฑ_3)},`$
where $`K_2`$ is the kernel defined by Eqs. (48) and Eq. (49), while $`K_3(๐ฑ๐ฑ_1,๐ฑ๐ฑ_2,๐ฑ๐ฑ_3)`$ is the triple inverse Fourier transform of the expression
$`\stackrel{~}{K_3}(๐ค_1,๐ค_2,๐ค_3)=(b_{\mathrm{NL}}1)+(a_{\mathrm{NL}}1)๐(๐ค_1,๐ค_2,๐ค_3)+๐(๐ค_1,๐ค_2,๐ค_3),`$ (60)
with
$`๐(๐ค_1,๐ค_2,๐ค_3)`$ $`=`$ $`{\displaystyle \frac{6}{5}}[{\displaystyle \frac{๐ค_1๐ค_4(๐ค_3๐ค_4+๐ค_2๐ค_4)+(๐ค_2๐ค_4)(๐ค_3๐ค_4)}{k^4}}`$ (61)
$``$ $`{\displaystyle \frac{1}{3}}{\displaystyle \frac{๐ค_1(๐ค_2+๐ค_3)+๐ค_2๐ค_3}{k^2}}],`$
$`๐(๐ค_1,๐ค_2,๐ค_3)`$ $`=`$ $`{\displaystyle \frac{54}{25}}{\displaystyle \frac{(๐ค_4๐ค_3)[(๐ค_1+๐ค_2)๐ค_4]}{k^4}}[{\displaystyle \frac{(๐ค_1(๐ค_1+๐ค_2))(๐ค_2(๐ค_1+๐ค_2))}{|๐ค_1+๐ค_2|^4}}`$ (62)
$``$ $`{\displaystyle \frac{1}{3}}{\displaystyle \frac{๐ค_1๐ค_2}{|๐ค_1+๐ค_2|^2}}]+\mathrm{cycl}.,`$
where $`k=|๐ค_4|`$ and $`๐ค_4=(๐ค_1+๐ค_2+๐ค_3)`$. In Eq. (62) one has to take cyclic terms by an exchange of the wavenumbers $`๐ค_1,๐ค_2,๐ค_3`$.
In order to compute $`\stackrel{~}{K_3}(๐ค_1,๐ค_2,๐ค_3)`$ we have applied the iterative procedure described in Section IV taking into account also a possible primordial non-Gaussian contribution by expanding the curvature perturbation as
$$\zeta =\zeta _\mathrm{L}+(a_{\mathrm{NL}}1)\zeta _\mathrm{L}^2+(b_{\mathrm{NL}}1)\zeta _\mathrm{L}^3.$$
(63)
The value of $`b_{\mathrm{NL}}`$ will depend on the different scenarios for the generation of the cosmological perturbations. For example, for standard single-field models of inflation $`b_{\mathrm{NL}}=1`$ (plus tiny contributions proportional to powers of the slow-roll parameters), while for other scenarios it might well be non-negligible. For simplicity from now on we will remove the subscript โLโ.
Similarly to the bispectrum, also for the 4-point connected correlation function there exists a specific limit for which the kernel $`๐ฆ`$ tends to zero, corresponding to take two wavenumbers much smaller than the other ones, e.g. $`k_1,k_2k_3,k_4`$. For this limit and in the case of single-field models of inflation ($`a_{\mathrm{NL}}=1,b_{\mathrm{NL}}=1`$), one can compute an exact expression of the trispectrum by Fourier transforming the 4-point connected correlation function which, in this limit, reads
$`W^{(4)}(๐ฑ_1,๐ฑ_2,๐ฑ_3,๐ฑ_4)=(e^{\zeta (๐ฑ_1)/5}1)(e^{\zeta (๐ฑ_2)/5}1)(e^{\zeta (๐ฑ_3)/5}1)(e^{\zeta (๐ฑ_4)/5}1)_{\mathrm{conn}.}`$
$`=W_0^{(2)}(๐ฑ_1,๐ฑ_2)W_0^{(2)}(๐ฑ_1,๐ฑ_3)W_0^{(2)}(๐ฑ_1,๐ฑ_4)W_0^{(2)}(๐ฑ_2,๐ฑ_3)W_0^{(2)}(๐ฑ_2,๐ฑ_4)W_0^{(2)}(๐ฑ_3,๐ฑ_4)`$
$`+W_0^{(2)}(๐ฑ_1,๐ฑ_4)W_0^{(2)}(๐ฑ_2,๐ฑ_3)W_0^{(2)}(๐ฑ_2,๐ฑ_4)W_0^{(2)}(๐ฑ_3,๐ฑ_4)\left(W_0^{(2)}(๐ฑ_1,๐ฑ_2)+W_0^{(2)}(๐ฑ_1,๐ฑ_3)\right)`$
$`+W_0^{(2)}(๐ฑ_1,๐ฑ_2)W_0^{(2)}(๐ฑ_1,๐ฑ_3)W_0^{(2)}(๐ฑ_2,๐ฑ_4)W_0^{(2)}(๐ฑ_3,๐ฑ_4)\left(W_0^{(2)}(๐ฑ_1,๐ฑ_4)+W_0^{(2)}(๐ฑ_2,๐ฑ_3)\right)`$
$`+W_0^{(2)}(๐ฑ_2,๐ฑ_3)W_0^{(2)}(๐ฑ_2,๐ฑ_4)W_0^{(2)}(๐ฑ_3,๐ฑ_4)\left(W_0^{(2)}(๐ฑ_1,๐ฑ_2)+W_0^{(2)}(๐ฑ_1,๐ฑ_3)+W_0^{(2)}(๐ฑ_1,๐ฑ_4)\right)`$
$`+W_0^{(2)}(๐ฑ_1,๐ฑ_4)W_0^{(2)}(๐ฑ_2,๐ฑ_4)W_0^{(2)}(๐ฑ_3,๐ฑ_4)\left(W_0^{(2)}(๐ฑ_1,๐ฑ_2)+W_0^{(2)}(๐ฑ_1,๐ฑ_3)\right)`$
$`+(W_0^{(2)}(๐ฑ_2,๐ฑ_4)+W_0^{(2)}(๐ฑ_3,๐ฑ_4))(W_0^{(2)}(๐ฑ_1,๐ฑ_2)W_0^{(2)}(๐ฑ_1,๐ฑ_4)W_0^{(2)}(๐ฑ_2,๐ฑ_3)+`$
$`+W_0^{(2)}(๐ฑ_2,๐ฑ_3)W_0^{(2)}(๐ฑ_1,๐ฑ_4)W_0^{(2)}(๐ฑ_1,๐ฑ_3)+W_0^{(2)}(๐ฑ_1,๐ฑ_3)W_0^{(2)}(๐ฑ_2,๐ฑ_3)W_0^{(2)}(๐ฑ_1,๐ฑ_2))`$
$`+W_0^{(2)}(๐ฑ_1,๐ฑ_3)W_0^{(2)}(๐ฑ_1,๐ฑ_4)W_0^{(2)}(๐ฑ_1,๐ฑ_2)\left(W_0^{(2)}(๐ฑ_2,๐ฑ_3)+W_0^{(2)}(๐ฑ_2,๐ฑ_4)+W_0^{(2)}(๐ฑ_3,๐ฑ_4)\right)`$
$`+W_0^{(2)}(๐ฑ_1,๐ฑ_2)W_0^{(2)}(๐ฑ_1,๐ฑ_3)W_0^{(2)}(๐ฑ_2,๐ฑ_4)W_0^{(2)}(๐ฑ_3,๐ฑ_4)`$
$`+W_0^{(2)}(๐ฑ_1,๐ฑ_2)W_0^{(2)}(๐ฑ_1,๐ฑ_3)\left(W_0^{(2)}(๐ฑ_1,๐ฑ_4)+W_0^{(2)}(๐ฑ_2,๐ฑ_4)+W_0^{(2)}(๐ฑ_3,๐ฑ_4)\right)`$
$`+W_0^{(2)}(๐ฑ_1,๐ฑ_4)W_0^{(2)}(๐ฑ_2,๐ฑ_3)\left(W_0^{(2)}(๐ฑ_1,๐ฑ_3)+W_0^{(2)}(๐ฑ_2,๐ฑ_4)+W_0^{(2)}(๐ฑ_3,๐ฑ_4)\right)`$
$`+W_0^{(2)}(๐ฑ_2,๐ฑ_4)W_0^{(2)}(๐ฑ_3,๐ฑ_4)\left(W_0^{(2)}(๐ฑ_1,๐ฑ_2)+W_0^{(2)}(๐ฑ_1,๐ฑ_3)+W_0^{(2)}(๐ฑ_1,๐ฑ_4)\right)`$
$`+\left(W_0^{(2)}(๐ฑ_1,๐ฑ_2)+W_0^{(2)}(๐ฑ_1,๐ฑ_3)\right)\left(W_0^{(2)}(๐ฑ_2,๐ฑ_3)W_0^{(2)}(๐ฑ_2,๐ฑ_4)+W_0^{(2)}(๐ฑ_2,๐ฑ_3)W_0^{(2)}(๐ฑ_3,๐ฑ_4)\right)`$
$`+W_0^{(2)}(๐ฑ_1,๐ฑ_2)W_0^{(2)}(๐ฑ_1,๐ฑ_4)W_0^{(2)}(๐ฑ_3,๐ฑ_4)+W_0^{(2)}(๐ฑ_1,๐ฑ_3)W_0^{(2)}(๐ฑ_1,๐ฑ_4)W_0^{(2)}(๐ฑ_2,๐ฑ_4).`$ (64)
On the other hand, for generic momenta configurations and for models for which $`a_{\mathrm{NL}}`$ and $`b_{\mathrm{NL}}`$ are sizeable we can expand the exponential entering in the expression (32) of the temperature anisotropies up to third order and use the kernel $`๐ฆ(\zeta )`$ in Eq. (59). In this way we are able to determine the non-linearity parameter $`g_{\mathrm{NL}}`$ which enters into the trispectrum of the CMB anisotropies according, for example, to the analysis of Refs. hu ; jul . The parameter $`g_{\mathrm{NL}}`$ is defined through the expansion of the (Bardeen) gravitational potential $`\mathrm{\Phi }`$ up to third-order as
$$\mathrm{\Phi }=\mathrm{\Phi }_L+f_{\mathrm{NL}}(\mathrm{\Phi }_L)^2+g_{\mathrm{NL}}(\mathrm{\Phi }_L)^3,$$
(65)
where $`\mathrm{\Phi }_L=\varphi _1`$, is the linear Gaussian part of $`\mathrm{\Phi }`$. We find the following expression
$`g_{\mathrm{NL}}(๐ค_\mathrm{๐},๐ค_2,๐ค_\mathrm{๐})`$ $`=`$ $`{\displaystyle \frac{25}{9}}(b_{\mathrm{NL}}1)+{\displaystyle \frac{25}{9}}(a_{\mathrm{NL}}1)๐(๐ค_1,๐ค_2,๐ค_3)+{\displaystyle \frac{25}{9}}๐(๐ค_1,๐ค_2,๐ค_3){\displaystyle \frac{5}{9}}(a_{\mathrm{NL}}1)`$ (66)
$`+`$ $`{\displaystyle \frac{1}{54}}{\displaystyle \frac{1}{3}}[{\displaystyle \frac{(๐ค_1(๐ค_1+๐ค_2))(๐ค_2(๐ค_1+๐ค_2))}{|๐ค_1+๐ค_2|^4}}{\displaystyle \frac{1}{3}}{\displaystyle \frac{๐ค_1๐ค_2}{|๐ค_1+๐ค_2|^2}}+\mathrm{cycl}.],`$
where $`๐(๐ค_1,๐ค_2,๐ค_3)`$ and $`๐(๐ค_1,๐ค_2,๐ค_3)`$ are defined through Eqs. (61) and (62).
## V Conclusions
In this paper we showed how to calculate exactly the $`n`$-point correlation function of CMB anisotropies in the case in which all wavelengths are beyond the horizon at last scattering. In this limit the Sachs-Wolfe effect is predominant and its contribution to higher-order correlation functions yields the most direct signal of non-Gaussianity in the primordial cosmological seeds. Of particular interest are the bispectrum and the trispectrum which can be used to assess the level of non-Gaussianity on cosmological scales. We have calculated the non-perturbative expressions for the bispectrum and the trispecturm as predicted within single-field models of inflation and in the so-called โsqueezedโ limit in which some of the wavenumbers are much smaller than the others. For other scenarios of generation of the cosmological perturbations, we have provided the non-linearity parameters $`f_{\mathrm{NL}}`$ and $`g_{\mathrm{NL}}`$ entering respectively the theoretical predictions of the bispectrum and trispectrum. Our results for the bispectrum and the trispectrum represent the essential input in order to obtain the predicted angular bispectrum $`B_{l_1l_2l_3}=a_{l_1m_1}a_{l_2m_2}a_{l_3m_3}`$ and the trispectrum of CMB anisotropies (and higher order correlation functions) according, for example, to the formalism developed in Refs. SG ; wk . In these works it is shown how to compute the angular bispectrum accounting for a non-trivial wavenumber dependence of the non-linearity parameter $`f_{\mathrm{NL}}(๐ค_1,๐ค_2)`$. The angular modulation of the quadratic non-linearity predicted by Eq. (58) is currently under investigation LMR , adopting the technique of Ref. SG , in order to look for specific signatures of inflationary non-Gaussianity in the CMB. Notice that for a $`\mathrm{\Lambda }`$CDM cosmology a late integrated Sachs-Wolfe effect arising from the explicit time dependence of the linear gravitational potential during the late accelerated phase would also give a contribution on large scales. The formalism developed in this paper can be extended to take into account also for this effect. On the other hand on smaller scales there will be other effects contributing to CMB non-Gaussianity such as gravitational lensing, Shapiro time-delays and Rees-Sciama effects produced at the non-linear level. One should be able to compute the angular connected correlation functions induced by these effects by using the techinique developed in Ref. CZ and to distinguish them from the large-scale Sachs-Wolfe effect provided here thanks to their specific angular dependence. Our predictions for the higher-order correlation functions should be compared model by model with the unavoidable contributions from various secondary anisotropies and systematic effects, such as astrophysical foregrounds.
## VI Appendix
In this Appendix we will show in detail how to compute the $`n`$-point connected correlation functions (IV.1) by making use of the generating functional $`w(๐ฑ_\mathrm{๐},\mathrm{},๐ฑ_๐ง)`$ and in particular how to get the result (IV.1).
Let us consider the quantity
$$\delta _T=e^{\phi +K(\phi )}1,$$
(A.1)
where $`\phi `$ is a generic Gaussian random field, and $`K(\phi )`$ is a generic functional of $`\phi `$ (apart from multiplicative coefficients $`\phi `$ will be identified with the comoving curvature perturbation $`\zeta `$ and $`K`$ will be given by the kernel $`๐ฆ`$ through the iterative procedure).
The generating functional for the correlated functions of $`\delta _T`$ is given by
$$Z[J]=๐[\phi ]๐ซ[\phi ]e^{i{\scriptscriptstyle ๐๐ฑJ(๐ฑ)(e^{\phi +K\left(\phi \right)}1)}},$$
(A.2)
where $`J(๐ฑ)`$ is an arbitrary external source. The functional integral is over all the $`\phi `$ configurations weighted by the Gaussian probability density functional
$$๐ซ[\phi ]=\frac{e^{\frac{1}{2}{\scriptscriptstyle ๐๐ฒ๐๐ฑ\phi (๐ฒ)๐ข(๐ฒ,๐ฑ)\phi (๐ฑ)}}}{๐[\phi ]e^{\frac{1}{2}{\scriptscriptstyle ๐๐ฒ๐๐ฑ\phi (๐ฒ)๐ข(๐ฒ,๐ฑ)\phi (๐ฑ)}}},$$
(A.3)
which has been properly normalized in such a way that the total probability equals unity, $`๐[\phi ]๐ซ[\phi ]=1`$.
To compute the functional derivatives with respect to $`J`$ we find it convenient to use an additional arbitrary source $`\lambda (๐ฑ)`$. We introduce the following generating functional
$$Z[J,\lambda ]=๐[\phi ]๐ซ[\phi ]e^{i{\scriptscriptstyle ๐๐ฑJ(๐ฑ)(e^{\phi +K\left(\phi \right)}1)}}e^{i{\scriptscriptstyle ๐๐ฑ\lambda (๐ฑ)\phi (๐ฑ)}},$$
(A.4)
which reduces to the expression (A.2) when $`\lambda =0`$. Functionals of the form in Eq. (A.4) are common in field theory when computing correlation functions of composite operators (in which case $`J(๐ฑ)`$ represents a โlocal couplingโ; see for example Ref. Zinn ). The correlation functions generated by $`Z[J,\lambda ]`$ are given by
$`(e^{\phi (๐ฒ_1)+K(\phi (๐ฒ_1))}1)\mathrm{}(e^{\phi (๐ฒ_n)+K(\phi (๐ฒ_n))}1)\phi (๐ฑ_1)\mathrm{}\phi (๐ฑ_m)`$
$`=i^{n+m}{\displaystyle \frac{\delta ^{n+m}Z[J,\lambda ]}{\delta J(๐ฒ_1)\mathrm{}\delta J(๐ฒ_n)\delta \lambda (๐ฑ_1)\mathrm{}.\delta \lambda (๐ฑ_m)}}|_{(J,\lambda )=0}.`$
Thus we will take only derivatives with respect to $`J`$ to get the correlation functions of $`e^{\phi +K(\phi )}`$. If we now write
$$e^{\phi +K(\phi )}1=\underset{n=1}{\overset{\mathrm{}}{}}\frac{K^n(\phi )}{n!}e^\phi +(e^\phi 1),$$
(A.6)
we can rewrite Eq. (A.4) as
$$Z[J,\lambda ]=e^{i{\scriptscriptstyle ๐๐ฑJ(๐ฑ)\left(_{n=1}{\scriptscriptstyle \frac{K^n(\frac{1}{i}\frac{\delta }{\delta \lambda })}{n!}}e^{{\scriptscriptstyle \frac{1}{i}}{\scriptscriptstyle \frac{\delta }{\delta \lambda }}}\right)}}Z_0[J,\lambda ],$$
(A.7)
where $`Z_0[J,\lambda ]`$ corresponds to the generating functional $`Z[J,\lambda ]`$ when the function $`K(\phi )=0`$
$`Z_0[J,\lambda ]`$ $`=`$ $`{\displaystyle ๐[\phi ]๐ซ[\phi ]e^{i{\scriptscriptstyle ๐๐ฑJ(๐ฑ)(e^\phi 1)}}e^{i{\scriptscriptstyle ๐๐ฑ\lambda (๐ฑ)\phi (๐ฑ)}}}`$ (A.8)
$``$ $`e^{W_0[J,\lambda ]}.`$
In writing Eq. (A.7) we have made use of the property $`i^1\delta e^{i{\scriptscriptstyle ๐๐ฒ\lambda (๐ฒ)\phi (๐ฑ)}}/\delta J(๐ฒ)=\phi (๐ฑ)e^{i{\scriptscriptstyle ๐๐ฒ\lambda (๐ฒ)\phi (๐ฒ)}}`$, in order to isolate the โinteraction termโ (see, for example, Ref. Ramond ).
Now let us write
$$W[J,\lambda ]\mathrm{ln}Z[J,\lambda ]=W_0[J,\lambda ]+\mathrm{ln}\left[1+e^{W_0}\left(e^{i{\scriptscriptstyle ๐๐ฑJ(๐ฑ)\left(_{n=1}{\scriptscriptstyle \frac{K^n(\frac{1}{i}\frac{\delta }{\delta \lambda })}{n!}}e^{{\scriptscriptstyle \frac{1}{i}}{\scriptscriptstyle \frac{\delta }{\delta \lambda }}}\right)}}1\right)e^{W_0}\right].$$
The derivatives with respect to $`J`$ evaluated for $`(J,\lambda )=0`$ will give the connected correlation functions we are looking for, accounting also for the kernel $`K`$. We are using the standard procedure to evaluate the connected correlation functions, for example, for an interacting scalar field (see e.g. Ref. Ramond ), except that in our case the โinteraction termโ is related to the kernel $`K`$ and we will make a perturbative expansion around $`W_0[J,\lambda ]`$, since the derivatives of $`W_0[J,\lambda ]`$ with respect to $`J`$ (evaluated for $`(J,\lambda )=0`$) give the connected correlated functions for $`\left(e^\phi 1\right)`$, see Eq. (33).
At this point we have to perform a perturbative expansion. We suppose that the perturbation parameter is the small r.m.s amplitude of the perturbations, $`\phi _{\mathrm{rms}}1`$ and use the fact that the function $`K(\phi )`$ is obtained through the iterative procedure described in Section III. Let us suppose that
$$K(\phi )=a\phi ^2+b\phi ^3+\mathrm{}.,$$
(A.9)
where the star denotes a convolution operation in configuration space. For the sake of simplicity, in the following we will neglect the convolutions, and treat $`a,b,..`$ as constant coefficients. In Eq. (VI) we will focus on the term $`\mathrm{ln}(1+\mathrm{\Delta })`$ where the definition of $`\mathrm{\Delta }`$ is obtained by comparison with Eq. (VI). Thus $`\mathrm{ln}(1+\mathrm{\Delta })=1+\mathrm{\Delta }\mathrm{\Delta }^2/2+\mathrm{\Delta }^3/3+\mathrm{}`$.
Now let us consider how to compute, for example, the three-point correlated function given in Eq. (IV.1). At lowest order in our approximation, the only term that in this case is relevant is just $`\mathrm{\Delta }`$ and moreover we keep only the first term coming from the expansion of the exponential
$$\mathrm{\Delta }e^{W_0}i๐๐ฑJ(๐ฑ)\left(\underset{n=1}{}\frac{K^n(\frac{1}{i}\frac{\delta }{\delta \lambda })}{n!}e^{\frac{1}{i}\frac{\delta }{\delta \lambda }}\right)e^{W_0}.$$
(A.10)
Here we have to expand once more, and looking at Eq. (A.9) we just need for the bispectrum to take $`K=a\phi ^2`$, and from $`(e^{\frac{1}{i}\frac{\delta }{\delta \lambda }}=1+\delta /(i\delta \lambda )+\mathrm{}.)`$ we just pick up the factor 1..
Thus we are left with
$`\mathrm{ln}(1+\mathrm{\Delta })`$ $``$ $`\mathrm{\Delta }iae^{W_0}{\displaystyle ๐๐ฑJ(๐ฑ)\left(\frac{1}{i}\frac{\delta }{\delta \lambda }\right)^2e^{W_0}}`$
$`=`$ $`i^1a{\displaystyle ๐๐ฑJ(๐ฑ)\left(\frac{\delta W_0}{\delta \lambda (๐ฑ)}\right)^2}+{\displaystyle \frac{\delta ^2W_0}{\delta \lambda ^2(๐ฑ)}}.`$
Therefore the bispectrum is given by
$`W^{(3)}(๐ฒ_1,๐ฒ_2,๐ฒ_3)=i^3{\displaystyle \frac{\delta ^3W[J,\lambda ]}{\delta J(๐ฒ_1)\delta J(๐ฒ_2)\delta J(๐ฒ_3)}}|_{(J,\lambda )=0}`$
$`=(e^{\phi (๐ฒ_1)}1)(e^{\phi (๐ฒ_2)}1)(e^{\phi (๐ฒ_3)}1)_{\mathrm{conn}.}`$
$`+i^3{\displaystyle \frac{\delta ^3\mathrm{\Delta }[J,\lambda ]}{\delta J(๐ฒ_1)\delta J(๐ฒ_2)\delta J(๐ฒ_3)}}|_{(J,\lambda )=0},`$ (A.12)
where we have to compute the last contribution in Eq. (VI). We need an expression for $`W_0[J,\lambda ]`$. Since the derivatives of $`W_0`$ with respect to $`J`$ and $`\lambda `$ (evaluated in $`(J,\lambda )=0`$) give the connected correlation functions for $`(e^\phi 1)`$ and $`\phi `$, we can write
$$W_0[J,\lambda ]=\underset{n=1}{\overset{\mathrm{}}{}}\frac{i^n}{n!}d๐ฑ_1\mathrm{}.d๐ฑ_n\stackrel{~}{z}(๐ฑ_1\mathrm{}.๐ฑ_n)\stackrel{~}{J}(๐ฑ_1)\mathrm{}\stackrel{~}{J}(๐ฑ_n),$$
(A.13)
where $`\stackrel{~}{J}(๐ฑ_i)`$ can be either the source $`J(๐ฑ_i)`$ or $`\lambda (๐ฑ_i)`$, and correspondingly $`\stackrel{~}{w}_n(๐ฑ_1\mathrm{}.๐ฑ_n)`$ are the connected correlation functions. For example $`(e^{\phi (๐ฑ_1)}1)\phi (๐ฑ_2)\phi (๐ฑ_3)`$ for the choice $`J(๐ฑ_1)\lambda (๐ฑ_2)\lambda (๐ฑ_3)`$, and one has to consider different combinations.
Using the usual operations for functional derivatives (see, for example, Mosel ) we find
$$i^3\frac{\delta ^3\mathrm{\Delta }[J,\lambda ]}{\delta J(๐ฒ_1)\delta J(๐ฒ_2)\delta J(๐ฒ_3)}|_{(J,\lambda )=0}=a\underset{p}{}[\stackrel{~}{w}_2(๐ฒ_{p_1},๐ฒ_{p_2}\stackrel{~}{w}_2(๐ฒ_{p_3},๐ฒ_{p_2})+\frac{1}{2}\stackrel{~}{z}_4(๐ฒ_{p_1},๐ฒ_{p_2},๐ฒ_{p_3},๐ฒ_{p_3})]$$
(A.14)
where the sum is over the permutations $`p_1,p_2,p_3`$ of indices $`(1,2,3)`$ and we have used the following notations
$`\stackrel{~}{w}_2(๐ฑ,๐ฒ)`$ $``$ $`(e^{\phi (๐ฑ)}1)\phi (๐ฒ)_{\mathrm{conn}.}`$ (A.15)
$`\stackrel{~}{w}_4(๐ฑ_1,๐ฑ_2,๐ฑ,๐ฑ)`$ $``$ $`(e^{\phi (๐ฑ_1)}1)(e^{\phi (๐ฑ_2)}1)\phi (๐ฑ)\phi (๐ฑ)_{\mathrm{conn}}.`$
Notice that it is possible to compute the connected correlation functions appearing in Eq. (A.15) and (VI) in a similar manner to what one does for $`e^\phi `$.
The result in Eq. (A.14) is a sum of two terms that correspond to the two pieces in Eq. (VI). One has two evaluate $`\phi `$ at two equal points because we have taken the derivative twice with respect to $`\lambda `$. Finally the product of the two-point connected correlation functions and the presence of the fourth-order connected correlation function are due to the fact that we have the product of two first derivatives w.r.t $`\lambda `$ and a second order derivative w.r.t $`\lambda `$, respectively. Taking then the three derivatives w.r.t. $`J`$ involves $`(e^\phi 1)`$ in the correlations. Now it is easy to generalize the result (A.14) when $`a`$ is not a constant coefficient but we have a kernel in configuration space such that
$$K(\phi )=๐\overline{๐ฑ}_1๐\overline{๐ฑ}_2K(๐ฑ\overline{๐ฑ}_\mathrm{๐},๐ฑ\overline{๐ฑ}_\mathrm{๐})\phi (\overline{๐ฑ}_1)\phi (\overline{๐ฑ}_2)+\mathrm{}.$$
(A.16)
As one can guess from Eq. (A.14) one has
$`i^3`$ $`{\displaystyle \frac{\delta ^3\mathrm{\Delta }[J,\lambda ]}{\delta J(๐ฒ_1)\delta J(๐ฒ_2)\delta J(๐ฒ_3)}}|_{(J,\lambda )=0}=`$
$`{\displaystyle \underset{p}{}}`$ $`{\displaystyle ๐\overline{๐ฑ}_1๐\overline{๐ฑ}_2K(๐ฒ_{p_2}\overline{๐ฑ}_\mathrm{๐},๐ฒ_{p_2}\overline{๐ฑ}_\mathrm{๐})\stackrel{~}{w}_2(๐ฒ_{p_1},\overline{๐ฑ}_1)\stackrel{~}{w}_2(๐ฒ_{p_3},\overline{๐ฑ}_2)}`$
$`+`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{p}{}}{\displaystyle ๐\overline{๐ฑ}_1๐\overline{๐ฑ}_2K(๐ฒ_{p_3}\overline{๐ฑ}_\mathrm{๐},๐ฒ_{p_3}\overline{๐ฑ}_\mathrm{๐})\stackrel{~}{w}_4(๐ฒ_{p_1},๐ฒ_{p_2},\overline{๐ฑ}_1,\overline{๐ฑ}_2)}.`$
In fact we have explicitly verified that this is the correct generalization of Eq. (A.14).
Thus from Eq. (VI) and Eq. (VI) the connected three-point correlation function for $`(e^{\phi +K(\phi )}1)`$ reads
$`W^{(3)}(๐ฒ_1,๐ฒ_2,๐ฒ_3)=i^3{\displaystyle \frac{\delta ^3Z[J,\lambda ]}{\delta J(๐ฒ_1)\delta J(๐ฒ_2)\delta J(๐ฒ_3)}}|_{(J,\lambda )=0}=`$
$`(e^{\phi (๐ฒ_1)}1)(e^{\phi (๐ฒ_2)}1)(e^{\phi (๐ฒ_3)}1)_{\mathrm{connected}}`$
$`+`$ $`{\displaystyle \underset{p}{}}{\displaystyle ๐\overline{๐ฑ}_1๐\overline{๐ฑ}_2K(๐ฒ_{p_2}\overline{๐ฑ}_\mathrm{๐},๐ฒ_{p_2}\overline{๐ฑ}_\mathrm{๐})\stackrel{~}{w}_2(๐ฒ_{p_1},\overline{๐ฑ}_1)\stackrel{~}{w}_2(๐ฒ_{p_3},\overline{๐ฑ}_2)}`$
$`+`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{p}{}}{\displaystyle ๐\overline{๐ฑ}_1๐\overline{๐ฑ}_2K(๐ฒ_{p_3}\overline{๐ฑ}_\mathrm{๐},๐ฒ_{p_3}\overline{๐ฑ}_\mathrm{๐})\stackrel{~}{w}_4(๐ฒ_{p_1},๐ฒ_{p_2},\overline{๐ฑ}_1,\overline{๐ฑ}_2)},`$
where one has to use the definitions in Eqs. (A.15) and (VI).
## Acknowledgments
N.B. would like to thank James Babington for useful discussions on techniques of functional-integral analysis in quantum field theory.
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# ๐-adic multiple zeta values II โ tannakian interpretations
## 0. Introduction
This paper is the continuation of our previous paper \[F1\]. Let $`p`$ be a prime number. Let $`m,k_1,\mathrm{},k_m`$. In \[F1\] we constructed $`p`$-adic multiple polylogarithm $`Li_{k_1,\mathrm{},k_m}(z)`$ ($`z_p`$), that is, a $`p`$-adic analogue of the (one-variable) complex multiple polylogarithm defined locally near $`0`$ by
(0.1)
$$Li_{k_1,\mathrm{},k_m}(z)=\underset{0<n_1<\mathrm{}<n_m}{}\frac{z^{n_m}}{n_1^{k_1}\mathrm{}n_m^{k_m}}$$
and $`p`$-adic multiple zeta value $`\zeta _p(k_1,\mathrm{},k_m)`$, that is, a $`p`$-adic analogue of multiple zeta value
(0.2)
$$\zeta (k_1,\mathrm{},k_m)=\underset{0<n_1<\mathrm{}<n_m}{}\frac{1}{n_1^{k_1}\mathrm{}n_m^{k_m}}(k_m>1).$$
This was achieved by giving a $`p`$-adic Drinfelโd associator $`\mathrm{\Phi }_{\mathrm{KZ}}^p(A,B)`$ of the $`p`$-adic KZ equation and its fundamental solution $`G_0(A,B)(z)`$ ($`z_p`$). The purpose of this paper is to relate these constructions to the motivic fundamental torsors of the projective line minus three points. To do this we will give a tannakian interpretation of the constructions of \[F1\]. Using the Frobenius action on the rigid realization of the above torsor, we will define the overconvergent $`p`$-adic multiple polylogarithm. This function can be expressed by the $`p`$-adic multiple polylogarithm constructed in \[F1\]. Using Hodge and $`l`$-adic รฉtale realizations we will recover the complex multiple polylogarithm and the $`l`$-adic polylogarithms of Wojtkowiak.
Our main tool is a comparison isomorphism (Lemma 2.2)
(0.3)
$$\pi _1^{\mathrm{DR}}(__p^1\backslash \{0,1,\mathrm{}\}:\stackrel{}{01},z)\pi _1^{p,\mathrm{rig}}(_{๐ฝ_p}^1\backslash \{0,1,\mathrm{}\}:\stackrel{}{01},z_0)$$
between the de Rham (ยง1.1) and the rigid (ยง1.2) fundamental torsor where $`z_0^1(๐ฝ_p)\backslash \{0,1,\mathrm{}\}`$ is the reduction of $`z^1(_p)\backslash \{0,1,\mathrm{}\}`$ and $`\stackrel{}{01}`$ is a tangential basepoint. There is a canonical de Rham path $`d_z`$ (Lemma 1.2) in LHS constructed from the canonical extension of unipotent connections and a canonical rigid path $`c_{z_0}`$ (Lemma 1.8) in RHS, Besser-Vologodskyโs Frobenius invariant path \[Bes, Vol\]. Identifying these two torsors by (0.3) we get a de Rham loop $`d_z^1c_{z_0}\pi _1^{\mathrm{DR}}(__p^1\backslash \{0,1,\mathrm{}\}:\stackrel{}{01})(_p)`$. Our tannakian interpretation of $`p`$-adic multiple polylogarithms is
Theorem 2.3. By the embedding $`i:\pi _1^{\mathrm{DR}}(__p^1\backslash \{0,1,\mathrm{}\}:\stackrel{}{01})(_p)_pA,B`$ (2.1), the loop $`d_z^1c_{z_0}`$ corresponds to the non-commutative formal power series $`G_0(A,B)(z)`$ of the $`p`$-adic multiple polylogarithms.
Similarly by letting $`z`$ to be a tangential basepoint $`\stackrel{}{10}`$ we get a tannakian interpretation of $`p`$-adic multiple zeta values:
Theorem 2.5. By the embedding $`i`$ the loop $`d_{\stackrel{}{10}}^1c_{\stackrel{}{10}}`$ corresponds to the non-commutative formal power series $`\mathrm{\Phi }_{\mathrm{KZ}}^p(A,B)`$ of the $`p`$-adic multiple zeta values.
RHS of (0.3) admits the Frobenius action $`\varphi _p`$ (Definition 1.7). Deligne \[De2\] introduced a variant of $`p`$-adic multiple zeta values to be a coefficient of the series $`\mathrm{\Phi }_{\mathrm{De}}^p(A,B)=i(d_{\stackrel{}{10}}^1\varphi _p(d_{\stackrel{}{10}}))`$ which we call the $`p`$-adic Deligne associator.
Theorem 2.8. There is the following explicit relationship between the $`p`$-adic Drinfelโd associator and the $`p`$-adic Deligne associator:
$$\mathrm{\Phi }_{\mathrm{KZ}}^p(A,B)=\mathrm{\Phi }_{\mathrm{De}}^p(A,B)\mathrm{\Phi }_{\mathrm{KZ}}^p(\frac{A}{p},\mathrm{\Phi }_{\mathrm{De}}^p(A,B)^1\frac{B}{p}\mathrm{\Phi }_{\mathrm{De}}^p(A,B)).$$
This gives formulae expressing our $`p`$-adic multiple zeta values in terms of Deligneโs $`p`$-adic multiple zeta values and vice versa. Namely his $`p`$-adic multiple zeta values are equivalent to ours. We introduce a pro-unipotent overconvergent iso-crystal $`(๐ฑ_{\mathrm{KZ}}^{},_{\mathrm{KZ}}^{})`$ on $`_{๐ฝ_p}^1\backslash \{0,1,\mathrm{}\}`$ associated with the $`p`$-adic KZ equation and show
Proposition 2.11. The pro-object $`(๐ฑ_{\mathrm{KZ}}^{},_{\mathrm{KZ}}^{})`$ naturally admits a Frobenius structure.
The canonical section $`1`$ of $`๐ฑ^{}`$ under the Frobenius action will map to a non-commutative power series $`G_0^{}(z)`$ whose coefficient will be denoted by $`Li_{k_1,\mathrm{},k_m}^{}(z)`$. This series satisfies the differential equation in \[U2, Y\]
$$dg=\left(\frac{A}{z}+\frac{B}{z1}\right)gdzg\left(\frac{dz^p}{z^p}\frac{A}{p}+\frac{dz^p}{z^p1}\mathrm{\Phi }_{\mathrm{De}}^p(A,B)^1\frac{B}{p}\mathrm{\Phi }_{\mathrm{De}}^p(A,B)\right)$$
and $`Li_{k_1,\mathrm{},k_m}^{}(z)`$ are an overconvergent analogue of our $`p`$-adic multiple polylogarithms. The special value of these functions at $`z=1`$ will be the Deligneโs $`p`$-adic multiple zeta values. The relation between $`G_0^{}(A,B)(z)`$ and our $`G_0(A,B)(z)`$ is given by the following:
Theorem 2.14. The overconvergent $`p`$-adic multiple polylogarithm $`Li_{k_1,\mathrm{},k_m}^{}(z)`$ is expressed as a combination of our $`p`$-adic multiple zeta values, $`p`$-adic multiple polylogarithms and the $`p`$-adic logarithm by
$$G_0^{}(A,B)(z)=G_0(A,B)(z)G_0(\frac{A}{p},\mathrm{\Phi }_{\mathrm{De}}^p(A,B)^1\frac{B}{p}\mathrm{\Phi }_{\mathrm{De}}^p(A,B))(z^p)^1.$$
The function $`Li_{k_1,\mathrm{},k_m}(z)`$ is not overconvergent but a Coleman function with log poles around $`z=1`$ and $`\mathrm{}`$. The above formulae is a way to erase these log poles. We will give calculations in Example 2.16.
The organization of this paper is as follows. ยง1 is preliminary for the rest of paper. We will review the definition of various (de Rham, rigid, Betti and รฉtale) fundamental groups, torsors with their various additional structures, i.e. the Frobenius action, the infinity Frobenius action and the Galois group action. We also discuss tangential basepoints. The Deligneโs canonical path in the de Rham fundamental torsor in Lemma 1.2 and Besser-Vologodskyโs Frobenius invariant path in the rigid fundamental torsor in Lemma 1.8 will be recalled.
ยง2 is devoted to tannakian formalisms which contains our main results. In Theorem 2.3 and Theorem 2.5 we will clarify tannakian origins of the $`p`$-adic multiple polylogarithms and $`p`$-adic multiple zeta values. An explicit relationship between our $`p`$-adic multiple zeta values and Deligneโs $`p`$-adic multiple zeta values \[De2\] will be stated in Theorem 2.8. A formula to express the overconvergent $`p`$-adic multiple polylogarithms in terms of our $`p`$-adic multiple polylogarithms will be stated in Theorem 2.14. In ยง2.2 the Hodge analogue of these results will be discussed. The tannakian origin of complex multiple polylogarithms will be discussed in Proposition 2.20 and for multiple zeta values in Proposition 2.21. The infinity Frobenius action on the Betti fundamental torsor of $`^1()\backslash \{0,1,\mathrm{}\}`$ will be used to introduce another variant of multiple polylogarithms (2.20). In Theorem 2.27 it will be shown that these functions are single-valued and real-analytic. In Proposition 2.26 we show a formula analogous to Theorem 2.14 to express them in terms of complex multiple polylogarithms. ยง2.3 is a brief explanation of analogous story in the รฉtale side. We discuss an รฉtale analogue of polylogarithm (2.22) expressed in terms of Wojtkowiakโs $`l`$-adic polylogarithm and an รฉtale analogue of Riemann zeta values (Example 2.33) expressed in term of Soulรฉ characters and cyclotomic characters following \[NW, I\].
ยง3 is motivic. The algebra generated by $`p`$-adic multiple zeta values will be related to Drinfelโdโs \[Dr\] pro-algebraic group $`\underset{ยฏ}{GRT}_1`$, Racinetโs \[R\] pro-algebraic group $`\underset{ยฏ}{DMR}_0`$ and the motivic Galois group $`\pi _1(๐ฏ())`$. In ยง3.1 we recall Drinfelโdโs pro-algebraic bi-torsor of Grothendieck-Teichmรผller. In ยง3.2, we recall Racinetโs \[R\] pro-algebraic bi-torsor made by double shuffle relations and briefly explain a partial analogous story to ยง3.1. ยง3.3 is a review of the formalism of mixed Tate motives in \[DG\]. The torsor of motivic Galois group is related with the torsors in ยง3.1 and ยง3.2. We give some motivic interpretations on Zagierโs conjecture on multiple zeta values and Iharaโs conjecture on Galois image. We also explain a motivic way of proving double shuffle relations for $`p`$-adic multiple zeta values using \[DG\] and \[Y\].
###### Acknowledgments .
This work was supported by the NSF grants DMS-0111298. The author would like to express particular thanks to Prof. Deligne for many suggestions. He is also grateful to Amir Jafari and Go Yamashita for many useful discussion. He is also grateful to the referee for comments.
## 1. Preliminaries
This section is preliminaries for the next section. We give materials which are indispensable to our main results by recalling the definitions of various (de Rham, rigid, Betti and รฉtale) fundamental groups and torsors with tangential base points and discussing additional structures there.
In this section, we concentrate on a curve $`X_K=\overline{X}_KD_K`$, where $`\overline{X}_K`$ is a proper smooth, geometrically connected curve $`\overline{X}_K`$ over a field $`K`$ and $`D_K`$ is its divisor over $`K`$ (sometimes we omit $`K`$).
### 1.1. de Rham setting
In this subsection, we will recall the definitions of the de Rham fundamental group, torsor, path space, the tangential base point and the canonical base point and see how we obtain the canonical de Rham path in the de Rham fundamental torsor by the canonical base point, all of which are developed in \[De1\]ยง12.
###### Notation 1.1.
In this subsection we assume that $`K`$ is a field of characteristic $`0`$. We denote the category of the nilpotent part <sup>1</sup><sup>1</sup>1 It means the full subcategory consisting of objects which are iterated extensions of the unit objects. of the category of the pair $`(๐ฑ,)`$ of a coherent $`๐ช_X`$-module sheaf $`๐ฑ`$ on $`X_K`$ and an integrable (i.e. $`=0`$) connection $`:๐ฑ๐ฑ\mathrm{\Omega }_{X_K}^1`$ by $`๐ฉ๐^{\mathrm{DR}}(X_K)`$, It forms a neutral tannakian category (for the basics, consult \[DM\]) by \[De1\]. For our quick review of these materials, see \[S1\].
In the de Rham realization, we consider three fiber functors. Suppose that $`x`$ is a $`K`$-valued point of $`X`$, i.e. $`x:SpecKX_K`$. The first one is $`\omega _x:๐ฉ๐^{\mathrm{DR}}(X_K)Vec_K`$ ($`Vec_K`$: the category of finite dimensional $`K`$-vector spaces) which is the pull-back of $`๐ฉ๐^{\mathrm{DR}}(X_K)`$ by $`x`$. Actually this forms a fiber functor (consult \[DM\] for its definition) by \[S1\]ยง3.1. The second fiber functor is the one introduced in \[De1\]ยง15 (and developed to higher dimensional case in \[BF\]): Let $`sD`$ and put $`T_s=T_s\overline{X}_K`$ and $`T_s^\times =T_s\backslash \{0\}`$($`๐พ_m`$), where $`T_s\overline{X}_K`$ is the tangent vector space of $`\overline{X}_K`$ at $`s`$. By the construction of the extension of unipotent integrable connections to tangent bundles in \[De1\]ยง15.28-15.36, we have a morphism of tannakian categories (see also \[BF\])
(1.1)
$$Res_s:๐ฉ๐^{\mathrm{DR}}(X_K)๐ฉ๐^{\mathrm{DR}}(T_s^\times )$$
which we call the tangential morphism. By pulling back, each $`K`$-valued point $`t_s`$ of $`T_s^\times `$ determines a fiber functor $`\omega _{t_s}:๐ฉ๐^{\mathrm{DR}}(X_K)Vec_K`$ which we call the tangential base point. The last fiber functor is the special one under the assumption
(1.2)
$$H^1(\overline{X},๐ช_{\overline{X}})=0.$$
(i.e. $`\overline{X}`$ is a curve with genus $`0`$): Let $`\omega _\mathrm{\Gamma }:๐ฉ๐^{\mathrm{DR}}(X_K)Vec_K`$ be a functor sending each object $`(๐ฑ,)๐ฉ๐^{\mathrm{DR}}(X_K)`$ to $`\mathrm{\Gamma }(\overline{X}_K,๐ฑ_{\mathrm{can}})`$, the global section of its canonical extension $`(๐ฑ_{\mathrm{can}},_{\mathrm{can}})`$ into $`\overline{X}_K`$ \[De1\] ยง12.2. This functor forms a fiber functor (\[De1\] ยง12.4) because $`๐ฑ_{\mathrm{can}}`$ forms a trivial bundle over $`\overline{X}_K`$ due to the condition (1.2) by \[De1\] Proposition 12.3. We call $`\omega _\mathrm{\Gamma }`$ the canonical base point. We note that the assumption (1.2) is necessary. The functor $`\omega _\mathrm{\Gamma }`$ would be no longer a fiber functor if $`\overline{X}`$ was a proper smooth curve with genus $`g>0`$.
###### Lemma 1.2 (\[De1\]ยง15.52).
Let $`\omega _{}`$ and $`\omega _{^{}}`$ be any fiber functors above. Then under the assumption (1.2) there is a canonical isomorphism $`d_{^{}}:\omega _{}\omega _{^{}}`$.
###### Proof .
It is because $`๐ฑ_{\mathrm{can}}`$ forms a trivial bundle for $`(๐ฑ,)๐ฉ๐^{\mathrm{DR}}(X_K)`$ that we get a canonical isomorphism
$$\omega _x(๐ฑ,)=๐ฑ_{(x)}=๐ฑ_{\mathrm{can},(x)}๐ช_{\overline{X}_K,(x)}\underset{๐พ}{}\mathrm{\Gamma }(\overline{X}_K,๐ฑ_{\mathrm{can}})๐ช_{\overline{X}_K,(x)}\underset{๐พ}{}\omega _\mathrm{\Gamma }(๐ฑ,)\omega _\mathrm{\Gamma }(๐ฑ,).$$
and similarly
$$\omega _{t_s}(๐ฑ,)๐ช_{T_s^\times ,(t_s)}\underset{๐พ}{}\mathrm{\Gamma }(\overline{X}_K,๐ฑ_{\mathrm{can}})\omega _\mathrm{\Gamma }(๐ฑ,).$$
Here <sub>(x)</sub> stands for the fiber of sheaves at $`x`$. By factoring through $`\omega _\mathrm{\Gamma }`$, we obtain a canonical isomorphism between any two of them. โ
The fiber functors $`\omega _x`$ and $`\omega _{t_s}`$ (and $`\omega _\mathrm{\Gamma }`$ under the assumption (1.2)) play a role of โbase pointsโ.
###### Definition 1.3 (\[De1\]).
Let $`\omega _{}`$ and $`\omega _{^{}}`$ be any fiber functors above. The de Rham fundamental group $`\pi _1^{\mathrm{DR}}(X_K:)`$ is $`\underset{ยฏ}{Aut}^{}(\omega _{})`$ and the de Rham fundamental torsor $`\pi _1^{\mathrm{DR}}(X_K:,^{})`$ is $`\underset{ยฏ}{Isom}^{}(\omega _{},\omega _{^{}})`$ (for $`\underset{ยฏ}{Aut}^{}`$ and $`\underset{ยฏ}{Isom}^{}`$, see \[DM\]).
The former is a pro-algebraic group over $`K`$ and the latter is a pro-algebraic bitorsor over $`K`$ where $`\pi _1^{\mathrm{DR}}(X_K:)`$ acts from the right and $`\pi _1^{\mathrm{DR}}(X_K:^{})`$ acts from the left <sup>2</sup><sup>2</sup>2 For any two loops $`\gamma `$ and $`\gamma ^{}`$, the symbol $`\gamma ^{}\gamma `$ stands for the composition of two which takes $`\gamma `$ at first and takes $`\gamma ^{}`$ next. Throughout this paper, we keep this direction. and by definition $`\pi _1^{\mathrm{DR}}(X_K:,)=\pi _1^{\mathrm{DR}}(X_K:)`$. Particularly under the assumption (1.2), we have a canonical de Rham path $`d_,^{}\pi _1^{\mathrm{DR}}(X_K:,^{})(K)`$ by Lemma 1.2. This path is compatible with the composition: $`d_{^{},^{\prime \prime }}d_,^{}=d_{,^{\prime \prime }}`$. We also note that for any field extension $`K^{}/K`$ the category $`๐ฉ๐^{\mathrm{DR}}(X_K^{})`$ is equivalent to $`๐ฉ๐^{\mathrm{DR}}(X_K)_KK^{}`$ which is the category whose set of morphisms, $`K^{}`$-linear space, is replaced by the extension of the set of morphisms of $`๐ฉ๐^{\mathrm{DR}}(X_K)`$, $`K`$-linear space, by $`K^{}`$ (\[De1\]ยง10.41), whence $`\pi _1^{\mathrm{DR}}(X_K^{}:)\pi _1^{\mathrm{DR}}(X_K:)\times _KK^{}`$ and $`\pi _1^{\mathrm{DR}}(X_K^{}:,^{})\pi _1^{\mathrm{DR}}(X_K:,^{})\times _KK^{}`$.
To establish the notion of the de Rham path space, we recall the notion of affine group $`๐ฏ`$-schemes for a tannakian category $`๐ฏ`$ (\[De1\]ยง5.4 and \[DG\]ยง2.6). The category of affine $`๐ฏ`$-schemes is the dual of the category of ind-objects of $`๐ฏ`$ which are unitary commutative algebras, that is, each object is an ind-object $`A`$ of $`๐ฏ`$ with a product $`AAA`$ and a unit $`1A`$ verifying usual axioms. We denote the corresponding affine $`๐ฏ`$-schemes with $`A`$ by $`SpecA`$. An affine group $`๐ฏ`$-scheme is a group object of the category of affine $`๐ฏ`$-schemes, i.e. the spectrum of a commutative Hopf algebra.
###### Definition 1.4.
The de Rham path space <sup>3</sup><sup>3</sup>3In \[De1\] ยง6.13, it is called the fundamental groupoid. $`๐ซ_{X_K}^{\mathrm{DR}}`$ is an affine group $`๐ฉ๐^{\mathrm{DR}}(X_K)๐ฉ๐^{\mathrm{DR}}(X_K)`$ <sup>4</sup><sup>4</sup>4 That is the tensor product of tannakian categories \[De1\]ยง5.18. -scheme introduced in \[De1\]ยง6.13 and \[DG\]ยง4, which satisfies
$$(\omega \omega ^{})(๐ซ_{X_K}^{\mathrm{DR}})=\underset{ยฏ}{Isom}^{}(\omega ,\omega ^{})(K)$$
for any fiber functors $`\omega ,\omega ^{}:๐ฉ๐^{\mathrm{DR}}(X_K)Vec_K`$ and is represented by an ind-object $`๐^{\mathrm{DR}}`$ of $`๐ฉ๐^{\mathrm{DR}}(X_K)๐ฉ๐^{\mathrm{DR}}(X_K)`$ i.e. $`๐ซ_{X_K}^{\mathrm{DR}}=Spec๐^{\mathrm{DR}}`$. The de Rham fundamental path space $`๐ซ_{X_K,\omega }^{\mathrm{DR}}`$ with the base point $`\omega `$ is its pull-back by $`\omega id`$, which is represented by the ind-object $`๐_\omega ^{\mathrm{DR}}:=(\omega id)(๐^{\mathrm{DR}})`$ of $`๐ฉ๐^{\mathrm{DR}}(X_K)`$.
By definition, we have
(1.3)
$$(\omega _{}\omega _{^{}})๐ซ_{X_K}^{\mathrm{DR}}=\omega _{^{}}(๐ซ_{X_K,\omega _{}}^{\mathrm{DR}})=\pi _1^{\mathrm{DR}}(X_K:,^{})(K).$$
### 1.2. rigid setting
In this paper we take rigid fundamental groups (torsors) as a crystalline realization of motivic fundamental groups (torsors). We will recall the definition of the rigid fundamental group \[CLS, S2\], torsor, path space, the tangential base point \[BF, De1\] and the canonical Frobenius invariant path \[Bes, Vol\].
###### Notation 1.5.
In this subsection we assume that $`K`$ is a non-archimedean local field of characteristic $`0`$. We denote $`V`$ to be its valuation ring and $`k`$ to be the residue field of characteristic $`p>0`$ with $`q(=p^r)`$-elements. We also assume that $`\overline{X}_K`$ with $`D_K`$ is a generic fiber of a proper smooth geometrically connected curve $`๐ณ`$ with a divisor $`๐`$ over $`V`$. We denote its special fiber over $`k`$ by $`\overline{X}_0`$ with $`D_0`$. Put $`X_0=\overline{X}_0\backslash D_0`$ and its natural embedding $`j:X_0\overline{X}_0`$. We denote the category of the nilpotent part of the category of (overconvergent <sup>5</sup><sup>5</sup>5 We can omit it because it is shown in \[CLS\] that unipotent (nilpotent) isocrystals are overconvergent. ) isocrystals, which consists of the pair $`(๐ฑ^{},^{})`$ of a coherent $`j^{}๐ช_{]\overline{X}_0[}`$ <sup>6</sup><sup>6</sup>6 We do not call $``$ โplusโ but โdaggerโ. For the definition of $`j^{}๐ช_{]\overline{X}_0[}`$, see \[Bes\]. -module sheaf $`๐ฑ^{}`$ and an integrable connection $`^{}:๐ฑ^{}๐ฑ^{}_{๐ช_{]\overline{X}_0[}}\mathrm{\Omega }_{]\overline{X}_0[}^1`$, by $`๐ฉ^{}(X_0)`$. Here $`]\overline{X}_0[`$ stands for the tubular neighborhood \[Ber\] of $`\overline{X}_0`$. This category depends only on $`X_0`$ and forms a neutral tannakian category and depends only on $`X_0`$ (cf. \[Ber, CLS, S2\]). For our quick review of these materials, see \[Bes\].
In the rigid realization, we consider three fiber functors. Suppose that $`x_0`$ is a $`k`$-valued point of $`X_0`$, i.e. $`x_0:SpeckX_0`$. The first one is $`\omega _{x_0}:๐ฉ^{}(X_0)Vec_{K_0}`$ which is the pull-back of $`๐ฉ^{}(X_0)`$ by $`x_0`$. This forms a fiber functor and is described by $`\omega _{x_0}(๐ฑ^{},^{})=๐ฑ^{}(]x_0[)^{^{}}=\{v๐ฑ^{}(]x_0[)|^{}(v)=0\}`$ for $`(๐ฑ^{},^{})๐ฉ^{}(X_0)`$ on the tubular neighborhood of $`]x_0[`$ \[Bes\]. The second fiber functor is $`\omega _{s_0}:๐ฉ^{}(X_0)Vec_K`$ ($`s_0D_0`$) introduced in \[BF\] Definition 1.3: This is described by $`\omega _{s_0}(๐ฑ^{},^{})=\{v๐ฑ^{}(]s_0[)[\mathrm{log}^az_{s_0}]|^{}(v)=0\}`$, where $`\mathrm{log}^a`$ is a branch of $`p`$-adic logarithm such that $`\mathrm{log}^ap=aK`$ and $`z_{s_0}`$ is a local parameter of $`]s_0[`$ (cf.\[F1\] ยง2.1). The last fiber functor is the rigid version of the tangential base point: Put $`T_{s_0}=T_{s_0}\overline{X}_0`$ and $`T_{s_0}^\times =T_{s_0}\backslash \{0\}(๐พ_m)`$ where $`T_{s_0}\overline{X}_0`$ is the tangent space of $`\overline{X}_0`$ at $`s_0`$. By the construction in \[BF\]ยง2, we have a morphism of tannakian categories
(1.4)
$$Res_{s_0}:๐ฉ^{}(X_0)๐ฉ^{}(T_{s_0}^\times ),$$
which we call the tangential morphism. Thus each $`k`$-valued point $`t_{s_0}`$ of $`T_{s_0}^\times `$ determines a fiber functor $`\omega _{t_{s_0}}:๐ฉ^{}(X_0)Vec_K`$ by pulling back. Shortly the construction of $`Res_{s_0}`$ developed in \[BF\] is the transmission of the morphism (1.1) into the rigid setting by the categorical equivalences $`๐ฉ๐^{\mathrm{DR}}(X_K)๐ฉ^{}(X_0)`$ and $`๐ฉ๐^{\mathrm{DR}}(T_s^\times )๐ฉ^{}(T_{s_0}^\times )`$ in \[CLS, S2\], where $`s_0`$ is the reduction of $`s`$.
The fiber functors, $`\omega _{x_0}`$, $`\omega _{s_0}`$ and $`\omega _{t_{s_0}}`$, play a role of โbase pointsโ.
###### Definition 1.6.
Let $`\omega _{}`$ and $`\omega _{^{}}`$ be any fiber functors above. The rigid fundamental group $`\pi _1^{p,\mathrm{rig}}(X_0:)`$ is $`\underset{ยฏ}{Aut}^{}(\omega _{})`$ and the rigid fundamental torsor $`\pi _1^{p,\mathrm{rig}}(X_0:,^{})`$ is $`\underset{ยฏ}{Isom}^{}(\omega _{},\omega _{^{}})`$ (cf. \[CLS, S2\]). The rigid path space $`๐ซ_{X_0}^{}`$ is an affine group $`๐ฉ^{}(X_0)๐ฉ^{}(X_0)`$-scheme which satisfies
$$(\omega \omega ^{})(๐ซ_{X_0}^{})=\underset{ยฏ}{Isom}^{}(\omega ,\omega ^{})(K)$$
for any fiber functors $`\omega ,\omega ^{}:๐ฉ^{}(X_0)Vec_K`$ and is represented by an ind-object $`๐^{}`$ of $`๐ฉ^{}(X_0)๐ฉ^{}(X_0)`$, i.e. $`๐ซ_{X_0}^{}=Spec๐^{}`$. The rigid path space $`๐ซ_{X_0,\omega }^{}`$ with the base point $`\omega `$ is its pull-back by $`\omega id`$, which is represented by the ind-object $`๐_\omega ^{}:=(\omega id)(๐^{})`$ of $`๐ฉ^{}(X_0)`$.
We note that $`\pi _1^{p,\mathrm{rig}}(X_0:,^{})`$ is a pro-algebraic bitorsor over $`K`$ where $`\pi _1^{p,\mathrm{rig}}(X_0:)`$ acts from the right and $`\pi _1^{p,\mathrm{rig}}(X_0:^{})`$ acts from the left. By definition we have
$$(\omega _{}\omega _{^{}})๐ซ_{X_0}^{}=\omega _{^{}}(๐ซ_{X_0,\omega _{}}^{})=\pi _1^{p,\mathrm{rig}}(X_0:,^{})(K).$$
Let $`F_q`$ be $`r`$-th power of the absolute Frobenius automorphism on $`X_0`$ and its induced automorphism on $`T_{s_0}^\times `$ $`(s_0D_0)`$. The pull-back functors $`F_q^{}:๐ฉ^{}(X_0)๐ฉ^{}(X_0)`$ and $`F_q^{}:๐ฉ^{}(T_{s_0}^\times )๐ฉ^{}(T_{s_0}^\times )`$ determine equivalences of tannakian categories. We note that they are compatible with (1.4) (\[BF\] ยง2). They induce isomorphisms (\[De1\] ยง10.44)
(1.5)
$$F_q:๐ซ_{X_0}F_q^{}๐ซ_{X_0},$$
(1.6)
$$F_q:๐ซ_{X_0,\omega }F_q^{}(๐ซ_{X_0,F_q^{}\omega })$$
(1.7)
$$F_q:\pi _1^{p,\mathrm{rig}}(X_0:,^{})\pi _1^{p,\mathrm{rig}}(X_0:,^{}).$$
###### Definition 1.7.
(\[De1\] ยง10.44) The Frobenius actions $`\varphi _q`$ stands for the inverse <sup>7</sup><sup>7</sup>7 We take the inverse to follow \[De1\] ยง13.6. $`F_q^1`$ in (1.5) $``$ (1.7).
By (1.5) and (1.6) the ind-objects $`๐^{}`$ and $`๐_\omega ^{}`$ in Definition 1.6 naturally admit structures of ind-overconvergent $`F`$-isocrystals, where an overconvergent $`F`$-isocrystal means an overconvergent isocrystal $`V`$ with a Frobenius structure $`\varphi `$ (i.e. a horizontal isomorphism $`\varphi :F_q^{}VV`$). The description of the Frobenius structure in case of $`X=^1\backslash \{0,1,\mathrm{}\}`$ is one of our main topics in this paper. The following lemma which takes a place of Lemma 1.2 in our rigid setting is indispensable to our study.
###### Lemma 1.8.
Let $`\omega _{}`$ and $`\omega _{^{}}`$ be any fiber functors defined above. Then there exists a unique path $`c_,^{}\pi _1^{p,\mathrm{rig}}(X_0:,^{})(K)`$ which is invariant under the Frobenius action, i.e. $`F_q(c_,^{})=c_,^{}`$.
This result is shown by Besser \[Bes\] and Vologodsky \[Vol\] for usual base points (more generally in higher dimensional setting) but their proofs also work for tangential base points and points on $`D_0`$ directly. We call $`c_,^{}`$ the Frobenius invariant path. We have compatibilities, such as functorialities and $`c_{^{},^{\prime \prime }}c_,^{}=c_{,^{\prime \prime }}`$ (for more details, see \[Bes\]).
### 1.3. Betti setting
We will briefly recall the definition of the Betti fundamental group, torsor and the tangential base point in \[De1\]ยง15.
###### Notation 1.9.
In this subsection, we assume $`k=`$. We denote the category of the nilpotent part of the category of local systems <sup>8</sup><sup>8</sup>8 We mean covariant functors $`V:\mathrm{\Pi }_{X()}Vec_{}`$ where $`\mathrm{\Pi }_{X()}`$ stands for the category whose objects are points on $`X()`$ and whose set of morphisms between two points $``$ and $`^{}`$ on $`X()`$ is the usual topological fundamental torsor $`\pi _1^{\mathrm{top}}(X():,^{})`$. of finite dimensional $``$-vector spaces over the topological space $`X()`$ by $`๐ฉ^{\mathrm{Be}}(X())`$. It also forms a neutral tannakian category over $``$ by \[De1\]ยง10.
In the Betti realization, we consider two fiber functors. Suppose that $`xX()`$. The first one is $`\omega _x:๐ฉ^{\mathrm{Be}}(X())Vec_{}`$ which associates each nilpotent local system $`V`$ with the pullback $`V(x)`$. Actually this forms a fiber functor of $`๐ฉ^{\mathrm{Be}}(X())`$ by \[De1\]ยง10. The second one is the tangential basepoint: Let $`sD()`$ and put $`T_s=T_s\overline{X}()`$ and $`T_s^\times =T_s\backslash \{0\}`$($`๐พ_m`$), where $`T_s\overline{X}()`$ is the tangent vector space of $`\overline{X}()`$ at $`s`$. By \[De1\]ยงยง15.3-15.12, we have a morphism of tannakian categories
(1.8)
$$Res_s:๐ฉ^{\mathrm{Be}}(X())๐ฉ^{\mathrm{Be}}(T_s^\times ()).$$
which we call the tangential morphism. By pulling back, each point $`t_sT_s^\times ()`$ determines a fiber functor $`\omega _{t_s}:๐ฉ^{\mathrm{Be}}(X())Vec_{}`$.
The above fiber functors, $`\omega _x`$ and $`\omega _{t_s}`$, play a role of โbase pointsโ and the following is a reformulation of the Betti fundamental group and torsor in \[De1\]ยง13.
###### Definition 1.10.
Let $`\omega _{}`$ and $`\omega _{}^{}`$ be any fiber functors above. The Betti fundamental group $`\pi _1^{\mathrm{Be}}(X():)`$ is $`\underset{ยฏ}{Aut}^{}(\omega _{})`$ and the Betti fundamental torsor $`\pi _1^{\mathrm{Be}}(X():,^{})`$ is $`\underset{ยฏ}{Isom}^{}(\omega _{},\omega _{^{}})`$.
The former is a pro-algebraic group over $``$ and the latter is a pro-algebraic bi-torsor over $``$ where $`\pi _1^{\mathrm{Be}}(X():)`$ acts on the right and $`\pi _1^{\mathrm{Be}}(X():^{})`$ acts on the left. We note that we have a natural morphism from the topological fundamental group $`\pi _1^{\mathrm{top}}(X():)`$ (resp. torsor $`\pi _1^{\mathrm{top}}(X():,^{})`$) to the set of $``$-valued points on the Betti fundamental group $`\pi _1^{\mathrm{Be}}(X():)`$ (resp. torsor $`\pi _1^{\mathrm{Be}}(X():,^{})`$), which induces an isomorphism between $`\pi _1^{\mathrm{Be}}(X():)`$ and the Malcev (unipotent) completion \[De1\]ยง9 over $``$ of $`\pi _1^{\mathrm{top}}(X():)`$ (\[De1\] ยง10.24). As for the topological fundamental groups and torsors with tangential base points, see \[De1\]ยงยง15.3-15.12. Although we have a canonical path, both in the de Rham setting $`d_,^{}`$ Lemma 1.2 and in the rigid setting $`c_,^{}`$ Lemma 1.8, the author do not know whether we have any canonical path connecting any two base points $``$ and $`^{}`$ in the Betti setting or not.
###### Assumption 1.11.
Further we assume that $`(\overline{X},D)`$ has a real structure $`(\overline{X^{}},D^{})`$, i.e. $`(\overline{X},D)=(\overline{X^{}}\times _{},D^{}\times _{})`$,
Let $`F_{\mathrm{}}`$ be an involution of $`X()`$ induced from the complex conjugation $`z\overline{z}`$ and its induced automorphism on $`T_s^\times `$ ($`sD`$). The pull-back functor $`F_{\mathrm{}}^{}:๐ฉ^{\mathrm{Be}}(X())๐ฉ^{\mathrm{Be}}(X())`$ and $`F_{\mathrm{}}^{}:๐ฉ^{\mathrm{Be}}(T_s^\times ())๐ฉ^{\mathrm{Be}}(T_s^\times ())`$ determine equivalences of tannakian categories. We note that they are compatible with (1.8). They induce the isomorphism
$$F_{\mathrm{}}:\pi _1^{\mathrm{Be}}(X():,^{})\pi _1^{\mathrm{Be}}(X():\overline{},\overline{}^{})$$
where $`\overline{}`$ (resp. $`\overline{}^{}`$) means $`F_{\mathrm{}}()`$ (resp. $`F_{\mathrm{}}(^{})`$).
###### Definition 1.12.
The infinity Frobenius action $`\varphi _{\mathrm{}}`$ stands for the inverse of $`F_{\mathrm{}}`$ above.
The description of the infinity Frobenius action in case of $`X=^1\backslash \{0,1,\mathrm{}\}`$ will be discussed in ยง2.
### 1.4. รฉtale setting
We will briefly recall the definition of the รฉtale fundamental group, torsor and the tangential base point in \[De1\]ยง15.
###### Notation 1.13.
In this subsection, we assume that $`K`$ is a field of characteristic $`0`$. We denote $`\overline{K}`$ to be its algebraic closure and $`X_{\overline{K}}`$ to be the scalar extension of $`X_K`$ by $`\overline{K}`$. Let $`l`$ be a prime (we do not need to assume that $`l`$ is inequal to the prime $`p`$ in ยง1.2). We denote the category of the nilpotent part of the category of รฉtale local systems <sup>9</sup><sup>9</sup>9 We mean covariant functors $`V:\widehat{\mathrm{\Pi }}_{X_{\overline{K}}}Vec__l`$ where $`\widehat{\mathrm{\Pi }}_{X_{\overline{K}}}`$ stands for the category whose objects are points on $`X(\overline{K})`$ and whose set of morphisms between two points $``$ and $`^{}`$ on $`X(\overline{K})`$ is the pro-finite fundamental torsor $`\widehat{\pi }_1(X_{\overline{K}}:,^{})`$ in \[G\]. of finite dimensional $`_l`$-vector spaces over $`X(\overline{K})`$ by $`๐ฉ^l(X_{\overline{K}})`$. It also forms a neutral tannakian category over $`_l`$.
In the รฉtale realization, we consider two fiber functors. Suppose that $`xX(\overline{K})`$. The first one is the pullback $`\omega _x:๐ฉ^l(X_{\overline{K}})Vec__l`$ which associates each nilpotent รฉtale local system $`V`$ with the pullback $`V(x)`$. The second one is the tangential base point: Let $`sD(\overline{K})`$ and put $`T_s=T_s\overline{X}_{\overline{K}}`$ and $`T_s^\times =T_s\backslash \{0\}`$($`๐พ_m`$), where $`T_s\overline{X}_{\overline{K}}`$ is the tangent vector space of $`\overline{X}_{\overline{K}}`$ at $`s`$. By \[De1\]ยงยง15.13-15.26, we have a morphism of tannakian categories
(1.9)
$$Res_s:๐ฉ^l(X_{\overline{K}})๐ฉ^l(T_s^\times )$$
which we call the tangential morphism. By pulling back, each $`\overline{K}`$-valued point $`t_s`$ in $`T_s^\times `$ determines a fiber functor $`\omega _{t_s}:๐ฉ^l(X_{\overline{K}})Vec__l`$.
The above fiber functors, $`\omega _x`$ and $`\omega _{t_s}`$, play a role of โbasepointsโ and the following is a reformulation of the $`l`$-adic รฉtale fundamental group and torsor in \[De1\] ยง13.
###### Definition 1.14.
Let $`\omega _{}`$ and $`\omega _{^{}}`$ be any fiber functors above. The $`l`$-adic รฉtale fundamental group $`\pi _1^{l,\text{รฉt}}(X_{\overline{K}}:)`$ is $`\underset{ยฏ}{Aut}^{}(\omega _{})`$ and the $`l`$-adic รฉtale fundamental torsor $`\pi _1^{l,\text{รฉt}}(X_{\overline{K}}:,^{})`$ is $`\underset{ยฏ}{Isom}^{}(\omega _{},\omega _{^{}})`$.
The former is a pro-algebraic group over $`_l`$ and the latter is a pro-algebraic bitorsor over $`_l`$ where $`\pi _1^{l,\text{รฉt}}(X_{\overline{K}}:)`$ acts on the right and $`\pi _1^{l,\text{รฉt}}(X_{\overline{K}}:^{})`$ acts on the left. Let $`\overline{K^{}}`$ be an algebraically closed field containing $`\overline{K}`$. The scalar extension induces a categorical equivalence $`๐ฉ^l(X_{\overline{K}^{}})๐ฉ^l(X_{\overline{K}})`$, from which we get isomorphisms $`\pi _1^{l,\text{รฉt}}(X_{\overline{K}^{}}:)\pi _1^{l,\text{รฉt}}(X_{\overline{K}}:)`$ and $`\pi _1^{l,\text{รฉt}}(X_{\overline{K}^{}}:,^{})\pi _1^{l,\text{รฉt}}(X_{\overline{K}}:,^{})`$ (\[De1\] ยง10.19).
Suppose that $``$ and $`^{}`$ are $`K`$-valued points on $`X`$ or $`T_s^\times `$ for $`sD`$ (this assumption is attained by enlarging the base field $`K`$). Let $`\sigma Gal(\overline{K}/K)`$. We get Galois action
$$\sigma :\pi _1^{l,\text{รฉt}}(X_{\overline{K}}:,^{})\pi _1^{l,\text{รฉt}}(X_{\overline{K}}:,^{})$$
from the Galois action $`id\sigma :X_{\overline{K}}X_{\overline{K}}`$ in a similar way to our construction of $`F_q`$ in ยง1.2 and $`F_{\mathrm{}}`$ in ยง1.3. We note that the pro-finite fundamental group $`\widehat{\pi _1}(X_{\overline{K}}:)`$ (resp. torsor $`\widehat{\pi _1}(X_{\overline{K}}:,^{})`$) also admits Galois action and we have a morphism from $`\widehat{\pi _1}(X_{\overline{K}}:)\pi _1^{l,\text{รฉt}}(X_{\overline{K}}:)(_l)`$ (resp. $`\widehat{\pi _1}(X_{\overline{K}}:,^{})\pi _1^{l,\text{รฉt}}(X_{\overline{K}}:,^{})(_l)`$) which is compatible with its group (resp. torsor) structure and their Galois group actions. Actually this morphism induces an isomorphism between $`\pi _1^{l,\text{รฉt}}(X_{\overline{K}}:)`$ and the Malcev (unipotent) completion over $`_l`$ of $`\widehat{\pi _1}(X_{\overline{K}}:)`$ \[De1\] ยง10.
###### Lemma 1.15.
Assume that $`K`$ is a number field with ring of integers $`๐ช_K`$. Let $``$ and $`^{}`$ be any $`K`$-valued points. Let $`๐ญ`$ be a prime ideal, coprime to $`l๐ช_K`$, where the triple $`(X,,^{})`$ has a proper smooth model over $`๐ช_{K,๐ญ}`$ with a good reduction. Then there exists a unique path $`e_,^{}^๐ญ\pi _1^{l,\text{รฉt}}(X_{\overline{K}}:,^{})(_l)`$ which is invariant under the action of Frobenius element $`Frob_๐ญ`$ of $`๐ญ`$, i.e. $`Frob_๐ญ(e_,^{}^๐ญ)=e_,^{}^๐ญ`$.
###### Proof .
By the above assumption $`\pi _1^{l,\text{รฉt}}(X_{\overline{K}}:,^{})`$ is unramified at $`๐ญ`$, whence we have an action of the geometric Frobenius $`Frob_๐ญ`$ of $`๐ญ`$. Since the eigenvalues of the $`Frob_๐ญ`$-action on $`H_{\text{รฉt}}^1(X_{\overline{K}},_l)`$ are positive weight Weil numbers, we obtain a unique Frobenius invariant path $`e_,^{}^๐ญ`$ by an inductive construction same to \[Bes\]ยง3 (or \[Vol\]ยง4.3). โ
We note that this Frobenius invariant path $`e_,^{}^๐ญ`$ equips nice compatibilities such as functorialities and $`e_{^{},^{\prime \prime }}^๐ญe_,^{}^๐ญ=e_{,^{\prime \prime }}^๐ญ`$ similar to $`d_,^{}`$ in ยง1.1 and $`c_,^{}`$ in ยง1.2. As Besser \[Bes\] gave the Frobenius invariant path $`c_,^{}`$ an interpretation of Colemanโs $`p`$-adic iterated integration in the rigid setting, the author thinks that there may be also a significance on the above special path $`e_,^{}^๐ญ`$ in the $`l`$-adic รฉtale setting.
## 2. Tannakian interpretations
This section is divided into three sides. The first side is the Berthelot-Ogus side (ยง2.1) where we shall give our main results. A tannakian interpretation of $`p`$-adic MPLโs and $`p`$-adic MZVโs ($`p`$: prime) <sup>10</sup><sup>10</sup>10 MPL stands for multiple polylogarithm and MZV stands for multiple zeta values. is given in Theorem 2.3 and Theorem 2.5. By describing Frobenius action on the rigid fundamental group of the projective line minus three points we deduce an explicit formula between our $`p`$-adic MZVโs and Deligneโs $`p`$-adic MZVโs \[De2, DG\] in Theorem 2.8. Further by considering Frobenius action on a rigid path space we introduce an overconvergent version of $`p`$-adic MPL in Definition 2.13. Its expression in terms of our $`p`$-adic MPLโs is given in Theorem 2.14. The second side is the Hodge side (ยง2.2), where results analogous to ยง2.1 are presented. A tannakian reinterpretation of MPLโs and MZVโs are given in Proposition 2.20 and Proposition 2.21 following \[BD\]. Using the infinity Frobenius action on the Betti fundamental torsor of the projective line minus three points we introduce a new version of MPL (2.20). Its monodromy-free property is shown in Theorem 2.27. Its description is given in Proposition 2.26. The third side is the Artin side (ยง2.3). We discuss $`l`$-adic รฉtale analogues of MPLโs and MZVโs defined in a way that Wojtkowiak \[W\] introduced the $`l`$-adic polylogarithms.
Throughout this and next section, we fix $`X=^1\backslash \{0,1,\mathrm{}\}=Spec[t,\frac{1}{t},\frac{1}{t1}]\overline{X}=^1`$ and $`D=\{0,1,\mathrm{}\}`$. We denote the local parameter of $`T_0^\times `$ induced from $`t`$ by $`t_0`$ and the point $`t_0=1`$ on $`T_0^\times `$ by $`\stackrel{}{01}`$. Similarly we denote the local parameter of $`T_1^\times `$ induced from $`1t`$ by $`t_1`$ and the point $`t_1=1`$ on $`T_1^\times `$ by $`\stackrel{}{10}`$.
### 2.1. Berthelot-Ogus side โ main results
In this side, we shall use a Berthelot-Ogus-type comparison isomorphism (Lemma 2.2) between a de Rham fundamental torsor and a rigid fundamental torsor for $`X`$. Our first results are a tannakian interpretation of $`p`$-adic MPLโs and $`p`$-adic MZVโs in Theorem 2.3 and Theorem 2.5 respectively. Secondly we consider Frobenius action on a rigid fundamental group. The definition of Deligneโs $`p`$-adic MZVโs \[De2, DG\] which describes the Frobenius action is recalled. By using a tangential morphism we give an explicit formula between his and ours in Theorem 2.8. Thirdly we consider Frobenius action on a rigid path space. Overconvergent $`p`$-adic MPL is introduced in Definition 2.13. It describes the Frobenius action and takes Deligneโs $`p`$-adic MZV as โa special valueโ at $`z=1`$. A formula to express it in terms of our $`p`$-adic MPLโs is given in Theorem 2.14.
Before we state our first results we give notations and basic lemma below which will be frequently employed in this subsection.
###### Notation 2.1.
We denote the special path $`d_{\stackrel{}{01},z}`$1.1) in the de Rham fundamental torsor $`\pi _1^{\mathrm{DR}}(X:\stackrel{}{01},z)(_p)`$ for $`zX(_p)`$ by $`d_z`$. In particular we simply denote $`d_{\stackrel{}{01},\stackrel{}{10}}`$ in by $`d`$. Put $`X_0=_{๐ฝ_p}^1\backslash \{0,1,\mathrm{}\}\overline{X_0}=_{๐ฝ_p}^1`$. Similarly we denote the special path $`c_{\stackrel{}{01},z_0}`$1.2) in the rigid fundamental torsor $`\pi _1^{p,\mathrm{rig}}(X_0:\stackrel{}{01},z_0)(_p)`$ for $`z_0\overline{X_0}(๐ฝ_p)`$ by $`c_{z_0}`$. In particular we simply denote $`c_{\stackrel{}{01},\stackrel{}{10}}`$ by $`c`$. The composition of the embedding $`\pi _1^{\mathrm{DR}}(T_0^\times :\stackrel{}{01})\pi _1^{\mathrm{DR}}(X:\stackrel{}{01})`$ by (1.1) and the exponential map $`H_1^{\mathrm{DR}}(T_0^\times ,)\pi _1^{\mathrm{DR}}(T_0^\times ,\stackrel{}{01})()`$ sends the dual of $`\frac{dt_0}{t_0}`$ to a loop denoted by $`x`$. Similarly a composite morphism to $`\pi _1^{\mathrm{DR}}(T_0^\times ,\stackrel{}{10})()`$ sends the dual of $`\frac{dt_1}{t_1}`$ to a loop denoted by $`y^{}`$. Put $`y=d^1y^{}d\pi _1^{\mathrm{DR}}(X:\stackrel{}{01})()`$ Then it is easy to see that $`x`$ and $`y`$ are free generators of the pro-algebraic group $`\pi _1^{\mathrm{DR}}(X:\stackrel{}{01})`$. In this section, we consider the embedding into the non-commutative formal power series with two variables $`A`$ and $`B`$
(2.1)
$$i:\pi _1^{\mathrm{DR}}(X:\stackrel{}{01})()A,B$$
which sends $`xe^A:=1+\frac{A}{1!}+\frac{A^2}{2!}+\mathrm{}`$ and $`ye^B:=1+\frac{B}{1!}+\frac{B^2}{2!}+\mathrm{}`$. We note that this morphism is isomorphic to the embedding of $`\pi _1^{\mathrm{DR}}(X:\stackrel{}{01})()`$ into the universal embedding algebra of $`Lie\pi _1^{\mathrm{DR}}(X:\stackrel{}{01})()`$. By abuse of notation we also employ the same symbol $`i`$ when we mention its scalar extension.
The following Berthelot-Ogus type comparison isomorphism for fundamental torsors is our basics in this subsection.
###### Lemma 2.2.
Let $`z`$ be a $`_p`$-valued point of $`X`$ or $`T_s^\times `$ (for $`sD`$) and $`z_0`$ be its modulo $`p`$ reduction.
1. If $`z_0`$ lies on $`X_0(๐ฝ_p)`$ or $`T_{s_0}^\times (๐ฝ_p)`$,
$$\pi _1^{\mathrm{DR}}(X__p:\stackrel{}{01},z)\pi _1^{p,\mathrm{rig}}(X_0:\stackrel{}{01},z_0).$$
2. If $`z_0`$ does not lie on $`X_0(๐ฝ_p)`$ nor $`T_{s_0}^\times (๐ฝ_p)`$,
$$\pi _1^{\mathrm{DR}}(X__p:\stackrel{}{01},z)\times _{p,st}\pi _1^{p,\mathrm{rig}}(X_0:\stackrel{}{01},z_0)\times _{p,\mathrm{st}}$$
where $`_{p,\mathrm{st}}`$ stands for the polynomial algebra $`_p[l(p)]`$ with one variable $`l(p)`$.
###### Proof .
An equivalence of tannakian categories
(2.2)
$$๐ฉ๐^{\mathrm{DR}}(X__p)๐ฉ^{}(X_{๐ฝ_p})$$
is shown in \[CLS, S2\] in more general situation. In case of $`z_0X_0(๐ฝ_p)`$ it gives rise to a natural isomorphism
$$\omega _z(๐ฑ,)\omega _{z_0}(๐ฑ^{},^{})$$
for $`(๐ฑ,)๐ฉ๐^{\mathrm{DR}}(X__p)`$, where $`(๐ฑ^{},^{})๐ฉ^{}(X_{๐ฝ_p})`$ is associated with $`(๐ฑ,)`$ via (2.2). The equivalences (2.2) and $`๐ฉ๐^{\mathrm{DR}}(T_{s,_p}^\times )๐ฉ^{}(T_{s_0}^\times )`$ are compatible with tangential morphisms (cf. \[BF\]). They also give rise to an natural isomorphism
$$\omega _z\left(Res(๐ฑ,)\right)\omega _{z_0}\left(Res(๐ฑ^{},^{})\right)$$
for $`(๐ฑ,)๐ฉ๐^{\mathrm{DR}}(X__p)`$ in case of $`z_0T_0^\times (๐ฝ_p)`$. The equivalence (2.2) and the above two isomorphisms of fiber functors follow the isomorphism in (1). In the rest case a natural isomorphism
$$\omega _z(๐ฑ,)_p[l(p)]\omega _{z_0}(๐ฑ^{},^{})_p[l(p)]$$
is explained in \[Vol\] ยง4.4, which follows the isomorphism in (2). โ
We note that this variable $`l(p)`$ reflects a branch parameter $`a_p`$ of the $`p`$-adic polylogarithm $`\mathrm{log}^a`$. By identifying the above two fundamental torsors by Lemma 2.2 we get a de Rham loop $`d_z^1c_{z_0}`$, which lies in $`\pi _1^{\mathrm{DR}}(X:\overline{01})(_p)`$ in case of (1) and in $`\pi _1^{\mathrm{DR}}(X:\overline{01})(_{p,\mathrm{st}})`$ in case of (2).
###### Theorem 2.3.
Let $`z`$ be a $`_p`$-valued point of $`X`$ and $`z_0`$ be its modulo $`p`$ reduction on $`\overline{X_0}`$. The de Rham loop $`d_z^1c_{z_0}`$ corresponds to the special value of the fundamental solution $`G_0`$ <sup>11</sup><sup>11</sup>11 In case of $`z_0X_0(๐ฝ_p)`$ it means the fundamental solution of the $`p`$-adic KZ (Knizhnik-Zamolodchikov) equation in $`_pA,B`$ constructed in \[F1\]Theorem 3.3 and in case of $`z_0X_0(๐ฝ_p)`$ it means the series in $`_{p,\mathrm{st}}A,B`$ whose specialization at $`l(p)=a`$ is $`G_0^a(z)`$ in loc.cit. at $`z`$ of the $`p`$-adic KZ equation by the embedding $`i`$:
$$i(d_z^1c_{z_0})=G_0(z).$$
###### Proof .
Our main tools to prove the theorem are a tannakian description of the KZ equation and Besserโs tannakian interpretation of Colemanโs $`p`$-adic iterated integration theory.
We start with the KZ equation. Let $`(๐ฑ_{\mathrm{KZ}},_{\mathrm{KZ}})`$ be a pro-object of $`๐ฉ๐^{\mathrm{DR}}(X)`$ associated with the KZ equation. Namely $`๐ฑ_{\mathrm{KZ}}=๐ช_XA,B`$ and $`_{\mathrm{KZ}}(g)=dg(\frac{A}{t}+\frac{B}{t1})gdt`$ for local section $`g`$. By \[De1\] Proposition 12.10, the Lie algebra $`๐ญ^{\mathrm{DR}}`$ of $`\pi _1^{\mathrm{DR}}(X:\mathrm{\Gamma })`$ is free generated by two representatives $`(\frac{dt}{t})^{}`$ and $`(\frac{dt}{t1})^{}`$ in $`H_1^{\mathrm{DR}}(X,)`$ which are dual basis of $`\frac{dt}{t}`$ and $`\frac{dt}{t1}`$ in $`H_{\mathrm{DR}}^1(X,)H^0(\overline{X},\mathrm{\Omega }_{\overline{X}}^1(\mathrm{log}D))`$ respectively because $`H^1(\overline{X},๐ช_{\overline{X}})=0`$. By the basic theorem \[DM\] of tannakian categories, the fiber functor $`\omega _\mathrm{\Gamma }`$ induces an equivalence of categories
(2.3)
$$๐ฉ๐^{\mathrm{DR}}(X_K)๐ฉRep๐ญ^{\mathrm{DR}}$$
where the right hand side stands for the category of nilpotent representations of the Lie algebra on finite dimensional $``$-vector spaces.
###### Lemma 2.4.
By (2.3) the pro-object $`(๐ฑ_{\mathrm{KZ}},_{\mathrm{KZ}})`$ corresponds to the vector space $`A,B`$ with the left multiplication $`L:๐ญ^{\mathrm{DR}}๐คl(A,B)`$ such that $`L\left((\frac{dt}{t})^{}\right)=L_A`$ and $`L\left((\frac{dt}{t1})^{}\right)=L_B`$. Here $`L_F`$ for $`FA,B`$ means the left multiplication by $`F`$.
###### Proof .
In \[De1\]ยง12.5, Deligne give a recipe of the correspondence (2.3) (in more general situation): Let $`(๐ฑ,)๐ฉ๐^{\mathrm{DR}}(X)`$. We put $`V=\mathrm{\Gamma }(\overline{X},๐ฑ_{\mathrm{can}})=\omega _\mathrm{\Gamma }(๐ฑ,)`$ and decompose as $`_{\mathrm{can}}=d+\omega `$ where $`d:V๐ช_{\overline{X}}V\mathrm{\Omega }_{\overline{X}}^1(\mathrm{log}D)`$ is the differential induced from $`d:๐ช_{\overline{X}}\mathrm{\Omega }_{\overline{X}}^1(\mathrm{log}D)`$ and $`\omega H^0(\overline{X},\mathrm{\Omega }_{\overline{X}}^1(\mathrm{log}D))EndV`$. By the integrability of $``$, we have the Lie algebra homomorphism $`\rho :๐ญ^{\mathrm{DR}}๐ค๐ฉV`$ such that $`\rho |_{H_1^{\mathrm{DR}}(X,)}=\omega `$, which gives an object $`(V,\rho )๐ฉRep๐ญ^{\mathrm{DR}}`$. Conversely let $`(V,\rho )๐ฉRep๐ญ^{\mathrm{DR}}`$. We put $`\omega =\rho |_{H_1^{\mathrm{DR}}(X,)}(\alpha )H_{\mathrm{DR}}^1(X,)EndVH^0(\overline{X},\mathrm{\Omega }_{\overline{X}}^1(\mathrm{log}D))EndV`$ where $`\alpha H_{\mathrm{DR}}^1(X,)H_1^{\mathrm{DR}}(X,)`$ is a canonical tensor representing the identity on $`H_1^{\mathrm{DR}}(X,)`$. This gives $`(๐ฑ,)=(V๐ช_X,d_X+\omega |_X)๐ฉ๐^{\mathrm{DR}}(X)`$.
Following his recipe we see that $`L`$ corresponds to
$$\omega =\frac{dt}{t}L_A+\frac{dt}{t1}L_BH_1^{\mathrm{DR}}(X,)\widehat{}๐คl(A,B)$$
which gives the claim. โ
By abuse of notation, we denote its induced group homomorphism by
$$L:\pi _1^{\mathrm{DR}}(X:\mathrm{\Gamma })AutA,B.$$
By the constructions of $`i`$ and $`L`$ above, we have
(2.4)
$$i=ev_1LIntd_{\stackrel{}{01},\mathrm{\Gamma }}.$$
Here $`Intd_{\stackrel{}{01},\mathrm{\Gamma }}:\pi _1^{\mathrm{DR}}(X:\stackrel{}{01})()\pi _1^{\mathrm{DR}}(X:\mathrm{\Gamma })()`$ is the automorphism sending $`\gamma d_{\stackrel{}{01},\mathrm{\Gamma }}\gamma d_{\stackrel{}{01},\mathrm{\Gamma }}^1`$ and $`ev_1:AutA,BA,B`$ is the evaluation at $`1`$ sending $`\sigma \sigma (1)`$.
Our second tool is Besserโs tannakian interpretation of the Colemanโs $`p`$-adic iterated integration theory \[C\] in \[Bes\]: He called a set $`v=\{v_{x_0}\omega _{x_0}(๐ฑ^{},^{})|x_0\overline{X_0}(\overline{๐ฝ_p})\}`$ for $`(๐ฑ^{},^{})๐ฉ^{}(X_0)`$ a collection of analytic continuation along Frobenius of horizontal section of $`(๐ฑ^{},^{})`$ if $`c_{x_0,x_0^{}}(v_{x_0})=v_{x_0^{}}`$ for all $`x_0,x_0^{}\overline{X_0}(\overline{๐ฝ_p})`$ where $`c_{x_0,x_0^{}}`$ is a Frobenius invariant path in Lemma 1.8 and showed that each Coleman function of $`X(_p)`$ (in the sense of \[Bes\] ยง5) is expressed as a set $`\{\theta (v_{x_0})\}_{x_0}`$ with a $`j^{}๐ช_{]\overline{X_0}[}_p`$-module homomorphism $`\theta :๐ฑ^{}_pj^{}๐ช_{]\overline{X_0}[}_p`$ and a collection $`v=\{v_{x_0}\}_{x_0}`$ of analytic continuation along Frobenius of a pair $`(๐ฑ^{},^{})๐ฉ^{}(X_0)`$. In \[BF\] we further extend Coleman functions to normal bundles by making use of Frobenius invariant path to tangential basepoints.
In our previous paper \[F1\] we consider the $`p`$-adic KZ equation and construct its fundamental solution $`G_0(t)`$, a two variable non-commutative formal power series with Coleman function coefficients. In his terminologies it is a collection $`\left\{G_0|_{]x_0[}\right\}_{x_0}`$ of analytic continuation along Frobenius of horizontal section of $`(๐ฑ_{\mathrm{KZ}}^{},_{\mathrm{KZ}}^{})`$ (the associated pro-object of $`๐ฉ^{}(X_0)`$ with $`(๐ฑ_{\mathrm{KZ}},_{\mathrm{KZ}})`$), i.e. $`c_{x_0,y_0}(G_0|_{]x_0[})=G_0|_{]y_0[}`$ for $`x_0`$ and $`y_0`$. And since in our case the tangential morphism (1.1) for $`s=0`$ sends $`(๐ฑ_{\mathrm{KZ}},_{\mathrm{KZ}})`$ to $`(๐ช_{T_0^\times }A,B,dA\frac{dt_0}{t_0})`$, $`t_0^A=1+\frac{\mathrm{log}^at_0}{1!}A+\frac{(\mathrm{log}^at_0)^2}{2!}A^2+\frac{(\mathrm{log}^at_0)^3}{3!}A^3+\mathrm{}`$ might be a collection of analytic continuation along Frobenius of horizontal section of $`Res_0(๐ฑ_{\mathrm{KZ}}^{},_{\mathrm{KZ}}^{})`$. By our tangential base point interpretation (\[BF\] Proposition 2.11) of the notion of the constant term (\[BF\] Definition 2.1) of horizontal sections, we see that $`t_0^A`$ is an analytic continuation along Frobenius of $`G_0(t)`$ to $`T_0^\times `$ because both the constant term of $`G_0(t)`$ at $`t=0`$ and the constant term of $`t_0^A`$ at $`t_0=0`$ are equal to $`1`$. Namely $`c_{\stackrel{}{01},y_0}(t_0^A|_{]\stackrel{}{01}[})=G_0(t)|_{]y_0[}`$.
Let $`z`$ be a $`_p`$-valued point $`X`$ with reduction $`z_0`$. We regard $`c_{z_0}`$ to be a de Rham path via the identification in Lemma 2.4. The isomorphism $`d_{,\mathrm{\Gamma }}:\omega _{}(๐ฑ_{\mathrm{KZ}},_{\mathrm{KZ}})\omega _\mathrm{\Gamma }(๐ฑ_{\mathrm{KZ}},_{\mathrm{KZ}})=_pA,B`$ for $`=\stackrel{}{01}`$ and $`z`$ let us rewrite
(2.5)
$$d_{z,\mathrm{\Gamma }}c_{z_0}d_{\mathrm{\Gamma },\stackrel{}{01}}(1)=G_0(z).$$
By (2.4), (2.5) and $`d_{z,\mathrm{\Gamma }}=d_{\stackrel{}{01},\mathrm{\Gamma }}d_z^1`$, we prove Theorem 2.3. โ
In our previous paper \[F1\] Theorem 3.15 we showed that in each coefficient of $`G_0(z)`$ there appears $`p`$-adic MPL $`Li_{k_1,\mathrm{},k_m}(z)`$, which is a Coleman function admitting the expansion (0.1) on $`|z|_p<1`$. In precise the coefficient of $`A^{k_m1}B\mathrm{}A^{k_11}B`$ in $`G_0(z)`$ is $`(1)^mLi_{k_1,\mathrm{},k_m}(z)`$:
$$G_0(z)=1+\mathrm{log}zALi_1(z)B+\mathrm{}+(1)^mLi_{k_1,\mathrm{},k_m}(z)A^{k_m1}B\mathrm{}A^{k_11}B+\mathrm{}.$$
By Theorem 2.3 we may say that the loop $`d_z^1c_{z_0}`$ is a tannakian origin of $`p`$-adic MPLโs.
In case when $`z`$ is also a tangential base point $`\stackrel{}{01}`$, similarly we have
###### Theorem 2.5.
The de Rham loop $`d^1c\pi _1^{\mathrm{DR}}(X:\stackrel{}{01})(_p)`$ corresponds to the $`p`$-adic Drinfelโd associator $`\mathrm{\Phi }_{\mathrm{KZ}}^p_pA,B`$ by the embedding $`i`$, i.e.
$$i(d^1c)=\mathrm{\Phi }_{\mathrm{KZ}}^p.$$
###### Proof .
The $`p`$-adic Drinfelโd associator $`\mathrm{\Phi }_{\mathrm{KZ}}^p`$ is the series constructed in \[F1\] Definition 3.12 which is equal to the limit value (\[F1\] Lemma 3.27)
(2.6)
$$\underset{ฯต0}{lim}^{}\mathrm{exp}\left(\mathrm{log}^aฯตB\right)G_0^a(1ฯต)_pA,B.$$
In other word it is the constant term of $`G_0^a(z)`$ at $`z=1`$. Then by a similar argument to the proof of Theorem 2.3 we have $`c_{\stackrel{}{01},\stackrel{}{10}}(t_0^A|_{]\stackrel{}{01}[})=\mathrm{\Phi }_{\mathrm{KZ}}^pt_1^B|_{]\stackrel{}{10}[}`$, whence
(2.7)
$$d_{\stackrel{}{10},\mathrm{\Gamma }}cd_{\mathrm{\Gamma },\stackrel{}{01}}(1)=\mathrm{\Phi }_{\mathrm{KZ}}^p.$$
By (2.4), (2.16) and $`d_{\stackrel{}{10},\mathrm{\Gamma }}=d_{\stackrel{}{01},\mathrm{\Gamma }}d^1`$, we prove Theorem 2.5. โ
In our previous paper \[F1\] Theorem 3.30 we showed that in each coefficient of the $`p`$-adic Drinfelโd associator $`\mathrm{\Phi }_{\mathrm{KZ}}^p`$ there appears $`p`$-adic MZV $`\zeta _p(k_1,\mathrm{},k_m)`$ ($`k_m>1`$) introduced in \[F1\], which is a $`p`$-adic analogue of (0.2). In precise the coefficient of $`A^{k_m1}B\mathrm{}A^{k_11}B`$ in $`\mathrm{\Phi }_{\mathrm{KZ}}^p`$ is $`(1)^m\zeta _p(k_1,\mathrm{},k_m)`$:
$$\mathrm{\Phi }_{\mathrm{KZ}}^p=1+\mathrm{}+(1)^m\zeta _p(k_1,\mathrm{},k_m)A^{k_m1}B\mathrm{}A^{k_11}B+\mathrm{}.$$
By Theorem 2.5 we may say that the loop $`d^1c`$ is a tannakian origin of $`p`$-adic MZVโs.
###### Note 2.6.
The $`p`$-adic Drinfelโd associator $`\mathrm{\Phi }_{\mathrm{KZ}}^p`$ is group-like \[F1\], that means $`\mathrm{\Delta }\mathrm{\Phi }_{\mathrm{KZ}}^p=\mathrm{\Phi }_{\mathrm{KZ}}^p\widehat{}\mathrm{\Phi }_{\mathrm{KZ}}^p`$ where $`\mathrm{\Delta }:_pA,B_pA,B`$ is the linear map induced from $`\mathrm{\Delta }(A)=A1+1A`$ and $`\mathrm{\Delta }(B)=B1+1B`$. Hence each coefficient of $`\mathrm{\Phi }_{\mathrm{KZ}}^p`$ must satisfy (integral) shuffle product formula (cf. \[F1\], \[FJ\]). We can recover and express general coefficients of $`\mathrm{\Phi }_{\mathrm{KZ}}^p`$ in terms of $`\zeta _p(k_1,\mathrm{},k_m)`$โs ($`k_m>1`$) by the following method: Let $`W`$ be a word, i.e. a monic and monomial element in $`_pA,B`$. Put $`C_W`$ to be its coefficient of $`\mathrm{\Phi }_{\mathrm{KZ}}^p`$. We have $`C_A=C_B=0`$. For a convergent word $`W`$ written as $`A^{k_m1}B\mathrm{}A^{k_11}B`$ ($`k_m>1`$), $`C_W=(1)^m\zeta _p(k_1,\mathrm{},k_m)`$. For a divergent word $`W^{}`$ written as $`B^rW`$ ($`r>0`$) with $`W`$: convergent, shuffle product formula gives
$$C_B^rC_W=r!C_{B^rW}+\text{other terms}.$$
By an induction with respect to $`r`$, $`C_W^{}`$ is calculated and is expressed in terms of $`C_W`$โs with $`W`$: convergent. For a word $`W^{\prime \prime }=B^rWA^s`$ ($`s>0`$) with $`W`$: convergent, a similar induction argument let us able to calculate $`C_{W^{\prime \prime }}`$.
Our second result is on a description of a Frobenius action. We consider a Frobenius action $`F_p`$ on $`\pi _1^{\mathrm{DR}}(X:\stackrel{}{01},\stackrel{}{10})(_p)`$ by transmitting that on $`\pi _1^{p,\mathrm{rig}}(X_0:\stackrel{}{01},\stackrel{}{10})(_p)`$ by Lemma 2.2. It gives a new path $`F_p(d)`$ in $`\pi _1^{\mathrm{DR}}(X:\stackrel{}{01},\stackrel{}{10})(_p)`$ which is our central object to discuss here. In Arizona Winter School 2002, Deligne \[De2\] introduced another version of $`p`$-adic MZV which has a different tannakian origin from ours (see also \[DG\] ยง5).
###### Definition 2.7.
The $`p`$-adic Deligne associator $`\mathrm{\Phi }_{\mathrm{De}}^p`$ is the series in $`_pA,B`$ which corresponds to $`d^1\varphi _p(d)`$:
$$i(d^1\varphi _p(d))=\mathrm{\Phi }_{\mathrm{De}}^p.$$
To adapt the notations to ours above, we denote $`\zeta _p^{\mathrm{De}}(k_1,\mathrm{},k_m)`$ ($`k_m>1`$) to be the coefficient of $`A^{k_m1}B\mathrm{}A^{k_11}B`$ in $`\mathrm{\Phi }_{\mathrm{De}}^p`$ multiplied by $`(1)^m`$:
$$\mathrm{\Phi }_{\mathrm{De}}^p=1+\mathrm{}+(1)^m\zeta _p^{\mathrm{De}}(k_1,\mathrm{},k_m)A^{k_m1}B\mathrm{}A^{k_11}B+\mathrm{}$$
and call it Deligneโs $`p`$-adic MZV. Since $`\mathrm{\Phi }_{\mathrm{De}}^p`$ is group-like, we can recover and express general coefficients of $`\mathrm{\Phi }_{\mathrm{De}}^p`$ in terms of $`\zeta _p^{\mathrm{De}}(k_1,\mathrm{},k_m)`$โs ($`k_m>1`$) as Note 2.6. Deligneโs $`p`$-adic MZVโs come from the loop $`d^1\varphi _p(d)`$ while ours come from the loop $`d^1c`$. So his are different from ours. But we can give a close relationship between them:
###### Theorem 2.8.
(2.8)
$$\mathrm{\Phi }_{\mathrm{KZ}}^p(A,B)=\mathrm{\Phi }_{\mathrm{De}}^p(A,B)\mathrm{\Phi }_{\mathrm{KZ}}^p(\frac{A}{p},\mathrm{\Phi }_{\mathrm{De}}^p(A,B)^1\frac{B}{p}\mathrm{\Phi }_{\mathrm{De}}^p(A,B)).$$
Here the last term means the series substituting $`(\frac{A}{p},\mathrm{\Phi }_{\mathrm{De}}^p(A,B)^1\frac{B}{p}\mathrm{\Phi }_{\mathrm{De}}^p(A,B))`$ in $`\mathrm{\Phi }_{\mathrm{KZ}}^p(A,B)`$.
This formula looks nasty. But according to the convention <sup>12</sup><sup>12</sup>12 Let $`k`$ be a field with characteristic $`0`$. Let $`c_1,c_2k`$ and $`g_1(A,B),g_2(A,B)kA,B`$. The product $`(c_2,g_2)(c_1,g_1)`$ is defined by $`(c_1c_2,g_2(A,B)g_1(\frac{A}{c_2},g_2(A,B)^1\frac{B}{c_2}g_2(A,B)))`$. in \[F0\] of the Grothendieck-Teichmรผller group \[Dr\] we rewrite simpler $`(p,\mathrm{\Phi }_{\mathrm{KZ}}^p)=(p,\mathrm{\Phi }_{\mathrm{De}}^p)(1,\mathrm{\Phi }_{\mathrm{KZ}}^p)`$, i.e.
(2.9)
$$(p,\mathrm{\Phi }_{\mathrm{De}}^p)=(p,\mathrm{\Phi }_{\mathrm{KZ}}^p)(1,\mathrm{\Phi }_{\mathrm{KZ}}^p)^1.$$
###### Proof .
We extends the Frobenius action $`\varphi _p`$ on $`\pi _1^{\mathrm{DR}}(X,\stackrel{}{01})(_p)`$ to $`_pA,B`$ by $`i`$.
###### Lemma 2.9.
$`\varphi _p(A)=\frac{A}{p}`$, $`\varphi _p(B)=\mathrm{\Phi }_{\mathrm{De}}^{p}{}_{}{}^{1}\frac{B}{p}\mathrm{\Phi }_{\mathrm{De}}^p`$.
###### Proof .
By the compatibility of the Frobenius action on the tangential morphism $`_p(1)=\pi _1^{p,\mathrm{rig}}(T_0^\times :\stackrel{}{01})(_p)\pi _1^{p,\mathrm{rig}}(X_0:\stackrel{}{01})(_p)`$ we have $`\varphi _p(x)=x^{\frac{1}{p}}`$. Similarly we have $`\varphi _p(y^{})=y^{\frac{1}{p}}`$. Since Frobenius action is compatible with the torsor structure we have $`\varphi _p(y)=\varphi _p(d^1y^{}d)=(d^1\varphi _p(d))^1d^1y^{\frac{1}{p}}d(d^1\varphi _p(d))=(d^1\varphi _p(d))^1y^{\frac{1}{p}}(d^1\varphi _p(d)).`$
By $`d^1\varphi _p(d)=d^1c\varphi _p(d^1c)^1`$ we have
$$\mathrm{\Phi }_{\mathrm{De}}^p=\mathrm{\Phi }_{\mathrm{KZ}}^p\varphi _p(\mathrm{\Phi }_{\mathrm{KZ}}^p)^1.$$
By Lemma 2.9 we get Theorem 2.8. โ
By expanding our explicit formula, we can express our $`p`$-adic MZVโs arising from Frobenius invariant path in terms of Deligneโs $`p`$-adic MZVโs describing Frobenius action very explicitly and vice versa. The following are the easiest examples.
###### Examples 2.10.
1. $`\zeta _p^{\mathrm{De}}(k)=(1\frac{1}{p^k})\zeta _p(k)`$ ($`k>1`$).
2. $`\zeta _p^{\mathrm{De}}(a,b)=(1\frac{1}{p^{a+b}})\zeta _p(a,b)(\frac{1}{p^b}\frac{1}{p^{a+b}})\zeta _p(a)\zeta _p(b)\underset{r=0}{\overset{a1}{}}(1)^r(\frac{1}{p^{ar}}\frac{1}{p^{a+b}})\left(\genfrac{}{}{0pt}{}{b1+r}{b1}\right)\zeta _p(ar)\zeta _p(b+r)(1)^a\underset{s=0}{\overset{b1}{}}(\frac{1}{p^{bs}}\frac{1}{p^{a+b}})\left(\genfrac{}{}{0pt}{}{a1+s}{a1}\right)\zeta _p(a+s)\zeta _p(bs)`$ ($`b>1`$).
We may say that Deligneโs $`p`$-adic MZVโs are not equal to but equivalent to ours.
Our third result is on Frobenius structures. We give a Frobenius structure on $`(๐ฑ_{\mathrm{KZ}}^{},_{\mathrm{KZ}}^{})`$. We introduce and discuss an overconvergent variant of the $`p`$-adic MPL which describes this structure.
###### Proposition 2.11.
The pro-object $`(๐ฑ_{\mathrm{KZ}}^{},_{\mathrm{KZ}}^{})`$ naturally admits a Frobenius structure
(2.10)
$$\varphi :F_p^{}(๐ฑ_{\mathrm{KZ}}^{},_{\mathrm{KZ}}^{})(๐ฑ_{\mathrm{KZ}}^{},_{\mathrm{KZ}}^{}).$$
###### Proof .
We just derive it from the Frobenius structure on $`๐_{\stackrel{}{01}}^{}`$ (Definition 1.6). The following arguments in the de Rham setting would help our understandings.
###### Lemma 2.12.
The pro-object $`(๐ฑ_{\mathrm{KZ}},_{\mathrm{KZ}})`$ is isomorphic to the pro-object $`(๐_{\stackrel{}{01}}^{\mathrm{DR}})^{}`$, the dual of the ind-object $`๐_{\stackrel{}{01}}^{\mathrm{DR}}`$ (Definition 1.4).
###### Proof .
By the categorical equivalence (2.3), $`(๐_\mathrm{\Gamma }^{\mathrm{DR}})^{}`$ corresponds to a natural representation $`\rho _\mathrm{\Gamma }`$ of $`Lie\pi _1^{\mathrm{DR}}(X:\mathrm{\Gamma })`$ on the universal enveloping algebra $`U\pi _1^{\mathrm{DR}}(X:\mathrm{\Gamma })`$ of $`Lie\pi _1^{\mathrm{DR}}(X:\mathrm{\Gamma })`$ (we note that $`U\pi _1^{\mathrm{DR}}(X:\mathrm{\Gamma })`$ is isomorphic to the dual $`H`$ of the coordinate ring $`\omega _\mathrm{\Gamma }(๐_\mathrm{\Gamma }^{\mathrm{DR}})`$ of $`\pi _1^{\mathrm{DR}}(X:\mathrm{\Gamma })`$ by \[A\] Theorem 2.5.3 because $`Lie\pi _1^{\mathrm{DR}}(X:\mathrm{\Gamma })`$ is the set of primitive elements in $`H`$ by loc.cit.ยง4.3). By (1.3) this representation $`\rho _\mathrm{\Gamma }`$ is induced from the left action of $`\pi _1^{\mathrm{DR}}(X:\mathrm{\Gamma })`$ on $`\pi _1^{\mathrm{DR}}(X:\mathrm{\Gamma },\mathrm{\Gamma })`$, thus $`U\pi _1^{\mathrm{DR}}(X:\mathrm{\Gamma })`$ is a free $`U\pi _1^{\mathrm{DR}}(X:\mathrm{\Gamma })`$-module of rank $`1`$ by $`\rho _\mathrm{\Gamma }`$. Therefore we say that the corresponding representations $`\rho _\mathrm{\Gamma }`$ and $`L`$ (Lemma 2.4) are equivalent, whence we have $`(๐ฑ_{\mathrm{KZ}},_{\mathrm{KZ}})(๐_\mathrm{\Gamma }^{\mathrm{DR}})^{}`$. On the other hand by the canonical isomorphism $`d_{\stackrel{}{01},\mathrm{\Gamma }}:\omega _{\stackrel{}{01}}\omega _\mathrm{\Gamma }`$ it is easy to see $`(๐_{\stackrel{}{01}}^{\mathrm{DR}})^{}(๐_\mathrm{\Gamma }^{\mathrm{DR}})^{}`$. Whence we get the claim. โ
By the categorical equivalence (2.2) we deduce $`(๐ฑ_{\mathrm{KZ}}^{},_{\mathrm{KZ}}^{})(๐_{\stackrel{}{01}}^{})^{}`$ from Lemma 2.12. Consequently we get the Frobenius structure on $`(๐ฑ_{\mathrm{KZ}}^{},_{\mathrm{KZ}}^{})`$ by transmitting from $`\varphi :๐_{\stackrel{}{01}}^{}F_p^{}๐_{\stackrel{}{01}}^{}`$. We get Proposition 2.11. โ
In the following section we introduce single-valued real-analytic MPL $`Li_{k_1,\mathrm{},k_m}^{}(z)`$. Here we introduce the following $`p`$-adic analogue $`Li_{k_1,\mathrm{},k_m}^{}(z)`$. From now on we fix a lift of Frobenius by $`\stackrel{~}{F}_p(t)=t^p`$. By $`๐ฑ_{\mathrm{KZ}}=๐ช_XA,B`$, we have $`๐ฑ_{\mathrm{KZ}}^{}=j^{}๐ช_{]\overline{X}_0[}A,B`$.
###### Definition 2.13.
The overconvergent $`p`$-adic MPL $`Li_{k_1,\mathrm{},k_m}^{}(z)j^{}๐ช_{]\overline{X}_0[}(]\overline{X}_0[)`$ is the coefficient of $`A^{k_m1}B\mathrm{}A^{k_11}B`$ in $`G_0^{}(z)`$ multiplied by $`(1)^m`$:
$$G_0^{}(z)=1+\mathrm{}+(1)^mLi_{k_1,\mathrm{},k_m}^{}(z)A^{k_m1}B\mathrm{}A^{k_11}B+\mathrm{}$$
where $`G_0^{}(z)๐ฑ_{\mathrm{KZ}}^{}(]\overline{X}_0[)`$ is the image of $`1\stackrel{~}{F}_p^{}๐ฑ_{\mathrm{KZ}}^{}(]\overline{X}_0[)=๐ฑ_{\mathrm{KZ}}^{}(]\overline{X}_0[)=j^{}๐ช_{]\overline{X}_0[}(]\overline{X}_0[)A,B`$ by (2.10).
We note that since $`G_0^{}(z)`$ is group-like we can recover and express general coefficients of $`G_0^{}(z)`$ in terms of $`Li_{k_1,\mathrm{},k_m}^{}(z)`$ as Note 2.6. Our choice of $`\stackrel{~}{F}_p`$ is a good lifting of the Frobenius of $`_{๐ฝ_p}^1\backslash \{1\}`$. So the overconvergent $`p`$-adic MPL is analytic on $`]0[`$ and $`]\mathrm{}[`$, i.e. $`Li_{k_1,\mathrm{},k_m}^{}(z)j^{}๐ช_{]\overline{X}_0[}(]\overline{X}_0[)`$ with $`j^{}:_{๐ฝ_p}^1\backslash \{1\}\overline{X}_0`$ while our $`p`$-adic MPL $`Li_{k_1,\mathrm{},k_m}(z)`$ is a Coleman function, i.e. belong to the ring $`A_{\mathrm{Col}}`$ of Coleman functions of $`X`$ which contains $`j^{}๐ช_{]\overline{X}_0[}(]\overline{X}_0[)`$. In Lemma 2.15 we see that overconvergent $`p`$-adic MPLโs have a different tannakian origin from ours, so that they are different from ours. But we have a close relationship between them in $`A_{\mathrm{Col}}A,B`$:
###### Theorem 2.14.
(2.11)
$$G_0^{}(A,B)(z)=G_0(A,B)(z)G_0(\frac{A}{p},\mathrm{\Phi }_{\mathrm{De}}^p(A,B)^1\frac{B}{p}\mathrm{\Phi }_{\mathrm{De}}^p(A,B))(z^p)^1.$$
###### Proof .
Let $`z`$ be a $`_p`$-valued point of $`X`$ whose modulo $`p`$ reduction $`z_0`$ lies on $`X(๐ฝ_p)`$ (If $`p=2`$, we need to enlarge the base field $`_p`$). We have a de Rham path $`d_z`$ in $`\pi _1^{\mathrm{DR}}(X:\stackrel{}{01},z)(_p)`$ while we also have a de Rham path $`\varphi _p(d_{z^p})`$ by $`\pi _1^{\mathrm{DR}}(X:\stackrel{}{01},z)\times _p\pi _1^{p,\mathrm{rig}}(X_0:\stackrel{}{01},z_0)\pi _1^{\mathrm{DR}}(X:\stackrel{}{01},z^p)\times _p`$ in Lemma 2.2. Identifying them we get a de Rham loop $`d_z^1\varphi _p(d_{z^p})`$ in $`\pi _1^{\mathrm{DR}}(X:\stackrel{}{01})(_p)`$.
###### Lemma 2.15.
$`i(d_z^1\varphi _p(d_{z^p}))=G_0^{}(z)`$.
###### Proof .
By the categorical equivalence (2.2), we identify $`๐ซ_{\stackrel{}{01}}^{\mathrm{DR}}\times _p`$, $`๐_{\stackrel{}{01}}^{\mathrm{DR}}_p`$ and $`๐ฑ_{\mathrm{KZ}}_p`$ with $`๐ซ_{\stackrel{}{01}}^{}`$, $`๐_{\stackrel{}{01}}^{}`$ and $`๐ฑ_{\mathrm{KZ}}^{}`$ respectively. Because $`๐ซ_{\stackrel{}{01}}^{\mathrm{DR}}=Spec๐_{\stackrel{}{01}}^{\mathrm{DR}}`$, we have a natural injection $`๐ซ_{\stackrel{}{01}}^{\mathrm{DR}}(๐_{\stackrel{}{01}}^{\mathrm{DR}})^{}`$ of pro-objects. By Lemma 2.12 we get a morphism $`๐ซ_{\stackrel{}{01}}^{\mathrm{DR}}๐ฑ_{\mathrm{KZ}}`$, whose associated morphism $`๐ซ_{\stackrel{}{01}}^{}๐ฑ_{\mathrm{KZ}}^{}`$ commutes with Frobenius actions, i.e.
$$\begin{array}{ccc}\stackrel{~}{F}_p^{}๐ซ_{\stackrel{}{01}}^{}& \stackrel{\varphi =F_p^1}{}& ๐ซ_{\stackrel{}{01}}^{}\\ & & & & \\ \stackrel{~}{F}_p^{}๐ฑ_{\mathrm{KZ}}^{}& \stackrel{\varphi =F_p^1}{}& ๐ฑ_{\mathrm{KZ}}^{}\end{array}$$
commutes. By taking the fiber $`\omega _z`$, we easily see that $`d_{z^p}\pi _1^{\mathrm{DR}}(X:\stackrel{}{01},z^p)(_p)=\omega _{z^p}(๐ซ_{\stackrel{}{01}}^{})=\omega _z(\stackrel{~}{F}_p^{}๐ซ_{\stackrel{}{01}}^{})`$ corresponds to $`1_{(z^p)}\omega _{z^p}(๐ฑ_{\mathrm{KZ}}^{})=\omega _z(\stackrel{~}{F}_p^{}๐ฑ_{\mathrm{KZ}}^{})`$. Hence the image $`\varphi _p(d_{z^p})\pi _1^{\mathrm{DR}}(X:\stackrel{}{01},z)(_p)=\omega _z(๐ซ_{\stackrel{}{01}}^{})`$ is mapped to $`G_0^{}(z)๐ฑ_{\mathrm{KZ},(z)}^{}=\omega _z(๐ฑ_{\mathrm{KZ}}_p)`$. The second row induces morphism $`\omega _{\stackrel{}{01}}(๐ซ_{\stackrel{}{01}}^{\mathrm{DR}})=\pi _1^{\mathrm{DR}}(X:\stackrel{}{01})(_p)\omega _{\stackrel{}{01}}(๐ฑ_{\mathrm{KZ}})_pA,B`$, so that the loop $`d_z^1\varphi _p(d_{z^p})`$ corresponds to $`G_0^{}(z)`$. โ
By $`d_z^1\varphi _p(d_{z^p})=d_z^1c_{z_0}\varphi _p(d_{z^p}^1c_{z_0})^1`$ we get the equality (2.11) whenever we fix each $`_p`$-point $`z`$. By letting $`z`$ varies for all $`_p`$-points on $`]z_0[`$ and restricting $`G_0^{}`$ and $`G_0`$ into $`]z_0[`$, we see that the equality (2.11) holds in $`_p[[t_{z_0}]]A,B`$ ($`t_{z_0}`$: a local parameter of $`]z_0[X(_p)`$), hence for $`A_{Col}|_{]z_0[}`$. By the uniqueness principle of Coleman functions \[C\], the equality (2.11) holds for the whole space. After all we get Theorem 2.14. โ
By Lemma 2.15 we may say that $`p`$-adic MZV ร la Deligne is โa special value โ of the overconvergent $`p`$-adic MPL. By Theorem 2.14 the overconvergent $`p`$-adic MPL is is described very explicitly as a combination of our $`p`$-adic MPLโs. The following are the easiest examples.
###### Examples 2.16.
1. $`Li_k^{}(z)=Li_k(z)\frac{1}{p^k}Li_k(z^p)=\underset{(n,p)=1}{}\frac{z^n}{n^k}`$.
2. $`Li_{a,b}^{}(z)=Li_{a,b}(z)\frac{1}{p^{a+b}}Li_{a,b}(z^p)(\frac{1}{p^b}\frac{1}{p^{a+b}})\zeta _p(a)Li_b(z^p)`$
$`\underset{r=0}{\overset{a1}{}}(1)^r\frac{1}{p^{ar}}\left(\genfrac{}{}{0pt}{}{b1+r}{r}\right)Li_{ar}(z^p)\{Li_{b+r}(z)\frac{1}{p^{b+r}}Li_{b+r}(z^p)\}`$
$`(1)^a\underset{s=0}{\overset{b1}{}}\left(\genfrac{}{}{0pt}{}{a1+s}{a1}\right)(\frac{1}{p^{bs}}\frac{1}{p^{a+b}})\zeta _p(a+s)Li_{bs}(z^p)`$.
We note that the coefficient of $`A`$ of $`G_0(z)`$ is $`\mathrm{log}^az`$ while in the overconvergent side the coefficient of $`G_0^{}(z)`$ is $`\mathrm{log}^az\frac{1}{p}\mathrm{log}^az^p=0`$ and $`Li_k^{}(z)`$ has been studied by Coleman \[C\]. Unver nearly got the same formula (2) in \[U2\] 5.16. In Example 2.28, we will give a Hodge analogue of the above two formulae. Our $`p`$-adic MPL has log poles around $`0`$ and $`\mathrm{}`$ however the overconvergent $`p`$-adic MPL does not. The formula (2.11) is an algorithm to erase log poles at $`\mathrm{}`$ of our $`p`$-adic MPL. We remark that the overconvergent $`p`$-adic MPL is rigid analytic on an open rigid analytic subspace $`]\overline{X}_0[`$ containing $`]^1\backslash \{1\}[`$.
###### Remark 2.17.
The $`p`$-adic KZ equation
$$dg=(\frac{A}{z}+\frac{B}{z1})gdz$$
does not have a solution with overconvergent function coefficients but have a solution with Coleman function coefficients (for instance $`G_0(z)`$). However if we modify it as follows
(2.12)
$$dg=\left(\frac{A}{z}+\frac{B}{z1}\right)gdzg\left(\frac{dz^p}{z^p}\frac{A}{p}+\frac{dz^p}{z^p1}\mathrm{\Phi }_{\mathrm{De}}^p(A,B)^1\frac{B}{p}\mathrm{\Phi }_{\mathrm{De}}^p(A,B)\right),$$
it has a solution with overconvergent function coefficients (for instance $`G_0^{}(z)`$). We note that the similar differential equation for the loop $`d_{z^p}^1F_p(d_z)`$ is introduced in \[U2, Y\] before.
### 2.2. Hodge side โ analogous results
In this side, we use a Hodge-type comparison isomorphism (Lemma 2.19) between a de Rham fundamental torsor and a Betti fundamental torsor for $`X`$. A tannakian interpretation of MPLโs and MZVโs are given in Proposition 2.20 and Proposition 2.21. By using the infinity Frobenius action on a Betti path space we introduce a new version of MPLโs (2.20). Its single-valuedness is shown in Theorem 2.27. A formula to express them in terms of usual (multi-valued) MPLโs is shown in Proposition 2.26.
###### Notation 2.18.
Let $`z`$ be a point in $`X()`$. We fix a topological path $`b_z\pi _1^{\mathrm{top}}(X():\stackrel{}{01},z)`$. By abuse of notations we also denote $`b_z`$ to be the corresponding Betti path in $`\pi _1(X():\stackrel{}{01},z)()`$. For $`z=\stackrel{}{10}`$ we consider a special path $`b\pi _1(X():\stackrel{}{01},\stackrel{}{10})()`$ which comes from a one-point set $`\pi _1^{\mathrm{top}}(X():\stackrel{}{01},\stackrel{}{10})`$.
The following Hodge type comparison isomorphism for fundamental torsors is our basics in this subsection.
###### Lemma 2.19.
Let $`z`$ be a $``$-valued point of $`X`$. Then we have
(2.13)
$$\pi _1^{\mathrm{DR}}(X:\stackrel{}{01},z)\times \pi _1^{\mathrm{Be}}(X:\stackrel{}{01},z)\times .$$
This follows from an equivalence of tannakian categories
(2.14)
$$๐ฉ๐^{\mathrm{DR}}(X_{})๐ฉ^{\mathrm{Be}}(X())$$
(see also \[De1\]). By identifying the above two fundamental torsors by (2.13) we get a de Rham loop $`d_z^1b_z`$, which lies in $`\pi _1^{\mathrm{DR}}(X:\overline{01})()`$.
###### Proposition 2.20.
Let $`z`$ be a $``$-valued point of $`X`$. The de Rham loop $`d_z^1b_z`$ corresponds to the special value at $`z`$ of the analytic continuation along $`b_z`$ of the fundamental solution $`G_0`$ <sup>13</sup><sup>13</sup>13 It means the fundamental solution of the KZ equation satisfying $`G_0(z)z^A`$ constructed by Drinfelโd \[Dr\]. of the KZ equation by the embedding $`i`$:
$$i(d_z^1b_z)=G_0(z).$$
###### Proof .
The proof is given in a similar way to the proof of Theorem 2.3. By \[De1\] ยง12.15-12.16 we have a tannakian interpretation of horizontal sections of $`๐ซ_{\stackrel{}{01}}^{\mathrm{DR}}`$, By Lemma 2.12 we have a tannakian interpretation of horizontal sections of $`(๐ฑ_{\mathrm{KZ}},_{\mathrm{KZ}})`$. Because both the constant term of $`G_0(t)`$ at $`t=0`$ and the constant term of $`t_0^A`$ at $`t_0=0`$ are equal to $`1`$, we have
$$b_{\stackrel{}{01},z}(t_0^A|_{\stackrel{}{01}})=\text{the analytic continuation of }G_0(t)\text{ along }b_z$$
which implies
(2.15)
$$d_{z,\mathrm{\Gamma }}b_{\stackrel{}{01},z}d_{\mathrm{\Gamma },\stackrel{}{01}}(1)=G_0(z).$$
By (2.4), (2.15) and $`d_{z,\mathrm{\Gamma }}=d_{\stackrel{}{01},\mathrm{\Gamma }}d_z^1`$, we get Proposition 2.20. โ
In each coefficient of $`G_0(z)`$ there appears MPL (0.1). In precise the coefficient of $`A^{k_m1}B\mathrm{}A^{k_11}B`$ in $`G_0(z)`$ is $`(1)^mLi_{k_1,\mathrm{},k_m}(z)`$:
$$G_0(z)=1+\mathrm{log}zALi_1(z)B+\mathrm{}+(1)^mLi_{k_1,\mathrm{},k_m}(z)A^{k_m1}B\mathrm{}A^{k_11}B+\mathrm{}.$$
By Proposition 2.20 we may say that the loop $`d_z^1b_z`$ is a tannakian origin of MPLโs. We remark that a choice of the Betti path $`b_z`$ corresponds to a choice of each branch of analytic continuation of MPL.
For $`z=\stackrel{}{10}`$ similarly we have
###### Proposition 2.21.
The de Rham loop $`d^1b\pi _1^{\mathrm{DR}}(X:\stackrel{}{01})()`$ corresponds to the Drinfelโd associator $`\mathrm{\Phi }_{\mathrm{KZ}}A,B`$ by the embedding $`i`$, i.e.
$$i(d^1b)=\mathrm{\Phi }_{\mathrm{KZ}}.$$
###### Proof .
The Drinfelโd associator $`\mathrm{\Phi }_{\mathrm{KZ}}`$ is the series constructed in \[De1\], which is also equal to the limit value
$$\underset{ฯต0}{lim}\mathrm{exp}\left(\mathrm{log}ฯตB\right)G_0(1ฯต)A,B$$
as (2.6). In other word it is the constant term of $`G_0(z)`$ at $`z=1`$. Then by a similar argument to the proof of Theorem 2.5 we have $`b_{\stackrel{}{01},\stackrel{}{10}}(t_0^A|_{\stackrel{}{01}})=\mathrm{\Phi }_{\mathrm{KZ}}t_1^B|_{\stackrel{}{10}}`$, whence
(2.16)
$$d_{\stackrel{}{10},\mathrm{\Gamma }}bd_{\mathrm{\Gamma },\stackrel{}{01}}(1)=\mathrm{\Phi }_{\mathrm{KZ}}.$$
By (2.4), (2.16) and $`d_{\stackrel{}{10},\mathrm{\Gamma }}=d_{\stackrel{}{01},\mathrm{\Gamma }}d^1`$, we prove Proposition 2.21. โ
In each coefficient of the Drinfelโd associator $`\mathrm{\Phi }_{\mathrm{KZ}}`$ there appears MZV (0.2). In precise the coefficient of $`A^{k_m1}B\mathrm{}A^{k_11}B`$ in $`\mathrm{\Phi }_{\mathrm{KZ}}`$ is $`(1)^m\zeta (k_1,\mathrm{},k_m)`$:
$$\mathrm{\Phi }_{\mathrm{KZ}}=1+\mathrm{}+(1)^m\zeta (k_1,\mathrm{},k_m)A^{k_m1}B\mathrm{}A^{k_11}B+\mathrm{}.$$
By Proposition 2.21 we may say that the loop $`d^1b`$ is a tannakian origin of MZVโs.
The above proposition enables us to calculate the period map
(2.17)
$$p:\pi _1^{\mathrm{Be}}(X():\stackrel{}{01})()\pi _1^{\mathrm{DR}}(X:\stackrel{}{01})().$$
In order to do that, as in ยง2.1 we fix generator of $`\pi _1^{\mathrm{Be}}(X():\stackrel{}{01})()`$ and its parameterization.
###### Notation 2.22.
We denote the loop in $`\pi _1^{\mathrm{Be}}(X():\stackrel{}{01})()`$ (resp. in $`\pi _1^{\mathrm{Be}}(X():\stackrel{}{10})()`$) which comes from a generator in $`\pi _1^{\mathrm{top}}(T_0^\times ():\stackrel{}{01})`$ (resp. in $`\pi _1^{\mathrm{top}}(T_1^\times ():\stackrel{}{10})`$) going around the point $`0T_0()`$ (resp. $`0T_1()`$) counterclockwisely by $`x_{\mathrm{Be}}`$ (resp. $`y_{\mathrm{Be}}^{}`$). Put $`y_{\mathrm{Be}}=b^1y_{\mathrm{Be}}^{}b`$ (for $`b`$ see Notation 2.18). Easily we see that two loops $`x_{\mathrm{Be}}`$ and $`y_{\mathrm{Be}}`$ are free generators of the pro-algebraic group $`\pi _1^{\mathrm{Be}}(X():\stackrel{}{01})()`$. Wojtkowiak \[W\] consider the embedding
(2.18)
$$j:\pi _1^{\mathrm{Be}}(X():\stackrel{}{01})()A,B$$
which sends $`x_{\mathrm{Be}}e^A:=1+\frac{A}{1!}+\frac{A^2}{2!}+\mathrm{}`$ and $`y_{\mathrm{Be}}e^B:=1+\frac{B}{1!}+\frac{B^2}{2!}+\mathrm{}`$. We note that this morphism is isomorphic to the embedding of $`\pi _1^{\mathrm{Be}}(X():\stackrel{}{01})()`$ into the universal embedding algebra of $`Lie\pi _1^{\mathrm{Be}}(X():\stackrel{}{01})()`$.
We note that this embedding $`j`$ is not compatible with $`i`$ in (2.1) under the Hodge comparison isomorphism (2.13). As is similar to the proof of Theorem 2.8, we extend the period map (2.17) to
$$p:A,BA,B$$
via the embedding $`i`$ (2.1) and $`j`$ (2.18).
###### Lemma 2.23.
$`p(A)=2\pi iA`$, $`p(B)=\mathrm{\Phi }_{\mathrm{KZ}}^1(2\pi iB)\mathrm{\Phi }_{\mathrm{KZ}}`$.
###### Proof .
By the compatibility of the period map on the tangential morphisms $`_{\mathrm{Be}}(1)\stackrel{exp}{=}\pi _1^{\mathrm{Be}}(T_0^\times ():\stackrel{}{01})\pi _1^{\mathrm{Be}}(X():\stackrel{}{01})`$ and $`_{\mathrm{DR}}(1)\stackrel{exp}{=}\pi _1^{\mathrm{DR}}(T_0^\times :\stackrel{}{01})\pi _1^{\mathrm{DR}}(X:\stackrel{}{01})`$, we have $`p(x_{\mathrm{Be}})=x^{2\pi i}`$. Similarly we have $`p(y_{\mathrm{Be}}^{})=y^{2\pi i}`$. Since the period map is compatible with the torsor structure we have $`p(y_{\mathrm{Be}})=p(b^1y_{\mathrm{Be}}^{}b)=b^1y^{2\pi i}b=(d^1b)^1y^{2\pi i}d^1b`$. By Proposition 2.21, we get the claim. โ
Our next work is on a description of a Frobenius action. To begin with we consider the infinity Frobenius action $`\varphi _{\mathrm{}}`$ on $`\pi _1^{\mathrm{DR}}(X:\stackrel{}{01},\stackrel{}{10})()`$ by transmitting the infinity Frobenius action on $`\pi _1^{\mathrm{Be}}(X():\stackrel{}{01},\stackrel{}{10})()`$ by Lemma 2.19. It gives a new path $`\varphi _{\mathrm{}}(d)`$ in $`\pi _1^{\mathrm{DR}}(X:\stackrel{}{01},\stackrel{}{10})()`$. We denote $`\mathrm{\Phi }_{\mathrm{KZ}}^{}`$ to be the series in $`A,B`$ which corresponds to $`d^1\varphi _{\mathrm{}}(d)`$:
$$i(d^1\varphi _{\mathrm{}}(d))=\mathrm{\Phi }_{\mathrm{KZ}}^{}.$$
It is a Hodge counterpart of the $`p`$-adic Deligne associator $`\mathrm{\Phi }_{\mathrm{De}}^p`$. By same arguments to the proof of Lemma 2.9 and Theorem 2.8 the extension of the infinity Frobenius action on $`\pi _1^{\mathrm{DR}}(X:\stackrel{}{01})()`$ into $`A,B`$ by $`i`$ is described as follows:
###### Lemma 2.24.
$`\varphi _{\mathrm{}}(A)=A`$, $`\varphi _{\mathrm{}}(B)=\mathrm{\Phi }_{\mathrm{KZ}}^{}{}_{}{}^{1}(B)\mathrm{\Phi }_{\mathrm{KZ}}^{}`$.
This gives the following formula analogous to (2.8).
###### Lemma 2.25.
$$\mathrm{\Phi }_{\mathrm{KZ}}(A,B)=\mathrm{\Phi }_{\mathrm{KZ}}^{}(A,B)\mathrm{\Phi }_{\mathrm{KZ}}(A,\mathrm{\Phi }_{\mathrm{KZ}}^{}(A,B)^1(B)\mathrm{\Phi }_{\mathrm{KZ}}^{}(A,B)).$$
This formula looks nasty. But according to the convention in \[F0\] of the Grothendieck-Teichmรผller group \[Dr\] we rewrite simpler
$$(1,\mathrm{\Phi }_{\mathrm{KZ}}^{})=(1,\mathrm{\Phi }_{\mathrm{KZ}})(1,\mathrm{\Phi }_{\mathrm{KZ}})^1.$$
Let $`z`$ be a $``$-valued point of $`X`$. By a same argument to above, we get a de Rham loop $`d_z^1\varphi _{\mathrm{}}(d_{\overline{z}})\pi _1^{\mathrm{DR}}(X:\overline{01})()`$. Put $`G_0^{}(z):=i(d_z^1\varphi _{\mathrm{}}(d_{\overline{z}}))`$. The following is an analogue to Theorem 2.14.
###### Proposition 2.26.
(2.19)
$$G_0^{}(A,B)(z)=G_0(A,B)(z)G_0(A,\mathrm{\Phi }_{\mathrm{KZ}}^{}(A,B)^1(B)\mathrm{\Phi }_{\mathrm{KZ}}^{}(A,B))(\overline{z})^1.$$
###### Proof .
By $`d_z^1\varphi _{\mathrm{}}(d_{\overline{z}})=d_z^1b_z\varphi _{\mathrm{}}(d_{\overline{z}}^1\overline{b_z})^1`$ and Lemma 2.24 we get the equality. โ
We introduce a new version of MPL $`Li_{k_1,\mathrm{},k_m}^{}(z)`$ to be the coefficient of $`A^{k_m1}B`$ $`\mathrm{}`$ $`A^{k_11}B`$ in $`G_0^{}(z)`$ multiplied by $`(1)^m`$:
(2.20)
$$G^{}(z)=:1+\mathrm{}+(1)^mLi_{k_1,\mathrm{},k_m}^{}(z)A^{k_m1}B\mathrm{}A^{k_11}B+\mathrm{}.$$
We note that since $`G^{}(z)`$ is group-like we can recover and express general coefficients of $`G^{}(z)`$ in terms of $`Li_{k_1,\mathrm{},k_m}^{}(z)`$ and $`\mathrm{log}|z|`$ as Note 2.6.
###### Theorem 2.27.
The coefficient $`Li_{k_1,\mathrm{},k_m}^{}(z)`$ is single-valued and real-analytic on $`X()=^1()\backslash \{0,1,\mathrm{}\}`$.
###### Proof .
The real-analyticity follows from (2.19) because both $`G_0(z)`$ and $`G_0(\overline{z})`$ are real-analytic. The single-valuedness is immediate by definition. โ
The single-valued MPLโs are not equal to usual MPLโs because they have a different tannakian origins. But by Proposition 2.26 the single-valued MPL is described very explicitly as a combination of usual MPLโs. The following are the easiest examples:
###### Examples 2.28.
1. $`Li_k^{}(z)=Li_k(z)\underset{a=0}{\overset{k1}{}}(1)^{ka}\frac{(\mathrm{log}|z|^2)^a}{a!}Li_{ka}(\overline{z})`$
2. $`Li_{a,b}^{}(z)=Li_{a,b}(z)\underset{r=0}{\overset{a1}{}}\underset{s=0}{\overset{r}{}}[(1)^{a+r+s}\left(\genfrac{}{}{0pt}{}{b1+s}{s}\right)\frac{(\mathrm{log}|z|^2)^{rs}}{(rs)!}Li_{ar}(\overline{z})`$
$`\{Li_{b+s}(z)\underset{w=0}{\overset{b+s1}{}}(1)^{b+s+w}\frac{(\mathrm{log}|z|^2)^w}{w!}Li_{b+sw}(\overline{z})\}]_{u=0}^{b1}\frac{(\mathrm{log}|z|^2)^u}{u!}`$
$`[(1)^{a+b+u}Li_{a,bu}(\overline{z})+\{(1)^{b+u}(1)^{a+b+u}\}\zeta (a)Li_{bu}(\overline{z})`$
$`+\underset{v=0}{\overset{bu1}{}}\{(1)^{a+b+u+v}(1)^{b+u}\}\left(\genfrac{}{}{0pt}{}{a+v1}{a1}\right)\zeta (a+v)Li_{buv}(\overline{z})]`$.
The coefficient of $`A`$ of $`G_0(z)`$ is $`\mathrm{log}z`$ while the coefficient of $`G_0^{}(z)`$ is $`\mathrm{log}\overline{z}+\mathrm{log}z=\mathrm{log}|z|^2`$. We also remark that Zhao \[Zh\] constructed a single-valued version of several variable MPL $`Li_{k_1,\mathrm{},k_m}(z_1,\mathrm{},z_m)=\underset{0<n_1<\mathrm{}<n_m}{}\frac{z_1^{n_1}\mathrm{}z_m^{n_m}}{n_1^{k_1}\mathrm{}n_m^{k_m}}`$.
###### Remark 2.29.
In \[Za\] Zagier studied another variant of polylogarithm
$$P_k(z)=\mathrm{}_k\underset{a=0}{\overset{k1}{}}\frac{B_a}{a!}(\mathrm{log}|z|^2)^aLi_{ka}(z)$$
where $`\mathrm{}_k`$ denote $`\mathrm{}`$ or $`\mathrm{}`$ depending whether $`k`$ is odd or even and $`\frac{t}{e^t1}=\underset{n=0}{\overset{\mathrm{}}{}}B_n\frac{t^n}{n!}`$. He showed that its single-valuedness and real-analyticity. Beilinson and Deligne \[BD\] give its interpretations in terms of variants of mixed Hodge structures. Their computations \[BD\] ยง1.5 says that $`P_k(z)`$ appears as the coefficient of $`A^{k1}B`$ of $`\mathrm{log}G_0^{}(z)`$ multiplied by $`\frac{1}{2}`$, whence the tannakian origin of $`P_k(z)`$ is also $`d_z^1\varphi _{\mathrm{}}(d_{\overline{z}})`$. We notice that $`P_k(z)`$ and $`Li_k^{}(z)`$ are related with each other by
(2.21)
$$P_k(z)=\frac{1}{2}\underset{i=0}{\overset{k1}{}}\frac{B_i}{i!}\{\mathrm{log}|z|^2\}^iLi_{ki}^{}(z)$$
and
$$Li_k^{}(z)=2\underset{i=0}{\overset{k1}{}}\frac{(\mathrm{log}|z|^2)^i}{(i+1)!}P_{ki}(z).$$
Our MPL has monodromies around $`0`$, $`1`$ and $`\mathrm{}`$ however the variant does not. The formula (2.19) gives an algorithm to erase monodromies of our MPL at $`0`$, $`1`$ and $`\mathrm{}`$.
###### Remark 2.30.
The KZ equation
$$dg=(\frac{A}{z}+\frac{B}{z1})gdz$$
have a solution with multi-valued complex analytic function coefficients (for instance $`G_0(z)`$) whereas the following modification analogous to (2.12)
$$dg=\left(\frac{A}{z}+\frac{B}{z1}\right)gdzg\left(\frac{d\overline{z}}{\overline{z}}(A)+\frac{d\overline{z}}{\overline{z}1}\mathrm{\Phi }_{\mathrm{KZ}}^{}(A,B)^1(B)\mathrm{\Phi }_{\mathrm{KZ}}^{}(A,B)\right)$$
does have a solution with single-valued real-analytic function coefficients(for instance $`G_0^{}(z)`$).
### 2.3. Artin side โ reviews
In this side we will discuss $`l`$-adic รฉtale analogues of MPLโs and MZVโs which we call $`l`$-adic MPLโs and $`l`$-adic multiple Soulรฉ elements respectively and whose definitions are imitations of Wojtkowiakโs $`l`$-adic polylogarithms \[W\].
Let $`l`$ be a prime (we do not assume whether $`l`$ is equal to the previous prime $`p`$ in ยง2.1 or not). Let $`\overline{}`$ be the algebraic closure of $``$ in $``$. The following Artin-type comparison isomorphism take place of Lemma 2.2 and Lemma 2.19.
###### Lemma 2.31.
Suppose that $`x`$ and $`y`$ are $`\overline{}`$-valued points of $`X`$ or $`T_s^\times `$ ($`s\{0,1,\mathrm{}\}`$) then there exists an isomorphism
$$\pi _1^{\mathrm{Be}}(X():x,y)\times _{}_l\pi _1^{l,\text{รฉt}}(X_\overline{}:x,y).$$
It follows from the comparison isomorphism \[G\] XII Corollary 5.2 (see also \[De1\] ยง13.11).
Let $`i:\pi _1^{l,\text{รฉt}}(X_{\overline{_l}}:\stackrel{}{01})(_l)_lA,B`$ be the embedding associated with $`j`$ (2.18). Suppose that $`zX()`$ and $`b_z\pi _1^{\mathrm{top}}(X():\stackrel{}{01},z)`$. It determines a Betti path in $`\pi _1^{\mathrm{Be}}(X():\stackrel{}{01},z)()`$ and we regard this to be an $`l`$-adic รฉtale path (which we denote $`b_z`$ by abuse of notation) in $`\pi _1^{l,\text{รฉt}}(X_\overline{}:\stackrel{}{01},z)(_l)`$ Lemma 2.31. Suppose that $`\sigma Gal(\overline{}/)`$. We obtain another path $`\sigma (b_z)`$ by the Galois group action on $`\pi _1^{l,\text{รฉt}}(X_\overline{}:\stackrel{}{01},z)(_l)`$. Combining $`\sigma (b_z)`$ with the inverse $`b_z^1`$, we get an $`l`$-adic รฉtale loop $`b_z^1\sigma (b_z)\pi _1^{l,\text{รฉt}}(X_\overline{}:\stackrel{}{01})(_l)`$. Wojtkowiak \[W\] introduced $`l`$-adic polylogarithm $`\mathrm{}_k^z(\sigma )`$ ($`k1`$) to be the coefficient of $`A^{k1}B`$ of $`\mathrm{log}j(b_z^1\sigma (b_z))`$ multiplied by $`(1)^{k1}`$ (cf. \[NW\]). In accordance with expressions in ยง2.1 and ยง2.2, we consider the following $`l`$-adic รฉtale analogues of MPLโs and MZVโs.
###### Definition 2.32.
Let $`m,k_1,\mathrm{},k_m1`$, $`zX()`$ and $`\sigma Gal(\overline{}/)`$. The $`l`$-adic multiple polylogarithm (shortly the $`l`$-adic MPL) $`Li_{k_1,\mathrm{},k_m}^l(z)(\sigma )`$ denotes the coefficient of $`A^{k_m1}B\mathrm{}A^{k_11}B`$ in $`j(b_z^1\sigma (b_z))`$ multiplied by $`(1)^m`$ and $`l`$-adic multiple Soulรฉ element (shortly the $`l`$-adic MSE) $`\zeta _l(k_1,\mathrm{},k_m)(\sigma )`$ ($`k_m>1`$) denotes the coefficient of $`A^{k_m1}B\mathrm{}A^{k_11}B`$ in the $`l`$-adic Ihara associator $`\mathrm{\Phi }_\sigma ^l(A,B)=j(b^1\sigma (b))`$ multiplied by $`(1)^m`$.
By making a same computations similar to \[BD\] ยง1.5, we have
(2.22)
$$\mathrm{}_k^z(\sigma )=(1)^k\underset{i=0}{\overset{k1}{}}\frac{B_i}{i!}\{\rho _z(\sigma )\}^iLi_{ki}^l(z)(\sigma ).$$
Here $`\rho _z(\sigma )`$ is the generalized $`l`$-adic Kummer 1-cocycle in \[NW\] Definition 3. By comparing it with (2.21), we may say that $`Li_k^l(z)(\sigma )`$ is a $`l`$-adic analogue of the usual (Leibnitzโs) polylogarithm $`Li_k(z)=\underset{n=1}{\overset{\mathrm{}}{}}\frac{z^n}{n!}`$ and Wojtkowiakโs $`\mathrm{}_k^z(\sigma )`$ is a $`l`$-adic analogue of the Beilinson-Deligneโs polylogarithm $`D_k(z)`$. By asking for $`p`$-adic Hodge theory, we may give new $`p`$-adic MPLโs valued on the ring of $`B_{\mathrm{crys}}`$ of Fontaineโs $`p`$-adic periods, which might unify $`p`$-adic MPLโs in ยง2.1 and those above (for $`l=p`$).
The following may explain why we call the multiple Soulรฉ elements.
###### Examples 2.33 (\[I\] ยง6.3.Theorem).
1. When $`k`$ is odd, $`\zeta _l(k)(\sigma )=\frac{\kappa _k^{(l)}(\sigma )}{(1l^{k1})\dot{(}k1)!}`$. Here $`\kappa _k^{(l)}(\sigma )`$ is $`l`$-adic Soulรฉ element (see \[I\] ยง6.2 for its definition).
2. When $`k`$ is even, $`\zeta _l(k)(\sigma )=\frac{B_k}{2(k!)}\{1\chi _l(\sigma )^k\}`$. Here $`B_k`$ is the Bernoulli number defined by $`\frac{t}{e^t1}=\underset{n=0}{\overset{\mathrm{}}{}}B_n\frac{t^n}{n!}`$ and $`\chi _l(\sigma )`$ is the $`l`$-adic cyclotomic character.
The author expects that MSE is a kind of โmultiple Euler systemโ which helps our understanding of multiple zeta functions.
We end this section by giving a brief observation on $`l`$-adic behavior of $`\zeta _l(k)(\sigma )`$ and $`p`$-adic behavior of $`\zeta _p(k)`$ with respect to weights.
###### Remark 2.34.
Let $`\sigma Gal(\overline{}/)`$. By the expression $`\kappa _k^l(\sigma )=_{_l^\times }x^{k1}๐\mu ^\sigma `$ ($`\mu ^\sigma `$: the measure associated with the Kummer distribution \[NW\] ยง2) we see that $`(k1)!(1l^{k1})\zeta _l(k)(\sigma )`$ admits a nice $`l`$-adic behavior with respect to $`k`$, that is, for $`kk^{}mod(l1)l^M`$ ($`M`$) we have
(2.23)
$$(k1)!(1l^{k1})\zeta _l(k)(\sigma )(k^{}1)!(1l^{k^{}1})\zeta _l(k^{})(\sigma )modl^{M+1}_l.$$
On the other hand in the $`p`$-adic setting, by $`L_p(1k,\omega ^{1k})=(1p^{k1})\frac{B_k}{k}=(1p^{k1})\zeta (1k)`$ and $`L_p(k,\omega ^{1k})=(1\frac{1}{p^k})\zeta _p(k)`$ we see that $`(1\frac{1}{p^k})\zeta _p(k)`$ admits a nice $`p`$-adic behavior with respect to $`k`$, that is, for $`kk^{}mod(p1)p^M`$ ($`M`$) we have
(2.24)
$$(1\frac{1}{p^k})\zeta _p(k)(1\frac{1}{p^k^{}})\zeta _p(k^{})modp^{M+1}_p.$$
The author is not sure whether we have โmultiple analoguesโ of (2.23) and (2.24). He is grateful to the referee who suggested him that Colemanโs integral formula (\[C\] Lemma 7.2) $`Li_k^{}(z)=_{_p^\times }x^k๐\mu _z(x)`$ for $`z_p\backslash \{|z1|_p<1\}`$ might help look for such relations. Here $`\mu _z`$ is a measure on $`_p`$ given by $`\mu _z(a+p^n_p)=\frac{z^a}{1z^{p^n}}`$ with $`n`$ and $`0a<p^n`$.
## 3. Motivic views
This section is complementary. We will consider three boxes related to the algebra generated by $`p`$-adic MZVโs. In ยง3.1 our special box is Drinfelโdโs \[Dr\] pro-algebraic bi-torsor where the Grothendieck-Teichmรผller pro-algebraic groups act. In ยง3.2 our special box is Racinetโs \[R\] pro-algebraic torsor defined by double shuffle relations. In ยง3.3 our special box is pro-algebraic bi-torsor of Deligne-Goncharovโs \[DG\] motivic Galois group. The relationship between these three box will be discussed.
### 3.1. Drinfelโdโs context
In this subsection we will add the Berthelot-Ogus part into the story of \[F0\] owing Unverโs result \[U1\] and discuss special elements in Iharaโs stable derivation algebra which are $`p`$-adic analogues of the Drinfelโdโs element $`\stackrel{~}{\psi }`$ in \[Dr\] Proposition 6.3.
In his celebrated paper \[Dr\], Drinfelโd introduced the Grothendieck-Teichmรผller torsor, a triple $`(\underset{ยฏ}{GRT}_1,\underset{ยฏ}{M}_1,\underset{ยฏ}{GT}_1)`$, where $`\underset{ยฏ}{M}_1`$ is the pro-algebraic bi-torsor with a left (resp. right) action of the Grothendieck-Teichmรผller pro-algebraic group $`\underset{ยฏ}{GRT}_1`$ (resp. $`\underset{ยฏ}{GT}_1`$) (for their presentations see also \[F0\]ยง2). In \[F0\]ยง4 we explained that for a prime $`l`$ the $`l`$-adic Galois image pro-algebraic group $`\underset{ยฏ}{Gal}__l^{(l)}`$ (which is defined in loc.cit. to be the Zariski closure of the image of the Galois group $`Gal(\overline{}/(\mu _l^{\mathrm{}}))`$-action in $`\underset{ยฏ}{Aut}\pi _1^{l,\text{รฉt}}(X_\overline{}:\stackrel{}{01})(_l)`$) lies in $`\underset{ยฏ}{GT}_1\times _l`$, i.e. we have
(3.1)
$$\mathrm{\Phi }__l^{(l)}:\underset{ยฏ}{Gal}__l^{(l)}\underset{ยฏ}{GT}_1\times _l,$$
whereas in \[F0\] ยง3 we also saw that the spectrum of the $``$-algebra $`()`$ generated by $`(2\pi i)^{(k_1+\mathrm{}+k_m)}\zeta (k_1,\mathrm{},k_m)`$ ($`m,k_i`$, $`k_m>1`$) is embedded into $`\underset{ยฏ}{M}_1`$, i.e. we have
(3.2)
$$\mathrm{\Phi }_{\mathrm{Hod}}:Spec\underset{ยฏ}{M}_1.$$
We note that (3.1) and (3.2) are expected to be isomorphism \[F0\]. In \[F1\] we announced to make up still lacking Berthelot-Ogus part in our next upcoming paper โ$`p`$-adic multiple zeta values IIIโ; we planned to construct an embedding
(3.3)
$$\mathrm{\Phi }_{\mathrm{crys}}^{(p)}:Spec^{(p)}\underset{ยฏ}{GRT}_1$$
where $`^{(p)}`$ is a subalgebra of $`_p`$ generated by all $`p`$-adic MZVโs. But we would cancel it and give an explanation below instead because a recent Unverโs paper \[U1\] nearly achieved it.
###### Proposition 3.1.
The $`p`$-adic Drinfelโd associator $`\mathrm{\Phi }_{\mathrm{KZ}}^p(A,B)`$ satisfies the following defining equations of $`\underset{ยฏ}{GRT}_1`$:
$$\{\begin{array}{cc}(0)\mathrm{\Phi }_{\mathrm{KZ}}^pexp[๐__p^{},๐__p^{}]\hfill & \\ (\mathrm{i})\mathrm{\Phi }_{\mathrm{KZ}}^p(A,B)\mathrm{\Phi }_{\mathrm{KZ}}^p(B,A)=1\hfill & \\ (\mathrm{ii})\mathrm{\Phi }_{\mathrm{KZ}}^p(C,A)\mathrm{\Phi }_{\mathrm{KZ}}^p(B,C)\mathrm{\Phi }_{\mathrm{KZ}}^p(A,B)=1\text{for}A+B+C=0\hfill & \\ (\mathrm{iii})\mathrm{\Phi }_{\mathrm{KZ}}^p(X_{1,2},X_{2,3})\mathrm{\Phi }_{\mathrm{KZ}}^p(X_{3,4},X_{4,5})\mathrm{\Phi }_{\mathrm{KZ}}^p(X_{5,1},X_{1,2})\hfill & \\ \mathrm{\Phi }_{\mathrm{KZ}}^p(X_{2,3},X_{3,4})\mathrm{\Phi }_{\mathrm{KZ}}^p(X_{4,5},X_{5,1})=1\text{in}\underset{ยฏ}{U๐_{}^{(5)}}(_p).\hfill & \end{array}$$
Here $`๐__p^{}`$ stands for the completed free Lie algebra generated by $`A`$ and $`B`$ and $`\underset{ยฏ}{U๐_{}^{(5)}}`$ stands for the universal enveloping algebra of the pure sphere 5-braid Lie algebra $`๐^{(5)}`$ with standard generators $`X_{i,j}`$ ($`1i,j5`$).
###### Proof .
By $`(p,\mathrm{\Phi }_{\mathrm{De}}^p)=(p,1)(1,\mathrm{\Phi }_{\mathrm{De}}^p(pA,pB))`$ and $`(p,\mathrm{\Phi }_{\mathrm{KZ}}^p)=(p,1)(1,\mathrm{\Phi }_{\mathrm{KZ}}^p(pA,pB))`$ in our convention of $`\underset{ยฏ}{GRT}_1`$, (2.9) is rewritten into
(3.4)
$$(1,\mathrm{\Phi }_{\mathrm{De}}^p(pA,pB))=(1,\mathrm{\Phi }_{\mathrm{KZ}}^p(pA,pB))(1,\mathrm{\Phi }_{\mathrm{KZ}}^p(A,B))^1.$$
In \[U1\] Unver showed that the $`p`$-adic Deligne associator $`\mathrm{\Phi }_{\mathrm{De}}^p(A,B)`$ satisfies the above relations, in other words $`(1,\mathrm{\Phi }_{\mathrm{De}}^p(A,B))\underset{ยฏ}{GRT}_1(_p)`$. Thus $`(1,\mathrm{\Phi }_{\mathrm{De}}^p(p^nA,p^nB))\underset{ยฏ}{GRT}_1(_p)`$ for $`n`$. Put
$$(1,a_n)=(1,\mathrm{\Phi }_{\mathrm{De}}^p(pA,pB))^1(1,\mathrm{\Phi }_{\mathrm{De}}^p(p^2A,p^2B))^1\mathrm{}(1,\mathrm{\Phi }_{\mathrm{De}}^p(p^nA,p^nB))^1.$$
Because $`\underset{ยฏ}{GRT}_1`$ is closed under the multiplication, $`(1,a_n)`$ must belongs to $`\underset{ยฏ}{GRT}_1(_p)`$. By (3.4),
$$(1,\mathrm{\Phi }_{\mathrm{De}}^p(p^nA,p^nB))=(1,\mathrm{\Phi }_{\mathrm{KZ}}^p(p^nA,p^nB))(1,\mathrm{\Phi }_{\mathrm{KZ}}^p(p^{n1}A,p^{n1}B))^1.$$
So we have $`(1,a_n)=(1,\mathrm{\Phi }_{\mathrm{KZ}}^p(A,B))(1,\mathrm{\Phi }_{\mathrm{KZ}}^p(p^nA,p^nB))^1`$. As $`n`$ goes $`\mathrm{}`$, $`\mathrm{\Phi }_{\mathrm{De}}^p(p^nA,p^nB)`$ converges to $`1`$ in $`_pA,B`$ which is a projective limit of finite dimensional topological $`_p`$-vector space. Hence the series $`a_n`$ ($`n1`$) should converge to $`\mathrm{\Phi }_{\mathrm{KZ}}^p(A,B)`$ in $`_pA,B`$. By the locally-compactness of the topological group $`\underset{ยฏ}{GRT}_1(_p)`$ the limit $`(1,\mathrm{\Phi }_{\mathrm{KZ}}^p(A,B))`$ must lie on $`\underset{ยฏ}{GRT}_1(_p)`$, whence we get the claim. โ
By Proposition 3.1 we get a morphism (3.3). The injectivity is easy to check.
Since there appear $`p`$-adic MZVโs in each coefficient of the $`p`$-adic Drinfelโd associator $`\mathrm{\Phi }_{\mathrm{KZ}}^p`$, we get algebraic relations among them. It is remarkable that we get $`\zeta _p(2k)=0`$ ($`k1`$) from (0), (i) and (ii) (proved in the same way to \[De1\]ยง18.16), that is, we have a geometric proof for $`L_p(2k,\omega ^{12k})=0`$.
In \[F0\] ยง6 by discussing a weight filtration on our torsor, we associated (3.1) and (3.2) with two surjections between their coordinate rings, $`Gr\mathrm{\Phi }__l^{(l)\mathrm{}}:๐ช(\underset{ยฏ}{GRT}_1)_{}\widehat{}_lGr_{}^W๐ช(\underset{ยฏ}{Gal}__l^{(l)})`$ and $`Gr\mathrm{\Phi }_{\mathrm{Hod}}^{\mathrm{}}:๐ช(\underset{ยฏ}{GRT}_1)_{}\left(Z_{}/(\pi ^2)Z_{}\right)`$ where $`Z_{}`$ is a graded $``$-algebra whose degree $`w`$-component ($`w1`$) is a $``$-vector space $`Z_w()`$ generated by all MZVโs with weight $`w`$ and $`Z_0=`$. By (3.3), we also obtain
$$\mathrm{\Phi }_{\mathrm{crys}}^{(p)\mathrm{}}:๐ช(\underset{ยฏ}{GRT}_1)_{}Z_{}^{(p)}$$
where $`Z_{}^{(p)}`$ is a graded $``$-algebra whose degree $`w`$-component is a $``$-vector space $`Z_w^{(p)}(_p)`$ generated by all $`p`$-adic MZVโs with weight $`w`$ and $`Z_0^{(p)}=`$.
###### Remark 3.2.
1. In \[F0\]ยง6.6.1, we introduced an invertible element $`\mathrm{\Phi }_{GRT}`$ in $`๐ช(\underset{ยฏ}{GRT}_1)A,B`$ and showed $`Gr\mathrm{\Phi }_{\mathrm{Hod}}^{\mathrm{}}(\mathrm{\Phi }_{GRT})=\mathrm{\Phi }_{\mathrm{KZ}}mod\pi ^2`$ and $`Gr\mathrm{\Phi }__l^{(l)\mathrm{}}(\mathrm{\Phi }_{GRT})`$ $`=Gr\mathrm{\Phi }_{\mathrm{Ih}}^{(l)}`$ (a series associated with the $`l`$-adic Ihara associator $`\mathrm{\Phi }_\sigma ^l`$ in loc.cit). By (3.3) we get their Berthelot-Ogus counterpart
$$\mathrm{\Phi }_{\mathrm{crys}}^{(p)\mathrm{}}(\mathrm{\Phi }_{GRT})=\mathrm{\Phi }_{\mathrm{KZ}}^p.$$
2. In loc.cit ยง6.6.2, we consider the meta-abelian quotient $`B_{GRT}๐ช(\underset{ยฏ}{GRT}_1)[[A,B]]`$ of $`\mathrm{\Phi }_{GRT}`$ and showed $`Gr\mathrm{\Phi }_{\mathrm{Hod}}^{\mathrm{}}(B_{GRT})=\frac{\mathrm{\Gamma }(1A)\mathrm{\Gamma }(1B)}{\mathrm{\Gamma }(1AB)}mod\pi ^2`$ and $`Gr\mathrm{\Phi }__l^{(l)\mathrm{}}(B_{GRT})`$ $`=GrB^{(l)}(A,B)`$ where $`\mathrm{\Gamma }(1z)=\mathrm{exp}\{\gamma z+\underset{n=2}{\overset{\mathrm{}}{}}\zeta (n)\frac{z^n}{n}\}`$ ($`\gamma `$: Euler constant) is the classical gamma function and $`GrB^{(l)}(A,B)`$ is the associated graded quotient of the Iharaโs universal power series for Jacobi sums \[I\]. By (3.3) we get their Berthelot-Ogus counterpart
$$\mathrm{\Phi }_{\mathrm{crys}}^{(p)\mathrm{}}(B_{GRT})=\underset{k=1}{\overset{\mathrm{}}{}}\frac{\mathrm{\Gamma }_p(1+p^kA)\mathrm{\Gamma }_p(1+p^kB)}{\mathrm{\Gamma }_p(1+p^k(A+B))}$$
where $`\mathrm{\Gamma }_p(1z)=\mathrm{exp}\{\gamma _pz+\underset{n=2}{\overset{\mathrm{}}{}}L_p(n,\omega ^{1n})\frac{z^n}{n}\}`$ ($`\gamma _p`$: a $`p`$-adic analogue of Euler constant) is the Moritaโs $`p`$-adic gamma function (\[M\] Theorem 1).
Next we discuss a special element of Drinfelโdโs $`๐ค๐ฏ๐ฑ_1=\underset{n}{}๐ค๐ฏ๐ฑ_1^n`$ <sup>14</sup><sup>14</sup>14 Actually it is isomorphic to the graded completion of Iharaโs \[I\] stable derivation algebra $`๐_{}`$. which is the graded Lie algebra of $`\underset{ยฏ}{GRT}_1`$ \[Dr\] ยง5. The original Drinfelโd associator $`\mathrm{\Phi }_{\mathrm{KZ}}(A,B)`$ satisfies the defining equations of $`\underset{ยฏ}{M}`$ \[Dr\] which is similar to but different from those of $`\underset{ยฏ}{GRT}_1`$ \[Dr\] however we say the following.
###### Lemma 3.3.
The โ$``$โ part $`\mathrm{\Phi }_{\mathrm{KZ}}^{}(A,B)`$ of the Drinfelโd associator satisfies the same equations to those of Proposition 3.1.
###### Proof .
Drinfelโd showed $`(1,\mathrm{\Phi }_{\mathrm{KZ}}(\frac{A}{2\pi i},\frac{B}{2\pi i}))\underset{ยฏ}{M}_1()`$ in \[Dr\]. By the complex conjugate action we get $`(1,\mathrm{\Phi }_{\mathrm{KZ}}(\frac{A}{2\pi i},\frac{B}{2\pi i}))\underset{ยฏ}{M}_1()`$. By the torsor structure of $`\underset{ยฏ}{M}_1()`$, there is a unique element $`(1,\phi )\underset{ยฏ}{GRT}_1()`$ such that $`(1,\phi )(1,\mathrm{\Phi }_{\mathrm{KZ}}(\frac{A}{2\pi i},\frac{B}{2\pi i}))=(1,\mathrm{\Phi }_{\mathrm{KZ}}(\frac{A}{2\pi i},\frac{B}{2\pi i}))`$, equivalently
$$\mathrm{\Phi }_{\mathrm{KZ}}(\frac{A}{2\pi i},\frac{B}{2\pi i})=\phi (A,B)\mathrm{\Phi }_{\mathrm{KZ}}(\frac{A}{2\pi i},\phi (A,B)^1(\frac{B}{2\pi i})\phi (A,B)).$$
By comparing the equations in Lemma 2.25, we get $`\phi =\mathrm{\Phi }_{\mathrm{KZ}}^{}(\frac{A}{2\pi i},\frac{B}{2\pi i})`$. So $`(1,\phi )=(1,\mathrm{\Phi }_{\mathrm{KZ}}^{}(\frac{A}{2\pi i},\frac{B}{2\pi i}))\underset{ยฏ}{GRT}_1`$, from which we get the claim. โ
###### Remark 3.4.
By Lemma 3.3, $`(1,\mathrm{\Phi }_{\mathrm{KZ}}^{})\underset{ยฏ}{GRT}_1()`$. By the logarithmic morphism $`Log:\underset{ยฏ}{GRT}_1()๐ค๐ฏ๐ฑ_1()`$, Drinfelโd obtained in \[Dr\] Proposition 6.3 a canonical element $`๐ฃ(\mathrm{})`$ in $`๐ค๐ฏ๐ฑ_1()`$, the image of $`(1,\mathrm{\Phi }_{\mathrm{KZ}}^{})`$, with the following presentations:
(3.5)
$$๐ฃ(\mathrm{})=\underset{m3:\text{ odd}}{}๐ฃ(\mathrm{})_m\text{where}๐ฃ(\mathrm{})_m=2\zeta (m)(adA)^{m1}(B)+\mathrm{}๐ค๐ฏ๐ฑ_1^m().$$
On the other hand, $`(1,\mathrm{\Phi }_{\mathrm{KZ}}^p)\underset{ยฏ}{GRT}_1(_p)`$ by Proposition 3.2. By the logarithmic morphism $`Log:\underset{ยฏ}{GRT}_1(_p)๐ค๐ฏ๐ฑ_1(_p)`$, we also get a canonical element $`๐ฃ_p`$ in $`๐ค๐ฏ๐ฑ_1(_p)`$, the image of $`(1,\mathrm{\Phi }_{\mathrm{KZ}}^p)`$, with the following presentation:
$$๐ฃ(p)=\underset{m3}{}๐ฃ(p)_m\text{where}๐ฃ(p)_m=\zeta _p(m)(adA)^{m1}(B)+\mathrm{}๐ค๐ฏ๐ฑ_1^m(_p).$$
(N.B. $`๐ฃ(\mathrm{})_m=0`$ for $`m=1`$ or even.) We stress that the summation in $`๐ฃ(\mathrm{})`$ is taken for odd number $`m`$ greater than or equal to 3 while the summation in $`๐ฃ(p)`$ is taken for natural number $`m`$ greater than or equal to 3.
### 3.2. Racinetโs context
We will recall Racinetโs \[R\] pro-algebraic torsor $`\underset{ยฏ}{DMR}_1`$ which is defined by two shuffle relations and discuss a story for $`\underset{ยฏ}{DMR}_1`$ analogous to previous section.
MZV $`\zeta (k_1,\mathrm{},k_m)`$ ($`m,k_1,\mathrm{},k_m`$, $`k_m>1`$) satisfies two types of so called shuffle product formula, expressing a product of two MZVโs as a linear combination of other such values. The first type, known as series shuffle product formula and for which the easiest example is the relation
$$\zeta (k_1)\zeta (k_2)=\zeta (k_1,k_2)+\zeta (k_2,k_1)+\zeta (k_1+k_2),$$
is easily obtained from the expression (0.2). The second type of shuffle product formula, known as integral shuffle product formula, come from their iterated integral expressions. The easiest example is the formula
$$\zeta (k_1)\zeta (k_2)=\underset{i=0}{\overset{k_11}{}}\left(\genfrac{}{}{0pt}{}{k_21+i}{i}\right)\zeta (k_1i,k_2+i)+\underset{j=0}{\overset{k_21}{}}\left(\genfrac{}{}{0pt}{}{k_11+j}{j}\right)\zeta (k_2j,k_1+j).$$
The double shuffle relations are linear relations combining series shuffle product formula and integral shuffle product formula. The above two formulae give the simplest example
$`\zeta (k_1,k_2)`$ $`+\zeta (k_2,k_1)+\zeta (k_1+k_2)`$
$`={\displaystyle \underset{i=0}{\overset{k_11}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{k_21+i}{i}}\right)\zeta (k_1i,k_2+i)+{\displaystyle \underset{j=0}{\overset{k_21}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{k_11+j}{j}}\right)\zeta (k_2j,k_1+j).`$
We have two extended notions of MZVโs to non-admissible indices, the case when the last component $`k_m`$ is equal to $`1`$, which we call series regularization and integral regularization here. The usual and series regularized MZVโs satisfy series shuffle product formulae and the usual and integral regularized MZVโs satisfy integral shuffle product formulae. These two kinds of regularized MZVโs are related by so called regularized relations.
Racinet constructed a pro-algebraic variety <sup>15</sup><sup>15</sup>15 For our convenience slightly we change the definition. $`\underset{ยฏ}{DMR}`$ whose set $`\underset{ยฏ}{DMR}(k)`$ of $`k`$-valued points ($`k`$: a field of characteristic $`0`$) consists of a series $`\phi (A,B)kA,B`$ such that $`\phi (A,B)`$ satisfies (3.2.1.1) and (3.2.1.2) in \[R\], which are relations reflecting the above relations. Let $`\lambda k`$. He denoted $`\underset{ยฏ}{DMR}_\lambda (k)`$ to be the subset consisting of the series in $`\underset{ยฏ}{DMR}(k)`$ whose coefficient of $`AB`$ is equal to $`\frac{\lambda ^2}{24}`$. He showed in \[R\] Theorem 1 that $`\underset{ยฏ}{DMR}_1(k)`$ (resp. $`\underset{ยฏ}{DMR}_0(k)`$) has a structure of the pro-algebraic torsor (resp. pro-algebraic group) where $`\underset{ยฏ}{DMR}_0(k)`$ acts from the left (resp. whose group law is defined by the same way to the Grothendieck-Teichmรผller group $`\underset{ยฏ}{GRT}_1`$).
We embed $`\pi _1^{\mathrm{DR}}(X_{}:\stackrel{}{01})(k)`$ into $`kA,B`$ via $`i`$ (2.1) and $`\pi _1^{\mathrm{Be}}(X():\stackrel{}{01})(k)`$ into $`kA,B`$ via $`j`$ (2.18). The group automorphism of $`\pi _1^{\mathrm{DR}}(X_{}:\stackrel{}{01})(k)`$ sending $`AA`$ and $`B\phi ^1B\phi `$ induces an embedding of both $`\underset{ยฏ}{DMR}_0`$ and $`\underset{ยฏ}{GRT}_1`$ into $`\underset{ยฏ}{Aut}\pi _1^{\mathrm{DR}}(X_{}:\stackrel{}{01})`$ as pro-algebraic groups. And the group isomorphism from $`\pi _1^{\mathrm{DR}}(X_{}:\stackrel{}{01})(k)`$ to $`\pi _1^{\mathrm{Be}}(X():\stackrel{}{01})(k)`$ sending $`AA`$ and $`B\phi ^1B\phi `$ induces an embedding of both $`\underset{ยฏ}{DMR}_1`$ and $`\underset{ยฏ}{M}_1`$ into $`\underset{ยฏ}{Isom}(\pi _1^{\mathrm{DR}}(X_{}:\stackrel{}{01}),\pi _1^{\mathrm{Be}}(X():\stackrel{}{01}))`$ as pro-algebraic torsors. From now on we regard $`\underset{ยฏ}{DMR}_0`$ and $`\underset{ยฏ}{GRT}_1`$ to be sub-pro-algebraic groups of $`\underset{ยฏ}{Aut}\pi _1^{\mathrm{DR}}(X_{}:\stackrel{}{01})`$ and regard $`\underset{ยฏ}{DMR}_1`$ and $`\underset{ยฏ}{M}_1`$ to be sub-pro-algebraic torsors of $`\underset{ยฏ}{Isom}(\pi _1^{\mathrm{DR}}(X_{}:\stackrel{}{01}),\pi _1^{\mathrm{Be}}(X():\stackrel{}{01}))`$.
We note that $`\mathrm{\Phi }_{\mathrm{KZ}}(\frac{A}{2\pi i},\frac{B}{2\pi i})`$ belongs to $`\underset{ยฏ}{DMR}_1()`$, which gives a morphism
(3.6)
$$\mathrm{\Psi }_{\mathrm{Hod}}:Spec\underset{ยฏ}{DMR}_1$$
analogous to (3.2) and conjectured to be isomorphic \[R\]. The paper \[BF\] is a trial to give a morphism
(3.7)
$$\mathrm{\Psi }_{\mathrm{crys}}^{(p)}:Spec^{(p)}\underset{ยฏ}{DMR}_0,$$
which is an analogous map to (3.3). We showed double shuffle relations, i.e. series and integral shuffle product formulae, for usual $`p`$-adic MZVโs but not for extended values. In \[FJ\] we will give a complete proof to give an embedding $`\mathrm{\Psi }_{\mathrm{crys}}^{(p)}`$, namely show series and integral shuffle product formulae and regularized relations for regularized $`p`$-adic MZVโs.
### 3.3. Deligne-Goncharovโs context
We will discuss the pro-algebraic bi-torsor of the motivic Galois group in \[DG\] by recalling the category of mixed Tate motives over $``$. This torsor is related with the Drinfelโdโs torsor (ยง3.1) in Proposition 3.8 and the Racinetโs torsor (ยง3.2) in Proposition 3.9. We also explain how the motivic formalism fits into our story in Note 3.10: In Hodge side we give a motivic interpretation of Zagierโs dimension conjecture on MZVโs. We recall how the upper-bounding part of this conjecture follows as a consequence. In Artin side we see how Iharaโs conjecture \[I2\] on Galois image is motivically related with Zagierโs conjecture via Problem 3.7. In Berthelot-Ogus side we deduce another motivic way, different from \[FJ\], of proving double (series and integral) shuffle relations and regularized relations for (extended) $`p`$-adic multiple zeta values by using Yamashitaโs \[Y\] fiber functor of the crystalline realization.
Let $`k`$ be a field with characteristic $`0`$. Levine \[L2\] and Voevodsky \[Voe\] constructed a triangulated category of mixed motives over $`k`$. Levine \[L2\] showed an equivalence of these two categories. This category denoted by $`DM(k)_{}`$ has Tate objects $`(n)`$ ($`n`$). Let $`DMT(k)_{}`$ be the triangulated sub-category of $`DM(k)_{}`$ generated by $`(n)`$ ($`n`$). Levine \[L1\] extracted a neutral tannakian $``$-category $`MT(k)_{}`$ of mixed Tate motives over $`k`$ from $`DMT(k)_{}`$ by taking a heart with respect to a $`t`$-structure under the Beilinson-Soulรฉ vanishing conjecture which says $`gr_i^\gamma K_n(k)=0`$ for $`n>2i`$. Here LHS is the graded quotient of the algebraic $`K`$-theory for $`k`$ with respect to $`\gamma `$-filtration.
From now on we assume that $`k`$ is a number field. In this case the Beilinson-Soulรฉ vanishing conjecture holds and we have $`MT(k)_{}`$. This category satisfies the following expected properties: Each object $`M`$ has an increasing filtration of subobjects called weight filtration, $`W:\mathrm{}W_{m1}MW_mMW_{m+1}M\mathrm{},`$ whose intersection is $`0`$ and union is $`M`$. The quotient $`Gr_{2m+1}^WM:=W_{2m+1}M/W_{2m}M`$ is trivial and $`Gr_{2m}^WM:=W_{2m}M/W_{2m+1}M`$ is a direct sum of finite copies of $`(m)`$ for each $`m`$. Morphisms of $`MT(k)_{}`$ are strictly compatible with weight filtration. The extension group is related to $`K`$-theory as follows
$$Ext_{MT(k)_{}}^i((0),(m))=\{\begin{array}{cc}K_{2mi}(k)_{}\hfill & \text{for }i=1,\hfill \\ 0\hfill & \text{for }i>1.\hfill \end{array}$$
There are realization fiber functors (\[L2\] and \[H\]) corresponding to usual cohomology theories.
Let $`S`$ be a finite set of finite places of $`k`$. Let $`๐ช_S`$ be the ring of $`S`$-integers in $`k`$. Deligne and Goncharov \[DG\] defined the full subcategory $`MT(๐ช_S)`$ of mixed Tate motives over $`๐ช_S`$, whose objects are mixed Tate motives $`M`$ in $`MT(k)_{}`$ such that for each subquotient $`E`$ of $`M`$ which is an extension of $`(n)`$ by $`(n+1)`$ for $`n`$, the extension class of $`E`$ in $`Ext_{MT(k)_{}}^1((n),(n+1))=Ext_{MT(k)_{}}^1((0),(1))=k_{}^\times `$ lies in $`๐ช_S^\times `$. In this category the following hold:
$`Ext_{MT(๐ช_S)}^1((0),(m))`$ $`=\{\begin{array}{cc}0\hfill & \text{for }m<1,\hfill \\ ๐ช_S^\times \hfill & \text{for }m=1,\hfill \\ K_{2m1}(k)_{}\hfill & \text{for }m>1,\hfill \end{array}`$
$`Ext_{MT(๐ช_S)}^2((0),(m))`$ $`=0.`$
Let $`\omega _{\mathrm{can}}:MT(๐ช_S)Vect_{}`$ ($`Vect_{}`$: the category of $``$-vector spaces) be the fiber functor which sends each motive $`M`$ to $`_nHom((n),Gr_{2n}^WM)`$. Let $`G_{\mathrm{can}}`$ be the motivic Galois group $`\underset{ยฏ}{Aut}^{}(MT():\omega _{\mathrm{can}})`$. The action of $`G_{\mathrm{can}}`$ on $`\omega _{\mathrm{can}}((1))=`$ defines a surjection $`G_{\mathrm{can}}๐พ_m`$ and its kernel $`U_{\mathrm{can}}`$ is the unipotent radical of $`G_{\mathrm{can}}`$. There is a canonical splitting $`\tau :๐พ_mG_{\mathrm{can}}`$ which gives a negative grading on the Lie algebra $`LieU_{\mathrm{can}}`$ (consult \[De1\] ยง8 for the full story). The above computations of $`Ext`$-groups follows
###### Proposition 3.5.
The graded Lie algebra $`LieU_{\mathrm{can}}`$ is free and its degree $`n`$-part of $`LieU_{\mathrm{can}}^{\mathrm{ab}}=U_{\mathrm{can}}^{\mathrm{ab}}`$ is isomorphic to the dual of $`Ext_{MT(๐ช_S)}^1((0),(n))`$.
###### Proof .
See \[De1\] ยง8 and \[DG\] ยง2. โ
Let us restrict in the case of $`k=`$, $`S=\mathrm{}`$, $`๐ช_S=`$. Let $`\omega _{}:MM_{}`$ ($`=`$ Be, DR) be the fiber functor which associates each mixed Tate motive $`MMT()`$ with the underlying vector space of its Betti, De Rham realization respectively. For $`,^{}=\{\text{Be, DR}\}`$ we denote the corresponding tannakian fundamental group $`\underset{ยฏ}{Aut}^{}(๐ฏ():\omega _{})`$ by $`G_{}`$ and the corresponding tannakian fundamental torsor $`\underset{ยฏ}{Isom}^{}(๐ฏ():\omega _{},\omega _{^{}})`$ by $`G_{^{}}`$. Note that the latter is a $`(G_{},G_{^{}})`$-bi-torsor and $`G_{}=G_{}`$. Let $`U_{}`$ be the sub-pro-algebraic group of $`G_{}`$ whose action on $`\omega _{}((1))=`$ is trivial and $`U_{^{}}`$ be the sub-pro-algebraic torsor of $`G_{^{}}`$ which induces a trivial map from $`\omega _{}((1))=`$ to $`\omega _{^{}}((1))=`$. Then $`U_{^{}}`$ is a $`(U_{},U_{^{}})`$-bi-torsor and $`U_{}=U_{}`$. Since $`M_{\mathrm{DR}}=\omega _{\mathrm{can}}(M)`$ we have $`G_{\mathrm{can}}=G_{\mathrm{DR}}`$. By Proposition 3.5 the Lie algebra $`LieU_{\mathrm{DR}}`$ of the unipotent part $`U_{\mathrm{DR}}`$ of $`G_{\mathrm{DR}}`$ should be a graded free Lie algebra generated by one element in each degree $`m`$ ($`m3`$: odd).
###### Remark 3.6.
In \[De1\] ยง8.12, Deligne constructed a free basis of the Lie algebra $`LieU_{\mathrm{DR}}()`$ as follows: The infinity Frobenius action $`\varphi _{\mathrm{}}`$ on $`\omega _{\mathrm{Be}}(M)`$ for $`M๐ฏ()`$ determines a point in $`G_{\mathrm{Be}}()`$ and by the Hodge comparison isomorphism (the period map) $`\omega _{\mathrm{Be}}(M)\omega _{\mathrm{DR}}(M)`$ it gives a point in $`G_{\mathrm{DR}}()`$ which we denote by the same symbol $`\varphi _{\mathrm{}}`$. Put $`\psi _{\mathrm{}}=\varphi _{\mathrm{}}\tau (1)`$. Then $`\psi _{\mathrm{}}U_{\mathrm{DR}}()`$ because $`\varphi _{\mathrm{}}G_{\mathrm{DR}}()`$ goes to $`1๐พ_m`$. The element $`๐ข_m(\mathrm{})`$ ($`m3`$: odd) generates freely the Lie algebra $`LieU_{\mathrm{DR}}()`$ where $`๐ข_m(\mathrm{})`$ means a degree $`m`$-part of the image $`๐ข(\mathrm{})=_m๐ข_m(\mathrm{})`$ of $`\psi _{\mathrm{}}`$ by the logarithmic morphism $`Log:U_{\mathrm{DR}}()LieU_{\mathrm{DR}}()`$. (N.B. $`๐ข_m(\mathrm{})=0`$ for $`m=1`$ or even.)
In \[DG\] ยง4 they constructed the motivic fundamental group $`\pi _1^{}(X:\stackrel{}{01})`$ with $`X=^1\backslash \{0,1,\mathrm{}\}`$, which is an ind-object of $`MT()`$. This is an affine group $`MT()`$-scheme (cf. ยง1.1), whose de Rham realization and Betti realization agree with the de Rham fundamental group $`\pi _1^{\mathrm{DR}}(X_{}:\stackrel{}{01})`$ in ยง1.1 and the Betti fundamental group $`\pi _1^{\mathrm{Be}}(X():\stackrel{}{01})`$ in ยง1.3 respectively. Namely $`\omega _{\mathrm{DR}}(\pi _1^{}(X:\stackrel{}{01}))=\pi _1^{\mathrm{DR}}(X_{}:\stackrel{}{01})`$ and $`\omega _{\mathrm{Be}}(\pi _1^{}(X:\stackrel{}{01}))=\pi _1^{\mathrm{Be}}(X():\stackrel{}{01})`$. Since all the structure morphism of $`\pi _1^{}(X:\stackrel{}{01})`$ belong to the set of morphisms of $`MT()`$ we have
(3.8)
$$\phi :G_{^{}}\underset{ยฏ}{Isom}(\omega _{}(\pi _1^{}(X:\stackrel{}{01})),\omega _{^{}}(\pi _1^{}(X:\stackrel{}{01})))$$
for $`,^{}\{\text{Be,DR}\}`$. On this map $`\phi `$ the following is one of the basic problems.
###### Problem 3.7.
Is $`\phi `$ injective?
This might be said a question which asks a validity of a unipotent variant of the so-called โBelyฤญโs theoremโ in \[Bel\] in the pro-finite setting. Equivalently this asks if the motivic fundamental group $`\pi _1^{}(X:\stackrel{}{01})`$ is a generator of the tannakian category $`MT()`$.
The $`(U_{\mathrm{DR}},U_{\mathrm{Be}})`$-bi-torsor $`U_{\mathrm{DR},\mathrm{Be}}`$ is related to the $`(\underset{ยฏ}{GRT}_1,\underset{ยฏ}{GT}_1`$)-bi-torsor $`\underset{ยฏ}{M}_1`$ in ยง3.1 as follows.
###### Proposition 3.8.
$`\phi (U_{\mathrm{DR}})\underset{ยฏ}{GRT}_1`$, $`\phi (U_{\mathrm{DR},\mathrm{Be}})\underset{ยฏ}{M}_1`$, $`\phi (U_{\mathrm{Be}})\underset{ยฏ}{GT}_1`$.
###### Proof .
By Lemma 2.24 the automorphism which corresponds to $`\varphi _{\mathrm{}}\tau (1)`$ is described by $`(1,\mathrm{\Phi }_{\mathrm{KZ}}^{})`$. So the induced homomorphism $`\mathrm{\Phi }:LieU_{\mathrm{DR}}DerLie\pi _1^{\mathrm{DR}}(X_{}:\stackrel{}{01})`$ sends $`๐ข(\mathrm{})`$ to $`๐ฃ(\mathrm{})`$ (cf. Remark 3.4). Since $`๐ฃ(\mathrm{})_m`$ is the image of the free generator $`๐ข(\mathrm{})_m`$, the image $`\mathrm{\Phi }(LieU_{\mathrm{DR}})`$ must lies in $`๐ค๐ฏ๐ฑ_1(=Lie\underset{ยฏ}{GRT}_1)DerLie\pi _1^{\mathrm{DR}}(X_{}:\stackrel{}{01})`$. So $`\phi (U_{\mathrm{DR}})`$ should lie in $`\underset{ยฏ}{GRT}_1\underset{ยฏ}{Aut}\pi _1^{\mathrm{DR}}(X_{}:\stackrel{}{01})`$.
The morphism $`\tau (2\pi i)^1p`$ lies in $`U_{\mathrm{DR},\mathrm{Be}}()`$. Its image $`\phi (\tau (2\pi i)^1p)`$, which is a morphism from $`\pi _1^{\mathrm{DR}}(X_{}:\stackrel{}{01})()`$ to $`\pi _1^{\mathrm{Be}}(X():\stackrel{}{01})()`$, is described by $`(1,\mathrm{\Phi }_{\mathrm{KZ}}(\frac{A}{2\pi i},\frac{B}{2\pi i}))`$ by Lemma 2.23. By \[Dr\] this element belongs to $`\underset{ยฏ}{M}_1()`$, which implies $`\phi (U_{\mathrm{DR},\mathrm{Be}})\underset{ยฏ}{M}_1`$. It is because $`\phi `$ is a morphism from the $`U_{\mathrm{DR}}`$-torsor $`U_{\mathrm{DR},\mathrm{Be}}`$ to the $`\underset{ยฏ}{GRT}_1`$-torsor $`\underset{ยฏ}{M}_1`$, which sends $`\tau (2\pi i)^1pU_{\mathrm{DR},\mathrm{Be}}()`$ to $`(1,\mathrm{\Phi }_{\mathrm{KZ}}(\frac{A}{2\pi i},\frac{B}{2\pi i}))\underset{ยฏ}{M}_1()`$. โ
The $`(U_{\mathrm{DR}},U_{\mathrm{Be}})`$-bi-torsor $`U_{\mathrm{DR},\mathrm{Be}}`$ is also related to the $`\underset{ยฏ}{DMR}_0`$-torsor $`\underset{ยฏ}{DMR}_1`$ in ยง3.2 as follows.
###### Proposition 3.9.
$`\phi (U_{\mathrm{DR}})\underset{ยฏ}{DMR}_0`$, $`\phi (U_{\mathrm{DR},\mathrm{Be}})\underset{ยฏ}{DMR}_1`$.
###### Proof .
The proof is similar to the previous proposition. As for the proof for $`\phi (U_{\mathrm{DR}})\underset{ยฏ}{DMR}_0`$, we use Racinetโs result \[R\] 5.3.2, which is the same argument to \[Dr\] Proposition 6.3, that the Lie algebra $`๐ก๐ช๐ฏ_0`$ of his $`\underset{ยฏ}{DMR}_0`$ contains the Drinfelโdโs element $`๐ฃ(\mathrm{})_m`$ (3.5). โ
At present we do not know the relationship between Drinfelโdโs $`\underset{ยฏ}{GRT}_1`$ and Racinetโs $`\underset{ยฏ}{DMR}_0`$ although it might be expected that they are equal. However the recent Terasomaโs work (to appear) might suggests a direction $`\underset{ยฏ}{DMR}_0\underset{ยฏ}{GRT}_1`$.
Below we explain how the motivic formalism fits in the scheme of our subject in the previous section.
###### Note 3.10.
1. In Hodge side: as is shown in the proof of Proposition 3.8 the pair $`(1,\mathrm{\Phi }_{\mathrm{KZ}}(\frac{A}{2\pi i},\frac{B}{2\pi i}))`$ is the image of $`\tau (2\pi i)^1p`$ by $`\phi `$, that is, $`(1,\mathrm{\Phi }_{\mathrm{KZ}}(\frac{A}{2\pi i},\frac{B}{2\pi i}))\phi (U_{\mathrm{DR},\mathrm{Be}})()`$. Since $`(2\pi i)^{(k_1+\mathrm{}+k_m)}\zeta (k_1,\mathrm{},k_m)`$ ($`m,k_i`$,$`k_m>1`$) appears among each coefficient of $`\mathrm{\Phi }_{\mathrm{KZ}}(\frac{A}{2\pi i},\frac{B}{2\pi i})`$, it gives rise a morphism
(3.9)
$$\mathrm{\Gamma }_{\mathrm{Hod}}:Spec๐ต\phi (U_{\mathrm{DR},\mathrm{Be}})$$
analogous to (3.2) and (3.6). There is a conjecture on the dimension of the vector space of MZVโs at each weight which is called Zagier conjecture (see \[Za\] and also \[F0\]) and partly proved (for upper-bounding part) by Terasoma \[T\]. By using the embedding (3.9) and Proposition 3.5, Deligne and Goncharov also get a partial (upper-bounding part) proof of Zagier conjecture in \[DG\]. We also get to know that to say Zagier conjecture holds is equivalent to say the surjectivity of $`\mathrm{\Gamma }_{\mathrm{Hod}}`$ and the injectivity of $`\phi `$ (i.e. the validity of Problem 3.7).
2. In Artin side: we consider the absolute Galois group $`Gal(\overline{}/)`$ action on the $`l`$-adic ($`l`$: a prime) รฉtale fundamental group $`\phi _1:Gal(\overline{}/)\underset{ยฏ}{Aut}\pi _1^{l,\text{รฉt}}(X_\overline{}:\stackrel{}{01})`$ and the $`l`$-adic Galois image pro-algebraic group $`\underset{ยฏ}{Gal}__l^{(l)}`$ defined to be the Zariski closure of $`\phi _1(Gal(\overline{}/(\mu _l^{\mathrm{}}))`$ in \[F0\] ยง4. The category $`MT()`$ has a fiber functor $`\omega _{l,\text{รฉt}}:MM_{l,\text{รฉt}}`$ of the $`l`$-adic รฉtale realization which associates each motives $`M`$ with its underlying $`_l`$-vector space of its $`l`$-adic รฉtale realization. The absolute Galois group $`Gal(\overline{}/)`$ acts functorially on $`\omega _{l,\text{รฉt}}(M)`$ and hence it induces a morphism $`Gal(\overline{}/)G_{l,\text{รฉt}}(_l)`$. Hereafter we fix an embedding $`\overline{}`$. Then there is the Artin comparison isomorphism $`\omega _{l,\text{รฉt}}(M)\omega _{\mathrm{Be}}(M)_l`$ functorial with respect to $`M`$, which gives $`G_{l,\text{รฉt}}=G_{\mathrm{Be}}\times _l`$. We might have a morphism $`\phi _2:Gal(\overline{}/)G_{\mathrm{Be}}(_l)`$. Because of $`\pi _1^{l,\text{รฉt}}(X_\overline{}:\stackrel{}{01})=\omega _{l,\text{รฉt}}(\pi _1^{}(X:\stackrel{}{01}))=\omega _{\mathrm{Be}}(\pi _1^{}(X:\stackrel{}{01}))\times _l`$, the map $`\phi _1`$ might factor through $`G_{\mathrm{Be}}`$-action on $`\omega _{Be}(\pi _1^{}(X:\stackrel{}{01}))`$, that means $`\phi _1=\phi \phi _2`$ (consult \[De1\] ยง8 for full story). In loc.cit. ยง8.14 it is also shown that the image of $`\phi _2`$ is open in the topology of $`G_{\mathrm{Be}}(_l)`$. So by restricting into unipotent parts we get an isomorphism
(3.10)
$$\mathrm{\Gamma }_{l,\text{รฉt}}:\underset{ยฏ}{Gal}__l^{(l)}\stackrel{=}{}\phi (U_{\mathrm{Be}})\times _l$$
which might be said an analogue of (3.1). This is the way how the common $``$-structure of the $`l`$-adic Galois image pro-algebraic group for all prime $`l`$ is detected in \[DG\] Remark 6.13. The graded Lie algebra $`๐ค_{}^{(l)}`$ (cf.\[I2\] and also \[F0\]) associated with $`\underset{ยฏ}{Gal}__l^{(l)}`$ is conjectured by Ihara (loc.cit.) to be a free Lie algebra generated by one element in each degree 3,5,7,9,.. As in Hodge case we can say that this Iharaโs conjecture is equivalent to the injectivity of $`\phi `$ (i.e. the validity of Problem 3.7).
3. In Berthelot-Ogus side: Yamashita will construct in \[Y\] the fiber functor $`\omega _{p,\mathrm{crys}}:MM_{p,\mathrm{crys}}`$ of the crystalline realization of $`MT()`$ which is compatible with $`\omega _{\mathrm{DR}}`$. It associates each mixed Tate motive $`M`$ with its underlying $`_l`$-vector space of its crystalline realization by using Fontaineโs functor to give a conjectured dimension bounding for $`Z_w^{(p)}`$. It admits a functorial crystalline Frobenius action $`\varphi _{p,\mathrm{crys}}`$ on $`\omega _{p,\mathrm{crys}}(M)`$ and a functorial comparison isomorphism $`\omega _{p,\mathrm{crys}}(M)=\omega _{\mathrm{DR}}(M)_p`$ for $`MT()`$, which gives a point in $`G_{\mathrm{DR}}(_p)`$ denoted by the same symbol $`\varphi _{p,\mathrm{crys}}`$. Since $`\varphi _{p,\mathrm{crys}}`$ goes to $`\frac{1}{p}๐พ_m`$, the element $`\psi _p=\varphi _{p,\mathrm{crys}}\tau (p)`$ must lie in $`U_{\mathrm{DR}}(_p)`$. As is similar to Remark 3.4 and Remark 3.6 we also get in the $`p`$-adic setting elements $`e(p)_m`$ ($`m3`$) which is on a degree $`m`$-part of $`LieU_{\mathrm{DR}}(_p)`$. The action of $`\psi _p`$ on $`\pi _1^{\mathrm{DR}}(X_{}:\stackrel{}{01})`$ is described as $`AA`$ and $`B\mathrm{\Phi }_{\mathrm{De}}^p(A,B)^1B\mathrm{\Phi }_{\mathrm{De}}^p(A,B)`$ by Lemma 2.9. Therefore $`(1,\mathrm{\Phi }_{\mathrm{De}}^p)\phi (U_{\mathrm{DR}}(_p))`$. By Proposition 3.9, $`\mathrm{\Phi }_{\mathrm{De}}^p\underset{ยฏ}{DMR}_0(_p)`$. Since Deligneโs $`p`$-adic MZVโs $`\zeta _p^{\mathrm{De}}(k_1,\mathrm{},k_m)`$ are coefficients of $`\mathrm{\Phi }_{\mathrm{De}}^p(A,B)`$, they must satisfy double shuffle relations. By the same argument to Proposition 3.1, we get $`(1,\mathrm{\Phi }_{\mathrm{KZ}}^p)\phi (U_{\mathrm{DR}}(_p))`$. It follows that we also get an embedding
(3.11)
$$\mathrm{\Gamma }_{\mathrm{crys}}^{(p)}:Spec^{(p)}\phi (U_{\mathrm{DR}}).$$
Combining (3.11) with Proposition 3.9, we get the embedding (3.7). Therefore our $`p`$-adic MZVโs $`\zeta _p(k_1,\mathrm{},k_m)`$ must satisfy (series and integral) double shuffle relations and regularization relations, which are defining relations of $`\underset{ยฏ}{DMR}_0`$. This is a motivic proof of double shuffle relations for $`p`$-adic MZVโs, which was a project posed in \[De2\].
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# Local existence of dynamical and trapping horizons
thanks: Supported in part by the NSF, contract no. DMS 0104402.thanks: Supported in part by the Spanish Ministerio de Educaciรณn y Ciencia, project BFME2003-02121.thanks: Supported in part by the Austrian FWF, project no. P14621-N05.
## Abstract
Given a spacelike foliation of a spacetime and a marginally outer trapped surface S on some initial leaf, we prove that under a suitable stability condition S is contained in a โhorizonโ, i.e. a smooth 3-surface foliated by marginally outer trapped slices which lie in the leaves of the given foliation. We also show that under rather weak energy conditions this horizon must be either achronal or spacelike everywhere. Furthermore, we discuss the relation between โboundingโ and โstabilityโ properties of marginally outer trapped surfaces.
The application of numerical relativity to black hole spacetimes is, together with the role played by black hole thermodynamics as a testing ground for quantum gravity, among the factors that have caused a shift of interest from global properties of black holes such as the event horizon, knowledge of which requires information about the infinite future, towards quasilocal properties. By quasilocal properties one means such properties that can at least in principle be measured by an observer with a finite life span and hence also can be studied during the course of a numerical evolution of a black hole spacetime.
A closed spacelike surface $`S`$ in a spacetime $`(M,g_{\alpha \beta })`$ is called trapped if future directed null rays emanating from $`S`$ are converging. If $`M`$ contains a trapped surface, satisfies the null energy condition and some causality condition, then $`M`$ is singular RP . Suppose $`M`$ is foliated by a family of spacelike Cauchy surfaces $`\{\mathrm{\Sigma }_t\}`$. The apparent horizon, defined as the family of boundaries of the regions containing trapped surfaces in the $`\{\mathrm{\Sigma }_t\}`$, is a quasilocally defined object that plays an important role in black hole thermodynamics, as well as in numerical evolutions of black holes. It should be noted however that the apparent horizon depends on the choice of the reference foliation $`\{\mathrm{\Sigma }_t\}`$. If sufficiently smooth, the apparent horizon is foliated by marginally outer trapped surfaces (MOTS) KH . The latter are defined to have vanishing outgoing null expansion, (while the ingoing one is not restricted).
In numerical evolution of black hole spacetimes, it is now standard to avoid the singular behavior of both gravitational field and gauge conditions in the interior of black holes, by excising a suitable region inside the boundary of the black holes, as defined by a collection of MOTS, from the computational domain. However, tracking a family of MOTS during an evolution one encounters the occasional โsuddenโ appearance of new MOTS and by โjumpsโ of the MOTS (see, e.g. BH ). It thus becomes important to study such an evolution analytically as far as possible. In this Letter we prove, in Theorem 1, existence of a horizon, i.e. a hypersurface foliated by MOTS, provided the initial surface $`S`$ satisfies a natural stability condition, and in Theorem 2 we give causal properties of $`H`$. The condition of โstrictly stably outermostโ, which is crucial for these results, means that there is an outward deformation of $`S`$ such that the corresponding variation of the outgoing null expansion is nonnegative and positive somewhere (c.f. Definition 2 for details).
###### Theorem 1.
Let $`(M,g_{\alpha \beta })`$ be a smooth spacetime foliated by smooth spacelike hypersurfaces $`\mathrm{\Sigma }_t`$. Assume that some leaf $`\mathrm{\Sigma }=\mathrm{\Sigma }_0`$ contains a smooth marginally outer trapped surface $`S`$ which is strictly stably outermost.
Then, $`S`$ is contained in a smooth horizon $`H`$ whose marginally outer trapped leaves lie in $`\mathrm{\Sigma }_t`$, and which exists at least as long as these marginally trapped leaves remain strictly stably outermost.
In Theorem 2 below we use the same notation as in Theorem 1, and we denote by $`l^\alpha `$ the null vector for which the expansion $`\theta _l`$ vanishes on $`S`$. Recall that the null energy condition holds if $`G_{\alpha \beta }j^\alpha j^\beta 0`$ for any null vector $`j^\alpha `$, where $`G_{\alpha \beta }`$ is the Einstein tensor.
###### Theorem 2.
If, in addition to the hypotheses of Theorem 1, the null energy condition holds, the horizon $`H`$ is locally achronal. If, moreover, $`G_{\alpha \beta }l^\alpha l^\beta >0`$ somewhere on $`S`$ or if $`S`$ has nonvanishing shear with respect to $`l^\alpha `$ somewhere, then $`H`$ is spacelike everywhere near $`\mathrm{\Sigma }`$.
The term horizon in this Letter is closely related to dynamical horizons introduced by Ashtekar and Krishnan AK and to trapping horizons introduced by Hayward SH . The latter two are more restrictive in the sense that they also require $`\theta _k<0`$ for the expansion along the other future null direction $`k`$. While dynamical horizons are spacelike by definition, trapping horizons may have any causal character a priori, but they are required to satisfy an additional stability condition, namely that the variation of $`\theta _l`$ along $`k`$ is negative. Our condition of โstrictly stably outermostโ can be generalized to variations in the outward past null cone $`C_{}`$ of $`S`$. Haywardโs stability condition is then closely related to our stability condition along null directions. Most of the results stated in this Letter, in particular the existence theorem, extend to the null case. Details will be given elsewhere AMS .
Since the location in spacetime of apparent, dynamical and trapping horizons depends in general on the foliation $`\mathrm{\Sigma }_t`$, it is clear that the same applies to the horizons obtained above. However, for generic dynamical horizons, this dependence is limited by the uniqueness results of Ashtekar and Galloway AG . We also note that the result in Theorem 2 on the causal character of $`H`$ is stronger than the one known SH for trapping horizons.
The example illustrated in Fig. 1 will shed light on the possible behavior of our horizons. It shows a horizon (the thick line) in a spacetime foliated by spacelike hypersurfaces $`\mathrm{\Sigma }_t`$ (thin lines). The horizon separates a region in which the $`\mathrm{\Sigma }_t`$ contain outer trapped surfaces (which we call โtrapped regionโ, shaded in Fig. 1) from a region where the $`\mathrm{\Sigma }_t`$ are free of them; the intersection โpointsโ with the foliation are MOTS. Note that Fig. 1 incorporates naturally the observed sudden appearance of MOTS during the evolution (e.g. point e). If the numerical analysis only looks for globally outermost MOTS (as it is usually done) it is clear that they jump (e.g., from b to e), while they are in fact connected by a horizon interpolating between both which โruns downwardโ in some places (in the interval (d,e)). Thus Fig. 1 makes compatible a smooth horizon with the jumps observed numerically. Examples like Fig. 1 can, in particular, be constructed in spherically symmetric spacetimes by choosing the spacelike foliation $`\mathrm{\Sigma }_t`$ suitably. Thus one may expect that the situation described by Fig. 1 is typical, though the causal character of the โdownwardโ part of the horizon is yet unclear; see the discussion in AG .
We now introduce some notation needed for the precise statements of our results. All fields and manifolds will be assumed to be $`C^{\mathrm{}}`$ unless otherwise stated. Let $`(M,g_{\alpha \beta })`$ be a spacetime with signature $`+++`$. Given a spacelike surface $`S`$ in $`M`$ we may choose two future directed null fields $`l^\alpha ,k^\alpha `$. Recall that the variation $`\delta _p\nu `$ of the geometric object $`\nu `$ defined on a surface $`T`$ in the direction of the vector $`p^\alpha `$ is defined by $`\delta _p\nu =\nu /\tau `$ for any one-parameter family of surfaces $`T_\tau `$ with $`T_0=T`$ and $`p^\alpha _{x^\alpha }=/\tau |_{\tau =0}`$. The null expansion $`\theta _l`$ is defined by $`\mu \theta _l=\delta _l\mu `$, where $`\mu `$ is the volume form on $`S`$. It should be noted that the variation is additive in the sense that for example $`\delta _{\psi k+l}\theta _l=\delta _{\psi k}\theta _l+\delta _l\theta _l`$, for some function $`\psi `$, but in general $`\delta _{\psi k}\theta _l\psi \delta _k\theta _l`$.
A closed spacelike surface $`S`$ is called outer trapped (weakly outer trapped, marginally outer trapped) if one of the null expansions, say $`\theta _l`$, is negative (non-positive, zero) everywhere on $`S`$. (For an alternative terminology, c.f. MS ). The usual definition of โtrapped surfaceโ and โmarginally trapped surfaceโ requires additional conditions on the expansion with respect to the other null vector $`k^\alpha `$. Corresponding definitions for untrapped surfaces are made by reversing the signs. Let a reference foliation $`\{\mathrm{\Sigma }_t\}`$ of $`M`$ by spacelike hypersurfaces be given, and select one such surface $`\mathrm{\Sigma }=\mathrm{\Sigma }_0`$. For a MOTS $`S\mathrm{\Sigma }`$, we define the โoutwardโ direction within $`\mathrm{\Sigma }`$ as the one to which the projection to $`\mathrm{\Sigma }`$ of the null vector $`l^\alpha `$ selected above points. This definition of โoutwardโ need not coincide with the intuitive one in asymptotically flat spacetimes. However, all our results hold for arbitrary spacetimes (not necessarily asymptotically flat) and $`l^\alpha `$ defines a local concept of โoutwardโ for MOTS. The unit outward normal to $`S`$ tangent to $`\mathrm{\Sigma }_t`$ is called $`m^\alpha `$, the future pointing unit normal to $`\mathrm{\Sigma }_t`$ is $`n^\alpha `$, and we scale the null vectors $`l^\alpha `$ and $`k^\alpha `$ such that $`l^\alpha =n^\alpha +m^\alpha `$ and $`k^\alpha =n^\alpha m^\alpha `$.
The following definitions are, apart from later use, motivated by similar definitions of Newman RN and of Kriele and Hayward KH , and by results in these papers.
###### Definition 1.
A marginally outer trapped surface $`S`$ is called locally outermost in $`\mathrm{\Sigma }`$, iff there exists a two-sided neighbourhood of $`S`$ such that its exterior part does not contain any weakly outer trapped surface.
###### Definition 2.
A marginally outer trapped surface $`S`$ is called stably outermost iff there exists a function $`\psi 0`$, $`\psi 0`$, on $`S`$ such that $`\delta _{\psi m}\theta _l0`$. $`S`$ is called strictly stably outermost if, moreover, $`\delta _{\psi m}\theta _l0`$ somewhere on $`S`$.
In Fig. 1, the points in the interval \[a,d\] represent stably outermost surfaces, those in \[a,c) locally outermost ones, and those in (a,c) strictly stably outermost ones. This example, and the result in KH suggest the implications: Strictly stably outermost $``$ Locally outermost $``$ Stably outermost, and the picture also suggests counterexamples for the opposite directions. We now give the tools required to show these results and the Theorems.
For a function $`\psi `$ on $`S`$, we define a linear elliptic operator $`L_\mathrm{\Sigma }`$ by $`L_\mathrm{\Sigma }\psi =\delta _{\psi m}\theta _l`$. Explicitly, we obtain
$`L_\mathrm{\Sigma }\psi `$ $`=\mathrm{\Delta }_S\psi +2s^AD_A\psi +({\displaystyle \frac{1}{2}}R_Ss_As^A+`$
$`+D_As^A{\displaystyle \frac{1}{2}}K_{AB}^\mu K^{\nu AB}l_\mu l_\nu G_{\alpha \beta }l^\alpha n^\beta )\psi .`$ (1)
Here $`D_A`$ is the covariant derivative on $`S`$, $`\mathrm{\Delta }_S`$ is the corresponding Laplacian, $`R_S`$ is the scalar curvature, $`K_{AB}^\mu `$ is the second fundamental form vector (defined by $`K_{AB}^\mu v_\mu =_Av_B`$ for any normal $`v_\alpha `$ to $`S`$, where $`_\alpha `$ is the covariant derivative on $`(M,g_{\alpha \beta })`$) and $`s_A`$ is the torsion of $`l^\alpha `$ (the 1-form $`s_A=\frac{1}{2}k_\alpha _Al^\alpha `$ on $`S`$).
$`L_\mathrm{\Sigma }`$ is analogous to the stability operator for minimal surfaces. In general, $`L_\mathrm{\Sigma }`$ is not self-adjoint but the eigenvalues of $`L_\mathrm{\Sigma }`$ have their real part bounded from below. The eigenvalue with smallest real part is called the principal eigenvalue. The following holds for second order elliptic operators of the form of $`L_\mathrm{\Sigma }`$.
###### Lemma 1.
The principal eigenvalue $`\lambda `$ of $`L_\mathrm{\Sigma }`$ is real. Moreover, the corresponding principal eigenfunction $`\varphi `$ (which satisfies $`L_\mathrm{\Sigma }\varphi =\lambda \varphi `$) is either everywhere positive or everywhere negative.
This Lemma is a consequence of the Krein-Rutman theorem which can be applied to second order elliptic operators along the lines in corollary A3 of Smoller JS . The discussion in Smollerโs corollary can be adapted straightforwardly to the case without boundary.
We now restate Definition 2 in terms of $`\lambda `$ as follows.
###### Lemma 2.
Let $`S\mathrm{\Sigma }`$ be a MOTS and let $`\lambda `$ be the principal eigenvalue of the corresponding operator $`L_\mathrm{\Sigma }`$. Then $`S`$ is stably outermost iff $`\lambda 0`$ and strictly stably outermost iff $`\lambda >0`$.
###### Proof.
If $`\lambda 0`$, choose $`\psi `$ in the definition of (strictly) stably outermost as a positive eigenfunction $`\varphi `$ corresponding to $`\lambda `$. Then $`\delta _{\varphi m}\theta _l=L_\mathrm{\Sigma }\varphi =\lambda \varphi 0`$. For the converse, we note that the adjoint $`L_\mathrm{\Sigma }^{}`$ (with respect to the standard $`L^2`$ inner product $`,`$ on $`S`$) has the same principal eigenvalue as $`L_\mathrm{\Sigma }`$, and a positive principal eigenfunction $`\varphi ^{}`$. Thus, for $`\psi `$ as in the definition of (strictly) stably outermost,
$$\lambda \varphi ^{},\psi =L_\mathrm{\Sigma }^{}\varphi ^{},\psi =\varphi ^{},L_\mathrm{\Sigma }\psi 0,$$
with strict inequality in the strictly stable case. Since $`\varphi ^{},\psi >0`$, the Lemma follows. โ
###### Proposition 1.
* A strictly stably outermost surface $`S`$ is locally outermost. Moreover, $`S`$ has a two-sided neighbourhood $`U`$ such that no weakly outer trapped surfaces contained in $`U`$ enter the exterior of $`S`$ and no weakly outer untrapped surfaces contained in $`U`$ enter the interior of $`S`$.
* A locally outermost surface $`S`$ is stably outermost.
###### Proof.
The first statement of (i) is in fact contained in the second one. To show the latter, let $`\varphi `$ be the positive principal eigenfunction of $`L_\mathrm{\Sigma }`$. Since $`L_\mathrm{\Sigma }\varphi >0`$ by assumption, flowing $`S`$ in $`\mathrm{\Sigma }`$ along any extension of $`\varphi m^\alpha `$ produces a family $`S_\sigma `$, $`\sigma (ฯต,ฯต)`$ for some $`ฯต>0`$. By choosing $`ฯต`$ small enough, the $`S_\sigma `$ have $`\theta _l|_{S_\sigma }>0`$ for $`\sigma (ฯต,0)`$ and $`\theta _l|_{S_\sigma }<0`$ for $`\sigma (0,ฯต)`$. We can now take $`U`$ to be the neighbourhood of $`S`$ given by $`U=_{\sigma (ฯต,ฯต)}S_\sigma `$.
Now let $`B`$ be a weakly outer trapped surface contained in $`U`$ which enters the exterior part of $`U`$. Then the function $`\sigma `$ has a maximum $`\sigma _p>0`$ at some point $`p`$ in $`B`$. At $`p`$, $`B`$ is tangent to $`S_{\sigma _p}`$ and we have
$$\theta _l|_B0<\theta _l|_{S_{\tau _p}}.$$
(2)
For a surface represented as a graph with respect to a function $`f`$, the map $`f\theta _l`$ is a quasilinear elliptic operator and the strong maximum principle GT applies to show that the inequality (2) holds only if $`B`$ coincides with $`S_{\tau _p}`$ and hence $`\theta _l|_B=\theta _l|_{S_{\tau _p}}`$ which gives a contradiction. Hence $`B`$ cannot enter the exterior part of $`U`$.
To show (ii), assume $`S`$ is locally outermost but not stably outermost. From Lemma 2, the principal eigenvalue $`\lambda `$ is then negative. Arguing as above one constructs a foliation outside $`S`$ with leaves which are outer trapped near $`S`$, contradicting the assumption. โ
Theorem 5.1 of Ashtekar and Galloway AG implies that the domain exterior to $`S`$ to which outer trapped surfaces cannot enter is determined by the past domain of dependence of any dynamical horizon through $`S`$, provided that some genericity conditions hold. Outer trapped surfaces โfar outsideโ of a locally outermost MOTS might exist in general (as Fig. 1 suggests for the surface in the interval (b,c)). To exclude this, one could define โglobally outermostโ surfaces (in particular in an asymptotically flat context). We can now prove our main theorem.
Proof of Theorem 1. Consider a foliation $`\mathrm{\Sigma }_t`$ with a MOTS $`S`$ on $`\mathrm{\Sigma }=\mathrm{\Sigma }_0`$. Let $`C_+`$ be the null cone generated by null rays starting from $`S`$ in the direction of $`l^\alpha `$ and let $`\stackrel{~}{S}_t=C_+\mathrm{\Sigma }_t`$ for $`t`$ close to $`0`$. We now introduce coordinates $`(t,r,x^A)`$ in a neighborhood of $`S`$ such that at $`S`$, $`_r`$ is the normal $`m^\alpha `$ and $`_t`$ is parallel to $`l^\alpha `$.
On $`\mathrm{\Sigma }_t`$ we consider surfaces which are given as graphs $`r=f(x^A)`$ in this coordinate system and we define a functional $`\mathrm{\Theta }[t,f]`$ whose value is $`\theta _l`$ on the surface, and which acts on $`f`$ as a quasilinear elliptic operator of the form $`\mathrm{\Theta }[t,f]=a^{AB}(f,f)_A_Bf+b(f,f)`$ where the coefficients $`a^{AB}`$ and $`b`$ are smooth functions depending on $`x`$ and on $`t`$, $`f`$ and $`_Af`$. For integer $`k0`$, $`\alpha (0,1)`$, let $`C^{k,\alpha }`$ be Hรถlder spaces on $`S`$. Let $`I=(ฯต,ฯต)`$ for $`ฯต>0`$. One checks that for some $`ฯต>0`$, and for any $`k2`$, there are neighborhoods $`๐ฐ_1`$ and $`๐ฐ_2`$ of zero in $`C^{k,\alpha }`$ and $`C^{k2,\alpha }`$ respectively so that $`\mathrm{\Theta }[t,f]:I\times ๐ฐ_1๐ฐ_2`$ is a well-defined $`C^{\mathrm{}}`$ map. Let $`D_y\mathrm{\Theta }`$ be the derivative with respect to the second argument. Then we have from the definition of $`\mathrm{\Theta }`$, $`D_y\mathrm{\Theta }[0,0].\beta =L_\mathrm{\Sigma }\beta .`$ As $`S`$ is assumed to be strictly stably outermost, the principal eigenvalue of $`L_\mathrm{\Sigma }`$ is positive by Lemma 2. By the Fredholm alternative $`L_\mathrm{\Sigma }`$ is invertible as a map $`L_\mathrm{\Sigma }:C^{k,\alpha }C^{k2,\alpha }`$ for arbitrary $`k2`$, $`\alpha (0,1)`$. Now the implicit function theorem for Banach space maps SL applies to prove local existence of a smooth horizon $`H`$ which by construction has the property that the leaves $`S_t=H\mathrm{\Sigma }_t`$ near $`S`$ are MOTS. By patching charts together it is also clear that existence holds as long as the $`S_t`$ stay strictly stably outermost. $`\mathrm{}`$
A theorem of Schoen RS asserts the existence of a MOTS between barrier surfaces $`S_1,S_2`$ with $`S_1`$ trapped and $`S_2`$ untrapped if the dominant energy condition holds. Besides its clear interest, this result also suggests an alternative approach to existence of a horizon: Start from a weakly outer trapped surface $`S`$ on some initial slice and take the null cone emanating from it. By the Raychaudhuri equation, $`\delta _{\omega l}\theta _l=\omega W`$ for any function $`\omega `$ and $`W=K_{AB}^\mu K^{\nu AB}l_\mu l_\nu +G_{\mu \nu }l^\mu l^\nu `$. If $`W>0`$, the null cone cuts each subsequent slice on an outer trapped surface which gives the trapped barrier $`S_1`$, while an untrapped barrier $`S_2`$ always exists near infinity for asymptotically flat data. Schoenโs result then yields existence of a MOTS on every subsequent slice. The resulting โhorizonโ may in general jump (e.g. from b to e in Fig. 1), but it need not be outermost. On the other hand, our Theorem 1 requires a MOTS on the initial surface instead of just a trapped one, but we do not assume any asymptotic properties, and we obtain a smooth horizon.
Results on the causal character of the horizon can be obtained by combining again the Raychaudhuri equation on the null cone as sketched above with a maximum principle or a โbarrier argumentโ inside $`\mathrm{\Sigma }_t`$ for $`t`$ near $`\mathrm{\Sigma }_0`$. For the latter, we may use e.g. part (i) of Proposition 1. However, we can also do โboth steps at onceโ by using the following Lemma.
###### Lemma 3.
For a strictly stably outermost surface $`S`$, any normal variation $`\psi m^\alpha `$ of $`S`$ with $`\delta _{\psi m}\theta _l0`$ satisfies $`\psi 0`$, i.e. the variation cannot be directed to the interior anywhere on $`S`$. If, moreover, $`\delta _{\psi m}\theta _l0`$ somewhere then $`\psi >0`$, i.e. the variation is directed to the exterior everywhere on $`S`$.
###### Proof.
Let $`\varphi `$ be a positive principal eigenfunction of $`L_\mathrm{\Sigma }`$ and define $`\chi `$ by $`\psi =\chi \varphi `$. A computation shows that
$$L_\mathrm{\Sigma }\psi =\chi L_\mathrm{\Sigma }\varphi \varphi \mathrm{\Delta }_S\chi +2\left(\varphi s^AD^A\varphi \right)D_A\chi .$$
Since $`S`$ is strictly stably outermost, $`L_\mathrm{\Sigma }\varphi >0`$ and the strong maximum principle GT yields that $`\chi 0`$ if $`L_\mathrm{\Sigma }\psi 0`$, with strict inequality if $`L_\mathrm{\Sigma }\psi 0`$ somewhere. โ
We can now prove Theorem 2.
Proof of Theorem 2. By construction, the variation of $`\theta _l`$ vanishes along the vector $`q^\alpha _{x^\alpha }=_t+\mu _r=\omega l+\mu m`$ tangent to $`H`$, where $`\omega >0`$ is defined by $`_t=\omega l`$. Therefore,
$$0=\delta _q\theta _l=\delta _{\omega l+\mu m}\theta _l=\omega W+L_\mathrm{\Sigma }\mu .$$
If the null energy condition holds, $`W`$ is non-negative and under the condition in the second part of Theorem 2, $`W`$ is even positive somewhere. By Lemma 3 it follows that $`\mu 0`$ in the first case and $`\mu >0`$ everywhere in the second one, which proves the assertions. $`\mathrm{}`$
Assume that a 2-surface $`S`$ is strictly stably outermost and that the dominant energy condition is satisfied (i.e. that $`G_\beta ^\alpha u^\beta `$ is future directed for all future directed timelike vectors $`u^\alpha `$). Then $`S`$ is topologically $`S^2`$ (and hence the horizon through $`S`$ is $`S^2\times \mathrm{}`$). We recall here Newmanโs proof RN of this fact. Denoting by $`\varphi `$ the positive principal eigenfunction of $`L_\mathrm{\Sigma }`$, the stability condition implies that $`0<_S\varphi ^1L_\mathrm{\Sigma }\varphi `$. Then the result follows from (Local existence of dynamical and trapping horizons) after integrating by parts and using the Gauss-Bonnet theorem.
While we have restricted ourselves to local results in this Letter, it would clearly be desirable to determine the global evolution of the horizon in a given black hole spacetime. It would be much more ambitious to look at the global Cauchy evolution for asymptotically flat initial data with a MOTS. In vacuum, one expects this evolution to approach a Kerr spacetime and our horizon to approach the event horizon. Moreover, the area of the marginally trapped slices should approach a quantity not greater than $`16\pi m^2`$ where $`m`$ is the mass of the final Kerr black hole. This version of the Penrose inequality BHI would involve the area of the MOTS instead of a minimal surface and one could, for axially symmetric data, include angular momentum as well AK .
###### Acknowledgements.
We wish to thank Greg Galloway, Thomas Hoffmann-Ostenhof, Josรฉ Senovilla and Helmuth Urbantke for helpful discussions and correspondence.
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# Possible Optical Detection of the Anomalous X-ray Pulsar CXOU J010043.1-721134
## 1 Introduction
The anomalous X-ray pulsars (AXPs) are a class of neutron stars, numbering about half a dozen, which are radio-quiet, with periods of the order $`10`$s and estimated ages of $`10^3`$ to $`10^5`$yr. Like the soft gamma-ray repeaters, they are thought to be magnetars, whose emission is powered by the decay of a super-strong magnetic field ($`10^{15}`$G). See Woods & Thompson (2004) for a review of the known magnetars and their properties.
While energetically, the emission at X-ray energies dominates, optical and infrared photometry of AXPs is giving interesting constraints on the physical processes of the stellar magnetospheres. Particularly intriguing is that for the brightest object, 4U 0142+61, the optical spectral energy distribution is not just a power law. It shows, unique among neutron stars, a spectral break between V and B (Hulleman et al., 2004). Unfortunately, because of the uncertainty in the high amount of reddening, the precise shape cannot be measured.
In the magnetar model, the optical emission could be dues to ion cyclotron emission. If so, the spectral break should be a general feature (C. Thompson, 2004, priv comm.) due to the existence of a cooling radius in the magnetar magnetosphere from within which ions do not radiate (for a brief discussion, see Hulleman et al., 2004). The $`5`$ other AXPs known so far are, unfortunately, too highly reddened to be detected in V or B. Another prediction is that the spectra of different AXPs should be similar, but again uncertainties in the reddening do not allow us to test this (e.g. Durant & van Kerkwijk, 2005). As an alternative model, Eichler et al. (2002), considered the possibility of coherent optical and infrared emission from the lower magnetosphere of a magnetar, in analogy to some radio pulsar models. Unfortunately, no clear predictions for the spectral shape were made.
For the purposes of investigating the optical spectra of AXPs, the recent discovery of an AXP in the Small Magellanic Cloud (SMC), CXOU J010043.1$``$721134 (Lamb et al., 2002; Majid et al., 2004) is particularly interesting. It is the only AXP found so far, that is not confined to the disc of the Milky Way. The reddening to this source is, therefore, much less than for the other AXPs. Furthermore, its distance is relatively well known at 60.6(1.0) kpc (e.g. Hilditch et al., 2005). It thus presents a unique opportunity to study an AXP in the blue/UV.
## 2 Archival Observation and Analysis
Seeking imaging data on CXOU J010043.1$``$721134, we searched all the archives available t us. We found that the field was observed on 20 April 2004 with the Wide Field and Planetary Camera 2 (WFPC2) on board the Hubble Space Telescope (HST), as part of a snapshot programme for three-colour photometry of several patches of the SMC (Tolstoy, 1999). Single exposures were taken of 230s in the near-ultraviolet F300W, 180s in the โbroad Vโ F606W and 300s in Cousins I-like F814W filters. The position of our object of interest is on chip WF2 of the WFPC2 array.
We determined an astrometric solution by matching sources off the WF2 image to objects in the USNO B1.0 catalogue (Monet et al., 2003), and fitting for offset, rotation and scale. Eight stars were matched, after rejecting 7 objects which had poorly measured positions or which corresponded to multiple sources on the WF2 image. With these eight sources, the uncertainty in the astrometric fit is $`0\stackrel{}{\mathrm{.}}19/\sqrt{6}=0\stackrel{}{\mathrm{.}}08`$ in each co-ordinate for the F606W frame. The uncertainty in applying the astrometry to the other two bands was negligible in comparison. The systematic uncertainty in connecting the USNO astrometry to the International Celestial Reference System is $`0\stackrel{}{\mathrm{.}}2`$ in each co-ordinate, and the uncertainty in the Chandra position of CXOU J010043.1$``$721134 is a radius $`r=0\stackrel{}{\mathrm{.}}6`$ at 90% confidence. Note that the latter is from the nominal Chandra performance, despite being somewhat off-axis (Lamb et al., 2002). The above numbers, combined in quadrature, give a total uncertainty in the AXPโs position on our images of $`r=0\stackrel{}{\mathrm{.}}72`$ at 90% confidence. Photometry was performed using HSTphot 1.1 (Dolphin, 2000).
Figure 1 shows the F606W image of the field immediately around CXOU J010043.1$``$721134, with the positional error circle indicated. Stars X and Y have positions consistent with that of the AXP, with Star Z being a nearby, much brighter source. Their positions and magnitudes are listed in Table 1, and indicated in a colour-magnitude diagramme of all stars detected in the WFPC2 images in Figure 2.
From the photometry, Star Y is consistent with being a G5V star at the distance and reddening of the SMC, and Star Z an early B-type star. The colours and magnitudes of Star X do not correspond to any known stellar type, and make it a clear out-lier in Figure 2, suggesting a very blue, possibly hot object. Based on its position and unusual colours, we therefore consider Star X a likely counterpart to CXOU J010043.1$``$721134.
As a caveat, however , it should be remembered that this measurement is based on a single F606W exposure. The source in Figure 1 does not appear like a cosmic ray hit, and the HSTphot $`\chi `$ and $`Sharp`$ parameters are within reasonable limits for a point source: $`\chi =1.18`$ (goodness of fit parameter; reasonable values: $`<2.5`$) and $`Sharp=0.425`$ (where 0 corresponds to a stellar point spread function, positive values to more peaked profiles and negative values to more diffuse ones; reasonable values: -0.5โ 0.5). There were no bad pixels within Star Xโs profile.
In order to test the robustness of this identification, we estimated the likelihood of such a detection in an error circle of this size at any point on this chip of the detector. We searched the photometry for all objects in the F606W image which are classified as stellar ($`\chi <2.5`$, $`0.5<sharp<0.5`$), and with $`m_{606}m_{814}<0.3`$, i.e. at least as blue as Star X. Forty such objects are found on the same chip as Star X (including those near Star X in Figure 2), giving the probability of one falling within a circle of radius $`0\stackrel{}{\mathrm{.}}72`$ of $`1.5`$%. We note that the majority of these are within 3 pixels of brighter sources in the F814W image, and consequently were not detected in that band. Since this does not apply to Star X, and there do not appear to be any artifacts close to it (i.e. the location of Star X appears like sky in the F814W image), the chance of it being a false detection is somewhat smaller, but how much smaller is hard to quantify.
We also calculated the likelihood of our putative detection being due to a cosmic ray hit or instrumental effects. We searched for objects which are classified as stars in terms of their Chi and Sharp parameters as above, which were detected in F300W but not in F606W. We find thirteen objects, which implies that the probability of Star X being due to a cosmic rays hit or purely instrumental effects is 0.4% (after correcting for the difference in exposure time between F300W and F606W).
Seeking to confirm this detection, we searched other archives for optical images. We found a V-band image from the Wide Field Imager (WFI) on the 2.2m ESO telescope, La Silla, Chile. This demonstrated that the area of sky was very crowded, and extremely good seeing would be required to separate and securely detect Star X. In this case, the seeing was poor. We also obtained Gemini DDT observations with GMOS-S (Crampton & Murowinski, 2004) at Cerro Pachon, Chile. Unfortunately, the seeing was also not good enough in these images to distinguish between the sources in the crowded field. A proposal was also accepted at Magellan, Las Campanas, Chile, but conditions have not been good enough to obtain images so far. Unfortunately, the presence of Star Z means that only the most exceptional seeing conditions will allow further measurements of this object from the ground.
## 3 Discussion and Conclusions
Taking Star X as the true optical counterpart, CXOU J010043.1$``$721134 has an X-ray to optical flux ratio $`F_X/F_V=1.0\times 10^{13}/5.5\times 10^{15}=18`$ (un-absorbed X-ray flux in the 2โ10keV range from Woods & Thompson, 2004; visual $`\nu F_\nu `$ flux is de-reddened using $`A_V=0.3`$ \[Hilditch et al., 2005\], and assumes $`m_{606}=V`$). This compares with $`F_X/F_V=460`$ for 4U 0142$`+`$61 (for $`A_V=5.1`$, the nominal reddening), the only other AXP with an optical detection (Hulleman et al., 2004). Clearly the two ratios are very different.
It has been observed that infrared to X-ray flux ratios are similar for those AXPs with secure measurements (4U 0142$`+`$61, 1E 1048.1$``$5937 and 1E 2259$`+`$586; Durant & van Kerkwijk, 2005). Variations have, however, been observed to be very large, of orders of magnitude in some cases. For example the transient AXP XTE J1810-197 (Ibrahim et al, 2004) increased dramatically in both X-rays and infrared flux before slowly dimming again.
It is possible that the difference in V-band to X-ray flux ratio above arises because the measurements for CXOU J010043.1$``$721134 were not simultaneous; the AXP could have been brighter by a large factor at the time of the HST observation. CXOU J010043.1$``$721134 was observed to be $`sim50`$% brighter in X-rays by Majid et al. (2004) than Lamb et al. (2002), but they attribute this to the different instruments used to make the observations rather than genuine variability. 4U 0142$`+`$61 has been the most stable of the AXPs in both X-ray and optical flux (Hulleman et al., 2004). This could, in principle, mean that the intrinsic spectra of the two objects are very different, possibly indicating differing magnetic field configurations.
The limit in F814W already provides some constraints on the shape of the optical spectrum. Whilst a Rayleigh-Jeans form $`\nu F_\nu \nu ^n`$, $`n=3`$ is possible, a flat spectrum ($`n=0`$) is excluded. The 90% confidence limit is $`n2`$. Since the spectrum should not increase steeper than Rayleigh-Jeans (in the absence of an emission feature), we predict that the I-band magnitude is not much below the limit we have established. The F300W limit is not constraining in this respect.
In summary, we present Star X, with $`m_{606}=24.19(15)`$, as the probable optical counterpart to CXOU J010043.1$``$721134. It is at the right location and has colours unlike normal stellar sources. Although based on a detection in a single exposure, HSTphot diagnostics point to it being a real detection, with only a $`1.5`$% probability of a false detection. If confirmed, this discovery will enable the measurement of AXP properties in the blue and UV.
Acknowledgements: This work made use of archival observations made with the NASA/ESA Hubble Space Telescope and with observations from ESO Telescopes at the La Silla Observatories. We thank Slavek Rucinski and the Gemini Observatories for attempting follow-up observations. We thank an anonymous referee for very useful comments which much improved the presentation of our results. We acknowledge financial support from NSERC.
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# Particle Acceleration and the Production of Relativistic Outflows in Advection-Dominated Accretion Disks with Shocks
## 1 INTRODUCTION
A large body of observational evidence has established that extragalactic relativistic jets are commonly associated with radio-loud active galactic nuclei (AGNs), which may contain hot, advection-dominated accretion disks. However, the precise nature of the mechanism responsible for transferring the gravitational potential energy from the infalling matter to the small population of nonthermal particles that escape to form the jet is not yet clear (see, e.g., Livio 1999). The most promising jet acceleration scenarios proposed so far are the Blandford-Znajek mechanism (Blandford & Znajek 1977) and the electromagnetic cocoon model (Lovelace 1976; Blandford & Payne 1982), both of which involve the extraction of energy from the rotation of the black hole in order to power the outflow. While conceptually attractive, one finds that the complex physics involved in these models tends to obscure the nature of the fundamental microphysical processes. In particular, the introduction of the relativistic particles that escape to form the jet is usually made in an ad hoc manner without any reference to the dynamics of the associated accretion disk, although recent magnetohydrodynamical simulations carried out by De Villiers et al. (2005) and McKinney & Gammie (2004) have achieved a higher level of self-consistency. Given the relative complexity of the electromagnetic models, it is natural to ask whether the outflows can be explained in terms of well-understood microphysical processes operating in the hot, tenuous disk, such as the possible acceleration of the jet particles at a standing accretion shock.
It has been known for some time that inviscid accretion disks can display both shocked and shock-free (i.e., smooth) solutions depending on the values of the energy and angular momentum per unit mass in the gas supplied at a large radius (e.g., Chakrabarti 1989a; Chakrabarti & Molteni 1993; Kafatos & Yang 1994; Lu & Yuan 1997; Das, Chattopadhyay, & Chakrabarti 2001). Shocks can also exist in viscous disks if the viscosity is relatively low (Chakrabarti 1996; Lu, Gu, & Yuan 1999), although smooth solutions are always possible for the same set of upstream parameters (Narayan, Kato, & Honma 1997; Chen, Abramowicz, & Lasota 1997). Hawley, Smarr, & Wilson (1984a, 1984b) have shown through general relativistic simulations that if the gas is falling with some rotation, then the centrifugal force can act as a โwall,โ triggering the formation of a shock. Furthermore, the possibility that shock instabilities may generate the quasi-periodic oscillations (QPOs) observed in some sources containing black holes has been pointed out by Chakrabarti, Acharyya, & Molteni (2004), Lanzafame, Molteni, & Chakrabarti (1998), Molteni, Sponholz, & Chakrabarti (1996), and Chakrabarti & Molteni (1995). Nevertheless, shocks are โoptionalโ even when they are allowed, and one is always free to construct models that avoid them. However, in general the shock solution possesses a higher entropy content than the shock-free solution, and therefore we argue based on the second law of thermodynamics that when possible, the shocked solution represents the preferred mode of accretion (Becker & Kazanas 2001; Chakrabarti & Molteni 1993).
Our primary objective in this paper is to explore the consequences of the presence of a shock in an ADAF disk for the acceleration of the nonthermal particles in the observed jets. The question of whether or not viscosity needs to be included in the disk model is difficult to answer in general. Several authors have shown that shock solutions are possible in viscous (e.g., Chakrabarti 1990, 1996; Lu, Gu, & Yuan 1999; Chakrabarti & Das 2004) as well as inviscid disks (e.g., Chakrabarti 1989a, 1989b; Abramowicz & Chakrabarti 1990; Yang & Kafatos 1995; Chakrabarti 1996; Das, Chattopadhyay, & Chakrabarti 2001). In particular, Chakrabarti (1990) and Chakrabarti & Das (2004) demonstrated that shocks can exist in viscous disks if the angular momentum and the viscosity are relatively low. Since the acceleration of particles in shocked disks has never been investigated before, in this first study we shall focus on inviscid flows containing isothermal shocks (e.g., Chakrabarti 1989a; Kafatos & Yang 1994; Lu & Yuan 1997), while deferring the treatment of viscous disks to future work. However, it is clearly important to address the potential connection between this idealized, inviscid calculation and the physical properties of real accretion disks, which undoubtedly have nonzero viscosity. We argue that the results presented here should be qualitatively similar to those obtained in a viscous disk provided a shock is present, in which case efficient first-order Fermi acceleration is expected to occur. While the possible existence of standing shocks in viscous disks is a controversial issue at the present time, we believe that the work of Chakrabarti (1990, 1996), Lu, Gu, & Yuan (1999), and Chakrabarti & Das (2004) provides sufficient support for the possibility to motivate the present investigation.
Although the effect of a standing shock in heating the gas in the post-shock region has been examined by a number of previous authors for both viscid (Chakrabarti & Das 2004; Lu, Gu, & Yuan 1999; Chakrabarti 1990) and inviscid (e.g., Lu & Yuan 1997, 1998; Yang & Kafatos 1995; Abramowicz & Chakrabarti 1990) disks, the implications of the shock for the acceleration of nonthermal particles in the disk have not been considered in detail before. However, a great deal of attention has been focused on particle acceleration in the vicinity of supernova-driven shock waves as a possible explanation for the observed cosmic-ray energy spectrum (Blandford & Ostriker 1978; Jones & Ellison 1991). In the present paper we consider the analogous process occurring in hot, advection-dominated accretion flows (ADAFs) around black holes. These disks are ideal sites for first-order Fermi acceleration at shocks because the plasma is collisionless and therefore a small fraction of the particles can gain a great deal of energy by repeatedly crossing the shock. Shock acceleration in the disk therefore provides an intriguing possible explanation for the powerful outflows of relativistic particles observed in many radio-loud systems (Le & Becker 2004).
The dynamical model for the disk/shock/outflow employed here was discussed by Le & Becker (2004), who demonstrated that the predicted kinetic power in the jets agrees with the observational estimates for M87 and Sgr A$`^{}`$. Here we present a more detailed development of the dynamical model, including a careful examination of the implications of the shock acceleration process for the evolution of the relativistic particle distribution in the disk and the jet. The number and energy densities of the relativistic particles are determined along with the hydrodynamical structure of the disk in a self-consistent manner by solving the fluid dynamical conservation equations and the transport equation simultaneously using a rigorous mathematical approach. In this sense, the model presented here represents a new type of synthesis between studies of accretion dynamics and particle transport.
The remainder of the paper is organized as follows. In ยง 2 we discuss the ADAF model assumptions and the possibility of shock acceleration in ADAF disks, and the general structure of the disk/shock model is examined in ยง 3. The isothermal shock jump conditions and the asymptotic variations of the physical parameters at both large and small radii are discussed in ยง 4. In ยง 5 we analyze the steady-state transport equation governing the distribution of the relativistic particles in the disk and the jet. Solutions for the number and energy density distributions of the relativistic particles are obtained in ยง 6, and detailed applications to the disks/outflows in M87 and Sgr A$`^{}`$are presented in ยง 7. The astrophysical implications of our results are discussed in ยง 8.
## 2 MODEL BACKGROUND
Accretion onto a black hole involves differentially-rotating flows in which the viscosity plays an essential role in transporting angular momentum outward, thereby allowing the accreting gas to spiral in toward the central mass (Pringle 1981). In the ADAF model, it is assumed that the mass accretion rate is much smaller than the Eddington rate,
$$\dot{M}_\mathrm{E}c^2\beta ^1L_\mathrm{E}=2.2\times 10^9\beta ^1\left(\frac{M}{M_{}}\right)\mathrm{M}_{}\mathrm{yr}^1,$$
(1)
where the efficiency parameter $`\beta `$ is of order $`10\%`$, and the Eddington luminosity is defined by $`L_\mathrm{E}4\pi GMm_pc/\sigma __\mathrm{T}`$ for pure, fully-ionized hydrogen, with $`\sigma __\mathrm{T}`$, $`M`$, $`m_p`$, and $`c`$ denoting the Thomson cross section, the black-hole mass, the proton mass, and the speed of light, respectively. Due to the sub-Eddington accretion rates in these systems, the plasma is rather tenuous, and this strongly inhibits the efficiency of two-body radiative processes such as free-free emission. The gas is therefore unable to cool effectively within an accretion time, and consequently the gravitational potential energy dissipated by viscosity is stored in the gas as thermal energy instead of being radiated away (e.g., Narayan, Kato, & Honma 1997). The low density also reduces the level of Coulomb coupling between the ions and the electrons, resulting in a two-temperature configuration with the ion temperature ($`T_i10^{12}\mathrm{K}`$) close to the virial value, and a much lower electron temperature ($`T_e10^9\mathrm{K}`$). In this scenario, most of the energy is advected across the horizon into the black hole, and the resulting X-ray luminosity is far below the Eddington value (Becker & Le 2003; Becker & Subramanian 2005).
When the ion temperature is close to the virial temperature, as in ADAFs, the disk is gravitationally unbound (e.g., Narayan, Kato, & Honma 1997; Blandford & Begelman 1999; Becker, Subramanian, & Kazanas 2001). It follows that the original ADAF model was not entirely self-consistent since it neglected outflows. This motivated Blandford & Begelman (1999) to propose the self-similar advection-dominated inflow-outflow solution (ADIOS) to address the question of self-consistency by including the possibility of powerful winds that carry away mass, energy, and angular momentum. In this Newtonian, nonrelativistic model, the dynamical solutions are not applicable near the event horizon, and therefore the ADIOS approach cannot be used to obtain a global understanding of the disk structure. This led Becker, Subramanian, & Kazanas (2001) to modify the ADIOS scenario to include general relativistic effects by replacing the Newtonian potential with the pseudo-Newtonian form (Paczyลski & Wiita 1980)
$$\mathrm{\Phi }(r)\frac{GM}{rr__\mathrm{S}},$$
(2)
where $`r__\mathrm{S}=2GM/c^2`$ is the Schwarzschild radius for a black hole of mass $`M`$. This modified model is known as the self-similar relativistic advection-dominated inflow-outflow solution (RADIOS). Despite the success of the self-similar RADIOS model in describing the general features of the disk/outflow structure, it does not provide a comprehensive picture since no explicit microphysical acceleration mechanism is included. It is therefore natural to explore possible extensions to the ADAF scenario that incorporate a concrete acceleration mechanism capable of powering the outflows.
The idea of shock acceleration in the environment of AGNs was first suggested by Blandford & Ostriker (1978). Subsequently, Protheroe & Kazanas (1983) and Kazanas & Ellison (1986) investigated shocks in spherically-symmetric accretion flows as a possible explanation for the energetic radiation emitted by many AGNs. However, in these papers the acceleration of the particles was studied without the benefit of a detailed transport equation, and the assumption of spherical symmetry precludes the treatment of acceleration in disks. The state of the theory was advanced by Webb & Bogdan (1987) and Spruit (1987), who employed a transport equation to solve for the distribution of energetic particles in a spherical accretion flow characterized by a self-similar velocity profile terminating at a standing shock. While more quantitative in nature than the earlier models, these solutions are not applicable to disks since the geometry is spherical and the velocity distribution is inappropriate. Hence none of these previous models can be used to develop a single, global, self-consistent picture for the acceleration of relativistic particles in an accretion disk containing a shock.
The success of the diffusive (first-order Fermi) shock acceleration model in the cosmic-ray context suggests that the same mechanism may be responsible for powering the outflows commonly observed in radio-loud systems containing black holes. As a preliminary step in evaluating the potential relevance of shock acceleration as a possible explanation for the observed outflows, we need to consider the basic physical properties of the hot plasma in ADAF disks. One of the critical issues for determining the efficiency of shock acceleration in accretion disks is the role of particle-particle collisions in thermalizing the high-energy ions. The mean free path for ion-ion collisions is given in cgs units by (Subramanian, Becker, & Kafatos 1996)
$$\lambda _{ii}=1.8\times 10^5\frac{T_i^2}{N_i\mathrm{ln}\mathrm{\Lambda }},$$
(3)
where $`N_i`$ and $`T_i`$ denote the thermal ion number density and temperature, respectively, and $`\mathrm{ln}\mathrm{\Lambda }`$ is the Coulomb logarithm. In ADAF disks, $`\lambda _{ii}`$ greatly exceeds the vertical thickness of the disk, and therefore the shock and the flow in general are collisionless. However, the mean free path $`\lambda _{\mathrm{mag}}`$ for collisions between ions and magnetohydrodynamical (MHD) waves is much shorter than $`\lambda _{ii}`$ for the thermal particles, and it is much longer than $`\lambda _{ii}`$ for the relativistic particles (Ellison & Eichler 1984; Subramanian, Becker, & Kafatos 1996). The increase in $`\lambda _{\mathrm{mag}}`$ with increasing particle energy reflects the fact that the high-energy particles will interact only with the highest-energy MHD waves. The low-energy background particles therefore tend to thermalize the energy they gain in crossing the shock due to collisions with magnetic waves. Conversely, the relativistic particles are able to diffuse back and forth across the shock many times, gaining a great deal of energy while avoiding thermalization due to the longer mean free path.
The probability of multiple shock crossings decreases exponentially with the number of crossings, and the mean energy of the particles increases exponentially with the number of crossings. This combination of factors naturally gives rise to a power-law energy distribution, which is a general characteristic of Fermi processes (Fermi 1954). Two effects limit the maximum energy that can be achieved by the particles. First, at very high energies the particles will tend to lose energy to the waves due to recoil. Second, the mean free path $`\lambda _{\mathrm{mag}}`$ will eventually exceed the thickness of the disk as the particle energy is increased, resulting in escape from the disk without further acceleration.
## 3 TRANSONIC FLOW STRUCTURE
As discussed in ยง 1, various authors have established that shocks can exist in both viscid and inviscid disks. In this first study of particle acceleration in shocked disks, we shall focus on the inviscid case because it is the most straightforward to analyze from a mathematical viewpoint, and also because it serves to illustrate the basic physical principles involved. Moreover, we expect that the results obtained in the viscous case will be qualitatively similar to those presented here since efficient Fermi acceleration will occur whether or not viscosity is present, provided the flow contains a shock. The equations governing the disk structure can yield solutions that include three possible types of standing shocks, namely (i) Rankine-Hugoniot shocks, where the effective cooling processes are so inefficient that no energy is lost from the surface of the disk, (ii) isentropic shocks, where the entropy generated at the shock is comparable to the amount radiated away, and (iii) isothermal shocks, where the cooling processes are so efficient that the post-shock sound speed and disk thickness remain the same as the pre-shock values. In the isothermal case, the shock must radiate away both energy and entropy through the upper and lower surfaces of the disk (e.g., Chakrabarti 1989a, 1989b; Abramowitz $`\&`$ Chakrabarti 1990). This renders the isothermal shock model particularly useful from the point of view of modeling outflows, since the energy lost from the shock can be identified with that powering the jet. On the other hand, Rankine-Hugoniot shocks cannot be used if we are interested in any kind of escape. The isentropic shock is an intermediate case. In this paper, we shall focus exclusively on the isothermal shock model since this case provides the strongest potential connection with the observed outflows.
The model considered here is depicted schematically in Figure 1. In this scenario, the gas is accelerated gravitationally toward the central mass, and experiences a shock transition due to an obstruction near the event horizon. The obstruction is provided by the โcentrifugal barrier,โ which is located between the inner and outer sonic points. Particles from the high-energy tail of the background Maxwellian distribution are accelerated at the shock discontinuity via the first-order Fermi mechanism, resulting in the formation of a nonthermal, relativistic particle distribution in the disk. The spatial transport of the energetic particles within the disk is a stochastic process based on a three-dimensional random walk through the accreting background gas. Consequently, some of the accelerated particles diffuse to the disk surface and become unbound, escaping through the upper and lower edges of the cylindrical shock to form the outflow, while others diffuse outward radially through the disk or advect across the event horizon into the black hole.
In order to analyze the connection between the disk/shock model and the transport/acceleration of the relativistic particles, we consider the set of physical conservation equations employed by Chakrabarti (1989a) and Abramowicz & Chakrabarti (1990), who investigated the structure of a one-dimensional, steady-state, axisymmetric, inviscid accretion flow based on the vertically-averaged conservation equations. The effects of general relativity are incorporated in an approximate manner by utilizing the pseudo-Newtonian form for the gravitational potential per unit mass given by equation (2). The use of such a potential allows one to investigate the complicated physical processes taking place in the accretion disk within the context of a semi-classical framework while maintaining good agreement with fully relativistic calculations (see, e.g., Narayan, Kato, & Honma 1997; Becker & Subramanian 2005). The pseudo-Newtonian potential correctly reproduces the radius of the event horizon, the marginally bound orbit, and the marginally stable orbit (Paczyลski $`\&`$ Wiita 1980). Furthermore, the dynamics of freely-falling particles near the event horizon computed using this potential agree perfectly with the results obtained using the Schwarzschild metric, although time dilation is not included (Becker & Le 2003).
### 3.1 Transport Rates
Becker & Le (2003) and Becker & Subramanian (2005) demonstrated that three integrals of the flow are conserved in viscous ADAF disks, namely, the mass transport rate
$$\dot{M}=4\pi rH\rho v,$$
(4)
the angular momentum transport rate
$$\dot{J}=\dot{M}r^2\mathrm{\Omega }๐ข,$$
(5)
and the energy transport rate
$$\dot{E}=๐ข\mathrm{\Omega }+\dot{M}\left(\frac{1}{2}v_\varphi ^2+\frac{1}{2}v^2+\frac{P+U}{\rho }+\mathrm{\Phi }\right),$$
(6)
where $`\rho `$ is the mass density, $`v`$ is the radial velocity (defined to be positive for inflow), $`\mathrm{\Omega }`$ is the angular velocity, $`๐ข`$ is the torque, $`H`$ is the disk half-thickness, $`v_\varphi =r\mathrm{\Omega }`$ is the azimuthal velocity, $`U`$ is the internal energy density, and $`P=(\gamma 1)U`$ is the pressure. Each of the various quantities represents a vertical average over the disk structure. We also assume that the ratio of specific heats, $`\gamma `$, maintains a constant value throughout the flow. Note that all of the transport rates $`\dot{M}`$, $`\dot{J}`$, and $`\dot{E}`$ are defined to be positive for inflow.
The torque $`๐ข`$ is related to the gradient of the angular velocity $`\mathrm{\Omega }`$ via the usual expression (e.g., Frank et al. 2002)
$$๐ข=4\pi r^3H\rho \nu \frac{d\mathrm{\Omega }}{dr},$$
(7)
where $`\nu `$ is the kinematic viscosity. The disk half-thickness $`H`$ is given by the standard hydrostatic prescription
$$H(r)=\frac{a}{\mathrm{\Omega }_\mathrm{K}},$$
(8)
where $`a`$ represents the adiabatic sound speed, defined by
$$a(r)\left(\frac{\gamma P}{\rho }\right)^{1/2},$$
(9)
and $`\mathrm{\Omega }_\mathrm{K}`$ denotes the Keplerian angular velocity of matter in a circular orbit at radius $`r`$ in the pseudo-Newtonian potential (eq. ), defined by
$$\mathrm{\Omega }_\mathrm{K}^2\frac{GM}{r(rr__\mathrm{S})^2}=\frac{1}{r}\frac{d\mathrm{\Phi }}{dr}.$$
(10)
The quantities $`\dot{M}`$ and $`\dot{J}`$ are constant throughout the flow, and therefore they represent the rates at which mass and angular momentum, respectively, enter the black hole. The energy transport rate $`\dot{E}`$ generally remains constant, although it will jump at the location of an isothermal shock if one is present in the disk. We can eliminate the torque $`๐ข`$ between equations (5) and (6) and combine the result with equation (9) to express the energy transport per unit mass as
$$ฯต\frac{\dot{E}}{\dot{M}}=\frac{1}{2}v^2\frac{1}{2}\frac{\mathrm{}^2}{r^2}+\frac{\mathrm{}_0\mathrm{}}{r^2}+\frac{a^2}{\gamma 1}+\mathrm{\Phi },$$
(11)
where $`\mathrm{}(r)r^2\mathrm{\Omega }(r)`$ and $`\mathrm{}_0\dot{J}/\dot{M}`$ denote the specific angular momentum at radius $`r`$ and the (constant) angular momentum transport per unit mass, respectively. Note that equation (11) is valid for both viscid and inviscid flows.
### 3.2 Inviscid Flow Equations
In the present application, viscosity is neglected, and therefore $`๐ข=0`$ and the specific angular momentum is given by
$$\mathrm{}(r)=\mathrm{}_0=\mathrm{constant}$$
(12)
throughout the disk. It follows that the flow is purely adiabatic, except at the location of a possible isothermal shock (Becker & Le 2003). In the inviscid case, equation (11) reduces to
$$ฯต=\frac{1}{2}v^2+\frac{1}{2}\frac{\mathrm{}^2}{r^2}+\frac{a^2}{\gamma 1}+\mathrm{\Phi }.$$
(13)
The resulting disk/shock model depends on three free parameters, namely the energy transport per unit mass $`ฯต`$, the specific heats ratio $`\gamma `$, and the specific angular momentum $`\mathrm{}`$. The value of $`ฯต`$ will jump at the location of an isothermal shock if one exists in the disk, but the value of $`\mathrm{}`$ remains constant throughout the flow. This implies that the specific angular momentum of the particles escaping through the upper and lower surfaces of the cylindrical shock must be equal to the average value of the specific angular momentum for the particles remaining in the disk, and therefore the outflow exerts no torque on the disk (Becker, Subramanian, & Kazanas 2001). Since the flow is purely adiabatic in the absence of viscosity, the pressure and density variations are coupled according to
$$P=D_0\rho ^\gamma ,$$
(14)
where $`D_0`$ is a parameter related to the specific entropy that remains constant except at the location of the isothermal shock if one is present.
By combining equations (4), (8), (9), (10), and (14), we find that the quantity
$$Kr^{3/2}(rr__\mathrm{S})va^{(\gamma +1)/(\gamma 1)}$$
(15)
is conserved throughout an adiabatic disk, except at the location of an isothermal shock. Following Becker & Le (2003), we refer to $`K`$ as the โentropy parameter,โ and we note that the entropy per particle $`S`$ is related to $`K`$ via
$$S=k\mathrm{ln}K+c_0,$$
(16)
where $`k`$ is the Boltzmann constant and $`c_0`$ is a constant that depends on the composition of the gas but is independent of its state.
### 3.3 Critical Conditions
By combining equations (13) and (15), one can solve for the flow velocity $`v`$ as a function of $`r`$ using a simple root-finding procedure. However, in order to understand the implications of the transonic (critical) nature of the accretion flow, we must also analyze the properties of the โwind equation,โ which is the first-order differential equation governing the flow velocity $`v`$. By differentiating equation (13) with respect to $`r`$, we obtain the steady-state radial momentum equation
$$v\frac{dv}{dr}=\frac{\mathrm{}^2}{r^3}\frac{d\mathrm{\Phi }}{dr}\left(\frac{2a}{\gamma 1}\right)\frac{da}{dr}.$$
(17)
The derivative of the sound speed appearing on the right-hand side of this expression can be evaluated by using equations (10) and (15) to write
$$\frac{1}{a}\frac{da}{dr}=\left(\frac{\gamma 1}{\gamma +1}\right)\left[\frac{1}{\mathrm{\Omega }_\mathrm{K}}\frac{d\mathrm{\Omega }_\mathrm{K}}{dr}\frac{1}{v}\frac{dv}{dr}\frac{1}{r}\right].$$
(18)
We can now construct the wind equation by combining equations (10), (17), and (18), which yields
$$\frac{1}{v}\frac{dv}{dr}=\frac{N}{D},$$
(19)
where the numerator and denominator functions $`N`$ and $`D`$ are given by
$$N=\frac{GM}{(rr__\mathrm{S})^2}\frac{\mathrm{}^2}{r^3}+\frac{a^2}{\gamma +1}\left[\frac{3r__\mathrm{S}5r}{r(rr__\mathrm{S})}\right],D=\frac{2a^2}{\gamma +1}v^2.$$
(20)
The simultaneous vanishing of $`N`$ and $`D`$ yields the critical conditions
$$\frac{GM}{(r_cr__\mathrm{S})^2}\frac{\mathrm{}^2}{r_c^3}+\frac{a_c^2}{\gamma +1}\left[\frac{3r__\mathrm{S}5r_c}{r_c(r_cr__\mathrm{S})}\right]=0,$$
(21)
and
$$\frac{2a_c^2}{\gamma +1}v_c^2=0,$$
(22)
where $`v_c`$ and $`a_c`$ denote the values of the velocity and the sound speed at the critical radius, $`r=r_c`$.
### 3.4 Critical Point Analysis
Equations (21) and (22) can be solved simultaneously to express $`v_c`$ and $`a_c`$ as explicit functions of the critical radius $`r_c`$, which yields
$$v_c^2=2\left[\frac{GMr_c^3\mathrm{}^2(r_cr__\mathrm{S})^2}{(5r_c3r__\mathrm{S})(r_cr__\mathrm{S})r_c^2}\right],$$
(23)
and
$$a_c^2=(\gamma +1)\left[\frac{GMr_c^3\mathrm{}^2(r_cr__\mathrm{S})^2}{(5r_c3r__\mathrm{S})(r_cr__\mathrm{S})r_c^2}\right].$$
(24)
The corresponding value of the entropy parameter $`K`$ at the critical point is given by (see eq. )
$$K_c=r_c^{3/2}(r_cr__\mathrm{S})v_ca_c^{(\gamma +1)/(\gamma 1)}.$$
(25)
By using equations (23) and (24) to substitute for $`v`$ and $`a`$ in equation (13), we can express the energy transport parameter $`ฯต`$ in terms of $`r_c`$, $`\mathrm{}`$, and $`\gamma `$, obtaining
$$ฯต=\frac{1}{2}\frac{\mathrm{}^2}{r_c^2}\frac{GM}{r_cr__\mathrm{S}}+\frac{2\gamma }{\gamma 1}\left[\frac{GMr_c^3\mathrm{}^2(r_cr__\mathrm{S})^2}{r_c^2(r_cr__\mathrm{S})(5r_c3r__\mathrm{S})}\right].$$
(26)
This expression can be rewritten as a quartic equation for $`r_c`$ of the form
$$๐ฉr_c^4๐ชr_c^3+๐ซr_c^2๐ฌr_c+=0,$$
(27)
where
$`๐ฉ`$ $`=`$ $`5ฯต,๐ช=16ฯต3+{\displaystyle \frac{2}{\gamma 1}},`$
$`๐ซ`$ $`=`$ $`12ฯต+{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{5\gamma }{\gamma 1}}\right)\mathrm{}^26,`$ (28)
$`๐ฌ`$ $`=`$ $`\left({\displaystyle \frac{8}{\gamma 1}}\right)\mathrm{}^2,`$
$``$ $`=`$ $`\left({\displaystyle \frac{2\gamma +6}{\gamma 1}}\right)\mathrm{}^2,`$
and we have utilized natural gravitational units with $`GM=c=1`$ and $`r__\mathrm{S}=2`$. These equations agree with the corresponding results derived by Das, Chattopadhyay, & Chakrabarti (2001). The four solutions for $`r_c`$ in terms of the three fundamental parameters $`ฯต`$, $`\mathrm{}`$, and $`\gamma `$ can be obtained analytically using the standard formulas for quartic equations (e.g., Abramowitz & Stegun 1970). We refer to the roots using the notation $`r_{c1}`$, $`r_{c2}`$, $`r_{c3}`$, and $`r_{c4}`$ in order of decreasing radius.
The critical radius $`r_{c4}`$ always lies inside the event horizon and is therefore not physically relevant, but the other three are located outside the horizon. The type of each critical point is determined by computing the two possible values for the derivative $`dv/dr`$ at the corresponding location using LโHรดpitalโs rule and then checking to see whether they are real or complex. We find that both values are complex at $`r_{c2}`$, and therefore this is an O-type critical point, which does not yield a physically acceptable solution. The remaining roots $`r_{c1}`$ and $`r_{c3}`$ each possess real derivatives, and are therefore physically acceptable sonic points, although the types of accretion flows that can pass through them are different. Specifically, $`r_{c3}`$ is an X-type critical point, and therefore a smooth, global, shock-free solution always exists in which the flow is transonic at $`r_{c3}`$ and then remains supersonic all the way to the event horizon. On the other hand, $`r_{c1}`$ is an $`\alpha `$-type critical point, and therefore any accretion flow originating at a large distance that passes through this point must display a shock transition below $`r_{c1}`$ (Abramowicz & Chakrabarti 1990). After crossing the shock, the subsonic gas must pass through another ($`\alpha `$-type) critical point before it enters the black hole since the flow has to be supersonic at the event horizon (Weinberg 1972).
### 3.5 Shock-Free Solutions
Even when a shock can exist in the flow, it is always possible to construct a globally smooth (shock-free) solution using the same set of parameters. Smooth flows must pass through the inner critical point located at radius $`r_{c3}`$, which is calculated using the quartic equation (27) for given values of $`ฯต`$, $`\mathrm{}`$, and $`\gamma `$. The corresponding values for the critical velocity, $`v_{c3}`$, the critical sound speed, $`a_{c3}`$, and the critical entropy, $`K_{c3}`$, are computed using equations (23), (24), and (25), respectively. Since the flows treated here are inviscid, they have a conserved value for the entropy parameter $`K`$ (eq. ) unless a shock is present. Hence in a smooth, shock-free flow, the value of $`K`$ is everywhere equal to the critical value $`K_{c3}`$. The structure of the velocity profile in a shock-free disk can therefore be determined using a simple root-finding procedure as follows. By eliminating the sound speed $`a`$ between equations (13) and (15), we obtain the equivalent expression
$$ฯต=\frac{1}{2}v^2+\frac{1}{2}\frac{\mathrm{}^2}{r^2}+\frac{1}{\gamma 1}\left[\frac{K_{c3}^2}{r^3(rr__\mathrm{S})^2v^2}\right]^{(\gamma 1)/(\gamma +1)}+\mathrm{\Phi },$$
(29)
where we have set $`K=K_{c3}`$. In general, at any radius $`r`$, equation (29) yields one subsonic root and one supersonic root for the velocity. The subsonic solution is chosen for $`r>r_{c3}`$, and the supersonic solution is selected for $`r<r_{c3}`$. Once the velocity profile $`v(r)`$ has been computed, we can obtain the corresponding sound speed distribution $`a(r)`$ by utilizing equation (15) with $`K=K_{c3}`$. Note that the velocity and sound speed solutions can also be calculated by integrating the wind equation (19) numerically, and the results obtained using this approach agree with the root-finding method.
## 4 ISOTHERMAL SHOCK MODEL
Our primary goal in this paper is to analyze the acceleration of relativistic particles due to the presence of a standing, isothermal shock in an accretion disk. Hence we are interested in flows that pass smoothly through the outer critical radius $`r_{c1}`$, and then experience a velocity discontinuity at the shock location, which we refer to as $`r_{}`$. In order to form self-consistent global models, we will first need to understand how the structure of the disk responds to the presence of a shock. This requires analysis of the shock jump conditions, which are based on the standard fluid dynamical conservation equations. Since shocks are always optional even when they can occur, we will compare our results for the relativistic particle acceleration with those obtained when there is no shock and the flow is globally smooth.
The values of the energy transport parameter $`ฯต`$ on the upstream and downstream sides of the isothermal shock at $`r=r_{}`$ are denoted by $`ฯต_{}`$ and $`ฯต_+`$, respectively. Note that $`ฯต_{}>ฯต_+`$ as a consequence of the loss of energy through the upper and lower surfaces of the disk at the shock location. It is important to emphasize that the drop in $`ฯต`$ at the shock has the effect of altering the transonic structure of the flow in the post-shock region. Hence, although the post-shock flow must pass through another critical point and become supersonic before crossing the event horizon, the new (inner) critical point is not equal to any of the four roots computed using the upstream energy transport parameter $`ฯต_{}`$. Instead, the new inner critical radius, which we refer to as $`\widehat{r}_{c3}`$, must be computed using the downstream value of the energy transport parameter, $`ฯต_+`$. We point out that the total energy inflow rate across the horizon, including the rest-mass contribution, must be positive since no energy can escape from the black hole, and therefore we require that $`c^2+ฯต_+>0`$.
### 4.1 Isothermal Shock Jump Conditions
We shall assume that the escape of the relativistic particles from the disk results in a negligible amount of mass loss because the Lorentz factor of the escaping particles is much greater than unity (see Table 2). This is confirmed ex post facto by comparing the rate of mass loss, $`\dot{M}_{\mathrm{loss}}`$, with the accretion rate $`\dot{M}`$. We find that for the models analyzed here, $`\dot{M}_{\mathrm{loss}}/\dot{M}\stackrel{<}{}10^3`$, and therefore our assumption of negligible mass loss is justified. Hence the accretion rate $`\dot{M}`$ is conserved as the gas crosses the shock, which is represented by the condition
$$\mathrm{\Delta }\dot{M}\underset{\epsilon 0}{lim}\dot{M}(r_{}\epsilon )\dot{M}(r_{}+\epsilon )=0,$$
(30)
where the symbol โ$`\mathrm{\Delta }`$โ will be used to denote the difference between post- and pre-shock quantities. The specific angular momentum $`\mathrm{}\dot{J}/\dot{M}`$ is also conserved throughout the flow, and therefore we find that
$$\mathrm{\Delta }\dot{J}=0.$$
(31)
Furthermore, the radial momentum transport rate, defined by
$$\dot{I}4\pi rH(P+\rho v^2),$$
(32)
must remain constant across the shock, and consequently
$$\mathrm{\Delta }\dot{I}=0.$$
(33)
Based on equations (13) and (31), we find that the jump condition for the energy transport rate $`\dot{E}`$ is given by
$$\mathrm{\Delta }\dot{E}=\dot{M}\left(\frac{1}{2}\mathrm{\Delta }v^2+\frac{1}{\gamma 1}\mathrm{\Delta }a^2\right).$$
(34)
Equations (4), (8), and (13) can be combined with equations (30), (33), and (34) to obtain
$`\rho _+v_+a_+`$ $`=`$ $`\rho _{}v_{}a_{},`$ (35)
$`a_+P_+a_{}P_{}`$ $`=`$ $`a_{}\rho _{}v_{}^2a_+\rho _+v_+^2,`$ (36)
$`ฯต__+ฯต_{_{}}`$ $`=`$ $`{\displaystyle \frac{v_+^2v_{}^2}{2}}+{\displaystyle \frac{a_+^2a_{}^2}{\gamma 1}},`$ (37)
where the subscripts โ-โ and โ+โ refer to quantities measured just upstream and just downstream from the shock, respectively. In the case of an isothermal shock, $`a_+=a_{}`$, and therefore the shock jump conditions reduce to
$`\rho _+v_+`$ $`=`$ $`\rho _{}v_{},`$ (38)
$`P_+P_{}`$ $`=`$ $`\rho _{}v_{}^2\rho _+v_+^2,`$ (39)
$`ฯต_+ฯต_{}`$ $`=`$ $`{\displaystyle \frac{v_+^2v_{}^2}{2}}.`$ (40)
Combining equations (38) and (39) and substituting for the pressure $`P`$ using equation (9) yields the velocity jump condition
$$\frac{v_+}{v_{}}=\gamma ^1_{}^2<1,_{}\frac{v_{}}{a_{}},$$
(41)
where $`_{}`$ is the incident Mach number of the shock. The corresponding result for the shock compression ratio $`R_{}`$ is
$$R_{}\frac{\rho _+}{\rho _{}}=\gamma _{}^2>1.$$
(42)
Hence the gas density increases across the shock as expected. Based on equations (15) and (41) and the fact that $`a_+=a_{}`$ in an isothermal shock, we find that the jump condition for the entropy parameter $`K`$ is given by
$$\frac{K_+}{K_{}}=\gamma ^1_{}^2<1,$$
(43)
which indicates that entropy is lost from the disk at the shock location due to the escape of the particles that form the outflow (jet).
We can also make use of equation (41) to rewrite the jump condition for the energy transport parameter (eq. ) as
$$\mathrm{\Delta }ฯตฯต_+ฯต_{}=\frac{v_{}^2}{2}\left(\frac{1}{\gamma ^2_{}^4}1\right)<0.$$
(44)
The associated rate at which energy escapes from the disk at the isothermal shock location (the โshock luminosityโ) is given by (see eq. )
$$L_{\mathrm{shock}}\mathrm{\Delta }\dot{E}=\dot{M}\mathrm{\Delta }ฯต\mathrm{ergs}\mathrm{s}^1.$$
(45)
Eliminating $`\mathrm{\Delta }ฯต`$ between equations (44) and (45) yields the alternative result
$$\frac{L_{\mathrm{shock}}}{\dot{M}}=\frac{v_{}^2}{2}\left(1\frac{1}{\gamma ^2_{}^4}\right).$$
(46)
### 4.2 Shock Point Analysis
For a given value of $`\gamma `$, it is known that smooth, shock-free global flow solutions exist only within a restricted region of the $`(ฯต_{},\mathrm{})`$ parameter space, and isothermal shocks can occur only in a subset of the smooth-flow region. In order for a shock to exist in the flow, it must be located between two critical points, and it must also satisfy the jump conditions given by equations (41), (43), and (44). The procedure for determining the disk/shock structure is summarized below.
The process begins with the selection of values for the fundamental parameters $`ฯต_{}`$, $`\mathrm{}`$, and $`\gamma `$. The values of $`ฯต_{}`$ and $`\mathrm{}`$ are ultimately constrained by the observations of a specific object, as discussed in ยง 7. Following Narayan, Kato, & Honma (1997), we shall assume an approximate equipartition between the gas and magnetic pressures, and therefore we set $`\gamma =1.5`$. The first step in the determination of the shock location is the computation of the outer critical point location, $`r_{c1}`$, using the quartic equation (27). The associated values for the critical velocity, $`v_{c1}`$, the critical sound speed, $`a_{c1}`$, and the critical entropy, $`K_{c1}`$, are then calculated using equations (23), (24), and (25), respectively. Note that since the flow is adiabatic everywhere in the pre-shock region, it follows that
$$K_{}=K_{c1}.$$
(47)
The profiles of the velocity $`v(r)`$ and the sound speed $`a(r)`$ in the pre-shock region can therefore be calculated using a root-finding procedure based on equation (29), or, alternatively, by integrating numerically the wind equation (19). The next step is the selection of an initial guess for the shock radius, $`r_{}`$, and the calculation of the associated shock Mach number $`_{}v_{}/a_{}`$ using the pre-shock dynamical solutions for $`v(r)`$ and $`a(r)`$. Based on the value of $`_{}`$, we can compute the jump in the entropy parameter $`K`$ using equation (43), and consequently we find that the entropy in the downstream region is given by
$$K_+=\frac{K_{c1}}{\gamma _{}^2},$$
(48)
where we have also used equation (47).
In order to determine whether the initial guess for $`r_{}`$ is self-consistent, we employ a second, independent procedure for calculating the entropy in the downstream region based on the critical nature of the flow. In this approach, the downstream energy parameter $`ฯต_+`$ is calculated using the jump condition given by equation (44), which yields
$$ฯต_+=ฯต_{}+\frac{v_{}^2}{2}\left(\frac{1}{\gamma ^2_{}^4}1\right).$$
(49)
We utilize this value to compute the downstream critical point radius $`\widehat{r}_{c3}`$ based on the quartic equation (27). The associated values for the critical velocity, $`\widehat{v}_{c3}`$, the critical sound speed, $`\widehat{a}_{c3}`$, and the critical entropy, $`\widehat{K}_{c3}`$, are then calculated using equations (23), (24), and (25), respectively. The final step is to compare the value of $`\widehat{K}_{c3}`$ with that obtained for $`K_+`$ using equation (48). If these two quantities are equal, then the shock radius $`r_{}`$ is correct and the disk/shock model is therefore dynamically self-consistent. Otherwise, the value for $`r_{}`$ must be iterated and the search continued. Roots for $`r_{}`$ can be found only in certain regions of the $`(ฯต_{},\mathrm{},\gamma )`$ parameter space, as discussed below.
By combining the analysis of the shock location discussed above with the critical conditions developed in ยง 3, we are able to compute the structure of shocked and shock-free (smooth) disk solutions for a given set of parameters $`ฯต_{}`$, $`\mathrm{}`$, and $`\gamma `$. The resulting topology of the parameter space is depicted in Figure 2 for the case with $`\gamma =1.5`$, which is the main focus of this paper. Within region I, only smooth flows are possible, and in regions II and III both smooth and shocked solutions are available. No global flow solutions (either smooth or shocked) exist in region IV, with $`\mathrm{}>\mathrm{}_{\mathrm{max}}`$. Inside region II, one root for the shock radius $`r_{}`$ can be found, and in region III two shock solutions are available, although only one actual shock can occur in a given flow. It is unclear which of the two roots for $`r_{}`$ in region III is preferred since the stability properties of the shocks are not completely understood (e.g., Chakrabarti 1989a, 1989b; Abramowicz & Chakrabarti 1990). However, it is worth noting that the inner shock is always the stronger of the two because the Mach number diverges as the gas approaches the horizon. The larger compression ratio associated with the inner location leads to more efficient particle acceleration and enhanced entropy generation, and therefore we expect that the inner solution is preferred in nature. Based on this argument, we will focus on the inner shock location in our subsequent analysis.
Before we proceed to examine the transport equation for the relativistic particles, it is important to analyze the asymptotic solutions obtained for the dynamical variables near the event horizon and also at large radii. This is a crucial issue since the nature of the global solutions to the transport equation depends sensitively on the boundary conditions imposed at large and small radii. The asymptotic solutions for the dynamical variables near the event horizon were fully discussed by Becker & Le (2003). We shall briefly review their results and then perform a similar analysis in order to determine the asymptotic properties of the solutions as $`r\mathrm{}`$.
### 4.3 Asymptotic Behavior Near the Horizon
Becker & Le (2003) demonstrated that the variation of the global solutions in a viscous ADAF disk becomes purely adiabatic close to the event horizon, and therefore the asymptotic solutions that they obtain can be directly applied to our inviscid model. Using their equations (47) and (51), we find that the asymptotic variations of the radial velocity $`v`$ and the sound speed $`a`$ near the horizon are given by
$$v^2(r)(rr__\mathrm{S})^1,a^2(r)(rr__\mathrm{S})^{(1\gamma )/(1+\gamma )},rr__\mathrm{S}.$$
(50)
The divergence of $`v`$ as $`rr__\mathrm{S}`$ implies that it cannot represent the standard velocity in the region near the horizon. However, our dynamical model is consistent with relativity if we interpret $`v`$ as the radial component of the four-velocity in this region (Becker & Le 2003; Becker & Subramanian 2005). By combining these relations with equations (4), (8), and (10), we find that the corresponding results for the asymptotic variations of the disk half-thickness $`H`$ and the density $`\rho `$ become
$$H(r)(rr__\mathrm{S})^{(\gamma +3)/(2\gamma +2)},\rho (r)(rr__\mathrm{S})^{1/(\gamma +1)},rr__\mathrm{S}.$$
(51)
### 4.4 Asymptotic Behavior at Infinity
We can use the energy transport equation (13) and the entropy equation (15) to obtain the asymptotic solutions for $`v`$ and $`a`$ at infinity as follows. In the limit $`r\mathrm{}`$, the two dominant terms in equation (13) are $`ฯต`$ and $`a^2/(\gamma 1)`$, and therefore we find that
$$a^2(r)(\gamma 1)ฯต,r\mathrm{}.$$
(52)
Recalling that $`K`$ is constant in the adiabatic upstream flow, we can combine equations (15) and (52) to conclude that the asymptotic variation of the inflow velocity is given by
$$vr^{5/2},r\mathrm{}.$$
(53)
Finally, based on equations (8), (10), and (52), we find that the disk half-height varies as
$$Hr^{3/2},r\mathrm{}.$$
(54)
We can also combine equations (9), (14), and (52) to conclude that the asymptotic behavior of the density is given by
$$\rho \mathrm{const}.,r\mathrm{}.$$
(55)
## 5 STEADY-STATE PARTICLE ACCELERATION
Our goal in this paper is to analyze the transport and acceleration of relativistic particles (ions) in a disk governed by the dynamical model developed in ยง 3 โ 4. For fixed values of the theory parameters $`ฯต_{}`$, $`\mathrm{}`$, and $`\gamma `$, we will study the transport of particles in disks with and without shocks. The particle transport model utilized here includes advection, spatial diffusion, Fermi energization, and particle escape. In order to maintain consistency with the dynamics of the disk, we will need to equate the energy escape rate for the relativistic particles with the โshock luminosityโ $`L_{\mathrm{shock}}`$ given by equation (45). Our treatment of Fermi energization includes both the general compression related to the overall convergence of the accretion flow, as well as the enhanced compression that occurs at the shock. In the scenario under consideration here, the escape of particles from the disk occurs via vertical spatial diffusion in the tangled magnetic field, as depicted in Figure 1. To avoid unnecessary complexity, we will utilize a simplified model in which only the radial ($`r`$) component of the spatial particle transport is treated in detail. In this approach, the diffusion and escape of the particles in the vertical ($`z`$) direction is modeled using an escape-probability formalism. We will treat the relativistic ions as test particles, meaning that their contribution to the pressure in the flow is neglected. This assumption is valid provided the pressure of the relativistic particles turns out to be a small fraction of the thermal pressure, as discussed in ยง 8. The ions accelerated at the shock are energized via collisions with MHD waves advected along with the background (thermal) flow, and therefore the shock width is expected to be comparable to the magnetic coherence length, $`\lambda _{\mathrm{mag}}`$. This approximation will be used to determine the rate at which particles escape from the disk in the vicinity of the shock.
### 5.1 Transport Equation
The Greenโs function, $`f__\mathrm{G}(E_0,E,r_{},r)`$, represents the particle distribution resulting from the continual injection of $`\dot{N}_0`$ particles per second, each with energy $`E_0`$, from a source located at the shock radius, $`r=r_{}`$. In a steady-state situation, the Greenโs function satisfies the transport equation (Becker 1992)
$$\frac{f__\mathrm{G}}{t}=0=\stackrel{}{}\stackrel{}{F}\frac{1}{3E^2}\frac{}{E}\left(E^3\stackrel{}{v}\stackrel{}{}f__\mathrm{G}\right)+\dot{f}_{\mathrm{source}}\dot{f}_{\mathrm{esc}},$$
(56)
where the specific flux $`\stackrel{}{F}`$ is evaluated using
$$\stackrel{}{F}=\kappa \stackrel{}{}f__\mathrm{G}\frac{\stackrel{}{v}E}{3}\frac{f__\mathrm{G}}{E},$$
(57)
and the source and escape terms are given by
$$\dot{f}_{\mathrm{source}}=\frac{\dot{N}_0\delta (EE_0)\delta (rr_{})}{(4\pi E_0)^2r_{}H_{}},\dot{f}_{\mathrm{esc}}=A_0c\delta (rr_{})f__\mathrm{G}.$$
(58)
The quantities $`E`$, $`\kappa `$, $`\stackrel{}{v}`$, and $`H_{}H(r_{})`$ represent the particle energy, the spatial diffusion coefficient, the vector velocity, and the disk half-thickness at the shock location, respectively, and the dimensionless parameter $`A_0`$ determines the rate of particle escape through the surface of the disk at the shock location. The vector velocity $`\stackrel{}{v}`$ has components given by $`\stackrel{}{v}=v_r\widehat{r}+v_z\widehat{z}+v_\varphi \widehat{\varphi }`$, where $`v=v_r`$ is the positive inflow speed.
The total number and energy densities of the relativistic particles, denoted by $`n_r`$ and $`U_r`$, respectively, are related to the Greenโs function via
$$n_r(r)=_0^{\mathrm{}}4\pi E^2f__\mathrm{G}๐E,U_r(r)=_0^{\mathrm{}}4\pi E^3f__\mathrm{G}๐E,$$
(59)
which determine the normalization of $`f__\mathrm{G}`$. Equations (56) and (57) can be combined to obtain the alternative form
$$\stackrel{}{v}\stackrel{}{}f__\mathrm{G}=\frac{\stackrel{}{}\stackrel{}{v}}{3}E\frac{f__\mathrm{G}}{E}+\stackrel{}{}\left(\kappa \stackrel{}{}f__\mathrm{G}\right)+\dot{f}_{\mathrm{source}}\dot{f}_{\mathrm{esc}},$$
(60)
where the left-hand side represents the co-moving (advective) time derivative and the terms on the right-hand side describe first-order Fermi acceleration, spatial diffusion, the particle source, and the escape of particles from the disk at the shock location, respectively. Note that escape and particle injection are localized to the shock radius due to the presence of the $`\delta `$-functions in equations (58). Our focus here is on the first-order Fermi acceleration of relativistic particles at a standing shock in an accretion disk, and therefore equation (60) does not include second-order Fermi processes that may also occur in the flow due to MHD turbulence (e.g., Schlickeiser 1989a,b; Subramanian, Becker, & Kazanas 1999).
Under the assumption of cylindrical symmetry, equations (58) and (60) can be rewritten as
$`v_r{\displaystyle \frac{f__\mathrm{G}}{r}}+v_z{\displaystyle \frac{f__\mathrm{G}}{z}}`$ $``$ $`{\displaystyle \frac{1}{3}}\left[{\displaystyle \frac{1}{r}}{\displaystyle \frac{}{r}}\left(rv_r\right)+{\displaystyle \frac{dv_z}{dz}}\right]E{\displaystyle \frac{f__\mathrm{G}}{E}}{\displaystyle \frac{1}{r}}{\displaystyle \frac{}{r}}\left(r\kappa {\displaystyle \frac{f__\mathrm{G}}{r}}\right)`$ (61)
$`=`$ $`{\displaystyle \frac{\dot{N}_0\delta (EE_0)\delta (rr_{})}{(4\pi E_0)^2r_{}H_{}}}A_0c\delta (rr_{})f__\mathrm{G},`$
where the escape of particles from the disk is described by the final term on the right-hand side. In Appendix A, we demonstrate that the vertically integrated transport equation is given by (see eq. \[A9\])
$`Hv_r{\displaystyle \frac{f__\mathrm{G}}{r}}`$ $`=`$ $`{\displaystyle \frac{1}{3r}}{\displaystyle \frac{}{r}}\left(rHv_r\right)E{\displaystyle \frac{f__\mathrm{G}}{E}}+{\displaystyle \frac{1}{r}}{\displaystyle \frac{}{r}}\left(rH\kappa {\displaystyle \frac{f__\mathrm{G}}{r}}\right)`$ (62)
$`+`$ $`{\displaystyle \frac{\dot{N}_0\delta (EE_0)\delta (rr_{})}{(4\pi E_0)^2r_{}}}A_0cH_{}f__\mathrm{G}\delta (rr_{}),`$
where the symbols $`f__\mathrm{G}`$, $`v_r`$, and $`\kappa `$ represent vertically averaged quantities. We establish in Appendix B that within the context of our one-dimensional spatial model, the dimensionless escape parameter $`A_0`$ appearing in equation (62) is given by (see eqs. \[B8\] and \[B10\])
$$A_0=\left(\frac{3\kappa _{}}{cH_{}}\right)^2<1,$$
(63)
where $`\kappa _{}(\kappa _{}+\kappa _+)/2`$ denotes the mean value of the diffusion coefficient at the shock location. The condition $`A_0<1`$ is required for the validity of the diffusive picture we have employed in our model for the vertical transport.
### 5.2 Number and Energy Densities
The energy moments of the Greenโs function, $`I_n(r)`$, are defined by
$$I_n(r)_0^{\mathrm{}}4\pi E^nf__\mathrm{G}๐E,$$
(64)
so that (cf. eqs. )
$$n_r(r)=I_2(r),U_r(r)=I_3(r).$$
(65)
By operating on equation (62) with $`_0^{\mathrm{}}4\pi E^n๐E`$ and integrating by parts once, we find that the function $`I_n`$ satisfies the differential equation
$`Hv_r{\displaystyle \frac{dI_n}{dr}}`$ $`=`$ $`\left({\displaystyle \frac{n+1}{3}}\right){\displaystyle \frac{I_n}{r}}{\displaystyle \frac{d}{dr}}\left(rHv_r\right)+{\displaystyle \frac{1}{r}}{\displaystyle \frac{d}{dr}}\left(rH\kappa {\displaystyle \frac{dI_n}{dr}}\right)`$ (66)
$`+`$ $`{\displaystyle \frac{\dot{N}_0E_0^{n2}\delta (rr_{})}{4\pi r_{}}}A_0cH_{}I_n\delta (rr_{}),`$
which can be expressed in the flux-conservation form
$$\frac{d}{dr}\left(4\pi rHF_n\right)=4\pi rH\left[\left(\frac{2n}{3}\right)v\frac{dI_n}{dr}+\frac{\dot{N}_0E_0^{n2}\delta (rr_{})}{4\pi r_{}H_{}}A_0c\delta (rr_{})I_n\right],$$
(67)
where $`4\pi rHF_n`$ represents the rate of transport of the $`n`$th moment, and the flux $`F_n`$ is defined by
$$F_n\left(\frac{n+1}{3}\right)vI_n\kappa \frac{dI_n}{dr},$$
(68)
and $`v=v_r`$ denotes the positive inflow speed.
In order to close the system of equations and solve for the relativistic particle number and energy densities $`I_2(r)`$ and $`I_3(r)`$ using equation (66), we must also specify the radial variation of the diffusion coefficient $`\kappa `$. The behavior of $`\kappa `$ can be constrained by considering the fundamental physical principles governing accretion onto a black hole. First, we note that near the event horizon, particles are swept into the black hole at the speed of light, and therefore advection must dominate over diffusion. This condition applies to both the thermal (background) and the nonthermal (relativistic) particles. Second, we note that in the outer region ($`r\mathrm{}`$), diffusion is expected to dominate over advection. Focusing on the flux equation for the particle number density $`n_r`$, obtained by setting $`n=2`$ in equation (68), we can employ scale analysis to conclude based on our physical constraints that
$$\underset{rr__\mathrm{S}}{lim}\frac{\kappa (r)}{(rr__\mathrm{S})v(r)}=0,\underset{r\mathrm{}}{lim}\frac{rv(r)}{\kappa (r)}=0.$$
(69)
The precise functional form for the spatial variation of $`\kappa `$ is not completely understood in the accretion disk environment. In order to obtain a mathematically tractable set of equations with a reasonable physical behavior, we shall utilize the general form
$$\kappa (r)=\kappa _0v(r)r__\mathrm{S}\left(\frac{r}{r__\mathrm{S}}1\right)^\alpha ,$$
(70)
where $`\kappa _0`$ and $`\alpha `$ are dimensionless constants. Due to the appearance of the inflow speed $`v`$ in equation (70), we note that $`\kappa `$ exhibits a jump at the shock. This is expected on physical grounds since the MHD waves that scatter the ions are swept along with the thermal background flow, and therefore they should also experience a density compression at the shock.
As discussed above, close to the event horizon, inward advection at the speed of light must dominate over outward diffusion. Conversely, in the outer region, we expect that diffusion will dominate over advection as $`r\mathrm{}`$. By combining equation (70) with the asymptotic velocity variations expressed by equations (50) and (53), we find that the conditions given by equations (69) are satisfied if $`\alpha >1`$, and in our work we set $`\alpha =2`$. Note that the escape parameter $`A_0`$ is related to $`\kappa _0`$ via equation (63), which can be combined with equation (70) to write
$$A_0=\left(\frac{3\kappa _0v_{}r__\mathrm{S}}{cH_{}}\right)^2\left(\frac{r_{}}{r__\mathrm{S}}1\right)^4<\mathrm{\hspace{0.17em}1},$$
(71)
where $`v_{}(v_{}+v_+)/2`$ represents the mean velocity at the shock location $`r=r_{}`$. The value of the diffusion parameter $`\kappa _0`$ is constrained by the inequality in equation (71). In ยง 7 we demonstrate that $`\kappa _0`$ can be computed for a given source based on energy conservation considerations. With the introduction of equations (70) and (71), we have completely defined all of the quantities in the transport equation, and we can now solve for the number and energy densities of the relativistic particles. The particle distribution Greenโs function $`f__\mathrm{G}`$ and its applications will be discussed in a separate paper.
## 6 SOLUTIONS FOR THE ENERGY MOMENTS
Once the disk/shock dynamics have been computed based on the selected values for the free parameters $`ฯต_{}`$, $`\mathrm{},`$ and $`\gamma `$ using the results in ยงยง 3 and 4, the associated solutions for the number and energy densities of the relativistic particles in the disk can be obtained by solving equation (66). In the case of the number density, $`n_r=I_2`$, an exact solution can be derived based on the linear first-order differential equation describing the conservation of particle flux. However, in order to understand the variation of the energy density, $`U_r=I_3`$, we must numerically integrate a second-order equation. The solutions obtained below are applied in ยง 7 to model the outflows observed in M87 and Sgr A$`^{}`$.
### 6.1 Relativistic Particle Number Density
The equation governing the transport of the particle number density, $`n_r`$, is obtained by setting $`n=2`$ in equation (67), which yields
$$\frac{d\dot{N}_r}{dr}=\dot{N}_0\delta (rr_{})4\pi r_{}H_{}A_0cn_r\delta (rr_{}),$$
(72)
where the relativistic particle transport rate, $`\dot{N}_r(r)`$, is defined by (cf. eq. )
$$\dot{N}_r(r)4\pi rH\left(vn_r+\kappa \frac{dn_r}{dr}\right)\mathrm{s}^1,$$
(73)
and $`\dot{N}_r>0`$ if the transport is in the outward direction. Since the source is located at the shock, there are two spatial domains of interest in our calculation of the particle transport, namely domain I ($`r>r_{}`$), and domain II ($`r<r_{}`$). Note that the number density $`n_r(r)`$ must be continuous at $`r=r_{}`$ in order to avoid generating an infinite diffusive flux according to equation (73). Away from the shock location, $`rr_{}`$, and therefore equation (72) reduces to $`\dot{N}_r=\mathrm{const}.`$, which reflects the fact that particle injection and escape are localized at the shock. We can therefore write
$$\dot{N}_r(r)=\{\begin{array}{cc}\dot{N}_\mathrm{I},\hfill & r>r_{}\text{ ,}\hfill \\ \dot{N}_{\mathrm{II}},\hfill & r<r_{}\text{ ,}\hfill \end{array}$$
(74)
where the constant $`\dot{N}_\mathrm{I}>0`$ denotes the rate at which particles are transported outward, radially, from the source location, and the constant $`\dot{N}_{\mathrm{II}}<0`$ represents the rate at which particles are transported inward towards the event horizon.
The magnitude of the jump in the particle transport rate at the shock is obtained by integrating equation (72) with respect to radius in a very small region around $`r=r_{}`$, which yields
$$\dot{N}_\mathrm{I}\dot{N}_{\mathrm{II}}=\dot{N}_0\dot{N}_{\mathrm{esc}},$$
(75)
where
$$\dot{N}_{\mathrm{esc}}4\pi r_{}H_{}A_0cn_{}$$
(76)
represents the (positive) rate at which particles escape from the disk at the shock location to form the outflow (jet), and $`n_{}n_r(r_{})`$. If no shock is present in the flow, then $`A_0=0`$ and therefore $`\dot{N}_{\mathrm{esc}}`$=0. Note that the discontinuity in $`\dot{N}_r`$ at the shock produces a jump in the derivative $`dn_r/dr`$ via equation (73).
We can rewrite equation (73) for the number density in the form
$$\frac{dn_r}{dr}+\frac{v}{\kappa }n_r=\frac{\dot{N}_r}{4\pi rH\kappa },$$
(77)
which is a linear, first-order differential equation for $`n_r(r)`$. Using the standard integrating factor technique and employing equation (70) for $`\kappa `$ yields the exact solution
$$n_r(r)=e^{J(r)}\left[n_{}\frac{\dot{N}_r(r)}{4\pi }_r_{}^r\frac{e^{J(r^{})}}{r^{}H\kappa }๐r^{}\right],$$
(78)
where $`\dot{N}_r(r)`$ is given by equation (74) and the function $`J(r)`$ is defined by
$$J(r)_r_{}^rv\kappa ^1๐r^{}=\kappa _0^1\left[\left(\frac{r_{}}{r__\mathrm{S}}1\right)^1\left(\frac{r}{r__\mathrm{S}}1\right)^1\right].$$
(79)
According to equation (78), $`n_r(r)`$ is continuous at the shock/source location as required. Far from the black hole, diffusion dominates the particle transport, and therefore $`n_r`$ should vanish as $`r\mathrm{}`$. In order to ensure this behavior, we must have
$$n_{}=\dot{N}_\mathrm{I}C_\mathrm{I},$$
(80)
where
$$C_\mathrm{I}\frac{1}{4\pi }_r_{}^{\mathrm{}}\frac{e^{J(r^{})}}{r^{}H\kappa }๐r^{}.$$
(81)
Furthermore, in order to avoid exponential divergence of $`n_r`$ as $`rr__\mathrm{S}`$ in domain II, we also require that
$$n_{}=\dot{N}_{\mathrm{II}}C_{\mathrm{II}},$$
(82)
where
$$C_{\mathrm{II}}\frac{1}{4\pi }_{r__\mathrm{S}}^r_{}\frac{e^{J(r^{})}}{r^{}H\kappa }๐r^{}.$$
(83)
By combining equations (75), (76), (80), and (82), we can develop explicit expressions for the quantities $`n_{}`$, $`\dot{N}_\mathrm{I}`$, $`\dot{N}_{\mathrm{II}}`$, and $`\dot{N}_{\mathrm{esc}}`$ based on the values of $`r_{}`$ and $`\dot{N}_0`$ and the profiles of the inflow velocity $`v(r)`$ and the diffusion coefficient $`\kappa (r)`$. The results obtained are
$`n_{}`$ $`=`$ $`{\displaystyle \frac{\dot{N}_0}{C_\mathrm{I}^1+C_{\mathrm{II}}^1+4\pi r_{}H_{}A_0c}},`$
$`\dot{N}_\mathrm{I}`$ $`=`$ $`{\displaystyle \frac{n_{}}{C_\mathrm{I}}},`$ (84)
$`\dot{N}_{\mathrm{II}}`$ $`=`$ $`{\displaystyle \frac{n_{}}{C_{\mathrm{II}}}},`$
$`\dot{N}_{\mathrm{esc}}`$ $`=`$ $`4\pi r_{}H_{}A_0cn_{}.`$
These relations, along with equation (78), complete the formal solution for the relativistic particle number density $`n_r(r)`$. The solution is valid in both shocked and shock-free disks (the shock-free case is treated by setting $`A_0=0`$). When a shock is present, the particle escape rate $`\dot{N}_{\mathrm{esc}}`$ is proportional to $`\dot{N}_0`$ but is independent of $`E_0`$ by virtue of equations (84).
It is interesting to examine the asymptotic variation of $`n_r`$ near the event horizon and also at large distances from the black hole. Far from the hole, advection is negligible and the particle transport in the disk is dominated by outward-bound diffusion. In this case we can use equation (77) to conclude that
$$\frac{dn_r}{dr}\frac{\dot{N}_\mathrm{I}}{4\pi rH\kappa },r\mathrm{},$$
(85)
where we have used the fact that $`\dot{N}_r=\dot{N}_\mathrm{I}`$ for $`r>r_{}`$. By combining equations (70) and (85) with the asymptotic relations given by equations (53) and (54), we find upon integration that
$$n_r(r)\frac{1}{r},r\mathrm{}.$$
(86)
In order to study the behavior of $`n_r`$ near the event horizon, we take the limit as $`rr__\mathrm{S}`$ in equation (78), obtaining after some algebra
$$n_r(r)\frac{\dot{N}_{\mathrm{II}}}{4\pi rHv},rr__\mathrm{S},$$
(87)
where we have set $`\dot{N}_r=\dot{N}_{\mathrm{II}}`$. Comparing this relation with equation (4), we find that
$$n_r(r)\rho (r),rr__\mathrm{S},$$
(88)
where $`\rho `$ is the density of the background (thermal) gas. Equation (88) is a natural consequence of the fact that the particle transport near the horizon is dominated by inward-bound advection. We can also combine equations (51) and (88) to obtain the explicit asymptotic form
$$n_r(r)(rr__\mathrm{S})^{1/(\gamma +1)},rr__\mathrm{S}.$$
(89)
### 6.2 Relativistic Particle Energy Density
The differential equation satisfied by the relativistic particle energy density, $`U_r=I_3`$, is obtained by setting $`n=3`$ in equation (66), which yields
$`Hv_r{\displaystyle \frac{dU_r}{dr}}`$ $`=`$ $`{\displaystyle \frac{4}{3}}{\displaystyle \frac{U_r}{r}}{\displaystyle \frac{d}{dr}}\left(rHv_r\right)+{\displaystyle \frac{1}{r}}{\displaystyle \frac{d}{dr}}\left(rH\kappa {\displaystyle \frac{dU_r}{dr}}\right)`$ (90)
$`+`$ $`{\displaystyle \frac{\dot{N}_0E_0\delta (rr_{})}{4\pi r_{}}}A_0cH_{}U_r\delta (rr_{}).`$
By analogy with equations (72) and (73), we can recast this expression in the flux-conservation form
$$\frac{d\dot{E}_r}{dr}=4\pi rH\left[\frac{v}{3}\frac{dU_r}{dr}+\frac{\dot{N}_0E_0\delta (rr_{})}{4\pi r_{}H_{}}A_0cU_r\delta (rr_{})\right],$$
(91)
where the relativistic particle energy transport rate, $`\dot{E}_r(r)`$, is defined by
$$\dot{E}_r(r)4\pi rH\left(\frac{4}{3}vU_r+\kappa \frac{dU_r}{dr}\right)\mathrm{ergs}\mathrm{s}^1,$$
(92)
and $`\dot{E}_r>0`$ for outwardly directed transport. Note that unlike the number transport rate $`\dot{N}_r`$, the energy transport rate $`\dot{E}_r`$ does not remain constant within domains I and II due to the appearance of the first term on the right-hand side of equation (91), which expresses the compressional work done on the relativistic particles by the background flow.
Although the energy density $`U_r`$ must be continuous at the shock/source location in order to avoid generating an infinite diffusive flux, the derivative $`dU_r/dr`$ displays a discontinuity at $`r=r_{}`$, which is related to the jump in the energy transport rate via equation (92). By integrating equation (91) in a very small region around $`r=r_{}`$, we find that
$$\mathrm{\Delta }\dot{E}_r=L_{\mathrm{esc}}\dot{N}_0E_0,$$
(93)
where
$$L_{\mathrm{esc}}4\pi r_{}H_{}A_0cU_{}\mathrm{ergs}\mathrm{s}^1$$
(94)
denotes the rate of escape of energy from the disk into the outflow (jet) at the shock location, and $`U_{}U_r(r_{})`$. If no shock is present, then $`A_0=0`$ and therefore $`L_{\mathrm{esc}}=0`$. We remind the reader that the symbol โ$`\mathrm{\Delta }`$โ refers to the difference between post- and pre-shock quantities (see eq. ). Equations (92) and (93) can be combined to show that the derivative jump is given by
$$\mathrm{\Delta }\left(\kappa \frac{dU_r}{dr}\right)=\frac{\dot{N}_0E_0L_{\mathrm{esc}}}{4\pi r_{}H_{}}\frac{4}{3}U_{}\mathrm{\Delta }v.$$
(95)
The differential equation (90) governing the relativistic particle energy density is second-order in radius, and therefore we will need to establish two boundary conditions in order to solve for $`U_r(r)`$. These can be obtained by analyzing the behavior of $`U_r`$ close to the event horizon and at large distances from the black hole. Far from the hole, advection is negligible and the particle transport in the disk is dominated by outward-bound diffusion. In this regime, Fermi acceleration is negligible, and consequently we find that $`U_rn_r`$. We can therefore use equation (86) to conclude that
$$U_r(r)\frac{1}{r},r\mathrm{}.$$
(96)
Close to the event horizon, the particle transport is dominated by advection, and therefore $`U_r`$ and $`n_r`$ obey the standard adiabatic relation
$$U_rn_r^{4/3},rr__\mathrm{S}.$$
(97)
Combining this result with equations (88) then yields
$$U_r(rr__\mathrm{S})^{4/(3\gamma +3)},rr__\mathrm{S}.$$
(98)
The global solution for $`U_r(r)`$ can now be expressed as
$$U_r(r)=\{\begin{array}{cc}AQ_\mathrm{I}(r),\hfill & r>r_{}\text{ ,}\hfill \\ BQ_{\mathrm{II}}(r),\hfill & r<r_{}\text{ ,}\hfill \end{array}$$
(99)
where $`A`$ and $`B`$ are constants and the functions $`Q_\mathrm{I}(r)`$ and $`Q_{\mathrm{II}}(r)`$ satisfy the homogeneous differential equation (see eq. )
$$Hv_r\frac{dQ}{dr}=\frac{4}{3}\frac{Q}{r}\frac{d}{dr}\left(rHv_r\right)+\frac{1}{r}\frac{d}{dr}\left(rH\kappa \frac{dQ}{dr}\right).$$
(100)
along with the boundary conditions (see eqs. and )
$$Q_\mathrm{I}(r_{\mathrm{out}})=\left(\frac{r_{\mathrm{out}}}{r__\mathrm{S}}\right)^1,Q_{\mathrm{II}}(r_{\mathrm{in}})=\left(\frac{r_{\mathrm{in}}}{r__\mathrm{S}}1\right)^{4/(3\gamma +3)},$$
(101)
where $`r_{\mathrm{in}}`$ and $`r_{\mathrm{out}}`$ denote the radii at which the inner and outer boundary conditions are applied, respectively. The constants $`A`$ and $`B`$ are computed by requiring that $`U_r`$ be continuous at $`r=r_{}`$ and that the derivative $`dU_r/dr`$ satisfy the jump condition given by equation (95). The results obtained are
$`A`$ $`=`$ $`B{\displaystyle \frac{Q_{\mathrm{II}}}{Q_\mathrm{I}}}|_{r=r_{}},`$ (102)
$`B`$ $`=`$ $`{\displaystyle \frac{\dot{N}_0E_0}{4\pi r_{}H_{}}}\left[{\displaystyle \frac{4}{3}}(v_{}v_+)Q_{\mathrm{II}}+Q_{\mathrm{II}}^{}\kappa _+{\displaystyle \frac{Q_{\mathrm{II}}Q_\mathrm{I}^{}\kappa _{}}{Q_\mathrm{I}}}+A_0cQ_{\mathrm{II}}\right]^1|_{r=r_{}},`$ (103)
where primes denote differentiation with respect to radius. The solutions for the functions $`Q_\mathrm{I}(r)`$ and $`Q_{\mathrm{II}}(r)`$ are obtained by integrating numerically equation (100) subject to the boundary conditions given by equations (101). Once the constants $`A`$ and $`B`$ are computed using equations (102) and (103), the global solution for $`U_r(r)`$ is evaluated using equation (99). The solution applies whether or not a shock is present in the flow. The shock-free case is treated by setting $`A_0=0`$. This completes the solution procedure for the relativistic particle energy density. The results derived in this section are used in ยง 7 to model the outflows observed in M87 and Sgr A$`^{}`$.
## 7 ASTROPHYSICAL APPLICATIONS
Our goal is to determine the properties of the integrated disk/shock/outflow model based on the observed values for the black hole mass $`M`$, the mass accretion rate $`\dot{M}`$, and the jet kinetic power $`L_{\mathrm{jet}}`$ associated with a given source. The fundamental free parameters for the theoretical model are $`ฯต_{}`$, $`\mathrm{}`$, and $`\gamma `$. Since we set $`\gamma =1.5`$ in order to represent an approximate equipartition between the gas and magnetic pressures (e.g., Narayan, Kato, & Honma 1997), only $`ฯต_{}`$ and $`\mathrm{}`$ remain to be determined. Here we describe how global energy conservation considerations can be used to solve for the various theoretical parameters in the model based on observations.
### 7.1 Energy Conservation Conditions
Once the values of $`M`$, $`\dot{M}`$, and $`L_{\mathrm{jet}}`$ have been specified for a source based on observations, we select a value for the free parameter $`ฯต_{}`$ and then compute $`\mathrm{}`$ by satisfying the relation
$$L_{\mathrm{shock}}=L_{\mathrm{jet}},$$
(104)
where $`L_{\mathrm{shock}}`$ is the shock luminosity given by equation (46). This result ensures that the jump in the energy transport rate at the isothermal shock location is equal to the observed jet kinetic luminosity. The procedure for determining $`\mathrm{}`$ also includes solving for the shock location and the critical structure using results from ยงยง 3 and 4. The velocity profile $`v(r)`$ is computed either by numerically integrating equation (19) or by using a root-finding procedure based on equation (29), and the associated solution for the adiabatic sound speed $`a(r)`$ is obtained using equation (13).
After the velocity profile has been determined, we can compute the number and energy density distributions for the relativistic particles in the disk using equations (78) and (99), respectively. This requires the specification of the injection energy of the seed particles $`E_0`$ as well as their injection rate $`\dot{N}_0`$. We set the injection energy using $`E_0=0.002`$ergs, which corresponds to an injected Lorentz factor $`\mathrm{\Gamma }_0E_0/(m_pc^2)1.3`$, where $`m_p`$ is the proton mass. Particles injected with energy $`E_0`$ are subsequently accelerated to much higher energies due to repeated shock crossings. We find that the speed of the injected particles, $`v_0=c(1\mathrm{\Gamma }_0^2)^{1/2}`$, is about three to four times higher than the mean ion thermal velocity at the shock location, $`v_{\mathrm{rms}}=(3kT_{}/m_p)^{1/2}`$, where $`T_{}`$ is the ion temperature at the shock. The seed particles are therefore picked up from the high-energy tail of the Maxwellian distribution for the thermal ions. With $`E_0`$ specified, we can compute the particle injection rate $`\dot{N}_0`$ using the energy conservation condition
$$\dot{N}_0E_0=L_{\mathrm{shock}},$$
(105)
which ensures that the rate at which energy is injected into the flow in the form of the relativistic seed particles is equal to the energy loss rate for the background gas at the isothermal shock location.
In order to maintain agreement between the transport model and the observations, we must also require that the rate at which particle energy escapes from the disk due to vertical diffusion is equal to the observed jet power. This condition can be written as
$$L_{\mathrm{esc}}=L_{\mathrm{jet}},$$
(106)
where $`L_{\mathrm{esc}}`$ is the energy escape rate given by equation (94). The escape constant, $`A_0`$, appearing in the transport equation is independent of the particle energy in our model, and consequently the escaping particles will have exactly the same mean energy as those in the disk at the shock location. The mean energy of the escaping particles is therefore given by
$$E_{\mathrm{esc}}\frac{U_{}}{n_{}},$$
(107)
where $`n_{}`$ and $`U_{}`$ denote the number and energy densities of the relativistic particles at the shock location, respectively. Hence $`E_{\mathrm{esc}}`$ is proportional to $`E_0`$ but it is independent of $`\dot{N}_0`$. We note that equations (76), (94), and (107) can be combined to show that
$$L_{\mathrm{esc}}=\dot{N}_{\mathrm{esc}}E_{\mathrm{esc}},$$
(108)
where $`\dot{N}_{\mathrm{esc}}`$ is the particle escape rate (eq. ). By satisfying equations (104), (105), and (106), we ensure that energy is properly conserved in our model. Taken together, these relations allow us to solve for the various theoretical parameters based on observational values for $`M`$, $`\dot{M}`$, and $`L_{\mathrm{jet}}`$, as explained below.
### 7.2 Model Parameters
Our simulations of the disk structure and particle transport in M87 and Sgr A$`^{}`$are based on various published observational estimates for $`M`$, $`\dot{M}`$, and $`L_{\mathrm{jet}}`$. For M87, we set $`M=3\times 10^9\mathrm{M}_{}`$ (e.g., Ford et al. 1994), $`\dot{M}=1.3\times 10^1\mathrm{M}_{}\mathrm{yr}^1`$ (e.g., Reynolds et al. 1996), and $`L_{\mathrm{jet}}=5.5\times 10^{43}\mathrm{ergs}\mathrm{s}^1`$ (e.g., Reynolds et al. 1996; Bicknell & Begelman 1996; Owen, Eilek, & Kassim 2000). For Sgr A$`^{}`$, we use the values $`M=2.6\times 10^6\mathrm{M}_{}`$ (e.g., Schรถdel et al. 2002) and $`\dot{M}=8.8\times 10^7\mathrm{M}_{}\mathrm{yr}^1`$ (e.g., Yuan, Markoff, & Falcke 2002; Quataert 2003). Although the kinetic luminosity of the jet in Sgr A$`^{}`$is rather uncertain (see, e.g., Yuan 2000; Yuan et al. 2002), we will adopt the value quoted by Falcke & Biermann (1999), and therefore we set $`L_{\mathrm{jet}}=5\times 10^{38}\mathrm{ergs}\mathrm{s}^1`$.
We study both shocked and shock-free solutions spanning the computational domain between the inner radius $`r_{\mathrm{in}}=2.001`$ and the outer radius $`r_{\mathrm{out}}=5000`$, where $`r_{\mathrm{in}}`$ and $`r_{\mathrm{out}}`$ are the radii at which the boundary conditions are applied (see eqs. ). Six different accretion/shock scenarios are explored in detail, with the values for the various parameters $`ฯต_{}`$, $`\mathrm{}`$, $`ฯต_+`$, $`r_{c1}`$, $`r_{c3}`$, $`r_{}`$, $`\widehat{r}_{c3}`$, $`H_{}`$, $`_{}`$, $`R_{}`$, and $`T_{}`$ reported in Table 1. Models 1, 2, and 3 are associated with M87, while models 4, 5, and 6 are used to study Sgr A$`^{}`$. In our numerical examples, we utilize natural gravitational units ($`GM=c=1`$ and $`r__\mathrm{S}=2`$), except as noted. Based on the observational values for $`\dot{M}`$ and $`L_{\mathrm{jet}}`$ associated with the two sources, we can use equations (45) and (104) to conclude that $`\mathrm{\Delta }ฯต=0.007`$ for M87 and $`\mathrm{\Delta }ฯต=0.01`$ for Sgr A$`^{}`$. These results are consistent with the values for $`ฯต_{}`$ and $`ฯต_+`$ reported in Table 1, and therefore $`L_{\mathrm{shock}}=L_{\mathrm{jet}}`$ as required (see eq. ).
Next we use the energy conservation condition $`L_{\mathrm{esc}}=L_{\mathrm{jet}}`$ (eq. ) to determine the value of the diffusion constant $`\kappa _0`$ (eq. ) for a shocked disk. In Figure 3, we plot $`L_{\mathrm{esc}}/L_{\mathrm{jet}}`$, $`\mathrm{\Gamma }_{\mathrm{esc}}`$, and $`\dot{N}_{\mathrm{esc}}/\dot{N}_0`$ as functions of $`\kappa _0`$ for the M87 and Sgr A$`^{}`$parameters, where $`\mathrm{\Gamma }_{\mathrm{esc}}E_{\mathrm{esc}}/(m_pc^2)`$ is the mean Lorentz factor of the escaping particles. The treatment of energy conservation in our disk/shock model is self-consistent when the condition $`L_{\mathrm{esc}}/L_{\mathrm{jet}}=1`$ is satisfied, which corresponds to specific values of $`\kappa _0`$ as indicated in Figures 3a and 3d. We find that two $`\kappa _0`$ roots exist for models 3 and 6, one root is possible for models 2 and 5, and no roots exist for models 1 and 4. Hence the values of $`ฯต_{}`$ associated with models 2 and 5 represent the maximum possible values for $`ฯต_{}`$ that yield self-consistent solutions based on the M87 and Sgr A$`^{}`$data, respectively. For illustrative purposes, we shall focus on the details of the disk structure and particle transport obtained in models 2 and 5.
### 7.3 Disk Structure and Particle Transport
In order to illustrate the importance of the shock for the acceleration of high-energy particles, we shall examine the structure of the accretion disk either with and without a shock based on the values of the upstream parameters $`ฯต_{}`$ and $`\mathrm{}`$ utilized in models 2 and 5 (see Table 1). In Figures 4a and 4b we plot the inflow speed $`v(r)`$ and the adiabatic sound speed $`a(r)`$ for the shocked and smooth (shock-free) solutions associated with models 2 and 5, respectively. Since we are working within the isothermal shock scenario, the sound speed $`a`$ is continuous at the shock location. In Figures 4c and 4d we plot the specific internal energy $`U/\rho =(\gamma 1)^1kT/m_p`$ for the background (thermal) gas along with the specific gravitational potential (binding) energy $`GM/(rr__\mathrm{S})`$ as functions of radius for $`\gamma =1.5`$. These results demonstrate that the gas is marginally bound in the absence of a shock, and strongly bound when a shock is present. The increased binding of the thermal gas in the disk results from the escape of energy in the outflow, which reduces the sound speed compared with the shock-free case. The enhanced cooling allows the accretion to proceed, thereby removing one of the major objections to the original ADAF scenario (Narayan & Yi 1994, 1995). We emphasize that these new results represent the first fully self-consistent calculations of the structure of an ADAF disk coupled with a shock-driven outflow, hence extending the heuristic work of Blandford & Begelman (1999) and Becker, Subramanian, & Kazanas (2001).
Next we study the solutions obtained for the relativistic particle number and energy density distributions in the disk based on the flow structures associated with models 2 and 5. The related transport parameters are listed in Table 2. In Figures 5 and 6 we plot the global number and energy density distributions obtained in a shocked disk using the model 2 and 5 parameters, respectively. We also include the corresponding results obtained in a shock-free (smooth) disk for the same values of the upstream energy transport rate $`ฯต_{}`$ and the specific angular momentum $`\mathrm{}`$. In each case the densities decrease monotonically with increasing radius. The increase near the horizon is a consequence of advection, while the decline as $`r\mathrm{}`$ reflects the fact that the particles injected at the shock have a very small chance of diffusing to large distances from the black hole. Note that the shocked disk has a lower value for the number density $`n_r`$ at all radii as a consequence of particle escape. However, the shocked disk also displays a higher value for the energy density $`U_r`$, which reflects the central role of shock in accelerating the relativistic test particles.
The kinks in the energy and number density distributions at the shock radius $`r=r_{}`$ indicated in Figures 5 and 6 reflect the derivative jump conditions given by equations (75) and (95). The values for the ratios $`\dot{N}_\mathrm{I}/\dot{N}_{\mathrm{II}}`$ and $`\dot{N}_{\mathrm{esc}}/\dot{N}_0`$ reported in Table 2 indicate that most of the injected particles are advected into the black hole, with $`20`$% escaping to form the outflow (see Figs. 3c and 3f). In order to validate the accuracy of the numerical solutions for $`n_r(r)`$ and $`U_r(r)`$, we also compare the profiles obtained with the asymptotic relations developed in ยง 6. We demonstrate in Figure 7 (model 2) and Figure 8 (model 5) that the solutions for both $`n_r(r)`$ and $`U_r(r)`$ agree closely with the asymptotic expressions given by equations (89) and (98) for small radii and by equations (86) and (96) for large radii. Note that the values reported by Le & Becker (2004) for $`n_{}`$ and $`U_{}`$ were expressed in incorrect units and are given correctly in our Table 2.
### 7.4 Jet Formation in M87 and Sgr A$`^{}`$
The mean energy of the relativistic particles in the disk is given by (cf. eq. )
$$E\frac{U_r(r)}{n_r(r)},$$
(109)
so that $`E=E_{\mathrm{esc}}`$ at $`r=r_{}`$. In Figure 9 we plot the mean energy as a function of radius in shocked and shock-free disks based on the parameters used for models 2 and 5. The results demonstrate that when a shock is present in the flow, the relativistic particle energy is boosted by a factor of $`56`$ at the shock location. By contrast, we find that in the shock-free models with the same values for $`ฯต_{}`$, $`\mathrm{}`$, and $`\kappa _0`$, the energy is boosted by a factor of only $`1.41.5`$. This clearly demonstrates the essential role of the shock in efficiently accelerating particles up to very high energies, far above the energy required to escape from the disk. Note that close to the event horizon, the mean energy of the relativistic particles is further enhanced by the strong compression of the accretion flow, as indicated by the sharp increase in $`E`$ as $`rr__\mathrm{S}`$.
The material in the outflow is initially ejected from the disk in the vicinity of the shock as a hot plasma which cools as it expands, with its outward acceleration powered by the pressure gradient in the surrounding plasma. Based on our results for models 2 and 5, we find that the shock/jet locations are given by $`r_{}22`$ and $`r_{}16`$ for M87 and Sgr A$`^{}`$, respectively. The terminal (asymptotic) Lorentz factor of the jet, $`\mathrm{\Gamma }_{\mathrm{}}`$, can be estimated by writing
$$\mathrm{\Gamma }_{\mathrm{}}=\mathrm{\Gamma }_{\mathrm{esc}}=\frac{E_{\mathrm{esc}}}{m_pc^2},$$
(110)
which is based on the assumption that the jet starts off โslowโ and โhotโ and subsequently expands to become โfastโ and โcold.โ Adopting the $`\mathrm{\Gamma }_{\mathrm{esc}}`$ values listed in Table 2 for M87 and Sgr A$`^{}`$, we obtain $`\mathrm{\Gamma }_{\mathrm{}}=7.92`$ (see Fig. 3b) and $`\mathrm{\Gamma }_{\mathrm{}}=7.26`$ (see Fig. 3e), respectively.
We can now compare our model predictions for the shock/jet location and the asymptotic Lorentz factor with the observations of M87 and Sgr A$`^{}`$. According to Biretta, Junor, & Livio (2002), the M87 jet forms in a region no larger than $`30`$ gravitational radii from the black hole, which agrees rather well with our predicted shock/jet location $`r_{}22`$ for this source. Turning now to the asymptotic (terminal) Lorentz factor, we note that Biretta, Sparks, & Macchetto (1999) estimated $`\mathrm{\Gamma }_{\mathrm{}}6`$ for the M87 jet, which is comparable to the result $`\mathrm{\Gamma }_{\mathrm{}}=7.92`$ obtained using our model. In the case of Sgr A$`^{}`$, our model indicates that the shock forms at $`r_{}16`$ which is fairly close to the value suggested by Yuan (2000). However, future observational work will be needed to test our prediction for the asymptotic Lorentz factor of Sgr A$`^{}`$, since no reliable observational estimate for that quantity is currently available.
### 7.5 Radiative Losses from the Jet
It is still unclear whether the outflows observed to emanate from many radio-loud systems containing black holes are composed of an electron-proton plasma or electron-positron pairs, or a mixture of both. Whichever is the case, the particles must maintain sufficient energy during their journey from the nucleus in order to power the observed radio emission, unless some form of reacceleration takes place along the way, due to shocks propagating along the jet (e.g., Atoyan & Dermer 2004a). Proton-electron outflows, such as those studied here, have a distinct advantage in this regard since most of the kinetic power is carried by the ions, which do not radiate much and are not strongly coupled to the electrons under the typical conditions in a jet (e.g., Felten 1968; Felten, Arp, & Lynds 1970; Anyakoha et al. 1987; Aharonian 2002). We therefore suggest that if the observed outflows are proton-driven, then they may be powered directly by the shock acceleration mechanism operating in the disk, with no requirement for additional in situ reacceleration in the jet. In this section we confirm this conjecture by considering the energy losses experienced by the protons in the outflow. The ions in the jet lose energy via two distinct channels, namely (1) direct radiative losses due to the production of synchrotron and inverse-Compton emission, and (2) indirect radiative losses via Coulomb coupling with the electrons. We will evaluate these two possibilities by estimating the corresponding cooling timescales for the outflows in M87 and Sgr A$`^{}`$and comparing the results with the jet propagation timescales for these sources.
The energy loss rate due to the production of synchrotron and inverse-Compton emission by the relativistic protons escaping from the disk with mean energy $`\mathrm{\Gamma }_{\mathrm{esc}}m_pc^2`$ is given by (see eqs. \[7.17\] and \[7.18\] from Rybicki & Lightman 1979)
$$\left(\frac{dE}{dt}\right)|_{\mathrm{rad}}=\frac{4\sigma __\mathrm{T}c\mathrm{\Gamma }_{\mathrm{esc}}^2}{3}\left(\frac{m_e}{m_p}\right)^2\left(U_B+U_{\mathrm{ph}}\right),$$
(111)
where $`U_{\mathrm{ph}}`$ and $`U_B=B^2/(8\pi )`$ denote the energy densities of the soft radiation and the magnetic field with strength $`B`$, respectively. The associated energy loss timescale is therefore
$$t_{\mathrm{rad}}\frac{\mathrm{\Gamma }_{\mathrm{esc}}m_pc^2}{(dE/dt)|_{\mathrm{rad}}}=\frac{3m_pc}{4\sigma __\mathrm{T}\mathrm{\Gamma }_{\mathrm{esc}}}\left(\frac{m_p}{m_e}\right)^2\left(U_B+U_{\mathrm{ph}}\right)^1.$$
(112)
In our application to M87, we take $`B0.1`$G based on estimates from Biretta, Stern, & Harris (1991), and we set $`\mathrm{\Gamma }_{\mathrm{esc}}8`$ (see Table 2, model 2). Assuming equipartition between the magnetic field and the soft radiation, this yields for the radiative cooling time $`t_{\mathrm{rad}}10^{12}`$yrs, which suggests that the protons can easily maintain their energy for many millions of parsecs without being seriously effected by synchrotron or inverse-Compton losses, as expected. For Sgr A$`^{}`$, we assume equipartition with $`B10`$G (Atoyan & Dermer 2004b) and $`\mathrm{\Gamma }_{\mathrm{esc}}7`$ (see Table 2, model 5). The radiative cooling time for the escaping protons is therefore $`t_{\mathrm{rad}}10^8`$yrs. Hence synchrotron and inverse-Compton losses have virtually no effect on the energy of the protons in the Sgr A$`^{}`$jet.
In addition to synchrotron and inverse-Compton radiation, the protons in the jet will also lose energy due to Coulomb coupling with the thermal electrons, which radiate much more efficiently than the protons. The energy loss rate for this process can be estimated using equation (4.16) from Mannheim & Schlickeiser (1994), which yields
$$\left(\frac{dE}{dt}\right)|_{\mathrm{coul}}=\mathrm{\hspace{0.17em}30}n_e\sigma __\mathrm{T}m_ec^3,$$
(113)
where $`n_e`$ represents the electron number density in the jet. The associated loss timescale for a proton escaping from the disk with mean energy $`\mathrm{\Gamma }_{\mathrm{esc}}m_pc^2`$ is
$$t_{\mathrm{coul}}\frac{\mathrm{\Gamma }_{\mathrm{esc}}m_pc^2}{(dE/dt)|_{\mathrm{coul}}}=\frac{\mathrm{\Gamma }_{\mathrm{esc}}m_p}{30n_e\sigma __\mathrm{T}cm_e}.$$
(114)
The electron number density $`n_e`$ decreases rapidly as the jet expands from the disk into the external medium. Hence the most conservative estimate (based on the strongest Coulomb coupling) is obtained by adopting conditions at the base of the jet, where $`n_e`$ has its maximum value. To estimate the electron number density at the base of the outflow, we begin by calculating the rate at which protons escape from the disk at the shock location. By using equation (B8) to eliminate $`A_0`$ in equation (76), we find that the proton escape rate is given by
$$\dot{N}_{\mathrm{esc}}=\frac{4\pi r_{}\lambda _{\mathrm{mag}}^2cn_{}}{H_{}},$$
(115)
where $`r_{}`$, $`n_{}`$, $`H_{}`$, and $`\lambda _{\mathrm{mag}}`$ denote the radius, the proton number density, the vertical half-thickness, and the magnetic mean free path inside the disk at the shock location, respectively. The shock is expected to have a width comparable to $`\lambda _{\mathrm{mag}}`$, and therefore the sum of the upper and lower face areas of the shock annulus is equal to $`4\pi r_{}\lambda _{\mathrm{mag}}`$. We also note that the flux of the relativistic protons escaping from the disk into the outflow is given by $`cn_p`$, where $`n_p`$ is the proton number density at the base of the jet. Combining these relations, we can write the proton escape rate in terms of $`n_p`$ using
$$\dot{N}_{\mathrm{esc}}=4\pi r_{}\lambda _{\mathrm{mag}}cn_p.$$
(116)
By equating the two expressions for $`\dot{N}_{\mathrm{esc}}`$ given by equations (115) and (116), we find that $`n_p`$ is related to $`n_{}`$ via
$$\frac{n_p}{n_{}}=\frac{\lambda _{\mathrm{mag}}}{H_{}}<1.$$
(117)
Since the electron-proton jet must be charge neutral, the electron number density at the base of the jet, $`n_e`$, is equal to the proton number density $`n_p`$, and therefore we obtain
$$n_e=\frac{\lambda _{\mathrm{mag}}}{H_{}}n_{}.$$
(118)
Using the relation $`\lambda _{\mathrm{mag}}/H_{}=A_0^{1/2}`$ (see eq. \[B8\]) along with the results for $`A_0`$ and $`n_{}`$ reported in Table 1 for M87 (model 2), we obtain $`n_e=0.11n_{}=2.2\times 10^3\mathrm{cm}^3`$. Setting $`\mathrm{\Gamma }_{\mathrm{esc}}8`$, we find that equation (114) yields for the electron-proton Coulomb coupling timescale $`t_{\mathrm{coul}}3.5\times 10^5`$yrs. Note that this is an extremely conservative estimate since it is based on conditions at the bottom of the jet, and therefore it suggests that Coulomb coupling between the protons and the electrons is insufficient to seriously degrade the energy of the accelerated ions escaping from the disk as they propagate out to the radio lobes via the jet. For Sgr A$`^{}`$, we use the model 5 data in Table 2 to obtain $`n_e=0.13n_{}=5.6\times 10^4\mathrm{cm}^3`$. Setting $`\mathrm{\Gamma }_{\mathrm{esc}}7`$ yields for the Coulomb coupling timescale $`t_{\mathrm{coul}}1.2\times 10^4`$yrs, which implies that the length of the jet can be as large several thousand parsecs before much energy is drained from the protons, assuming the material in the jet travels at half the speed of light. We emphasize that these numerical estimates of the importance of radiative and Coulomb losses experienced by the relativistic protons are based on the โworst-caseโ assumption that the conditions at the base of the outflow prevail throughout the jet. In reality, the jet density will drop rapidly as the gas expands, and therefore the true values for the proton energy loss timescales will be much larger than the results obtained here. This strongly suggests that shock acceleration of the protons in the disk, as investigated here, is sufficient to power the observed outflows without requiring any reacceleration in the jets.
### 7.6 Radiative Losses from the Disk
In the ADAF scenario that we have focused on, radiative losses from the disk are ignored. The self-consistency of this approximation can be evaluated by computing the free-free emissivity due to the thermal gas in the disk. The total X-ray luminosity can be estimated by integrating equation (5.15b) from Rybicki & Lightman (1979) over the disk volume to obtain for pure, fully-ionized hydrogen
$$L_{\mathrm{rad}}=_{r__\mathrm{S}}^{\mathrm{}}1.4\times 10^{27}T_e^{1/2}\rho ^2m_p^2๐V,$$
(119)
where $`dV=4\pi rHdr`$ represents the differential (cylindrical) volume element, and $`T_e`$ denotes the electron temperature. We can obtain an upper limit on the X-ray luminosity by assuming that the electron temperature is equal to the ion temperature. Based on the detailed disk structures associated with models 2 and 5, we find that $`L_{\mathrm{rad}}/L_{\mathrm{jet}}10^2`$ and $`L_{\mathrm{rad}}/L_{\mathrm{jet}}10^5`$, respectively. However, in an actual ADAF disk, the X-ray luminosity will of course be substantially smaller than these values because the electron temperature is roughly three orders of magnitude lower than the ion temperature. Hence our neglect of radiative losses is completely justified, as expected for ADAF disks.
## 8 CONCLUSION
In this paper we have demonstrated that particle acceleration at a standing, isothermal shock in an ADAF accretion disk can energize the relativistic protons that power the jets emanating from radio-loud sources containing black holes. The work presented here represents a new type of synthesis that combines the standard model for a transonic ADAF flow with a self-consistent treatment of the relativistic particle transport occurring in the disk. The energy lost from the background (thermal) gas at the isothermal shock location results in the acceleration of a small fraction of the background particles to relativistic energies. One of the major advantages of our coupled, global model is that it provides a single, coherent explanation for the disk structure and the formation of the outflow based on the well-understood concept of first-order Fermi acceleration in shock waves. The theory employs an exact mathematical approach in order to solve simultaneously the combined hydrodynamical and particle transport equations.
The analysis presented here closely parallels the early studies of cosmic-ray shock acceleration. As in those first investigations (e.g., Blandford & Ostriker 1978), we have employed an idealized model in which the pressure of the accelerated particles is assumed to be negligible compared with that of the thermal background gas (the โtest particleโ approximation). In order to check the self-consistency of this assumption, we have confirmed that the total pressure is dominated by the pressure of the background (thermal) gas throughout most of the disk. However, in the vicinity of the shock the two pressures can become comparable and this suggests that the dynamical results will change slightly if the test particle approximation is relaxed. We plan to consider this question in future work by developing a โtwo-fluidโ version of our model that includes the particle pressure, in analogy with the โcosmic-ray modified shockโ scenario for cosmic-ray acceleration (Becker & Kazanas 2001; Drury & Vรถlk 1981).
We have presented detailed results that confirm that the general properties of the jets observed in M87 and Sgr A$`^{}`$can be understood within the context of our disk/shock/outflow model. In particular, our results indicate that the shock acceleration mechanism can produce relativistic outflows with terminal Lorentz factors and total powers comparable to those observed in M87 and Sgr A$`^{}`$. However, in principle even higher efficiencies can be achieved by varying the upstream energy transport rate $`ฯต_{}`$ which is the fundamental free parameter in our model. The buildup of the particle pressure in such high-efficiency situations would require relaxation of the test-particle approximation, as discussed above. In this paper we have focused on inviscid disks, which are the simplest to analyze. While the inviscid model provides useful insight into the importance of shock acceleration in ADAF disks, this restriction clearly must be lifted in the future, since viscosity plays a key role in determining the structure of an actual accretion disk. We are currently developing a self-consistent viscous disk model in order to explore shock formation and particle acceleration in a more rigorous context. However, we do not expect the presence of viscosity to alter any of the basic conclusions reached in this paper because significant particle acceleration will occur regardless of the viscosity, provided a shock is present. The existence of shocks in viscous disks is a controversial issue, but several studies suggest that shock formation is possible provided the viscosity is relatively low. In the absence of a consensus regarding the possible presence of shocks in accretion disks, we believe that it is important to study models with shocks in order to develop theoretical predictions that can be tested observationally.
The shock acceleration mechanism analyzed in this paper is effective only in rather tenuous, hot disks, and therefore we conclude that our model may help to explain the observational fact that the brightest X-ray AGNs do not possess strong outflows, whereas the sources with low X-ray luminosities but high levels of radio emission do. We suggest that the gas in the luminous X-ray sources is too dense to allow efficient Fermi acceleration of a relativistic particle population, and therefore in these systems, the gas simply heats as it crosses the shock. Conversely, in the tenuous ADAF accretion flows studied here, the relativistic particles are able to avoid thermalization due to the long collisional mean free path, resulting in the development of a significant nonthermal component in the particle distribution which powers the jets and produces the strong radio emission. We therefore conclude that the coupled, self-consistent theory for the disk structure and the particle acceleration investigated here is capable of powering the outflows observed in many radio-loud systems containing black holes.
The authors are grateful to Dr. Lev Titarchuk for providing a number of useful comments on the manuscript, and also to the anonymous referee for several insightful suggestions that significantly improved the paper.
## APPENDICES
## Appendix A Treatment of the Vertical Structure
In principle, the pressure $`P`$, density $`\rho `$, diffusion coefficient $`\kappa `$, Greenโs function $`f__\mathrm{G}`$, and velocity components $`v_r`$ and $`v_z`$ in the disk all display significant variations in the vertical ($`z`$) direction. Following Abramowicz & Chakrabarti (1990), we will use the first five quantities to represent vertical averages over the disk structure at radius $`r`$. However, the vertical variation of the velocity component $`v_z`$ must be treated differently. Here, we assume for simplicity that the vertical expansion is homologous, and therefore the vertical velocity variation is given by
$$v_z(r,z)=B(r)z.$$
(A1)
It follows that the vertical velocity at the surface of the disk, $`z=H(r)`$, can be written as
$$v_z(r,z)|_{z=H}=B(r)H(r).$$
(A2)
In a steady-state situation, we can also express the vertical velocity at the disk surface using
$$v_z(r,z)|_{z=H}=v_r\frac{dH}{dr}.$$
(A3)
By combining the two previous expressions, we find that the function $`B(r)`$ is given by
$$B(r)=v_r\frac{d\mathrm{ln}H}{dr}.$$
(A4)
This result will prove useful when we vertically integrate the transport equation. Note that in terms of $`B(r)`$, we can write the divergence of the flow velocity $`\stackrel{}{v}`$ in cylindrical coordinates as
$$\stackrel{}{}\stackrel{}{v}=\frac{1}{r}\frac{}{r}\left(rv_r\right)+\frac{v_z}{z}=\frac{1}{r}\frac{}{r}\left(rv_r\right)+B(r),$$
(A5)
where we have assumed azimuthal symmetry. Application of equation (A4) now yields
$$\stackrel{}{}\stackrel{}{v}=\frac{1}{Hr}\frac{}{r}\left(rHv_r\right).$$
(A6)
The steady-state transport equation expressed in cylindrical coordinates is (see eq. )
$`v_r{\displaystyle \frac{f__\mathrm{G}}{r}}+v_z{\displaystyle \frac{f__\mathrm{G}}{z}}`$ $`=`$ $`{\displaystyle \frac{1}{3}}\left[{\displaystyle \frac{1}{r}}{\displaystyle \frac{}{r}}\left(rv_r\right)+{\displaystyle \frac{dv_z}{dz}}\right]E{\displaystyle \frac{f__\mathrm{G}}{E}}+{\displaystyle \frac{1}{r}}{\displaystyle \frac{}{r}}\left(r\kappa {\displaystyle \frac{f__\mathrm{G}}{r}}\right)`$ (A7)
$`+`$ $`{\displaystyle \frac{\dot{N}_0\delta (EE_0)\delta (rr_{})}{(4\pi E_0)^2r_{}H_{}}}A_0c\delta (rr_{})f__\mathrm{G}.`$
Operating on equation (A7) with $`_0^{\mathrm{}}๐z`$ and applying equation (A1) yields, after partially integrating the term containing $`v_z`$ on the left-hand side,
$`v_r{\displaystyle \frac{}{r}}(Hf__\mathrm{G})HBf__\mathrm{G}`$ $`=`$ $`{\displaystyle \frac{1}{3}}\left[{\displaystyle \frac{1}{r}}{\displaystyle \frac{d}{dr}}\left(rv_r\right)+B\right]HE{\displaystyle \frac{f__\mathrm{G}}{E}}+{\displaystyle \frac{1}{r}}{\displaystyle \frac{}{r}}\left(rH\kappa {\displaystyle \frac{f__\mathrm{G}}{r}}\right)`$ (A8)
$`+`$ $`{\displaystyle \frac{\dot{N}_0\delta (EE_0)\delta (rr_{})}{(4\pi E_0)^2r_{}}}A_0cH_{}\delta (rr_{})f__\mathrm{G},`$
where the symbols $`f__\mathrm{G}`$, $`v_r`$, and $`\kappa `$ now refer to vertically averaged quantities. Using equations (A4), (A5), and (A6), we can rewrite the vertically integrated transport equation as
$`Hv_r{\displaystyle \frac{f__\mathrm{G}}{r}}`$ $`=`$ $`{\displaystyle \frac{1}{3r}}{\displaystyle \frac{}{r}}\left(rHv_r\right)E{\displaystyle \frac{f__\mathrm{G}}{E}}+{\displaystyle \frac{1}{r}}{\displaystyle \frac{}{r}}\left(rH\kappa {\displaystyle \frac{f__\mathrm{G}}{r}}\right)`$ (A9)
$`+`$ $`{\displaystyle \frac{\dot{N}_0\delta (EE_0)\delta (rr_{})}{(4\pi E_0)^2r_{}}}A_0cH_{}\delta (rr_{})f__\mathrm{G}.`$
## Appendix B Derivation of the Escape Parameter
The dimensionless parameter $`A_0`$ appearing in equation (58) determines the rate of particle escape through the surface of the disk due to random walks occurring near the shock location. Since the particles are accelerated as a consequence of collisions with magnetic waves, we will assume that the thickness of the shock is comparable to the magnetic mean free path, $`\lambda _{\mathrm{mag}}`$. In order to estimate $`A_0`$, we model the escape of the particles from the disk using the analogy of โleakageโ through an opening in a cylindrical pipe with radius equal to the half-thickness of the disk at the shock location, $`H_{}`$. The length of the open section of the pipe is set equal to the shock thickness $`\lambda _{\mathrm{mag}}`$. The particle number density in the open section is governed by the equation
$$v_x\frac{dn_r}{dx}=\frac{n_r}{t_{\mathrm{esc}}},$$
(B1)
where $`v_x`$, $`n_r`$, and $`t_{\mathrm{esc}}`$ represent the flow velocity, the relativistic particle number density, and the average time for the particles to escape through the open walls of the pipe via diffusion. Upon integration, the solution to equation (B1) is given by
$$n_r(x)=n_0\mathrm{exp}\left(\frac{x}{v_xt_{\mathrm{esc}}}\right),$$
(B2)
where $`n_0`$ is the incident number density as the flow encounters the opening in the pipe, at $`x=0`$. We can approximate the solution for $`n_r(x)`$ by performing a Taylor expansion around $`x=0`$, which yields
$$n_r(x)n_0\left(1\frac{x}{v_xt_{\mathrm{esc}}}\right).$$
(B3)
The fraction of particles that escape from the pipe can therefore be estimated by setting $`x=\lambda _{\mathrm{mag}}`$ to obtain
$$f_{\mathrm{esc}}=1\frac{n_r}{n_0}=\frac{\lambda _{\mathrm{mag}}}{v_xt_{\mathrm{esc}}}.$$
(B4)
In order to make contact with the disk application, we note that according to equations (75) and (76), the fraction of particles that escape as the gas crosses the isothermal shock is given by
$$f_{\mathrm{esc}}=A_0\frac{c}{v_{}},$$
(B5)
where $`v_{}(v_++v_{})/2`$ is the mean velocity at the shock, and we have assumed that advection dominates over diffusion. Eliminating $`f_{\mathrm{esc}}`$ between equations (B4) and (B5), and setting $`v_x=v_{}`$, we find that
$$A_0=\frac{\lambda _{\mathrm{mag}}}{ct_{\mathrm{esc}}}.$$
(B6)
Within the context of our one-dimensional model for the particle transport in the disk, the mean escape time $`t_{\mathrm{esc}}`$ is related to $`\lambda _{\mathrm{mag}}`$ and the disk half-thickness at the shock $`H_{}`$ via
$$t_{\mathrm{esc}}=\frac{H_{}}{v_{\mathrm{diff}}}=\frac{H_{}^2}{c\lambda _{\mathrm{mag}}},$$
(B7)
where $`v_{\mathrm{diff}}=c\lambda _{\mathrm{mag}}/H_{}`$ denotes the vertical diffusion velocity of the protons in the tangled magnetic field near the shock, which is valid provided $`H/\lambda _{\mathrm{mag}}>1`$. Eliminating $`t_{\mathrm{esc}}`$ between equations (B6) and (B7) then yields
$$A_0=\left(\frac{\lambda _{\mathrm{mag}}}{H_{}}\right)^2<1.$$
(B8)
The diffusion coefficient at the shock is related to the magnetic mean free path by the standard expression (e.g., Reif 1965)
$$\kappa =\frac{c\lambda _{\mathrm{mag}}}{3},$$
(B9)
and therefore equation (B8) can be rewritten as
$$A_0=\left(\frac{3\kappa _{}}{cH_{}}\right)^2,$$
(B10)
where $`\kappa _{}(\kappa _{}+\kappa _+)/2`$ denotes the average of the upstream and downstream values of $`\kappa `$ on either side of the shock.
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# Laser nanotraps and nanotweezers for cold atoms: 3D gradient dipole force trap in the vicinity of Scanning Near-field Optical Microscope tip
## Abstract
Using a two-dipole model of an optical near-field of Scanning Near-field Optical Microscope tip, i. e. taking into account contributions of magnetic and electric dipoles, we propose and analyze a new type of 3D optical nanotrap found for certain relations between electric and magnetic dipoles. Electric field attains a minimum value in vacuum in the vicinity of the tip and hence such a trap is quite suitable for manipulations with cold atoms.
Recent enormous progress in the study of laser-cooled atoms and molecules put forward the problem of their using for quantum computing, frequency standard construction, and other technological applications (see e.g. \[1-3\]). To achieve this goal, as well as for the further progress of fundamental experiments in the field, new compact traps and new methods of โhandlingโ the cold atoms, for example their transportation to/from a trap or between different traps, should be elaborated. An example of successful work in this direction is given e.g. by โatom-chip technologyโ experiments , where guiding of cold atoms along the wires has been demonstrated. Further miniaturization of such devices is desirable. Ideally, one should have at his/her disposal true โsingle cold atom nanotweezersโ which have very small sizes and only slightly perturb the trap. Evidently, the near-field optical configurations, based on subwavelength-size aperture in the apex of a sharp fiber tip or local field enhancement in the vicinity of sharp conducting tips (see e.g. for review on near-field optics) look rather promising. Hence it is not surprising that a number of near-field optical traps/tweezers have been proposed \[7 - 15\].
However, all of them have certain drawbacks and, we believe, this is the main reason why, to the best of our knowledge, none has been realized up to now. Putting aside the configurations where an extremum of the optical field is achieved on the surface of the tip (configuration obviously inappropriate for atom trap), among these drawbacks we could mention that the traps proposed needed different hard-to-control non-optical interactions (centripetal potential, gravitational interaction, van der Waals force, etc.) to be closed, can be realized only in the non practical light reflection mode from the subwavelength aperture, and so on.
Here we propose and analyze the true โfree standingโ, or โsupport-freeโ purely optical 3D trap emerging in vacuum in the vicinity of an aperture of the Scanning Near-field Optical Microscope (SNOM) tip. The characteristic size of our trap is small in comparison with laser wavelength , so one indeed can speak about nanotrap. Because methods of Angstrom-precision motion of SNOM tip are well elaborated, the same construction is cold atom nanotweezers.
Our analysis is based on the two-dipole model of optical near field occurring in the vicinity of this tip, see Fig. 1. Such a model has been established recently, when it has been shown that in addition to the classic Bethe and Bouwkamp consideration, where in the case of normal incidence optical near field is modeled by one magnetic dipole M (see e.g. Refs. \[16-18\]), a field of an electric dipole P should be added to describe correctly optical near-field of a real fiber tip Drezet1 -Drezet2 . This effect is due mainly to the conical shape of the end area of such tip, and for a variety of tips and light polarizations different relations between the values and mutual orientations of M, P dipoles can be anticipated.
In near field of tip as in any laser fields, resonant atoms are subject to the optical dipole force with the potential Ashkin ,Letokhov :
$$U=\frac{\mathrm{}\mathrm{\Omega }}{2}\mathrm{ln}(1+\frac{\mu ^2E_{las}^2}{\mathrm{}^2\gamma ^2})\frac{\mathrm{}\mathrm{\Omega }}{2}\frac{\mu ^2E_{las}^2}{\mathrm{}^2\gamma ^2}$$
(1)
Here $`\mu `$ is an atomic transition dipole moment, $`\gamma ^2=\mathrm{\Omega }^2+(\mathrm{\Gamma }/2)^2`$ where $`\mathrm{\Gamma }`$ is a natural line width, and $`\mathrm{\Omega }=\omega \omega _0`$ is detuning between laser frequency $`\omega `$ and resonant frequency of an atom $`\omega _0`$. Throughout the paper we consider the blue detuning $`\mathrm{\Omega }>0`$, which results in atom trapping at the minimum of laser field intensity.It is well known that blue-detuned gradient force optical traps based on minimum of an electric field have essential advantages (they do not heat trapped atoms, etc.) in comparison with those based on maximum of an electric field using red-detuned light Letokhov .
To characterize the proposed 3D nanotrap let us consider electric field of light in the vicinity of a SNOM tip (i. e. in the near-field regions of both dipoles). Corresponding electric field of an electric dipole P has the form:
$$๐_e=\frac{๐}{R^3}+3\frac{๐\left(\mathrm{๐๐}\right)}{R^5}$$
(2)
where R is the radius vector from the dipole position to an observation point. (SGCE units are used throughout the paper). Vector potential of a magnetic dipole M has the following form
$$๐=\frac{\left[\mathrm{๐๐}\right]}{R^3}$$
(3)
and according to Faradayโs law electric field of this magnetic dipole in the near-field region has the form
$$๐_m=ik\frac{\left[\mathrm{๐๐}\right]}{R^3}$$
(4)
where $`k={\displaystyle \frac{\omega }{c}}`$ is wavevector in free space. When two dipoles are located at the same point {0,0,0}, intensity of total electric field $`๐=๐_m+๐_e`$ behind an aperture can be presented in the form
$$\begin{array}{c}E^2=\left|๐_e+๐_m\right|^2=\hfill \\ \\ \frac{k^2\left|๐\right|^2R^2k^2\left(\mathrm{๐๐}\right)\left(๐^{}๐\right)+\left|๐\right|^2+3\left(\mathrm{๐๐}\right)\left(๐^{}๐\right)/R^2ik๐\left(\left[๐^{}๐\right]\left[\mathrm{๐๐}^{}\right]\right)}{R^6}\hfill \end{array}$$
(5)
Below we will speak about this value as about an โintensityโ having in mind its obvious connection with the intensity of laser light $`I_{las}`$ โseepingโ through the aperture of a SNOM tip: $`I_{las}={\displaystyle \frac{c}{8\pi }}E^2`$, where $`c`$ is speed of light .
It is naturally to look first for an axially symmetric trap. Such a trap occurs in particularly if
$$\begin{array}{c}M_z=P_z=0\hfill \\ M_xM_y^{}+M_yM_x^{}=0\hfill \\ \left|M_x\right|=\left|M_y\right|=M_0\hfill \\ P_xP_y^{}+P_yP_x^{}=0\hfill \\ \left|P_x\right|=\left|P_y\right|=P_0\hfill \end{array}$$
(6)
If these conditions hold true, expression for the intensity can be rewritten in the form
$$E^2=\frac{2k^2M_0^2R^2k^2M_0^2\rho ^2+2P_0^2+3P_0^2\rho ^2/R^2zkQ}{R^6}$$
(7)
where $`Q=i\left(\left[๐^{}๐\right]_z\left[\mathrm{๐๐}^{}\right]_z\right)`$ and $`\rho =\sqrt{x^2+y^2}`$ is the distance from the symmetry axis to an observation point.
Analysis shows that if we put the following additional condition on the dipole momenta
$$Q^2=16M_0^2P_0^2$$
(8)
intensity (7) becomes equal to zero at the points
$$(x=y=0,kz=\frac{P_0}{M_0})\mathrm{or}(x=y=0,kz=\frac{P_0}{M_0})$$
(9)
In its turn, it is possible to show that (8) will be satisfied in the case
$$\begin{array}{c}๐=P_0\{i,1,0\},\hfill \\ ๐=\pm M_0\{i,1,0\}\hfill \end{array}$$
(10)
It means that dipole momenta should be collinear to ensure the 3D trap.
As intensity is a positive function of coordinates, zeros of intensity (9)correspond to true 3D minimum of an electric field. In Fig.2 the distribution of intensity in $`xz`$ plane is shown for the case $`P_0=20;M_0=10`$.
Intensity (7) has another point of extremum
$$\left(x=0,y=0,kz=\frac{3}{2}\frac{P_0}{M_0}\right)\mathrm{or}\left(x=0,y=0,kz=\frac{3}{2}\frac{P_0}{M_0}\right)$$
(11)
which is a saddle point (see Fig. 2.). The value of intensity at this point is
$$E^2=\frac{32}{729}k^6\frac{M_0^6}{P_0^4}$$
(12)
This quantity can be used to estimate the potential well depth. For such an estimation one can use $`M_0=P_0=E_0a^3`$, where $`E_0`$is an amplitude of the incoming light wave in aperture plane and $`a`$ is the radius of an aperture \[16-19\]. Substituting these values into (12), we get for the potential well depth
$$\mathrm{\Delta }E^2\frac{32}{729}\left(ka\right)^6E_0^2$$
(13)
For modern tips $`ka1/2`$ and hence $`\mathrm{\Delta }E^210^3E_0^2`$. For typical experimental conditions, an intensity of optical near field at the aperture of SNOM tip is about $`I_0=10^310^4W/cm^2`$ . This means that e. g. for alkaline atoms, which are characterized by resonant dipole moments of the order of 10<sup>-17</sup> SGSE and natural line widths $`\mathrm{\Gamma }10^710^8s^1`$ (for example, for the 6S<sub>1/2</sub>โ6P<sub>3/2</sub> D2 transition of cesium atom at $`\lambda `$=852 nm, $`\mu `$=8.01$``$10<sup>-18</sup> CGSE and $`\mathrm{\Gamma }=3.0710^7s^1`$ Rafac ), traps with the depth of the order of a few milliKelvin, what is quite standard for optical dipole gradient force โ based traps, can be realized when using $`\mathrm{\Omega }100\mathrm{\Gamma }`$, see (1). It is very important that under these conditions the trap will have about 10 energy levels of atomic motion with lowest level being about $`10^4K`$.
It is worthwhile to note, that existence of such a 3D trap is highly nontrivial, because the relevant field components have very complicated structure (see Figs. 3, 4, 5), and it seems very difficult to provide a minimum for their sum.
Mathematically, by varying P and M ratio and other parameters of the problem, our trap can be placed at any point on the symmetry axis. However, the two dipole approximation used is not valid very close to the aperture plane $`z=0`$. Its validity starts from $`za`$, where $`a`$ is the radius of an aperture. Hence we should consider only such parameters where the trap position occurs not too close to the aperture plane. Besides the condition $`za`$ allows us to neglect van der Waals attractive force which is always important in close vicinity of tip surface.
For our trap the minimum of intensity is stable against small perturbations. For example, if relation between the dipole values is
$$\begin{array}{c}๐=P_0\{i,1,0\},\hfill \\ ๐=\pm M_0e^{i\phi }\{i,1,0\}\hfill \end{array}$$
(14)
then true 3D minimum still exists provided $`\mathrm{sin}\left|\phi \right|<\frac{1}{5}`$ or $`0.2<\phi <0.2`$. Small variations of z-components of momenta ($`P_z`$ and $`M_z)`$ also result in small variations of trap position. Small mutual displacements (splitting) of dipoles in axial and/or radial directions also result in variation of the trap. In the case when magnetic dipole is moving away from the trap, that is in the case when electric dipole is placed between the trap and the magnetic one, the trap suffers only minor shifts. In the case when magnetic dipole is placed between an electric one and the trap, the intensity minimum disappears for large enough splitting, see Fig.6.
The retardation effects in the near field region generally are small. Nevertheless to estimate their influence let us consider the full electric field of electric and magnetic dipoles:
$$\begin{array}{c}๐_e=๐\left(\frac{1}{R^3}+\frac{ik}{R^2}+\frac{k^2}{R}\right)+\frac{๐\left(\mathrm{๐๐}\right)}{R^2}\left(\frac{3}{R^3}\frac{3ik}{R^2}\frac{k^2}{R}\right)\hfill \\ ๐_m=\frac{\left[\mathrm{๐๐}\right]}{R}\left(\frac{ik}{R^2}+\frac{k^2}{R}\right)\hfill \end{array}$$
(15)
Analysis shows that for such a case 3d trap survives only when $`\xi =P_0/M_0<1.5035`$. Positions of the trapping region and the saddle point are given respectively by the formulae
$$(x=y=0,kz=\xi \xi ^3+3\xi ^5+\mathrm{})\mathrm{or}(x=y=0,kz=(\xi \xi ^3+3\xi ^5+\mathrm{}))$$
(16)
and
$$(x=y=0,kz=\frac{3}{2}\xi \frac{9}{16}\xi ^3\frac{189}{128}\xi ^5+\mathrm{})\mathrm{or}(x=y=0,kz=(\frac{3}{2}\xi \frac{9}{16}\xi ^3\frac{189}{128}\xi ^5\mathrm{}))$$
(17)
Intensity at the bottom of the potential well is now nonzero
$$E_{\mathrm{min}}^2=2k^6M_0^2\left(\xi ^2\xi ^4+\mathrm{}\right)$$
(18)
while the intensity at the saddle point is
$$E_{saddle}^2=k^6M_0^2\left(\frac{32}{729\xi ^4}+\frac{40}{81\xi ^2}+\mathrm{}\right)$$
(19)
These values determine the depth of our trap. Comparing results (16)-(19), where retardation effects are taken into account, with the quasistatic results (11) we see that in the case of magnetic dipole domination $`\left(\xi 1\right)`$ the position of the trap remains in the near field region and hence the retardation effects have only minor influence. On the other hand, in the case of substantially large amplitude of electric dipole, the retardation effects destroy our trap.
Hence we have shown that minimum of an electric field with the size smaller than the light wavelength and depth of a few milliKelvin did occur in vacuum in the vicinity of a SNOM tip. This minimum is stable against perturbations and can be tuned both in position and depth by changing relation between $`P_0`$ and $`M_0`$ and incident light electric field amplitude E0. This attests the trap proposed as a very promising basic element for future cold atom nanotraps and nanotweezers. Finally, we would like to note the following. Despite the two dipole model of optical near-field, whose using is inherent to obtain the reported results, nowadays seems is well established and supported by experiments \[19-21\], complete and rigorous analysis of all possible dipole configurations occurring for near field of a real tip is still lacking. In particular, the conditions to be imposed on experimental setup to obtain the trapping configuration of magnetic and electric dipoles also remain to be understood. Nevertheless, we believe that broad possibilities to vary SNOM tip shapes and coatings as well as to vary other parameters (e. g. incoming light polarization) give enough hope for practical realization of such a configuration. Indeed, it can quite happen that this is already done at least for one of plenty of the tips demonstrated up to date.
###### Acknowledgements.
The authors are grateful to Swiss National Science Foundation and Russian Foundation for Basic Research (V.V.K., grant # 04-02-16211) for financial support of this work.
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# 0. Introduction
## 0. Introduction
A sequence of positive integers $`w_1,w_2,\mathrm{},w_n`$ is called an ascending wave if $`w_{i+1}w_iw_iw_{i1}`$ for $`2in1`$. For $`k,r^+`$, let $`AW(k;r)`$ be the least positive integer such that under any $`r`$-coloring of $`[1,AW(k;r)]`$ there exists a $`k`$-term monochromatic ascending wave. Although guaranteed by van der Waerdenโs theorem, the existence of $`AW(k;r)`$ can be guaranteed independently, as we will show.
Bounds on $`AW(k;2)`$ have appeared in the literature. Brown, Erdลs, and Freedman showed that for all $`k1`$,
$$\begin{array}{cc}k^2k+1AW(k;2)\frac{k^3}{3}\frac{4k}{3}+3.\hfill & \end{array}$$
Soon after, Alon and Spencer showed that for sufficiently large $`k`$,
$$\begin{array}{cc}AW(k;2)>\frac{k^3}{10^{21}}\frac{k^2}{10^{20}}\frac{k}{10}+4.\hfill & \end{array}$$
Recently, Landman and Robertson proposed the refinement of the bounds on $`AW(k;2)`$ and the study of $`AW(k;r)`$ for $`r3`$. Here, we offer bounds on $`AW(k;r)`$ for all $`r1`$, improving upon the previous upper bound for $`AW(k;2)`$.
## 1. An Upper Bound
To show that $`AW(k;r)O(k^{2r1})`$ is straightforward. We will first show that $`AW(k;r)k^{2r1}`$ by induction on $`r`$ (and hence prove the existence of $`AW(k;r)`$ without appealing to van der Waerdenโs theorem.) The case $`r=1`$ is trivial, thus, for $`r2`$, assume $`AW(k;r1)k^{2r3}`$ and consider any $`r`$-coloring of $`[1,k^{2r1}]`$. Set $`w_1=1`$ and let the color of $`1`$ be red. In order to avoid a $`k`$-term monochromatic ascending wave there must exist an integer $`w_2[2,k^{2r3}+1]`$ that is colored red, lest the inductive hypothesis guarantee a $`k`$-term monochromatic ascending wave of some color other than red (and we are done). Similarly, there must be an integer $`w_3[w_2+(w_2w_1),w_2+(w_2w_1)+k^{2r3}1]`$ that is colored red to avoid a monochromatic $`k`$-term ascending wave. Iterating this argument defines a monochromatic (red) $`k`$-term ascending wave $`w_1,w_2,\mathrm{},w_k`$, provided that $`w_kk^{2r1}`$. Since for $`i2`$, $`w_{i+1}w_i+(w_iw_{i1})+k^{2r3}`$ we see that $`w_{i+1}w_iik^{2r3}`$ for $`i1`$. Hence, $`w_kw_1=_{i=1}^{k1}(w_{i+1}w_i)_{i=1}^{k1}ik^{2r3}k^{2r1}1`$ and we are done.
In this section we provide a better upper bound. Our main theorem in this section follows.
Theorem 1 For fixed $`r1`$,
$$\begin{array}{cc}AW(k;r)\frac{k^{2r1}}{(2r1)!}(1+o(1)).\hfill & \end{array}$$
We will prove Theorem 1 via a series of lemmas, but first we introduce some pertinent notation.
Notation For $`k2`$ and $`MAW(k;r)`$, let $`\mathrm{\Psi }^M(k;r)`$ be the collection of all $`r`$-colorings of $`[1,M]`$. For $`\psi \mathrm{\Psi }^M(k;r)`$, let $`\chi _k(\psi )`$ be the set of all monochromatic $`k`$-term ascending waves under $`\psi `$. For each monochromatic $`k`$-term ascending wave $`w=\{w_1,w_2,\mathrm{},w_k\}\chi _k(\psi )`$, define the $`i^{\mathrm{th}}`$ difference, $`d_i(w)=w_{i+1}w_i`$, for $`1ik1`$. For $`\psi \mathrm{\Psi }^M(k;r)`$, define
$$\delta _k(\psi )=\mathrm{min}\{d_{k1}(w)|w\chi _k(\psi )\},$$
i.e., the minimum last difference over all monochromatic $`k`$-term ascending waves under $`\psi `$. Lastly, define
$$\mathrm{\Delta }^M(k;r)=\mathrm{max}\{\delta _k(\psi )|\psi \mathrm{\Psi }^M(k;r)\}.$$
These concepts will provide us with the necessary tools to prove Theorem 1.
We begin with a recursive bound on $`AW(k;r)`$, which also proves the existence of $`AW(k;r)`$ without appealing to van der Waerdenโs theorem on arithmetic progressions.
Lemma 1.1 For $`k,r1`$, let $`M(k;1)=k`$, $`M(1;r)=1`$, $`M(2;r)=r+1`$, and define, for $`k3`$ and $`r2`$,
$$M(k;r)=M(k1;r)+\mathrm{\Delta }^{M(k1;r)}(k1;r)+M(k;r1)1.$$
Then, for all $`k,r1`$, $`AW(k;r)M(k;r).`$
Proof. Noting that the cases $`k+r=2,3,`$ and $`4`$ are, by definition, true, we proceed by induction on $`k+r`$ using $`k+r=5`$ as our basis. We have $`M(3;2)=7`$. An easy calculation shows that $`AW(3;2)=7`$. So, for some $`n5`$, we assume Lemma 1.1 holds for all $`k,r1`$ such that $`k+r=n`$. Now, consider $`k+r=n+1`$. The result is trivial when $`k=1\mathrm{or}\mathrm{\hspace{0.33em}2}`$, or if $`r=1`$, thus we may assume $`k3`$ and $`r2`$. Let $`\psi `$ be an $`r`$-coloring of $`[1,M(k;r)]`$. We will show that $`\psi `$ admits a monochromatic $`k`$-term ascending wave, thereby proving Lemma 1.1.
By the inductive hypothesis, under $`\psi `$ there must be a monochromatic $`(k1)`$-term ascending wave $`w=\{w_1,w_2,\mathrm{},w_{k1}\}[1,M(k1;r)]`$ with $`d_{k2}(w)\mathrm{\Delta }^{M(k1;r)}(k1;r)`$. Let
$$N=[w_{k1}+\mathrm{\Delta }^{M(k1;r)}(k1;r),w_{k1}+\mathrm{\Delta }^{M(k1;r)}(k1;r)+M(k;r1)1].$$
If there exists $`qN`$ colored identically to $`w`$, then $`w\{q\}`$ is a monochromatic $`k`$-term ascending wave, since $`qw_{k1}\mathrm{\Delta }^{M(k1;r)}(k1;r)d_{k2}(w)`$. If there is no such $`qN`$, then $`N`$ contains integers of at most $`r1`$ colors. Thus, since $`|N|=M(k;r1)`$, the inductive hypothesis guarantees that we have a monochromatic $`k`$-term ascending wave in $`N`$. As
$$w_{k1}+\mathrm{\Delta }^{M(k1;r)}(k1;r)+M(k;r1)1M(k;r),$$
this completes the proof. $`\mathrm{}`$
Corollary 1.2 Let $`k3`$ and $`r2`$. Let $`M(k;r)`$ be as in Lemma 1.1. Then
$$\mathrm{\Delta }^{M(k;r)}(k;r)\mathrm{\Delta }^{M(k1;r)}(k1;r)+M(k;r1)1.$$
Proof. Let $`\psi `$, $`w`$, and $`N`$ be as defined in the proof of Lemma 1.1. If there exists $`qN`$ colored identically to $`w`$, then
$$\delta _k(\psi )d_{k1}(w\{q\})\mathrm{\Delta }^{M(k1;r)}(k1;r)+M(k;r1)1.$$
If there is no such $`qN`$, then there exists a monochromatic $`k`$-term ascending wave, say $`v`$, in $`N`$. Hence, $`\delta _k(\psi )d_{k1}(v)M(k;r1)(k1)`$. Since $`\psi `$ was chosen arbitrarily, it follows that
$$\mathrm{\Delta }^{M(k;r)}(k;r)\mathrm{\Delta }^{M(k1;r)}(k1;r)+M(k;r1)1.$$
$`\mathrm{}`$
Now, as it is easily seen that $`\mathrm{\Delta }^{M(2;r)}(2;r)=r`$ for all $`r1`$, Lemma 1.1 and Corollary 1.2 can be iterated to provide bounds on $`M(k;r)`$ for any $`k3`$ and $`r2`$. The following lemma will provide a means for recursively bounding $`M(k;r)`$, and thus $`AW(k;r)`$, by a function of $`k`$ for any $`r2`$.
Lemma 1.3 Let $`k3`$ and $`r2`$. Let $`M(k;r)`$ be as in Lemma 1.1. Then
$$\begin{array}{cc}M(k;r)\hfill & _{i=0}^{k3}((i+1)M(ki;r1))\frac{k^2}{2}+\frac{3k}{2}+(k1)r\hfill \end{array}$$
Proof. We proceed by induction on $`k`$. Consider $`M(3;r)`$. We have
$$M(3;r)=M(2;r)+\mathrm{\Delta }^{M(2;r)}(2;r)+M(3;r1)1.$$
Since $`M(2;r)=r+1`$ and $`\mathrm{\Delta }^{M(2;r)}(2;r)=r`$, we have
$$\begin{array}{cc}M(3;r)=M(3;r1)+2r=M(3;r1)\frac{3^2}{2}+\frac{3(3)}{2}+2r,\hfill & \end{array}$$
thereby finishing the case $`k=3`$ and arbitrary $`r`$. Now assume that Lemma 1.3 holds for some $`k3`$. Lemma 1.1, the inductive hypothesis, and Corollary 1.2 give us
$$\begin{array}{cc}M(k+1;r)\hfill & =M(k;r)+\mathrm{\Delta }^{M(k;r)}(k;r)+M(k+1;r1)1\hfill \\ & _{i=0}^{k3}((i+1)M(ki;r1))\frac{k^2}{2}+\frac{3k}{2}+(k1)r\hfill \\ & +\mathrm{\Delta }^{M(k;r)}(k;r)+M(k+1;r1)1\hfill \\ & _{i=0}^{k3}((i+1)M(ki;r1))\frac{k^2}{2}+\frac{3k}{2}+(k1)r\hfill \\ & +\mathrm{\Delta }^{M(2;r)}(2;r)+_{i=0}^{k3}M(ki;r1)\hfill \\ & +M(k+1;r1)(k2)1\hfill \\ & _{i=0}^{k2}((i+1)M(k+1i;r1))\frac{(k+1)^2}{2}+\frac{3(k+1)}{2}+kr\hfill \end{array}$$
as desired. $`\mathrm{}`$
Now, for $`r2`$, an upper bound on $`M(k;r)`$, and thus on $`AW(k;r)`$, can be obtained by iterating Lemma 1.3. We offer one additional lemma, from which Theorem 1 will follow by application of Lemma 1.1.
Lemma 1.4 For $`k3`$ and $`r1`$, there exists a polynomial $`p_r(k)`$ of degree at most $`2r2`$ such that
$$\begin{array}{cc}M(k;r)\frac{k^{2r1}}{(2r1)!}+p_r(k).\hfill & \end{array}$$
Proof. We have $`M(k;1)=k`$, so we can take $`p_1(k)=1`$, having degree $`0`$. We proceed by induction on $`r`$. Assume Lemma 1.4 holds for some $`r1`$ so that $`M(k;r)\frac{k^{2r1}}{(2r1)!}+p_r(k)`$, where $`p_r(k)`$ is a polynomial of degree at most $`2r2`$. Lemma 1.3 gives
$$\begin{array}{cc}M(k;r+1)\hfill & _{j=3}^k((kj+1)M(j;r))\frac{k^2}{2}+\frac{3k}{2}+(k1)(r+1)\hfill \\ & k_{j=3}^k\left(\frac{j^{2r1}}{(2r1)!}+p_r(j)\right)_{j=3}^k\left((j1)\left(\frac{j^{2r1}}{(2r1)!}+p_r(j)\right)\right)\hfill \\ & \frac{k^2}{2}+\frac{3k}{2}+(k1)(r+1).\hfill \end{array}$$
By Faulhaberโs formula , for some polynomial $`p_{r+1}(k)`$ of degree at most $`2r`$, we now have
$$M(k;r+1)k\frac{\frac{k^{2r}}{2r}}{(2r1)!}\frac{\frac{k^{2r+1}}{2r+1}}{(2r1)!}+p_{r+1}(k)=\frac{k^{2r+1}}{(2r+1)!}+p_{r+1}(k)$$
and the proof is complete. $`\mathrm{}`$
As $`AW(k;r)M(k;r)`$, Theorem 1 now follows, giving that for fixed $`r1`$,
$$\begin{array}{cc}AW(k;r)\frac{k^{2r1}}{(2r1)!}(1+o(1)).\hfill & \end{array}$$
Interestingly, Lemma 1.3 can also be used to show the following corollary.
Corollary 1.5 For fixed $`k3`$,
$$\begin{array}{cc}AW(k;r)\frac{2^{k2}}{(k1)!}r^{k1}(1+o(1)).\hfill & \end{array}$$
Proof. In analogy to Lemma 1.4, we show that for $`k3`$ and $`r2`$, there exists a polynomial $`s_k(r)`$ of degree at most $`k2`$ such that
$$\begin{array}{cc}M(k;r)\frac{2^{k2}}{(k1)!}r^{k1}+s_k(r).\hfill & \end{array}$$
We proceed by induction on $`k`$. Let $`r2`$ be arbitrary. Lemma 1.3 gives that
$$M(3;r)=M(3;r1)+2r.$$
As $`M(3;1)=3`$, we now have a recursive definition of $`M(3;r)`$ for all $`r2`$. We get, for $`r2`$,
$$\begin{array}{cc}M(3;r)=M(3;1)+_{i=2}^r2i=r^2+r+1,\hfill & \end{array}$$
which serves as our basis. Now, for given $`k4`$, let $`\widehat{s}_3(r)=(k1)r\frac{k^2}{2}+\frac{3k}{2}`$ and assume Corollary 1.5 holds for all integers $`3jk1`$ and for all $`r2`$. Lemma 1.3 yields
$$\begin{array}{cc}M(k;r)\hfill & _{i=0}^{k3}((i+1)M(ki;r1))+\widehat{s}_3(r)\hfill \\ & =M(k;r1)+_{i=1}^{k3}((i+1)M(ki;r1))+\widehat{s}_3(r).\hfill \end{array}$$
Now, by the inductive hypothesis, for $`1ik3`$, we have that
$$\begin{array}{cc}M(ki;r1)\hfill & \frac{2^{ki2}}{(ki1)!}(r1)^{ki1}+s_{ki}(r1)\hfill \\ & =\frac{2^{ki2}}{(ki1)!}r^{ki1}+\stackrel{~}{s}_{ki}(r),\hfill \end{array}$$
where $`\stackrel{~}{s}_{ki}(r)`$ is polynomial of degree at most $`ki2k3`$. This gives us that
$$\begin{array}{cc}_{i=1}^{k3}(i+1)M(ki;r1)+\widehat{s}_3(r)\hfill & _{i=1}^{k3}\left((i+1)\left(\frac{2^{ki2}}{(ki1)!}r^{ki1}+\stackrel{~}{s}_{ki}(r)\right)\right)+\widehat{s}_3(r)\hfill \\ & =2\frac{2^{k3}}{(k2)!}r^{k2}+\stackrel{ห}{s}_{k1}(r),\hfill \end{array}$$
where $`\stackrel{ห}{s}_{k1}(r)`$ is a polynomial of degree at most $`k3`$. Hence, we have
$$\begin{array}{cc}M(k;r)\hfill & M(k;r1)+2\frac{2^{k3}}{(k2)!}r^{k2}+\stackrel{ห}{s}_{k1}(r)\hfill \\ & =M(k;r1)+\frac{2^{k2}}{(k2)!}r^{k2}+\stackrel{ห}{s}_{k1}(r).\hfill \end{array}$$
As $`M(k;1)=k`$, we have a recursive bound on $`M(k;r)`$ for $`r2`$. Faulhaberโs formula yields
$$\begin{array}{cc}M(k;r)M(k;1)+_{i=2}^r\left(\frac{2^{k2}}{(k2)!}i^{k2}+\stackrel{ห}{s}_{k1}(r)\right)\frac{2^{k2}}{(k1)!}r^{k1}+s_k(r),\hfill & \end{array}$$
where $`s_k(r)`$ is a polynomial of degree at most $`k2`$. By Lemma 1.1, the result follows. $`\mathrm{}`$
## 2. A Lower Bound for more than Three Colors
We now provide a lower bound on $`AW(k;r)`$ for arbitrary fixed $`r1`$. We generalize an argument of Alon and Spencer to provide our lower bound.
We will use $`\mathrm{log}(x)=\mathrm{log}_2(x)`$ throughout. Also, by $`k=x`$ for $`x^+`$ we mean $`k=x`$.
Theorem 2 For fixed $`r1`$ and any $`ฯต>0`$, for $`k`$ sufficiently large,
$$AW(k;r)\frac{k^{2r1ฯต}}{2^{r1}(40r)^{r^21}}.$$
The result is trivial for $`r=1`$ (since $`AW(k;1)=k`$). We will assume that for $`r2`$ the inequality holds for $`r1`$ and show that it holds for $`r`$.
We proceed by defining a certain type of random coloring. To this end, consider the $`r\times 2r`$ matrix $`A_0=(a_{ij})`$:
$$\left[\begin{array}{ccccccccc}0& 0& 1& 1& 2& 2& \mathrm{}& (r1)& (r1)\\ 0& 1& 1& 2& 2& 3& \mathrm{}& (r1)& 0\\ 0& 2& 1& 3& 2& 4& \mathrm{}& (r1)& 1\\ \mathrm{}& & & & & \mathrm{}& & \mathrm{}& \mathrm{}\\ 0& (r1)& 1& 0& 2& 1& \mathrm{}& (r1)& (r2)\end{array}\right]$$
where, for $`j[0,r1]`$, we have $`a_{i,(2j+1)}=j`$, for all $`1ir`$, and $`a_{i,2j+2}=i+j1(modr).`$
Next, we define $`A_j=A_0๐ฃ`$ where $``$ means addition modulo $`r`$ and $`๐ฃ`$ is the $`r\times 2r`$ matrix with all entries equal to $`j`$.
Consider the $`r^2\times 2r`$ matrix $`A=[A_0A_1A_2\mathrm{}A_{r1}]^t.`$
In the sequel, we will use the following notation.
Notation For $`r1`$, let
$$N_r=\frac{1}{2^{r1}(40r)^{r^21}}.$$
Fix $`ฯต>0`$. Let
$$b=AW(\frac{k}{10(4r4)};r1)1,$$
so that by the inductive hypothesis, we have
$$b>N_{r1}\left(\frac{k}{10(4r4)}\right)^{2r3ฯต/2}$$
for $`k`$ sufficiently large. Using the colors $`0,1,\mathrm{},r1`$, let $`\gamma _i`$ be an $`(r1)`$-coloring of $`b`$ consecutive integers with no monochromatic $`\frac{k}{10(4r4)}`$-term ascending wave, where the $`r1`$ colors used are $`\{0,1,\mathrm{},i1,i+1,i+2,\mathrm{},r1\}`$ (i.e., color $`i`$ is not used, and hence the subscript on $`\gamma `$).
We next describe how we randomly $`r`$-color $`[1,M_ฯต]`$, where
$$M_ฯต=N_rk^{2r1ฯต}.$$
We partition the interval $`[1,M_ฯต]`$ into consecutive intervals of length $`b`$ and denote the $`i^{\mathrm{th}}`$ such interval by $`B_i`$ and call it a block (note that the last block may be a partial block). Let $`C_i`$ be the $`2r`$ consecutive blocks $`B_{2r(i1)+1},B_{2r(i1)+2},\mathrm{},B_{2ri}`$, $`i=1,2,\mathrm{},\frac{M_ฯต}{2rb}`$. For each $`C_i`$, we randomly choose a row in $`A`$, say $`(s_1,s_2,\mathrm{},s_{2r})`$. We color the $`j^{\mathrm{th}}`$ block of $`C_i`$ by $`\gamma _{s_j}`$. By $`col(B_i)`$ we mean the coloring of the $`i^{\mathrm{th}}`$ block, $`1i\frac{M_ฯต}{b}`$, which is one of $`\gamma _0,\gamma _1,\mathrm{},\gamma _{r1}`$.
The following is immediate by construction.
Lemma 2.1
(i) For all $`1i2rb`$, $`P(col(B_i)=\gamma _c)=\frac{1}{r}`$ for each $`c=0,1,\mathrm{},r1`$.
(ii) For any $`i`$, $`P(col(B_i)=\gamma _c\mathrm{and}col(B_{i+1})=\gamma _d)=\frac{1}{r^2}`$ for any $`c`$ and $`d`$.
(iii) The colorings of blocks whose pairwise distances are at least $`2r`$ are mutually independent.
The approach we take, following Alon and Spencer , is to show that there exists a coloring such that for any monochromatic $`\frac{k}{2}`$-term ascending wave $`w_1,w_2,\mathrm{},w_{k/2}`$ we have $`w_{k/2}w_{k/21}O(k^{2r2ฯต/2})`$. The following definition and lemma, which are generalizations of those found in , will give us the desired result.
Definition An arithmetic progression $`x_1<x_2<\mathrm{}<x_t`$ is called a good progression if for each $`c\{0,1,\mathrm{},r1\}`$, there exists $`x_iB_j`$ such that $`col(B_j)=col(B_{j+1})=\gamma _c`$. An arithmetic progression that is not good is called a bad progression.
Lemma 2.2 For $`k,r2`$, let $`t=\frac{(4r2)(2r+1)}{\mathrm{log}(r^2/(r^21))}\mathrm{log}(k)+\frac{(2r+1)(\mathrm{log}(r)+1)}{\mathrm{log}(r^2/(r^21))}`$. For $`k`$ sufficiently large, the probability that there is a bad progression in a random coloring of $`[1,M_ฯต]`$ with difference greater than $`b`$ of $`t`$ terms is at most $`\frac{1}{2}`$.
Proof. Let $`x_1<x_2<\mathrm{}<x_t`$ be a bad progression with $`x_2x_1>b`$. Then no 2 elements belong to the same block. For each $`i`$, $`1i\frac{t}{2r+1}`$, let $`C_i`$ be the block in which $`x_{(2r+1)i}`$ resides, and let $`D_i`$ be the consecutive block. Then, the probability that the progression is bad is at most
$$p=\underset{j=1}{\overset{r}{}}P(\mathrm{}i[1,\frac{t}{2r+1}]:Col(C_i)=Col(D_i)=\gamma _j).$$
We have
$$\begin{array}{cc}p\hfill & rP(\mathrm{}i[1,\frac{t}{2r+1}]:Col(C_i)=Col(D_i)=\gamma _0)\hfill \\ & =r\left(\frac{r^21}{r^2}\right)^{\frac{t}{2r+1}}\hfill \\ & r\left(\frac{r^21}{r^2}\right)^{\frac{(4r2)}{\mathrm{log}(r^2/(r^21))}\mathrm{log}(k)+\frac{\mathrm{log}(r)+1}{\mathrm{log}(r^2/(r^21))}}\hfill \\ & \frac{2^1}{k^{4r2}}\hfill \end{array}$$
for $`k`$ sufficiently large.
Since the number of $`t`$-term arithmetic progressions in $`[1,M_ฯต]`$ is less than $`M_ฯต^2<k^{4r2}`$, the probability that there is a bad progression is less than
$$k^{4r2}\frac{2^1}{k^{4r2}}=\frac{1}{2},$$
thereby completing the proof. $`\mathrm{}`$
Lemma 2.3 Let $`C`$ be an $`r`$-coloring of $`[1,M_ฯต]`$ having no bad progression with difference greater than $`b`$ of $`t`$ terms ($`t`$ from Lemma 2.2). For any $`ฯต>0`$, for $`k`$ sufficiently large, any monochromatic $`\frac{k}{2}`$-term ascending wave $`w_1,w_2,\mathrm{},w_{k/2}`$ has $`w_{k/2}w_{k/21}bk^{1ฯต/2}=O(k^{2r2ฯต/2})`$.
Proof. At most $`4r4`$ consecutive blocks can have a specific color in all of them. (To achieve this say the color is $`0`$. The random coloring must have chosen row $`1`$ followed by row $`r+1`$, to have $`\gamma _0\gamma _0\gamma _1\gamma _1\mathrm{}\gamma _{r1}\gamma _{r1}\gamma _1\gamma _1\gamma _2\gamma _2\mathrm{}\gamma _0\gamma _0`$.) Since each block has a monochromatic ascending wave of length at most $`\frac{k}{10(4r4)}1`$, any $`4r4`$ consecutive blocks contribute less than $`\frac{k}{10}`$ terms to a monochromatic ascending wave. After that, the next difference must be more than $`b`$.
Let the monochromatic ascending wave be $`A=a_1,a_2,\mathrm{},a_{k/2}`$. Then, there exists $`i<\frac{k}{10}`$ such that $`a_{i+1}a_ib+1`$. Now let $`X=x_1,x_2,\mathrm{},x_t`$ be a $`t`$-term good progression with $`x_1=a_i`$ and $`d=x_2x_1=a_{i+1}a_ib+1`$.
Assume, without loss of generality, that the color of $`A`$ is $`0`$. Since $`X`$ is a good progression, there exists $`x_jB_{\mathrm{}}`$ with $`Col(B_{\mathrm{}})=Col(B_{\mathrm{}+1})=\gamma _0`$ for some block $`B_{\mathrm{}}`$. Since $`a_{i+j}x_j`$ as $`A`$ is an ascending wave, we see that $`a_{i+j}a_ijd+b+1`$. We conclude that $`a_{i+t}a_itd+b+1`$ so that, since $`A`$ is an ascending wave, $`a_{i+t+1}a_{i+t}d+\frac{b+1}{t}`$. Now, let $`X=x_1,x_2,\mathrm{},x_t`$ be the $`t`$-term good progression with $`x_1=a_{i+t}`$ and $`d^{}=x_2x_1=a_{i+t+1}a_{i+t}d+\frac{b+1}{t}(b+1)\left(1+\frac{1}{t}\right)`$. Repeating the above argument, we see that $`a_{i+2t}a_{i+t}td^{}+b+1`$ so that $`a_{i+2t}a_{i+2t1}d^{}+\frac{b+1}{t}(b+1)\left(1+\frac{2}{t}\right)`$. In general,
$$a_{i+st}a_{i+st1}(b+1)\left(1+\frac{s}{t}\right)$$
for $`s=1,2,\mathrm{}\frac{2k5t}{5t}`$. Thus, we have (with $`s=(k^{1ฯต/2}1)t\frac{2k5t}{5t}`$ for $`k`$ sufficiently large)
$$a_{k/2}a_{k/21}(b+1)\left(1+\frac{(k^{1ฯต/2}1)t}{t}\right)=(b+1)k^{1ฯต/2}.$$
$`\mathrm{}`$
Using Lemma 2.3, there exists an $`r`$-coloring of $`[1,M_ฯต]`$ such that any $`\frac{k}{2}`$-term monochromatic ascending wave has last difference at least $`(b+1)k^{1ฯต/2}`$. This implies that the last term of any $`k`$-term monochromatic ascending wave must be at least $`\frac{k}{2}+(b+1)k^{1ฯต/2}\frac{k}{2}>\frac{1}{2}(b+1)k^{2ฯต/2}.`$
We have
$$b+1N_{r1}\frac{1}{40^{2r3ฯต/2}(r1)^{2r3ฯต/2}}k^{2r3ฯต/2}N_{r1}\frac{k^{2r3ฯต/2}}{40^{2r1}r^{2r1}}.$$
Hence, for $`k`$ sufficiently large, the last term of any monochromatic $`k`$-term ascending wave must be greater than
$$N_{r1}\frac{1}{40^{2r1}r^{2r1}}k^{2r3ฯต/2}\frac{k^{2ฯต/2}}{2}=N_rk^{2r1ฯต}=M_ฯต.$$
Hence, we have an $`r`$-coloring of $`[1,M_ฯต]`$ with no $`k`$-term monochromatic ascending wave, for $`k`$ sufficiently large, thereby proving Theorem 2.
## 3. A Lower Bound for Three Colors
We believe that $`AW(k;r)=O(k^{2r1})`$, however, we have thus far been unable to prove this. The approach of Alon and Spencer , which is to show that there exists an $`r`$-coloring (under a random coloring scheme) such that every monochromatic $`\frac{3k}{4}`$-term ascending wave has $`d_{3k/41}>ck^{2r2}`$ does not work for an arbitrary number of colors with our generalization. However, for 3 colors, we can refine their argument to prove that $`AW(k;3)=O(k^5)`$.
Theorem 3
$$\frac{k^5}{2^{13}10^{39}}AW(k;3)\frac{k^5}{120}(1+o(1))$$
The upper bound comes from Theorem 1, hence we need only prove the lower bound. We use the same coloring scheme as in Section 2 and proceed with a series of lemmas.
Let $`w_1,w_2,\mathrm{},w_{k/2}`$ be a $`\frac{k}{2}`$-term ascending wave. From Lemmas 2.2 and 2.3, there exist (many) colorings such that for $`k`$ sufficiently large, $`w_{k/2}w_{k/21}>6b(=2rb)`$. Before proving Theorem 3, we introduce the following definition.
Definition We call a sequence $`x_1,x_2,\mathrm{},x_n`$ with $`x_2x_11`$ an almost ascending wave if, for $`2in1`$, we have $`d_i=x_{i+1}x_i`$ with $`d_id_{i1}1`$, with equality for at least one such $`i`$ and with the property that if $`d_i=d_{i1}1`$ and $`d_j=d_{j1}1`$ with $`j>i`$ there must exist $`s`$, $`i<s<j`$, such that $`d_sd_{s1}+1`$.
The upper bound of the following proposition is a slight refinement of a result of Alon and Spencer \[1, Lemma 1.7\].
Proposition 3.1 Denote by $`aw(n)`$ the number of ascending waves of length $`n`$ with first term given and $`d_{n1}<\frac{n}{10^{14}}`$. Analogously, let $`aaw(n)`$ be the number of almost ascending waves of length $`n`$ with first term given and $`d_{n1}\frac{n}{10^{14}}`$. Then, for all $`n`$ sufficiently large,
$$2^{\frac{n}{2}1}<aw(n)+aaw(n)2^{\frac{13n}{25}}\left(\frac{3}{2}\right)^{n/100}.$$
Proof. We start with the lower bound. We start by constructing a sequence where all of $`\frac{n}{2}1`$ slots contain 2 terms of a sequence. From a list of $`\frac{n}{2}1`$ slots, choose $`j`$, $`0j\frac{n}{2}1`$, of them. In these slots place the pair $`1,1`$. In the remaining slots put the pair $`0,0`$. We now have a sequence of length $`n2`$ or $`n3`$. If the length is $`n2`$, put a $`2`$ at the end; if the length is $`n3`$, put $`2,2`$ at the end. We now have, for each $`j`$ and each choice of $`j`$ slots, a distinct sequence of length $`n1`$. Denote one such sequence by $`s_1,s_2,\mathrm{},s_{n1}`$. Now, let $`d_1=1`$ and $`d_i=d_{i1}+s_{i1}`$ for $`i=2,3,\mathrm{},n`$. Since we have the first term of an almost ascending, or ascending, wave $`w_1,\mathrm{},w_n`$ given, such a wave is determined by its sequence of differences $`w_{i+1}w_i`$. Above, we have constructed a sequence $`\{d_i\}`$ of differences that adhere to the rules of an almost ascending, or ascending, wave. Hence, $`aw(n;r)+aaw(n;r)>_{j=0}^{\frac{n}{2}1}\left(\genfrac{}{}{0pt}{}{\frac{n}{2}1}{j}\right)=2^{\frac{n}{2}1}`$.
For the upper bound, we follow the proof of Alon and Spencer \[1, Lemma 1.7\], improving the bound enough to serve our purpose. Their lemma includes the term $`\left(\genfrac{}{}{0pt}{}{n+10^6n1}{n1}\right)`$ which we will work on to refine their upper bound on $`aw(n)+aaw(n)`$.
First, we have
$$\left(\genfrac{}{}{0pt}{}{n+10^6n1}{n1}\right)\left(\genfrac{}{}{0pt}{}{(1+10^5)n}{n}\right)$$
for $`n`$ sufficiently large.
Let $`q=(1+10^5)^1`$, $`m=\frac{n}{q}`$ and let $`H(x)=x\mathrm{log}(x)(1x)\mathrm{log}(1x)`$ for $`0x1`$ be the binary entropy function. Then we have<sup>2</sup><sup>2</sup>2Hereโs a quick derivation: For all $`n1`$, we have $`\sqrt{2\pi n}e^{1/(12n+1)}(n/e)^nn!\sqrt{2\pi n}e^{1/(12n)}(n/e)^n`$ (see ). Hence, $`\left(\genfrac{}{}{0pt}{}{m}{qm}\right)\frac{c}{\sqrt{m(1q)}}\left(q^q(1q)^{(1q)}\right)^m`$ for some positive $`c<e^2`$ (so that $`\frac{c}{\sqrt{m(1q)}}<1`$ for $`m`$ sufficiently large). Using the base 2 $`\mathrm{log}`$, this gives $`\left(\genfrac{}{}{0pt}{}{m}{qm}\right)2^{mH(q)}`$.
$$\left(\genfrac{}{}{0pt}{}{m}{qm}\right)2^{mH(q)}.$$
Applying this, we have
$$\begin{array}{cc}H(q)\hfill & =\frac{1}{(1+10^5)}\mathrm{log}((1+10^5))\left(\frac{10^5}{(1+10^5)}\right)\mathrm{log}\left(\frac{10^5}{(1+10^5)}\right)\hfill \end{array}$$
so that
$$\begin{array}{cc}mH(q)\hfill & =\left[\mathrm{log}(1+10^5)\frac{1}{10^5}\mathrm{log}\left(\frac{10^5}{1+10^5}\right)\right]n\hfill \\ & =\left[\frac{1}{10^5}\mathrm{log}\left(10^5(1+10^5)^{10^5+1}\right)\right]n\hfill \\ & \left[\frac{1}{10^5}\mathrm{log}(e(10^5+1))\right]n.\hfill \end{array}$$
We proceed by noting that
$$\left[\frac{\mathrm{log}(e(10^5+1))}{10^5}\right]n\left[\frac{1}{100}\mathrm{log}\left(\frac{3}{2}\right)\right]n.$$
Hence, $`2^{mH(q)}2^{\frac{n}{100}\mathrm{log}\left(\frac{3}{2}\right)}=\left(\frac{3}{2}\right)^{\frac{n}{100}}.`$ Now, using Alon and Spencerโs result \[1, Lemma 1.7\], the result follows. $`\mathrm{}`$
We are now in a position to prove the fundamental lemma of this section. In the proof we refer to the following definition.
Definition Let $`a_1,\mathrm{},a_n`$ be an ascending wave and let $`b^+`$. We call $`\frac{a_1}{b},\frac{a_2}{b},\mathrm{},\frac{a_n}{b}`$ the associated $`b`$-floor wave.
Lemma 3.2 Let
$$Q=\frac{k^5}{2^{13}10^{39}}.$$
The probability that in a random $`3`$-coloring of $`[1,Q]`$ there is a monochromatic ascending wave of length $`\frac{k}{4}`$ whose first difference is greater than $`6b(=2rb)`$ and whose last difference is smaller than $`\frac{kb}{410^{14}}=O(k^4)`$ is less than $`\frac{1}{2}`$ for $`k`$ sufficiently large.
Proof. Let $`A=a_1<a_2<\mathrm{}<a_{k/4}`$ be an ascending wave and let $`\frac{a_1}{b}<\frac{a_2}{b}<\mathrm{}<\frac{a_{k/4}}{b}`$ be the associated $`b`$-floor wave. Note that this $`b`$-floor wave is either an ascending wave or an almost ascending wave with last difference at most $`\frac{k/4}{10^{14}}`$. Hence, by Proposition 3.1, the number of such $`b`$-floor waves is at most, for $`k`$ sufficiently large,
$$k^22^{\frac{13k}{100}}\left(\frac{3}{2}\right)^{k/400}2^{\frac{14k}{100}}\left(\frac{3}{2}\right)^{k/400}$$
(we have less than $`k^2`$ choices for $`\frac{a_1}{b}`$).
Note that $`A`$ is monochromatic of color, say $`c`$, only if none of the blocks $`B_{\frac{a_i}{b}}`$, $`1i\frac{k}{4}`$, is colored by $`\gamma _c`$. Note that all of these blocks are at least $`6(=2r)`$ blocks from each other. We use Lemma 2.1 to give us that the probability that $`A`$ is monochromatic is no more than
$$3\left(\frac{2}{3}\right)^{k/4}.$$
Thus, the probability that in a random $`3`$-coloring of $`[1,Q]`$ we have a monochromatic ascending wave with last difference less than $`\frac{kb}{410^{14}}`$ is at most
$$32^{\frac{14k}{100}}\left(\frac{2}{3}\right)^{99k/400}.$$
We have $`3<\left(\frac{3}{2}\right)^{3k/400}`$ for $`k`$ sufficiently large, so that the above probability is less than
$$2^{\frac{14k}{100}}\left(\frac{2}{3}\right)^{24k/100}.$$
The above quantity is, in particular, less than $`1/2`$ for $`k`$ sufficiently large. $`\mathrm{}`$
To finish proving Theorem 3, we apply Corollary 2.4 and Lemma 3.2 to show that, for $`k`$ sufficiently large, there exists a $`3`$-coloring of $`[1,Q]`$ such that both of the following hold:
1) Any $`\frac{k}{2}`$-term monochromatic ascending wave has last difference greater than $`6b(=2rb)`$.
2) Any $`\frac{k}{4}`$-term monochromatic ascending wave with first difference greater than $`6b(=2rb)`$ has last difference greater than $`\frac{kb}{410^{14}}`$.
Hence, we conclude that there is a $`3`$-coloring of $`[1,Q]`$ such that any $`\frac{3k}{4}`$-term monochromatic ascending wave has last difference greater than $`\frac{kb}{410^{14}}`$, for $`k`$ sufficiently large. This implies that the last term of such a monochromatic ascending wave must be at least $`\frac{3k}{4}+\frac{kb}{410^{14}}\frac{k}{4}.`$
We have $`b=AW(\frac{k}{10(4r4)};r1)1`$ with $`r=3`$. By Alon and Spencerโs result , this gives us
$$b\frac{k^3}{10^{25}8^3}$$
for $`k`$ sufficiently large.
Hence, for $`k`$ sufficiently large, the last term of an ascending wave must be at least
$$\frac{3k}{4}+\frac{k^2}{4^210^{14}}\frac{k^3}{10^{25}8^3}>\frac{k^5}{2^{13}10^{39}}=Q.$$
Since we have the existence of a $`3`$-coloring of $`[1,Q]`$ with no monochromatic $`k`$-term ascending wave, this completes the proof of Theorem 3.
Remark From the lower bound given in Proposition 3.1, it is not possible to show that $`AW(k;r)O(k^{2r1})`$ for $`r4`$, which we believe is correct, by using the argument presented in Sections 2 and 3.
## References
N. Alon and J. Spencer, Ascending waves, J. Combin. Theory, Series A 52 (1989), 275-287.
T. Brown, P. Erdลs, and A. Freedman, Quasi-progressions and descending waves, J. Combin. Theory Series A 53 (1990), 81-95.
J. Conway and R. Guy, The Book of Numbers, Springer-Verlag, New York, p. 106, 1996.
B. Landman and A. Robertson, Ramsey Theory on the Integers, American Math. Society, Providence, RI, 317pp., 2003.
H. Robbins, A remark on Stirlingโs formula, American Math. Monthly 62 (1955), 26-29.
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# On (non-Hermitian) Lagrangeans in (particle) physics and their dynamical generation
## 1 Dynamical generation of Lagrangeans
### 1.1 The concept of dynamical generation
The concept and terminology of โdynamical generationโ occurred to us for the first time explicitly in the context of the (one-loop) โdynamical generationโ of the Quark-Level Linear Sigma Model by M.D. Scadron and R. Delbourgo .
A particularly important issue in the process of quantizing a theory given by some classical Lagrangean is the aspect of renormalization and renormalizability . The process of renormalization is typically performed โ after choosing some valid regularization scheme (See e.g. Ref. ) โ by adding to the classical Lagrangean divergent counterterms, which subtract divergencies which would otherwise show up in the unrenormalized effective action. Naively one might think that renormalization affects only terms belonging to the same order of perturbation theory in some coupling constant, while other parameters of the same Lagrangean do not interfere. The underlying philosophy would here be that in a quantum theory distinct parameters (e.g. masses, couplings) in a Lagrangean can be considered โ like in a classical Lagrangean โ to a great extent uncorrelated, as long as the Lagrangean is renormalizable. It appears that this philosophy seems to work quite well, when it is to renormalize logarithmic divergencies. That the situation is not so easy can be seen from the formalism needed to renormalize non-Abelian vector fields . In such theories the values of the coupling constants responsible for the self-interaction of three vector fields and of four vector fields are highly correlated due to the need to cancel appearing quadratic divergencies in the process of summing up diagrams of different loop order (in particular to achieve here the fundamental principle of gauge invariance). If this were not like that, their values could be chosen independently and therefore also renormalized independently. We see here a first example of โdynamicalโ generation or interrelation of two otherwise independent parameters in a Lagrangean due to the requirement of renormalizability, which affects here also the cancellation of quadratic divergencies. Furthermore we learn that โdynamical generationโ typically interrelates seemingly uncorrelated parameters of the Lagrangean and different loop orders <sup>1</sup><sup>1</sup>1Most probably the most outstanding example for dynamically generated theories are theories containing supersymmetry. This is reflected by the fact that supersymmetric theories typically contain a minimum of parameters, quadratic divergencies cancel exactly without extra renormalization (See e.g. Ref. ), and the renormalization of logarithmic divergencies at one-loop order yields simultaneously an automatic renormalization of all higher-loop orders. That observation led already to (non-conclusive) speculations about the question, whether all theories cancelling quadratic divergencies must be supersymmetric (See e.g. Refs. ). In certain situations some โ not necessarily supersymmetric โ theories may display even strong cancellations on the level of logarithmic divergencies. In such โbootstrappingโ theories physics is determined already at โtree-levelโ, as cancelling loop-contributions show up to be marginal.. Renormalizable theories with scalar fields only seem naively to have the priviledge, not to be affected by the problem faced by non-Abelian gauge theories, as the quadratic divergencies seem to be subtractable before entering the renormalization of logarithmic divergencies. Hence it seems naively, that โ as long as a Lagrangean with scalar fields only is in a classical sense considered to be renormalizable โ different parameters of the Lagrangean can be renormalized individually (up to constraints resulting from multiplicative renormalization). It is exactly this misbelief, which leads indeed to the triviality of scalar field theories like the text book $`\varphi ^4`$ theory or even to intimately related Abelian gauge theories like QED, if not โdynamically generatedโ. If instead the respective theories are โdynamically generatedโ one does find โ besides the trivial solution โ also non-trivial choices of the their parameter space, which survive the renormalization process without running into triviality. Interestingly in many cases such non-trivial solutions are found in the sector of the parameter space related to a PT-symmetric , yet not necessarily to a Hermitian non-trivial theory <sup>2</sup><sup>2</sup>2Before proceeding we want to deliver here also some warning about some common regularization schemes used which must not to be used in the context of โdynamical generationโ: Most important information about divergencies underlying a theory is contained in tadpole diagrams; hence any kind of artificial normal ordering or suppression of important surface terms will erase information needed to dynamical generate the theory and will lead therefore to wrong conclusions (See e.g. the discussion in Refs. ). As dimensional regularization erases or changes several important divergent diagrams like the massless tadpole (See e.g. Ref. ) or the quadratic divergence in the sunset/sunrise graph (See e.g. the dimensional regularization calculations performed in Refs. , or on p. 114 ff in Ref. ), it should not be used to dynamically generate a theory! According to our experience cutoff regularization โ if correctly used โ seems to yield always correct and most compact results compared to other regularization schemes.. In order to โdynamically generateโ a theory (e.g. like the supersymmetric Wess-Zumino model ) on the basis of some tentative classical Lagrangean we have to perform two steps: first we have to construct the terms in the effective action which are causing non-logarithmic divergencies (i.e. linear, quadratic, and higher divergencies) in all connected Feynman-diagrams, which can be constructed from the theory; then we have to relate and choose the parameters entering these terms of the effective action such, that all non-logarithmic divergencies cancel.<sup>3</sup><sup>3</sup>3One feels the need to remark that the very existence of a dynamically generated theory is not always guaranteed, as the procedure of dynamical generation is intimately related to renormalization and โ even more โ is strongly constraining the parameters of the effective action.
### 1.2 New method for the derivation of the effective action and its Lagrangean
A powerful method to construct the effective action has been known at least since the benchmarking work of S. Coleman & E. Weinberg and R. Jackiw . Unfortunately it is for our purposes not very convenient, as the determination of desired terms of the effective action responsible for leading singularities requires typically the simultaneous tedious evaluation of many other terms, which do not alter the discussion. This is why we want to propose here a different โ to our best knowledge โ new and more pragmatic approach yielding equivalent results compared to the formalism of S. Coleman, E. Weinberg, and R. Jackiw. Without loss of generality we want to explain our simple method here on the basis of some example, the generalization of which is quite straight forward.
Letโs start with the interaction part $`S_{int}=d^4zL_{int}(\stackrel{}{\varphi }(z),_z\stackrel{}{\varphi }(z))`$ of an action $`S`$ of $`N`$ interacting Klein-Gordon fields $`\varphi _1(z)`$, $`\mathrm{}`$, $`\varphi _N(z)`$. Then the interaction part of the effective action responsible for a process involving $`n`$ external legs is calculated by the connected ($`\mathrm{}_c`$) time-ordered vacuum expectation value of the Dyson-operator, where contractions are to be performed over all fields except $`n`$ fields (โexcept $`\varphi ^n`$โ), which remain to be contracted with creation or annihilation operators appearing in initial or final states, i.e.:
$`{\displaystyle \frac{i}{1!}}S_{eff}`$ $`=`$ $`0\left|T[\mathrm{exp}(iS_{int})]\right|0_c|_{\text{except}\varphi ^n}`$ (1)
$`=`$ $`{\displaystyle \frac{i}{1!}}0\left|T[S_{int}]\right|0_c|_{\text{except}\varphi ^n}+{\displaystyle \frac{i^2}{2!}}0\left|T[S_{int}S_{int}]\right|0_c|_{\text{except}\varphi ^n}+`$
$`+`$ $`{\displaystyle \frac{i^3}{3!}}0\left|T[S_{int}S_{int}S_{int}]\right|0_c|_{\text{except}\varphi ^n}+\mathrm{}.`$
The method is proved by making heavy use of the following identity (inserted between initial and final states $`|i`$ and $`f|`$, respectively) found e.g. on p. 44 in a well known book by C. Nash , i.e.:
$`f\left|T[\mathrm{exp}(iS_{int})]\right|i=f|\mathrm{exp}\left(\left[{\displaystyle \frac{1}{2}}\varphi ^2{\displaystyle \frac{\delta ^2}{\delta \varphi ^2}}\right]\right):\mathrm{exp}(iS_{int}):|i=`$ (2)
$`=`$ $`f\left|\mathrm{exp}\left(\left[{\displaystyle \frac{1}{2}}\varphi ^2{\displaystyle \frac{\delta ^2}{\delta \varphi ^2}}\right]\right)(:{\displaystyle \frac{i}{1!}}S_{int}:+:{\displaystyle \frac{i^2}{2!}}S_{int}S_{int}:+\mathrm{})\right|i`$
$`=`$ $`f|[{\displaystyle \frac{i}{1!}}\{:S_{int}:+{\displaystyle \frac{1}{1!}}\left[{\displaystyle \frac{1}{2}}\varphi ^2{\displaystyle \frac{\delta ^2}{\delta \varphi ^2}}\right]:S_{int}:+{\displaystyle \frac{1}{2!}}\left[{\displaystyle \frac{1}{2}}\varphi ^2{\displaystyle \frac{\delta ^2}{\delta \varphi ^2}}\right]^2:S_{int}:+\mathrm{}\}`$
$`+{\displaystyle \frac{i^2}{2!}}\{:S_{int}^2:+{\displaystyle \frac{1}{1!}}\left[{\displaystyle \frac{1}{2}}\varphi ^2{\displaystyle \frac{\delta ^2}{\delta \varphi ^2}}\right]:S_{int}^2:+{\displaystyle \frac{1}{2!}}\left[{\displaystyle \frac{1}{2}}\varphi ^2{\displaystyle \frac{\delta ^2}{\delta \varphi ^2}}\right]^2:S_{int}^2:+\mathrm{}\}`$
$`+{\displaystyle \frac{i^3}{3!}}\{:S_{int}^3:+{\displaystyle \frac{1}{1!}}\left[{\displaystyle \frac{1}{2}}\varphi ^2{\displaystyle \frac{\delta ^2}{\delta \varphi ^2}}\right]:S_{int}^3:+{\displaystyle \frac{1}{2!}}\left[{\displaystyle \frac{1}{2}}\varphi ^2{\displaystyle \frac{\delta ^2}{\delta \varphi ^2}}\right]^2:S_{int}^3:+\mathrm{}\}`$
$`+\mathrm{}]|i,`$
where we have defined for convenience the short-hand notation
$$\left[\frac{1}{2}\varphi ^2\frac{\delta ^2}{\delta \varphi ^2}\right]\frac{1}{2}\underset{i_1,i_2=1}{\overset{N}{}}d^4z_1d^4z_20\left|T[\varphi _{i_1}(z_1)\varphi _{i_2}(z_2)]\right|0\frac{\delta ^2}{\delta \varphi _{i_2}(z_2)\delta \varphi _{i_1}(z_1)}.$$
(3)
The identity (See e.g. p. 49 in Ref. ) and method is easily extended to Fermions, i.e. Grassmann fields $`\psi _1(z)`$, โฆ, $`\psi _N(z)`$, by replacing $`\left[\frac{1}{2}\varphi ^2\frac{\delta ^2}{\delta \varphi ^2}\right]`$ by
$$\left[\psi \overline{\psi }\frac{\delta ^2}{\delta \overline{\psi }\delta \psi }\right]\underset{i_1,i_2=1}{\overset{N}{}}d^4z_1d^4z_20\left|T[\psi _{i_1}(z_1)\overline{\psi }_{i_2}(z_2)]\right|0\frac{\delta ^2}{\delta \overline{\psi }_{i_2}(z_2)\delta \psi _{i_1}(z_1)}.$$
(4)
Convince yourself, that the method reproduces S. Colemanโs and E. Weinbergโs loop-expansion for a simple massless $`\varphi ^4`$-theory with $`S_{int}=d^4z(\frac{\lambda }{4!})\varphi ^4(z)`$ .<sup>4</sup><sup>4</sup>4We show here only the most important steps of the derivation: $`{\displaystyle \frac{i}{1!}}S_{eff}`$ $`=`$ $`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{i^n}{n!}}0\left|T[S_{int}^n]\right|0_c|_{\text{except}\varphi ^{2n}}={\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{i^n}{n!}}{\displaystyle \frac{1}{n!}}\left[{\displaystyle \frac{1}{2}}\varphi ^2{\displaystyle \frac{\delta ^2}{\delta \varphi ^2}}\right]^nS_{int}^n`$ (5) $`=`$ $`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{i^n}{n!}}{\displaystyle \frac{1}{n!}}{\displaystyle }d^4z_1\mathrm{}d^4z_n{\displaystyle \frac{n!(n1)!}{2}}({\displaystyle \frac{\lambda }{2!}})^n\varphi ^2(z_1)\mathrm{}\varphi ^2(z_n)\times `$ $`\times 0\left|T[\varphi (z_1)\varphi (z_2)]\right|00\left|T[\varphi (z_2)\varphi (z_3)]\right|0\mathrm{}0\left|T[\varphi (z_n)\varphi (z_1)]\right|0`$ $`=`$ $`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\left({\displaystyle \frac{\lambda }{2!}}\right)^n{\displaystyle \frac{1}{2n}}{\displaystyle }d^4z_1\mathrm{}d^4z_n\varphi ^2(z_1)\mathrm{}\varphi ^2(z_n)\times `$ $`\times {\displaystyle }{\displaystyle \frac{d^4p_{12}}{(2\pi )^4}}{\displaystyle \frac{d^4p_{23}}{(2\pi )^4}}\mathrm{}{\displaystyle \frac{d^4p_{n1}}{(2\pi )^4}}{\displaystyle \frac{e^{ip_{12}(z_1z_2)}e^{ip_{23}(z_2z_3)}\mathrm{}e^{ip_{n1}(z_nz_1)}}{(p_{12}^2+i\epsilon )(p_{23}^2+i\epsilon )\mathrm{}(p_{n1}^2+i\epsilon )}}`$ $`=`$ $`{\displaystyle d^4z\left(\underset{n=1}{\overset{\mathrm{}}{}}\left(\frac{\lambda }{2!}\right)^n\frac{1}{2n}\frac{d^4p}{(2\pi )^4}\left(\frac{\varphi ^2(0)}{p^2+i\epsilon }\right)^n+\text{non-local terms}\right)}.`$ Some of the resulting non-local terms are nicely discussed e.g. in Ref. .
## 2 Applications
### 2.1 A.S. Wightmanโs (non-)trivial and K. Symanzikโs precarious $`\varphi ^4`$ theory
In this section we want to shortly sketch the steps to dynamically generate the โScalar Wightman Theory in 4 Space-Time Dimensionsโ (See also Ref. ). As we will see below, the dynamical generation of this so-called $`\varphi ^4`$ theory yields โ besides the well known โtrivialโ solution โ the โprecariousโ non-trivial solution suggested by K. Symanzik being non-Hermitian and โ under certain circumstances also โ PT-symmetric .
To dynamically generate a $`\varphi ^N`$-theory upto $`N=4`$ we start from the following lowest order action containing just a three-point interaction:
$`S_{_{(0)}}`$ $`=`$ $`{\displaystyle d^4z\left\{\frac{1}{2}\left((\varphi _{_{(0)}}(z))^2m_{_{(0)}}^2\varphi _{_{(0)}}^2(z)\right)\frac{1}{3!}g_{_{(0)}}\varphi _{_{(0)}}^3(z)\right\}}`$ (6)
$`=`$ $`S_{_{(0)}}[(\varphi )^2]+S_{_{(0)}}[\varphi ^2]+S_{_{(0)}}[\varphi ^3].`$
In a first step we want to absorb by dynamical generation the finite one-loop correction to the $`\varphi ^3`$-coupling into a renormalization of the three-point coupling, i.e.:
$`{\displaystyle \frac{i}{1!}}S_{_{(1)}}[\varphi ^3]`$ $`=`$ $`{\displaystyle \frac{i}{1!}}0|T\left[S_{_{(0)}}[\varphi ^3]\right]|0_c|_{\text{except }\varphi _{_{\left(0\right)}}^3}`$ (7)
$`+`$ $`{\displaystyle \frac{i^3}{3!}}0|T\left[S_{_{(0)}}[\varphi ^3]S_{_{(0)}}[\varphi ^3]S_{_{(0)}}[\varphi ^3]\right]|0_c|_{\text{except }\varphi _{_{\left(0\right)}}^3}.`$
The next step is to dynamically generate on the basis of $`S_{_{(1)}}[\varphi ^3]`$ the term of the effective action quadratic in the fields $`\varphi _{_{(0)}}(z)`$ assuming the absence of quadratically divergent terms. <sup>5</sup><sup>5</sup>5I.e. we consider: $`{\displaystyle \frac{i}{1!}}\left(S_{_{(1)}}[(\varphi )^2]+S_{_{(1)}}[\varphi ^2]\right)={\displaystyle \frac{i}{1!}}0|T\left[S_{_{(0)}}[(\varphi )^2]\right]|0_c|_{\text{except }\varphi _{_{\left(0\right)}}^2}+`$ (8) $`+`$ $`{\displaystyle \frac{i}{1!}}0|T\left[S_{_{(0)}}[\varphi ^2]\right]|0_c|_{\text{except }\varphi _{_{\left(0\right)}}^2}+{\displaystyle \frac{i^2}{2!}}0|T\left[S_{_{(1)}}[\varphi ^3]S_{_{(1)}}[\varphi ^3]\right]|0_c|_{\text{except }\varphi _{_{\left(0\right)}}^2}.`$ The result of the previous steps is simple multiplicative coupling, wave function and mass renormalization, as we obtain as a whole (The omissions (โ$`\mathrm{}`$โ) denote here non-local terms not relevant for our present discussion.):
$`S_{_{(1)}}[(\varphi )^2]+S_{_{(1)}}[\varphi ^2]+S_{_{(1)}}[\varphi ^3]=`$ (9)
$`=`$ $`{\displaystyle d^4z\left(\frac{1}{2}\left((\varphi _{_{(1)}}(z))^2m_{_{(1)}}^2\varphi _{_{(1)}}^2(z)\right)\frac{1}{3!}g_{_{(1)}}\varphi _{_{(1)}}^3(z)\right)}+\mathrm{},`$
with
$`g_{_{(1)}}`$ $`=`$ $`\overline{g}_{_{(0)}}/\left(1{\displaystyle \frac{1}{32\pi ^2}}{\displaystyle \frac{\overline{g}_{_{(0)}}^2}{m_{_{(0)}}^2}}\right)^{3/2},\overline{g}_{_{(0)}}=g_{_{(0)}}\left(1+{\displaystyle \frac{1}{32\pi ^2}}{\displaystyle \frac{g_{_{(0)}}^2}{m_{_{(0)}}^2}}\right),`$
$`\varphi _{_{(1)}}^2(z)`$ $`=`$ $`\varphi _{_{(0)}}^2(z)\left(1{\displaystyle \frac{1}{32\pi ^2}}{\displaystyle \frac{\overline{g}_{_{(0)}}^2}{m_{_{(0)}}^2}}\right),`$
$`m_{_{(1)}}^2`$ $`=`$ $`m_{_{(0)}}^2\left(1+{\displaystyle \frac{i}{2}}{\displaystyle \frac{\overline{g}_{_{(0)}}^2}{m_{_{(0)}}^2}}{\displaystyle \frac{d^4p}{(2\pi )^4}\frac{1}{(p^2m_{_{(0)}}^2)^2}}\right)/\left(1{\displaystyle \frac{1}{32\pi ^2}}{\displaystyle \frac{\overline{g}_{_{(0)}}^2}{m_{_{(0)}}^2}}\right).`$
If we renormalize this result through a suitable mass counter term yielding a log.-divergent gap-equation promoted e.g. by M.D. Scadron , i.e. by applying
$$\frac{d^4p}{(2\pi )^4}\frac{1}{(p^2m_{_{(0)}}^2)^2}+\frac{i}{16\pi ^2},$$
(11)
then we have a bootstrapping situation for the mass, as there holds then $`m_{_{(1)}}^2=m_{_{(0)}}^2`$. Recall that the result has been obtained by assuming the absence, i.e. the cancellation of quadratically divergent terms in $`S_{_{(1)}}[(\varphi )^2]+S_{_{(1)}}[\varphi ^2]`$. In order to show now the absence of quadratically divergent terms for self-consistency reasons, we have first to dynamically generate on the basis of $`g_{_{(1)}}`$ and $`m_{_{(1)}}`$ the effective action for a four-point interaction of the field $`\varphi _{_{(0)}}(z)`$, and then test the cancellations of quadratic divergencies on the level of tadpoles and selfenergies. The effective action for a four-point interaction of the field $`\varphi _{_{(0)}}(z)`$ (expressed in terms of $`\varphi _{_{(1)}}(z)`$) is here dynamically generated for simplicity just up to order $`g_{_{(1)}}^4`$ assuming again the absence of quadratically divergent terms, i.e.:
$`{\displaystyle \frac{i}{1!}}S_{_{(1)}}[\varphi ^4]={\displaystyle \frac{i}{1!}}\left(S_{_{(1)}}^{tree}[\varphi ^4]+S_{_{(1)}}^{loop}[\varphi ^4]\right)=`$ (12)
$`=`$ $`{\displaystyle \frac{i^2}{2!}}0|T\left[S_{_{(1)}}[\varphi ^3]S_{_{(1)}}[\varphi ^3]\right]|0_c|_{\text{except }\varphi _{_{\left(1\right)}}^4}`$
$`+`$ $`{\displaystyle \frac{i^4}{4!}}0|T\left[S_{_{(1)}}[\varphi ^3]S_{_{(1)}}[\varphi ^3]S_{_{(1)}}[\varphi ^3]S_{_{(1)}}[\varphi ^3]\right]|0_c|_{\text{except }\varphi _{_{\left(1\right)}}^4}`$
$`=`$ $`{\displaystyle \frac{i^2}{2!}}{\displaystyle d^4z_1d^4z_2\left(\frac{1}{3!}g_{_{(1)}}\right)^2\mathrm{\hspace{0.17em}3}^2\varphi _{_{(1)}}^2(z_1)\varphi _{_{(1)}}^2(z_2)i\frac{d^4p}{(2\pi )^4}\frac{e^{ip(z_1z_2)}}{(p^2m_{_{(1)}}^2)}}`$
$`+`$ $`{\displaystyle \frac{i^4}{4!}}{\displaystyle }d^4z_1d^4z_2d^4z_3d^4z_4({\displaystyle \frac{1}{3!}}g_{_{(1)}})^43(3!)^4\varphi _{_{(1)}}(z_1)\varphi _{_{(1)}}(z_2)\varphi _{_{(1)}}(z_3)\varphi _{_{(1)}}(z_4)i^4\times `$
$`\times `$ $`{\displaystyle \frac{d^4p_{_{12}}}{(2\pi )^4}\frac{d^4p_{_{23}}}{(2\pi )^4}\frac{d^4p_{_{34}}}{(2\pi )^4}\frac{d^4p_{_{41}}}{(2\pi )^4}\frac{e^{ip_{_{12}}(z_1z_2)}e^{ip_{_{23}}(z_2z_3)}e^{ip_{_{34}}(z_3z_4)}e^{ip_{_{41}}(z_4z_1)}}{(p_{_{12}}^2m_{_{(1)}}^2)(p_{_{23}}^2m_{_{(1)}}^2)(p_{_{34}}^2m_{_{(1)}}^2)(p_{_{41}}^2m_{_{(1)}}^2)}}`$
$`=`$ $`{\displaystyle \frac{i}{1!}}{\displaystyle d^4z_1d^4z_2\left(\frac{1}{4!}\right)\mathrm{\hspace{0.17em}3}g_{_{(1)}}^2\varphi _{_{(1)}}^2(z_1)\varphi _{_{(1)}}^2(z_2)\frac{d^4p}{(2\pi )^4}\frac{e^{ip(z_1z_2)}}{(p^2m_{_{(1)}}^2)}}`$
$`+`$ $`{\displaystyle \frac{i}{1!}}{\displaystyle d^4z\left(\frac{1}{4!}\right)\mathrm{\hspace{0.17em}3}ig_{_{(1)}}^4\varphi _{_{(1)}}^4(z)\frac{d^4p}{(2\pi )^4}\frac{1}{(p^2m_{_{(1)}}^2)^4}}+\mathrm{}`$
$`=`$ $`{\displaystyle \frac{i}{1!}}{\displaystyle d^4z\left(\frac{1}{4!}\right)\left((\mathrm{\hspace{0.17em}3})\frac{g_{_{(1)}}^2}{m_{_{(1)}}^2}\varphi _{_{(1)}}^4(z)+\left(\frac{1}{32\pi ^2}\right)\frac{g_{_{(1)}}^4}{m_{_{(1)}}^4}\varphi _{_{(1)}}^4(z)\right)}+\mathrm{}`$
$`=`$ $`{\displaystyle \frac{i}{1!}}{\displaystyle d^4z\left(\left(\frac{1}{4!}\right)(\mathrm{\hspace{0.17em}3})\frac{g_{_{(1)}}^2}{m_{_{(1)}}^2}\varphi _{_{(1)}}^4(z)+\left(\frac{1}{4!}\lambda _{_{(1)}}\right)\varphi _{_{(1)}}^4(z)\right)}+\mathrm{}.`$
As a result of this consideration we have
$$S_{_{(1)}}=d^4z\left(\frac{1}{2}\left((\varphi _{_{(1)}}(z))^2m_{_{(1)}}^2\varphi _{_{(1)}}^2(z)\right)\frac{1}{3!}g_{_{(1)}}\varphi _{_{(1)}}^3(z)\frac{1}{4!}\lambda _{_{(1)}}\varphi _{_{(1)}}^4(z)\right)+\mathrm{},$$
(13)
with $`\lambda _{_{(1)}}=g_{_{(1)}}^4/(32\pi ^2m_{_{(1)}}^4)`$ and the replacements made in Eq. (LABEL:repleqns1). Letโs see now on the basis of this action, in how far quadratic divergencies cancel, as assumed in our approach from the beginning. Therefore we dynamically generate โ for convenience โ e.g. the effective action describing the sum of quadratically divergent tadpoles:
$`{\displaystyle \frac{i}{1!}}S_{_{(1)}}[\varphi ]`$ $`=`$ $`{\displaystyle \frac{i}{1!}}0|T\left[S_{_{(1)}}[\varphi ^3]\right]|0_c|_{\text{except }\varphi _{_{\left(1\right)}}}`$ (14)
$`+`$ $`{\displaystyle \frac{i^2}{2!}}\mathrm{\hspace{0.17em}2}!0|T\left[S_{_{(1)}}^{loop}[\varphi ^4]S_{_{(1)}}[\varphi ^3]\right]|0_c|_{\text{except }\varphi _{_{\left(1\right)}}}`$
$`=`$ $`{\displaystyle \frac{i}{1!}}{\displaystyle d^4z\left(\frac{1}{3!}g_{_{(1)}}\right)\mathrm{\hspace{0.33em}3}\varphi _{_{(1)}}(z)i\frac{d^4p}{(2\pi )^4}\frac{1}{(p^2m_{_{(1)}}^2)}}`$
$`+`$ $`{\displaystyle \frac{i}{1!}}{\displaystyle }d^4z({\displaystyle \frac{1}{3!}}g_{_{(1)}})(1)\lambda _{_{(1)}}\varphi _{_{(1)}}(z)\times `$
$`\times `$ $`{\displaystyle \frac{d^4p__1}{(2\pi )^4}\frac{d^4p__2}{(2\pi )^4}\frac{d^4p__3}{(2\pi )^4}\frac{(2\pi )^4\delta ^4(p__1+p__2+p__3)}{(p__1^2m_{_{(1)}}^2)(p__2^2m_{_{(1)}}^2)(p__3^2m_{_{(1)}}^2)}}.`$
To proceed further we extract shortly in the footnote the leading singularity structure of the occuring massive sunset/sunrise diagram, being particularly complicated due to the overlap of one quadratic divergence with three logarithmic divergences (See e.g. p. 78 ff in Ref. ).<sup>6</sup><sup>6</sup>6The safest and most compact discussion of the sunset/sunrise diagram is achieved in cutoff regularization, even though the full diagram in cutoff regularization has โ to our present knowledge โ never been calculated in a closed form. For a discussion of the finite part of the sunset/sunrise integral for non-zero external four-momentum on the basis of implicit renormalization see e.g. Ref. . The leading divergent parts of the sunset/sunrise diagram for zero external four-momentum and equal masses have been determined in cutoff regularization in Ref. to be: $`{\displaystyle ^\mathrm{\Lambda }}{\displaystyle \frac{d^4p_1}{(2\pi )^4}}{\displaystyle ^\mathrm{\Lambda }}{\displaystyle \frac{d^4p_2}{(2\pi )^4}}{\displaystyle ^\mathrm{\Lambda }}{\displaystyle \frac{d^4p_3}{(2\pi )^4}}{\displaystyle \frac{(2\pi )^4\delta ^4(p_1+p_2+p_3)}{(p_1^2m^2)(p_2^2m^2)(p_3^2m^2)}}=`$ (15) $`=`$ $`\left({\displaystyle \frac{1}{16\pi ^2}}\right)^2\left(2\mathrm{\Lambda }^2+{\displaystyle \frac{3}{2}}m^2\mathrm{ln}^2\left({\displaystyle \frac{\mathrm{\Lambda }^2}{m^2}}\right)3m^2\mathrm{ln}\left({\displaystyle \frac{\mathrm{\Lambda }^2}{m^2}}\right)+Cm^2\right)+O(\mathrm{\Lambda }^2),`$ while the integration constant $`C`$ was numerically estimated in Ref. to be approximately $`C4`$. After recalling $`^\mathrm{\Lambda }\frac{d^4p}{(2\pi )^4}\frac{1}{(p^2m^2)^2}=\frac{i}{16\pi ^2}\left(\mathrm{ln}\frac{\mathrm{\Lambda }^2+m^2}{m^2}\frac{\mathrm{\Lambda }^2}{\mathrm{\Lambda }^2+m^2}\right)`$ and $`^\mathrm{\Lambda }\frac{d^4p}{(2\pi )^4}\frac{1}{(p^2m^2)}=\frac{i}{16\pi ^2}m^2\left(\frac{\mathrm{\Lambda }^2}{m^2}\mathrm{ln}\frac{\mathrm{\Lambda }^2+m^2}{m^2}\right)`$ Eq. (15) is replaced for $`\mathrm{\Lambda }\mathrm{}`$ and in the local limit by $`I_{sunset/sunrise}={\displaystyle \frac{d^4p_1}{(2\pi )^4}\frac{d^4p_2}{(2\pi )^4}\frac{d^4p_3}{(2\pi )^4}\frac{(2\pi )^4\delta ^4(p_1+p_2+p_3)}{(p_1^2m^2)(p_2^2m^2)(p_3^2m^2)}}=`$ (16) $`=`$ $`2{\displaystyle \frac{i}{16\pi ^2}}{\displaystyle \frac{d^4p}{(2\pi )^4}\frac{1}{(p^2m^2)}}+{\displaystyle \frac{2}{3}}m^2\left({\displaystyle \frac{3}{2}}{\displaystyle \frac{d^4p}{(2\pi )^4}\frac{1}{(p^2m^2)^2}}+{\displaystyle \frac{i}{16\pi ^2}}\right)^2`$ $`+\left({\displaystyle \frac{1}{16\pi ^2}}\right)^2m^2\left({\displaystyle \frac{1}{6}}C\right)+\mathrm{}.`$ The last line displays the most divergent part of the massive sunset/sunrise diagram at zero external four-momentum in a regularization scheme independent manner. The application of a renormalization scheme yielding the โbootstrappingโ log.-divergent gap-equation Eq. (11) reduces the foregoing equation finally to
$$I_{sunset/sunrise}2\frac{i}{16\pi ^2}\frac{d^4p}{(2\pi )^4}\frac{1}{(p^2m^2)}\left(\frac{1}{16\pi ^2}\right)^2m^2(4+C)+\mathrm{}.$$
(17) The expression for the leading divergence of the sunset/sunrise graph is then to be inserted in Eq. (14) yielding the following result for the local limit of the effective action describing tadpoles:
$`S_{_{(1)}}[\varphi ]={\displaystyle }d^4z({\displaystyle \frac{1}{3!}}g_{_{(1)}})3i\varphi _{_{(1)}}(z)\times `$ (18)
$`\times `$ $`\left\{\left(1+{\displaystyle \frac{2}{3}}{\displaystyle \frac{1}{16\pi ^2}}\lambda _{_{(1)}}\right){\displaystyle \frac{d^4p}{(2\pi )^4}\frac{1}{(p^2m_{_{(1)}}^2)}}i\left({\displaystyle \frac{1}{16\pi ^2}}\right)^2m_{_{(1)}}^2{\displaystyle \frac{(4+C)}{3}}\lambda _{_{(1)}}\right\}`$
$`+`$ $`\mathrm{}.`$
Simple inspection of this expression yields that the quadratic divergencies cancel on one hand for the well known โtrivialโ solution $`g_{_{(1)}}=0`$. On the other hand the dynamically generated theory displays a non-trivial, precarious solution in the spirit of K. Symanzik for $`\lambda _{_{(1)}}=(3/2)\mathrm{\hspace{0.17em}16}\pi ^2=24\pi ^2`$ implying due to $`\lambda _{_{(1)}}=g_{_{(1)}}^4/(32\pi ^2m_{_{(1)}}^4)`$ four solutions for the three-point coupling constant $`g_{_{(1)}}`$, i.e. $`g_{_{(1)}}=\pm 4\pi \mathrm{\hspace{0.17em}3}^{1/4}m_{_{(1)}}`$ and $`g_{_{(1)}}=\pm i\mathrm{\hspace{0.17em}4}\pi \mathrm{\hspace{0.17em}3}^{1/4}m_{_{(1)}}`$. Furthermore we notice that for the probable case of $`C4`$ and non-vanishing mass $`m_{_{(1)}}`$ the non-trivial theory develops already at this stage a finite non-vanishing vacuum expectation value (See also the discussion in Ref. ). Finally we mention in view of self-consistency without listing the explicit proof that the obtained non-trivial values for $`\lambda _{_{(1)}}`$ and $`g_{_{(1)}}`$ lead also to a cancellation of quadratic divergencies on the level of the selfenergy, consistent with our starting assumption that quadratic divergencies cancel.
### 2.2 A non-Hermitian and โPT-symmetricโ theory of strong interactions
The purpose of this section is to demonstrate on the basis of experimental โevidenceโ that a dynamically generated theory of strong interactions based on mesons and quarks has to be non-Hermitian and close to PT-symmetric . Starting point for our considerations โ inspired somehow by Ref. โ is the sum of the interaction Lagrangean of weak interactions containing (anti)leptons denoted by $`\mathrm{}_{}(x)`$, $`\overline{\mathrm{}_+^c}(x)`$ and (anti)quarks denoted by $`q_{}(z)`$, $`\overline{q_+^c}(z)`$ and a Yukawa-like interaction Lagrangean describing the strong interaction between (anti)quarks and scalar ($`S(z)`$), pseusoscalar ($`P(z)`$), vector ($`V(z)`$), and axialvector ($`Y(z)`$) $`U(6)\times U(6)`$ meson field matrices in flavour space inspired by Ref. (See also ) (The undetermined signs $`s_s`$, $`s_p`$, $`s_v`$, $`s_y\{1,+1\}`$ are here irrelevant!):
$`L_{int}^{strong}(z)=`$
$`=`$ $`\sqrt{2}g\overline{q_+^c}(z)\left(s_sS(z)+s_piP(z)\gamma _5+{\displaystyle \frac{e^{i\alpha }}{2}}\left(s_v\overline{)}V(z)+s_y\overline{)}Y(z)\gamma _5\right)\right)q_{}(z),`$
with $`g=|g|\mathrm{exp}(i\alpha )`$ being the eventually complex strong interaction coupling constant, while contrary to Refs. we do not allow any further extra direct meson-meson interaction terms in the Lagrangean, as they shall be generated dynamically through quark-loops only <sup>7</sup><sup>7</sup>7This follows the same philosophy as in the previous section, where the $`\varphi ^4`$-interaction was dynamically generated starting out just from a $`\varphi ^3`$-theory. It is an interesting possibility to be considered in future, whether in a similar manner the whole non-Fermionic part of the Lagrangean of the standard model of particle physics can be dynamically generated on the basis of Yukawa-like interaction terms coupling of Bosons (gauge-bosons, Higgs-(pseudo)scalars, $`\mathrm{}`$) to Fermions, i.e. (anti)quarks and (anti)leptons.. The first step is now to study leptonic decays of pseudoscalar mesons to extract the pseudoscalar decay constants $`f_P`$. By dynamical generation we obtain for the relevant part of the effective action $`S_{eff}`$ in the local limit ($`M_q\text{diag}[m_u,m_c,m_t,m_d,m_s,m_b]`$, โ$`\text{tr}__F`$โ= flavour trace) <sup>8</sup><sup>8</sup>8We assumed here without loss of generality for traditional reasons a colour factor $`N_c`$, which can be absorbed by a redefinition of the strong coupling constant $`g`$.:
$`{\displaystyle \frac{i}{1!}}S_{eff}={\displaystyle \frac{i}{1!}}0|T[S]|0_c|_{\text{except }P\overline{\mathrm{}}\mathrm{}}+{\displaystyle \frac{i^2}{2!}}0|T[SS]|0_c|_{\text{except }P\overline{\mathrm{}}\mathrm{}}+\mathrm{}`$ (30)
$`=`$ $`{\displaystyle }d^4z(\mathrm{\hspace{0.17em}2}{\displaystyle \frac{G_F}{\sqrt{2}}})\sqrt{2}s_pe^{i\alpha }\times `$
$`\times `$ $`\text{tr}__F[\mathrm{\hspace{0.17em}4}iN_c|g|{\displaystyle }{\displaystyle \frac{d^4p}{(2\pi )^4}}{\displaystyle \frac{1}{p^2M_q^2}}{\displaystyle \frac{1}{2}}\{M_q,(_\mu P(z))\}{\displaystyle \frac{1}{p^2M_q^2}}`$
$`\times (\overline{\mathrm{}_+^c}(z)\gamma ^\mu {\displaystyle \frac{1}{2}}(1\gamma _5)\left(\begin{array}{cc}0_3& 0_3\\ 1_3& 0_3\end{array}\right)\mathrm{}_{}(z)\left[\left(\begin{array}{cc}0_3& V_{_{CKM}}\\ 0_3& 0_3\end{array}\right)\right]`$
$`+\overline{\mathrm{}_+^c}(z)\gamma ^\mu {\displaystyle \frac{1}{2}}(1\gamma _5)\left(\begin{array}{cc}0_3& 1_3\\ 0_3& 0_3\end{array}\right)\mathrm{}_{}(z)\left[\left(\begin{array}{cc}0_3& 0_3\\ \overline{V}_{_{CKM}}& 0_3\end{array}\right)\right]`$
$`+\overline{\mathrm{}_+^c}(x)\gamma ^\mu {\displaystyle \frac{1}{2}}(T_3(1\gamma _5)2Q_{\mathrm{}}\mathrm{sin}^2\theta _W)\mathrm{}_{}(z)\left[\mathrm{\hspace{0.17em}2}T_3\right])]+\mathrm{}.`$
Inspection yields for the decay constant $`f_{\eta _{q_1\overline{q}_2}}`$ of a pseudoscalar meson $`\eta _{q_1\overline{q}_2}`$
$$if_{\eta _{q_1\overline{q}_2}}\mathrm{\hspace{0.33em}4}N_c|g|\frac{d^4p}{(2\pi )^4}\frac{(m_{q_1}+m_{\overline{q}_2})/2}{(p^2m_{q_1}^2)(p^2m_{\overline{q}_2}^2)},$$
(31)
being in accordance with the log.-divergent gap-equation Eq. (11) promoted by M.D. Scadron<sup>9</sup><sup>9</sup>9The log.-divergent gap-equation should be understood here as a prescription to renormalize the original unrenormalized Lagrangean in replacing originally divergent quantitites by finite experimental numbers through a suitable choice of counter terms implying Eq. (11). It is interesting to note that the previous result yields the extremly important sum-rule (resulting from the properties of the underlying integral) $`(m_{q_1}m_{\overline{q}_3})f_{\eta _{q_1\overline{q}_3}}=(m_{q_1}m_{\overline{q}_2})f_{\eta _{q_1\overline{q}_2}}+(m_{q_2}m_{\overline{q}_3})f_{\eta _{q_2\overline{q}_3}}`$ yielding e.g. $`(m_um_s)f_{K^+}=(m_um_d)f_{\pi ^+}+(m_dm_s)f_{K^0}`$.. As we will need it in the following we have now to dynamically generate the effective action describing the coupling of a scalar and two pseudoscalar mesons. The result is listed in the footnote <sup>10</sup><sup>10</sup>10In the considered local limit we obtain: $`{\displaystyle \frac{i}{1!}}S_{eff}={\displaystyle \frac{i^3}{3!}}{\displaystyle \frac{3!}{1!\mathrm{\hspace{0.17em}2}!}}0|T[S_{int}^{Sq\overline{q}}S_{int}^{Pq\overline{q}}S_{int}^{Pq\overline{q}}]|0_c|_{\text{except }SPP}=`$ $`\stackrel{!}{=}`$ $`{\displaystyle d^4z\sqrt{2}g^2e^{i\alpha }s_s(4iN_c|g|)\frac{d^4p}{(2\pi )^4}}`$ $`\times `$ $`\{\text{tr}_F\left[S(z){\displaystyle \frac{1}{(p^2M_q^2)}}\{P^2(z),M_q\}{\displaystyle \frac{1}{(p^2M_q^2)}}\right]`$ $`+\text{tr}_F\left[[S(z),P(z)]{\displaystyle \frac{1}{(p^2M_q^2)}}[P(z),M_q]{\displaystyle \frac{1}{(p^2M_q^2)}}\right]`$ $`\text{tr}_F\left[\{S(z),M_q\}{\displaystyle \frac{1}{(p^2M_q^2)}}[P(z),M_q]{\displaystyle \frac{1}{(p^2M_q^2)}}[P(z),M_q]{\displaystyle \frac{1}{(p^2M_q^2)}}\right]\}+\mathrm{}`$ Recalling our โdefiningโ equation for pseudoscalar decay constants Eq. (31) the first two terms on the right-hand side of Eq. (LABEL:sppeq1) are equivalent to a $`SPP`$-interaction term, which one would obtain from a โshiftedโ quartic interaction Lagrangean with quartic coupling $`\lambda `$. The โshiftedโ Lagrangean is $`L(x)=\frac{\lambda }{2}\text{tr}_F[((S(x)+iP(x)D)(S(x)iP(x)D))^2]=\lambda \text{tr}_F[(S(x)+iP(x))(S(x)iP(x))(\{S(x),D\}+i[P(x),D])]+\mathrm{}`$. The quantity $`D`$ is the matrix (identified with decay constants of neutral pseudoscalar mesons) leading to spontaneous symmetry breaking according to the shift $`S(x)S(x)D`$ and inducing quark-masses according to the relation $`M_q=\sqrt{2}gs_sD`$. The last term on the right-hand side of Eq. (LABEL:sppeq1) involving only commutators $`[P(z),M_q]`$ is proportional to the square of quark-mass differences and therefore small in the sense of the nonrenormalization theorem by M. Ademollo and R. Gatto .. In order to arrive at our final conclusions we can use the previous result to study the experimentally measured transition formfactors $`f_\pm ^{K^+\pi ^0}(0)`$ characterizing the process $`K^+\pi ^0e^+\nu _e`$ at zero four-momentum transfer. First we dynamically generate the respective effective action in the local limit displaying here only the for us relevant terms representing $`W`$-emission graphs and an exchange of a scalar $`\kappa ^+`$-meson due to Partial Conservation of Vector Currents (PCVC) <sup>11</sup><sup>11</sup>11The exchange of a vector meson $`K^{}`$ is here disregarded, as it can contribute only marginally to the transition formfactor $`f_+^{K^+\pi ^0}(0)`$ at zero four-momentum transfer, i.e. at most of the order of the nonrenormalization theorem by M. Ademollo & R. Gatto , as the charge of the $`K^+`$ is solely generated due to photon-quark interactions.:
$`S_{eff}`$ $`=`$ $`{\displaystyle d^4z(ie^{2i\alpha })\left(\frac{G_F}{\sqrt{2}}\right)\overline{V}_{us}\overline{e_+^c}(z)\gamma _\mu (1\gamma _5)\nu _e(z)}`$ (33)
$`\times `$ $`{\displaystyle \frac{1}{\sqrt{2}}}\{\pi ^0(z)\left({\displaystyle \frac{2|g|f_{K^+}}{m_u+m_s}}^\mu K^+(z)\right)K^+(z)\left({\displaystyle \frac{2|g|f_{\eta _{u\overline{u}}}}{m_u+m_u}}^\mu \pi ^0(z)\right)`$
$`+4iN_c|g|^2(m_sm_u)^2K^+(z)(^\mu \pi ^0(z)){\displaystyle \frac{d^4p}{(2\pi )^4}\frac{1}{(p^2m_s^2)(p^2m_u^2)^2}}`$
$`+{\displaystyle \frac{\lambda }{g^2}}m_s{\displaystyle \frac{(m_sm_u)}{m_{\kappa ^+}^2}}{\displaystyle \frac{2|g|f_{K^+}}{m_u+m_s}}(\pi ^0(z)(^\mu K^+(z))+K^+(z)(^\mu \pi ^0(z)))\}`$
$`+`$ $`K^{}\text{-exchange}+\mathrm{}.`$
From this result we can read off the desired transition formfactors $`f_\pm ^{K^+\pi ^0}(0)`$ at zero four momentum transfer. Displaying only terms being of relevant order in the scale $`\delta =(m_s/m_u)10.44`$ according to the nonrenormalization theorem of M. Ademollo and R. Gatto we obtain $`f_+^{K^+\pi ^0}(0)=1+O(\delta ^2)`$ and
$`f_{}^{K^+\pi ^0}(0)O(\delta ^2)={\displaystyle \frac{\lambda }{g^2}}m_s{\displaystyle \frac{(m_sm_u)}{m_{\kappa ^+}^2}}{\displaystyle \frac{2|g|f_{K^+}}{m_u+m_s}}=`$ (34)
$`=`$ $`e^{2i\alpha }{\displaystyle \frac{\lambda }{g^2}}{\displaystyle \frac{2\delta (1+\delta )}{(2+\delta )}}{\displaystyle \frac{|m_u||f_{K^+}|}{m_{\kappa ^+}^2}}|g|\stackrel{!}{=}e^{2i\alpha }{\displaystyle \frac{2\delta (1+\delta )}{(2+\delta )}}{\displaystyle \frac{|m_u||f_{K^+}|}{m_{\kappa ^+}^2}}{\displaystyle \frac{4\pi }{\sqrt{3}}}.`$
On the right-hand side of this equation we used that M.D. Scadronโs log.-div. gap-equation Eq. (11) in combination with Eq. (31) implies $`|g|=2\pi /\sqrt{N_c}=2\pi /\sqrt{3}`$, and that there holds $`\lambda 2g^2`$ according to a one-loop dynamical generation . In using the experimental values $`|f_{K^+}|159\text{MeV}/\sqrt{2}`$ and $`m_{\kappa ^+}797`$ MeV we produce with the help of the last line of Eq. (34) the following Table 1:
$`f_{}^{K^+\pi ^0}(0)/e^{2i\alpha }`$ 0.050 0.102 0.125 0.148 0.200 0.225 $`\delta `$ for $`|m_u|=337`$ MeV 0.1098 0.2149 0.2591 0.3023 0.3965 0.4404 $`\delta `$ for $`|m_u|=\mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}3}`$ MeV 7.274 14.11 17.12 20.12 26.89 30.14
Inspection of the constituent quark mass case $`|m_u|337`$ MeV reveils that the experimentally measured negative transition formfactor ratio $`f_{}^{K^+\pi ^0}(0)/f_+^{K^+\pi ^0}(0)0.125\pm 0.023`$ can be only accomodated for $`e^{2i\alpha }<0`$, while for reasonable values of $`\delta `$ and $`m_\kappa `$ experiment seems to suggest the extreme PT-symmetric case $`\alpha \pi /2+0`$ yielding an imaginary PT-symmetric Yukawa-coupling $`g=i\mathrm{\hspace{0.17em}2}\pi /\sqrt{3}`$ and a Symanzik-like quartic coupling $`\lambda 2g^2=8\pi ^2/3<0`$, as obtained already earlier by the author, when โderivingโ the Lagrangean of the Quark-Level Linear Sigma Model from the Lagrangean of QCD (See also Ref. ). Finally it is interesting to consider our rough estimate for the experimentally yet badly determined mass of the $`\kappa (800)`$ scalar resonance (biased by $`K_0^{}(1430)`$) as a function of $`\delta `$. For $`|m_u|337`$ MeV and $`f_{}^{K^+\pi ^0}(0)=0.125`$ we obtain the following Table 2:
$`\delta `$ 0.10 0.20 0.26 0.30 0.40 0.44 0.50 $`m_\kappa `$ \[MeV\] 480.0 692.7 798.5 863.6 1013.1 1068.69 1148.7
Hence, semileptonic decays of pseudoscalar mesons can not only be used to reveil the seemingly non-Hermitian nature of a theory of strong interaction with a sizable amount of scalar confinement, they also may be used to โmeasureโ badly known experimental quantities like the masses of light and heavy scalar resonances.
This work is dedicated to M.D. Scadron and G. Rupp. It has been supported by the Fundaรงรฃo para a Ciรชncia e a Tecnologia (FCT) of the Ministรฉrio da Ciรชncia e da Tecnologia (e do Ensino Superior) of Portugal, under Grants no. PRAXIS XXI/BPD/20186/99, SFRH/BDP/9480/2002, POCTI/FNU/49555/2002, and POCTI/FP/FNU/50328/2003.
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# 1 Introduction
## 1 Introduction
The physics of nucleon resonances has always been a hot matter of discussion in hadronic spectroscopy. It is well known that conventional scheme of classification, systematics and interpretation of all hadronic states is based on the constituent quark model in its many different versions. In the simplest variant of the constituent quark model constructed in the early years all known lightest hadrons made of just three $`u`$, $`d`$ and $`s`$ quarks were classified in according to irreducible representations of the $`SU(3)_f`$ group where mesons were made out of $`q\overline{q}`$, while baryons were built from $`qqq`$. By this way the lowest $`q\overline{q}`$ meson configuration can be decomposed as $`\mathrm{๐}\overline{\mathrm{๐}}=\mathrm{๐}\mathrm{๐}`$, while the lowest $`qqq`$ baryon configuration can be decomposed $`\mathrm{๐}\mathrm{๐}\mathrm{๐}=\mathrm{๐๐}\mathrm{๐}\mathrm{๐}\mathrm{๐}`$ as well. It is a remarkable fact that octet ($`\pi ^0,\pi ^+,\pi ^{},K^0,K^+,\overline{K^0},K^{},\eta `$) of mesons and octet ($`p,n,\mathrm{\Sigma }^0,\mathrm{\Sigma }^+,\mathrm{\Sigma }^{},\mathrm{\Xi }^0,\mathrm{\Xi }^{},\mathrm{\Lambda }`$) and decuplet ($`\mathrm{\Delta }^0,\mathrm{\Delta }^{},\mathrm{\Delta }^+,\mathrm{\Delta }^{++},\mathrm{\Sigma }^0,\mathrm{\Sigma }^+,\mathrm{\Sigma }^{},\mathrm{\Xi }^0,\mathrm{\Xi }^{},\mathrm{\Omega }^{}`$) of baryons have experimentally been observed. The experimental discovery of the predicted $`\mathrm{\Omega }^{}`$ hyperon was a shining confirmation of $`SU(3)_f`$ symmetry and of its important role in the classification and systematics of hadronic states. The addition of the $`c`$, $`b`$, and $`t`$ quarks to the above three light quarks extends, in principle, the flavor symmetry to $`SU(6)_f`$. However, earlier the $`SU(6)`$ group has been considered as an approximate spin-flavor symmetry for the baryons made from just $`u`$, $`d`$ and $`s`$ quarks (see e.g. ). In that case the baryons are classified by the multiplets arising in the decomposition $`\mathrm{๐}\mathrm{๐}\mathrm{๐}=\mathrm{๐๐}\mathrm{๐๐}\mathrm{๐๐}\mathrm{๐๐}`$. Here, the basic states are $`u_{},u_{},d_{},d_{},s_{},s_{}`$ where $``$ and $``$ denote spin up and down. Next, the $`SU(6)`$ multiplets decompose into $`SU(3)_f`$ multiplets $`\mathrm{๐๐}={}_{}{}^{4}\mathrm{๐๐}{}_{}{}^{2}\mathrm{๐}`$, $`\mathrm{๐๐}={}_{}{}^{2}\mathrm{๐๐}{}_{}{}^{4}\mathrm{๐}{}_{}{}^{2}\mathrm{๐}{}_{}{}^{2}\mathrm{๐}`$, $`\mathrm{๐๐}={}_{}{}^{2}\mathrm{๐}{}_{}{}^{4}\mathrm{๐}`$, where the superscript $`(2S+1)`$ represents the total spin $`S`$ of the quarks for all particles in the given $`SU(3)_f`$ multiplet. So, the above mentioned baryonโs octet containing the nucleon, and the decuplet containing $`\mathrm{\Delta }`$(1232) belong to one and the same $`SU(6)`$ multiplet (56-plet ) which might be considered as a lowest state where the orbital angular momenta of all quark pairs are zero. Then the $`\mathrm{๐๐}`$ and $`\mathrm{๐๐}`$ could refer to the states with nonzero orbital angular momenta of quark pairs or something else. In this case the states with nonzero orbital angular momenta may be classified by $`SU(6)O(3)`$ supermultiplets. Even though the $`SU(6)`$ symmetry is broken by spin dependent interactions the $`SU(6)O(3)`$ basis was a suitable one for representing the baryon states. However, here the problem of so called โmissingโ quark-model states arises, it has no solution so far. Of course, in that case one could imagine some selection rules which are responsible for the fact that many baryons have not been observed. At the same time, many recent experiments have reported the observation narrow structures which cannot be explained by the standard quark-model assignments for baryons as $`qqq`$ states. This, first of all, concerns the number of narrow baryon structures observed in the missing mass $`M_X`$ and in the $`p\pi ^+`$ invariant mass distribution in the reaction $`ppp\pi ^+X`$, which cannot be associated with the standard $`qqq`$ quark configurations .<sup>2</sup><sup>2</sup>2I thank B. Tatischeff for drawing my attention to Ref. 3 caused this study. Why these baryon states are less massive and so narrow than predicted in quark models is an open question so far. Certainly, this raises a challenge to the theory. The other non-$`qqq`$ baryons have been observed as sharp structures in the $`nK^+`$ and $`pK^0`$ invariant mass distribution, now denoted as the $`\mathrm{\Theta }^+`$ baryons , and in the $`\mathrm{\Xi }^{}\pi ^\pm `$, $`\overline{\mathrm{\Xi }}^+\pi ^\pm `$ invariant mass spectra as well โ all of that was interpreted as candidates for pentaquarks ($`qqqq\overline{q}`$). Here, it should also be mentioned an experimentally well established evidence for many non-$`q\overline{q}`$ mesonโs states (see e.g. discussion in Ref. and references therein) declared often and often as exotics. In fact, exotic simply includes all hadrons which cannot be explained in the framework of the simple valence picture of $`q\overline{q}`$ for mesons or $`qqq`$ for baryons. Note, that even though the simple valence picture operates degrees of freedom like QCD fields the valence quarks are not identical to QCD fields.
The recent PDG Baryon Particle Listings contain 22 Nucleon and 22 Delta states which are the excited states of the nucleon observed in a large number of formation and production experiments. The conventional masses, widths and other discrete quantum numbers of the $`N`$ and $`\mathrm{\Delta }`$ resonances in the PDG Baryon Listings have largely been defined from partial-wave analysis of $`\pi N`$ scattering data. However, any specific constituent quark models even though with a clear set of dynamical ingredients provide quite an another assignments for the quantum numbers of baryons as $`qqq`$ states and predict a much richer spectrum of baryon states than has been found in partial-wave analysis of $`\pi N`$ scattering data. That is why, many attempts have been undertaken to search for โmissingโ or โhiddenโ quark-model states from partial-wave analysis in the production processes of other final states such as $`N2\pi `$, $`N\rho `$, $`N\eta `$, $`N\omega `$, $`\mathrm{\Lambda }K`$, $`\mathrm{\Sigma }K`$ (see e.g. ). Besides, there is a serious problem to translate the results of a constituent quark model calculations and predictions into the standard partial-wave analysis conventions accepted by PDG. That translation cannot be constructed in the framework of a given constituent quark model without some additional assumptions and conventions. At any rate, to make such translation in a more clear way we have to consider pentaquark states even for the usual, non-exotic $`N`$ and $`\mathrm{\Delta }`$ baryons. The consistent dynamical description of pentaquark states is very tedious and hard labour which has not been done yet, if not impossible at all. That sheer drudgery will unlikely done in the near future because the consistent dynamical description of more simple three-quark states though has not been performed so far.
Concerning the dynamical content of constituent quark models one could say that any specific quark model even with some โQCD-inspiredโ improvements is a phenomenological, non-relativistic potential model without a reliable ground in quantum field theory, in particular, in QCD. For instance, the usual used constituent quark mass parameters of about 300 MeV for the light $`u`$ and $`d`$ quarks cannot be derived in the framework of the underlying QCD. Of course quarks can move relativistically and this also means that theoretical consistency of a non-relativistic potential model is most likely doubtful. We would like also to point out that constituent quark models with the current quark masses about a few MeV have a serious problem with the value of the nucleon sigma term measured in low energy $`\pi N`$ scattering. To resolve the $`\sigma `$-term problem the strong Chiral Symmetry Breaking is needed, and the mechanism for that is not clear so far. Nevertheless, there is understanding that such mechanism might be found in the framework of non-perturbative quantum field theory.
Moreover, there is a hope that lattice computations in QCD with a powerful computers could help us to eliminate all imperfections of non-relativistic quark potential models, if lattice studies can help us at all. Time will show.
Recently a new, very simple and quite general theoretical concept concerning the structure of hadron spectra has been formulated which allowed to construct the global solution of the spectral problem in hadron spectroscopy; see and references therein where some of our previous studies were partially summarized. It has been claimed that existence of the extra dimensions in the spirit of Kaluza and Klein together with some novel dynamical ideas may provide new conceptual issues and quite new scheme of systematics for hadron states to build the unified picture for hadron spectra up. The main advantage of the developed theoretical concept is that all calculated numbers for masses and widths of hadron states do not depend on a special dynamical model but follow from fundamental hypothesis on existence of the extra dimensions with a compact internal extra space. One very important fact has been established in a reliable way: the size of the internal compact extra space determines the global characteristics of the hadron spectra while the masses of the decay products are the fundamental parameters of the compound systems being the elements of the global structure. What is remarkable that all new hadron states experimentally discovered last years have been observed just at the masses predicted in our approach, and those states appeared to be narrow as predicted too. A thorough analysis of many different experiments reported the observation of a new very narrow, manifestly exotic $`\mathrm{\Theta }^+`$ (Q=1, S=1) baryon, with the simplest quark assignment ($`uudd\overline{s}`$) decaying into $`nK^+`$ and $`pK_S^0`$, taken together allowed us to claim that many different $`\mathrm{\Theta }`$ states have been discovered and all of them were excellently incorporated in the unified picture for hadron spectra developed. This concerns the newly discovered $`\mathrm{\Xi }_{3/2}^{}`$ baryon with strangeness $`S=2`$, isospin I = $`\frac{3}{2}`$ and a quark content of ($`dsds\overline{u}`$) , now denoted as $`\mathrm{\Phi }^{}`$ by PDG, as well. We have also shown that a large amount of experimental data may excellently be incorporated in the systematics provided by the created unified picture for hadron spectra. In this article we apply our approach to show what place in the unified picture for hadron spectra the $`N`$ and $`\mathrm{\Delta }`$ baryons take up.
## 2 Understanding the $`N`$ and $`\mathrm{\Delta }`$ baryons <br>in the unified picture for hadron spectra
According to the general, theoretical concept we calculate the Kaluza-Klein tower of KK-excitations for the $`N\pi `$ system by the formula
$$M_n^{N\pi }=\sqrt{m_N^2+\frac{n^2}{R^2}}+\sqrt{m_\pi ^2+\frac{n^2}{R^2}},(n=1,2,3,\mathrm{}),$$
(1)
where $`R`$ is the same fundamental scale established before (see and references therein for the details), $`N=(p,n)`$, $`\pi =(\pi ^0,\pi ^\pm )`$, and the masses of proton, neutron and pions have been taken from PDG. The such built Kaluza-Klein tower is shown in Tables 1โ5. For simplicity we have considered one-dimensional compact internal extra space and only diagonal elements of the Kaluza-Klein tower have been presented. The experimental data extracted from PDG and Refs. are shown in Tables 1โ5 as well. The data from Refs. only are shown in separate Table 5. As seen all narrow low mass baryons shown in the Tables are in excellent agreement with the calculated KK excitations. The other narrow low mass baryonโs structures found in Refs. , if any, might be identified with non-diagonal elements of the KK tower for the $`N\pi `$ system. As a rule, non-diagonal elements of the KK towers are suppressed in reality; see however for the details. We would like to emphasize that the so-called universal internal toroidal extra spaces might be considered as a natural explanation of suppression for non-diagonal elements of the KK towers by conservation low of KK numbers. In other words, an experimental observation of hadronic states corresponding to non-diagonal elements of the KK towers could be considered as an evidence for existence of generic internal compact extra spaces.
Our conservative estimate for the widths of KK excitations looks like
$$\mathrm{\Gamma }_n\frac{\alpha }{2}\frac{n}{R}0.4n\text{MeV},$$
(2)
where $`n`$ is the number of KK excitation, and $`\alpha 0.02`$, $`R^1=41.48\text{MeV}`$ are known from our previous studies . This model independent estimate is universal for all of the KK towers, it does not depend on a composition of the compound systems living there. Certainly, some model dependent dynamics might modify this estimate. However, one property of estimate (2) is an absolute evidence for the higher the KK excitation is, the larger is the width of the corresponding compound system. This property has to be fulfilled in any model. The broad peaks in the hadron spectra are interpreted in our approach as an envelope of the narrow KK excitations predicted by the Kaluza-Klein scenario.
The most of the nucleon resonances presented in the PDG Baryon Particle Listings have been extracted from partial wave analysis performed by a few groups: the Carnegie-Mellon Berkeley (CMB) group , the Karlsruhe-Helsinki (KH) group and the Kent State University (KSU) group are the most famous among of them. It should also be noted the article where the CMB analysis has significantly been extended with account of a larger data set including the modern data at the moment. In fact, the formalism in Ref. is identical to CMB but the data base used is similar to KSU. We would also like to mention the old paper and review article . Each group has an own โcookeryโ in preparing the results of the analysis, the performed analyses differ from each other often significantly in the methods and the data sets used to extract the resonances, that is why, there exist sometimes large enough discrepancies between different groups.
Tables 6-7 compare the results of the KSU, KH and CMB analyses with the values of KK excitations for the $`p\pi `$ system. The best agreements with the KK excitations values are shown in Tables 6-7 by the bold-face numbers.
We have also calculated the Kaluza-Klein towers of KK-excitations for the $`N\rho `$, $`N\eta `$, $`N\omega `$, $`\mathrm{\Lambda }K`$, $`\mathrm{\Sigma }K`$ systems by the formulae similar to (1), and for the $`N2\pi `$ system by the formula
$$M_n^{N2\pi }=\sqrt{m_N^2+\frac{n^2}{R^2}}+\sqrt{m_{\pi ^1}^2+\frac{n^2}{R^2}}+\sqrt{m_{\pi ^2}^2+\frac{n^2}{R^2}},(n=1,2,3,\mathrm{}),$$
(3)
as this has been prescribed in ; here, as above, $`N=(p,n)`$, $`\pi ^1(\pi ^2)=(\pi ^0,\pi ^\pm )`$. These KK towers are shown in Tables 8-15. As above, we have restricted ourselves by the simplest case of one-dimensional compact internal extra space and only diagonal elements of the Kaluza-Klein towers have been presented. The arrangement of the known Nucleon and Delta baryons in the Kaluza-Klein towers is presented in Tables 8-15<sup>3</sup><sup>3</sup>3Even though the $`\mathrm{\Delta }`$ states in the $`N\eta `$, $`N\omega `$, and $`\mathrm{\Lambda }K`$ systems are forbidden by isospin we have listed the $`\mathrm{\Delta }`$ baryons in Tables 11-13 for convenience too. as well.
## 3 Discussion of comparison with experiment
As seen from the Tables all experimentally observed Nucleon and Delta baryons including narrow low mass baryons are excellently accommodated within them. We see only one empty cell in the $`N\pi `$ KK tower corresponding to the $`M_{26}^{N\pi }`$(1517) storey (see Tables 1-4), but very probably that this fact relates to our incomplete knowledge of modern experimental data base. Sometimes one and the same storey in the $`N\pi `$ KK tower is occupied by several baryons with approximately equal masses within errors.
It should be emphasized one remarkable fact: we did not find a place for the $`P_{33}`$(1232) baryon in Tables 3-4, the $`M_5^{N\pi }`$(1208-1212) and $`M_6^{N\pi }`$(1254-1257) storeys are not comfortable for this state. However, we found that the first $`M_1^{N2\pi }`$(1222-1232) storey in the $`N2\pi `$ KK tower (see Tables 8-9) is just the place for the $`P_{33}`$(1232) baryon. This means that the $`P_{33}`$(1232) baryon may have the true three-body origin; really, the symbol $`\mathrm{\Delta }`$ is quite appropriate one to correspond to this fact. The other possibility is to search for the $`P_{33}`$(1232) baryon among the non-diagonal elements in the $`N\pi `$ KK tower. For example,
$$M_{nm}^{p\pi ^\pm }(n=3,m=6)=1231.84\text{MeV},M_{nm}^{p\pi ^\pm }(n=7,m=5)=1232.17\text{MeV},$$
$$M_{nm}^{p\pi ^\pm }(n=11,m=3)=1230.33\text{MeV},M_{nm}^{n\pi ^\pm }(n=11,m=3)=1231.5\text{MeV},$$
$$M_{nm}^{n\pi ^0}(n=3,m=6)=1230.9\text{MeV},M_{nm}^{n\pi ^0}(n=7,m=5)=1230.87\text{MeV},$$
where
$$M_{nm}^{N\pi }=\sqrt{m_N^2+\frac{n^2}{R^2}}+\sqrt{m_\pi ^2+\frac{m^2}{R^2}},(n,m=1,2,3,\mathrm{}).$$
(4)
Tables 8-9 contain the arrangement of the other Delta baryons in the $`N2\pi `$ KK tower too.
Obviously, that is noteworthy fact, which we would like to point out here, concerning the $`P_{11}`$ resonance at 1462$`\pm `$10 MeV extracted in Ref. . This resonance just occupies the $`M_{10}^{N\pi }`$-storey in the $`N\pi `$ KK tower; see Table 1. Note, that the quark-model calculations for the mass of this resonance give 1405 MeV and 1383 MeV that is in strong disagreement.
The new $`P_{31}`$ resonance at 1744$`\pm `$36 MeV found in is also incorporated in our approach; see $`M_{15}^{N\pi }`$-storey in the $`N\pi `$ KK tower in Table 2. The quark-model calculations for the mass of the $`P_{31}`$ resonance give 1875 MeV and 1906 MeV which are also in strong disagreement. The new $`F_{35}`$ resonance at 1752$`\pm `$32 MeV found in may occupy the same $`M_{15}^{N\pi }`$-storey in the $`N\pi `$ KK tower.
The third $`D_{13}`$ resonance at 1804$`\pm `$55 MeV found in has also an own place in the $`N\pi `$ KK tower and especially in the $`N\rho `$ and $`N\eta `$ KK towers; see Tables 10-11. Here, one of the quark-model predictions 1809 MeV for the masses of the $`D_{13}`$ resonances is in good agreement .
The second $`P_{13}`$ resonance at 1879$`\pm `$17 MeV, the third $`P_{11}`$ resonance at 1885$`\pm `$30 MeV, and the second $`F_{15}`$ resonance at 1903$`\pm `$87 MeV found in live on one and the same $`M_{17}^{N\pi }`$-storey in the $`N\pi `$ KK tower; see Table 1. These resonances have also comfortable places in the $`N\rho `$, $`N\eta `$, $`N\omega `$, $`\mathrm{\Lambda }K`$ and $`\mathrm{\Sigma }K`$ KK towers; see Tables 10-15.
Probably the $`M_{18}^{N\pi }`$-storey in the $`N\pi `$ KK tower is not so good place for the $`S_{31}`$ resonance at 1920$`\pm `$24 MeV, and for the $`S_{11}`$ resonance at 1928$`\pm `$59 MeV found in . However, the $`M_{12}^{\mathrm{\Lambda }K}`$-storey in the $`\mathrm{\Lambda }K`$ KK tower (Table 13) is quite suitable for these resonances.
We would like to especially emphasize that the lowest $`D_{35}`$ resonance at 1956$`\pm `$22 MeV, the third $`P_{33}`$ resonance at 2014$`\pm `$16 MeV, the second $`D_{33}`$ resonance at 2057$`\pm `$110 MeV, the lowest $`F_{17}`$ resonance at 2086$`\pm `$28 MeV, and the high-mass $`D_{35}`$ resonance at 2171$`\pm `$18 MeV found in are excellently incorporated in our approach; see Tables 1-4,10-15. The masses for all of these resonances do not correspond to the quark-model calculations .
It seems the $`M_{21}^{N\pi }`$-storey in the $`N\pi `$ KK tower is not so comfortable for the lowest $`G_{17}`$ resonance at 2127$`\pm `$9 MeV found in . Very probably that this resonance lives together with the N$`P_{11}`$ and with the $`F_{17}`$ resonance at 2086$`\pm `$28 MeV on the same storey, as itโs clear from Tables 1-4,10-13. In addition, the same $`G_{17}`$ resonance at 2168$`\pm `$18 MeV found in excellently corresponds to the $`M_{21}^{N\pi }`$-storey in the $`N\pi `$ KK tower, and the $`S_{31}`$ resonance at 1802$`\pm `$87 MeV extracted in is excellently incorporated in our approach too; see Tables 1-4,10-15. .
We would like to mention too the $`S_{11}`$(1535) resonance at 1542$`\pm `$3 MeV (in accordance with $`M_5^{N\eta }`$; see detail discussion in ), the $`P_{11}`$(1440) resonance at 1479$`\pm `$80 MeV (in accordance with $`M_{10}^{N\pi }`$), the $`P_{11}`$(1710) resonance at 1699$`\pm `$65 MeV (in accordance with $`M_{14}^{N\pi }`$), the $`P_{11}`$(2100) resonance at 2084$`\pm `$93 MeV (in accordance with $`M_{20}^{N\pi }`$), the $`D_{13}`$(1520) resonance at 1518$`\pm `$3 MeV (in accordance with $`M_{11}^{N\pi }`$), the $`G_{17}`$(2190) resonance at 2168$`\pm `$18 MeV (in accordance with $`M_{21}^{N\pi }`$), the $`S_{31}`$(1620) resonance at 1617$`\pm `$15 MeV (in accordance with $`M_{13}^{N\pi }`$), the $`S_{31}`$(1900) resonance at 1802$`\pm `$87 MeV (in accordance with $`M_{16}^{N\pi }`$), the $`P_{31}`$(1750) resonance at 1721$`\pm `$61 MeV (in accordance with $`M_{15}^{N\pi }`$), and the $`F_{35}`$(1905) resonance at 1873$`\pm `$77 MeV (in accordance with $`M_{17}^{N\pi }`$) found all in Ref. which are also excellently incorporated in the unified picture for hadron spectra.
## 4 Summary and conclusion
This work should be considered as a continuation of our previous studies concerning the structure of hadron spectra. We have established that the recent PDG Baryon Particle Listings of Nucleon and Delta states, including some evidence for new states, have excellently incorporated in the theoretical concept developed earlier . In particular, it was shown that new resonances found in Ref. , including the $`P_{31}`$ state at 1744$`\pm `$36, the $`F_{35}`$ state at 1752$`\pm `$32, and $`P_{13}`$ state at 1879$`\pm `$17 which did not predicted in the quark-model calculations, have excellently accommodated in the corresponding KK towers. Moreover, the recently reported narrow low mass baryons , which cannot, in principle, be explained in conventional quark-models, have found own comfortable places in the corresponding KK towers.
Our predictions concerning the masses of hadron states are model independent, they are related with the fundamental hypothesis on existence of the extra dimensions with a compact internal extra space only. In general, each storey in the KK towers is degenerated, i.e. it may contain several flats for hadron states with the different quantum numbers but with approximately equal masses. In addition, the hadron states with the same quantum numbers may have different masses depending on what KK tower they live in, or in other words depending on what decay mode the hadron states have been observed in. This difference in the masses might be measured in the experiments with a high mass resolution. We have already discussed this non-trivial fact in analysing the SELEX measurements; see details in .
It should be especially emphasized that KK excitation corresponding to a certain storey in the given KK tower may have exactly the same quantum numbers which have been extracted from partial wave analysis because the definite KK tower corresponds to the definite decay channel what the given partial wave analysis has been done for.
As mentioned above all KK excitations are very narrow, they have the widths about a few MeV; see Eq. (2). The broad peaks in hadron spectra may appear in our approach as an envelope of the narrow KK excitations predicted by the Kaluza-Klein scenario. We have an idea that an availability of non-diagonal elements in the KK towers might play the crucial role in understanding the broad peaks in hadron spectra. This idea has to be explored in the nearest future.
It is well known that the pole positions extracted form partial wave analyses have the least model dependence compared to other parameters such as widths, (in)elasticity, couplings and so on. Our predictions for the masses of KK excitations are strong, that is wy we have performed at the moment only the comparison of the calculated mass values for the KK excitations with the masses of the $`N`$ and $`\mathrm{\Delta }`$ baryons determined from the partial wave analyses.
In conclusion, we would like once again to claim that all well established $`N`$ and $`\mathrm{\Delta }`$ baryons are excellently incorporated in the created unified picture for hadron spectra. No doubt, new experiments with a higher mass resolution and sensitivity would be very helpful to obtain new, more accurate and more full data of high quality. In that case it would be possible to make a reavaluation of all known data to refit the total baryon spectrum and to remove the existing discrepancies. We expect that such new experiments will be set up in the near future for this goal.
Table 1. Kaluza-Klein tower of KK excitations for the $`p\pi `$ system and $`N`$ baryons.
| n | $`M_n^{p\pi ^\pm }`$MeV | $`M_n^{p\pi ^0}`$MeV | N BARYONS (PDG) \[St:{$``$}\] | $`M_{exp}^{N\pi }`$ MeV |
| --- | --- | --- | --- | --- |
| 1 | 1084.79 | 1080.39 | | 1080.6, 1086.5 MAMI |
| 2 | 1104.30 | 1100.37 | | 1104, 1113 SPES4 |
| 3 | 1133.48 | 1130.08 | | 1130.4(8), 1133.7 SPES3 |
| 4 | 1169.65 | 1166.72 | | 1166.9, 1170.4 SPES3 |
| 5 | 1210.92 | 1208.38 | | 1198, 1202 SPES4 |
| 6 | 1256.07 | 1253.85 | | 1251.1 SPES3 |
| 7 | 1304.35 | 1302.38 | | 1313, 1322 SPES4 |
| 8 | 1355.23 | 1353.48 | | 1347.2 MAMI |
| 9 | 1408.38 | 1406.80 | | 1394 SPES4 |
| 10 | 1463.54 | 1462.10 | N$`P_{11}`$ (1430-1470) \[4$``$\] | 1462$`\pm `$10 Manley 92 |
| 11 | 1520.50 | 1519.18 | N$`D_{13}`$ (1515-1530) \[4$``$\] | 1519$`\pm `$ 4 Hoehler 79 |
| 12 | 1579.11 | 1577.89 | N$`S_{11}`$ (1520-1555) \[4$``$\] | 1550$`\pm `$40 Cutkosky 80 |
| 13 | 1639.22 | 1638.09 | N$`S_{11}`$ (1640-1680) \[4$``$\] | 1650$`\pm `$30 Cutkosky 80 |
| | | | N$`D_{15}`$ (1670-1685) \[4$``$\] | 1679$`\pm `$8 Hoehler 79 |
| | | | N$`F_{15}`$ (1675-1690) \[4$``$\] | 1680$`\pm `$10 Cutkosky 80 |
| 14 | 1700.73 | 1699.67 | N$`D_{13}`$ (1650-1750) \[3$``$\] | 1675$`\pm `$25 Cutkosky 80 |
| | | | N$`P_{11}`$ (1680-1740) \[3$``$\] | 1717$`\pm `$28 Manley 92 |
| | | | N$`P_{13}`$ (1650-1750) \[4$``$\] | 1710$`\pm `$20 Hoehler 79 |
| 15 | 1763.52 | 1762.53 | $`P_{11}`$ | 1766$`\pm `$34 Batinic 95 |
| 16 | 1827.50 | 1826.57 | $`D_{13}`$ | 1804$`\pm `$55 Manley 92 |
| | | | N$`P_{13}`$ ($`1900`$) \[2$``$\] | 1879$`\pm `$17 Manley 92 |
| 17 | 1892.59 | 1891.71 | $`P_{11}`$ | 1885$`\pm `$30 Manley 92 |
| | | | $`F_{15}`$ | 1903$`\pm `$87 Manley 92 |
| 18 | 1958.70 | 1957.87 | N$`F_{17}`$ ($`1990`$) \[2$``$\] | 1970$`\pm `$50 Cutkosky 80 |
| | | | $`S_{11}`$ | 1928$`\pm `$59 Manley 92 |
| 19 | 2025.77 | 2024.98 | N$`F_{15}`$ ($`2000`$) \[2$``$\] | 2025 Ayed 76 |
| | | | N$`D_{13}`$ ($`2080`$) \[2$``$\] | 2081$`\pm `$20 Hoehler 79 |
| 20 | 2093.73 | 2092.98 | $`F_{17}`$ | 2086$`\pm `$28 Manley 92 |
| | | | N$`S_{11}`$ ($`2090`$) \[1$``$\] | 2180$`\pm `$80 Cutkosky 80 |
| | | | N$`P_{11}`$ ($`2100`$) \[1$``$\] | 2050$`\pm `$20 Hoehler 79 |
| 21 | 2162.52 | 2161.81 | N$`G_{17}`$ (2100-2200) \[4$``$\] | 2168$`\pm `$18 Vrana 00 |
| | | | $`G_{17}`$ | 2127$`\pm `$9 Manley 92 |
| | | | $`G_{17}`$ | 2140$`\pm `$40 Hendry 78 |
| | | | N$`D_{15}`$ ($`2200`$) \[2$``$\] | 2228$`\pm `$30 Hoehler 79 |
| 22 | 2232.08 | 2231.39 | N$`H_{19}`$ (2180-2310) \[4$``$\] | 2230$`\pm `$80 Cutkosky 80 |
| | | | N$`G_{19}`$ (2170-2310) \[4$``$\] | 2250$`\pm `$80 Cutkosky 80 |
| 23 | 2302.36 | 2301.70 | $`H_{19}`$ | 2300$`\pm `$100 Hendry 78 |
| 24 | 2373.30 | 2372.68 | | |
| 25 | 2444.88 | 2444.27 | | |
| 26 | 2517.03 | 2516.45 | | |
| 27 | 2589.73 | 2589.17 | N$`I_{1,11}`$ (2550-2750) \[3$``$\] | 2577$`\pm `$50 Hoehler 79 |
| 28 | 2662.94 | 2662.40 | N$`K_{1,13}`$ ($`2700`$) \[2$``$\] | 2612$`\pm `$45 Hoehler 79 |
| 29 | 2736.63 | 2736.10 | $`I_{1,11}`$ | 2700$`\pm `$100 Hendry 78 |
| 30 | 2810.76 | 2810.25 | | |
Table 2. Kaluza-Klein tower of KK excitations for the $`n\pi `$ system and $`N`$ baryons.
| n | $`M_n^{n\pi ^0}`$MeV | $`M_n^{n\pi ^\pm }`$MeV | N BARYONS (PDG) \[St:{$``$}\] | $`M_{exp}^{N\pi }`$ MeV |
| --- | --- | --- | --- | --- |
| 1 | 1081.69 | 1086.08 | | 1080.6, 1086.5 MAMI |
| 2 | 1101.66 | 1105.59 | | 1104, 1113 SPES4 |
| 3 | 1131.36 | 1134.76 | | 1130.4(8), 1133.7 SPES3 |
| 4 | 1168.00 | 1170.92 | | 1166.9, 1170.4 SPES3 |
| 5 | 1209.64 | 1212.18 | | 1198, 1202 SPES4 |
| 6 | 1255.10 | 1257.32 | | 1251.1 SPES3 |
| 7 | 1303.62 | 1305.58 | | 1313, 1322 SPES4 |
| 8 | 1354.70 | 1356.45 | | 1347.2 MAMI |
| 9 | 1408.00 | 1409.59 | | 1394 SPES4 |
| 10 | 1463.28 | 1464.72 | N$`P_{11}`$ (1430-1470) \[4$``$\] | 1462$`\pm `$10 Manley 92 |
| 11 | 1520.34 | 1521.67 | N$`D_{13}`$ (1515-1530) \[4$``$\] | 1519$`\pm `$4 Hoehler 79 |
| 12 | 1579.03 | 1580.25 | N$`S_{11}`$ (1520-1555) \[4$``$\] | 1550$`\pm `$40 Cutkosky 80 |
| 13 | 1639.21 | 1640.34 | N$`S_{11}`$ (1640-1680) \[4$``$\] | 1650$`\pm `$30 Cutkosky 80 |
| | | | N$`D_{15}`$ (1670-1685) \[4$``$\] | 1679$`\pm `$8 Hoehler 79 |
| | | | N$`F_{15}`$ (1675-1690) \[4$``$\] | 1680$`\pm `$10 Cutkosky 80 |
| 14 | 1700.77 | 1701.83 | N$`D_{13}`$ (1650-1750) \[3$``$\] | 1675$`\pm `$25 Cutkosky 80 |
| | | | N$`P_{11}`$ (1680-1740) \[3$``$\] | 1717$`\pm `$28 Manley 92 |
| | | | N$`P_{13}`$ (1650-1750) \[4$``$\] | 1710$`\pm `$20 Hoehler 79 |
| 15 | 1763.61 | 1764.60 | $`P_{11}`$ | 1766$`\pm `$34 Batinic 95 |
| 16 | 1827.62 | 1828.56 | $`D_{13}`$ | 1804$`\pm `$55 Manley 92 |
| | | | N$`P_{13}`$ ($`1900`$) \[2$``$\] | 1879$`\pm `$17 Manley 92 |
| 17 | 1892.74 | 1893.62 | $`P_{11}`$ | 1885$`\pm `$30 Manley 92 |
| | | | $`F_{15}`$ | 1903$`\pm `$87 Manley 92 |
| 18 | 1958.88 | 1959.71 | N$`F_{17}`$ ($`1990`$) \[2$``$\] | 1970$`\pm `$50 Cutkosky 80 |
| | | | $`S_{11}`$ | 1928$`\pm `$59 Manley 92 |
| 19 | 2025.98 | 2026.76 | N$`F_{15}`$ ($`2000`$) \[2$``$\] | 2025 Ayed 76 |
| | | | N$`D_{13}`$ ($`2080`$) \[2$``$\] | 2081$`\pm `$20 Hoehler 79 |
| 20 | 2093.95 | 2094.70 | $`F_{17}`$ | 2086$`\pm `$28 Manley 92 |
| | | | N$`S_{11}`$ ($`2090`$) \[1$``$\] | 2180$`\pm `$80 Cutkosky 80 |
| | | | N$`P_{11}`$ ($`2100`$) \[1$``$\] | 2050$`\pm `$20 Hoehler 79 |
| 21 | 2162.75 | 2163.47 | N$`G_{17}`$ (2100-2200) \[4$``$\] | 2168$`\pm `$18 Vrana 00 |
| | | | $`G_{17}`$ | 2127$`\pm `$9 Manley 92 |
| | | | $`G_{17}`$ | 2140$`\pm `$40 Hendry 78 |
| | | | N$`D_{15}`$ ($`2200`$) \[2$``$\] | 2228$`\pm `$30 Hoehler 79 |
| 22 | 2232.32 | 2233.01 | N$`H_{19}`$ (2180-2310) \[4$``$\] | 2230$`\pm `$80 Cutkosky 80 |
| | | | N$`G_{19}`$ (2170-2310) \[4$``$\] | 2250$`\pm `$80 Cutkosky 80 |
| 23 | 2302.61 | 2303.26 | $`H_{19}`$ | 2300$`\pm `$100 Hendry 78 |
| 24 | 2373.56 | 2374.19 | | |
| 25 | 2445.14 | 2445.74 | | |
| 26 | 2517.30 | 2517.88 | | |
| 27 | 2590.00 | 2590.56 | N$`I_{1,11}`$ (2550-2750) \[3$``$\] | 2577$`\pm `$50 Hoehler 79 |
| 28 | 2663.21 | 2663.75 | N$`K_{1,13}`$ ($`2700`$) \[2$``$\] | 2612$`\pm `$45 Hoehler 79 |
| 29 | 2736.90 | 2737.42 | $`I_{1,11}`$ | 2700$`\pm `$100 Hendry 78 |
| 30 | 2811.03 | 2811.54 | | |
Table 3. Kaluza-Klein tower of KK excitations for the $`p\pi `$ system and $`\mathrm{\Delta }`$ baryons.
| n | $`M_n^{p\pi ^\pm }`$MeV | $`M_n^{p\pi ^0}`$MeV | $`\mathrm{\Delta }`$ BARYONS (PDG) \[St:{$``$}\] | $`M_{exp}^{N\pi }`$ MeV |
| --- | --- | --- | --- | --- |
| 1 | 1084.79 | 1080.39 | | |
| 2 | 1104.30 | 1100.37 | | |
| 3 | 1133.48 | 1130.08 | | |
| 4 | 1169.65 | 1166.72 | | |
| 5 | 1210.92 | 1208.38 | | |
| 6 | 1256.07 | 1253.85 | | |
| 7 | 1304.35 | 1302.38 | | |
| 8 | 1355.23 | 1353.48 | | |
| 9 | 1408.38 | 1406.80 | | |
| 10 | 1463.54 | 1462.10 | | |
| 11 | 1520.50 | 1519.18 | | |
| 12 | 1579.11 | 1577.89 | $`\mathrm{\Delta }[1600]P_{33}`$ (1550-1700) \[3$``$\] | 1600$`\pm `$50 Cutkosky 80 |
| 13 | 1639.22 | 1638.09 | $`\mathrm{\Delta }[1620]S_{31}`$ (1615-1675) \[4$``$\] | 1620$`\pm `$20 Cutkosky 80 |
| 14 | 1700.73 | 1699.67 | $`\mathrm{\Delta }[1700]D_{33}`$ (1670-1770) \[4$``$\] | 1710$`\pm `$30 Cutkosky 80 |
| 15 | 1763.52 | 1762.53 | $`\mathrm{\Delta }[1750]P_{31}`$ ($`1750`$) \[1$``$\] | 1744$`\pm `$36 Manley 92 |
| | | | $`F_{35}`$ | 1752$`\pm `$32 Manley 92 |
| 16 | 1827.50 | 1826.57 | $`S_{31}`$ | 1802$`\pm `$87 Vrana 00 |
| | | | $`\mathrm{\Delta }[1900]S_{31}`$ (1850-1950) \[2$``$\] | 1890$`\pm `$50 Cutkosky 80 |
| 17 | 1892.59 | 1891.71 | $`\mathrm{\Delta }[1905]F_{35}`$ (1870-1920) \[4$``$\] | 1881$`\pm `$18 Manley 92 |
| | | | $`\mathrm{\Delta }[1910]P_{31}`$ (1870-1920) \[4$``$\] | 1888$`\pm `$20 Hoehler 79 |
| | | | $`S_{31}`$ | 1920$`\pm `$24 Manley 92 |
| | | | $`\mathrm{\Delta }[1920]P_{33}`$ (1900-1970) \[3$``$\] | 1920$`\pm `$80 Cutkosky 80 |
| 18 | 1958.70 | 1957.87 | $`\mathrm{\Delta }[1930]D_{35}`$ (1920-1970) \[3$``$\] | 1956$`\pm `$22 Manley 92 |
| | | | $`\mathrm{\Delta }[1940]D_{33}`$ ($`1940`$) \[1$``$\] | 1940$`\pm `$100 Cutkosky 80 |
| | | | $`\mathrm{\Delta }[1950]F_{37}`$ (1940-1960) \[4$``$\] | 1950$`\pm `$15 Cutkosky 80 |
| 19 | 2025.77 | 2024.98 | $`\mathrm{\Delta }[2000]F_{35}`$ ($`2000`$) \[2$``$\] | 1752$`\pm `$32 Manley 92 |
| | | | $`P_{33}`$ | 2014$`\pm `$16 Manley 92 |
| 20 | 2093.73 | 2092.98 | $`D_{33}`$ | 2057$`\pm `$110 Manley 92 |
| 21 | 2162.52 | 2161.81 | $`\mathrm{\Delta }[2150]S_{31}`$ ($`2150`$) \[1$``$\] | 2150$`\pm `$100 Cutkosky 80 |
| | | | $`D_{35}`$ | 2171$`\pm `$18 Manley 92 |
| 22 | 2232.08 | 2231.39 | $`\mathrm{\Delta }[2200]G_{37}`$ ($`2200`$) \[1$``$\] | 2215$`\pm `$60 Hoehler 79 |
| 23 | 2302.36 | 2301.70 | $`\mathrm{\Delta }[2300]H_{39}`$ ($`2300`$) \[2$``$\] | 2217$`\pm `$80 Hoehler 79 |
| 24 | 2373.30 | 2372.68 | $`\mathrm{\Delta }[2350]D_{35}`$ ($`2350`$) \[2$``$\] | 2305$`\pm `$26 Hoehler 79 |
| | | | $`\mathrm{\Delta }[2390]F_{37}`$ ($`2390`$) \[1$``$\] | 2350$`\pm `$100 Cutkosky 80 |
| | | | $`H_{39}`$ | 2450$`\pm `$100 Hendry 78 |
| 25 | 2444.88 | 2444.27 | $`\mathrm{\Delta }[2400]G_{39}`$ ($`2400`$) \[2$``$\] | 2468$`\pm 50`$ Hoehler 79 |
| | | | $`\mathrm{\Delta }[2420]H_{3,11}`$ (2300-2500) \[4$``$\] | 2416$`\pm 17`$ Hoehler 79 |
| 26 | 2517.03 | 2516.45 | | |
| 27 | 2589.73 | 2589.17 | | |
| 28 | 2662.94 | 2662.40 | $`I_{3,13}`$ | 2650$`\pm `$100 Hendry 78 |
| 29 | 2736.63 | 2736.10 | $`\mathrm{\Delta }[2750]I_{3,13}`$ ($`2750`$) \[2$``$\] | 2794$`\pm `$80 Hoehler 79 |
| 30 | 2810.76 | 2810.25 | $`\mathrm{\Delta }[2950]K_{3,15}`$ ($`2950`$) \[2$``$\] | 2850$`\pm `$100 Hendry 78 |
Table 4. Kaluza-Klein tower of KK excitations for the $`n\pi `$ system and $`\mathrm{\Delta }`$ baryons.
| n | $`M_n^{n\pi ^0}`$MeV | $`M_n^{n\pi ^\pm }`$MeV | $`\mathrm{\Delta }`$ BARYONS (PDG) \[St:{$``$}\] | $`M_{exp}^{N\pi }`$ MeV |
| --- | --- | --- | --- | --- |
| 1 | 1081.69 | 1086.08 | | |
| 2 | 1101.66 | 1105.59 | | |
| 3 | 1131.36 | 1134.76 | | |
| 4 | 1168.00 | 1170.92 | | |
| 5 | 1209.64 | 1212.18 | | |
| 6 | 1255.10 | 1257.32 | | |
| 7 | 1303.62 | 1305.58 | | |
| 8 | 1354.70 | 1356.45 | | |
| 9 | 1408.00 | 1409.59 | | |
| 10 | 1463.28 | 1464.72 | | |
| 11 | 1520.34 | 1521.67 | | |
| 12 | 1579.03 | 1580.25 | $`\mathrm{\Delta }[1600]P_{33}`$ (1550-1700) \[3$``$\] | 1600$`\pm `$50 Cutkosky 80 |
| 13 | 1639.21 | 1640.34 | $`\mathrm{\Delta }[1620]S_{31}`$ (1615-1675) \[4$``$\] | 1620$`\pm `$20 Cutkosky 80 |
| 14 | 1700.77 | 1701.83 | $`\mathrm{\Delta }[1700]D_{33}`$ (1670-1770) \[4$``$\] | 1710$`\pm `$30 Cutkosky 80 |
| 15 | 1763.61 | 1764.60 | $`\mathrm{\Delta }[1750]P_{31}`$ ($`1750`$) \[1$``$\] | 1744$`\pm `$36 Manley 92 |
| | | | $`F_{35}`$ | 1752$`\pm `$32 Manley 92 |
| 16 | 1827.62 | 1828.56 | $`S_{31}`$ | 1802$`\pm `$87 Vrana 00 |
| | | | $`\mathrm{\Delta }[1900]S_{31}`$ (1850-1950) \[2$``$\] | 1890$`\pm `$50 Cutkosky 80 |
| 17 | 1892.74 | 1893.62 | $`\mathrm{\Delta }[1905]F_{35}`$ (1870-1920) \[4$``$\] | 1881$`\pm `$18 Manley 92 |
| | | | $`\mathrm{\Delta }[1910]P_{31}`$ (1870-1920) \[4$``$\] | 1888$`\pm `$20 Hoehler 79 |
| | | | $`S_{31}`$ | 1920$`\pm `$24 Manley 92 |
| | | | $`\mathrm{\Delta }[1920]P_{33}`$ (1900-1970) \[3$``$\] | 1920$`\pm `$80 Cutkosky 80 |
| 18 | 1958.88 | 1959.71 | $`\mathrm{\Delta }[1930]D_{35}`$ (1920-1970) \[3$``$\] | 1956$`\pm `$22 Manley 92 |
| | | | $`\mathrm{\Delta }[1940]D_{33}`$ ($`1940`$) \[1$``$\] | 1940$`\pm `$100 Cutkosky 80 |
| | | | $`\mathrm{\Delta }[1950]F_{37}`$ (1940-1960) \[4$``$\] | 1950$`\pm `$15 Cutkosky 80 |
| 19 | 2025.98 | 2026.76 | $`\mathrm{\Delta }[2000]F_{35}`$ ($`2000`$) \[2$``$\] | 1752$`\pm `$32 Manley 92 |
| | | | $`P_{33}`$ | 2014$`\pm `$16 Manley 92 |
| 20 | 2093.95 | 2094.70 | $`D_{33}`$ | 2057$`\pm `$110 Manley 92 |
| 21 | 2162.75 | 2163.47 | $`\mathrm{\Delta }[2150]S_{31}`$ ($`2150`$) \[1$``$\] | 2150$`\pm `$100 Cutkosky 80 |
| | | | $`D_{35}`$ | 2171$`\pm `$18 Manley 92 |
| 22 | 2232.32 | 2233.01 | $`\mathrm{\Delta }[2200]G_{37}`$ ($`2200`$) \[1$``$\] | 2215$`\pm `$60 Hoehler 79 |
| 23 | 2302.61 | 2303.26 | $`\mathrm{\Delta }[2300]H_{39}`$ ($`2300`$) \[2$``$\] | 2217$`\pm `$80 Hoehler 79 |
| 24 | 2373.56 | 2374.19 | $`\mathrm{\Delta }[2350]D_{35}`$ ($`2350`$) \[2$``$\] | 2305$`\pm `$26 Hoehler 79 |
| | | | $`\mathrm{\Delta }[2390]F_{37}`$ ($`2390`$) \[1$``$\] | 2350$`\pm `$100 Cutkosky 80 |
| | | | $`H_{39}`$ | 2450$`\pm `$100 Hendry 78 |
| 25 | 2445.14 | 2445.74 | $`\mathrm{\Delta }[2400]G_{39}`$ ($`2400`$) \[2$``$\] | 2468$`\pm 50`$ Hoehler 79 |
| | | | $`\mathrm{\Delta }[2420]H_{3,11}`$ (2300-2500) \[4$``$\] | 2416$`\pm 17`$ Hoehler 79 |
| 26 | 2517.30 | 2517.88 | | |
| 27 | 2590.00 | 2590.56 | | |
| 28 | 2663.21 | 2663.75 | $`I_{3,13}`$ | 2650$`\pm `$100 Hendry 78 |
| 29 | 2736.90 | 2737.42 | $`\mathrm{\Delta }[2750]I_{3,13}`$ ($`2750`$) \[2$``$\] | 2794$`\pm `$80 Hoehler 79 |
| 30 | 2811.03 | 2811.54 | $`\mathrm{\Delta }[2950]K_{3,15}`$ ($`2950`$) \[2$``$\] | 2850$`\pm `$100 Hendry 78 |
Table 5. KK excitations for the $`N\pi `$ system and narrow exotic low mass baryons .
| n | $`M_n^{p\pi ^\pm }`$MeV | $`M_n^{p\pi ^0}`$MeV | $`M_n^{n\pi ^0}`$MeV | $`M_n^{n\pi ^\pm }`$MeV | $`M_{exp}^{N\pi }`$ MeV |
| --- | --- | --- | --- | --- | --- |
| 1 | 1084.79 | 1080.39 | 1081.69 | 1086.08 | 1080.6, 1086.5 MAMI |
| 2 | 1104.30 | 1100.37 | 1101.66 | 1105.59 | 1104, 1113 SPES4 |
| 3 | 1133.48 | 1130.08 | 1131.36 | 1134.76 | 1130.4(8), 1133.7 SPES3 |
| 4 | 1169.65 | 1166.72 | 1168.00 | 1170.92 | 1166.9, 1170.4 SPES3 |
| 5 | 1210.92 | 1208.38 | 1209.64 | 1212.18 | 1198, 1202 SPES4 |
| 6 | 1256.07 | 1253.85 | 1255.10 | 1257.32 | 1251.1 SPES3 |
| 7 | 1304.35 | 1302.38 | 1303.62 | 1305.58 | 1313, 1322 SPES4 |
| 8 | 1355.23 | 1353.48 | 1354.70 | 1356.45 | 1347.2 MAMI |
| 9 | 1408.38 | 1406.80 | 1408.00 | 1409.59 | 1394 SPES4 |
| 10 | 1463.54 | 1462.10 | 1463.28 | 1464.72 | 1477 SPES4 |
| 11 | 1520.50 | 1519.18 | 1520.34 | 1521.67 | 1517 SPES3 SPES4 |
| 12 | 1579.11 | 1577.89 | 1579.03 | 1580.25 | 1577 SPES4 |
| 13 | 1639.22 | 1638.09 | 1639.21 | 1640.34 | 1639 SPES4 |
| 14 | 1700.73 | 1699.67 | 1700.77 | 1701.83 | |
| 15 | 1763.52 | 1762.53 | 1763.61 | 1764.60 | |
| 16 | 1827.50 | 1826.57 | 1827.62 | 1828.56 | |
| 17 | 1892.59 | 1891.71 | 1892.74 | 1893.62 | |
| 18 | 1958.70 | 1957.87 | 1958.88 | 1959.71 | |
| 19 | 2025.77 | 2024.98 | 2025.98 | 2026.76 | |
| 20 | 2093.73 | 2092.98 | 2093.95 | 2094.70 | |
| 21 | 2162.52 | 2161.81 | 2162.75 | 2163.47 | |
| 22 | 2232.08 | 2231.39 | 2232.32 | 2233.01 | |
| 23 | 2302.36 | 2301.70 | 2302.61 | 2303.26 | |
| 24 | 2373.30 | 2372.68 | 2373.56 | 2374.19 | |
| 25 | 2444.88 | 2444.27 | 2445.14 | 2445.74 | |
| 26 | 2517.03 | 2516.45 | 2517.30 | 2517.88 | |
| 27 | 2589.73 | 2589.17 | 2590.00 | 2590.56 | |
| 28 | 2662.94 | 2662.40 | 2663.21 | 2663.75 | |
| 29 | 2736.63 | 2736.10 | 2736.90 | 2737.42 | |
| 30 | 2810.76 | 2810.25 | 2811.03 | 2811.54 | |
Table 6. Kaluza-Klein tower of KK excitations for the $`p\pi `$ system compared with the masses of the $`N`$ baryons determined from the KSU , KH and CMB analyses.
| n | $`M_n^{p\pi ^\pm }`$MeV | $`M_n^{p\pi ^0}`$MeV | N(PDG) \[St:{$``$}\] | KSU | KH | CMB |
| --- | --- | --- | --- | --- | --- | --- |
| 1 | 1084.79 | 1080.39 | | | | |
| 2 | 1104.30 | 1100.37 | | | | |
| 3 | 1133.48 | 1130.08 | | | | |
| 4 | 1169.65 | 1166.72 | | | | |
| 5 | 1210.92 | 1208.38 | | | | |
| 6 | 1256.07 | 1253.85 | | | | |
| 7 | 1304.35 | 1302.38 | | | | |
| 8 | 1355.23 | 1353.48 | | | | |
| 9 | 1408.38 | 1406.80 | | | | |
| 10 | 1463.54 | 1462.10 | N$`P_{11}`$ \[4$``$\] | 1462$`\pm `$10 | 1410$`\pm `$12 | 1440$`\pm `$30 |
| 11 | 1520.50 | 1519.18 | N$`D_{13}`$ \[4$``$\] | 1524$`\pm `$4 | 1519$`\pm `$4 | 1525$`\pm `$10 |
| 12 | 1579.11 | 1577.89 | N$`S_{11}`$ \[4$``$\] | 1534$`\pm `$7 | 1526$`\pm `$7 | 1550$`\pm `$40 |
| 13 | 1639.22 | 1638.09 | N$`S_{11}`$ \[4$``$\] | 1659$`\pm `$9 | 1670$`\pm `$8 | 1650$`\pm `$30 |
| | | | N$`D_{15}`$ \[4$``$\] | 1676$`\pm `$2 | 1679$`\pm `$8 | 1675$`\pm `$10 |
| | | | N$`F_{15}`$ \[4$``$\] | 1684$`\pm `$4 | 1684$`\pm `$3 | 1680$`\pm `$10 |
| 14 | 1700.73 | 1699.67 | N$`D_{13}`$ \[3$``$\] | 1737$`\pm `$44 | 1731$`\pm `$15 | 1675$`\pm `$25 |
| | | | N$`P_{11}`$ \[3$``$\] | 1717$`\pm `$28 | 1723$`\pm `$9 | 1700$`\pm `$50 |
| | | | N$`P_{13}`$ \[4$``$\] | 1717$`\pm `$31 | 1710$`\pm `$20 | 1700$`\pm `$50 |
| 15 | 1763.52 | 1762.53 | | | | |
| 16 | 1827.50 | 1826.57 | $`D_{13}`$ | 1804$`\pm `$55 | | |
| | | | N$`P_{13}`$ \[2$``$\] | 1879$`\pm `$17 | | |
| 17 | 1892.59 | 1891.71 | $`P_{11}`$ | 1885$`\pm `$30 | | |
| | | | $`F_{15}`$ | 1903$`\pm `$87 | | |
| 18 | 1958.70 | 1957.87 | N$`F_{17}`$ \[2$``$\] | 2086$`\pm `$28 | 2005$`\pm `$150 | 1970$`\pm `$50 |
| | | | $`S_{11}`$ | 1928$`\pm `$59 | | |
| 19 | 2025.77 | 2024.98 | N$`F_{15}`$ \[2$``$\] | 1903$`\pm `$97 | 1882$`\pm `$10 | |
| | | | N$`D_{13}`$ \[2$``$\] | 1804$`\pm `$55 | 2081$`\pm `$20 | 2060$`\pm `$80 |
| 20 | 2093.73 | 2092.98 | $`F_{17}`$ | 2086$`\pm `$28 | | |
| | | | N$`S_{11}`$ \[1$``$\] | 1928$`\pm `$59 | 1880$`\pm `$20 | 2180$`\pm `$80 |
| | | | N$`P_{11}`$ \[1$``$\] | | 2050$`\pm `$20 | |
| 21 | 2162.52 | 2161.81 | N$`G_{17}`$ \[4$``$\] | 2127$`\pm `$9 | 2140$`\pm `$12 | 2200$`\pm `$70 |
| | | | N$`D_{15}`$ \[2$``$\] | | 2228$`\pm `$30 | |
| 22 | 2232.08 | 2231.39 | N$`H_{19}`$ \[4$``$\] | | | 2230$`\pm `$80 |
| | | | N$`G_{19}`$ \[4$``$\] | | 2250$`\pm `$80 | |
| 23 | 2302.36 | 2301.70 | | | | |
| 24 | 2373.30 | 2372.68 | | | | |
| 25 | 2444.88 | 2444.27 | | | | |
| 26 | 2517.03 | 2516.45 | | | | |
| 27 | 2589.73 | 2589.17 | N$`I_{1,11}`$ \[3$``$\] | | 2577$`\pm `$50 | |
| 28 | 2662.94 | 2662.40 | N$`K_{1,13}`$ \[2$``$\] | | 2612$`\pm `$45 | |
| 29 | 2736.63 | 2736.10 | | | | |
| 30 | 2810.76 | 2810.25 | | | | |
Table 7. Kaluza-Klein tower of KK excitations for the $`p\pi `$ system compared with the masses of the $`\mathrm{\Delta }`$ baryons determined from the KSU , KH and CMB analyses.
| n | $`M_n^{p\pi ^\pm }`$MeV | $`M_n^{p\pi ^0}`$MeV | $`\mathrm{\Delta }`$(PDG) \[St:{$``$}\] | KSU | KH | CMB |
| --- | --- | --- | --- | --- | --- | --- |
| 1 | 1084.79 | 1080.39 | | | | |
| 2 | 1104.30 | 1100.37 | | | | |
| 3 | 1133.48 | 1130.08 | | | | |
| 4 | 1169.65 | 1166.72 | | | | |
| 5 | 1210.92 | 1208.38 | | | | |
| 6 | 1256.07 | 1253.85 | | | | |
| 7 | 1304.35 | 1302.38 | | | | |
| 8 | 1355.23 | 1353.48 | | | | |
| 9 | 1408.38 | 1406.80 | | | | |
| 10 | 1463.54 | 1462.10 | | | | |
| 11 | 1520.50 | 1519.18 | | | | |
| 12 | 1579.11 | 1577.89 | $`\mathrm{\Delta }[1600]P_{33}`$ \[3$``$\] | 1706$`\pm `$10 | 1522$`\pm `$13 | 1600$`\pm `$50 |
| 13 | 1639.22 | 1638.09 | $`\mathrm{\Delta }[1620]S_{31}`$ \[4$``$\] | 1672$`\pm `$7 | 1610$`\pm `$7 | 1620$`\pm `$20 |
| 14 | 1700.73 | 1699.67 | $`\mathrm{\Delta }[1700]D_{33}`$ \[4$``$\] | 1762$`\pm `$44 | 1680$`\pm `$70 | 1710$`\pm `$30 |
| 15 | 1763.52 | 1762.53 | $`\mathrm{\Delta }[1750]P_{31}`$ \[1$``$\] | 1744$`\pm `$36 | | |
| | | | $`F_{35}`$ | 1752$`\pm `$32 | | |
| 16 | 1827.50 | 1826.57 | | | | |
| | | | $`\mathrm{\Delta }[1900]S_{31}`$ \[2$``$\] | 1920$`\pm `$24 | 1908$`\pm `$30 | 1890$`\pm `$50 |
| 17 | 1892.59 | 1891.71 | $`\mathrm{\Delta }[1905]F_{35}`$ \[4$``$\] | 1881$`\pm `$18 | 1905$`\pm `$20 | 1910$`\pm `$30 |
| | | | $`\mathrm{\Delta }[1910]P_{31}`$ \[4$``$\] | 1882$`\pm `$10 | 1888$`\pm `$20 | 1910$`\pm `$40 |
| | | | $`S_{31}`$ | 1920$`\pm `$24 | | |
| | | | $`\mathrm{\Delta }[1920]P_{33}`$ \[3$``$\] | 2014$`\pm `$16 | 1868$`\pm `$10 | 1920$`\pm `$80 |
| 18 | 1958.70 | 1957.87 | $`\mathrm{\Delta }[1930]D_{35}`$ \[3$``$\] | 1956$`\pm `$22 | 1901$`\pm `$15 | 1940$`\pm `$30 |
| | | | $`\mathrm{\Delta }[1940]D_{33}`$ \[1$``$\] | 2057$`\pm `$110 | | 1940$`\pm `$100 |
| | | | $`\mathrm{\Delta }[1950]F_{37}`$ \[4$``$\] | 1945$`\pm `$2 | 1913$`\pm `$8 | 1950$`\pm `$15 |
| 19 | 2025.77 | 2024.98 | $`\mathrm{\Delta }[2000]F_{35}`$ \[2$``$\] | 1752$`\pm `$32 | | 2200$`\pm `$125 |
| | | | $`P_{33}`$ | 2014$`\pm `$16 | | |
| 20 | 2093.73 | 2092.98 | $`D_{33}`$ | 2057$`\pm `$110 | | |
| 21 | 2162.52 | 2161.81 | $`\mathrm{\Delta }[2150]S_{31}`$ \[1$``$\] | | | 2150$`\pm `$100 |
| | | | $`D_{35}`$ | 2171$`\pm `$18 | | |
| 22 | 2232.08 | 2231.39 | $`\mathrm{\Delta }[2200]G_{37}`$ \[1$``$\] | | 2215$`\pm `$60 | |
| 23 | 2302.36 | 2301.70 | $`\mathrm{\Delta }[2300]H_{39}`$ \[2$``$\] | | 2217$`\pm `$80 | |
| 24 | 2373.30 | 2372.68 | $`\mathrm{\Delta }[2350]D_{35}`$ \[2$``$\] | | 2305$`\pm `$26 | 2400$`\pm `$125 |
| | | | $`\mathrm{\Delta }[2390]F_{37}`$ \[1$``$\] | | | 2350$`\pm `$100 |
| 25 | 2444.88 | 2444.27 | $`\mathrm{\Delta }[2400]G_{39}`$ \[2$``$\] | | 2468$`\pm `$50 | |
| | | | $`\mathrm{\Delta }[2420]H_{3,11}`$ \[4$``$\] | | 2416$`\pm `$17 | |
| 26 | 2517.03 | 2516.45 | | | | |
| 27 | 2589.73 | 2589.17 | | | | |
| 28 | 2662.94 | 2662.40 | | | | |
| 29 | 2736.63 | 2736.10 | $`\mathrm{\Delta }[2750]I_{3,13}`$ \[2$``$\] | | 2794$`\pm `$80 | |
| 30 | 2810.76 | 2810.25 | | | | |
Table 8. Kaluza-Klein tower of KK excitations for the $`p\pi \pi `$ system and $`\mathrm{\Delta }`$ baryons.
| n | $`M_n^{p\mathrm{\hspace{0.17em}2}\pi ^0}`$MeV | $`M_n^{p\pi ^0\pi ^\pm }`$MeV | $`M_n^{p\mathrm{\hspace{0.17em}2}\pi ^\pm }`$MeV | $`\mathrm{\Delta }`$ BARYON | $`M_{exp}^\mathrm{\Delta }`$ MeV |
| --- | --- | --- | --- | --- | --- |
| 1 | 1221.60 | 1226.00 | 1230.40 | $`\mathrm{\Delta }[1232]`$ $`P_{33}`$ | 1231$`\pm `$1 Manley 92 |
| 2 | 1258.80 | 1262.73 | 1266.66 | | |
| 3 | 1313.67 | 1317.07 | 1320.47 | | |
| 4 | 1380.61 | 1383.54 | 1386.47 | | |
| 5 | 1455.84 | 1458.38 | 1460.91 | | |
| 6 | 1536.98 | 1539.20 | 1541.42 | $`\mathrm{\Delta }[1600]`$ $`P_{33}`$ | 1522$`\pm `$13 Hoehler 79 |
| 7 | 1622.59 | 1624.55 | 1626.52 | $`\mathrm{\Delta }[1620]`$ $`S_{31}`$ | 1620$`\pm `$20 Cutkosky 80 |
| 8 | 1711.73 | 1713.48 | 1715.24 | $`\mathrm{\Delta }[1700]`$ $`D_{33}`$ | 1710$`\pm `$30 Cutkosky 80 |
| 9 | 1803.78 | 1805.37 | 1806.95 | $`S_{31}`$ | 1802$`\pm `$87 Vrana 00 |
| | | | | $`\mathrm{\Delta }[1900]`$ $`S_{31}`$ | 1890$`\pm `$50 Cutkosky 80 |
| 10 | 1898.32 | 1899.76 | 1901.20 | $`\mathrm{\Delta }[1905]`$ $`F_{35}`$ | 1881$`\pm `$18 Manley 92 |
| | | | | $`\mathrm{\Delta }[1910]`$ $`P_{31}`$ | 1888$`\pm `$20 Hoehler 79 |
| 11 | 1995.02 | 1996.34 | 1997.66 | $`\mathrm{\Delta }(1910)`$ $`P_{31}`$ | 1995$`\pm `$12 Vrana 00 |
| | | | | $`\mathrm{\Delta }[1920]`$ $`P_{33}`$ | 1920$`\pm `$80 Cutkosky 80 |
| | | | | $`\mathrm{\Delta }[1930]`$ $`D_{35}`$ | 1956$`\pm `$22 Manley 92 |
| | | | | $`\mathrm{\Delta }[1940]`$ $`D_{33}`$ | 1940$`\pm `$100 Cutkosky 80 |
| | | | | $`\mathrm{\Delta }[1950]`$ $`F_{37}`$ | 1950$`\pm `$15 Cutkosky 80 |
| | | | | $`\mathrm{\Delta }[2000]`$ $`F_{35}`$ | 1752$`\pm `$32 Manley 92 |
| 12 | 2093.64 | 2094.86 | 2096.08 | $`\mathrm{\Delta }[2150]`$ $`S_{31}`$ | 2150$`\pm `$100 Cutkosky 80 |
| | | | | $`P_{33}`$ | 2065$`{}_{12.9}{}^{}{}_{}{}^{+13.6}`$ Chew 80 |
| 13 | 2193.98 | 2195.11 | 2196.25 | $`\mathrm{\Delta }[2150]`$ $`S_{31}`$ | 2150$`\pm `$100 Cutkosky 80 |
| | | | | $`\mathrm{\Delta }[2200]`$ $`G_{37}`$ | 2215$`\pm `$60 Hoehler 79 |
| 14 | 2295.89 | 2296.94 | 2298.00 | $`\mathrm{\Delta }[2300]`$ $`H_{39}`$ | 2217$`\pm `$80 Hoehler 79 |
| | | | | $`\mathrm{\Delta }[2350]`$ $`D_{35}`$ | 2305$`\pm `$26 Hoehler 79 |
| | | | | $`\mathrm{\Delta }[2390]`$ $`F_{37}`$ | 2350$`\pm `$100 Cutkosky 80 |
| 15 | 2399.22 | 2400.21 | 2401.20 | $`\mathrm{\Delta }[2400]`$ $`G_{39}`$ | 2468$`\pm 50`$ Hoehler 79 |
| | | | | $`\mathrm{\Delta }[2420]H_{3,11}`$ | 2416$`\pm 17`$ Hoehler 79 |
| 16 | 2503.85 | 2504.78 | 2505.72 | | |
| 17 | 2609.69 | 2610.57 | 2611.45 | $`\mathrm{\Delta }[2750]I_{3,13}`$ | 2650$`\pm `$100 Hendry 78 |
| 18 | 2716.64 | 2717.47 | 2718.30 | $`\mathrm{\Delta }[2750]I_{3,13}`$ | 2794$`\pm `$80 Hoehler 79 |
| 19 | 2824.60 | 2825.39 | 2826.18 | $`\mathrm{\Delta }[2950]K_{3,15}`$ | 2850$`\pm `$100 Hendry 78 |
| 20 | 2933.52 | 2934.27 | 2935.02 | $`\mathrm{\Delta }[2950]K_{3,15}`$ | 2990$`\pm `$100 Hoehler 79 |
| 21 | 3043.31 | 3044.02 | 3044.74 | | |
| 22 | 3153.91 | 3154.59 | 3155.28 | $`K_{3,13}`$ | 3200$`\pm `$200 Hendry 78 |
| 23 | 3265.27 | 3265.92 | 3266.58 | $`L_{3,17}`$ | 3300$`\pm `$200 Hendry 78 |
| 24 | 3377.33 | 3377.96 | 3378.59 | | |
| 25 | 3490.05 | 3490.65 | 3491.26 | | |
| 26 | 3603.38 | 3603.96 | 3604.54 | | |
| 27 | 3717.27 | 3717.83 | 3718.39 | $`M_{3,19}`$ | 3700$`\pm `$200 Hendry 78 |
| 28 | 3831.69 | 3832.23 | 3832.77 | | |
| 29 | 3946.61 | 3947.13 | 3947.65 | | |
| 30 | 4061.99 | 4062.49 | 4063.00 | $`N_{3,21}`$ | 4100$`\pm `$300 Hendry 78 |
Table 9. Kaluza-Klein tower of KK excitations for the $`n\pi \pi `$ system and $`\mathrm{\Delta }`$ baryons.
| n | $`M_n^{n2\pi ^0}`$MeV | $`M_n^{n\pi ^0\pi ^\pm }`$MeV | $`M_n^{n2\pi ^\pm }`$MeV | $`\mathrm{\Delta }`$ BARYON | $`M_{exp}^\mathrm{\Delta }`$ MeV |
| --- | --- | --- | --- | --- | --- |
| 1 | 1222.89 | 1227.29 | 1231.69 | $`\mathrm{\Delta }[1232]`$ $`P_{33}`$ | 1231$`\pm `$1 Manley 92 |
| 2 | 1260.09 | 1264.02 | 1267.95 | | |
| 3 | 1314.95 | 1318.35 | 1321.75 | | |
| 4 | 1381.89 | 1384.82 | 1387.74 | | |
| 5 | 1457.10 | 1459.64 | 1462.17 | | |
| 6 | 1538.23 | 1540.45 | 1542.67 | $`\mathrm{\Delta }[1600]`$ $`P_{33}`$ | 1522$`\pm `$13 Hoehler 79 |
| 7 | 1623.82 | 1625.79 | 1627.75 | $`\mathrm{\Delta }[1620]`$ $`S_{31}`$ | 1620$`\pm `$20 Cutkosky 80 |
| 8 | 1712.95 | 1714.70 | 1716.46 | $`\mathrm{\Delta }[1700]`$ $`D_{33}`$ | 1710$`\pm `$30 Cutkosky 80 |
| 9 | 1804.98 | 1806.57 | 1808.15 | $`S_{31}`$ | 1802$`\pm `$87 Vrana 00 |
| | | | | $`\mathrm{\Delta }[1900]`$ $`S_{31}`$ | 1890$`\pm `$50 Cutkosky 80 |
| 10 | 1899.50 | 1900.94 | 1902.39 | $`\mathrm{\Delta }[1905]`$ $`F_{35}`$ | 1881$`\pm `$18 Manley 92 |
| | | | | $`\mathrm{\Delta }[1910]`$ $`P_{31}`$ | 1888$`\pm `$20 Hoehler 79 |
| 11 | 1996.18 | 1997.50 | 1998.83 | $`\mathrm{\Delta }(1910)`$ $`P_{31}`$ | 1995$`\pm `$12 Vrana 00 |
| | | | | $`\mathrm{\Delta }[1920]`$ $`P_{33}`$ | 1920$`\pm `$80 Cutkosky 80 |
| | | | | $`\mathrm{\Delta }[1930]`$ $`D_{35}`$ | 1956$`\pm `$22 Manley 92 |
| | | | | $`\mathrm{\Delta }[1940]`$ $`D_{33}`$ | 1940$`\pm `$100 Cutkosky 80 |
| | | | | $`\mathrm{\Delta }[1950]`$ $`F_{37}`$ | 1950$`\pm `$15 Cutkosky 80 |
| | | | | $`\mathrm{\Delta }[2000]`$ $`F_{35}`$ | 1752$`\pm `$32 Manley 92 |
| 12 | 2093.64 | 2094.86 | 2096.08 | $`\mathrm{\Delta }[2150]`$ $`S_{31}`$ | 2150$`\pm `$100 Cutkosky 80 |
| | | | | $`P_{33}`$ | 2065$`{}_{12.9}{}^{}{}_{}{}^{+13.6}`$ Chew 80 |
| 13 | 2193.98 | 2195.11 | 2196.25 | $`\mathrm{\Delta }[2150]`$ $`S_{31}`$ | 2150$`\pm `$100 Cutkosky 80 |
| | | | | $`\mathrm{\Delta }[2200]`$ $`G_{37}`$ | 2215$`\pm `$60 Hoehler 79 |
| 14 | 2296.99 | 2298.04 | 2299.10 | $`\mathrm{\Delta }[2300]`$ $`H_{39}`$ | 2217$`\pm `$80 Hoehler 79 |
| | | | | $`\mathrm{\Delta }[2350]`$ $`D_{35}`$ | 2305$`\pm `$26 Hoehler 79 |
| | | | | $`\mathrm{\Delta }[2390]`$ $`F_{37}`$ | 2350$`\pm `$100 Cutkosky 80 |
| 15 | 2400.30 | 2401.29 | 2402.28 | $`\mathrm{\Delta }[2400]`$ $`G_{39}`$ | 2468$`\pm 50`$ Hoehler 79 |
| | | | | $`\mathrm{\Delta }[2420]H_{3,11}`$ | 2416$`\pm 17`$ Hoehler 79 |
| 16 | 2504.91 | 2505.84 | 2506.77 | | |
| 17 | 2610.73 | 2611.60 | 2612.48 | $`\mathrm{\Delta }[2750]I_{3,13}`$ | 2650$`\pm `$100 Hendry 78 |
| 18 | 2717.65 | 2718.49 | 2719.31 | $`\mathrm{\Delta }[2750]I_{3,13}`$ | 2794$`\pm `$80 Hoehler 79 |
| 19 | 2825.59 | 2826.38 | 2827.17 | $`\mathrm{\Delta }[2950]K_{3,15}`$ | 2850$`\pm `$100 Hendry 78 |
| 20 | 2934.48 | 2935.24 | 2935.98 | $`\mathrm{\Delta }[2950]K_{3,15}`$ | 2990$`\pm `$100 Hoehler 79 |
| 21 | 3044.25 | 3044.97 | 3045.68 | | |
| 22 | 3154.84 | 3155.52 | 3156.20 | $`K_{3,13}`$ | 3200$`\pm `$200 Hendry 78 |
| 23 | 3266.18 | 3266.83 | 3267.49 | $`L_{3,17}`$ | 3300$`\pm `$200 Hendry 78 |
| 24 | 3378.22 | 3378.85 | 3379.48 | | |
| 25 | 3490.92 | 3491.52 | 3492.12 | | |
| 26 | 3604.22 | 3604.80 | 3605.38 | | |
| 27 | 3718.10 | 3718.66 | 3719.22 | $`M_{3,19}`$ | 3700$`\pm `$200 Hendry 78 |
| 28 | 3832.50 | 3833.04 | 3833.58 | | |
| 29 | 3947.40 | 3947.93 | 3948.45 | | |
| 30 | 4062.77 | 4063.27 | 4063.78 | $`N_{3,21}`$ | 4100$`\pm `$300 Hendry 78 |
Table 10. Kaluza-Klein tower of KK excitations for the $`N\rho `$ system and $`N,\mathrm{\Delta }`$ baryons.
| n | $`M_n^{n\rho }`$ MeV | $`M_n^{p\rho }`$ MeV | $`M_{exp}^N`$ MeV | $`M_{exp}^\mathrm{\Delta }`$ MeV |
| --- | --- | --- | --- | --- |
| 1 | 1710.90 | 1709.61 | N$`P_{11}`$ | $`D_{33}`$(1710$`\pm `$30) |
| 2 | 1716.98 | 1715.69 | $`P_{11}`$(1717$`\pm `$28) | |
| 3 | 1727.07 | 1725.79 | N$`P_{13}`$ | |
| 4 | 1741.09 | 1739.82 | | $`P_{31}`$(1744$`\pm `$36) |
| 5 | 1758.95 | 1757.69 | $`P_{11}`$(1766$`\pm `$34) | $`\mathrm{\Delta }[1750]P_{31}`$ |
| | | | | $`F_{35}`$(1752$`\pm `$32) |
| 6 | 1780.53 | 1779.28 | | |
| 7 | 1805.69 | 1804.45 | $`D_{13}`$(1804$`\pm `$55) | $`S_{31}`$(1802$`\pm `$87) |
| 8 | 1834.27 | 1833.05 | | |
| 9 | 1866.12 | 1864.92 | $`P_{13}`$(1879$`\pm `$17) | $`S_{31}`$(1890$`\pm `$50) |
| 10 | 1901.07 | 1899.89 | N$`P_{13}`$ | $`\mathrm{\Delta }[1900]S_{31}`$ |
| | | | | $`\mathrm{\Delta }[1905]F_{35}`$ |
| | | | | $`\mathrm{\Delta }[1910]P_{31}`$ |
| | | | | $`\mathrm{\Delta }[1920]P_{33}`$ |
| | | | | $`\mathrm{\Delta }[1930]D_{35}`$ |
| 11 | 1938.94 | 1937.78 | | $`\mathrm{\Delta }[1940]D_{33}`$ |
| | | | | $`\mathrm{\Delta }[1950]F_{37}`$ |
| | | | N$`F_{17}`$ | |
| 12 | 1979.58 | 1978.43 | $`F_{17}`$(1970$`\pm `$50) | $`F_{37}`$(1950$`\pm `$15) |
| 13 | 2022.80 | 2021.68 | N$`F_{15}`$ | $`\mathrm{\Delta }[2000]F_{35}`$ |
| | | | $`F_{15}`$($``$2025) | $`P_{33}`$(2014$`\pm `$16) |
| 14 | 2068.44 | 2067.34 | N$`D_{13}`$ | $`D_{33}`$(2057$`\pm `$110) |
| | | | N$`S_{11}`$ | |
| 15 | 2116.35 | 2115.27 | N$`P_{11}`$ | |
| | | | N$`G_{17}`$ | $`\mathrm{\Delta }[2150]S_{31}`$ |
| 16 | 2166.37 | 2165.31 | $`G_{17}`$(2168$`\pm `$18) | $`D_{35}`$(2171$`\pm `$18) |
| | | | $`S_{11}`$(2180$`\pm `$80) | $`S_{31}`$(2150$`\pm `$100) |
| | | | N$`D_{15}`$ | $`\mathrm{\Delta }[2200]G_{37}`$ |
| 17 | 2218.36 | 2217.33 | $`D_{15}`$(2228$`\pm `$30) | $`G_{37}`$(2215$`\pm `$60) |
| | | | N$`H_{19}`$ | $`H_{39}`$(2217$`\pm `$80) |
| 18 | 2272.19 | 2271.17 | N$`G_{19}`$ | |
| 19 | 2327.72 | 2326.73 | $`H_{19}`$(2300$`\pm `$100) | $`\mathrm{\Delta }[2300]H_{39}`$ |
| | | | | $`\mathrm{\Delta }[2350]D_{35}`$ |
| 20 | 2384.83 | 2383.86 | | $`\mathrm{\Delta }[2390]F_{37}`$ |
| | | | | $`\mathrm{\Delta }[2400]G_{39}`$ |
| 21 | 2443.43 | 2442.48 | | $`\mathrm{\Delta }[2420]H_{3,11}`$ |
| | | | | $`G_{39}`$(2468$`\pm `$50) |
| 22 | 2503.39 | 2502.46 | | $`H_{39}`$(2450$`\pm `$100) |
| 23 | 2564.63 | 2563.72 | $`I_{1,11}`$(2577$`\pm `$50) | |
| 24 | 2627.06 | 2626.17 | $`K_{1,13}`$(2612$`\pm `$45) | $`I_{3,13}`$(2650$`\pm `$100) |
| 25 | 2690.58 | 2689.72 | $`I_{1,11}`$(2700$`\pm `$100) | |
| 26 | 2755.14 | 2754.29 | | $`\mathrm{\Delta }[2750]I_{3,13}`$ |
| 27 | 2820.66 | 2819.83 | | $`I_{3,13}`$(2794$`\pm `$80) |
| 28 | 2887.06 | 2886.25 | | $`K_{3,15}`$(2850$`\pm `$100) |
| 29 | 2954.31 | 2953.51 | | $`\mathrm{\Delta }[2950]K_{3,15}`$ |
| 30 | 3022.32 | 3021.54 | $`K_{1,13}`$(3000$`\pm `$100) | |
Table 11. Kaluza-Klein tower of KK excitations for the $`N\eta `$ system and $`N,\mathrm{\Delta }`$ baryons.
| n | $`M_n^{n\eta }`$ MeV | $`M_n^{p\eta }`$ MeV | $`M_{exp}^N`$ MeV | $`M_{exp}^\mathrm{\Delta }`$ MeV |
| --- | --- | --- | --- | --- |
| 1 | 1489.35 | 1488.06 | | |
| 2 | 1496.77 | 1495.48 | | |
| 3 | 1509.04 | 1507.76 | | |
| 4 | 1526.00 | 1524.73 | N$`D_{13}`$ | |
| 5 | 1547.47 | 1546.20 | N$`S_{11}`$ | |
| 6 | 1573.20 | 1571.95 | | |
| 7 | 1602.97 | 1601.73 | | $`\mathrm{\Delta }[1600]P_{33}`$ |
| 8 | 1636.50 | 1635.28 | N$`S_{11}`$ | $`\mathrm{\Delta }[1620]S_{31}`$ |
| 9 | 1673.52 | 1672.32 | N$`D_{15}`$ | |
| | | | N$`F_{15}`$ | |
| | | | N$`D_{13}`$ | $`\mathrm{\Delta }[1700]D_{33}`$ |
| 10 | 1713.80 | 1712.61 | N$`P_{11}`$ | |
| | | | N$`P_{13}`$ | |
| 11 | 1757.06 | 1755.90 | $`P_{11}`$(1766$`\pm `$34) | $`\mathrm{\Delta }[1750]P_{31}`$ |
| 12 | 1803.09 | 1801.94 | $`D_{13}`$(1804$`\pm `$55) | $`S_{31}`$(1802$`\pm `$87) |
| 13 | 1851.65 | 1850.53 | | |
| 14 | 1902.55 | 1901.45 | N$`P_{13}`$ | $`\mathrm{\Delta }[1900]S_{31}`$ |
| | | | | $`\mathrm{\Delta }[1905]F_{35}`$ |
| | | | | $`\mathrm{\Delta }[1910]P_{31}`$ |
| | | | | $`\mathrm{\Delta }[1920]P_{33}`$ |
| | | | | $`\mathrm{\Delta }[1930]D_{35}`$ |
| | | | | $`\mathrm{\Delta }[1940]D_{33}`$ |
| 15 | 1955.58 | 1954.51 | N$`F_{17}`$ | $`\mathrm{\Delta }[1950]F_{37}`$ |
| 16 | 2010.59 | 2009.54 | N$`F_{15}`$ | $`\mathrm{\Delta }[2000]F_{35}`$ |
| 17 | 2067.40 | 2066.37 | N$`D_{13}`$ | $`D_{33}`$(2057$`\pm `$110) |
| | | | N$`S_{11}`$ | |
| 18 | 2125.88 | 2124.87 | N$`P_{11}`$ | $`\mathrm{\Delta }[2150]S_{31}`$ |
| 19 | 2185.89 | 2184.90 | N$`G_{17}`$ | $`D_{35}`$(2171$`\pm `$18) |
| | | | N$`D_{15}`$ | $`\mathrm{\Delta }[2200]G_{37}`$ |
| 20 | 2247.31 | 2246.34 | N$`H_{19}`$ | $`H_{39}`$(2217$`\pm `$80) |
| | | | N$`G_{19}`$ | |
| 21 | 2310.02 | 2309.07 | $`H_{19}`$(2300$`\pm `$100) | $`\mathrm{\Delta }[2300]H_{39}`$ |
| | | | | $`D_{35}`$(2305$`\pm `$26) |
| 22 | 2373.93 | 2373.00 | | $`\mathrm{\Delta }[2390]F_{37}`$ |
| 23 | 2438.94 | 2438.03 | | $`\mathrm{\Delta }[2400]G_{39}`$ |
| | | | | $`\mathrm{\Delta }[2420]H_{3,11}`$ |
| 24 | 2504.98 | 2504.09 | | $`H_{39}`$(2450$`\pm `$100) |
| 25 | 2571.95 | 2571.09 | $`I_{1,11}`$(2577$`\pm `$50) | |
| 26 | 2639.81 | 2638.96 | $`K_{1,13}`$(2612$`\pm `$45) | $`I_{3,13}`$(2650$`\pm `$100) |
| 27 | 2708.47 | 2707.64 | $`I_{1,11}`$(2700$`\pm `$100) | |
| 28 | 2777.89 | 2777.07 | | $`\mathrm{\Delta }[2750]I_{3,13}`$ |
| 29 | 2848.00 | 2847.20 | | $`K_{3,15}`$(2850$`\pm `$100) |
| 30 | 2918.77 | 2917.99 | | $`\mathrm{\Delta }[2950]K_{3,15}`$ |
Table 12. Kaluza-Klein tower of KK excitations for the $`N\omega `$ system and $`N,\mathrm{\Delta }`$ baryons.
| n | $`M_n^{n\omega }`$ MeV | $`M_n^{p\omega }`$ MeV | $`M_{exp}^N`$ MeV | $`M_{exp}^\mathrm{\Delta }`$ MeV |
| --- | --- | --- | --- | --- |
| 1 | 1724.15 | 1722.86 | N$`P_{13}`$ | $`D_{33}`$(1710$`\pm `$30) |
| 2 | 1730.18 | 1728.89 | | |
| 3 | 1740.17 | 1738.89 | | $`P_{31}`$(1744$`\pm `$36) |
| 4 | 1754.07 | 1752.80 | | $`\mathrm{\Delta }[1750]P_{31}`$ |
| | | | | $`F_{35}`$(1752$`\pm `$32) |
| 5 | 1771.77 | 1770.51 | $`P_{11}`$(1766$`\pm `$34) | |
| 6 | 1793.17 | 1791.92 | | |
| 7 | 1818.11 | 1816.88 | $`D_{13}`$(1804$`\pm `$55) | $`S_{31}`$(1802$`\pm `$87) |
| 8 | 1846.47 | 1845.25 | | |
| 9 | 1878.08 | 1876.88 | $`P_{13}`$(1879$`\pm `$17) | $`S_{31}`$(1890$`\pm `$50) |
| | | | N$`P_{13}`$ | $`\mathrm{\Delta }[1900]S_{31}`$ |
| | | | | $`\mathrm{\Delta }[1905]F_{35}`$ |
| 10 | 1912.77 | 1911.59 | | $`\mathrm{\Delta }[1910]P_{31}`$ |
| | | | | $`\mathrm{\Delta }[1920]P_{33}`$ |
| | | | | $`\mathrm{\Delta }[1930]D_{35}`$ |
| 11 | 1950.38 | 1949.22 | | $`\mathrm{\Delta }[1940]D_{33}`$ |
| | | | | $`\mathrm{\Delta }[1950]F_{37}`$ |
| 12 | 1990.75 | 1989.60 | N$`F_{17}`$ | |
| | | | $`F_{17}`$(1970$`\pm `$50) | $`F_{37}`$(1950$`\pm `$15) |
| 13 | 2033.69 | 2032.57 | N$`F_{15}`$ | $`\mathrm{\Delta }[2000]F_{35}`$ |
| | | | $`F_{15}`$($``$2025) | $`P_{33}`$(2014$`\pm `$16) |
| 14 | 2079.06 | 2077.96 | N$`D_{13}`$ | $`D_{33}`$(2057$`\pm `$110) |
| | | | N$`S_{11}`$ | |
| 15 | 2126.70 | 2125.62 | N$`P_{11}`$ | |
| | | | N$`G_{17}`$ | $`\mathrm{\Delta }[2150]S_{31}`$ |
| 16 | 2176.45 | 2175.40 | $`G_{17}`$(2168$`\pm `$18) | $`D_{35}`$(2171$`\pm `$18) |
| | | | $`S_{11}`$(2180$`\pm `$80) | $`S_{31}`$(2150$`\pm `$100) |
| | | | N$`D_{15}`$ | $`\mathrm{\Delta }[2200]G_{37}`$ |
| 17 | 2228.18 | 2227.15 | $`D_{15}`$(2228$`\pm `$30) | $`G_{37}`$(2215$`\pm `$60) |
| | | | N$`H_{19}`$ | $`H_{39}`$(2217$`\pm `$80) |
| 18 | 2281.75 | 2280.74 | N$`G_{19}`$ | |
| 19 | 2337.03 | 2336.04 | $`H_{19}`$(2300$`\pm `$100) | $`\mathrm{\Delta }[2300]H_{39}`$ |
| | | | | $`\mathrm{\Delta }[2350]D_{35}`$ |
| 20 | 2393.90 | 2392.93 | | $`\mathrm{\Delta }[2390]F_{37}`$ |
| | | | | $`\mathrm{\Delta }[2400]G_{39}`$ |
| 21 | 2452.25 | 2451.30 | | $`\mathrm{\Delta }[2420]H_{3,11}`$ |
| | | | | $`G_{39}`$(2468$`\pm `$50) |
| 22 | 2511.99 | 2511.06 | | $`H_{39}`$(2450$`\pm `$100) |
| 23 | 2573.00 | 2572.10 | $`I_{1,11}`$(2577$`\pm `$50) | |
| 24 | 2635.21 | 2634.32 | $`K_{1,13}`$(2612$`\pm `$45) | $`I_{3,13}`$(2650$`\pm `$100) |
| 25 | 2698.53 | 2697.67 | $`I_{1,11}`$(2700$`\pm `$100) | |
| 26 | 2762.89 | 2762.04 | | $`\mathrm{\Delta }[2750]I_{3,13}`$ |
| 27 | 2828.21 | 2827.38 | | $`I_{3,13}`$(2794$`\pm `$80) |
| 28 | 2894.44 | 2893.62 | | $`K_{3,15}`$(2850$`\pm `$100) |
| 29 | 2961.50 | 2960.70 | | $`\mathrm{\Delta }[2950]K_{3,15}`$ |
| 30 | 3029.34 | 3028.56 | $`K_{1,13}`$(3000$`\pm `$100) | |
Table 13. Kaluza-Klein tower of KK excitations for the $`\mathrm{\Lambda }K`$ system and $`N,\mathrm{\Delta }`$ baryons.
| n | $`M_n^{\mathrm{\Lambda }K^\pm }`$MeV | $`M_n^{\mathrm{\Lambda }K^0}`$MeV | $`M_{exp}^N`$ MeV | $`M_{exp}^\mathrm{\Delta }`$ MeV |
| --- | --- | --- | --- | --- |
| 1 | 1611.87 | 1615.85 | | $`\mathrm{\Delta }[1600]P_{33}`$ |
| 2 | 1619.36 | 1623.30 | | $`\mathrm{\Delta }[1620]S_{31}`$ |
| 3 | 1631.72 | 1635.60 | | |
| 4 | 1648.77 | 1652.56 | N$`S_{11}`$ | |
| 5 | 1670.27 | 1673.96 | N$`D_{15}`$ | |
| | | | N$`F_{15}`$ | |
| 6 | 1695.97 | 1699.54 | N$`D_{13}`$ | $`\mathrm{\Delta }[1700]D_{33}`$ |
| 7 | 1725.59 | 1729.04 | N$`P_{11}`$ | |
| | | | N$`P_{13}`$ | |
| 8 | 1758.84 | 1762.16 | $`P_{11}`$(1766$`\pm `$34) | $`\mathrm{\Delta }[1750]P_{31}`$ |
| 9 | 1795.43 | 1798.62 | $`D_{13}`$(1804$`\pm `$55) | $`S_{31}`$(1802$`\pm `$87) |
| 10 | 1835.12 | 1838.18 | | |
| 11 | 1877.64 | 1880.57 | $`P_{13}`$(1879$`\pm `$17) | |
| 12 | 1922.76 | 1925.58 | N$`P_{13}`$ | $`\mathrm{\Delta }[1900]S_{31}`$ |
| | | | | $`\mathrm{\Delta }[1905]F_{35}`$ |
| | | | | $`\mathrm{\Delta }[1910]P_{31}`$ |
| | | | | $`\mathrm{\Delta }[1920]P_{33}`$ |
| | | | | $`\mathrm{\Delta }[1930]D_{35}`$ |
| | | | | $`\mathrm{\Delta }[1940]D_{33}`$ |
| 13 | 1970.28 | 1972.98 | N$`F_{17}`$ | $`\mathrm{\Delta }[1950]F_{37}`$ |
| 14 | 2019.99 | 2022.59 | N$`F_{15}`$ | $`\mathrm{\Delta }[2000]F_{35}`$ |
| 15 | 2071.74 | 2074.22 | N$`D_{13}`$ | $`D_{33}`$(2057$`\pm `$110) |
| | | | N$`S_{11}`$ | |
| 16 | 2125.34 | 2127.73 | N$`P_{11}`$ | $`\mathrm{\Delta }[2150]S_{31}`$ |
| 17 | 2180.67 | 2182.97 | N$`G_{17}`$ | $`D_{35}`$(2171$`\pm `$18) |
| | | | N$`D_{15}`$ | $`\mathrm{\Delta }[2200]G_{37}`$ |
| 18 | 2237.59 | 2239.80 | N$`H_{19}`$ | $`H_{39}`$(2217$`\pm `$80) |
| | | | N$`G_{19}`$ | |
| 19 | 2295.98 | 2298.11 | $`H_{19}`$(2300$`\pm `$100) | $`\mathrm{\Delta }[2300]H_{39}`$ |
| 20 | 2355.73 | 2357.78 | | $`\mathrm{\Delta }[2350]D_{35}`$ |
| | | | | $`F_{37}`$(2350$`\pm `$100) |
| 21 | 2416.75 | 2418.72 | | $`\mathrm{\Delta }[2400]G_{39}`$ |
| | | | | $`\mathrm{\Delta }[2420]H_{3,11}`$ |
| 22 | 2478.94 | 2480.84 | | $`H_{39}`$(2450$`\pm `$100) |
| 23 | 2542.22 | 2544.06 | $`I_{1,11}`$(2577$`\pm `$50) | |
| 24 | 2606.51 | 2608.29 | $`K_{1,13}`$(2612$`\pm `$45) | |
| 25 | 2671.76 | 2673.48 | $`I_{1,11}`$(2700$`\pm `$100 | $`I_{3,13}`$(2650$`\pm `$100) |
| 26 | 2737.88 | 2739.55 | | $`\mathrm{\Delta }[2750]I_{3,13}`$ |
| 27 | 2804.83 | 2806.45 | | $`I_{3,13}`$(2794$`\pm `$80) |
| 28 | 2872.56 | 2874.13 | | $`K_{3,15}`$(2850$`\pm `$100) |
| 29 | 2941.00 | 2942.52 | | $`\mathrm{\Delta }[2950]K_{3,15}`$ |
| 30 | 3010.12 | 3011.60 | $`K_{1,13}`$(3000$`\pm `$100) | |
Table 14. KK excitations for the $`(\mathrm{\Sigma }K)^\pm `$ system and $`N,\mathrm{\Delta }`$ baryons.
| n | $`M_n^{\mathrm{\Sigma }^+K^0}`$MeV | $`M_n^{\mathrm{\Sigma }^{}K^0}`$MeV | $`M_n^{\mathrm{\Sigma }^0K^\pm }`$MeV | $`M_{exp}^N`$ MeV | $`M_{exp}^\mathrm{\Delta }`$ MeV |
| --- | --- | --- | --- | --- | --- |
| 1 | 1689.49 | 1697.59 | 1688.78 | N$`F_{15}`$ | |
| 2 | 1696.80 | 1704.88 | 1696.12 | N$`D_{13}`$ | $`\mathrm{\Delta }[1700]D_{33}`$ |
| 3 | 1708.86 | 1716.91 | 1708.24 | N$`P_{11}`$ | |
| 4 | 1725.49 | 1733.51 | 1724.94 | N$`P_{13}`$ | |
| 5 | 1746.48 | 1754.46 | 1746.02 | | $`\mathrm{\Delta }[1750]P_{31}`$ |
| | | | | | $`F_{35}`$(1752$`\pm `$32) |
| 6 | 1771.57 | 1779.50 | 1771.20 | $`P_{11}`$(1766$`\pm `$34) | |
| 7 | 1800.49 | 1808.36 | 1800.22 | $`D_{13}`$(1804$`\pm `$55) | $`S_{31}`$(1802$`\pm `$87) |
| 8 | 1832.96 | 1840.77 | 1832.80 | | |
| 9 | 1868.72 | 1876.45 | 1868.65 | $`P_{13}`$(1879$`\pm `$17) | |
| | | | | | $`\mathrm{\Delta }[1905]F_{35}`$ |
| 10 | 1907.51 | 1915.16 | 1907.54 | N$`P_{13}`$ | $`\mathrm{\Delta }[1900]S_{31}`$ |
| | | | | | $`\mathrm{\Delta }[1910]P_{31}`$ |
| | | | | | $`\mathrm{\Delta }[1920]P_{33}`$ |
| | | | | | $`\mathrm{\Delta }[1930]D_{35}`$ |
| 11 | 1949.08 | 1956.65 | 1949.20 | | $`\mathrm{\Delta }[1940]D_{33}`$ |
| | | | | | $`\mathrm{\Delta }[1950]F_{37}`$ |
| 12 | 1993.22 | 2000.70 | 1993.42 | N$`F_{17}`$ | $`P_{33}`$(2014$`\pm `$16) |
| 13 | 2039.72 | 2047.10 | 2039.99 | N$`F_{15}`$ | $`\mathrm{\Delta }[2000]F_{35}`$ |
| | | | | N$`D_{13}`$ | $`D_{33}`$(2057$`\pm `$110) |
| 14 | 2088.39 | 2095.67 | 2088.73 | N$`S_{11}`$ | |
| | | | | N$`P_{11}`$ | |
| 15 | 2139.06 | 2146.24 | 2139.47 | | $`\mathrm{\Delta }[2150]S_{31}`$ |
| 16 | 2191.58 | 2198.66 | 2192.05 | N$`G_{17}`$ | $`\mathrm{\Delta }[2200]G_{37}`$ |
| | | | | N$`D_{15}`$ | $`D_{35}`$(2171$`\pm `$18) |
| | | | | N$`H_{19}`$ | |
| 17 | 2245.82 | 2252.79 | 2246.34 | N$`G_{19}`$ | |
| 18 | 2301.64 | 2308.50 | 2302.20 | $`H_{19}`$(2300) | $`\mathrm{\Delta }[2300]H_{39}`$ |
| 19 | 2358.92 | 2365.68 | 2359.53 | | $`\mathrm{\Delta }[2350]D_{35}`$ |
| | | | | | $`\mathrm{\Delta }[2390]F_{37}`$ |
| 20 | 2417.58 | 2424.23 | 2418.21 | | $`\mathrm{\Delta }[2400]G_{39}`$ |
| | | | | | $`\mathrm{\Delta }[2420]H_{3,11}`$ |
| 21 | 2477.50 | 2484.04 | 2478.17 | | $`H_{39}`$(2450$`\pm `$100) |
| 22 | 2538.61 | 2545.04 | 2539.30 | | |
| 23 | 2600.81 | 2607.14 | 2601.53 | N$`I_{1,11}`$ | |
| 24 | 2664.05 | 2670.27 | 2664.78 | | $`I_{3,13}`$(2650$`\pm `$100) |
| 25 | 2728.25 | 2734.36 | 2728.99 | N$`K_{1,13}`$ | |
| 26 | 2793.35 | 2799.36 | 2794.10 | | $`\mathrm{\Delta }[2750]I_{3,13}`$ |
| 27 | 2859.29 | 2865.20 | 2860.06 | | $`K_{3,15}`$(2850$`\pm `$100) |
| 28 | 2926.02 | 2931.83 | 2926.80 | | $`\mathrm{\Delta }[2950]K_{3,15}`$ |
| 29 | 2993.49 | 2999.20 | 2994.27 | $`K_{1,13}`$(3000) | $`K_{3,15}`$(2990$`\pm `$100) |
| 30 | 3061.66 | 3067.27 | 3062.45 | | |
Table 15. KK excitations for the $`(\mathrm{\Sigma }K)^{0,++,}`$ system and $`N,\mathrm{\Delta }`$ baryons.
| n | $`M_n^{\mathrm{\Sigma }^+K^\pm }`$MeV | $`M_n^{\mathrm{\Sigma }^{}K^\pm }`$MeV | $`M_n^{\mathrm{\Sigma }^0K^0}`$MeV | $`M_{exp}^N`$ MeV | $`M_{exp}^\mathrm{\Delta }`$ MeV |
| --- | --- | --- | --- | --- | --- |
| 1 | 1685.51 | 1693.60 | 1692.76 | N$`F_{15}`$ | |
| 2 | 1692.86 | 1700.94 | 1700.06 | N$`D_{13}`$ | $`\mathrm{\Delta }[1700]D_{33}`$ |
| 3 | 1704.98 | 1713.04 | 1712.11 | N$`P_{11}`$ | |
| 4 | 1721.70 | 1729.73 | 1728.73 | N$`P_{13}`$ | |
| 5 | 1742.79 | 1750.78 | 1749.70 | | $`\mathrm{\Delta }[1750]P_{31}`$ |
| | | | | | $`F_{35}`$(1752$`\pm `$32) |
| 6 | 1768.00 | 1775.93 | 1774.77 | $`P_{11}`$(1766$`\pm `$34) | |
| 7 | 1797.04 | 1804.91 | 1803.67 | $`D_{13}`$(1804$`\pm `$55) | $`S_{31}`$(1802$`\pm `$87) |
| 8 | 1829.64 | 1837.45 | 1836.12 | | |
| 9 | 1865.53 | 1873.26 | 1871.84 | $`P_{13}`$(1879$`\pm `$17) | |
| | | | | | $`\mathrm{\Delta }[1905]F_{35}`$ |
| 10 | 1904.45 | 1912.10 | 1910.60 | N$`P_{13}`$ | $`\mathrm{\Delta }[1900]S_{31}`$ |
| | | | | | $`\mathrm{\Delta }[1910]P_{31}`$ |
| | | | | | $`\mathrm{\Delta }[1920]P_{33}`$ |
| | | | | | $`\mathrm{\Delta }[1930]D_{35}`$ |
| 11 | 1946.14 | 1953.71 | 1952.14 | | $`\mathrm{\Delta }[1940]D_{33}`$ |
| | | | | | $`\mathrm{\Delta }[1950]F_{37}`$ |
| 12 | 1990.40 | 1997.88 | 1996.24 | N$`F_{17}`$ | $`P_{33}`$(2014$`\pm `$16) |
| 13 | 2037.01 | 2044.40 | 2042.70 | N$`F_{15}`$ | $`\mathrm{\Delta }[2000]F_{35}`$ |
| | | | | N$`D_{13}`$ | $`D_{33}`$(2057$`\pm `$110) |
| 14 | 2085.79 | 2093.08 | 2091.33 | N$`S_{11}`$ | |
| | | | | N$`P_{11}`$ | |
| 15 | 2136.57 | 2143.75 | 2141.96 | | $`\mathrm{\Delta }[2150]S_{31}`$ |
| 16 | 2189.19 | 2196.27 | 2194.44 | N$`G_{17}`$ | $`\mathrm{\Delta }[2200]G_{37}`$ |
| | | | | N$`D_{15}`$ | $`D_{35}`$(2171$`\pm `$18) |
| | | | | N$`H_{19}`$ | |
| 17 | 2243.52 | 2250.49 | 2248.63 | N$`G_{19}`$ | |
| 18 | 2299.43 | 2306.29 | 2304.41 | $`H_{19}`$(2300) | $`\mathrm{\Delta }[2300]H_{39}`$ |
| 19 | 2356.80 | 2363.56 | 2361.65 | | $`\mathrm{\Delta }[2350]D_{35}`$ |
| | | | | | $`\mathrm{\Delta }[2390]F_{37}`$ |
| 20 | 2415.53 | 2422.18 | 2420.26 | | $`\mathrm{\Delta }[2400]G_{39}`$ |
| | | | | | $`\mathrm{\Delta }[2420]H_{3,11}`$ |
| 21 | 2475.52 | 2482.07 | 2480.14 | | $`H_{39}`$(2450$`\pm `$100) |
| 22 | 2536.70 | 2543.13 | 2541.20 | | |
| 23 | 2598.97 | 2605.30 | 2603.37 | N$`I_{1,11}`$ | |
| 24 | 2662.27 | 2668.49 | 2666.56 | | $`I_{3,13}`$(2650$`\pm `$100) |
| 25 | 2726.53 | 2732.64 | 2730.72 | N$`K_{1,13}`$ | |
| 26 | 2791.68 | 2797.69 | 2795.77 | | $`\mathrm{\Delta }[2750]I_{3,13}`$ |
| 27 | 2857.67 | 2863.58 | 2861.67 | | $`K_{3,15}`$(2850$`\pm `$100) |
| 28 | 2924.45 | 2930.26 | 2928.36 | | $`\mathrm{\Delta }[2950]K_{3,15}`$ |
| 29 | 2991.97 | 2997.68 | 2995.80 | $`K_{1,13}`$(3000) | $`K_{3,15}`$(2990$`\pm `$100) |
| 30 | 3060.18 | 3065.79 | 3063.93 | | |
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# Locating the peaks of semilinear elliptic systems
## 1 Introduction and main results
In the asymptotic analysis of the singularly perturbed elliptic equation
$$\epsilon ^2\mathrm{\Delta }u+u=f(x,u)\text{in }^n,u>0\text{in }^n,$$
($`P_\epsilon `$)
there are well known situations where the associated ground energy function $`\mathrm{\Sigma }`$ (cf. ) is $`C^1`$-smooth and around its nondegenerate critical points the solutions $`u_\epsilon `$ of $`P_\epsilon `$ exhibit a spike-like profile as $`\epsilon `$ goes to zero. This is the case, for instance, for the power nonlinearity
$$f(x,u)=K(x)u^q,1<q<\frac{n+2}{n2},n3,$$
where $`K(x)`$ is a suitable $`C^1`$ function (see e.g. and references therein). It turns out that the $`C^1`$ (and higher) smoothness of $`\mathrm{\Sigma }`$ is related to the crucial fact that, for every fixed $`z^n`$, the limiting autonomous equation
$$\mathrm{\Delta }u+u=f(z,u)\text{in }^n,u>0\text{in }^n,$$
($`P_0`$)
admits a unique solution, up to translations . However, unfortunately, the uniqueness feature for $`P_0`$ is a delicate matter and it is currently available only under rather restrictive assumptions on $`f`$ (cf. e.g. ). What it is know, in general, is only that $`\mathrm{\Sigma }`$ is a locally Lipschitz continuous function which admits representation formulas for the left and right derivatives (cf. \[30, Lemma 2.3\]). Motivated by these facts, recently, some conditions for locating the concentration points for $`P_\epsilon `$ in presence of a more general nonlinearity $`f`$, not necessarily of power type, have been investigated (see and also ). The underlying philosophy is that when the limit problem $`P_0`$ lacks of uniqueness up to translations, then the ground energy function $`\mathrm{\Sigma }`$ could loose its additional regularity properties.
Nevertheless, in this (possibly nonsmooth) framework, it turns out that a necessary condition for the solutions $`u_\epsilon `$ to concentrate (in a suitable sense) around a given point $`z`$ is that it is critical for $`\mathrm{\Sigma }`$ in the sense of the Clarke subdifferential $`_C`$, that is $`0_C\mathrm{\Sigma }(z)`$, or in a even weaker sense. The main theme of this note is the search of suitable conditions for locating the spikes, as $`\epsilon 0`$, of the solutions to the semilinear model system
$$\{\begin{array}{cc}\epsilon ^2\mathrm{\Delta }u+u=K(x)v^q,\hfill & \text{in }^n,\hfill \\ \multicolumn{2}{c}{}\\ \epsilon ^2\mathrm{\Delta }v+v=Q(x)u^p,\hfill & \text{in }^n,\hfill \\ \multicolumn{2}{c}{}\\ u,v>0,\hfill & \text{in }^n,\hfill \end{array}$$
($`S_\epsilon `$)
where $`p>1`$ and $`q>1`$ are lying below the so called โcritical hyperbolaโ
$$๐_n=\{(p,q)(1,\mathrm{})\times (1,\mathrm{}):\frac{1}{p+1}+\frac{1}{q+1}=1\frac{2}{n}\},n3,$$
which naturally arises in the study of this problem and constitutes the borderline between existence and nonexistence results (cf. e.g. ).
Now, according to what discussed right above, the interest in looking for conditions for the spike location of the solutions to ($`S_\epsilon `$) is mainly motivated by the following simple observation: contrary to the scalar case, there is no uniqueness result available in the literature for the (radial) solutions to the (limiting) system associated with ($`S_\epsilon `$)
$$\{\begin{array}{cc}\mathrm{\Delta }u+u=K(z)v^q,\hfill & \text{in }^n,\hfill \\ \multicolumn{2}{c}{}\\ \mathrm{\Delta }v+v=Q(z)u^p,\hfill & \text{in }^n,\hfill \\ \multicolumn{2}{c}{}\\ u,v>0,\hfill & \text{in }^n,\hfill \end{array}$$
($`S_z`$)
where $`z^n`$ is frozen and acts as a parameter. As a consequence, in the vectorial case, we do not know whether the (suitably defined) ground energy map $`\mathrm{\Sigma }`$ associated with ($`S_\epsilon `$) (cf. Definition 1.2) is $`C^1`$-smooth and admits an explicit representation formula. Hence, the necessary conditions in terms of Clarke subdifferential (or weaker) appear here even more natural than in the case of a single equation. See Section 1.2 for the statements of the main results, Theorems 1.1 and 1.2. As far as we are aware, other criteria for the concentration have been established so far, but all of them consider the scalar case. We refer the reader e.g. to for the case of power-like nonlinearities and to for more general classes of nonlinearities.
Semilinear systems like ($`S_\epsilon `$) naturally arise in the study of various kinds of nonlinear phenomena such as population evolution, pattern formation, chemical reaction, etc., being $`u`$ and $`v`$ the concentrations of different species in the process (see also and references therein). Visibly, the interest in the study of the various qualitative properties of ($`S_\epsilon `$) has steadily increased in recent times. In a smooth bounded domain $`\mathrm{\Omega }`$, $`(S_1)`$ was extensively studied by Clement, Costa, De Figueiredo, Felmer, Hulshof, Magalhรฃes, van der Vorst in . The asymptotic analysis with respect to $`\epsilon `$ has been very recently performed both with Dirichlet and Neumann boundary conditions by Pistoia-Ramos and Ramos-Yang . In the whole space $`^n`$, the existence of least energy solutions to ($`S_\epsilon `$) has been investigated by Alves-Carriรฃo-Miyagaki, De Figueiredo, Yang and Sirakov in , whereas the asymptotic behavior with respect to $`\epsilon `$ has been pursued by Alves-Soares-Yang in . Finally, for the exponential decay, the radial symmetry and the regularity properties of the solutions to ($`S_z`$), we refer the reader to the quite recent achievements of Busca-Sirakov and Sirakov .
The outline of the paper is as follows: in Sections 1.1-1.2 we provide preliminary stuff such as the (dual) variational framework and the (dual) ground energy function $`\mathrm{\Sigma }`$ and we state the main results of the paper. Throughout Section 2 we deal with the $`\mathrm{Lip}_{\mathrm{loc}}`$ regularity and the representation formulas of the directional derivatives for $`\mathrm{\Sigma }`$. Finally, in Section 3 we end up the proofs of the main results.
### 1.1 The dual variational functional
As it is known, if e.g. $`p`$ and $`q`$ are both less than $`\frac{n+2}{n2}`$, then system ($`S_\epsilon `$) admits a natural variational structure (of Hamiltonian type) which is based on the strongly indefinite functional $`f_\epsilon :H^1(^n)\times H^1(^n)`$,
$$f_\epsilon (u,v)=_^n\epsilon ^2uv+uv\frac{1}{q+1}_^nK(x)|v|^{q+1}\frac{1}{p+1}_^nQ(x)|u|^{p+1}.$$
However, as already done in , for our purposes, as well as for dealing with possibly supercritical values of $`p`$ or $`q`$, we consider a corresponding dual variational structure, mainly relying on the Legendre-Fenchel transformation (see e.g. and references therein). In the following, we just briefly recall some of the core ingredients, referring to \[1, Section 2\] for expanded details on this framework. For $`\frac{1}{p+1}+\frac{1}{q+1}>\frac{n2}{n}`$, consider the linear operators
$`T_1:L^{\frac{q+1}{q}}(^n)W^{2,\frac{q+1}{q}}(^n)L^{p+1}(^n),`$
$`T_2:L^{\frac{p+1}{p}}(^n)W^{2,\frac{p+1}{p}}(^n)L^{q+1}(^n),`$
defined as
$$T_1=T_2=(\mathrm{\Delta }+\mathrm{Id})^1.$$
Notice that $`T_1`$ and $`T_2`$ are continuous. Then, we consider the linear operator (take into account the proper Sobolev embeddings)
$$T:L^{\frac{p+1}{p}}(^n)\times L^{\frac{q+1}{q}}(^n)L^{p+1}(^n)\times L^{q+1}(^n),T=\left[\begin{array}{cc}0& T_1\\ T_2& 0\end{array}\right],$$
explicitly defined by
$$T\eta ,\xi =\xi _1T_1\eta _2+\xi _2T_2\eta _1,\eta =(\eta _1,\eta _2),\xi =(\xi _1,\xi _2).$$
Finally we introduce the Banach space $`(,_{})`$,
$$=L^{\frac{p+1}{p}}(^n)\times L^{\frac{q+1}{q}}(^n),\eta _{}^2=\eta _1_{L^{\frac{p+1}{p}}(^n)}^2+\eta _2_{L^{\frac{q+1}{q}}(^n)}^2$$
and the (dual) $`C^1`$ functional $`J_\epsilon :`$ defined as
$$J_\epsilon (\eta )=\frac{p}{p+1}_^n\frac{|\eta _1|^{\frac{p+1}{p}}}{Q^{\frac{1}{p}}(\epsilon x)}+\frac{q}{q+1}_^n\frac{|\eta _2|^{\frac{q+1}{q}}}{K^{\frac{1}{q}}(\epsilon x)}\frac{1}{2}_^nT\eta ,\eta .$$
If $`\eta ^\epsilon =(\eta _1^\epsilon ,\eta _2^\epsilon )`$ is a critical point of $`J_\epsilon `$, then $`(u_\epsilon (x),v_\epsilon (x))=(\overline{u}_\epsilon (\frac{x}{\epsilon }),\overline{v}_\epsilon (\frac{x}{\epsilon }))`$, with
$$(\overline{u}_\epsilon ,\overline{v}_\epsilon )=(T_1\eta _2^\epsilon ,T_2\eta _1^\epsilon )W^{2,\frac{q+1}{q}}L^{p+1}\times W^{2,\frac{p+1}{p}}L^{q+1},$$
(1.1)
corresponds to a solution to ($`S_\epsilon `$) with $`u_\epsilon (x),v_\epsilon (x)0`$ for $`|x|\mathrm{}`$ (see \[1, p.677\]). In light of the above summability, we have $`f_\epsilon (u_\epsilon ,v_\epsilon )`$ for all $`\epsilon >0`$. Analogously, associated with ($`S_z`$), we introduce the limiting functional $`I_z:`$
$$I_z(\eta )=\frac{p}{p+1}_^n\frac{|\eta _1|^{\frac{p+1}{p}}}{Q^{\frac{1}{p}}(z)}+\frac{q}{q+1}_^n\frac{|\eta _2|^{\frac{q+1}{q}}}{K^{\frac{1}{q}}(z)}\frac{1}{2}_^nT\eta ,\eta .$$
From the viewpoint of our investigation, the main advantage of exploiting the dual variational functional $`I_z`$ is that it admits a mountain-pass geometry and the mountain-pass value corresponds to the least possible energy of system ($`S_z`$). As we shall see in the next section, this allows to provide in the vectorial framework a suitable definition of ground energy function with nice features, similar to those available in the scalar case.
### 1.2 Preliminaries and the main results
In order to state the main achievements of the paper, we need some preparatory stuff. For the sake of self-containedness we shall also recall a few pretty well known notions from nonsmooth calculus (see e.g. ).
###### Definition 1.1
Let $`f:^n`$ be a locally Lipschitz function near a point $`z^n`$. The Clarke subdifferential of $`f`$ at $`z`$ is defined by
$$_Cf(z):=\{\eta ^n:f^0(z,w)\eta w,\text{for every }w^n\},$$
where $`f^0(z,w)`$ is the generalized derivative of $`f`$ at $`z`$ along $`w^n`$, defined by
$$f^0(z;w):=\underset{\begin{array}{c}\xi z\\ \lambda 0+\end{array}}{lim\; sup}\frac{f(\xi +\lambda w)f(\xi )}{\lambda }.$$
###### Definition 1.2
The (dual) ground energy function $`\mathrm{\Sigma }:^n`$ of ($`S_z`$) is given by
$$\mathrm{\Sigma }(z):=\underset{\eta ๐ฉ_z}{inf}I_z(\eta ),$$
where $`๐ฉ_z`$ is the Nehari manifold of $`I_z`$, that is
$$๐ฉ_z=\{\eta :\eta (0,0)\text{ and }I_z^{}(\eta )[\eta ]=0\}.$$
We shall denote by $`๐ฆ^n`$ the set of Clarke critical points of $`\mathrm{\Sigma }`$, namely
$$๐ฆ:=\{z^n:\mathrm{\hspace{0.17em}0}_C\mathrm{\Sigma }(z)\}.$$
###### Definition 1.3
We say that the pair $`(u_\epsilon ,v_\epsilon )`$ is a strong solution to system ($`S_\epsilon `$) if it is a distributional solution and $`(u_\epsilon ,v_\epsilon )W^{2,(q+1)/q}(^n)\times W^{2,(p+1)/p}(^n)`$. We say that the pair $`\eta ^\epsilon =(\eta _1^\epsilon ,\eta _2^\epsilon )`$ corresponding to $`(u_\epsilon ,v_\epsilon )`$ through (1.1) is the related dual solution.
###### Definition 1.4
We set
$`:=\{`$ zn:there exists a sequence of strong solutions (uฮตh,vฮตh) of (Sฮต) with:๐งsuperscript๐there exists a sequence of strong solutions (uฮตh,vฮตh) of (Sฮต) with\displaystyle z\in{\mathbb{R}}^{n}:\,\text{there exists a sequence of strong solutions
$(u_{\varepsilon_{h}},v_{\varepsilon_{h}})$ of~{}\eqref{problema} with}
$`|u_{\epsilon _h}(z)|,|v_{\epsilon _h}(z)|\delta `$ for some $`\delta >0`$, $`|u_{\epsilon _h}(z+\epsilon _hx)|,|v_{\epsilon _h}(z+\epsilon _hx)|0`$
$`\text{as }|x|\mathrm{}\text{ uniformly w.r.t. }h\text{, and }\epsilon _{h}^{}{}_{}{}^{n}f_{\epsilon _h}(u_{\epsilon _h},v_{\epsilon _h})\mathrm{\Sigma }(z)\text{ as }h\mathrm{}\}.`$
We say that $``$ is the energy concentration set for ($`S_\epsilon `$).
Assume that $`K,QC^1(^n)`$ and
$$\alpha K(x)\beta ,\alpha Q(x)\beta ,\text{for all }x^n,$$
(1.2)
$$|K(x)|,|Q(x)|Ce^{M|x|},\text{for all }x^n\text{ with }|x|\text{ large}.$$
(1.3)
for some positive constants $`\alpha ,\beta ,C`$ and $`M`$.
The main result of the paper, linking the energy concentration set $``$ with the set $`๐ฆ`$ of Clarke critical set of $`\mathrm{\Sigma }`$, is provided by the following
###### Theorem 1.1
Assume that $`K,QC^1(^n)`$ and that (1.2)-(1.3) hold. Then $`๐ฆ`$.
###### Remark 1.1
By \[2, Theorem 1\], under suitable assumptions, if there exists an absolute minimum (or maximum) point $`z_{}`$ for $`\mathrm{\Sigma }`$, then $`z_{}\mathrm{}`$.
###### Remark 1.2
As a straightforward combination of Theorem 1.1 with the well known convex hull characterization of $`_C\mathrm{\Sigma }(z)`$, if $`z`$ is a concentration point for ($`S_\epsilon `$), then
$$0\mathrm{Co}\{\underset{j}{lim}\mathrm{\Sigma }(\xi _j):\xi _j\mathrm{\Omega }\text{ and }\xi _jz\},$$
where $`\mathrm{Co}\{X\}`$ denotes the convex hull of $`X`$ and $`\mathrm{\Omega }`$ is any null set containing the set of points at which $`\mathrm{\Sigma }`$ fails to be differentiable.
###### Corollary 1.1
Under the (unproved) assumption that, for all $`z^n`$, system ($`S_z`$) admits a unique positive solution (up to translations), $`\mathrm{\Sigma }`$ is $`C^1`$-smooth and
$$\mathrm{Crit}\left(Q^{\frac{q+1}{pq1}}K^{\frac{p+1}{pq1}}\right),$$
where Crit(f) denotes the set of (classical) critical points of $`f`$.
In the following definition we consider solutions which concentrate close to a point $`z`$, with bounded energy but not necessary stabilizing towards $`\mathrm{\Sigma }(z)`$.
###### Definition 1.5
Let $`m1`$. We set
$`_m:=\{`$ zn:there exists a sequence of strong solutions (uฮตh,vฮตh) of (Sฮต) with:๐งsuperscript๐there exists a sequence of strong solutions (uฮตh,vฮตh) of (Sฮต) with\displaystyle z\in{\mathbb{R}}^{n}:\,\text{there exists a sequence of strong solutions
$(u_{\varepsilon_{h}},v_{\varepsilon_{h}})$ of~{}\eqref{problema} with}
$`|u_{\epsilon _h}(z)|,|v_{\epsilon _h}(z)|\delta `$ for some $`\delta >0`$, $`|u_{\epsilon _h}(z+\epsilon _hx)|,|v_{\epsilon _h}(z+\epsilon _hx)|0`$
$`\text{as }|x|\mathrm{}\text{ uniformly w.r.t. }h\text{, and }\epsilon _{h}^{}{}_{}{}^{n}f_{\epsilon _h}(u_{\epsilon _h},v_{\epsilon _h})m\text{ as }h\mathrm{}\}.`$
We say that $`_m`$ is the concentration set for ($`S_\epsilon `$) at the energy level $`m`$.
###### Definition 1.6
Let $`m1`$ and $`z^n`$. For every $`w^n`$ we define $`\mathrm{\Gamma }_{z,m}^{}(w)`$ by
$`\mathrm{\Gamma }_{z,m}^{}(w)`$ $`:=\underset{\eta ๐พ_m(z)}{sup}\left[{\displaystyle \frac{1}{p+1}}{\displaystyle \frac{Q}{w}}(z){\displaystyle _^n}{\displaystyle \frac{|\eta _1|^{\frac{p+1}{p}}}{Q^{\frac{p+1}{p}}(z)}}{\displaystyle \frac{1}{q+1}}{\displaystyle \frac{K}{w}}(z){\displaystyle _^n}{\displaystyle \frac{|\eta _2|^{\frac{q+1}{q}}}{K^{\frac{q+1}{q}}(z)}}\right],`$
$`\mathrm{\Gamma }_{z,m}^+(w)`$ $`:=\underset{\eta ๐พ_m(z)}{inf}\left[{\displaystyle \frac{1}{p+1}}{\displaystyle \frac{Q}{w}}(z){\displaystyle _^n}{\displaystyle \frac{|\eta _1|^{\frac{p+1}{p}}}{Q^{\frac{p+1}{p}}(z)}}{\displaystyle \frac{1}{q+1}}{\displaystyle \frac{K}{w}}(z){\displaystyle _^n}{\displaystyle \frac{|\eta _2|^{\frac{q+1}{q}}}{K^{\frac{q+1}{q}}(z)}}\right],`$
where $`๐พ_m(z)`$ denotes the set of all the nontrivial, radial, exponentially decaying solutions of ($`S_z`$) having energy equal to $`m`$.
It is readily seen that $`\mathrm{\Gamma }_{z,m}^{}(w)`$ for all $`z,w`$ in $`^n`$ (see the proof of (2.11)). It is also straightforward to check that, for any $`z^n`$, the functions $`\{w\mathrm{\Gamma }_{z,m}^{}(w)\}`$ are convex.
###### Definition 1.7
Let $`m1`$. We set
$$๐ฆ_m:=\{z^n:\mathrm{\hspace{0.17em}\hspace{0.17em}0}\mathrm{\Gamma }_{z,m}^{}(0)\mathrm{\Gamma }_{z,m}^+(0)\},$$
where $``$ stands for the subdifferential of convex functions,
$$\mathrm{\Gamma }_{z,m}^{}(0)=\{\xi ^n:\mathrm{\Gamma }_{z,m}^{}(w)\xi w\text{,\hspace{0.17em}\hspace{0.17em} for every }w^n\}.$$
It is known by standard convex analysis that $`\mathrm{\Gamma }_{z,m}^{}(0)\mathrm{}`$, for every $`z^n`$. Observe that $`z๐ฆ_m`$ if and only if $`0`$ is a critical point for both $`\mathrm{\Gamma }_{z,m}^{}`$ and $`\mathrm{\Gamma }_{z,m}^+`$. Of course, if $`๐พ_m(z)=\{\eta _0\}`$ was a singleton, then $`z๐ฆ_m`$ if and only if
$$\mathrm{\Gamma }_{z,m}^{}(w)=\mathrm{\Gamma }_{z,m}^+(w)=\frac{I_z}{w}(\eta _0)=0,w^n.$$
Without forcing the energy levels of the solutions to approach the least energy of the limit system, we get the following correlation between the sets $`_m`$ and $`๐ฆ_m`$.
###### Theorem 1.2
Assume that $`K,QC^1(^n)`$ and (1.2)-(1.3) hold. Then $`_m๐ฆ_m`$.
## 2 Properties of the ground energy function
Before coming to the proof of the results, we need some preliminary stuff.
### 2.1 Some preparatory lemmas
The next proposition is well known (see e.g. ); on the other hand, for the sake of completeness and self-containedness, we report a brief proof.
###### Proposition 2.1
Let $`z^n`$. Then $`(u,v)W^{2,\frac{q+1}{q}}(^n)\times W^{2,\frac{p+1}{p}}(^n)`$ is a solution to ($`S_z`$) if and only if $`\eta =(\eta _1,\eta _2)=(T_2^1v,T_1^1u)`$ is a critical point of $`I_z`$. Moreover, there holds $`f_z(u,v)=I_z(\eta _1,\eta _2)`$, where $`f_z`$ is the functional defined as
$$f_z(u,v)=_^nuv+uv\frac{1}{q+1}_^nK(z)|v|^{q+1}\frac{1}{p+1}_^nQ(z)|u|^{p+1}.$$
Proof. Observe first that, if $`(u,v)W^{2,\frac{q+1}{q}}(^n)\times W^{2,\frac{p+1}{p}}(^n)`$ solves ($`S_z`$), taking into account the Sobolev embedding, the value $`f_z(u,v)`$ is indeed finite (cf. (1.1)). Let $`(u,v)`$ be a solution to ($`S_z`$). Then, since
$$\eta _1=T_2^1v,\eta _2=T_1^1u,$$
we have
$$\{\begin{array}{cc}\eta _2=T_1^1u=\mathrm{\Delta }u+u=K(z)v^q,\hfill & \\ \multicolumn{2}{c}{}\\ \eta _1=T_2^1v=\mathrm{\Delta }v+v=Q(z)u^p.\hfill & \end{array}$$
Therefore, we get
$$T_2\eta _1=v=\frac{\eta _2^{\frac{1}{q}}}{K^{\frac{1}{q}}(z)}\mathrm{and}T_1\eta _2=u=\frac{\eta _1^{\frac{1}{p}}}{Q^{\frac{1}{p}}(z)},$$
(2.1)
and so $`(\eta _1,\eta _2)`$ is a critical point of $`I_z`$. Vice versa, if $`(\eta _1,\eta _2)`$ is a critical point of $`I_z`$, it is readily seen that (2.1) hold, so that $`(T_1\eta _2,T_2\eta _1)=(u,v)`$ is a solution to ($`S_z`$) (cf. \[1, p.677\]). Furthermore, on the solutions to ($`S_z`$), we have
$$f_z(u,v)=\left(\frac{1}{2}\frac{1}{p+1}\right)_^nQ(z)u^{p+1}+\left(\frac{1}{2}\frac{1}{q+1}\right)_^nK(z)v^{q+1}.$$
Then, in light of (2.1), we have
$`I_z(\eta )`$ $`=\frac{p}{p+1}{\displaystyle _^n}{\displaystyle \frac{|\eta _1|^{\frac{p+1}{p}}}{Q^{\frac{1}{p}}(z)}}+\frac{q}{q+1}{\displaystyle _^n}{\displaystyle \frac{|\eta _2|^{\frac{q+1}{q}}}{K^{\frac{1}{q}}(z)}}\frac{1}{2}{\displaystyle _^n}T\eta ,\eta `$
$`=\left(\frac{p}{p+1}\frac{1}{2}\right){\displaystyle _^n}\eta _1T_1\eta _2+\left(\frac{q}{q+1}\frac{1}{2}\right){\displaystyle _^n}\eta _2T_2\eta _1`$
$`=\left(\frac{p}{p+1}\frac{1}{2}\right){\displaystyle _^n}(\mathrm{\Delta }v+v)u+\left(\frac{q}{q+1}\frac{1}{2}\right){\displaystyle _^n}(\mathrm{\Delta }u+u)v`$
$`=\left(\frac{1}{2}\frac{1}{p+1}\right){\displaystyle _^n}Q(z)u^{p+1}+\left(\frac{1}{2}\frac{1}{q+1}\right){\displaystyle _^n}K(z)v^{q+1}=f_z(u,v),`$
which concludes the proof.
###### Definition 2.1
We say that $`\eta `$ is a dual solution to ($`S_z`$) if it is a critical point of $`I_z`$. We say that $`\eta `$ is a dual least energy solution to ($`S_z`$) if it is a dual solution and, in addition, $`I_z(\eta )=\mathrm{\Sigma }(z)`$.
The next property, classical in the scalar case, will be pretty useful for our purposes.
###### Lemma 2.1
For every $`z^n`$, let us set
$`b_1(z)`$ $`:=\underset{\eta \{0\}}{inf}\underset{t0}{sup}I_z(t\eta ),`$
$`b_2(z)`$ $`:=\underset{\eta ๐ฉ_z}{inf}I_z(\eta )=\mathrm{\Sigma }(z),`$
$`b_3(z)`$ :=inf{Iz(ฮท):ฮท{0} is a dual solution to (Sz)}.assignabsentinfimumconditional-setsubscript๐ผ๐ง๐ฮท{0} is a dual solution to (Sz)\displaystyle:=\inf\big{\{}I_{z}(\eta):\text{$\eta\in{{\mathscr{H}}}\setminus\{0\}$ is a dual solution to \eqref{limit-z}}\big{\}}.
Then $`b_1(z)=b_2(z)=b_3(z)`$. Moreover $`\{z\mathrm{\Sigma }(z)\}`$ is continuous.
Proof. The first equality follows from \[1, Lemma 2\]. Moreover in it is proved that $`b_1(z)=b_2(z)`$ is a critical value so that also $`b_2(z)=b_3(z)`$ follows. Finally, by virtue of \[2, Lemma 1\], we know that $`\mathrm{\Sigma }`$ is continuous.
###### Lemma 2.2
Let $`z^n`$ and define the (nonempty) set
$$_+:=\{\eta :_^nT\eta ,\eta >0\}.$$
Then, for every $`\eta _+`$, there exists a unique maximum point $`t_\eta >0`$ of the map $`\varphi :t(0,\mathrm{})I_z(t\eta )`$. In particular, $`t_\eta \eta ๐ฉ_z`$.
Proof. Let us observe that if $`\varphi ^{}(t)=0`$, then
$$_^nT\eta ,\eta =t^{\frac{1p}{p}}_^n\frac{|\eta _1|^{\frac{p+1}{p}}}{Q^{\frac{1}{p}}(z)}+t^{\frac{1q}{q}}_^n\frac{|\eta _2|^{\frac{q+1}{q}}}{K^{\frac{1}{q}}(z)}.$$
Since the function $`g(t)=At^{\frac{1q}{q}}+Bt^{\frac{1p}{p}}`$ with $`A,B>0`$ is strictly decreasing for $`t>0`$, then $`\varphi `$ has at most one critical value. It is easy to see that for all $`\eta `$, $`\varphi (t)>0`$ for $`t`$ small, while if $`\eta _+`$, it is readily seen that $`\varphi (t)<0`$ for big $`t`$โs.
### 2.2 Conjecturing the representation of $`\mathrm{\Sigma }`$
Consider for a moment the equation
$$\epsilon ^2\mathrm{\Delta }u+V(x)u=K(x)u^p,\text{in }^n,$$
(2.2)
with $`p`$ subcritical and $`V`$ and $`K`$ potentials functions bounded away from zero. By the results of , we know that there is uniqueness (up to translation) of positive solutions for
$$\mathrm{\Delta }u+u=u^p,\text{in }^n,$$
and, by a suitable change of variable, also for the โlimitโ problem at $`x=z`$ of (2.2)
$$\mathrm{\Delta }u+V(z)u=K(z)u^p,\text{in }^n.$$
This allows to give an explicit representation for the ground state function associated with (2.2), merely depending on the potentials $`V`$ and $`K`$ (see for example ):
$$\mathrm{\Sigma }(z)=\mathrm{\Gamma }\frac{V^{\frac{p+1}{p1}\frac{n}{2}}(z)}{K^{\frac{2}{p1}}(z)},$$
(2.3)
for a suitable positive constant $`\mathrm{\Gamma }`$. On the contrary, as already observed, up to our knowledge there is no (known) uniqueness result for the elliptic system
$$\mathrm{\Delta }\xi +\xi =\zeta ^q,\mathrm{\Delta }\zeta +\zeta =\xi ^p,\text{in }^n,$$
(2.4)
and so, in general, we cannot provide an explicit expression for $`\mathrm{\Sigma }`$ for ($`S_z`$). Slightly more in general, if $`V`$ is smooth and $`\alpha V(x)\beta `$, consider the system
$$\{\begin{array}{cc}\mathrm{\Delta }u+V(z)u=K(z)v^q,\hfill & \text{in }^n,\hfill \\ \multicolumn{2}{c}{}\\ \mathrm{\Delta }v+V(z)v=Q(z)u^p,\hfill & \text{in }^n,\hfill \\ \multicolumn{2}{c}{}\\ u,v>0,\hfill & \text{in }^n.\hfill \end{array}$$
(2.5)
Assuming for a moment that (2.4) has a unique solution $`(\xi ,\zeta )`$, then we claim that
$$\mathrm{\Sigma }(z)=\mathrm{\Gamma }\frac{V^{\frac{(p+1)(q+1)}{pq1}\frac{n}{2}}(z)}{Q^{\frac{q+1}{pq1}}(z)K^{\frac{p+1}{pq1}}(z)},$$
(2.6)
for a suitable positive constant $`\mathrm{\Gamma }`$. Indeed, by rescaling
$$u(x)=\varpi _1\xi (\mu x)\text{and}v(x)=\varpi _2\zeta (\mu x),$$
where we have set
$`\mu =\mu (z)`$ $`:=V^{\frac{1}{2}}(z),`$
$`\varpi _1=\varpi _1(z)`$ $`:={\displaystyle \frac{V^{\frac{q+1}{pq1}}(z)}{Q^{\frac{q}{pq1}}(z)K^{\frac{1}{pq1}}(z)}},`$
$`\varpi _2=\varpi _2(z)`$ $`:={\displaystyle \frac{V^{\frac{p+1}{pq1}}(z)}{Q^{\frac{1}{pq1}}(z)K^{\frac{p}{pq1}}(z)}},`$
it is easy to see that $`(u,v)`$ is the unique solution of the system (2.5). Hence, by a straightforward calculation, we reach (2.6). Let us observe that the exponent of $`V(z)`$ in (2.6) is equal to zero if, and only if, the pair $`(p,q)`$ belongs to $`๐_n`$. Then, for problems with powers $`p,q`$ close to the set $`๐_n`$, the potential $`V`$ is expected to have a weak influence in the location of concentration points. Notice that the same phenomenon appears in the scalar case (cf. formula (2.3)), since $`\frac{p+1}{p1}\frac{n}{2}0`$ if and only if $`p\frac{n+2}{n2}=2^{}1`$, where $`2^{}`$ is the critical Sobolev exponent for $`H^1`$. Finally, we just wish to mention that, incidentally, the exponents
$$\theta _1=\frac{p+1}{pq1},\theta _2=\frac{q+1}{pq1}$$
in formula (2.6) also arise in the study of the blow-up rates for the parabolic system
$$u_t=\mathrm{\Delta }u+v^q,v_t=\mathrm{\Delta }v+u^p,x\mathrm{\Omega },t>0,$$
with initial data $`u(x,0)=u_0(x)0`$, $`v(x,0)=v_0(x)0`$ and Dirichlet boundary conditions $`u=v=0`$ on $`\mathrm{\Omega }`$. Here $`\mathrm{\Omega }`$ is a ball in $`^n`$ and $`u_0`$ and $`v_0`$ are continuous which vanish on the boundary. If $`u_0`$, $`v_0`$ are nontrivial $`C^1`$ functions, the solution $`(u,v)`$ blows up at a finite time $`T<\mathrm{}`$, and $`u_t0`$, $`v_t0`$ on $`\mathrm{\Omega }\times (0,T)`$, then there exist two constants $`C>c>0`$ with
$$\frac{c}{(Tt)^{\theta _1}}\underset{\overline{\mathrm{\Omega }}}{\mathrm{max}}u(x,t)\frac{C}{(Tt)^{\theta _1}},\frac{c}{(Tt)^{\theta _2}}\underset{\overline{\mathrm{\Omega }}}{\mathrm{max}}v(x,t)\frac{C}{(Tt)^{\theta _2}},$$
for all $`t(0,T)`$. We refer the interested reader, e.g., to .
### 2.3 Local lipschitzianity of $`\mathrm{\Sigma }`$
In the case of a single semilinear elliptic equation, it is known that the ground energy map enjoys a basic regularity property, in addition to the continuity, namely it is locally Lipschitz continuous (hence differentiable a.e. by virtue of Rademacherโs theorem). Analogously, for system ($`S_\epsilon `$), we obtain the following
###### Theorem 2.1
$`\mathrm{\Sigma }\mathrm{Lip}_{\mathrm{loc}}(^n)`$.
Proof. Let $`\rho _0>0`$ and $`\mu ^n`$ with $`|\mu |\rho _0`$ and let $`\eta ^\mu `$ be a (dual) solution to ($`S_\mu `$) such that $`I_\mu (\eta ^\mu )=\mathrm{\Sigma }(\mu )`$ (we already know that such a solution does exist, see ). Then, the corresponding (direct) solution $`(u_\mu ,v_\mu )`$ satisfies
$$\mathrm{\Delta }u+u=K(\mu )v^q,\mathrm{\Delta }v+v=Q(\mu )u^p,\text{in }^n.$$
(2.7)
We also know that $`u_\mu `$ and $`v_\mu `$ are radially symmetric, radially decreasing with respect to, say, the origin, and exponentially decaying (see , in particular \[6, Theorem 2\] and \[27, Theorem 1(a)\]). We claim that there exist $`\varpi _1>0`$ and $`\varpi _2=\varpi _2(\rho _0)>0`$ independent of $`\mu `$ such that
$$\varpi _1u_\mu _{L^{p+1}}\varpi _2\text{and}\varpi _1v_\mu _{L^{q+1}}\varpi _2.$$
(2.8)
Let us prove first the estimates from below. By multiplying the first equation of (2.7) by $`u_\mu `$ and taking into account (1.2), we get
$`u_\mu _{H^1}^2={\displaystyle _^n}K(\mu )v_\mu ^qu_\mu \beta v_\mu _{L^{q+1}}^qu_\mu _{L^{q+1}}\beta Sv_\mu _{L^{q+1}}^qu_\mu _{H^1},`$ (2.9)
where $`S`$ is the Sobolev constant. Now, by multiplying the first equation of system (2.7) by $`v_\mu `$ and the second equation by $`u_\mu `$, and comparing the resulting equations, we have
$$v_\mu _{L^{q+1}}^q\left(\frac{\beta }{\alpha }\right)^{q/(q+1)}u_\mu _{L^{p+1}}^{q(p+1)/(q+1)}.$$
(2.10)
By combining inequalities (2.9) and (2.10), and using again the Sobolev inequality, the assertion follows. The proof of the estimate from below for $`v_\mu _{L^{q+1}}`$ is similar. To prove the inequalities from above we simply observe that $`\mathrm{\Sigma }`$ is continuous and
$`\underset{|\mu |\rho _0}{\mathrm{max}}\mathrm{\Sigma }(\mu )`$ $`=\underset{|\mu |\rho _0}{\mathrm{max}}I_\mu (\eta ^\mu )=\underset{|\mu |\rho _0}{\mathrm{max}}f_\mu (u_\mu ,v_\mu )`$
$`\left(\frac{\alpha }{2}\frac{\alpha }{q+1}\right)v_\mu _{L^{q+1}}^{q+1}+\left(\frac{\alpha }{2}\frac{\alpha }{p+1}\right)u_\mu _{L^{p+1}}^{p+1}.`$
Thus (2.8) follows. As a consequence, according to the definition of the dual norm $`_{}`$, we immediately obtain
$$\alpha \sqrt{\varpi _1^{2p}+\varpi _1^{2q}}\underset{|\mu |\rho _0}{\mathrm{max}}\eta ^\mu _{}\beta \sqrt{\varpi _2^{2p}(\rho _0)+\varpi _2^{2q}(\rho _0)}.$$
(2.11)
Now, since $`\eta ^\mu ๐ฉ_\mu `$, we get
$$_^nT\eta ^\mu ,\eta ^\mu =_^n\frac{|\eta _1^\mu |^{\frac{p+1}{p}}}{Q^{\frac{1}{p}}(\mu )}+_^n\frac{|\eta _2^\mu |^{\frac{q+1}{q}}}{K^{\frac{1}{q}}(\mu )}>0.$$
(2.12)
Hence $`\eta ^\mu _+`$ and, by means of Lemma 2.2, there exists precisely one positive number $`\theta (\mu ,\xi )`$ such that $`\theta (\mu ,\xi )\eta ^\mu ๐ฉ_\xi `$. By definition, this means that
$$_^nT\eta ^\mu ,\eta ^\mu =\theta (\mu ,\xi )^{\frac{1p}{p}}_^n\frac{|\eta _1^\mu |^{\frac{p+1}{p}}}{Q^{\frac{1}{p}}(\xi )}+\theta (\mu ,\xi )^{\frac{1q}{q}}_^n\frac{|\eta _2^\mu |^{\frac{q+1}{q}}}{K^{\frac{1}{q}}(\xi )}.$$
(2.13)
Moreover, we have $`\theta (\mu ,\mu )=1`$. Collecting these facts, we see that, by the implicit function theorem, $`\theta `$ is differentiable with respect to the variable $`\xi `$. Moreover, in light of (2.11), it results that $`\theta (\mu ,\xi )`$ remains bounded for $`\mu `$ and $`\xi `$ varying in a bounded set. Indeed, by combining (2.12) and (2.13), supposing for example that $`pq`$, we have
$$\theta (\mu ,\xi )^{\frac{p1}{p}}\left[_^n\frac{|\eta _1^\mu |^{\frac{p+1}{p}}}{Q^{\frac{1}{p}}(\mu )}+_^n\frac{|\eta _2^\mu |^{\frac{q+1}{q}}}{K^{\frac{1}{q}}(\mu )}\right]=_^n\frac{|\eta _1^\mu |^{\frac{p+1}{p}}}{Q^{\frac{1}{p}}(\xi )}+\theta (\mu ,\xi )^{\frac{1}{q}\frac{1}{p}}_^n\frac{|\eta _2^\mu |^{\frac{q+1}{q}}}{K^{\frac{1}{q}}(\xi )}.$$
Then the (local) boundedness of $`\theta (\mu ,\xi )`$ follows immediately by (2.11) and by the fact that $`\frac{1}{q}\frac{1}{p}0`$. Let us now observe that
$`I_\xi (\theta (\mu ,\xi )\eta ^\mu )`$ $`=\theta (\mu ,\xi )^{\frac{p+1}{p}}{\displaystyle \frac{p}{p+1}}{\displaystyle _^n}{\displaystyle \frac{|\eta _1^\mu |^{\frac{p+1}{p}}}{Q^{\frac{1}{p}}(\xi )}}`$
$`+\theta (\mu ,\xi )^{\frac{q+1}{q}}{\displaystyle \frac{q}{q+1}}{\displaystyle _^n}{\displaystyle \frac{|\eta _2^\mu |^{\frac{q+1}{q}}}{K^{\frac{1}{q}}(\xi )}}`$
$`{\displaystyle \frac{\theta (\mu ,\xi )^2}{2}}{\displaystyle _^n}T\eta ^\mu ,\eta ^\mu .`$
The gradient of the function $`\left\{\xi I_\xi (\theta (\mu ,\xi )\eta ^\mu )\right\}`$ is thus given by
$`_\xi I_\xi (\theta (\mu ,\xi )\eta ^\mu )=`$ $`{\displaystyle \frac{\theta (\mu ,\xi )^{\frac{p+1}{p}}}{p+1}}_\xi Q(\xi ){\displaystyle _^n}{\displaystyle \frac{|\eta _1^\mu |^{\frac{p+1}{p}}}{Q^{\frac{p+1}{p}}(\xi )}}`$
$`{\displaystyle \frac{\theta (\mu ,\xi )^{\frac{q+1}{q}}}{q+1}}_\xi K(\xi ){\displaystyle _^n}{\displaystyle \frac{|\eta _2^\mu |^{\frac{q+1}{q}}}{K^{\frac{q+1}{q}}(\xi )}}`$
$`+_\xi \theta (\mu ,\xi )[\theta (\mu ,\xi )^{\frac{1}{p}}{\displaystyle _^n}{\displaystyle \frac{|\eta _1^\mu |^{\frac{p+1}{p}}}{Q^{\frac{1}{p}}(\xi )}}`$
$`+\theta (\mu ,\xi )^{\frac{1}{q}}{\displaystyle _^n}{\displaystyle \frac{|\eta _2^\mu |^{\frac{q+1}{q}}}{K^{\frac{1}{q}}(\xi )}}`$
$`\theta (\mu ,\xi ){\displaystyle _^n}T\eta ^\mu ,\eta ^\mu ],`$
and so, since $`\theta (\mu ,\xi )\eta ^\mu ๐ฉ_\xi `$, in turn we get
$`_\xi I_\xi (\theta (\mu ,\xi )\eta ^\mu )=`$ $`{\displaystyle \frac{\theta (\mu ,\xi )^{\frac{p+1}{p}}}{p+1}}_\xi Q(\xi ){\displaystyle _^n}{\displaystyle \frac{|\eta _1^\mu |^{\frac{p+1}{p}}}{Q^{\frac{p+1}{p}}(\xi )}}`$ (2.14)
$`{\displaystyle \frac{\theta (\mu ,\xi )^{\frac{q+1}{q}}}{q+1}}_\xi K(\xi ){\displaystyle _^n}{\displaystyle \frac{|\eta _2^\mu |^{\frac{q+1}{q}}}{K^{\frac{q+1}{q}}(\xi )}}.`$
From this representation formula, the Mean-Value Theorem and the local boundedness of $`\theta `$, the assertion readily follows.
### 2.4 Left and right derivatives of $`\mathrm{\Sigma }`$
Let us define $`๐(z)`$ as the set of all the positive (dual) solutions $`(\eta _1,\eta _2)`$ of ($`S_z`$) at the energy level $`\mathrm{\Sigma }(z)`$. The representation formulas for the (left and right) directional derivatives of $`\mathrm{\Sigma }`$ are provided in the following
###### Theorem 2.2
The directional derivatives from the left and the right of $`\mathrm{\Sigma }`$ at every point $`z^n`$ along any $`w^n`$ exist and it holds
$`\left({\displaystyle \frac{\mathrm{\Sigma }}{w}}\right)^{}(z)=\underset{\eta ๐(z)}{sup}_zI_z(\eta )w,`$
$`\left({\displaystyle \frac{\mathrm{\Sigma }}{w}}\right)^+(z)=\underset{\eta ๐(z)}{inf}_zI_z(\eta )w.`$
Explicitly, we have
$`\left({\displaystyle \frac{\mathrm{\Sigma }}{w}}\right)^{}(z)`$ $`=\underset{\eta ๐(z)}{sup}\left[{\displaystyle \frac{1}{p+1}}{\displaystyle \frac{Q}{w}}(z){\displaystyle _^n}{\displaystyle \frac{|\eta _1|^{\frac{p+1}{p}}}{Q^{\frac{p+1}{p}}(z)}}{\displaystyle \frac{1}{q+1}}{\displaystyle \frac{K}{w}}(z){\displaystyle _^n}{\displaystyle \frac{|\eta _2|^{\frac{q+1}{q}}}{K^{\frac{q+1}{q}}(z)}}\right]`$
$`\left({\displaystyle \frac{\mathrm{\Sigma }}{w}}\right)^+(z)`$ $`=\underset{\eta ๐(z)}{inf}\left[{\displaystyle \frac{1}{p+1}}{\displaystyle \frac{Q}{w}}(z){\displaystyle _^n}{\displaystyle \frac{|\eta _1|^{\frac{p+1}{p}}}{Q^{\frac{p+1}{p}}(z)}}{\displaystyle \frac{1}{q+1}}{\displaystyle \frac{K}{w}}(z){\displaystyle _^n}{\displaystyle \frac{|\eta _2|^{\frac{q+1}{q}}}{K^{\frac{q+1}{q}}(z)}}\right].`$
for every $`z,w^n`$.
Proof. Let $`\{\mu _j\}^n`$ be a sequence converging to $`\mu _0`$ and let $`\eta ^j=\eta ^{\mu _j}`$ be a sequence of (dual) solutions of least energy $`\mathrm{\Sigma }(\mu _j)`$. We want to prove that, up to a subsequence,
$$\eta ^{\mu _j}\eta ^0,\text{strongly in },\eta ^0๐(\mu _0).$$
(2.15)
Consider the corresponding (direct) solutions $`(u_{\mu _j},v_{\mu _j})`$ (resp. $`(u_0,v_0)`$) of (2.7) with $`\mu =\mu _j`$ (resp. $`\mu =\mu _0`$). Since $`(u_{\mu _j},v_{\mu _j})`$ is bounded in $`W^{2,(q+1)/q}(^n)\times W^{2,(p+1)/p}(^n)`$ (cf. ), up to a subsequence, it converges weakly to a pair $`(u_0,v_0)`$. In addition, since $`K`$ and $`Q`$ are uniformly bounded, by virtue of the Schauder local regularity estimates (cf. ), $`(u_{\mu _j},v_{\mu _j})`$ is bounded in $`C_{\mathrm{loc}}^{2,\beta }(^n)`$ for some $`\beta >0`$ and
$$u_{\mu _j}u_0\text{and}v_{\mu _j}v_0,\text{locally in }C^2\text{-sense},$$
(2.16)
so that $`(u_0,v_0)`$ solves (2.7) with $`\mu =\mu _0`$. We claim that $`u_0>0`$ and $`v_0>0`$. By \[6, Theorem 2\], for every $`j1`$, $`u_{\mu _j}`$ and $`v_{\mu _j}`$ are radially symmetric and radially decreasing with respect to some point, say the origin, that is
$$u_{\mu _j}(x)=u_j(r),v_{\mu _j}(x)=v_j(r),\frac{d}{dr}u_j(r)<0,\frac{d}{dr}v_j(r)<0,$$
(2.17)
for every $`r>0`$. Hence, for every $`j1`$, we have
$`u_{\mu _j}(0)\mathrm{\Delta }u_{\mu _j}(0)+u_{\mu _j}(0)=K(\mu _j)v_{\mu _j}^q(0)\beta v_{\mu _j}^q(0),`$
$`v_{\mu _j}(0)\mathrm{\Delta }v_{\mu _j}(0)+v_{\mu _j}(0)=Q(\mu _j)u_{\mu _j}^p(0)\beta u_{\mu _j}^p(0).`$
It follows that, for every $`j1`$,
$$u_{\mu _j}(0)\beta ^{q+1}u_{\mu _j}^{pq}(0).$$
Then there exists $`\widehat{\delta }>0`$ such that $`u_{\mu _j}(0)\widehat{\delta }`$ for every $`j1`$. Similarly, $`v_{\mu _j}(0)\widehat{\delta }`$ for every $`j1`$. Hence, letting $`j\mathrm{}`$, by (2.16), we conclude that $`u_0(0)\widehat{\delta }`$ and $`v_0(0)\widehat{\delta }`$, which entails $`u_00`$ and $`v_00`$. Since we have $`u_00`$, $`v_00`$, $`K(\mu _0),Q(\mu _0)>0`$ and
$$\mathrm{\Delta }u_0+u_00\text{and}\mathrm{\Delta }v_0+v_00,$$
the claim just follows by a straightforward application of the maximum principle.
Observe that, by the continuity of $`\mathrm{\Sigma }`$ and by Fatouโs Lemma, we get
$$\mathrm{\Sigma }(\mu _0)=\underset{j\mathrm{}}{lim}\mathrm{\Sigma }(\mu _j)=\underset{j\mathrm{}}{lim}I_{\mu _j}(\eta ^j)I_{\mu _0}(\eta ^0)\mathrm{\Sigma }(\mu _0).$$
Hence
$$\underset{j\mathrm{}}{lim}I_{\mu _j}(\eta ^j)=I_{\mu _0}(\eta ^0)=\mathrm{\Sigma }(\mu _0),$$
which reads as
$$\underset{j\mathrm{}}{lim}_^n\frac{|\eta _1^j|^{\frac{p+1}{p}}}{Q^{\frac{1}{p}}(\mu _j)}=_^n\frac{|\eta _1^0|^{\frac{p+1}{p}}}{Q^{\frac{1}{p}}(\mu _0)},\underset{j\mathrm{}}{lim}_^n\frac{|\eta _2^j|^{\frac{q+1}{q}}}{K^{\frac{1}{q}}(\mu _j)}=_^n\frac{|\eta _2^0|^{\frac{q+1}{q}}}{K^{\frac{1}{q}}(\mu _0)}.$$
In particular, taking into account (1.2), for any $`\delta >0`$, there exists $`\rho >0`$ such that
$$_{\{|x|\rho \}}|\eta _1^j|^{\frac{p+1}{p}}<\delta ,_{\{|x|\rho \}}|\eta _2^j|^{\frac{q+1}{q}}<\delta ,$$
for every $`j1`$ sufficiently large. Moreover, of course
$$\underset{j\mathrm{}}{lim}_{\{|x|\rho \}}|\eta _1^j|^{\frac{p+1}{p}}=_{\{|x|\rho \}}|\eta _1^0|^{\frac{p+1}{p}},\underset{j\mathrm{}}{lim}_{\{|x|\rho \}}|\eta _2^j|^{\frac{q+1}{q}}=_{\{|x|\rho \}}|\eta _2^0|^{\frac{q+1}{q}}.$$
Then we have $`\eta ^{\mu _j}\eta ^0`$ strongly in $``$, namely (2.15) holds true.
Without loss of generality, we can prove the formula of the right derivative of $`\mathrm{\Sigma }`$ in the case $`n=1`$, $`z=0`$ and $`w=1`$. For any $`\eta ^0๐(0)`$, we get
$`\mathrm{\Sigma }(\rho )\mathrm{\Sigma }(0)`$ $`I_\rho (\vartheta (\rho ,0)\eta ^0)I_0(\eta ^0)`$
$`=\rho _\xi I_\xi (\vartheta (\xi ,0)\eta ^0)|_{\xi =\mu [0,\rho ]}.`$
Whence, by virtue of (2.14) and the arbitrariness of $`\eta ^0๐(0)`$,
$$\underset{\rho 0^+}{lim\; sup}\frac{\mathrm{\Sigma }(\rho )\mathrm{\Sigma }(0)}{\rho }\underset{\eta ^0๐(0)}{inf}\left[\frac{Q^{}(0)}{p+1}_^n\frac{|\eta _1^0|^{\frac{p+1}{p}}}{Q^{\frac{p+1}{p}}(0)}\frac{K^{}(0)}{q+1}_^n\frac{|\eta _2^0|^{\frac{q+1}{q}}}{K^{\frac{q+1}{q}}(0)}\right].$$
Moreover, similarly, we get
$`\mathrm{\Sigma }(\rho )\mathrm{\Sigma }(0)`$ $`I_\rho (\vartheta (\rho ,\rho )\eta ^\rho )I_0(\vartheta (0,\rho )\eta ^\rho )`$
$`=\rho _\xi I_\xi (\vartheta (\xi ,\rho )\eta ^\rho )|_{\xi =\mu [0,\rho ]},`$
so that, by exploiting (2.14) and (2.15), we conclude
$$\underset{\rho 0^+}{lim\; inf}\frac{\mathrm{\Sigma }(\rho )\mathrm{\Sigma }(0)}{\rho }\underset{\eta ^0๐(0)}{inf}\left[\frac{Q^{}(0)}{p+1}_^n\frac{|\eta _1^0|^{\frac{p+1}{p}}}{Q^{\frac{p+1}{p}}(0)}\frac{K^{}(0)}{q+1}_^n\frac{|\eta _2^0|^{\frac{q+1}{q}}}{K^{\frac{q+1}{q}}(0)}\right].$$
Then the desired formula for the right derivative of $`\mathrm{\Sigma }`$ follows. A very similar argument provides the corresponding formula for the left derivative.
###### Remark 2.1
Nowadays, the further regularity of $`\mathrm{\Sigma }`$ is, to our knowledge, an open problem. Actually, not even in the case of a single equation the situation is very well understood. For instance, on one hand, if we consider the problem
$$\epsilon ^2\mathrm{\Delta }u+V(x)u=K(x)u^p\text{in }^n,u>0\text{in }^n,$$
then $`\mathrm{\Sigma }C^m(^n)`$ provided that both the potentials $`V`$ and $`K`$ belong to $`C^m(^n)`$, with $`m1`$. On the other hand, if $`f`$ is not a power (and does not satisfy conditions ensuring uniqueness up to translations), for the equation
$$\epsilon ^2\mathrm{\Delta }u+V(x)u=K(x)f(u)\text{in }^n,u>0\text{in }^n,$$
we do not know which regularity beyond $`\mathrm{Lip}_{\mathrm{loc}}`$ can be achieved by $`\mathrm{\Sigma }`$. Even though we do not possess any specific counterexample, our feeling is that there exist functions $`f`$ for which the associated $`\mathrm{\Sigma }`$ fails to be $`C^1`$ smooth. It is evident by the (left and right) derivative formulas of $`\mathrm{\Sigma }`$ that its further regularity is related to the uniqueness of positive radial solutions to $`\mathrm{\Delta }u+u=f(u)`$ in $`^n`$, $`u>0`$ in $`^n`$, which occurs just for very particular nonlinearities $`f`$. Based upon these considerations, for semilinear systems, the further regularity of $`\mathrm{\Sigma }`$ seems an ever harder matter, since as already stressed nothing is known, so far, about the uniqueness of solutions to the system
$$\mathrm{\Delta }u+u=f(v),\mathrm{\Delta }v+v=g(u),\text{in }^n,u,v>0\text{in }^n,$$
not even with the particular choices $`f(v)=v^q`$ and $`g(u)=u^p`$.
## 3 Proof of the results
### 3.1 Proof of Theorem 1.1
Let $`z`$ and let $`(u_{\epsilon _h},v_{\epsilon _h})W^{2,\frac{q+1}{q}}(^n)\times W^{2,\frac{p+1}{p}}(^n)`$ be a corresponding a sequence of strong solutions to ($`S_\epsilon `$) with $`|u_{\epsilon _h}(z)|,|v_{\epsilon _h}(z)|\delta `$ for some $`\delta >0`$, $`|u_{\epsilon _h}(z+\epsilon _hx)|0,|v_{\epsilon _h}(z+\epsilon _hx)|0`$ as $`|x|\mathrm{}`$ uniformly w.r.t. $`h`$, and $`\epsilon _{h}^{}{}_{}{}^{n}f_{\epsilon _h}(u_{\epsilon _h},v_{\epsilon _h})\mathrm{\Sigma }(z)`$ as $`h\mathrm{}`$. Let us set:
$$\phi _h(x)=u_{\epsilon _h}(z+\epsilon _hx)\text{and}\psi _h(x)=v_{\epsilon _h}(z+\epsilon _hx),$$
for all $`h1`$. Then, since $`(u_{\epsilon _h},v_{\epsilon _h})`$ is a solution ($`S_\epsilon `$), $`(\phi _h,\psi _h)`$ is solution of
$$\mathrm{\Delta }\phi _h+\phi _h=K(z+\epsilon _hx)\psi _h^q,\mathrm{\Delta }\psi _h+\psi _h=Q(z+\epsilon _hx)\phi _h^p.$$
(3.1)
By arguing as in the proof of Theorem 2.2, it is readily proved that, up to a subsequence, $`(\phi _h)`$ and $`(\psi _h)`$ converge weakly in $`W^{2,(q+1)/q}(^n)\times W^{2,(p+1)/p}(^n)`$ to some $`\phi _0`$ and $`\psi _0`$ respectively. Let us now prove that there exist $`\mathrm{\Theta }>0`$, $`\rho >0`$ and $`h_01`$ such that
$$\phi _h(x)ce^{\mathrm{\Theta }|x|}\text{and}\psi _h(x)ce^{\mathrm{\Theta }|x|},\text{for all }|x|\rho \text{ and }hh_0.$$
(3.2)
We follow the line of . Since $`z`$, then the functions $`\phi _h`$ and $`\psi _h`$ decay to zero at infinity, uniformly with respect to $`h`$. Hence, since $`p,q>1`$, we can find $`\rho >0`$, $`\mathrm{\Theta }>0`$ and $`h_01`$ such that
$`K(z+\epsilon _hx)\psi _h^q`$ $`(1\mathrm{\Theta }^2)\psi _h,`$
$`Q(z+\epsilon _hx)\phi _h^p`$ $`(1\mathrm{\Theta }^2)\phi _h,`$
for all $`|x|>\rho `$ and $`hh_0`$. Let us set
$$\xi (x)=\mu e^{\mathrm{\Theta }(|x|\rho )},\mu =\underset{|x|=\rho }{\mathrm{max}}\underset{hh_0}{\mathrm{max}}(\psi _h+\phi _h),$$
and introduce the set
$$A=\underset{R>\rho }{}D_R,$$
where, for any $`R>\rho `$, we put
$$D_R=\{\rho <|x|<R:\psi _h(x)+\phi _h(x)>\xi (x)\text{for some }hh_0\}.$$
If $`A=\mathrm{}`$, we are done. Instead, if $`A`$ is nonempty, there exists $`R_{}>\rho `$ such that
$`\mathrm{\Delta }(\xi \psi _h\phi _h)`$ $`\left[\mathrm{\Theta }^2{\displaystyle \frac{\mathrm{\Theta }(n1)}{|x|}}\right]\xi (x)\mathrm{\Theta }^2\psi _h\mathrm{\Theta }^2\phi _h`$
$`\mathrm{\Theta }^2(\xi \psi _h\phi _h)<0,\text{on }D_R\text{ for all }RR_{}.`$
Hence, by the maximum principle, since $`(\xi \psi _h\phi _h)|_{\{|x|=\rho \}}0`$, we get
$$\xi \psi _h\phi _h\mathrm{min}\{0,\underset{|x|=R}{\mathrm{min}}(\xi \psi _h\phi _h)\},\text{for all }RR_{}$$
so that, letting $`R\mathrm{}`$, yields, for any $`\rho >0`$, $`\psi _h(x)+\phi _h(x)\xi (x)`$ for $`|x|>\rho `$, which contradicts the definition of $`D_R_{}\mathrm{}`$.
By virtue of the Schauder interior estimates (see e.g. ), $`\phi _h\phi _0`$ and $`\psi _h\psi _0`$ locally in $`C^2`$ sense, so that $`(\phi _0,\psi _0)`$ is a (nontrivial, radial, decaying) solution to ($`S_z`$). Moreover, in light of the exponential barriers provided by (3.2), since $`z`$, it is not difficult to see that $`(\phi _0,\psi _0)๐(z)`$, for we have
$`\mathrm{\Sigma }(z)`$ $`=\left(\frac{1}{2}\frac{1}{q+1}\right){\displaystyle _^n}K(z)|\psi _0|^{q+1}+\left(\frac{1}{2}\frac{1}{p+1}\right){\displaystyle _^n}Q(z)|\phi _0|^{p+1}`$
$`=f_z(\phi _0,\psi _0)=I_z(\eta ^0),`$
where $`\eta ^0`$ is the dual solution corresponding to $`(\phi _0,\psi _0)`$.
Let us now consider the Lagrangian $`:^n\times \times \times ^n\times ^n`$ defined as
$$(x,s_1,s_2,\xi _1,\xi _2)=\xi _1\xi _2+s_1s_2\frac{1}{q+1}K(z+\epsilon _hx)s_2^{q+1}\frac{1}{p+1}Q(z+\epsilon _hx)s_1^{p+1}.$$
Then system (3.1) rewrites as
$$\{\begin{array}{cc}\mathrm{div}\left(_{\xi _2}(x,\phi _h,\psi _h,\phi _h,\psi _h)\right)+_{s_2}(x,\phi _h,\psi _h,\phi _h,\psi _h)=0,\hfill & \text{in }^n,\hfill \\ \mathrm{div}\left(_{\xi _1}(x,\phi _h,\psi _h,\phi _h,\psi _h)\right)+_{s_1}(x,\phi _h,\psi _h,\phi _h,\psi _h)=0,\hfill & \text{in }^n,\hfill \\ \multicolumn{2}{c}{}\\ \phi _h,\psi _h>0,\hfill & \text{in }^n.\hfill \end{array}$$
By the Pucci-Serrin identity for systems \[21, see ยง5\], we have
$$\underset{i,l=1}{\overset{n}{}}_^n_i๐^l_{(\xi _2)_i}(x,\phi _h,\psi _h,\phi _h,\psi _h)_l\psi _h$$
$$+\underset{i,l=1}{\overset{n}{}}_^n_i๐^l_{(\xi _1)_i}(x,\phi _h,\psi _h,\phi _h,\psi _h)_l\phi _h$$
$$=_^n\left[(\mathrm{div}๐)(x,\phi _h,\psi _h,\phi _h,\psi _h)+๐_x(x,\phi _h,\psi _h,\phi _h,\psi _h)\right],$$
for all $`๐C_\mathrm{c}^1(^n,^n)`$. Let us take, for $`\lambda >0`$,
$$๐(x)=(\mathrm{{\rm Y}}(\lambda x),\mathrm{\hspace{0.17em}0},\mathrm{},\mathrm{\hspace{0.17em}0}),$$
and $`\mathrm{{\rm Y}}C_\mathrm{c}^1(^n)`$ such that $`\mathrm{{\rm Y}}(x)=1`$ if $`|x|1`$ and $`\mathrm{{\rm Y}}(x)=0`$ if $`|x|2`$. Then,
$$\underset{i=1}{\overset{n}{}}_^n\lambda _i\mathrm{{\rm Y}}(\lambda x)_i\phi _h_1\psi _h+\underset{i=1}{\overset{n}{}}_^n\lambda _i\mathrm{{\rm Y}}(\lambda x)_i\psi _h_1\phi _h$$
$$=_^n\lambda _1\mathrm{{\rm Y}}(\lambda x)(x,\phi _h,\psi _h,\phi _h,\psi _h)$$
$$+_^n\epsilon _h\mathrm{{\rm Y}}(\lambda x)\left[\frac{1}{q+1}_1K(z+\epsilon _hx)\psi _h^{q+1}\frac{1}{p+1}_1Q(z+\epsilon _hx)\phi _h^{p+1}\right].$$
By the arbitrariness of $`\lambda >0`$, letting $`\lambda 0`$ and keeping $`h`$ fixed, we obtain
$$_^n\left[\frac{1}{q+1}_1K(z+\epsilon _hx)\psi _h^{q+1}\frac{1}{p+1}_1Q(z+\epsilon _hx)\phi _h^{p+1}\right]=0.$$
Therefore, letting now $`h\mathrm{}`$, since in light of (1.3) we get
$$|K(z+\epsilon _hx)|,|Q(z+\epsilon _hx)|ce^{M\epsilon _h|x|},\text{for }|x|\text{ large,}$$
by virtue of (3.2), there holds
$$_^n\left[\frac{1}{q+1}_1K(z)\psi _0^{q+1}\frac{1}{p+1}_1Q(z)\phi _0^{p+1}\right]=0.$$
Analogously, we can show that, for all $`w^n`$,
$$_^n\left[\frac{1}{q+1}K(z)\psi _0^{q+1}\frac{1}{p+1}Q(z)\phi _0^{p+1}\right]w=0.$$
Hence
$$\frac{1}{p+1}\frac{Q}{w}(z)_^n\frac{|\eta _1^0|^{\frac{p+1}{p}}}{Q^{\frac{p+1}{p}}(z)}\frac{1}{q+1}\frac{K}{w}(z)_^n\frac{|\eta _2^0|^{\frac{q+1}{q}}}{K^{\frac{q+1}{q}}(z)}=0$$
(3.3)
Since $`\eta ^0๐(z)`$, by Theorem 2.2 we have
$`\left({\displaystyle \frac{\mathrm{\Sigma }}{w}}\right)^+(z)`$ $`=\underset{\eta ๐(z)}{inf}_zI_z(\eta )w`$
$`\frac{1}{p+1}{\displaystyle \frac{Q}{w}}(z){\displaystyle _^n}{\displaystyle \frac{|\eta _1^0|^{\frac{p+1}{p}}}{Q^{\frac{p+1}{p}}(z)}}\frac{1}{q+1}{\displaystyle \frac{K}{w}}(z){\displaystyle _^n}{\displaystyle \frac{|\eta _2^0|^{\frac{q+1}{q}}}{K^{\frac{q+1}{q}}(z)}}=0.`$
Then, by the very definition of $`(\mathrm{\Sigma })^0(z;w)`$ (see Definition 1.1), we get
$$(\mathrm{\Sigma })^0(z;w)\left(\frac{(\mathrm{\Sigma })}{w}\right)^+(z)0,\text{for every }w^n.$$
Then $`0_C(\mathrm{\Sigma })(z)`$ and, since $`_C(\mathrm{\Sigma })(z)=_C\mathrm{\Sigma }(z)`$ (cf. ), we obtain $`z๐ฆ`$.
### 3.2 Proof of Corollary 1.1
It suffices to combine Theorems 1.1 and 2.2, taking into account what discussed in Section 2.2 about the conjectured explicit representation formula for $`\mathrm{\Sigma }`$.
### 3.3 Proof of Theorem 1.2
Let $`m1`$ and $`z_m`$. The assertion follows by mimicking the various steps in the proof of Theorem 1.1 with $`_m`$ in place of $``$, and combining formula (3.3) with the definitions of $`\mathrm{\Gamma }_{z,m}^{}`$ and $`๐ฆ_m`$, taking into account that $`\eta ^0๐พ_m(z)`$, as it holds $`I_z(\eta ^0)=m`$, being $`\eta ^0`$ the strong limit of $`\eta ^{\epsilon _j}`$. Indeed, by (3.3), there holds
$$\mathrm{\Gamma }_{z,m}^+(w)0,w^n,\mathrm{\Gamma }_{z,m}^{}(w)0,w^n,$$
so that $`0\mathrm{\Gamma }_{z,m}^+(0)\mathrm{\Gamma }_{z,m}^{}(0)`$, yielding $`z๐ฆ_m`$.
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# Massive gauge field theory without Higgs mechanism
## I References
J. C. Su, Nuovo Cimento 117B (2002) 203.
C. N. Yang and R. L. Mills, Phys. Rev. 96 (1954) 191.
C. Itzykson and F-B, Zuber, Quantum Field Theory, McGraw-Hill, New York (1980).
L. D. Faddeev and V. N. Popov, Phys. Lett. B25 (1967) 29.
L. D. Faddeev, Theor. Math. Phys., 1 (1970) 1.
J. C. Su, hep.th/9805192; 9805193; 9805194.
H. Umezawa and S. Kamefuchi, Nucl. Phys. 23 (1961) 399.
A. Salam, Nucl. Phys. 18 (1960) 681; Phys. Rev. 127 (1962) 331.
D. G. Boulware, Ann. Phys. 56 (1970) 140.
P. Senjanovic, Ann. Phys. (N.Y.) 100 (1976) 227.
C. Grosse-Knetter, Phys. Rev. D48 (1993) 2854.
N. Banerjee, R. Banerjee and G. Subir, Ann. Phys. (N. Y.) 241 (1995) 237.
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# Exact soliton solution and inelastic two-soliton collision in spin chain driven by a time-dependent magnetic field
## I Introduction
Over the past three decades, an enormous amount of literature has appeared throughout soliton physics and the underlying completely integrable models. The classical Heisenberg spin chain which exhibits both coherent and chaotic structures depending on the nature of the magnetic interactions has attracted considerable attentions in nonlinear science and condensed-matter physics. Solitons in quasi one-dimensional magnetic systems have already been probed experimentally by neutron inelastic scattering , nuclear magnetic resonance , Mossbauer linewidth measurements , and electron spin resonance . The corresponding theoretical studies are based usually on the Landau-Lifshitz equation . The isotropic spin chain has been studied in various aspects and the construction of soliton solutions of Landau-Lifschitz equation with an easy axis has been also discussed. It is demonstrated that the inverse scattering transformation can be used to solve the Landau-Lifschitz equation for an anisotropic spin chain. Great efforts have been devoted to construct the soliton solution which are found by means of the Darboux transformation . The continuum spin chain in an external magnetic field is of great interest and multi-soliton solutions of Landau-Lifschitz equation for an isotropic spin chain have been reported . Using Darboux transformation the nonlinear dynamics of anisotropic Heisenberg spin chain in an external magnetic field is investigated and exact soliton solutions are obtained. Recently soliton interaction has been investigated. The main goal of this paper is to study the new effect of soliton-soliton interaction in spin chain driven by time oscillating magnetic field. We obtain exact solution of N-soliton trains in terms of an inverse scattering transformation. It is shown that inelastic collisions generally appear due to the time-varying field and the complete inelastic collisions which may lead to a interesting technique of soliton filter and switch can be achieved in special case.
The outline of this paper is organized as follows: In Sec. II the formalism obtained by an inverse scattering transformation is explained in detail and the general N-soliton solution for reflectionless case is obtained. Precession of nonlinear spin waves in the oscillating magnetic field is shown in Sec. III. Sec. IV is devoted to general two-soliton solution and soliton collisions. Finally, Sec. V will give our concluding remarks.
## II Exact solution of N-soliton train
Our starting Hamiltonian describing the spin chain in a time oscillating magnetic field with an arbitrary direction can be written as
$$\widehat{H}=J\underset{<n,n^{}>}{}\widehat{S}_n\widehat{S}_n^{}g\mu _B๐\left(t\right)\underset{n}{}\widehat{S}_n,$$
(1)
where $`\widehat{S}_n(\widehat{S}_n^x,\widehat{S}_n^y,\widehat{S}_n^z)`$ with $`n=1,2,\mathrm{},N`$ are spin operators, $`J>0`$ is the pair interaction parameter, $`g`$ the Lande factor and $`\mu _B`$ is the Bohr magneton, $`๐\left(t\right)=B\mathrm{cos}(\omega t)๐`$ is the external magnetic field with $`๐=(\mathrm{sin}\theta ,0,\mathrm{cos}\theta )`$ denoting the unit vector of field direction where chain axis and direction of magnetic field are assumed in x-z plane. The angle $`\theta `$ between direction of magnetic field and z-axis is arbitrary.
The equation of motion for the spin operator on the nth site is $`\frac{d}{dt}\widehat{S}_n=\frac{i}{\mathrm{}}[\widehat{S}_n,\widehat{H}]`$. At low temperature, the spin can be treated as a classical vector such that $`\widehat{S}_n๐(x)`$. So that the equation of motion in a continuum spin chain under a time-dependent magnetic field can be obtained as a Landau-Lifschitz type
$$\frac{}{t}๐=๐\times \left(\frac{^2}{x^2}๐+\epsilon \right),$$
(2)
with $`\epsilon =g\mu _B๐\left(t\right)/\left(2J\right)`$, where $`๐(x,t)=(S^x(x,t)`$, $`S^y(x,t)`$, $`S^z(x,t))`$. We set the length of the spin vector to unit for the sake of simplicity $`๐^2(x,t)=1`$. The dimensionless time $`t`$ and coordinate $`x`$ in Eq. (2) are scaled in unit $`\frac{1}{2J}`$ and $`d`$ respectively, where $`d`$ denotes the lattice constant.
The corresponding Lax equations for the equation of motion (2) are written as
$`{\displaystyle \frac{}{x}}\mathrm{\Psi }(x,t,\lambda )`$ $`=`$ $`L\left(\lambda \right)\mathrm{\Psi }(x,t,\lambda ),`$ (3)
$`{\displaystyle \frac{}{t}}\mathrm{\Psi }(x,t,\lambda )`$ $`=`$ $`M\left(\lambda \right)\mathrm{\Psi }(x,t,\lambda ),`$ (4)
where $`\lambda `$ is the spectral parameter, $`\mathrm{\Psi }(x,t,\lambda )`$ is eigenfunction corresponding to $`\lambda `$, and $`L`$ and $`M`$ are given in the form
$`L`$ $`=`$ $`i\lambda (๐\sigma ),`$ (5)
$`M`$ $`=`$ $`{\displaystyle \frac{i}{2}}(\epsilon \sigma )+i2\lambda ^2(๐\sigma )\lambda (๐\sigma )({\displaystyle \frac{}{x}}๐\sigma ).`$ (6)
Here $`\sigma `$ is Pauli matrix. Thus Eq. (2) can be recovered from the compatibility condition $`\frac{}{t}L\frac{}{x}M+[L,M]=0`$. Based on the Lax equations (4), we derive the exact $`N`$-soliton solution by employing the inverse scattering transformation. We consider the following natural boundary condition of initial time($`t=0`$), $`๐\left(x\right)(S^x,S^y,S^z)(\mathrm{sin}\theta ,0,\mathrm{cos}\theta )\text{ as}\left|x\right|\mathrm{}`$, namely, the spin vector is along the field direction. We then have the asymptotic form of Eq. (4) at $`\left|x\right|\mathrm{}`$,
$$_xE(x,\lambda )=L_0(\lambda )E(x,\lambda ),$$
(7)
where
$$E(x,\lambda )=Ue^{i\lambda x\sigma _3},L_0(\lambda )=i\lambda U_0,$$
(8)
and
$$U_0=\left(\begin{array}{cc}\mathrm{cos}\theta & \mathrm{sin}\theta \\ \mathrm{sin}\theta & \mathrm{cos}\theta \end{array}\right),U=\left(\begin{array}{cc}1& \mathrm{tan}\frac{\theta }{2}\\ \mathrm{tan}\frac{\theta }{2}& 1\end{array}\right).$$
(9)
The Jost solutions $`\mathrm{\Psi }_+(x,\lambda )`$ and $`\mathrm{\Psi }_{}(x,\lambda )`$ of Eq. (4) are defined as
$`\mathrm{\Psi }_+(x,\lambda )`$ $``$ $`E(x,\lambda )\text{ as }x\mathrm{},`$ (10)
$`\mathrm{\Psi }_{}(x,\lambda )`$ $``$ $`E(x,\lambda )\text{ as }x\mathrm{}.`$ (11)
With standard procedures, one finds the following integral representations of the Jost solutions in terms of the integration kernels $`K`$ and $`N`$ to be determined,
$`\mathrm{\Psi }_+(x,\lambda )=Ue^{i\lambda x\sigma _3}+\lambda {\displaystyle _x^{\mathrm{}}}๐yK(x,y)Ue^{i\lambda y\sigma _3},`$ (12)
$`K(x,\mathrm{})=0,K(x,y)=0\text{as}y<x.`$ (13)
and
$`\mathrm{\Psi }_{}(x,\lambda )=Ue^{i\lambda x\sigma _3}+\lambda {\displaystyle _{\mathrm{}}^x}๐yN(x,y)Ue^{i\lambda y\sigma _3},`$ (14)
$`N(x,\mathrm{})=0,N(x,y)=0\text{as}y<x.`$ (15)
where $`K`$ and $`N`$ are $`2\times 2`$ matrices. Substituting $`\mathrm{\Psi }_+(x,\lambda )`$ in Eq. (13) into Eq.(4) and noting $`U\sigma _3U^1=U_0`$, we obtain
$$๐\sigma =\left[IiK(x,x)U_0\right]U_0\left[IiK(x,x)U_0\right]^1$$
(16)
where $`I`$ is unit matrix. It is obvious that Eq. (16) gives rise to a relation between kernel $`K`$ and spin vector $`๐`$ to be obtained.
The scattering data for the operator $`L(x,\lambda )`$ are the set $`s=\{a(\lambda ),b(\lambda );\lambda _n,c_n,Im\lambda >0,n=1,\mathrm{},N\}`$, where $`|a(\lambda )|^2+|b(\lambda )|^2=1`$, and the function $`a(\lambda )`$ can be analytically continued to the half-plane $`Im\lambda >0`$. The discrete eigenvalues, $`\lambda _n`$, for the operator $`L(x,\lambda )`$ are zeroes of $`a(\lambda )`$ such that $`a(\lambda _n)=0`$ (for the simplicity we consider only simple zeroes). The functions $`a(\lambda )`$ and $`b(\lambda )`$ are seen to be transmission and reflection coefficients of the operator $`L`$ respectively. The parameter $`c_n`$ denotes the asymptotic characteristics of the eigenfunctions.
The time-dependence of the scattering data $`s\left(t\right)`$ can be obtained from the second Lax equation (4),
$`a(\lambda ,t)`$ $`=`$ $`a(\lambda ,0),`$ (17)
$`b(\lambda ,t)`$ $`=`$ $`\mathrm{exp}\left(4i\lambda ^2ti{\displaystyle \frac{g\mu _BB\mathrm{sin}\omega t}{J\omega }}\right)b(\lambda ,0),`$ (18)
$`\lambda _n(t)`$ $`=`$ $`\lambda _n(0),`$ (19)
$`c_n(t)`$ $`=`$ $`\mathrm{exp}\left(4i\lambda _n^2ti{\displaystyle \frac{g\mu _BB\mathrm{sin}\omega t}{J\omega }}\right)c_n(0).`$ (20)
where $`c_n(0)`$, $`b(\lambda ,0)`$ and $`a(\lambda ,0)`$ are constants determined by initial conditions. The Gelfand-Levitan-Marchenko equation establishes a relation between the kernel $`K(x,y,t)`$ and the scattering data $`s\left(t\right)`$ and has the form
$$K(x,y,t)U\left(\begin{array}{c}1\hfill \\ 0\hfill \end{array}\right)+F_1+\frac{1}{2\pi }_{\mathrm{}}^{\mathrm{}}\lambda ^1r\left(\lambda \right)F_2๐\lambda =0,$$
(21)
as $`y>x`$, where $`r\left(\lambda \right)=b\left(\lambda \right)/a\left(\lambda \right)`$ and
$`F_1=U\left(\begin{array}{c}0\hfill \\ 1\hfill \end{array}\right){\displaystyle \underset{n=1}{\overset{N}{}}}{\displaystyle \frac{c_n\left(t\right)}{\lambda _n}}e^{i\lambda _n\left(x+y\right)}`$ (24)
$`+{\displaystyle _x^{\mathrm{}}}K(x,z,t)U\left(\begin{array}{c}0\hfill \\ 1\hfill \end{array}\right){\displaystyle \underset{n=1}{\overset{N}{}}}c_n\left(t\right)e^{i\lambda _n\left(y+z\right)}dz,`$ (27)
$`F_2=U\left(\begin{array}{c}0\hfill \\ 1\hfill \end{array}\right)e^{i\lambda x}+\lambda {\displaystyle _x^{\mathrm{}}}K(x,z,t)U\left(\begin{array}{c}0\hfill \\ 1\hfill \end{array}\right)e^{i\lambda z}๐z.`$ (32)
For the reflectionless case, $`r\left(\lambda \right)=0`$, Eq. (21) becomes a set of algebraic equations and after tedious calculation the matrix elements of the kernel $`K`$ are obtained as
$`K_{11}(x,x,t)`$ $`=`$ $`\mathrm{cos}^2{\displaystyle \frac{\theta }{2}}[B_1+B_2\mathrm{tan}{\displaystyle \frac{\theta }{2}}],`$ (33)
$`K_{12}(x,x,t)`$ $`=`$ $`\mathrm{cos}^2{\displaystyle \frac{\theta }{2}}[B_1\mathrm{tan}{\displaystyle \frac{\theta }{2}}B_2].`$ (34)
with
$`B_1`$ $`=`$ $`{\displaystyle \frac{det[I+G^{\prime \prime }G^{}+D^T(C\mathrm{tan}\frac{\theta }{2}\overline{C}G^{})]}{det(I+G^{}G^{\prime \prime })}}1,`$ (35)
$`B_2`$ $`=`$ $`{\displaystyle \frac{det[I+G^{}G^{\prime \prime }\overline{D}^T(\overline{C}+CG^{\prime \prime }\mathrm{tan}\frac{\theta }{2})]}{det(I+G^{}G^{\prime \prime })}}1.`$ (36)
where $`C(x,t)`$, $`C^{}(x,t)`$, $`D\left(x\right)`$ are $`1\times N`$ matrices, $`G^{}(x,t)`$, $`G^{\prime \prime }(x,t)`$ $`N\times N`$ matrices, respectively. The superscript $`T`$ means the transposed matrix and the overbar denotes complex conjugate,
$`C(x,t)_n`$ $`=`$ $`c_n\left(t\right)\lambda _n^1D\left(x\right)_n,`$ (37)
$`C^{}(x,t)_n`$ $`=`$ $`c_n\left(t\right)D\left(x\right)_n,`$ (38)
$`D(x)_n`$ $`=`$ $`\mathrm{exp}(i\lambda _nx),`$ (39)
$`G^{}(x,t)_{nm}`$ $`=`$ $`{\displaystyle \frac{1}{i\left(\overline{\lambda }_n\lambda _m\right)}}\overline{D\left(x\right)}_nC^{}(x,t)_m,`$ (40)
$`G^{\prime \prime }(x,t)_{nm}`$ $`=`$ $`{\displaystyle \frac{1}{i\left(\lambda _n\overline{\lambda }_m\right)}}D\left(x\right)_n\overline{C^{}(x,t)}_m.`$ (41)
Substituting Eq. (34) into Eq. (16) , we obtain the general form of N-soliton trains,
$`S^x`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{\Delta }}}Re\{i2K_{12}[1iK_{11}\mathrm{cos}\theta ]`$ (43)
$`+[1+K_{11}^2K_{12}^2]\mathrm{sin}\theta \},`$
$`S^y`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{\Delta }}}Im\{i2K_{12}[1iK_{11}\mathrm{cos}\theta ]`$ (45)
$`+[1+K_{11}^2K_{12}^2]\mathrm{sin}\theta \},`$
$`S^z`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{\Delta }}}\{[1+|K_{11}|^2|K_{12}|^2]\mathrm{cos}\theta `$ (47)
$`+2Im\left[K_{11}(1+i\overline{K_{12}}\mathrm{sin}\theta )\right]\}.`$
where $`\overline{K_{12}}`$ is the complex conjugate of $`K_{12}`$.
$`\mathrm{\Delta }`$ $`=`$ $`\left|1i\left[K_{11}\mathrm{cos}\theta +K_{12}\mathrm{sin}\theta \right]\right|^2`$ (49)
$`+\left|K_{11}\mathrm{sin}\theta K_{12}\mathrm{cos}\theta \right|^2.`$
According to exact N-soliton solutions in Eq. (47), we, generally speaking, can investigate the dynamics of soliton trains and soliton interaction. The neighboring solitons may repulse or attract each other with a force depending on their phase difference. Particularly we in the following shall concentrate on the analyses of one-soliton dynamics and two-soliton collisions which may be of more interest.
## III One-soliton dynamics and Spin precession with time varying amplitude
When $`N=1`$, from Eq. (34) and Eq. (47) we obtain the general form of the exact one-soliton solution as follows
$`S^x`$ $`=`$ $`{\displaystyle \frac{R_1}{|\lambda _1|^4\mathrm{cosh}^2\mathrm{\Theta }_1}},`$ (50)
$`S^y`$ $`=`$ $`{\displaystyle \frac{R_2}{|\lambda _1|^4\mathrm{cosh}^2\mathrm{\Theta }_1}},`$ (51)
$`S^z`$ $`=`$ $`R_3\mathrm{cos}\theta +R_4\mathrm{sin}\theta ,`$ (52)
where
$`R_1`$ $`=`$ $`[|\lambda _1|^4\mathrm{cosh}^2\mathrm{\Theta }_1+\beta _1^2(\alpha _1^2\beta _1^2\mathrm{cos}2\theta )e^{2\mathrm{\Theta }_1}]\mathrm{sin}\theta `$ (57)
$`+\beta _1^2|\lambda _1|^2(2\mathrm{cos}^2\theta \mathrm{sin}^2\mathrm{\Phi }_11)\mathrm{sin}\theta `$
$`+2\beta _1^2|\lambda _1|(2\beta _1\mathrm{sin}\mathrm{\Phi }_1\mathrm{sin}^2\theta +\alpha _1\mathrm{cos}\mathrm{\Phi }_1)e^{\mathrm{\Theta }_1}\mathrm{cos}\theta `$
$`2\beta _1(\beta _1e^{\mathrm{\Theta }_1}\mathrm{sin}\theta +|\lambda _1|\mathrm{sin}\mathrm{\Phi }_1\mathrm{cos}\theta )[|\lambda _1|^2\mathrm{cosh}\mathrm{\Theta }_1`$
$`+\beta _1(|\lambda _1|\mathrm{sin}\mathrm{\Phi }_1\mathrm{sin}\theta \beta _1e^{\mathrm{\Theta }_1}\mathrm{cos}\theta )\mathrm{cos}\theta ],`$
$`R_2`$ $`=`$ $`2\alpha _1\beta _1^2|\lambda _1|e^{\mathrm{\Theta }_1}\mathrm{sin}\mathrm{\Phi }_1+2\beta _1|\lambda _1|^3\mathrm{cos}\mathrm{\Phi }_1\mathrm{cosh}\mathrm{\Theta }_1`$ (59)
$`2\beta _1^3|\lambda _1|\left(\mathrm{sin}\theta +\mathrm{cos}^2\theta \right)e^{\mathrm{\Theta }_1}\mathrm{cos}\mathrm{\Phi }_1,`$
$`R_3`$ $`=`$ $`1{\displaystyle \frac{2\beta _1^2}{|\lambda _1|^2\mathrm{cosh}^2\mathrm{\Theta }_1}},`$ (60)
$`R_4`$ $`=`$ $`{\displaystyle \frac{1}{|\lambda _1|^2\mathrm{cosh}^2\mathrm{\Theta }_1}}[2\beta _1^2\mathrm{cos}(\mathrm{\Phi }_1\varphi _1)\mathrm{sinh}\mathrm{\Theta }_1`$ (62)
$`+2\alpha _1\beta _1\mathrm{sin}(\mathrm{\Phi }_1\varphi _1)\mathrm{cosh}\mathrm{\Theta }_1]`$
with
$`\mathrm{\Theta }_1`$ $`=`$ $`2\beta _1(xV_1t)x_1,`$ (63)
$`V_1`$ $`=`$ $`4\alpha _1,x_1=\mathrm{ln}[(2\beta _1)^1c_1\left(0\right)],`$ (64)
$`\mathrm{\Phi }_1`$ $`=`$ $`2\alpha _1x4(\alpha _1^2\beta _1^2)t\left(\omega J\right)^1g\mu _BB\mathrm{sin}(\omega t)\varphi _1,`$ (65)
$`\mathrm{\Omega }_1`$ $`=`$ $`2\alpha _1^1\left(\alpha _1^2\beta _1^2\right)+\mathrm{\Omega }_B,`$ (66)
$`\mathrm{\Omega }_B`$ $`=`$ $`\left(2\alpha _1J\right)^1g\mu _BB\mathrm{cos}\omega t,`$ (67)
$`\varphi _1=\mathrm{arg}\lambda _1`$, $`\lambda _1=\alpha _1+i\beta _1`$ is eigenvalue parameter. The solution (52) describes a spin precession around magnetic field direction characterized by four real parameters: velocity $`V_1`$, frequency $`\mathrm{\Omega }_1`$, coordinate of the center of the solitary wave $`x_1`$ and initial phase $`\varphi _1`$. The center of solitary wave moves with a velocity $`V_1`$, while the wave depth and width vary periodically with time. The wave shape is modulated periodically by frequency $`\mathrm{\Omega }_1`$ depending on magnetic field. Therefore, the solution (52) can not be written as the form of separating variables. Amplitude $`A`$ and phase $`\mathrm{\Phi }_1`$ are complicated functions of $`J`$, $`B`$, $`\omega `$ and $`\lambda _1`$. When $`\alpha _1=\beta _1`$, the frequency $`\mathrm{\Omega }_1`$ depends on magnetic field only, and we have $`\mathrm{\Omega }_1=\mathrm{\Omega }_B`$. If $`\alpha _1=\beta _1`$ and $`B=0`$, the solution (52) reduces to the usual soliton without shape changing. Therefore, we can use magnetic field to adjust spin precession and the wave shape as well.
For a special case, $`\theta =0`$, namely the magnetic field is along the z-axis, $`S^z`$ is independent of magnetic field, $`S^z=R_3`$, while $`S^x`$ and $`S^y`$ precess around magnetic field (z-axis). The precession frequency $`\mathrm{\Omega }_1`$ is determined by magnetic field. As magnetic field rotates from $`\theta =0`$ (z-axis) to $`\theta =\pi /2`$ (x-axis), we can find the correspondence such that $`S^xS^z`$, $`S^yS^y`$, $`S^zS^x`$. The three components of spin vector satisfy โleft-hand ruleโ. When $`\theta =\pi /2`$, $`S^x`$ is independent of magnetic field, while $`S^y`$ and $`S^z`$ precess around magnetic field (x-axis). These results show that the magnetic field results in the motion of the center of solitary waves along the field direction and the spin vector rotates around the field in any case.
## IV Two-soliton collision
When $`N=2`$, from Eq. (34) and Eq. (47) the general form of the exact two-soliton solution is seen to be
$`S^x`$ $`=`$ $`Re[i2Q_2(1iQ_1\mathrm{cos}\theta )+(1+Q_1^2Q_2^2)\mathrm{sin}\theta ],`$ (68)
$`S^y`$ $`=`$ $`Im[i2Q_2(1iQ_1\mathrm{cos}\theta )(1+Q_1^2Q_2^2)\mathrm{sin}\theta ],`$ (69)
$`S^z`$ $`=`$ $`(1+\left|Q_1\right|^2\left|Q_2\right|^2)\mathrm{cos}\theta `$ (71)
$`+2Im[Q_1(1+i\overline{Q_2}\mathrm{sin}\theta )].`$
where
$`Q_1`$ $`=`$ $`{\displaystyle \frac{\mathrm{cos}^2\frac{\theta }{2}}{W}}\{(f_1f_3)f_6+(f_2f_4)f_5`$ (73)
$`+\mathrm{tan}{\displaystyle \frac{\theta }{2}}[(\overline{f_1}\overline{f_3})f_8+(\overline{f_2}\overline{f_4})f_7]\},`$
$`Q_2`$ $`=`$ $`{\displaystyle \frac{\mathrm{cos}^2\frac{\theta }{2}}{W}}\{[(f_1f_3)f_6+(f_2f_4)f_5]\mathrm{tan}{\displaystyle \frac{\theta }{2}}`$ (75)
$`(\overline{f_1}\overline{f_3})f_8(\overline{f_2}\overline{f_4})f_7\},`$
with
$`f_1`$ $`=`$ $`1+\left|q_1\right|^2+\chi _1\overline{\chi }_2q_1\overline{q}_2,\text{ }f_2=1+\left|q_2\right|^2+\overline{\chi }_1\chi _2\overline{q}_1q_2,`$ (76)
$`f_3`$ $`=`$ $`\overline{\chi }_1\left|q_1\right|^2+\chi _1q_1\overline{q}_2,\text{ }f_4=\overline{\chi }_2\left|q_2\right|^2+\chi _2\overline{q}_1q_2,`$ (77)
$`f_5`$ $`=`$ $`\xi _1\left(q_1\mathrm{tan}{\displaystyle \frac{\theta }{2}}\left|q_1\right|^2\right)\chi _1\xi _2q_1\overline{q}_2,`$ (78)
$`f_6`$ $`=`$ $`\xi _2\left(q_2\mathrm{tan}{\displaystyle \frac{\theta }{2}}\left|q_2\right|^2\right)\chi _2\xi _1\overline{q}_1q_2,`$ (79)
$`f_7`$ $`=`$ $`\xi _1\left(\overline{q}_1+\left|q_1\right|^2\mathrm{tan}{\displaystyle \frac{\theta }{2}}\right)\xi _2\overline{\chi }_1\overline{q}_1q_2\mathrm{tan}{\displaystyle \frac{\theta }{2}},`$ (80)
$`f_8`$ $`=`$ $`\xi _2\left(\overline{q}_2+\left|q_2\right|^2\mathrm{tan}{\displaystyle \frac{\theta }{2}}\right)\xi _1\overline{\chi }_2q_1\overline{q}_2\mathrm{tan}{\displaystyle \frac{\theta }{2}},`$ (81)
$`\chi _1`$ $`=`$ $`{\displaystyle \frac{2\beta _1\lambda _1}{i\left(\lambda _1\overline{\lambda }_2\right)\left|\lambda _1\right|}},\chi _2={\displaystyle \frac{2\beta _2\lambda _2}{i\left(\lambda _2\overline{\lambda }_1\right)\left|\lambda _2\right|}},`$ (82)
$`W`$ $`=`$ $`f_1f_2f_3f_4,\text{ }q_j=e^{\mathrm{\Theta }_j+i\mathrm{\Phi }_j},\text{ }\xi _j=2\beta _j\left|\lambda _j\right|^1,`$ (83)
and
$`\mathrm{\Theta }_j`$ $`=`$ $`2\beta _j(xV_jt)x_j,`$ (84)
$`V_j`$ $`=`$ $`4\alpha _j,x_j=\mathrm{ln}[(2\beta _j)^1c_j\left(0\right)],`$ (85)
$`\mathrm{\Phi }_j`$ $`=`$ $`2\alpha _jx4(\alpha _j^2\beta _j^2)t\left(\omega J\right)^1g\mu _BB\mathrm{sin}(\omega t)\varphi _j,`$ (86)
$`\mathrm{\Omega }_j`$ $`=`$ $`2\alpha _j^1\left(\alpha _j^2\beta _j^2\right)+\mathrm{\Omega }_B,`$ (87)
$`\mathrm{\Omega }_B`$ $`=`$ $`\left(2\alpha _jJ\right)^1g\mu _BB\mathrm{cos}\omega t,`$ (88)
here $`\varphi _j=\mathrm{arg}\lambda _j`$ and $`\lambda _j=\alpha _j+i\beta _j`$ is eigenvalue parameter, $`j=1,2`$. The solutions (71) describe a general inelastic scattering process of two solitary waves with different center velocities $`V_1`$ and $`V_2`$, different shape variation frequencies $`\mathrm{\Omega }_1`$ and $`\mathrm{\Omega }_2`$. Before collision, they move towards each other, one with velocity $`V_1`$ and shape variation frequency $`\mathrm{\Omega }_1`$, the other with $`V_2`$ and $`\mathrm{\Omega }_2`$. The interaction potential between two solitons is a complicated function of parameters $`J,B,\omega `$ and $`\lambda _j`$. When $`\alpha _j=\beta _j`$, two-soliton shape-variation frequencies $`\mathrm{\Omega }_j(j=1,2)`$ are determined by magnetic field. In the case of $`B=0`$, the solutions (71) reduce to that of the usual two-soliton with two center velocities while without shape change where a interesting process in the absence of magnetic field is that the collision can result in the interchange of amplitude $`A_j`$ and phase $`\mathrm{\Phi }_j(j=1,2)`$ like exactly in the case of elastic collision of two particles.
In order to understand the nature of two-soliton interaction, we analyze asymptotic behavior of two-soliton solutions (71). Asymptotically, the two-soliton waves (71) can be written as a combination of two one-soliton waves (52) with different amplitude and phase. The formation of two-soliton waves in the corresponding limits $`x\mathrm{}`$ and $`x\mathrm{}`$ is similar to that of one-soliton waves (52). Analysis reveals that there is an amplitude exchange among three components $`S^x`$, $`S^y`$ and $`S^z`$ of each soliton during collision, which can be described by a transition matrix $`T_l^k`$ such that $`A_l^{k+}=A_l^kT_l^k`$, where the subscript $`l=1,2`$ respectively represents the first and the second soliton, $`k=x,y,z`$ denote three components of each soliton, the sign $`\pm `$ denotes the asymptotic limits of the corresponding amplitude, $`A_l^{k\pm }`$, at $`x\pm \mathrm{}`$. As a consequence, amplitude change of the three components $`S_1^k`$ of the first soliton from $`A_1^k`$ to $`A_1^{k+}`$ is given by square of transition matrices $`|T_1^k|^2`$ along with phase shift $`\delta \mathrm{\Phi }_1^k`$ during collision. In a similar fashion, the three components $`S_2^k`$ of the second soliton also change amplitudes from $`A_2^k`$ to $`A_2^{k+}`$ with a quantity $`|T_2^k|^2`$. The associate phase shift for the second soliton is $`\delta \mathrm{\Phi }_2^k`$. We also note a net change of separation distance between two solitons by $`\delta X_{12}`$.
For the special case $`|T_l^k|=1`$, which is possible only when $`\lambda _2=\overline{\lambda }_1`$, we have the standard elastic collision. For all other cases, we have the quantity $`|T_l^k|1`$, which corresponds to relative change among three components of the spin vector leading to the deformation of soliton shape. However, the total amplitude of individual solitons $`S_1`$ and $`S_2`$ is conserved quantity i.e., $`_l|A_l^{k\pm }|^2`$ is constant for $`l=1,2`$.
It is interesting to show the inelastic collision graphically. The general inelastic head on collision is explained in Fig. 1 from which it is seen that the amplitudes of $`S_1`$, $`S_2`$ are respectively suppressed and enhanced after collision. Fig. 2 are devoted to the complete inelastic head on collisions. The amplitudes of $`S_1`$and $`S_2`$ are respectively suppressed after collision shown in Fig. 2a and 2b. The complete inelastic overtake-collision is shown in Fig. 3 with the amplitudes of $`S_1`$and $`S_2`$ suppressed, respectively.
## V Conclusion
In terms of an inverse scattering transformation the exact solution of N-soliton trains in a spin chain driven by a time oscillating magnetic field is obtained. From the general solution the dynamics and soliton interactions are analyzed. The one-soliton solution gives rise explicitly to the spin precession along with the soliton shape variation induced by the time varying field. It is also shown that the time varying field leads generally to the inelastic and particularly the complete inelastic two-soliton collisions which may be useful in developing a soliton-shape control technique.
## VI Acknowledgment
This work was supported by the NSF of China under Grant Nos. 10194095, 90103024 and 10075032.
Figure caption
Fig. 1
Inelastic head on collision between two solitons โ profiles of z-component $`S^z(x;t)`$ of spin vector in Eq. (22) in spin chain under a time-dependent magnetic field showing two different dramatic scenarios of the shape changing collision, where $`\theta =\frac{\pi }{36}`$, $`\lambda _1=0.2+i0.45`$, $`\lambda _2=0.3+i0.65`$, $`c_1\left(0\right)=0.2`$, $`c_2\left(0\right)=3.5`$, $`g\mu _BB/J=0.01`$, $`\omega =10`$, $`V_1=0.8`$, $`V_2=1.2`$. All quantities plotted are dimensionless. The same is in Fig. 2 and 3.
Fig. 2
(2a) Complete inelastic head on collision expressed by Eq. (22) when $`S_1`$ suppressed, where $`\theta =0`$, $`\lambda _1=0.35+i0.4`$, $`\lambda _2=0.2i0.6`$, $`c_1\left(0\right)=0.2`$, $`c_2\left(0\right)=2.5`$, $`g\mu _BB/J=0.01`$, $`\omega =10`$, $`V_1=1.4`$, $`V_2=0.8`$.
(2b) Complete inelastic head on collision expressed by Eq. (22) when $`S_2`$ suppressed, where $`\theta =0`$, $`\lambda _1=0.35i0.4`$, $`\lambda _2=0.2+i0.6`$, $`c_1\left(0\right)=0.2`$, $`c_2\left(0\right)=2.5`$, $`g\mu _BB/J=0.01`$, $`\omega =10`$, $`V_1=1.4`$, $`V_2=0.8`$.
Fig. 3
(3a) Complete inelastic overtake-collision expressed by Eq. (22) when $`S_1`$ suppressed, where $`\theta =0`$, $`\lambda _1=0.55+i0.4`$, $`\lambda _2=0.1i0.45`$, $`c_1\left(0\right)=0.2`$, $`c_2\left(0\right)=2.5`$, $`g\mu _BB/J=0.01`$, $`\omega =10`$, $`V_1=2.2`$, $`V_2=0.4`$.
(3b) Complete inelastic overtake-collision expressed by Eq. (22) when $`S_1`$ suppressed, where $`\theta =0`$, $`\lambda _1=0.55i0.4`$, $`\lambda _2=0.1+i0.45`$, $`c_1\left(0\right)=0.2`$, $`c_2\left(0\right)=2.5`$, $`g\mu _BB/J=0.01`$, $`\omega =10`$,$`V_1=2.2`$, $`V_2=0.4`$.
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# Lepton asymmetry and primordial nucleosynthesis in the era of precision cosmology
## I Introduction
The successful calculation of the primordial light-element abundances in the framework of the standard theory of big-bang nucleosynthesis (BBN) Wagoner:1966pv ; Olive:1999ij ; Burles:2000zk ; Eidelman:2004wy is among the main pillars supporting our modern understanding of the universe. Apart from the cross sections that are needed for the network of nuclear reactions, these predictions depend on the baryon-to-photon ratio $`\eta _B`$, on the radiation density at the BBN epoch that is traditionally parametrized by the equivalent number of neutrino flavors $`N_{\mathrm{eff}}`$, on the electron-neutrino degeneracy parameter $`\xi =\mu _{\nu _e}/T_{\nu _e}`$, and on the degeneracy parameters of the other neutrino flavors. Originally all of these parameters had to be fitted such as to reproduce the observed light-element abundances.
Meanwhile the number of neutrino flavors has been fixed by the $`Z^0`$ decay width Eidelman:2004wy . More recently, the baryon abundance has been determined by the cosmic microwave background (CMB) anisotropies to be Spergel:2003cb
$$\eta _B=\frac{n_Bn_{\overline{B}}}{n_\gamma }=6.14\times 10^{10}(1.00\pm 0.04).$$
(1)
Finally, the measured neutrino mixing parameters imply that neutrinos reach approximate chemical equilibrium before the BBN epoch so that all neutrino chemical potentials can be taken to be equal, i.e. they are all characterized by the same degeneracy parameter $`\xi `$ that applies to $`\nu _e`$ Dolgov:2002ab ; Wong:2002fa ; Abazajian:2002qx .
Of course, one may still assume new low-energy physics such as sizeable active-sterile neutrino mixings, the existence of new particles such as axions that could contribute to the primordial radiation density, a time-varying Newton constant, and a plethora of other non-standard options that would affect BBN Sarkar:1995dd . Barring these more exotic possibilities, the single remaining BBN parameter that is not otherwise fixed is $`\xi `$. In the light of these developments we re-examine BBN under the assumption that $`\xi `$ is the only free input parameter.
## II BBN with a lepton asymmetry
Over the years, BBN with a lepton asymmetry has been studied by many authors Wagoner:1966pv ; Freese:1982ci ; Kang:1991xa ; Esposito:2000hi ; Esposito:2000hh ; Barger:2003rt ; Cuoco:2003cu ; Kneller:2004jz . To achieve approximate agreement between the observed and predicted light-element abundances a possible lepton asymmetry must be small so that one may well constrain a new analysis, say, to the range $`|\xi |<0.5`$. For such small $`\xi `$ values the most important impact on BBN is a shift of the beta equilibrium between protons and neutrons. Subleading effects include a modification of the radiation density,
$$\mathrm{\Delta }N_{\mathrm{eff}}=3\left[\frac{30}{7}\left(\frac{\xi }{\pi }\right)^2+\frac{15}{7}\left(\frac{\xi }{\pi }\right)^4\right].$$
(2)
Moreover, the neutrino decoupling temperature is higher than in the standard case Freese:1982ci ; Kang:1991xa so that in principle one could get a non-standard $`T_\nu (T)`$ evolution, but such effects are completely negligible for our case. A non-zero $`\xi `$ slightly modifies the partial neutrino reheating following the $`e^+e^{}`$ annihilation Esposito:2000hi , again a completely negligible effect for our $`\xi `$ range of interest.
To include the modified $`np`$ weak rates we adopt a perturbative approach for the modified neutrino distribution functions, similar to the one adopted by us in a different context Serpico:2004nm . Therefore, the additional terms for the rates can be obtained by integrating the factors
$$\delta f_\nu (E,\xi )=\underset{n=1}{\overset{\mathrm{}}{}}\frac{f_\nu (E,\xi )}{\xi }|_{\xi =0}\frac{\xi ^n}{n!}.$$
(3)
Corrections up to fourth order were included into the code of Ref. Serpico:2004gx , thus reducing the error of the truncation on the nuclide yields to the fourth significant digit for $`|\xi |0.1`$, and to the third for $`|\xi |0.5`$. This also sets the level of the required fitting accuracy of the auxiliary functions. The corrections were calculated as relative changes with respect to the $`\xi =0`$ Born rates so that in the $`\xi 0`$ limit we recover the standard $`np`$ rates with finite mass, QED radiative and thermodynamic corrections, and partial neutrino reheating Serpico:2004gx . For $`\xi 0`$ these sub-leading effects are rescaled without introducing significant errors. We have checked the reliability of our approach by comparing with the code of Ref. Esposito:2000hh where a fully numerical and non-perturbative method is used.
In Fig. 1 we show our predictions for the primordial light-element abundances as a function of the neutrino degeneracy parameter $`\xi `$, taken to be equal for all flavors. The gray band is the $`1\sigma `$ predicted range. We used the nuclear reactions and uncertainties adopted in Ref. Serpico:2004gx . However, except for lithium the uncertainty is dominated by the error of the Wilkinson microwave anisotropy probe (WMAP)-implied Spergel:2003cb baryon abundance of Eq. (1). In Fig. 2 we show a blowup of the top panel, i.e. the predicted helium abundance.
We also show observations of the primordial abundances. Beginning with $`{}_{}{}^{7}\mathrm{Li}`$ (bottom panel of Fig. 1) we show the allowed range $`\mathrm{log}_{10}(X_7)+12=2.0`$โ2.3, where the uncertainty on $`T_{\mathrm{eff}}`$ has been included Lambert:2004kn . In the second to last panel we show Bania et al.โs possible upper limit on the <sup>3</sup>He abundance of $`1.3\times 10^5`$ Bania2002 (at $`1\sigma `$). The status and the interpretation of these observations is quite controversial, especially for $`{}_{}{}^{7}\mathrm{Li}`$; see Ref. Serpico:2004gx for further discussions and references. It is evident, however, that the $`{}_{}{}^{3}\mathrm{He}`$ upper limit is compatible with any realistic $`\xi `$ value whereas $`{}_{}{}^{7}\mathrm{Li}`$ is always incompatible. Rather than dismissing standard BBN as the correct theory for the light-element synthesis we follow the usual practice of dismissing the $`{}_{}{}^{7}\mathrm{Li}`$ observations as too uncertain to be useful. Putting this more positively, it is remarkable how close the $`{}_{}{}^{7}\mathrm{Li}`$ abundances observed on the Spite plateau in metal poor stars come to the BBN prediction.
Turning to deuterium (second panel in Fig. 1), we note that a new detection was reported, giving $`X_2{}_{}{}^{2}\mathrm{H}/\mathrm{H}=1.60_{0.30}^{+0.25}\times 10^5`$ Crighton:2004aj , significantly lower than the previous value towards Q1243+3047 of $`2.42_{0.25}^{+0.35}\times 10^5`$ Kirkman:2003uv . By adopting the conservative approach to symmetrize the errors to the higher value, and the wider allowed range (1.57โ$`4.00)\times 10^5`$ for the most uncertain determination towards Q0347-3819, one gets $`2.42\pm 0.27`$ (outer error box in Fig. 1) for the average value of the seven positive detections so far obtained in quasar Lyman absorption systems. As previously noted Kirkman:2003uv , since $`\chi ^2`$ per degree of freedom is approximately 4.4, even with these generous error estimates, underestimated systematics may be present.
Still, the observed primordial deuterium abundance looks perfectly compatible with the standard BBN predictions in the absence of a lepton asymmetry. Neglecting possible systematics one finds a lower limit $`\xi >0.2`$ at $`1\sigma `$, but no useful upper limit obtains. This remains true even if the precision of the deuterium observations improve by a factor of two (inner box in Fig. 1).
From Figs. 1 and 2 it is obvious that $`{}_{}{}^{4}\mathrm{He}`$ is the best probe of $`\xi `$. As the dependence on the baryon abundance $`\eta _B`$ is only logarithmic, the new WMAP data do not significantly improve the constraint. Therefore, the observed range for $`Y_\mathrm{p}`$ is the only truly significant measure for a possible lepton asymmetry.
Recently new data were published by Izotov and Thuan Izotov:2003xn , giving $`Y_\mathrm{p}=0.2421\pm 0.0021`$ (inner error box in Figs. 1 and 2). Slightly higher values arise for differently chosen samples and/or for other metallicity regressions. For the purpose of illustration we have performed a likelihood analysis of the BBN prediction for $`Y_\mathrm{p}(\eta _B,\xi )`$ using the $`\eta _B`$ prior of Eq. (1) and Izotov and Thuanโs $`Y_\mathrm{p}`$ with their quoted statistical error. The result suggests a positive value $`\xi =0.0245\pm 0.0092`$ ($`1\sigma `$), i.e. a hint for $`\xi 0`$ at $`2.7\sigma `$. This exercise illustrates the possible sensitivity that present determinations of <sup>4</sup>He could reach if the systematic uncertainties were fully understood.
In an independent analysis Olive and Skillman Olive:2004kq find $`Y_\mathrm{p}=0.249\pm 0.009`$. In an attempt to quantify systematic effects, they suggest as the most conservative range $`0.232Y_\mathrm{p}0.258`$ (outer box in Figs. 1 and 2), corresponding to $`0.044\xi 0.070`$ for the $`1\sigma `$ range of $`\eta _B`$ or $`0.046\xi 0.072`$ for the $`2\sigma `$ range.
To compare $`\xi `$ with $`\eta _B`$ on the same footing we note that the cosmic abundance of a lepton flavor $`\alpha `$ is given in terms of the corresponding degeneracy parameter as
$$\eta _\alpha =\frac{n_\alpha n_{\overline{\alpha }}}{n_\gamma }=\frac{1}{12\zeta (3)}\left(\frac{T_\alpha }{T_\gamma }\right)^3\left(\pi ^2\xi _\alpha +\xi _\alpha ^3\right).$$
(4)
Charge neutrality implies that $`\eta _e`$ is of the same order as $`\eta _B`$. A possible large lepton asymmetry thus resides in neutrinos,
$$\eta _L\underset{\nu }{}\eta _\nu 3\times 0.249\xi ,$$
(5)
where $`(T_\nu /T_\gamma )^3=4/11`$ was assumed. Note that this factor is missing in Eq. (1) of Ref. Barger:2003rt and Eq. (10) of Ref. Kneller:2004jz , that limits their validity to temperatures larger than the electron mass.
## III Discussion and Summary
We have argued that after the experimental determination of the neutrino mixing parameters and after the WMAP determination of the baryon abundance, of all parameters that characterize the cosmic thermal heat bath at the BBN epoch only the lepton asymmetry, i.e. a neutrino degeneracy parameter $`\xi `$ common to all flavors, remains undetermined and has to be fixed by the observed light-element abundances. For the first time we have produced a meaningful plot of the predicted light-element abundances as a function of $`\xi `$. One motivation for writing this short note was to present for the first time this plot as an alternative to the still common but outdated depiction of the light-element abundances as a function of $`\eta _B`$.
The errors of the predictions in our plot are typically dominated by the WMAP determination of the baryon fraction $`\eta _B`$ and thus will significantly improve with forthcoming CMB missions.
The deuterium abundance is very sensitive to $`\eta _B`$ so that deuterium is the best baryometer of all the light elements. On the other hand, this element does not respond much to a non-zero lepton asymmetry. The opposite applies to helium which is by far the most sensitive leptometer whereas it is virtually insensitive to future improvements of $`\eta _B`$ determinations.
The usual attitude towards the possibility of a large cosmic lepton asymmetry is that sphaleron effects before electroweak symmetry breaking equilibrate the cosmic lepton and baryon asymmetries to within a factor of order unity. In the standard model one finds that $`\eta _L=\frac{51}{28}\eta _B`$. From this perspective, BBN is a parameter-free theory and the cosmic helium abundance is fixed by the $`\xi =0`$ value shown in Fig. 2, i.e. $`Y_\mathrm{p}=0.2479\pm 0.0005`$ (1 $`\sigma `$). In this scenario it remains to understand the systematic errors in the spectroscopic $`Y_\mathrm{p}`$ determinations.
However, we think that this argument should be turned around. Sphaleron effects are a crucial ingredient in most baryogenesis scenarios Dine:2003ax , including leptogenesis Buchmuller:2005eh . No experimental evidence for or against these effects exists. Therefore, a possible indication for a nonzero $`\xi `$ by the observed helium abundance presents an approach for a possible falsification of the usual assumptions about baryon-lepton equilibration.
In principle, other cosmological probes of a non-zero $`\xi `$ exist in the form of the CMB anisotropies and large-scale structure (LSS) power spectra Kinney:1999pd ; Lesgourgues:1999wu . The former essentially feels a boost in the amplitude of the first peak and a shift towards higher multipoles, effects that are degenerate with $`N_{\mathrm{eff}}>3`$. The latter suffers a suppression very similar to the effect of a non-zero $`m_\nu `$. However, the current CMB and LSS data are only sensitive to $`\xi =๐ช(1)`$ Hannestad:2003xv .
A step further may be provided by future CMB limits on $`\mathrm{\Delta }N_{\mathrm{eff}}`$, possibly helped by accurate polarization maps Lopez:1998aq . Assuming Kinney and Riottoโs estimate Kinney:1999pd of the capability of PLANCK to detect one degeneracy parameter of order $`0.5`$, we can easily translate it into a marginal sensitivity of $`\xi 0.3`$ that is common to all three flavors. However, unless one is able to break the degeneracies with other cosmological parameters, such a level of accuracy is too optimistic.
To reach a sensitivity comparable with BBN, one has to wait for more ambitious methods studying the gravitational lensing distortions on both the temperature and polarization maps of the CMB Kaplinghat:2003bh . Emphasis is usually put on the sensitivity to neutrino masses, but these measurements are useful to constrain the relic neutrino asymmetry as well, especially if complementary information would be provided by improved LSS maps and direct experimental data on the neutrino masses. Incidentally, such an accurate satellite mission would be sensitive to $`Y_\mathrm{p}`$ at the recombination epoch at a level of $`\delta Y_\mathrm{p}0.005`$ Kaplinghat:2003bh , thus comparable or better than present astrophysical spectroscopic determinations, with a positive impact on the BBN constraint on $`\xi `$. In any event, BBN is the only probe sensitive to the sign of $`\xi `$.
In summary, now that the cosmic baryon abundance has been extremely well determined by CMB observations and now that much about neutrino properties has been learnt by experiments, the role of BBN as a baryometer has shifted to that of the best available cosmic leptometer. Therefore, a more reliable $`Y_\mathrm{p}`$ determination is of much greater fundamental interest than the next round of more precise CMB baryon determinations.
## ACKNOWLEDGMENTS
We acknowledge partial support by the Deutsche Forschungsgemeinschaft under Grant No. SFB-375 and by the European Union under the Ilias project, contract No. RII3-CT-2004-506222.
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# Hole and electron dynamics in the triangular-lattice antiferromagnet โ interplay of frustration and spin fluctuations
## I Introduction
There has been renewed interest in correlated electron systems on triangular lattices, as evidenced by recent studies of antiferromagnetism, superconductivity and metal-insulator transition in the organic systems $`\kappa (\mathrm{BEDT}\mathrm{TTF})_2\mathrm{X}`$,review1 ; review2 the discovery of superconductivity in $`\mathrm{Na}_\mathrm{x}\mathrm{CoO}_2.\mathrm{yH}_2\mathrm{O}`$,watersup the observation of low-temperature insulating phases in some $`\sqrt{3}`$-adlayer structures such as K on Si,weitering and quasi two-dimensional $`120^{}`$ spin ordering and spin-wave excitations in $`\mathrm{RbFe}(\mathrm{MoO}_4)_2`$ (Refs. 5,6) and the multiferroic material $`\mathrm{HoMnO}_3`$.holmium ; sw
The recent finding of finite $`U`$-induced competing interactions and frustration in the $`120^{}`$ ordered antiferromagnetic (AF) state of the Hubbard model on a triangular lattice,tri resulting in vanishing spin stiffness at $`U6`$ and a magnetic instability towards a F-AF state at $`U7`$, both in the insulating state, adds a new dimension to the intrinsic geometric frustration of the triangular-lattice antiferromagnet. Indeed, strongly enhanced quantum spin fluctuations associated with the magnetic instability may account for the absence of long-range magnetic order in the nearly isotropic organic antiferromagnet $`\kappa (\mathrm{BEDT}\mathrm{TTF})_2\mathrm{Cu}_2(\mathrm{CN})_3`$, as inferred from recent <sup>1</sup>H NMR and static susceptibility measurements down to 32 mK, well below the estimated exchange constant $`J250`$ K, suggesting the realization of a quantum spin-liquid state.kanoda The realization of a non-magnetic insulator at intermediate $`U`$ is also interesting as it allows, with decreasing $`U`$, for a Mott-type metal-insulator transition not accompanied by any magnetic symmetry breaking, as seen in the layered system $`\kappa (\mathrm{BEDT}\mathrm{TTF})_2\mathrm{Cu}[\mathrm{N}(\mathrm{CN})_2]\mathrm{Cl}`$,kagawa and currently of much theoretical interest.pirg2 ; pirg3
Recently, self-energy corrections due to multiple magnon processes in the AF state of the frustrated square-lattice $`tt^{}`$-Hubbard model were evaluated using a fluctuation approach.self Quasiparticle dispersion obtained with the same set of Hubbard model cuprate parameters as obtained from a recent magnon spectrum fitspectrum was found to yield excellent agreement with ARPES data for $`\mathrm{Sr}_2\mathrm{CuO}_2\mathrm{Cl}_2`$,self thus providing a unified description of magnetic and electronic excitations in cuprates.
It is therefore interesting to examine and contrast self-energy corrections and quasiparticle behaviour in the $`120^{}`$ ordered AF state of the triangular-lattice Hubbard model which involves non-collinear ordering, intrinsic geometric frustration, and also finite $`U`$-induced competing interactions and frustration. Indeed, we find that long-wavelength magnon modes yield finite contribution to the fermion-magnon scattering process, unlike the square-lattice case where this contribution was negligible.self
Frustration and spin fluctuations are involved in an interesting interplay with respect to quasiparticle behaviour. As neighbouring spins are not antiparallel in triangular-lattice and frustrated square-lattice AF states, frustration results in an O($`t`$) or O($`t^{}`$) direct hopping in addition to the O($`J`$) virtual hopping. Competition between the two dispersions results in band broadening/narrowing, which has a dramatic effect on self-energy corrections due to significantly different density of fermion scattering states. Competition also results in a reduced band gap, thus bringing the system closer to a metal-insulator (MI) transition. The $`tJ`$ model calculations do not involve this competition as the virtual hopping term is absent.
In this paper, we examine the hole and electron spectral functions in terms of self-energy corrections in the $`120^{}`$ ordered AF state of the Hubbard model on a triangular lattice. Hole and electron dynamics in an antiferromagnetic background is associated with multiple magnon emission and absorption processes corresponding to broken AF bonds, string states, and scrambling of the AF spin ordering. The fluctuation approach adopted in this paper in terms of a diagrammatical expansion provides a physically transparent picture in which the hole motion is renormalized due to the fluctuating transverse field associated with spin fluctuations. This fluctuation perspective is physically relevant in view of the substantial ordered-moment reduction, which approaches $`59\%`$ in the strong-coupling (Heisenberg) limit.capriolti
Single-particle dynamics in the AF state of the triangular-lattice $`tJ`$ model has been studied in the self-consistent Born approximation (SCBA),dombre ; trumper ; apel exact diagonalization,dombre and projection techniques.vojta ; dagotto One-electron density of states has been examined using the quantum Monte Carlo method,qmc\_green showing a pseudogap development for intermediate $`U`$.
We consider the Hubbard model
$$=t\underset{i\delta \sigma }{}a_{i,\sigma }^{}a_{i+\delta ,\sigma }+U\underset{i}{}a_i^{}a_ia_i^{}a_i$$
(1)
with nearest-neighbour (NN) hopping on a triangular lattice. The model has particle-hole symmetry under the transformation $`tt`$. In this paper we consider the case of positive $`t`$ and hole (electron) doping in the lower (upper) band; same result holds for negative $`t`$ and electron (hole) doping in the upper (lower) band. In the following we set the hopping energy $`t=1`$.
The organization of this paper is as follows. The three-sublattice basis is briefly reviewed in section II to introduce the notation and key features of the classical-level fermion dispersion. Transverse spin fluctuations are introduced in section III in terms of magnon amplitudes and energies. Intraband self-energy corrections due to multi-magnon processes are studied in section IV in the non-crossing approximation. Hole and electron spectral functions for the lower and upper bands are discussed in section V and conclusions are presented in section VI.
## II Three-sublattice representation
While the spiral-state description applies only to Bravais lattices, the sublattice-basis description applies to Kagomรฉ type non-Bravais lattices as well. In general, the $`120^{}`$ AF state is characterized by an ordering plane (normal $`\widehat{n}_1`$) and a planar direction ($`\widehat{n}_2`$) in spin space, with reference to which spin orientations are given by
$$\widehat{\alpha }=\mathrm{cos}\varphi _\alpha \widehat{n}_2+\mathrm{sin}\varphi _\alpha (\widehat{n}_1\times \widehat{n}_2)$$
(2)
corresponding to angles $`\varphi _\alpha =0^{}`$, $`120^{}`$, and $`240^{}`$ on the three sublattices $`\alpha =\mathrm{A},\mathrm{B},\mathrm{C}`$. A convenient choice is $`\widehat{n}_1=\widehat{z}`$ (spin-ordering in the $`xy`$ plane) and $`\widehat{n}_2=\widehat{x}`$, so that the local mean field $`๐ซ_\alpha =\frac{1}{2}U\mathrm{\Psi }_i^{}๐\mathrm{\Psi }_i_\alpha `$ in the $`120^{}`$ ordered Hartree-Fock state is given by
$$๐ซ_\alpha =\mathrm{\Delta }\widehat{\alpha }(\widehat{\alpha }=\widehat{a},\widehat{b},\widehat{c})$$
(3)
on the three sublattices in terms of the three lattice unit vectors
$$\widehat{a}=\widehat{x},\widehat{b}=\frac{1}{2}\widehat{x}+\frac{\sqrt{3}}{2}\widehat{y},\widehat{c}=\frac{1}{2}\widehat{x}\frac{\sqrt{3}}{2}\widehat{y}.$$
(4)
Fourier transformation within the sublattice basis yields
$$_{\mathrm{HF}}=\underset{๐ค}{}\mathrm{\Psi }_๐ค^{}\left[\begin{array}{ccc}๐.๐ซ_\mathrm{A}& \delta _๐ค& \delta _๐ค^{}\\ & & \\ \delta _๐ค^{}& ๐.๐ซ_\mathrm{B}& \delta _๐ค\\ & & \\ \delta _๐ค& \delta _๐ค^{}& ๐.๐ซ_\mathrm{C}\end{array}\right]\mathrm{\Psi }_๐ค,$$
(5)
where $`\mathrm{\Psi }_๐ค(a_๐คa_๐คb_๐คb_๐คc_๐คc_๐ค)`$ in terms of the fermion operators $`a_๐ค,b_๐ค,c_๐ค`$ defined on the three sublattices A, B, C. Wavevector $`๐ค`$ lies within the Magnetic Brillouin Zone (MBZ), corresponding to the three inter-penetrating triangular sublattices (lattice parameter $`\sqrt{3}a`$). The NN hopping term
$$\delta _๐ค=t\underset{\widehat{\delta }=\widehat{a},\widehat{b},\widehat{c}}{}e^{i๐ค.\widehat{\delta }}=t[e^{ik_x}+2e^{ik_x/2}\mathrm{cos}(\sqrt{3}k_y/2)]$$
(6)
mixes AB, BC, and CA sublattices, which are connected by the three lattice unit vectors. The lattice hopping term $`\delta _๐ค`$ yields the triangular-lattice free-fermion energy
$$ฯต_๐ค=\delta _๐ค+\delta _๐ค^{}$$
(7)
and transforms as
$$\delta _{๐ค\pm ๐}=\delta _๐คe^{\pm i2\pi /3}$$
(8)
under momentum translation by the ordering wavevector $`๐=(2\pi /3,2\pi /\sqrt{3})`$.
The $`[6\times 6]`$ Hamiltonian matrix obeys the cyclic property $`[_{\mathrm{HF}}]_{AB}=[_{\mathrm{HF}}]_{BC}=[_{\mathrm{HF}}]_{CA}`$ of the $`120^{}`$ ordered state, resulting in the following spin-sublattice structure of the normalized eigenvectors
$$|๐ค,l=\frac{1}{\sqrt{3}}\left(\begin{array}{c}\alpha _{๐ค,l}e^{i\varphi _\alpha }\hfill \\ \beta _{๐ค,l}\hfill \end{array}\right)_\sigma \left(\begin{array}{c}1\\ e^{i\theta _l}\\ e^{i\theta _l}\end{array}\right)_\alpha ,$$
(9)
where the planar spin orientations $`\varphi _\alpha =0^{},\mathrm{\hspace{0.33em}120}^{},\mathrm{\hspace{0.33em}240}^{}`$ for the three sublattices, and the sublattice phase angle $`\theta _l=2\pi l/3`$ for the three fermion branches $`l=0,\pm 1`$. Substituting the above eigenvector structure, and contracting over the sublattice sector $`\alpha `$, the eigenvalue equations of the $`[6\times 6]`$ Hamiltonian matrix reduce to three $`[2\times 2]`$ equations in the spin sector corresponding to $`l=0,\pm 1`$; the six eigenvalues $`E_{๐ค,l}^{}`$ and amplitudes $`(\alpha _{๐ค,l}^{}\beta _{๐ค,l}^{})`$ for the lower ($``$) and upper ($`+`$) AF bands are then given by
$`E_{๐ค,l}^{}`$ $`=`$ $`\eta _๐ค\pm \sqrt{\mathrm{\Delta }^2+\xi _๐ค^2}`$ (10)
$`\alpha _{๐ค,l}^{}`$ $`=`$ $`\pm {\displaystyle \frac{1}{\sqrt{2}}}\left(1\pm {\displaystyle \frac{\xi _๐ค}{\sqrt{\mathrm{\Delta }^2+\xi _๐ค^2}}}\right)^{1/2}`$
$`\beta _{๐ค,l}^{}`$ $`=`$ $`+{\displaystyle \frac{1}{\sqrt{2}}}\left(1{\displaystyle \frac{\xi _๐ค}{\sqrt{\mathrm{\Delta }^2+\xi _๐ค^2}}}\right)^{1/2},`$ (11)
where $`\eta _๐ค(ฯต_๐ค+ฯต_{๐ค๐})/2`$ and $`\xi _๐ค(ฯต_๐คฯต_{๐ค๐})/2`$, with momentum values $`๐ค`$, $`๐ค\pm ๐`$ corresponding to fermion branches $`l=0,\pm 1`$, respectively. The mean field $`\mathrm{\Delta }`$ and magnetization $`m`$ are related to $`U`$ through the self-consistency condition.tri
Typical of the AF state, upper band states yield a negative contribution to spin densities due to the negative sign of $`\alpha _๐ค^+`$. Consequently, when the two bands start overlapping on decreasing $`U`$, partial occupation of the upper band has an amplified effect on reduction of magnetic order โ reduced sublattice magnetization due to band overlap decreases the mean-field $`\mathrm{\Delta }`$, which further increases the overlap. Typically, the magnetic order therefore falls very rapidly after band overlap sets in.
The density of states (DOS) corresponding to the full HF (classical) dispersion in Eq. (10), shown in Fig. 2 for both positive and negative $`t`$, exhibits a key competition between the two dispersion terms associated with direct hopping $`(\eta _๐ค)`$ of order $`t`$ and virtual hopping $`(\sqrt{\mathrm{\Delta }^2+\xi _๐ค^2}\mathrm{\Delta })`$, which is of order $`J`$ in the strong-coupling limit. For positive $`t`$, both dispersion terms favour same state at the top of the lower band, while they favour different states at the bottom of the upper band, thereby causing band narrowing. The competition results in broadening and narrowing of the lower and upper bands, depending on the sign of $`t`$, which has a dramatic effect on the self-energy correction due to significantly different density of fermion scattering states. There is no such competition in the $`tJ`$ model studies where the virtual hopping dispersion term is absent.
From Eq. (8), it follows that
$$ฯต_๐ค=(ฯต_{๐ค+๐}+ฯต_{๐ค๐})$$
(12)
so that
$$\eta _{๐ค๐}=\frac{ฯต_๐ค}{2}.$$
(13)
Therefore, the direct hopping dispersion $`\eta _{๐ค๐}`$, to which the HF quasiparticle dispersion (10) reduces in the strong-coupling limit, is identical to the classical-level dispersion $`ฯต_๐ค/2`$ taken in earlier $`tJ`$ model studies, corresponding to effective hopping $`t\mathrm{cos}(120^{})=t/2`$ associated with the $`120^{}`$ ordering, where momentum translation by $`๐`$ connects the real and slave fermions.dombre Comparison of DOS corresponding to the full dispersion with that of $`\eta _๐ค`$ is also shown in Fig. 2.
## III Transverse spin fluctuations
Including both transverse and longitudinal spin fluctuations, the full spin-fluctuation propagator $`\mathrm{\Psi }_\mathrm{G}|\mathrm{T}[S_i^\mu (t)S_j^\nu (t^{})]\mathrm{\Psi }_\mathrm{G}`$ in the $`120^{}`$ ordered AF state has been studied recently in the random phase approximation (RPA) in the full $`U`$ range.tri Even in the intermediate-coupling regime, the magnitude of longitudinal fluctuation $`S_\alpha ^2`$ along the local ordering directions $`\widehat{\alpha }`$ was found to be quite negligible ($`S_\alpha ^210^4`$ at $`U7`$). We therefore focus on transverse spin fluctuations along the two locally normal spin directions.
Spin rotation about the $`\widehat{z}`$ direction
$$\left(\begin{array}{c}\sigma _x\\ \sigma _y\\ \sigma _z\end{array}\right)^{}=\left[\begin{array}{ccc}\mathrm{cos}\varphi _\alpha & \mathrm{sin}\varphi _\alpha & 0\\ \mathrm{sin}\varphi _\alpha & \mathrm{cos}\varphi _\alpha & 0\\ 0& 0& 1\end{array}\right]\left(\begin{array}{c}\sigma _x\\ \sigma _y\\ \sigma _z\end{array}\right)$$
(14)
by angles $`\varphi _\alpha =0^{},120^{},120^{}`$ for $`\alpha `$ = A,B,C renders $`x^{}`$ as the spin ordering direction for all three sublattices. In the $`23`$ spin-sublattice basis of the two transverse spin directions $`\mu ,\nu =y^{},z^{}`$ and the three sublattices $`\alpha ,\beta =A,B,C`$, the RPA-level spin-fluctuation propagator is then given by
$$[\chi (๐ช,\omega )]_{\alpha \beta }^{\mu \nu }=\frac{\frac{1}{2}[\chi ^0(๐ช,\omega )]}{\mathrm{๐}U[\chi ^0(๐ช,\omega )]},$$
(15)
where the bare particle-hole propagator
$`[\chi ^0(๐ช,\omega )]_{\alpha \beta }^{\mu \nu }=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{๐ค,l,m}{}}[{\displaystyle \frac{\sigma _\mu _\alpha ^+\sigma _\nu _\beta ^+}{E_{๐ค๐ช,m}^+E_{๐ค,l}^{}+\omega }}`$ (16)
$`+{\displaystyle \frac{\sigma _\mu _\alpha ^+\sigma _\nu _\beta ^+}{E_{๐ค,m}^+E_{๐ค๐ช,l}^{}\omega }}]`$
involves integrating out the fermions in the broken-symmetry state. In the particle-hole matrix elements of the rotated spins
$$\sigma _\mu _\alpha ^+๐ค๐ช,m|\sigma _\mu |๐ค,l_\alpha $$
(17)
the spin orientation angles $`\varphi _\alpha `$ in the fermion states (Eq. 9) are transformed out.
We now discuss the spin-sublattice structure of the $`[\chi ^0(๐ช,\omega )]`$ matrix and its eigenvectors. While the spin-diagonal blocks $`[\chi ^0(๐ช,\omega )]^{\mu \mu }`$ are Hermitian and the off-diagonal blocks $`[\chi ^0(๐ช,\omega )]^{\mu \overline{\mu }}`$ are anti-Hermitian, they obey the cyclic symmetry $`[\chi ^0(๐ช,\omega )]_{AB}=[\chi ^0(๐ช,\omega )]_{BC}=[\chi ^0(๐ช,\omega )]_{CA}`$ of the $`120^{}`$ ordered phase. Also, the sublattice diagonal elements $`[\chi ^0(๐ช,\omega )]_{\alpha \alpha }`$ are all identical due to sublattice symmetry. Consequently, the normalized eigenvectors of $`[\chi ^0(๐ช,\omega )]`$ have the following spin-sublattice structure
$$|\varphi _\lambda =\frac{1}{\sqrt{3}}\left(\begin{array}{c}\hfill iu\\ \hfill v\end{array}\right)_\mu \left(\begin{array}{c}1\\ e^{i\mathrm{\hspace{0.17em}2}\pi \lambda /3}\\ e^{i\mathrm{\hspace{0.17em}2}\pi \lambda /3}\end{array}\right)_\alpha $$
(18)
where $`\lambda =0,\pm 1`$ for the three magnon branches, and the real and normalized amplitudes $`u`$ and $`v`$ represent the fluctuation amplitudes in the $`y^{}`$ and $`z^{}`$ directions, respectively. Contracting over the sublattice index, the eigenvalue equation for $`|\varphi `$ therefore reduces to
$$[\chi ^0(๐ช,\omega )]|\varphi _\lambda =[\chi _\lambda ^0(๐ช,\omega )]\left(\begin{array}{c}u\\ v\end{array}\right)=\lambda _๐ช\left(\begin{array}{c}u\\ v\end{array}\right),$$
(19)
where $`[\chi _\lambda ^0(๐ช,\omega )]`$ is a $`[2\times 2]`$ real-symmetric matrix.
Solving the pole equation $`1U\lambda _๐ช(\omega )=0`$ for the magnon energy $`\omega _๐ช`$, and expanding $`\lambda _๐ช(\omega )`$ around the poles to obtain the magnon amplitudes, yields the magnon propagator
$$[\chi (๐ช,\omega )]=\underset{\lambda =0,\pm 1}{}\left[\frac{|๐ช,\lambda ๐ช,\lambda |_A}{\omega +\omega _{๐ช,\lambda }i\eta }\frac{|๐ช,\lambda ๐ช,\lambda |_R}{\omega \omega _{๐ช,\lambda }+i\eta }\right],$$
(20)
where the magnon eigenvectors for the advanced (A) and retarded (R) modes are given by
$$|๐ช,\lambda =\frac{1}{\sqrt{3}}\left(\begin{array}{c}\hfill i๐ด_{๐ช,\lambda }\\ \hfill ๐ต_{๐ช,\lambda }\end{array}\right)_\mu \left(\begin{array}{c}1\\ e^{i\mathrm{\hspace{0.17em}2}\pi \lambda /3}\\ e^{i\mathrm{\hspace{0.17em}2}\pi \lambda /3}\end{array}\right)_\alpha $$
(21)
in terms of the magnon amplitudes
$`๐ด_{๐ช,\lambda }=u_{๐ช,\lambda }/\sqrt{2U^2|d\lambda _{๐ช,\lambda }/d\omega |}\mathrm{and}`$
$`๐ต_{๐ช,\lambda }=v_{๐ช,\lambda }/\sqrt{2U^2|d\lambda _{๐ช,\lambda }/d\omega |}`$ (22)
in the $`y^{}`$ and $`z^{}`$ directions. Expressions for magnon energy and amplitudes in the strong-coupling limit are given in the Appendix.
## IV Self-energy correction
Due to multiple magnon emission and absorption processes associated with fermion motion in the AF state, the fermion self energy matrix in the spin-sublattice basis is obtained as
$`[\mathrm{\Sigma }(๐ค,l,\omega )]_{\alpha \beta }`$
$`=`$ $`U^2{\displaystyle \frac{d\mathrm{\Omega }}{2\pi i}\underset{\mu \nu }{}[\sigma _\mu ][G_{\alpha \beta }(๐ค๐ช,m,\omega \mathrm{\Omega })][\sigma _\nu ][\chi (๐ช,\mathrm{\Omega })]_{\alpha \beta }^{\mu \nu }}`$
in the non-crossing (rainbow) approximation.
We consider the intraband contribution involving hole (electron) scattering states in the lower (upper) band, which redistributes the spectral function, leaving the integrated spectral weight and the sublattice magnetization unchanged. We obtain for the hole self energy
$`\mathrm{\Sigma }_{๐ค,l}(\omega )`$ $``$ $`๐ค,l|[\mathrm{\Sigma }(๐ค,l,\omega )]|๐ค,l`$
$`=`$ $`{\displaystyle \underset{๐ช,m,\lambda }{}}{\displaystyle \frac{|M|^2}{\omega +\omega _{๐ช,\lambda }E_{๐ค๐ช,m}^{}\mathrm{\Sigma }_{๐ค๐ช,m}(\omega +\omega _{๐ช,\lambda })}}`$
where the fermion-magnon scattering matrix element
$$M=U\underset{\alpha \mu }{}๐ค,l|\sigma _\mu |๐ค๐ช,m_\alpha |๐ช,\lambda _{\alpha \mu }^\mathrm{A}$$
(25)
involves the advanced magnon mode. Substituting the sublattice structure of the fermion and magnon amplitudes, the sum over sublattice index $`\alpha `$ yields
$$1+e^{i(m+\lambda l)2\pi /3}+e^{i(m+\lambda l)2\pi /3}=3\delta _{m+\lambda l}$$
(26)
effectively amounting to a conservation of sublattice polarization at the fermion-magnon interaction vertex. Therefore, the fermion-magnon scattering matrix element reduces to a sum of the $`y^{}`$ and $`z^{}`$ fluctuation terms
$`M={\displaystyle \frac{U}{\sqrt{3}}}[`$ $``$ $`(\alpha _{๐ค,l}^{}\beta _{๐ค๐ช,m}^{}\beta _{๐ค,l}^{}\alpha _{๐ค๐ช,m}^{})๐ด_{๐ช,\lambda }`$ (27)
$`+`$ $`(\alpha _{๐ค,l}^{}\alpha _{๐ค๐ช,m}^{}\beta _{๐ค,l}^{}\beta _{๐ค๐ช,m}^{})๐ต_{๐ช,\lambda }].`$
For fermion and magnon states in the matrix element $`M`$, the sublattice-basis MBZ description of momentum translates to a BZ description according to the correspondence $`๐ค,l๐ค+l๐`$ and $`๐ช,\lambda ๐ช+\lambda ๐`$. With this equivalence, Eq. (26) simply corresponds to momentum conservation in the BZ. Analysis of the fermion-magnon matrix element $`M`$ in the strong-coupling limit and comparison with earlier results for the $`tJ`$ model is discussed in the Appendix.
It is important to note here that long-wavelength magnon modes yield finite contribution to the fermion-magnon scattering process in the triangular-lattice AF, unlike the square-lattice case where this contribution is negligible.self For the square-lattice AF, the small-$`q`$ contribution was suppressed because the fermion-magnon matrix element $`M^2q`$ due to destructive interference within sublattice summation. For the triangular lattice also, for $`q=0`$ and $`\lambda =0`$ (in-plane mode), the fermion matrix element in Eq. (25) reduces to an expectation value which identically vanishes for $`\mu =y^{}`$ as the spins are oriented in the $`x^{}`$ direction, yielding similar $`M^2q`$ behaviour for small $`q`$. However, for the out-of-plane $`z^{}`$ fluctuation modes $`(\lambda =\pm 1)`$, the fermion matrix element is finite, resulting in $`M^21/q`$ and a finite contribution of long-wavelength modes within the two-dimensional $`(q๐q)`$ momentum summation.
## V Results and Discussion
The self-consistent numerical evaluation of the self energy (24) was carried out on a $`30\times 40`$ grid in the MBZ $`๐ค`$ space and a frequency interval $`\mathrm{\Delta }\omega =0.025`$ for $`\omega `$ in the range $`15<\omega <10`$. The self energy was iteratively evaluated, starting with $`\mathrm{\Sigma }_๐ค(\omega )=0`$. Typically, self-consistency was achieved within ten iterations for the lower band (Fig. 4) and fifteen iterations for the upper band (Fig. 8).
Self-energy corrections for an added hole in the broad, flat band near the top of the lower band are qualitatively different from that of an added electron in the narrow, sharp peak near the bottom of the upper band. The low density of hole scattering states and the dominant band-energy denominator suppresses the hole self energy. However, the electron self energy is significantly enhanced due to the sharp peak and the small band energy compared with magnon energy, resulting in the characteristic signature of string states associated with multi-magnon processes. It is convenient to visualize these qualitatively different self-energy corrections in terms of hole (electron) motion in an effective ferromagnetic (antiferromagnetic) spin background projected out of the $`120^{}`$ spin ordering.
### V.1 Hole Dynamics (Lower Band)
Figure 5 shows renormalized quasiparticle dispersion along different symmetry directions in the MBZ. Comparison with the HF result shows nearly momentum independent shift in the quasiparticle energies, leaving the effective hole mass essentially unchanged. States in upper portion of the band are pushed up, while those in the lower portion are pulled down, in accordance with the formally second-order structure of the self-energy correction. The quasiparticle (hole) energy is maximum (minimum) at the $`\mathrm{\Gamma }`$ point $`๐ค=(0,0)`$ for $`l=1`$, corresponding to BZ momentum $`๐`$.
Figure 6 shows the spectral function $`A_๐ค(\omega )`$ for the special points $`\mathrm{\Gamma }`$, K, and M. The weak self-energy correction in the broadened lower band of the frustrated $`120^{}`$ ordered AF state results in no visible oscillatory structure, typically associated with the string of broken bonds as the hole moves in the AF background. As expected, the spectral function at the $`\mathrm{\Gamma }`$ point $`๐ค=(0,0)`$ and $`l=1`$ shows a coherent quasiparticle peak, as this state lies at the top of the lower band. All three branches are degenerate at point K, as at the HF level. The spectral functions at K and M points show well-defined peak structures, with finite quasiparticle damping as these states lie well within the band. The quasiparticle peak broadens and loses intensity along both the $`\mathrm{\Gamma }K`$ and $`\mathrm{\Gamma }M`$ directions, as seen in Fig. 7.
### V.2 Electron Dynamics (Upper Band)
Self-energy correction in the narrow upper band is relatively large and self consistency is noticeably slower (Fig. 8), illustrating the importance of multi-magnon processes and resulting in the characteristic oscillatory structure in the spectral functions (Fig. 10) associated with string states.
Figure 9 shows the quasiparticle dispersion along the $`\mathrm{\Gamma }KM\mathrm{\Gamma }`$ directions for the lowest-energy branch $`l=0`$. The lowest-energy state at point $`M`$ shows a well-defined quasiparticle peak at $`\omega =1.3`$, along with a long incoherent tail, as seen in Fig. 10. All branches are degenerate at $`K`$, second and third branches at $`M`$, first and second branches at $`\mathrm{\Gamma }`$, while at $`N=(2\pi /5,0)`$ all branches are non-degenerate, exactly as at the HF level.
The states $`K`$, $`\mathrm{\Gamma }`$, and $`N`$, which are well inside the band, are strongly damped and yield dominantly incoherent spectral functions, along with small quasiparticle peaks at nearly same frequency $`\omega 2`$ (Fig. 10). It is the strong negative peak in the self energy at $`\omega 2`$ (Fig. 8) which leads to nearly same quasiparticle energy for all $`๐ค`$ points, resulting in drastically reduced quasiparticle bandwidth (Fig. 9) and enhanced effective mass. Figure 11 shows that the well-defined quasiparticle peak at $`M`$ rapidly diminishes in intensity and disappears as one moves along the $`M\mathrm{\Gamma }`$ direction. However, in the $`MK`$ direction, the quasiparticle peak is discernible in the full $`k`$ range. In the $`K\mathrm{\Gamma }`$ direction, the quasiparticle peak marginally increases and then rapidly disappears as one moves towards $`\mathrm{\Gamma }`$.
The lowest-energy hole and electron pockets are shown in Fig. 12. The one-particle density of states is shown in Fig. 13 for $`\mathrm{\Delta }=4`$ and $`3`$. The classical-level asymmetry strongly influences the quantum corrections and the characteristic signature of multi-magnon processes is dominant for the upper band corresponding to the narrow sharp peak at the classical level. As expected within the many-body expansion, the string-state signature of multi-magnon processes is more prominent for higher $`U`$. The renormalized gap vanishes for $`U7`$, whereas the HF band gap $`2(\mathrm{\Delta }2)`$ vanishes at $`U5`$. The band gap is indirect, with the lowest-energy hole and electron states corresponding to momenta $`(4\pi /3,0)`$ and $`(2\pi /3,0)`$, respectively.
It is interesting to note that in the intermediate-coupling regime of interest, the interaction strength effectively controls the fermion-magnon scattering through three different aspects. With decreasing interaction strength, besides the explicit reduction in the fermion-magnon scattering matrix element due to $`U`$ in Eq. (25), the classical-level fermion and magnon dispersions are also significantly modified. Due to enhanced virtual hopping and competition with the direct hopping term the sharp peak in the upper band DOS (Fig. 2) shifts towards higher energy. Also, the magnon energy $`\omega _M`$ at momentum $`๐/2`$ decreases rapidly and vanishes at $`\mathrm{\Delta }=2.9`$.tri Fig. 13 shows the modifications in the upper band DOS due to these intermediate-$`U`$ effects.
## VI Conclusions
In conclusion, we have studied the hole and electron dynamics in the $`120^{}`$ ordered AF state of the Hubbard model on a triangular lattice using a physically transparent fluctuation approach involving the dynamical spin fluctuations, which interpolates between the weak and strong coupling limits. Finite-$`U`$, double occupancy effects, neglected in earlier $`tJ`$ model studies, have been incorporated naturally in terms of classical level fermion and spin dynamics. Intrinsic features of the frustration-induced direct hopping dispersion associated with the $`120^{}`$ ordering - the broad flat band and the narrow sharp peak in the fermion DOS corresponding to lowest-energy hole and electron states - are characteristics of ferromagnetic and antiferromagnetic ordering, respectively. The qualitatively different self-energy corrections for hole and electron can therefore be conveniently visualized in terms of hole (electron) motion in an effective ferromagnetic (antiferromagnetic) spin background projected out of the $`120^{}`$ spin ordering.
For an added hole in the broad lower band, the reduced density of scattering states suppresses the fermion-magnon interaction resulting in nearly coherent quasiparticle peak for all $`๐ค`$ states. No signature of string states in the spectral function reflects an effective F background seen by holes near the top of the lower band. Quasiparticle dispersion shows a nearly momentum-independent shift of hole energies, implying no mass renormalization.
For an electron in the lowest-energy branch $`l=0`$ of the upper band, quasiparticle peak is observed only near the MBZ boundary ($`MK`$), and rapidly vanishes away from it. Strong incoherent behaviour and clear signature of string states in the spectral function is a consequence of an effective AF background seen by electrons near the bottom of the upper band. The strong and nearly momentum-independent peak in the self energy leads to nearly same quasiparticle energy for all $`๐ค`$ points, resulting in drastically reduced bandwidth and enhanced effective mass.
The renormalized band gap was found to vanish for $`U7`$, yielding a first-order M-I transition, as also obtained earlier for the frustrated square-lattice AF.self On the other hand, for the unfrustrated AF the band gap never vanishes for any finite $`U`$. The vanishing of band gap at moderate $`U`$ for both frustrated antiferromagnets due to frustration-induced band broadening thus highlights the role of frustration in M-I transition.
Finally, frustration and spin fluctuations are involved in an interesting interplay with respect to stability of the insulating state. Frustration generally enhances spin fluctuations and magnetic disordering. It will therefore be interesting to also examine the interband self-energy contribution which reduces magnetic order due to interband spectral weight transfer, and also widens the band gap and thereby stabilizes the insulating state. The first-order interband contribution exactly cancels for negligible AF bandwidth, as for the unfrustrated AF in the strong-coupling limit,self but will survive for finite frustration-induced bandwidth, thus highlighting the interplay between frustration and spin fluctuations, and providing deeper insight into the nature of the Mott insulator.
*
## Appendix A Fermion-magnon matrix element in the strong-coupling limit
We show that to leading order in the $`U/t\mathrm{}`$ limit, the intraband fermion matrix elements reduce to order $`t/\mathrm{\Delta }`$, and the fermion-magnon matrix element explicitly reduces to order $`t`$, as within the $`tJ`$ model.
For the magnon energy and magnon amplitudes in the $`y^{}`$ and $`z^{}`$ directions we have
$`\omega _๐ช`$ $`=`$ $`3JS[(1\gamma _๐ช)(1+2\gamma _๐ช)]^{1/2}`$
$`๐ด_๐ช^2`$ $`=`$ $`{\displaystyle \frac{3JS}{4\omega _๐ช}}(1+2\gamma _๐ช)`$
$`๐ต_๐ช^2`$ $`=`$ $`{\displaystyle \frac{3JS}{4\omega _๐ช}}(1\gamma _๐ช)`$ (28)
$`\mathrm{where}\gamma _๐ช`$ $`=`$ $`{\displaystyle \frac{1}{3}}\left[\mathrm{cos}q_x+2\mathrm{cos}{\displaystyle \frac{q_x}{2}}\mathrm{cos}{\displaystyle \frac{\sqrt{3}}{2}}q_y\right].`$
For the lower $`()`$ and upper $`(+)`$ bands, the fermion amplitudes given in Eq. (11) reduce to
$$\alpha _๐ค^{}\pm \frac{1}{\sqrt{2}}\left(1\pm \frac{\xi _๐ค}{2\mathrm{\Delta }}\right),\beta _๐ค^{}+\frac{1}{\sqrt{2}}\left(1\frac{\xi _๐ค}{2\mathrm{\Delta }}\right).$$
(29)
Hence in terms of the spin raising and lowering operators defined below, the intraband fermion matrix elements
$`\sigma ^+`$ $``$ $`๐ค,l|\sigma _z^{}+i\sigma _y^{}|๐ค๐ช,m{\displaystyle \frac{\xi _{๐ค,l}}{\mathrm{\Delta }}}`$
$`\sigma ^{}`$ $``$ $`๐ค,l|\sigma _z^{}i\sigma _y^{}|๐ค๐ช,m{\displaystyle \frac{\xi _{๐ค๐ช,m}}{\mathrm{\Delta }}}`$ (30)
for lower band states are explicitly of order $`t/U`$.
In terms of corresponding magnon amplitudes
$$\mathrm{\Phi }_{๐ช,\lambda }^\pm \sqrt{3}[|๐ช,\lambda _z^{}\pm i|๐ช,\lambda _y^{}]=๐ต_{๐ช,\lambda }\pm ๐ด_{๐ช,\lambda }$$
(31)
for the advanced mode, the fermion-magnon matrix element
$`\sqrt{3}M`$ $`=`$ $`U[i\sigma _y^{}๐ด+\sigma _z^{}๐ต]`$
$`=`$ $`U[\sigma ^+\mathrm{\Phi }^{}+\sigma ^{}\mathrm{\Phi }^+]/2`$
$`=`$ $`(\xi _{๐ค,l}\xi _{๐ค๐ช,m})๐ด_{๐ช,\lambda }+(\xi _{๐ค,l}+\xi _{๐ค๐ช,m})๐ต_{๐ช,\lambda }`$
$`=`$ $`\xi _{๐ค,l}\mathrm{\Phi }_{๐ช,\lambda }^{}+\xi _{๐ค๐ช,m}\mathrm{\Phi }_{๐ช,\lambda }^+(\mathrm{where}m=l\lambda )`$
explicitly reduces to order $`t`$.
The above expression for the fermion-magnon matrix element has exactly same structure as the result $`\sqrt{3}tM(k,q)=\sqrt{3}t[v_๐ชh_๐คu_๐ชh_{๐ค+๐ช}]`$ obtained within the $`tJ`$ model.dombre Indeed, magnon amplitudes $`u_๐ช`$ and $`v_๐ช`$ in Ref. exactly correspond to $`๐ต_๐ช\pm ๐ด_๐ช`$, and $`(\sqrt{3}t)h_๐ค\sqrt{3}t_{\widehat{\delta }}\mathrm{sin}๐ค.\widehat{\delta }=\xi _{๐ค๐}`$. However, the Goldstone-mode contribution appears to be different. While our fermion-magnon matrix element vanishes for the Goldstone mode $`q0,\lambda =0`$, for which $`\gamma _๐ช1`$ and the magnon amplitude $`๐ต0`$ representing rigid spin rotation about $`z`$ axis, the matrix element $`M(๐ค,๐ช)`$ in Ref. does not vanish. The other two Goldstone modes $`q0,\lambda =\pm 1`$, for which $`\gamma _๐ช1/2`$ and $`๐ด0`$, do yield non-vanishing matrix elements with $`M^21/q`$, although resulting in a finite contribution $`q๐q/q`$ from long-wavelength modes. Thus, for the triangular-lattice AF, long-wavelength modes do contribute to the fermion-magnon scattering.
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# Multivariable Alexander invariants of hypersurface complements
## 1. Introduction
Alexander invariants in the form of Alexander modules, characteristic varieties and Alexander polynomials have been recently intensively studied, in particular in relation to the twisted cohomology of hypersurface arrangement complements, see for instance , , , , , , , , , , , , , .
In section ยง2, after giving the basic definitions introducing the Alexander modules $`A^q(๐ฐ)`$ of an affine hypersurface arrangement complement $`๐ฐ`$, we investigate in Proposition 2.4 the relation between the first non-trivial Alexander polynomial in one variable and the corresponding Alexander polynomial in several variables. Proposition 2.5 expresses the relation between the characteristic varieties defined using the Fitting ideals and the characteristic varieties defined using the jumping loci of the cohomology with rank one local coefficients. Example 2.8 treats the simplest local situations: the normal crossing case and the case of isolated non-normal crossing singularities, whose study was initiated by A. Libgober in .
In section ยง3, Theorem 3.1 relates the Alexander invariants of the affine hypersurface arrangement complement $`๐ฐ=^{n+1}X`$ to the Alexander invariants of the complement $`๐ฐ_{\mathrm{}}`$ of the corresponding link at infinity. Theorems 3.2, 3.6 and Corollary 3.5 estimate the support of the Alexander modules $`A^q(๐ฐ)`$ in terms of local properties of the projective closure $`V=\overline{X}`$.
In section ยง4, we recall and slightly extend the idea of combining the degeneration of the Hodge to de Rham spectral sequence to the purity of some cohomology groups (used first by Esnault, Schechtman and Viehweg in and by Schechtman, Terao and Varchenko in ), see Corollary 4.1 and Proposition 4.5. Examples 4.8 and 4.10 illustrate this approach by looking at some arrangements of lines and conics in the plane. Though these examples may be treated using the results by Cogolludo in , we feel that our approach is more general and hence more likely to extend to other situations.
In the last section we consider the complement $`๐ฐ_0`$ of an arbitrary projective hypersurface arrangement $`V`$, and, after a short general discussion, we revisit from a new perspective a useful result by Randell saying what happens to the twisted cohomology of a plane curve complement when we add an extra line, see Corollary 5.1. Coming back to dimension $`n2`$, Example 5.2 discusses the already interesting case when $`V`$ is irreducible and has only isolated singularities. This case leads, in particular, to examples where for $`m=n,n+1`$ and some rank one local coefficients $`_\beta `$ on $`๐ฐ_0`$ one has
$$dimH^m(๐ฐ_0,_\beta )>dimH^m(๐ฐ_0,).$$
By the minimality property of hyperplane arrangement complements, it is known that the above inequality is impossible for such type of complements, . We conclude by a detailed study of the case when $`V`$ has two irreducible components, each of them having only isolated singularities.
Throughout the paper we usually work with complex coefficients $``$, although the study of finite field coefficients is very important, due for instance to torsion open questions, see , . Our choice is imposed by the analytic tools used in the last two sections. Most of the results in the previous sections hold over arbitrary fields.
## 2. Multivariable Alexander invariants
### 2.1. Algebraic Preliminaries
Let $`R`$ be a commutative ring with unit, which is Noetherian and a unique factorization domain (e.g. the ring of complex Laurent polynomials in $`s`$ variables, $`s1`$). Let $`A`$ be a finitely generated $`R`$-module, and $`M`$ a $`(n\times m)`$ presentation matrix of $`A`$ associated to an exact sequence
$$R^mR^nA0.$$
The $`i`$-th *elementary ideal* $`_i(A)`$ of $`A`$ is the ideal in $`R`$ generated by the $`(ni)\times (ni)`$ minor determinants of $`M`$, with the convention that $`_i(A)=R`$ if $`in`$, and $`_i(A)=0`$ if $`ni>m`$. Let $`\mathrm{\Delta }_i(A)`$ be the generator of the smallest principal ideal in $`R`$ containing $`_i(A)`$, i.e. the greatest common divisor of all elements of $`_i(A)`$. $`\mathrm{\Delta }_i(A)`$ is called the $`i`$-th *characteristic polynomial* of $`A`$. Note that $`\mathrm{\Delta }_{i+1}(A)`$ divides $`\mathrm{\Delta }_i(A)`$ in $`R`$ for all $`i`$ since $`_i(A)_{i+1}(A)`$. In particular, if $`R`$ is a principal ideal domain (e.g. the ring of complex Laurent polynomials in one variable), then $`_i(A)`$ is a principal ideal generated exactly by $`\mathrm{\Delta }_i(A)`$.
As an example, for any ring $`R`$, assume that $`A=R^sR/(\lambda _1)\mathrm{}R/(\lambda _r)`$, where $`\lambda _j`$ ($`j=1,2,\mathrm{},r`$) are non-zero elements in $`R`$ such that $`\lambda _{j+1}|\lambda _j`$. Then we have $`\mathrm{\Delta }_i(A)`$ is $`0`$, $`\lambda _{is+1}\mathrm{}\lambda _r`$, or $`1`$, according to whether $`0is1`$, $`sis+r1`$, or $`s+ri`$.
The *support* $`\text{Supp}(A)`$ of $`A`$ is the reduced sub-scheme of $`\text{Spec}(R)`$ defined by (the *order ideal*) $`_0(A)`$. Since
$$\sqrt{_0(A)}=\sqrt{Ann(A)}$$
this is the usual notion of support in algebraic geometry based on the *annihilator ideal* $`\text{Ann}(A)`$ of the module $`A`$. In particular, for a prime ideal $`PR`$, $`P\text{Supp}(A)`$ if and only if the localized module $`A_P`$ is non-zero.
The support $`\text{Supp}(A)`$ is also called the *first characteristic variety* of $`A`$, and we define the $`i`$-th *characteristic variety* $`V_i(A)`$ of $`A`$ to be the reduced sub-scheme of $`\text{Spec}(R)`$ defined by the ($`i`$-th *Fitting ideal*) ideal $`_{i1}(A)`$.
Note that $`\text{codim}V_i(A)>1`$ implies $`\mathrm{\Delta }_{i1}(A)=1`$, i.e. the corresponding Alexander polynomial carries no information.
All definitions above are independent (up to multiplication by a unit of $`R`$) of the choices involved, thus the characteristic varieties and polynomials of $`A`$ are invariants of the $`R`$-isomorphism type of $`A`$.
We state for future reference the following โdivisibilityโ properties of the polynomials and characteristic varieties (for proofs, see and ):
###### Lemma 2.1.
* If $`A`$, $`B`$ are finitely generated $`R`$-modules, then: $`\mathrm{\Delta }_0(AB)=\mathrm{\Delta }_0(A)\times \mathrm{\Delta }_0(B)`$.
* If $`A`$ and $`B`$ are finitely generated $`R`$-modules then:
$`\mathrm{Supp}(A_RB)=\mathrm{Supp}(A)\mathrm{Supp}(B)`$.
* If $`A`$ is a submodule of $`B`$, then for all $`i`$, $`\mathrm{\Delta }_i(A)`$ divides $`\mathrm{\Delta }_i(B)`$.
* If $`0ABC0`$ is a short exact sequence of finitely generated $`R`$-modules, then the following hold:
1. $`\mathrm{\Delta }_0(B)=\mathrm{\Delta }_0(A)\times \mathrm{\Delta }_0(C)`$;
2. for all $`i`$, $`\mathrm{\Delta }_i(B)`$ divides $`\mathrm{\Delta }_i(A)\times \mathrm{\Delta }_0(C)`$;
3. If $`\mathrm{\Delta }_0(C)=1`$, then $`\mathrm{\Delta }_i(A)=\mathrm{\Delta }_i(B)`$ for all $`i`$;
4. $`\mathrm{Supp}(B)=\mathrm{Supp}(A)\mathrm{Supp}(C)`$;
5. For $`i2`$: $`V_i(C)V_i(B)V_i(C)(V_{i1}(C)\mathrm{Supp}(A))`$.
### 2.2. Alexander Invariants of Hypersurface Complements
Let $`V`$ be a reduced hypersurface in $`^{n+1}`$, defined by a homogeneous equation: $`f=f_1\mathrm{}f_s=0`$, where $`f_i`$ are the irreducible factors of $`f`$, and $`V_i=\{f_i=0\}`$ the irreducible components of $`V`$. We fix a hyperplane $`H`$ in $`^{n+1}`$ which we call โthe hyperplane at infinityโ. Let $`๐ฐ`$ be the (affine) hypersurface complement $`๐ฐ=^{n+1}(VH)`$. (Alternatively, $`๐ฐ`$ may be regarded as the complement of a hypersurface in the affine space $`^{n+1}`$.) Then $`H_1(๐ฐ)^s`$ (, (4.1.3), (4.1.4)), generated by the meridian loops $`\gamma _i`$ about the non-singular part of each irreducible component $`V_i`$, for $`i=1,\mathrm{},s`$. If $`\gamma _{\mathrm{}}`$ denotes the meridian about the hyperplane at infinity, then in $`H_1(๐ฐ)`$ there is a relation: $`\gamma _{\mathrm{}}+d_i\gamma _i=0`$, where $`d_i=deg(V_i)`$.
Note that $`๐ฐ`$ is affine, therefore has the homotopy type of a finite CW complex. Let $`๐ฐ^{ab}`$ be the universal abelian cover of $`๐ฐ`$, i.e. the covering associated to the commutator subgroup of $`\pi _1(๐ฐ)`$, or equivalently, the covering associated to the kernel of the linking number homomorphism $`lk:\pi _1(๐ฐ)^s`$, which maps a loop $`\alpha `$ to $`(\text{lk}(\alpha ,V_1d_1H),\mathrm{},\text{lk}(\alpha ,V_sd_sH))`$. The group of covering transformations of $`๐ฐ^{ab}`$ is isomorphic to $`^s`$ and acts on the covering space. By choosing fixed lifts of the cells of $`๐ฐ`$ to $`๐ฐ^{ab}`$, we obtain a free basis for $`C_{}`$, the cellular cell complex of $`๐ฐ^{ab}`$, as a $`[^s]`$-module. The isomorphism determined by the meridians $`\{\gamma _i\}`$ enables us to identify $`[^s]`$ with $`[t_1,t_1^1,\mathrm{},t_s,t_s^1]`$, the ring of integral Laurent polynomials in $`s`$ variables. When $`s=1`$ we set $`t_1=t`$.
For reasons that will become transparent later, our base ring will always be the ring of complex Laurent polynomials in $`s`$ variables, $`[t_1,t_1^1,\mathrm{},t_s,t_s^1]`$, which we denote by $`R_s`$. Note that $`R_s`$ is a regular Noetherian domain, and in particular it is factorial. As a groups ring, $`R_s`$ has a natural involution denoted by an overbar, sending each $`t_i`$ to $`\overline{t}_i:=t_i^1`$. To an $`R_s`$-module $`A`$, we associate the conjugate $`R_s`$-module, still denoted by $`A`$, with the same underlying abelian group but with the $`R_s`$-action given by $`(r,a)\overline{r}a`$, for $`aA`$ and $`rR_s`$.
###### Remark 2.2.
Though the ring $`R_s`$ is commutative, it should be regarded as a quotient ring of $`[\pi _1(๐ฐ)]`$, which is non-commutative in general. Because of that, one should be careful to distinguish the right from the left $`R_s`$-modules. If, for instance, $`A`$ is a left $`R_s`$-module, then the associated right $`R_s`$-module is the module conjugate to $`A`$, whose module structure is given by
$$ar:=\overline{r}a$$
for all $`aA`$ and $`rR_s`$. This corresponds to regarding any left $`[\pi _1(๐ฐ)]`$-module $`A`$ as a right $`[\pi _1(๐ฐ)]`$-module by setting $`a\gamma =\gamma ^1a`$, for all $`aA`$ and $`\gamma \pi _1(๐ฐ)`$, and extending by linearity. Following , p. 97, we regard in this paper $`C_{}^0=C_{}`$ as a complex of right $`R_s`$-modules.
Define a local coefficient system $``$ on $`๐ฐ`$, with stalk $`R_s`$ and action of a loop $`\alpha \pi _1(๐ฐ)`$ determined by (left) multiplication by $`_{j=1}^s(t_j)^{\text{lk}(\alpha ,V_jd_jH)}`$. In particular, the action of the meridian $`\gamma _i`$ is given by multiplication by $`t_i`$. Let $`^{}`$ be the dual local system, whose stalk at a point $`y๐ฐ`$ is $`_y^{}:=\text{Hom}(_y,R_s)`$, and let $`\alpha \pi _1(๐ฐ,y)`$ act on $`\phi _y^{}`$ by:
$$(\alpha \phi )(m):=\phi (\alpha ^1m),m_y.$$
We denote by $`\overline{}`$ the local system obtained from $``$ by composing all module structures with the involution of $`R_s`$ (i.e. by changing the stalks of $``$ from left into right $`R_s`$-modules). The perfect pairing
$$\overline{}_{R_s}R_s$$
given by
$$(f,g)\overline{f}g$$
on the stalk over a basepoint, tell us that there is an isomorphism of local systems on $`๐ฐ`$:
$$^{}\overline{}.$$
The *universal homology k-th Alexander invariant* $`A_k(๐ฐ)`$ of $`๐ฐ`$ is by definition the $`R_s`$-module $`H_k(C_{}^0)`$, or equivalently $`H_k(๐ฐ;)`$. This is the group $`H_k(๐ฐ^{ab};)`$ considered as a $`R_s`$-module via the covering transformations (see , Example 3H.2). Similarly, the *universal cohomology k-th Alexander invariant* $`A^k(๐ฐ)`$ of $`๐ฐ`$ is by definition the $`k`$-th cohomology module of the dual complex $`\text{Hom}_{R_s}(C_{}^0,R_s)`$. Here $`R_s`$ is considered with the induced right $`R_s`$-module structure as explained in Remark 2.2. Based on our previous considerations on local systems, $`A^k(๐ฐ)`$ is just $`H^k(๐ฐ;^{})`$. This may be also regarded as the $`k`$-th cohomology with compact support and complex coefficients of $`๐ฐ_b^{ab}`$, where $`๐ฐ_b`$ is the compact manifold with boundary obtained from $`^{n+1}`$ by removing a small open regular neighborhood of the divisor $`VH`$ (compare , Prop. 3H.5).
Note that, since $`๐ฐ`$ is a $`(n+1)`$-dimensional affine variety, the modules $`A^k(๐ฐ)`$ and resp. $`A_k(๐ฐ)`$ are trivial for $`k>n+1`$. Moreover, since the stalks of $``$ are torsion-free, $`A_{n+1}(๐ฐ)`$ is also a torsion-free $`R_s`$-module (see , Example 6.0.6).
As in the classical knot theory, by using a deformation retract argument, one could define the universal abelian invariants above after replacing $`๐ฐ`$ by the manifold with boundary $`๐ฐ_b`$, obtained from $`^{n+1}`$ by removing a small open regular neighborhood of the divisor $`VH`$. Now, since the chain complex $`C_{}(๐ฐ_b^{ab})`$ is of finite type, and since $`R_s`$ is Noetherian, this implies that all these universal Alexander modules are finitely generated. Hence their characteristic varieties and polynomials are well-defined. The associated characteristic varieties, in particular the supports, become sub-varieties of the $`s`$-dimensional torus $`๐^s=(^{})^s`$, which is regarded as the set of closed points in $`\text{Spec}(R_s)`$. More precisely, for $`\lambda =(\lambda _1,\mathrm{},\lambda _s)๐^s`$, we denote by $`m_\lambda `$ the corresponding maximal ideal in $`R_s`$ and by $`_\lambda `$ the quotient $`R_s/m_\lambda R_s`$. This quotient is isomorphic to $``$ and the canonical projection
(2.1)
$$\rho _\lambda :R_sR_s/m_\lambda R_s=_\lambda $$
corresponds to replacing $`t_j`$ by $`\lambda _j`$ for $`j=1,\mathrm{},s.`$ Here we regard $`_\lambda `$ as a (left) $`R_s`$-module, with an involution given by the complex conjugation (which is compatible with the one induced from $`R_s`$ since $`\lambda _j๐^1`$).
If $`A`$ is an $`R_s`$-module, we denote be $`A_\lambda `$ the localization of $`A`$ at the maximal ideal $`m_\lambda `$. For $`A=R_s`$, we use the simpler notation $`R_\lambda `$ when there is no danger of confusion. If $`A`$ is of finite type, then $`A=0`$ if and only if $`A_\lambda =0`$ for all $`\lambda ๐^s`$. More precisely
$$\text{Supp}(A)=\{\lambda ๐^s;A_\lambda 0\}$$
In particular $`A_0(๐ฐ)=_\mathrm{๐}`$, where $`\mathrm{๐}=(1,\mathrm{},1)`$ and hence
(2.2)
$$\text{Supp}(A_0(๐ฐ))=\{\mathrm{๐}\}.$$
We denote by $`V_{i,k}(๐ฐ)`$ the $`i`$-th characteristic variety associated to the homological Alexander module $`A_k(๐ฐ)`$, and similarly denote by $`\mathrm{\Delta }_{i,k}(๐ฐ)`$ the associated characteristic polynomials. The notations $`V^{i,k}(๐ฐ)`$ and $`\mathrm{\Delta }^{i,k}(๐ฐ)`$ denote the similar objects associated to the cohomological Alexander invariants $`A^k(๐ฐ)`$.
### 2.3. Homology versus Cohomology Alexander Modules
It is natural to ask what are the relations between the homology and the cohomology universal Alexander modules. Or to find the relations between $`V_{i,k}(๐ฐ)`$ and $`V^{i,k}(๐ฐ)`$; and between $`\mathrm{\Delta }_{i,k}(๐ฐ)`$ and $`\mathrm{\Delta }^{i,k}(๐ฐ)`$.
Some answers to this question can be given as follows. The cohomology modules may be related to the homology modules by the Universal Coefficient spectral sequence (see , p.20 or , Thm. 2.3).
(2.3)
$$\text{Ext}_{R_s}^q(A_p(๐ฐ),R_s)A^{p+q}(๐ฐ).$$
Using the exactness of the localization (see , p. 76), we get the following spectral sequence for any $`\lambda ๐^s`$.
(2.4)
$$\text{Ext}_{R_\lambda }^q(A_p(๐ฐ)_\lambda ,R_\lambda )A^{p+q}(๐ฐ)_\lambda .$$
For a fixed $`\lambda ๐^s`$, we define
(2.5)
$$k(\lambda )=\text{min}\{m;A_m(๐ฐ)_\lambda 0\}.$$
Then the spectral sequence 2.4 implies the following.
###### Proposition 2.3.
For any $`\lambda ๐^s`$, $`A^k(๐ฐ)_\lambda =0`$ for $`k<k(\lambda )`$ and
(2.6)
$$A^{k(\lambda )}(๐ฐ)_\lambda =\text{Hom}(A_{k(\lambda )}(๐ฐ)_\lambda ,R_\lambda )$$
This equality shows in particular that one may have $`A^{k(\lambda )}(๐ฐ)_\lambda =0`$, even when $`A_{k(\lambda )}(๐ฐ)_\lambda 0`$, e.g. when the last module is torsion, which is often the case, e.g. see 2.2.
### 2.4. Multivariable versus one variable Alexander Modules
Consider a family of integral weights $`๐=(e_1,\mathrm{},e_s)^s`$, and let
$$q:=\text{g.c.d.}(e_1,\mathrm{},e_s).$$
Consider the morphism $`p(๐):R_sR_1`$ defined by $`t_it^{e_i}`$, inducing a (left) $`R_s`$-module structure on $`R_1`$. Let $`(๐)`$ be the local system on $`๐ฐ`$ with stalk $`R_1`$ and monodromy action for a loop $`\alpha \pi _1(๐ฐ)`$ given by multiplication by $`t^{{\scriptscriptstyle }e_j\text{lk}(\alpha ,V_jd_jH)}`$.
The corresponding homology groups $`H_k(๐ฐ,(๐))=H_k(C_{}^0_{R_s}R_1)`$ are finite type $`R_1`$-modules, and hence they have associated characteristic varieties $`V_{i,k}(๐ฐ,๐)`$ and Alexander polynomials $`\mathrm{\Delta }_{i,k}(๐ฐ,๐)`$.
It is natural to ask under which conditions the equalities
$$\mathrm{\Delta }_{i,k}(๐ฐ,๐)(t)=(t^q1)\mathrm{\Delta }_{i,k}(๐ฐ)(t^{e_1},\mathrm{},t^{e_s})$$
do hold? Something like this works in classical knot theory, more precisely for oriented multilinks in $`S^3`$ with at least 2 components, where the case $`i=0`$, $`k=1`$ is considered (see , Prop. 5.1, and also , Lemma 10.1 for the case of weight $`(1,\mathrm{},1)`$).
For the weight $`\mathrm{๐}=(1,1,\mathrm{},1)`$, we call the coresponding Alexander polynomials the usual (or, univariable ) Alexander polynomials and we denote them by $`\mathrm{\Delta }_{i,k}^T(๐ฐ)`$ (see below for some explanation).
If the equality in Question 2 holds for all but finitely many multi-indices $`๐`$, then the 1-variable polynomials $`\mathrm{\Delta }_{i,k}(๐ฐ,๐)`$ determine (up-to a unit in $`R_s`$) the multi-variable polynomial $`\mathrm{\Delta }_{i,k}(๐ฐ)`$ (see , Lemma 2.2).
Some insight into this question can be obtained as follows. We consider only the simplest case, namely $`๐=\mathrm{๐}`$, and leave the other cases to the interested reader.
Note that the universal abelian covering $`๐ฐ^{ab}๐ฐ`$ corresponds to the kernel $`K^{ab}`$ of the abelianization morphism
$$\pi _1(๐ฐ)H_1(๐ฐ).$$
The total linking number covering $`๐ฐ^T๐ฐ`$ corresponds to the kernel $`K^T`$ of the morphism
$$\pi _1(๐ฐ)H_1(๐ฐ)=^s$$
where the second morphism is $`c_j\gamma _jc_j`$. It follows that $`๐ฐ^{ab}๐ฐ^T`$ is a covering with deck transformation group $`G=K^T/K^{ab}`$ identified to the subgroup
$$\{c^s;c_j=0\}.$$
The complex $`C_{}^0`$ is a complex of free $`R_s`$-modules of finite rank and the derivatives are $`R_s`$-liniar. It follows that we can regard this complex as being a complex $`๐_{}^0`$ of free $`๐ช_{๐^s}`$-modules on the affine variety $`๐^s`$.
Since $`๐ฐ^T=๐ฐ^{ab}/G`$, it follows that the complex of singular chains of $`๐ฐ^T`$ is
(2.7)
$$C_{}(๐ฐ^T)=C_{}(๐ฐ^{ab})_G=(C_{}^0)_G$$
(see , p.204). Here
(2.8)
$$(C_p^0)_G=C_p^0/<gmm;gG,mC_p^0>.$$
Using the fact that the group $`G`$ is generated by the elements having an 1 as the $`i`$-th coordinate, a $`1`$ as the $`j`$-th coordinate (for $`i<j`$) and all the other coordinates zero, we see that $`(C_p^0)_G`$ is the quotient of $`C_p^0`$ by the submodule
$$<(t_it_j)m;mC_p^0>.$$
It follows that the associated sheaf $`(๐_p^0)_G`$ is just the restriction (as a coherent sheaf) of $`๐_p^0`$ to the 1-dimensional subtorus $`S=\{(t,t,\mathrm{},t)๐^s\}`$, i.e. $`(๐_p^0)_G=๐_p^0_{๐ช_{๐^s}}๐ช_S.`$ Unfortunately, the inclusion $`S๐^s`$ is not a flat morphism (see , p. 254), and hence the restriction to $`S`$ does not commute to taking homology.
However, by our discussion above
$$(C_p^0)_G=C_p^0_{R_s}R_1,$$
with the (left) $`R_s`$-module structure on $`R_1`$ induced by $`p(\mathrm{๐})`$. Use now the Kรผnneth spectral sequence (see , p.143), and get
(2.9)
$$E_{p,q}^2=Tor_p^{R_s}(A_q(๐ฐ),R_1)H_{p+q}((C_{}^0)_G)=A_{p+q}^T(๐ฐ).$$
For $`a๐^1=S=\{(t,t,\mathrm{},t)๐^s\}`$, we get by localization a new Kรผnneth spectral sequence, namely
(2.10)
$$E_{p,q}^2=Tor_p^{R_a}(A_q(๐ฐ)_a,R_{1,a})H_{p+q}((C_{}^0)_G)_a.$$
In particular we get the following.
###### Proposition 2.4.
For any $`a๐^1`$, $`A_k^T(๐ฐ)_a=0`$ for $`k<k(a)`$ and
(2.11)
$$A_{k(a)}(๐ฐ)_a_{R_a}R_{1,a}=A_{k(a)}^T(๐ฐ)_a$$
In particular, for any $`a๐^1=S`$, the multiplicity of the root $`t=a`$ in the polynomials $`\mathrm{\Delta }_{i,k(a)}^T(๐ฐ)(t)`$ and $`\mathrm{\Delta }_{i,k(a)}(๐ฐ)(t,\mathrm{},t)`$ is the same.
###### Proof.
To get the second claim, note that any presentation
$$R_a^mR_a^nA_{k(a)}(๐ฐ)_a0$$
yields by tensor product a presentation
$$R_{1,a}^mR_{1,a}^nA_{k(a)}^T(๐ฐ)_a0.$$
### 2.5. Characteristic varieties as jumping loci of rank-1 local systems
Let $`\lambda =(\lambda _1,\mathrm{},\lambda _s)๐^s`$ and denote by $`_\lambda `$ the local coefficient system on $`๐ฐ`$ with stalk $`=_\lambda `$ and action of a loop $`\alpha \pi _1(๐ฐ)`$ determined by multiplication by $`_{j=1}^s(\lambda _j)^{\text{lk}(\alpha ,V_jd_jH)}`$. We let $`_\lambda ^{}_{\lambda ^1}`$ be the dual local system, where $`\lambda ^1:=(\lambda _1^1,\mathrm{},\lambda _s^1)๐^s`$.
One can define new *topological characteristic varieties* by setting
$$V_{i,k}^t(๐ฐ)=\{\lambda ๐^s;\text{dim}H_k(๐ฐ,_\lambda )>i\}$$
and
$$V_t^{i,k}(๐ฐ)=\{\lambda ๐^s;\text{dim}H^k(๐ฐ,_\lambda )>i\}.$$
It is natural to investigate the relations between the two types of characteristic varieties. Some cases are considered in , .
Here is a general approach to this question. It is known that
$$H_k(๐ฐ,_\lambda )=H_k(C_{}^0_{R_s}_\lambda ).$$
Using the Kรผnneth spectral sequence, we get
(2.12)
$$E_{p,q}^2=Tor_p^{R_s}(A_q(๐ฐ),_\lambda )H_{p+q}(๐ฐ,_\lambda ).$$
Now since the localization is exact, the base change for Tor under $`R_sR_\lambda `$ (see , p. 144), yields a new spectral sequence
(2.13)
$$E_{p,q}^2=Tor_p^{R_\lambda }(A_q(๐ฐ)_\lambda ,_\lambda )H_{p+q}(๐ฐ,_\lambda ).$$
This proves the first claim of the next result.
###### Proposition 2.5.
For any point $`\lambda ๐^s`$, one has the following.
(i) $`\mathrm{min}\{m,H_m(๐ฐ,_\lambda )0\}=\mathrm{min}\{m,\lambda \mathrm{Supp}(A_m(๐ฐ))\}=k(\lambda ).`$
(ii) $`dimH_{k(\lambda )}(๐ฐ,_\lambda )=\mathrm{max}\{m,\lambda V_{m,k(\lambda )}(๐ฐ)\}`$.
###### Proof.
To prove the second claim, note that the spectral sequence 2.13 yields
$$H_{k(\lambda )}(๐ฐ,_\lambda )=A_{k(\lambda )}(๐ฐ)_\lambda /m_\lambda A_{k(\lambda )}(๐ฐ)_\lambda .$$
Let $`n`$ be the dimension of these two vector spaces. Then by Nakayamaโs Lemma, the module $`A_{k(\lambda )}(๐ฐ)_\lambda `$ is generated by $`n`$ elements over the local ring $`R_\lambda `$. In other words, there is presentation
$$R_\lambda ^mR_\lambda ^nA_{k(\lambda )}(๐ฐ)_\lambda 0.$$
Moreover, the first morphism is given by a matrix $`M`$ whose entries $`m_{ij}`$ are all in the maximal ideal $`m_\lambda `$. The second claim follows now by the definition of the characteristic varieties.
###### Remark 2.6.
Note that there is also a spectral sequence
(2.14)
$$E_2^{p,q}=\text{Ext}_{R_\lambda }^q(A_p(๐ฐ)_\lambda ,_\lambda )H^{p+q}(๐ฐ,_{\lambda ^1}).$$
Here $`_\lambda `$ is considered with the right $`R_s`$-module structure as indicated in remark 2.2. This is why in the abutment of the spectral sequence 2.14, we obtain cohomology with coefficients in the dual local system $`_\lambda ^{}_{\lambda ^1}`$.
The above spectral sequence yields that $`H^m(๐ฐ,_{\lambda ^1})=0`$ for $`m<k(\lambda )`$ and $`H^{k(\lambda )}(๐ฐ,_{\lambda ^1})=Hom_{R_\lambda }(A_{k(\lambda )}(๐ฐ)_\lambda ,_\lambda )`$. However
$$Hom_{R_\lambda }(A_{k(\lambda )}(๐ฐ)_\lambda ,_\lambda )=Hom_{}(A_{k(\lambda )}(๐ฐ)_\lambda /m_\lambda A_{k(\lambda )}(๐ฐ)_\lambda ,_\lambda )$$
and hence
(2.15)
$$H_{k(\lambda )}(๐ฐ,_\lambda )^{}=H^{k(\lambda )}(๐ฐ,_{\lambda ^1}),$$
(compare , p.50 and p. 69). The case $`k=1`$ of this useful formula was established in , Remark 5.2. Note that this formula holds over arbitrary fields, with the same proof as above.
###### Remark 2.7.
All the results in this section so far hold for the local setting as well, i.e. when $`๐ฐ`$ is the complement of a hypersurface germ in a small ball. The first part of the example below corresponds to the germ of a normal crossing divisor. The second part of the example below corresponds to isolated non-normal crossing divisors (for short INNC), see , , .
Similarly, instead of localizing at a point, one may localize along the hyperplane at infinity $`H`$, i.e. replace $`๐ฐ`$ by $`๐ฐ_{\mathrm{}}=๐ฐS_{\mathrm{}}`$, where $`S_{\mathrm{}}`$ is a large enough sphere in $`^{n+1}`$, see Theorem 3.1 below.
###### Example 2.8.
(i) Let $`๐ฐ=(^{})^s\times ^{n+1s}`$ for some integer $`0sn+1`$. Then the universal abelian covering $`๐ฐ^{ab}`$ is contractible and then $`A_0(๐ฐ)=_\mathrm{๐}`$ and $`A_k(๐ฐ)=0`$ for $`k>0`$. Therefore, by the spectral sequence 2.3 we get $`A^k(๐ฐ)\text{Ext}_{R_s}^k(_\mathrm{๐},R_s)`$ for all $`k0`$. Using the free resolution of $`_\mathrm{๐}`$ given by the Koszul complex of the regular sequence $`\{x_j=t_j1\}_{j=1,\mathrm{},s}`$ in the ring $`R_s`$ (, Cor. 4.5.5), we obtain that $`A^k(๐ฐ)=0`$ for $`ks`$ and $`A^s(๐ฐ)=_\mathrm{๐}`$ (, Ex. 4.5.2 and Cor 4.5.4). Therefore, for any $`\lambda \mathrm{๐}`$, Proposition 2.3 shows that the corresponding cohomology Alexander modules satisfy $`A^k(๐ฐ)_\lambda =0`$ for any $`k`$. Moreover $`H_k(๐ฐ,_\lambda )=H^k(๐ฐ,_{\lambda ^1})=0`$ for any $`k`$.
(ii) Let $`(Y,0)`$ be an INNC singularity at the origin of $`^{n+1}`$. Set $`๐ฐ(Y,0)=BY`$, where $`B`$ is a small open ball centered at the origin in $`^{n+1}`$. Assume that $`n2`$. Then the universal abelian cover $`๐ฐ(Y,0)^{ab}`$ of $`๐ฐ(Y,0)`$ is $`(n1)`$-connected, see Libgober . More precisely, it is a bouquet of $`n`$-spheres, see , and hence $`A_0(๐ฐ(Y,0))=_\mathrm{๐}`$ and $`A_k(๐ฐ(Y,0))=0`$ for $`kn`$. As in (i) above, we get $`A^k(๐ฐ(Y,0))\text{Ext}_{R_s}^k(_\mathrm{๐},R_s)`$ for all $`k<n`$. For $`\lambda \mathrm{๐}`$ this yields $`A^k(๐ฐ(Y,0))_\lambda =0`$ for $`k<n`$, and therefore $`H^k(๐ฐ(Y,0),_\lambda )=0`$ for any $`k<n`$.
## 3. Divisibility Results and Characteristic Varieties
In this section we give an algebraic-geometrical interpretation for the multi-variable Alexander invariants of the hypersurface complement, similar in flavor to the one-variable case described in , but see also the reformulation of these results in . We will use an approach based on the general theory of perverse sheaves, close to the one presented in (see also and ). Note that the supports and characteristic polynomials $`\mathrm{\Delta }_0`$ of the multi-variable Alexander modules are the analogue of the set of roots of the Alexander polynomials and respectively Alexander polynomials in the one-variable case (cf. , ).
The first result is an extension of , Theorem 3.2, to arbitrary hypersurface singularities . Let $`S_{\mathrm{}}`$ be a sphere of sufficiently large radius in $`^{n+1}=^{n+1}H`$ (or equivalently, the boundary of a sufficiently small tubular neighborhood of $`H`$ in $`^{n+1}`$). Let $`V_{\mathrm{}}=S_{\mathrm{}}V`$ be the link of $`V`$ at infinity, and $`๐ฐ_{\mathrm{}}=S_{\mathrm{}}V_{\mathrm{}}`$ its complement.
###### Theorem 3.1.
For all $`i`$, and all $`kn`$: $`V_{i,k}(๐ฐ)V_{i,k}(๐ฐ_{\mathrm{}})`$, and $`\mathrm{\Delta }_{i,k}(๐ฐ)|\mathrm{\Delta }_{i,k}(๐ฐ_{\mathrm{}})`$. Moreover, for $`k<n`$, these inclusions and divisibility conditions are replaced by equalities.
###### Proof.
The case $`n=1`$ is considered in . In fact in this situation one sets, for $`i1`$ and $`k1`$, $`V_{i,k}(๐ฐ_{\mathrm{}})`$ to be the $`k`$-th characteristic variety of the $`i`$-th homology module of the covering space of $`๐ฐ_{\mathrm{}}`$ corresponding to the kernel of the composition
$$\pi _1(๐ฐ_{\mathrm{}})\pi _1(๐ฐ)H_1(๐ฐ).$$
For $`n2`$, the theorem is an easy consequence of the Lefschetz hyperplane theorem. Indeed, as in the proof of Theorem 4.5 of , it follows that $`\pi _1(๐ฐ)\pi _1(๐ฐ_{\mathrm{}})`$, and more generally $`\pi _k(๐ฐ,๐ฐ_{\mathrm{}})0`$ for all $`kn`$. Therefore, the same is true for any covering, in particular for the universal abelian coverings: $`\pi _k(๐ฐ^{ab},๐ฐ_{\mathrm{}}^{ab})0`$ for all $`kn`$. Hence, by Hurewicz Theorem, the vanishing also holds for the relative homology groups, i.e., the maps of groups $`H_k(๐ฐ_{\mathrm{}}^{ab})H_k(๐ฐ^{ab})`$ are isomorphism for $`k<n`$ and onto for $`k=n`$. Since these maps are induced by an embedding (recall $`n2`$), the above are morphisms of modules over the ring of Laurent polynomials in $`s`$ variables. The statement of the theorem follows now from Lemma 2.1.
From now on to the end of this section, we will make the assumption that the hyperplane at infinity $`H`$ is *transversal* (in the stratified sense) to the hypersurface $`V`$. With this assumption, we show that the global cohomological Alexander invariants of the hypersurface complement are entirely determined by the degrees of the irreducible components on one hand, and by the local topological information encoded by the singularities of $`V`$ on the other hand. In particular, these invariants depend on the local type of singularities of the hypersurface.
First, we need some notations. Recall from ยง 2.2 that $`A^q(๐ฐ)H^q(๐ฐ,^{})`$. For $`xV`$, we let $`๐ฐ_x=๐ฐB_x`$, for $`B_x`$ a small open ball at $`x`$ in $`^{n+1}`$. Denote by $`_x`$ the restriction of the local coefficient system $``$ to $`๐ฐ_x`$. Then the groups $`H^{}(๐ฐ_x,_x^{})`$ inherit a $`R_s`$-module structure.
###### Theorem 3.2.
Let $`\lambda =(\lambda _1,\mathrm{},\lambda _s)๐^s`$ and $`ฯต_0`$. Fix an irreducible component $`V_1`$ of $`V`$, and assume that $`\lambda \mathrm{Supp}(H^q(๐ฐ_x,_x^{}))`$ for all $`q<n+1ฯต`$ and all points $`xV_1`$. Then $`\lambda \mathrm{Supp}(A^q(๐ฐ))`$ for all $`q<n+1ฯต`$.
###### Proof.
Let $`๐ฐ_1=^{n+1}V_1`$, and let $`i:๐ฐ๐ฐ_1`$ and $`j:๐ฐ_1^{n+1}`$ be the two inclusions. Then $`^{}[n+1]\text{Perv}(๐ฐ)`$, since $`๐ฐ`$ is smooth. Moreover $`:=Ri_{}(^{}[n+1])\text{Perv}(๐ฐ_1)`$, since $`i`$ is a quasi-finite affine morphism (see , Theorem 6.0.4). But $`๐ฐ_1`$ is affine $`(n+1)`$-dimensional, and $`\text{Perv}(๐ฐ_1)`$, therefore by Artinโs vanishing theorem for perverse sheaves (see , Corollary 6.0.4), the following hold:
$$^k(๐ฐ_1,)=0,\text{for all}k>0,$$
$$_c^k(๐ฐ_1,)=0,\text{for all}k<0.$$
Let $`a:^{n+1}point`$ be the constant map. Then:
$$^k(๐ฐ_1,)H^{k+n+1}(๐ฐ,^{})H^k(Ra_{}Rj_{})$$
and
$$_c^k(๐ฐ_1,)H^k(Ra_!Rj_!)$$
Note that since $`a`$ is a proper map, we have $`Ra_!=Ra_{}`$.
Now consider the canonical morphism $`Rj_!Rj_{}`$ and extend it to the distinguished triangle:
$$Rj_!Rj_{}๐ข\stackrel{[1]}{}$$
in $`D_c^b(^{n+1})`$. Since $`j^{}j_!idj^{}j_{}`$, the complex $`๐ข`$ is supported on $`V_1`$. Apply $`Ra_!=Ra_{}`$ to the above distinguished triangle and obtain:
$$Ra_!Rj_!Ra_{}Rj_{}Ra_{}๐ข\stackrel{[1]}{}$$
Upon applying the cohomology functor to this triangle, and using the above vanishing, we obtain that:
$$H^{k+n+1}(๐ฐ,^{})^k(^{n+1},๐ข)^k(V_1,๐ข)\text{for}k<1,$$
and $`H^n(๐ฐ,^{})`$ is a sub-module of $`^1(V_1,๐ข)`$.
Therefore, by Lemma 2.1, in order to prove the theorem it suffices to show that, under our assumptions, $`\lambda \text{Supp}(^k(V_1,๐ข))`$ for all $`k<ฯต`$. This follows from the local calculation and the hypercohomology spectral sequence. Indeed, for $`xV_1`$, we have:
$`^q(๐ข)_x`$ $``$ $`^q(Rj_{})_x^{q+n+1}(Rj_{}Ri_{}^{})_x^{q+n+1}(B_x,R(ji)_{}^{})`$
$``$ $`H^{q+n+1}(๐ฐ_x,_x^{})`$
where $`๐ฐ_x=๐ฐB_x`$, for $`B_x`$ a small open ball at $`x`$ in $`^{n+1}`$, and $`_x`$ is the restriction of the local coefficient system $``$ to $`๐ฐ_x`$. Therefore, for a fixed $`xV_1`$ the assumption that $`\lambda \text{Supp}(H^q(๐ฐ_x,_x^{}))`$ for all $`q<n+1ฯต`$ is equivalent to the assumption $`\lambda \text{Supp}(^q(๐ข)_x)`$ for all $`q<ฯต`$. Next note that $`^k(V_1,๐ข)`$ is the abutment of a spectral sequence with the $`E_2`$-term defined by $`E_2^{p,q}=H^p(V_1,^q(๐ข))`$. Moreover, if $`\lambda \text{Supp}(^q(๐ข)_x)`$ for all $`q<ฯต`$ and for all $`xV_1`$, then $`\lambda \text{Supp}(H^p(V_1,^q(๐ข))`$ for $`p+q=k<ฯต`$ (since $`E_2^{p,q}`$ is non-trivial only if $`p0`$). Thus, from the spectral sequence, it follows that $`\lambda \text{Supp}(^k(V_1,๐ข))`$ for all $`k<ฯต`$. This finishes the proof of the theorem.
###### Remark 3.3.
In order to show that the universal cohomological modules depend only on the local information around the singularities of the hypersurface, it suffices to observe that the modules $`H^{}(๐ฐ_x,_x^{})`$, $`xV_1`$, are entirely determined by the local universal homological Alexander modules at $`x`$.
In order to see this, we first introduce some notation: let $`๐ฐ_0`$ denote the hypersurface complement $`^{n+1}V`$, and for $`xV_1`$ we set $`๐ฐ_x^{}=๐ฐ_0B_x`$, for $`B_x`$ a small open ball at $`x`$ in $`^{n+1}`$. Note that $`H_1(๐ฐ_x^{})=^k`$, where $`k`$ is the number of irreducible components of the hypersurface singularity germ $`(V,x)`$ (cf. , p.103). Let $`๐ฐ_x^{ab}`$ and $`(๐ฐ_x^{})^{ab}`$ be the universal abelian covers of $`๐ฐ_x`$ and $`๐ฐ_x^{}`$, respectively, and denote by $`A_{}(๐ฐ_x)`$ and respectively $`A_{}(๐ฐ_x^{})`$ the associated universal homological Alexander modules. The modules $`A_{}(๐ฐ_x^{})`$ will be called the local universal homological Alexander modules at $`x`$, as they depend only on the singularity germ $`(V,x)`$.
If $`i_x:๐ฐ_x๐ฐ`$ denotes the inclusion map, then the local system $`_x`$ on $`๐ฐ_x`$ is induced via the composition of maps
$$\varphi :\pi _1(๐ฐ_x)\stackrel{(i_x)_\mathrm{\#}}{}\pi _1(๐ฐ)\stackrel{\text{lk}}{}H_1(๐ฐ)\text{Aut}(R_s)$$
On the other hand, by the naturality of the Hurewicz morphism, $`\varphi `$ factors through $`\text{lk}_x:\pi _1(๐ฐ_x)H_1(๐ฐ_x)`$, $`R_s`$ becoming in this way a (left) $`[H_1(๐ฐ_x)]`$-module. Then, by p.50, it follows that $`H^{}(๐ฐ_x,_x^{})`$ is the homology of the equivariant Hom:
$$C^{}(๐ฐ_x,_x^{})=\text{Hom}_{[H_1(๐ฐ_x)]}(C_{}^0(๐ฐ_x^{ab}),R_s),$$
where $`R_s`$ is regarded now as a right $`[H_1(๐ฐ_x)]`$-module using the involution on the group ring as in Remark 2.2, and as a left $`R_s`$-module. By , p.6, there is a spectral sequence converging to $`H^{}(๐ฐ_x,_x^{})`$ with
(3.1)
$$E_2^{p,q}=\text{Ext}_{[H_1(๐ฐ_x)]}^q(A_p(๐ฐ_x),R_s).$$
In order to fully justify our claim, it remains to relate the local universal Alexander invariants $`A_{}(๐ฐ_x^{})`$ to the modules $`A_{}(๐ฐ_x)`$, at points $`xV_1`$.
For points $`xV_1(V_1H)`$ we have $`๐ฐ_x^{}=๐ฐ_x`$, thus our claim follows for such points by the above spectral sequence.
If $`xV_1H`$ then due to the transversality assumption, itโs easy to see that $`๐ฐ_x`$ is homotopy equivalent to $`๐ฐ_x^{}\times S^1`$. It follows that $`๐ฐ_x^{ab}(๐ฐ_x^{})^{ab}\times `$, thus by the homological Kรผnneth formula we obtain that the group $`A_p(๐ฐ_x)`$ is isomorphic to $`H_p((๐ฐ_x^{})^{ab},)H_0(,)A_p(๐ฐ_x^{})`$. When regarded as a $`[H_1(๐ฐ_x)]`$-module, the isomorphism can be written as (see , Prop. 1.8):
$$A_p(๐ฐ_x)(A_p(๐ฐ_x^{})_{[H_1(๐ฐ_x^{})]}[H_1(๐ฐ_x)])_{[]}.$$
Together with the spectral sequence 3.1 this finishes the proof of the claim.
###### Remark 3.4.
If $`S`$ is an $`s`$-dimensional stratum in a Whitney stratification of $`V`$ such that $`xS`$, then $`A_p(๐ฐ_x^{})=0`$ if $`p>ns`$. Indeed, $`๐ฐ_x^{}`$ has the homotopy type of the link complement $`S_x^{2n2s+1}L_x`$, where $`S_x^{2n2s+1}`$ is a small sphere at $`x`$ in a submanifold of $`^{n+1}`$ which meets $`S`$ transversally at $`x`$ (and no other point), and $`(S_x^{2n2s+1},L_x)`$ is the link pair of the stratum $`S`$ in the pair $`(^{n+1},V)`$. Since $`S_x^{2n2s+1}L_x`$ admits a cyclic cover which has the homotopy type of a CW complex of dimension $`ns`$ (i.e. the fiber of the Milnor fibration associated to the algebraic link $`(S_x^{2n2s+1},L_x)`$), it follows that the universal abelian cover $`(๐ฐ_x^{})^{ab}`$ has the homotopy type of a $`(ns)`$-dimensional CW complex, thus proving the claim.
The following consequence of Theorem 3.2, Remark 3.3, and of Example 2.8 is similar to some results in , , .
###### Corollary 3.5.
(i) (Case $`ฯต=0`$) With the notation in the above theorem, assume in addition that $`V`$ is a normal crossing divisor at any point of the component $`V_1`$. Then $`\mathrm{Supp}(A^k(๐ฐ))\{\mathrm{๐}\}`$ for any $`k<n+1.`$
(ii) (Case $`ฯต=1`$) With the notation in the above theorem, assume in addition that $`V`$ is an INNC divisor at any point of the component $`V_1`$. Then $`\mathrm{Supp}(A^k(๐ฐ))\{\mathrm{๐}\}`$ for any $`k<n.`$
Using a similar argument (see also ) we obtain the following result.
###### Theorem 3.6.
Assume that the hypersurface $`V`$ is transversal (in the stratified sense) to the hyperplane at infinity $`H`$. Then for $`kn`$, $`\mathrm{Supp}(A^k(๐ฐ))`$ is contained in the zero set of the polynomial $`t_1^{d_1}\mathrm{}t_s^{d_s}1`$, thus has positive codimension in $`๐^s`$.
The positive codimension property of supports in the universal abelian case should be regarded as the analogue of the torsion property in the infinite cyclic case (cf. , ). Example 5.5 below shows that transversality except finitely many points is not enough to get Theorem 3.6.
###### Proof.
As in the proof of the previous theorem, after replacing $`๐ฐ_1`$ by the affine space $`^{n+1}=^{n+1}H`$, it follows that for $`k1`$, $`H^{k+n+1}(๐ฐ,^{})`$ is a sub-module of $`^k(^{n+1},๐ข)`$, where $`๐ข`$ is now a complex of sheaves supported on $`H`$. Therefore, by Lemma 2.1, it suffices to prove the theorem for the supports of the modules $`^k(H,๐ข)`$ with $`k1`$.
As in the previous theorem, for $`xH`$, the local calculation on stalks yields $`^q(๐ข)_xH^{q+n+1}(๐ฐ_x,_x^{})`$, where $`๐ฐ_x=๐ฐB_x`$, for $`B_x`$ a small open ball at $`x`$ in $`^{n+1}`$. If $`xHHV`$, then $`๐ฐ_x`$ is homotopy equivalent to $`^{}`$, and the corresponding local system $`_x^{}`$ is defined by the action of $`\gamma _{\mathrm{}}`$, i.e. by multiplication by $`_{j=1}^s(t_j)^{d_j}`$. On the other hand, if $`xVH`$, then due to the transversality assumption, $`๐ฐ_x`$ is homotopy equivalent to a product $`(B_x^{}VB_x^{})\times ^{}`$, with $`B_x^{}`$ a small open ball centered at $`x`$ in $`H`$, and the local system $`_x^{}`$ is an external tensor product, the second factor being defined by the multiplication by $`_{j=1}^s(t_j)^{d_j}`$. Thus, by the Kunneth spectral sequence, the stalk cohomology groups of $`๐ข`$ along $`H`$, i.e. $`^q(๐ข)_{xH}`$, have supports contained in the zero set of the polynomial $`t_1^{d_1}\mathrm{}t_s^{d_s}1`$. Then by the hypercohomology spectral sequence, the same is true for the supports of the hypercohomology groups $`^k(H,๐ข)`$.
## 4. Explicit Computations via Logarithmic Connections
We review a general method used to determine the characteristic varieties in the case of hyperplane arrangements, see and , and show that essentially the same method applies to more general situations as well.
Let $`\pi :(Z,D)(^{n+1},VH)`$ be an embedded resolution of singularities for the reduced divisor $`VH`$. In particular
(i) $`D`$ is a normal crossing divisor with smooth irreducible components;
(ii) $`\pi :ZD๐ฐ`$ is an isomorphism.
In this setting there is a Hodge-Deligne spectral sequence
(4.1)
$$E_1^{p,q}=H^q(Z,\mathrm{\Omega }_Z^p(logD))H^{p+q}(๐ฐ,)$$
degenerating at $`E_1`$ and inducing the Hodge filtration $`F`$ of the Deligne mixed Hodge structure on $`H^{p+q}(๐ฐ,)`$, see .
###### Corollary 4.1.
If the Deligne mixed Hodge structure on some cohomology space $`H^m(๐ฐ)`$ is pure of type $`(m,m)`$, then
(i) $`H^0(Z,\mathrm{\Omega }_Z^m(logD))=H^m(๐ฐ)`$ and
(ii) $`H^q(Z,\mathrm{\Omega }_Z^p(logD))=0`$ for $`p+q=m`$ and $`q>0`$.
We list below several cases when this property holds.
###### Example 4.2.
(a) When $`V`$ is a hyperplane arrangement, the cohomology space $`H^m(๐ฐ)`$ is pure of type $`(m,m)`$ for all $`m0`$, see .
(b) When $`V`$ is a smooth rational curve arrangement in the projective plane (i.e. any irreducible component of $`V`$ is either a line or a smooth conic), the cohomology space $`H^m(๐ฐ)`$ is pure of type $`(m,m)`$ for all $`m0`$ (easy exercise for the reader).
(c) $`H^m(๐ฐ)`$ is always pure of type $`(m,m)`$ for all $`m1`$. This follows from the fact that $`g=(g_1,\mathrm{},g_s):๐ฐ๐^s`$ induces an isomorphism at the $`H^m`$-level all $`m1`$. Here we look at $`๐ฐ`$ as a subset of $`^{n+1}`$ and we set $`g_j(x_1,\mathrm{},x_{n+1})=f_j(1,x_1,\mathrm{},x_{n+1})`$.
For $`\lambda =(\lambda _1,\mathrm{},\lambda _s)๐^s`$, let $`_\lambda `$ be the corresponding local system on $`๐ฐ=ZD`$. Let $`\alpha _j`$ be such that $`\mathrm{exp}(2\pi i\alpha _j)=\lambda _j`$ for $`j=1,\mathrm{},s.`$ Then $`_\lambda `$ is the local system of horizontal sections of the connection
$$_\alpha :๐ช_๐ฐ\mathrm{\Omega }_๐ฐ^1$$
given by $`_\alpha (u)=du+u\omega _\alpha `$ where
$$\omega _\alpha =\underset{j=1,s}{}\alpha _j\frac{dg_j}{g_j}.$$
Alternatively, if we look at $`๐ฐ`$ as a subset of $`^{n+1}`$, then we can use the formula
$$\omega _\alpha =\underset{j=0,s}{}\alpha _j\frac{df_j}{f_j}$$
where we set $`\alpha _0=_{j=1,s}d_j\alpha _j`$. Recall that $`f_0=x_0`$.
Using the fact that $`๐ฐ`$ is affine and our connection is regular, it follows that
(4.2)
$$H^m(๐ฐ,_\lambda )=H^m(H^0(๐ฐ,\mathrm{\Omega }_๐ฐ^{}),_\alpha )$$
just as in , (Thm. 3.4.18) or, for complete proofs, . However, this result is not so useful to perform explicit computations since the groups $`H^0(๐ฐ,\mathrm{\Omega }_๐ฐ^{})`$ are too large.
There is a second approach to computing $`H^m(๐ฐ,_\lambda )`$, this time using logarithmic connections. It has the advantage of reducing the size of the spaces $`H^0(๐ฐ,\mathrm{\Omega }_๐ฐ^{})`$, but one has to be more careful about the residues $`\alpha _j`$. More precisely, the pull-back of the connection $`_\alpha `$ under the embedded resolution $`\pi `$ is a logarithmic connection $`\stackrel{~}{}_\alpha `$ on $`Z`$ with poles along $`D`$. Let $`\rho _i`$ be the residue of the connection $`\stackrel{~}{}_\alpha `$ along the irreducible component $`D_i`$ of $`D`$. When $`D_i`$ is the proper transform of some component $`V_j`$ of $`V`$ one has $`\rho _i=\alpha _j`$.
###### Definition 4.3.
A choice of residues $`\alpha =(\alpha _0,\alpha _1,\mathrm{},\alpha _s)`$ for $`_\lambda `$ as above is an admissible choice of residues for $`_\lambda `$ if $`\rho _i_{>0}`$ for all irreducible components $`D_i`$ of $`D`$. A rank one local system $`_\lambda `$ is admissible if there is some admissible choice of residues for it.
###### Remark 4.4.
It is easy to see using Hironakaโs embedded resolution of singularities by blowing-up smooth subvarieties, that for any $`i`$ there is a relation
$$\rho _i=\underset{j=1,s}{}n_{ij}\alpha _j$$
with $`n_{ij}`$ (see for similar formulas and note that negative coefficients occur due to the presence of the hyperplane at infinity). The condition $`\rho _i_{>0}`$ is clearly satisfied if all $`\alpha _j`$ are sufficiently small. In other words, there is a neighborhood $`U(\mathrm{๐})`$ of the trivial local system $`\mathrm{๐}๐^s`$ formed entirely by admissible local systems.
If we move away from the trivial local system , it is not clear whether all the local systems are admissible. The answer to this question is negative for some hyperplane arrangements, see , Example 4.4, , Example 3.4 and . On the other hand, for not very complicated arrangements, see Examples 4.8 and 4.10 below, the answer is positive.
For an admissible choice of residues one has an $`E_1`$-spectral sequence
(4.3)
$$E_1^{p,q}=H^q(Z,\mathrm{\Omega }_Z^p(logD))H^{p+q}(๐ฐ,_\lambda )$$
whose differential $`d_1`$ is induced by $`\stackrel{~}{}_\alpha `$, see , (Thm. 3.4.11 (i)). The above discussion proves the following.
###### Proposition 4.5.
Assume that $`\alpha =(\alpha _0,\alpha _1,\mathrm{},\alpha _s)`$ is an admissible choice of residues for $`_\lambda `$ and that the cohomology groups $`H^m(๐ฐ)`$ are pure of type $`(m,m)`$ for all $`mk`$. Then
$$H^m(๐ฐ,_\lambda )=H^m(H^{}(๐ฐ),\omega _\alpha )$$
for all $`mk`$ and $`H^{k+1}(H^{}(๐ฐ),\omega _\alpha )`$ is a subspace in $`H^{k+1}(๐ฐ,_\lambda )`$.
When $`๐ฐ`$ is a hyperplane arrangement complement, this is exactly the argument used in and . Proposition 4.5, Remark 4.4 and Example 4.2 yield the following.
###### Corollary 4.6.
If $`๐ฐ`$ is any affine hypersurface arrangement complement, then there is a neighborhood $`U(\mathrm{๐})`$ of the trivial local system $`\mathrm{๐}๐^s`$ such that
$$H^1(๐ฐ,_\lambda )=H^1(H^{}(๐ฐ),\omega _\alpha )$$
for any local system $`_\lambda U(\mathrm{๐})`$, $`\alpha `$ being an arbitrary choice of admissible residues for $`_\lambda `$.
###### Corollary 4.7.
If $`๐ฐ=M(๐)`$ is a hyperplane arrangement complement, then there is a neighborhood $`U(\mathrm{๐})`$ of the trivial local system $`\mathrm{๐}๐^s`$ such that
$$H^m(๐ฐ,_\lambda )=H^m(H^{}(๐ฐ),\omega _\alpha )$$
for any $`m`$, and any local system $`_\lambda U(\mathrm{๐})`$, $`\alpha `$ being an arbitrary choice of admissible residues for $`_\lambda `$.
###### Example 4.8.
In the projective plane $`^2`$ consider the hypersurface $`V`$ having as irreducible components $`V_1:x=0`$, $`V_2:y=0`$, $`V_3:x^2yz=0`$. Let $`H=V_0`$ be the line at infinity given by $`z=0`$ and note that $`H`$ is not transverse in a stratified sense to $`V`$. Consider the connection $`_\lambda `$ whose residues are $`\alpha =(\alpha _0,\alpha _1,\mathrm{},\alpha _3)`$ with
$$\alpha _0=\alpha _1\alpha _22\alpha _3.$$
Let $`A=V_1V_2V_3=(0:0:1)`$ and $`B=V_1V_0V_3=(0:1:0)`$. To construct the embedded resolution of $`VH`$ we first blow-up the points $`A`$ and $`B`$, creating thus two exceptional divisors, $`D_A`$ and respectively $`D_B`$. The corresponding residues along $`D_A`$ and $`D_B`$ are easily computable and we get $`\alpha _A=\alpha _1+\alpha _2+\alpha _3`$ and respectively $`\alpha _B=\alpha _1+\alpha _0+\alpha _3=\alpha _2\alpha _3.`$ Let $`P=D_AV_2^{}V_3^{}`$ and $`Q=D_BV_0^{}V_3^{}`$, where denotes the proper transform of a divisor. To get the embedded resolution of $`VH`$ we have just to blow-up the points $`P`$ and $`Q`$, creating thus two new exceptional divisors, $`D_P`$ and respectively $`D_Q`$. The corresponding residues are $`\alpha _P=\alpha _Q=\alpha _1+2\alpha _2+2\alpha _3`$. Therefore the choice of residues $`\alpha =(\alpha _0,\alpha _1,\mathrm{},\alpha _3)`$ is admissible if and only if none of the residues
$$\alpha _1,\alpha _2,\alpha _3,\alpha _1\alpha _22\alpha _3,\alpha _1+\alpha _2+\alpha _3,\alpha _2\alpha _3,\alpha _1+2\alpha _2+2\alpha _3,(\alpha _1+2\alpha _2+2\alpha _3)$$
is a strictly positive integer.
###### Lemma 4.9.
In the situation of Example 4.8, any rank one local system is admissible.
###### Proof.
It is clearly enough to consider the case of real residues $`\alpha _j`$. Otherwise, we just look at the corresponding real parts.
We divide the possibilities in the following two cases.
Case 1. ($`\alpha _1+2\alpha _2+2\alpha _3`$).
Suppose first that in addition $`\alpha _1+\alpha _2+\alpha _3`$. Then the choice with $`\alpha _j[0,1)`$ for $`j=1,2,3`$ is admissible.
Now suppose that $`\alpha _1+\alpha _2+\alpha _3`$. It follows that $`\alpha _2+\alpha _3`$. Then the choice with $`\alpha _j[0,1)`$ for $`j=2,3`$ and $`\alpha _1<0`$ such that $`\alpha _1+\alpha _2+\alpha _3=0`$ is admissible.
Case 2. ($`\alpha _1+2\alpha _2+2\alpha _3`$).
Then we have to choose $`\alpha _1=2\alpha _22\alpha _3`$. The residues in this case are just
$$2(\alpha _2+\alpha _3),(\alpha _2+\alpha _3),\alpha _2,\alpha _3.$$
Hence it is enough to take $`\alpha _j[0,1)`$ for $`j=2,3`$.
Now we continue Example 4.8 by applying Example 4.2 and Proposition 4.5 to get $`H^m(๐ฐ,_\lambda )=H^m(H^{}(๐ฐ),\omega _\alpha )`$ for all $`m`$. In order to perform this computation, we need a precise description of the cohomology algebra $`H^{}(๐ฐ)`$ (with $``$ coefficients) and this can be obtained in this example from the local considerations in , pp. 47-49. The result can be described as follows.
(i) $`H^0(๐ฐ)=`$ and the generator is $`1`$;
(ii) $`H^1(๐ฐ)=^3`$ and a basis is given by $`\eta _1=\frac{dx}{x}`$, $`\eta _2=\frac{dy}{y}`$ and $`\eta _3=\frac{d(x^2y)}{x^2y}`$;
(iii) $`H^2(๐ฐ)=^2`$ and a basis is given by $`\eta _{12}=\eta _1\eta _2`$ and $`\eta _{23}=\eta _2\eta _3`$. The multiplication is given by the relation
$$2\eta _1\eta _3=2\eta _{12}+\eta _{23}.$$
(iv) Since $`๐ฐ`$ is affine, $`H^m(๐ฐ)=0`$ for $`m>2`$.
The computation of $`H^m(H^{}(๐ฐ),\omega _\lambda )`$ falls into 3 cases.
Case 1. $`\alpha _1=\alpha _2=\alpha _3=0`$ and $`_\lambda =`$ is the constant local system. Then of course $`H^m(๐ฐ,_\lambda )=H^m(๐ฐ)`$ for all $`m`$.
Case 2. $`\alpha _1+2\alpha _2+2\alpha _3=0`$. Then a direct computation shows that $`H^0(๐ฐ,_\lambda )=0`$ and $`dimH^1(๐ฐ,_\lambda )=dimH^2(๐ฐ,_\lambda )=1.`$
Case 3. $`\alpha _1+2\alpha _2+2\alpha _30`$. Again a direct computation shows that $`H^0(๐ฐ,_\lambda )=H^1(๐ฐ,_\lambda )=H^2(๐ฐ,_\lambda )=0.`$
The above computations yield the following equalities.
$$V_t^{0,1}(๐ฐ)=V_t^{0,2}(๐ฐ)=\{\lambda ๐^3;\lambda _1\lambda _2^2\lambda _3^2=1\}$$
$$V_t^{1,1}(๐ฐ)=V_t^{2,1}(๐ฐ)=V_t^{1,2}(๐ฐ)=\{\mathrm{๐}\}$$
$`V_t^{m,1}(๐ฐ)=\mathrm{}`$ for $`m>2`$ and $`V_t^{m,2}(๐ฐ)=\mathrm{}`$ for $`m>1`$.
These results are consistent with the general results by Arapura . See also Suciu for a related discussion.
Note that the above 2-dimensional subtorus $`๐=\{\lambda ๐^3;\lambda _1\lambda _2^2\lambda _3^2=1\}`$ is different from the 2-dimensional subtorus predicted by Theorem 3.6 in the case of a divisor $`V`$ transversal to the line at infinity.
A special class of local systems is formed by the equimonodromical local systems $`_\lambda `$ such that $`\lambda _0=\lambda _1=\mathrm{}=\lambda _3`$. Then, for $`\lambda _0^5=1`$, the dimension of the cohomology space $`H^m(๐ฐ,_\lambda )`$ is exactly the multiplicity of the root $`t=\lambda _0`$ in the characteristic polynomial
$$\mathrm{\Delta }^m(t)=det(tIdh^m)$$
where $`F:xyz(x^2yz)=1`$ is the associated Milnor fiber in $`^3`$ and $`h:FF`$ is the monodromy operator, see for instance , (6.4.6). To compute the cohomology of such a equimonodromical local system $`_\lambda `$, one should start by an admissible choice for the residues $`\alpha =(\alpha _0,\alpha _1,\alpha _2,\alpha _3)`$. For instance, the obvious choice $`\alpha =(\frac{4}{5},\frac{1}{5},\frac{1}{5},\frac{1}{5})`$ is not admissible. A good choice here is $`\alpha =(\frac{1}{5},\frac{4}{5},\frac{1}{5},\frac{1}{5})`$. Using this choice, we get the following characteristic polynomials in this situation.
$$\mathrm{\Delta }^0(t)=t1,\mathrm{\Delta }^1(t)=(t1)^2(t^51),\mathrm{\Delta }^2(t)=(t1)(t^51).$$
The following example is similar to the previous one, but it exhibits a curve $`V`$ which is transversal to the line at infinity $`H`$ and it needs a different approach for the computation of the cohomology algebra $`H^{}(๐ฐ)`$. Moreover, in this case the cohomology algebra $`H^{}(๐ฐ)`$ is not spanned by the degree one part $`H^1(๐ฐ)`$.
###### Example 4.10.
In the projective plane $`^2`$ consider the hypersurface $`V`$ having as irreducible components $`V_1:x=0`$, $`V_2:y=0`$, $`V_3:x^2y^2+yz=0`$. Let $`H=V_0`$ be the line at infinity given by $`z=0`$ and note that $`H`$ is transverse in a stratified sense to $`V`$. Consider the connection $`_\lambda `$ whose residues are $`\alpha =(\alpha _0,\alpha _1,\mathrm{},\alpha _3)`$ with
$$\alpha _0=\alpha _1\alpha _22\alpha _3.$$
Let $`A=V_1V_2V_3=(0:0:1)`$. To construct the embedded resolution of $`VH`$ we first blow-up the point $`A`$, creating an exceptional divisor $`D_A`$. The corresponding residue along $`D_A`$ is $`\alpha _A=\alpha _1+\alpha _2+\alpha _3`$. Let $`P=D_AV_2^{}V_3^{}`$, where denotes the proper transform of a divisor. To get the embedded resolution of $`VH`$ we have just to blow-up the point $`P`$, creating a new exceptional divisor $`D_P`$. The corresponding residue is $`\alpha _P=\alpha _1+2\alpha _2+2\alpha _3`$. Therefore the choice of residues $`\alpha =(\alpha _0,\alpha _1,\mathrm{},\alpha _3)`$ is admissible in this case if and only if none of the residues
$$\alpha _1,\alpha _2,\alpha _3,\alpha _1\alpha _22\alpha _3,\alpha _1+\alpha _2+\alpha _3,\alpha _1+2\alpha _2+2\alpha _3$$
is a strictly positive integer. It can be shown, exactly as in Lemma 4.9 above, that in this situation any rank one local system is admissible.
It follows that we can apply Example 4.2 and Proposition 4.5 to get $`H^m(๐ฐ,_\lambda )=H^m(H^{}(๐ฐ),\omega _\lambda )`$ for all $`m`$. To get a precise description of the cohomology algebra $`H^{}(๐ฐ)`$ we can proceed as follows.
(i) $`H^0(๐ฐ)=`$ and the generator is $`1`$;
(ii) $`H^1(๐ฐ)=^3`$ and a basis is given by $`\eta _1=\frac{dx}{x}`$, $`\eta _2=\frac{dy}{y}`$ and $`\eta _3=\frac{d(x^2y^2+y)}{x^2y^2+y}`$;
(iii) To compute $`H^2(๐ฐ)`$ is the first difficulty. This can be done by setting $`๐ฐ^0=^2(V_0V_1V_2)`$, $`V_3^0=V_3(V_0V_1V_2)`$ and considering the Gysin sequence
$$H^1(๐ฐ)H^0(V_3^0)H^2(๐ฐ^0)H^2(๐ฐ)H^1(V_3^0)0.$$
The first morphism, given by the Poincarรฉ-Leray residue $`R`$, is clearly surjective, i.e. $`R(\eta _3)=1`$. Then $`dimH^2(๐ฐ^0)=1`$ and a generator is $`\eta _{12}=\eta _1\eta _2`$. The affine curve $`V_3^0`$ is isomorphic to $`\{1,0,1\}`$ under the parametrization
$$x=\frac{t}{t^21},y=\frac{t^2}{t^21}.$$
Using this parametrization, we can identify $`H^1(V_3^0)`$ to $`^3`$ by sending a rational differential form to its residues at the points $`\{1,0,1\}`$. Some explicit computations involving the last nonzero morphism in the exact sequence above (which is again given by the Poincarรฉ-Leray residue $`R`$) show that $`R(\eta _{13})`$ and $`R(\eta _{23})`$ are linearly independent in $`H^1(V_3^0)=^3`$, where $`\eta _{13}=\eta _1\eta _3`$ and $`\eta _{23}=\eta _2\eta _3`$. It follows that $`\eta _{12}`$, $`\eta _{13}`$ and $`\eta _{23}`$ are linearly independent in $`H^2(๐ฐ)`$, which is 4-dimensional.
It follows that the following cases are possible in this example.
Case 1. $`\alpha _1=\alpha _2=\alpha _3=0`$ and $`_\lambda =`$ is the constant local system. Then of course $`H^m(๐ฐ,_\lambda )=H^m(๐ฐ)`$ for all $`m`$.
Case 2. $`(\alpha _1,\alpha _2,\alpha _3)(0,0,0)`$. Then a direct computation shows that $`H^0(๐ฐ,_\lambda )=H^1(๐ฐ,_\lambda )=0`$ and $`dimH^2(๐ฐ,_\lambda )=2.`$
The above computations yield the following equalities.
$$V_t^{0,2}(๐ฐ)=V_t^{1,2}(๐ฐ)=๐^3$$
(hence here the support has 0-codimension)
$$V_t^{1,1}(๐ฐ)=V_t^{2,1}(๐ฐ)=V_t^{2,2}(๐ฐ)=V_t^{3,2}(๐ฐ)=\{\mathrm{๐}\}$$
$`V_t^{m,1}(๐ฐ)=\mathrm{}`$ for $`m>2`$ and $`V_t^{m,2}(๐ฐ)=\mathrm{}`$ for $`m>3`$. Note that the inclusion in Theorem 3.6 is strict in this case.
Consider as in the above example the associated Milnor fiber $`F:xyz(x^2y^2+y)=1`$ and $`h:FF`$ the monodromy operator. A good choice of residue is again given by $`\alpha =(\frac{1}{5},\frac{4}{5},\frac{1}{5},\frac{1}{5})`$. Using this choice, we get the following characteristic polynomials in this situation.
$$\mathrm{\Delta }^0(t)=t1,\mathrm{\Delta }^1(t)=(t1)^3,\mathrm{\Delta }^2(t)=(t1)^2(t^51)^2.$$
###### Remark 4.11.
In order to apply Theorem 3.2, we have to check the vanishing of some local cohomology groups. When the hypersurface germs occuring in these local complements are quasi-homogeneous, then we can globalise the local situation and compute the corresponding local cohomology groups using the ideas explained in this section. For instance, Example 4.8 covers the case of a plane curve singularity consisting of 3 smooth branches $`(C_1,0)`$, $`(C_2,0)`$ and $`(C_3,0)`$ such that the intersection multiplicities are given by $`(C_1,C_2)=1`$, $`(C_1,C_3)=1`$ and $`(C_2,C_3)=2`$. This follows from the topological classification of the plane curve germs, see , p. 45.
## 5. A More General Setting
In this section we define multi-variable Alexander invariants in a more general setting (see below) and attempt to relate them to the invariants previously defined.
Assume that the hypersurface $`V`$ in $`^{n+1}`$ has $`s`$ irreducible components $`V_i`$ with degrees $`\text{deg}(V_i)=d_i`$ for $`i=1,\mathrm{},s`$. Denote by $`๐ฐ_0`$ the complement $`^{n+1}V`$, and let $`d=g.c.d.(d_1,\mathrm{},d_s)`$. Then
$$H_1(๐ฐ_0)=^{s1}(/d)$$
is generated by the meridians $`\gamma _i`$ about the non-singular part of each component $`V_i`$, for $`i=1,\mathrm{},s`$ (cf. , (4.1.3)). These meridians satisfy a single relation, namely
$$\underset{i=1}{\overset{s}{}}d_i\gamma _i=0.$$
Now fix a hyperplane $`H`$ and set, as before, $`๐ฐ=^{n+1}(VH)`$. Recall that $`H_1(๐ฐ)=^s`$, freely generated by the meridians $`\gamma _i`$, $`i=1,\mathrm{},s`$. Let $`i:๐ฐ๐ฐ_0`$ be the inclusion map, and denote by $`๐ฐ_0^{ab}`$ and $`๐ฐ^{ab}`$ the universal abelian covers of $`๐ฐ_0`$ and $`๐ฐ`$ respectively, and by $`p_0`$ and $`p`$ the corresponding covering projections.
The invariants we are interested in are $`H_{}(๐ฐ_0^{ab};)`$, regarded as modules over the quotient ring $`[H_1(๐ฐ_0)]=[t_1^{\pm 1},\mathrm{},t_s^{\pm 1}]/(t_1^{d_1}\mathrm{}t_s^{d_s}1)`$.
It is a natural question to find the relation between the universal abelian invariants associated to the complement of $`V`$, and those associated to the complement of $`VH`$.
For a topological space $`X`$, let $`(X)`$ denote the set of rank one complex local systems on $`X`$. When $`X=๐ฐ`$, then $`(๐ฐ)`$ is naturally identified to the $`s`$-dimensional complex torus $`๐^s`$. For $`X=๐ฐ_0`$, the set $`(๐ฐ_0)`$ corresponds to the subset in $`๐^s`$ given by
$$\{\lambda =(\lambda _1,\mathrm{},\lambda _s)๐^s|\lambda _1^{d_1}\mathrm{}\lambda _s^{d_s}=1\}.$$
With the notation above, let $`d_j=dd_j^{}`$ and consider the $`(s1)`$-dimensional complex subtorus
$$๐=\{\lambda =(\lambda _1,\mathrm{},\lambda _s)๐^s|\lambda _1^{d_1^{}}\mathrm{}\lambda _s^{d_s^{}}=1\}.$$
For each $`d`$-root of unity $`\beta `$, let $`\lambda (\beta )`$ be one point in the hypersurface in $`๐^s`$ given by the equation
$$\lambda _1^{d_1^{}}\mathrm{}\lambda _s^{d_s^{}}=\beta .$$
Then $`(๐ฐ_0)`$ is precisely the disjoint union of translated tori given by
$$(๐ฐ_0)=\underset{\beta }{}\lambda (\beta )๐.$$
A different way of looking at a local system $``$ in $`(๐ฐ_0)`$ is by considering it as a local system in $`(๐ฐ)`$ (given by the obvious restriction $`|๐ฐ`$) such that the action of the elementary loop about the hyperplane $`H`$ is trivial. This view-point yields the following exact sequence
$$\mathrm{}H^k(๐ฐ_0,)H^k(๐ฐ,)H^{k1}(๐ฐ_0H,)H^{k+1}(๐ฐ_0,)\mathrm{}$$
for details on this see , pp. 221-222. The following obvious consequence should be compared to , , Proposition 1.3. The higher dimensional case, but with a generic hyperplane at infinity $`H`$, was considered in , Lemmas 1.5, 1.11 and 1.13.
###### Corollary 5.1.
Assume that $`V`$ is a plane curve arrangement, i.e. $`n=1`$. Then, for any rank one local system $`=_\lambda `$ on $`๐ฐ_0`$ and any choice of the line at infinity $`H`$, one has
$$dimH^1(๐ฐ,)=dimH^1(๐ฐ_0,)+ฯต.$$
Here $`ฯต\{0,1\}`$ and $`ฯต=0`$ if and only if there is a point $`pVH`$ such that
$$\underset{j=1,s}{}\lambda _j^{k_j}1$$
where $`k_j=mult_p(V_j,H)`$ is the intersection multiplicity of the component $`V_j`$ and the line $`H`$ at the point $`p`$.
One case which is already well-explored is the following.
###### Example 5.2.
Assume that $`n>1`$, $`s=1`$ and that $`V=V_1`$ is a hypersurface of degree $`d`$ having only isolated singularities. Then $`\pi _1(๐ฐ_0)=/d`$ and hence a local system $`=_\beta `$ corresponds to a choice of a $`d`$-root of unity $`\beta `$. For $`\beta =1`$ we get $`H^0(๐ฐ_0;)=`$ and $`H^j(๐ฐ_0;)=0`$ for $`0<j<n`$. When $`V`$ is a $``$-manifold, one also has $`H^n(๐ฐ_0;)=0`$. The computation of $`H^n(๐ฐ_0;)H_0^{n+1}(V)`$ is quite difficult in general, as it may depend on the position of singularities, see , . Here $`H_0^{}(V)`$ denotes the primitive cohomology of $`V`$, i.e. the cokernel of the natural monomorphism $`H^{}(^{n+1})H^{}(V)`$ induced by the inclusion of $`V`$ into $`^{n+1}`$.
For $`\beta 1`$, one can use the isomorphism $`H^m(๐ฐ_0,)=H^m(F,)_\beta `$, the $`\beta `$-eigenspace of the monodromy acting on the Milnor fiber $`F`$ associated to $`V`$. In particular, $`H^m(๐ฐ_0,)=0`$ for $`m<n`$. It is possible to construct examples such that for $`m\{n,n+1\}`$ one has
$$dimH^m(๐ฐ_0,)>dimH^m(๐ฐ_0,).$$
Indeed, consider the polynomials in , p. 148, which have a monodromy operator without the eigenvalue 1 on all the reduced cohomology groups $`\stackrel{~}{H}^m(F,)`$ (equivalently, $`V`$ has the same rational cohomology as $`^n`$). It is not possible that $`\stackrel{~}{H}^m(F,)=0`$ for all $`m`$, by AโCampoโs result on the Lefschetz number of the monodromy, see , p. 174. Hence there is some integer $`m`$ and some $`d`$-root of unity $`\beta 1`$ such that $`dimH^m(๐ฐ_0,_\beta )>0=dimH^m(๐ฐ_0,)`$. Using the Euler characteristic equality $`\chi (๐ฐ_0,)=\chi (๐ฐ_0,_\beta )`$, it follows that the inequality should hold for the two possible values of $`m`$. By the minimality property of hyperplane arrangement complements, it is known that the above inequality is impossible for such complements, .
The discussion in the previous section relating local systems to connections can be extended to this setting in an obvious way. For instance, we should use now the 1-form
$$\omega _\alpha =\underset{j=1,s}{}\alpha _j\frac{df_j}{f_j}$$
where the residues $`\alpha `$ satisfy the condition $`_{j=1,s}d_j\alpha _j=0`$, which is a necessary condition in order to have a 1-form on $`๐ฐ_0`$.
### 5.1. Some 2-component Arrangements
We consider now in detail the case of hypersurface arrangements $`V`$ with $`s=2`$ irreducible components. We assume moreover that:
(i) $`n>1`$ and each $`V_i`$ has at most isolated singularities and is a $``$-manifold;
(ii) $`V^{}=V_1V_2`$ has at most isolated singularities; this condition is automatically fulfilled when $`d_1<d_2`$ and $`V_2`$ is smooth, see .
Let $`๐ฐ_i=^{n+1}V_i`$. Then the Mayer-Vietoris sequence of the covering $`๐ฐ^{}=๐ฐ_1๐ฐ_2`$ reads like
(5.1)
$$\mathrm{}H^{k1}(๐ฐ_0)H^k(๐ฐ^{})H^k(๐ฐ_1)H^k(๐ฐ_2)H^k(๐ฐ_0)\mathrm{}$$
Here and in the sequel the constant coefficients $``$ are used unless stated otherwise. Using Example 5.2 to handle the cohomology groups $`H^{}(๐ฐ_i)`$ for $`i=1,2`$ and the Alexander duality isomorphism (which is compatible with the MHS due to the Tate twist $`(n1)`$)
(5.2)
$$H^k(๐ฐ^{})=H^{2n+2k}(^{n+1},V^{})^{}(n1)=H_0^{2n+1k}(V^{})^{}(n1)$$
we get the following result.
###### Proposition 5.3.
With the above notation and assumptions, the following hold.
(i) $`H^0(๐ฐ_0)=`$ is pure of type (0,0) and $`H^1(๐ฐ_0)=`$ is pure of type (1,1) and is spanned by the 1-form
$$\omega _1=d_2\frac{df_1}{f_1}d_1\frac{df_2}{f_2}.$$
(ii) $`H^k(๐ฐ_0)=0`$ for $`1<k<n`$.
(iii) $`H^n(๐ฐ_0)`$ is pure of weight $`n+2`$ and $`b_n(๐ฐ_0)dimH_0^n(V^{})`$. Moreover $`H^n(๐ฐ_0)=0`$ if $`d_1<d_2`$ and $`V_2`$ is smooth.
(iv) $`H^{n+1}(๐ฐ_0)`$ has weights $`n+2`$ and one has an isomorphism of MHS
$$H^{n+1}(๐ฐ_0)/W_{n+2}H^{n+1}(๐ฐ_0)=H_0^{n1}(V^{})^{}(n1).$$
###### Proof.
The vanishing of $`H^n(๐ฐ_0)`$ in the third claim follows from an unexpected source. Indeed, the Gysin sequence of the smooth divisor $`X_2=V_2๐ฐ_1`$ in $`๐ฐ_1`$ gives a monomorphism $`H^n(๐ฐ_0)H^{n1}(X_2)`$. But this latter group $`H^{n1}(X_2)`$ is trivial by some general connectivity results recently obtained by the first author, see . The examples given in show that the case $`d_1=d_2`$ is much more complicated, in particular the group $`H^{n1}(X_2)`$ can be non-zero. The example below shows that the assumption $`V_2`$ smooth cannot be relaxed to $`V_2`$ with isolated singularities and a $``$-manifold. The key point there is that the singularities of $`V_2`$ are situated on $`V_1`$, a situation not covered by the results in .
The only other claims that are not obvious are those on the MHS. They follow from the fact that $`H_0^n(V^{})`$ has a pure HS of weight $`n`$ (the singularities of $`V^{}`$ being isolated) and the following consequence of the Alexander duality (5.2)
(5.3)
$$h^{p,q}(H^k(๐ฐ^{}))=h^{n+1p,n+1q}(H_0^{2n+1k}(V^{})).$$
For a rank one local system $`(๐ฐ_0)`$, we can choose the corresponding form $`\omega _\alpha `$ to be a multiple $`a(\alpha )\omega _1`$ of the 1-form $`\omega _1`$ introduced above. Then Propositions 4.5 and 5.3 yield the following.
###### Corollary 5.4.
For a non-trivial rank one local system $`(๐ฐ_0)`$ for which an admissible choice of residues $`\alpha =(d_2a(\alpha ),d_1a(\alpha ))`$ exists, the following hold.
(i) $`H^k(๐ฐ_0,)=0`$ for $`k<n`$.
(ii) If $`H_0^n(V^{})=0`$ or if $`d_1<d_2`$ and $`V_2`$ is smooth, then $`H^n(๐ฐ_0,)=0`$.
Note that the first claim above holds by Corollary 3.5, since $`V^{}`$ has only INNC singularities.
The vanishing of $`H_0^n(V^{})`$ holds when $`V^{}`$ is a $``$-homology manifold, but also in many other cases, see for instance the discussion in , pp. 207-216. There one considers only the case when $`V_1`$ is a hyperplane. Indeed, any hypersurface $`W`$ having only isolated singularities in $`^n`$ can be obtained as the intersection of a smooth hypersurface $`V_2`$ in $`^{n+1}`$ with the hyperplane $`H=^n`$, see , p. 206. However, this situation is usually uninteresting according to the second claim of the above corollary.
We conclude with an example where $`V_1`$ is a hyperplane, and $`V_2`$ is singular, so that it may have been considered already in the previous section (in such a case $`๐ฐ_0`$ from this section is exactly $`๐ฐ`$ from the previous section, but for the hypersurface $`V=V_2`$ !).
###### Example 5.5.
In $`^3`$ (with homogeneous coordinates $`(x:y:z:t)`$) consider the hyperplane $`H=V_1:t=0`$ and the surface $`V_2:xyzt^3=0`$. Then $`V_2`$ has exactly 3 singularities of type $`A_2`$, hence it is a $``$-manifold. Moreover, $`H`$ is transversal to $`V_2`$, except at the 3 singular points of $`V_2`$.
To compute the cohomology of the complement $`๐ฐ_0`$, we use the Gysin exact sequence of the smooth divisor $`D=V_2V_1`$ in the affine space $`^3V_1`$ (with coordinates $`(x,y,z)`$) and get
$$H^k(๐ฐ_0)=H^{k1}(D)(1)$$
for $`k=2,3`$, where $`(1)`$ denotes the Tate twist. Now $`D`$ is given by the equation $`xyz=1`$, hence it is a 2-dimensional torus. It follows that
(i) $`H^2(๐ฐ_0)=^2`$ is pure of type (2,2);
(ii) $`H^3(๐ฐ_0)=`$ is pure of type (3,3). Moreover, as explained in the forth section, the 1-form $`\omega _1`$ is a multiple of $`\frac{dg}{g}`$ with $`g=xyz1`$.
Let $`g_0=g+1=xyz`$, and note that $`F_{g_0}=D`$ is the Milnor fiber of the homogeneous polynomial $`g_0`$. Since $`๐ฐ_0=^3F_{g_0}`$, we may use the description of the cohomology groups of $`๐ฐ_0`$ using Remark (2.11) in , p.192. By taking
$$\eta _i=A_ixdydzB_idxdz+C_idxdy$$
in the formula (2.12) loc.cit. with $`(A_1,B_1,C_1)=(x,y,0)`$ and $`(A_2,B_2,C_2)=(0,y,z)`$ we get a basis of $`H^2(๐ฐ_0)`$. A direct computation then shows that $`\omega _1\eta _i=0`$ in $`H^3(๐ฐ_0)=H^2(D)`$. To see this, note that $`d\eta _i=dg\eta _i=0`$ and hence the Poincarรฉ-Leray residue of the form
$$\omega _1\eta _i=\frac{dg}{g}\eta _i$$
is the form $`\eta _i`$. Since $`d\eta _i=0`$ on $`^3`$, it follows the $`\eta _i=d\eta _i^{}`$, for some 1-forms $`\eta _i^{}`$ on $`^3`$. Hence the cohomology class of $`\eta _i`$ in $`H^2(D)`$ is trivial.
It follows that, for a non-trivial rank one local system $`(๐ฐ_0)`$ for which an admissible choice of residues $`\alpha `$ exists, one has $`H^{}(๐ฐ_0,)=H^{}(H^{}(๐ฐ_0),\omega _\alpha ).`$ Therefore we get the following equalities.
$$dimH^0(๐ฐ_0,)=dimH^1(๐ฐ_0,)=0,dimH^2(๐ฐ_0,)=2,\text{and}dimH^3(๐ฐ_0,)=1.$$
In particular, $`\text{Supp}A^2(๐ฐ_0)`$ coincides to the character torus $`๐^1`$. Indeed, $`\text{Supp}A^2(๐ฐ_0)`$ is a Zariski closed subset, with a non-empty interior by Remark 4.4, in the irreducible algebraic variety $`๐^1`$. The reader should compare this fact to Theorem 3.6 above.
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# Populations of AGB stars and LPVs in the Galaxy and Local Group
## 1. Introduction
All main-sequence stars born with masses below $`<`$ 8 $`M_{}`$ have or will go through the evolutionary phase called Asymptotic Giant Branch (AGB). The lower limit in initial mass is set by the age of the Galactic Disc, the upper limit is set by the mass where carbon can be ignited in the stellar core. The AGB is the final phase where intermediate-mass stars have nuclear burning in the form of alternate Hydrogen and Helium shell burning, before they cross the Herzsprung-Russell diagram to become Planetary Nebulae and then White Dwarfs.
A summary of the interior structure and stellar evolution up to and on the AGB can be found in the recent textbook โAsymptotic Giant Branch Starsโ (Habing & Oloffson 2004).
AGB stars are luminous (\[$``$0.1 โ a few\] 10<sup>4</sup> $`L_{}`$) and cool, with effective temperatures in the range 3850 to $``$2500 K (for M0 to M10 giants, e.g. Fluks et al. 1994). From this it follows that AGB stars are big (up to a few hundred $`R_{}`$), and combining this with the classical pulsation equation<sup>1</sup><sup>1</sup>1With the pulsation constant $`Q`$ being of order 0.038. $`P=QR^{1.5}M^{0.5}`$ it follows that any radial pulsations that would occur would have periods between tens and hundreds of days.
Equally important, and typical for the AGB, are the chemical peculiarities that occur during this evolutionary phase (see Chapter 2 in the abovementioned book). Depending in a complex way on initial mass, metallicity, mass-loss, mixing and burning in the envelope \[hot bottom-burning\], an AGB star may go through several third dredge-up events whereby mainly carbon, nitrogen, oxygen and $`s`$-process elements are mixed ultimately into the stellar photosphere. Depending on the C/O-ratio different molecules form in the cool atmospheres (VO, TiO, C<sub>2</sub>, CN) and a star can be classified as M-star (C/O $`<`$ 0.95), S-star (0.95 $`<`$ C/O $`<`$ 1.0) or C-star (C/O $``$ 1.0). Intermediate classes, MS, SC, also exist. These type of AGB C-stars are also identified with the N-type carbon stars, and they are different from the R-type carbon stars which are likely not on the AGB, and possibly related to a merger of stars on the RGB (see McClure 1997).
The low effective temperatures already make AGB stars redder than all their Main-Sequence progenitors. In addition, the formation of different molecules depending on chemical type makes that the infrared colours of M- and C-stars are different, as will be discussed later.
Furthermore for spectral types later than $``$M4-M5 (e.g. Glass & Schultheis 2002) the region close to the star has the right combination of temperature and density for dust grains to form. Dust absorbs efficiently in the optical and radiates in the infra-red. This implies that AGB stars surrounded by dust shells are even redder.
So far, the binary chanel has not been discussed. Main-sequence, sub-giants and giants that have been polluted when a present-day WD was on the AGB and polluted his companion with material enriched in carbon and $`s`$-process elements. Examples of these classes are the carbon-dwarfs (e.g. Steinhardt & Sasselov 2005), the CH-stars (e.g. Bartkevicius 1996, McClure & Woodsworth 1990), the Barium stars (e.g. Jorissen et al. 1998) and the โextrinsicโ S-stars (e.g. Van Eck & Jorissen, 1999). Recently, it has become clear that about 20% of Carbon-enhanced Extreme Metal-Poor (CEMP) stars are binaries and show s-process enhancements (Lucatello et al. 2005). These classes will not be discussed further here.
Based on the main properties of AGB stars mentioned above, the following recent results will be discussed: 2MASS/DENIS and MSX in connection with the infrared colours, the result of the micro-lensing surveys in connection with variability, and the results of dedicated narrow-band surveys in connection with the chemical characteristics of AGB stars.
## 2. AGB stars in the 2MASS/DENIS and MSX surveys
### 2.1. 2 Micron All Sky Survey โ DEep Near-Infrared Survey
In May 2003 the DENIS team released their second data release containing $`IJK`$ photometry for 195 million point sources in the southern sky down to $`K`$ = 14.0 (available at VIZIER at http://vizier.u-strasbg.fr/viz-bin/VizieR?-source=B/denis)
In March 2003 the final all-sky 2MASS $`JHK`$-database was released containing 470 million point sources down to a typical limiting magnitude of $`K`$ = 14.3 (Cutri et al. 2003). This implies that a typical carbon star even with no mass loss and with a luminosity of 7000 L and effective temperature of 2650 K can be viewed to 150 kpc, and mass-losing AGB stars to even larger distances.
This suggests that AGB stars in many Local Group galaxies have been detected by 2MASS. Groenewegen (2005) presents a list of almost 100 candidates in 16 Local Group (LG) galaxies (excluding the Magellanic Clouds, M31, M33), selected on $`(JK)`$-colour and $`K`$-magnitude. We hope to confirm the 34 candidates in Fornax dSph by NIR spectroscopy in the fall of 2005 using the VLT.
In the near future much deeper NIR survey data of LG galaxies is within reach of wide-field imaging facilities as UKIDSS (http://www.ukidss.org) and VISTA (http://www.vista.ac.uk).
Up to now, the 2MASS and DENIS data have extensively been used to study AGB stars in the Magellanic Clouds:
* Weinberg & Nikolaev (2001) who select stars with $`9.5<K<11.5`$ and $`1.4<(JK)<2`$ (assumed to be C-rich LPVs) to study the 3D structure.
* Cioni & Habing (2003) present results on the C/M ratio across the face of the SMC and LMC (results on NGC 6822 are presented in Cioni & Habing, 2005) using DENIS data. Carbon and O-rich AGB stars are selected from the $`(IJ)(JK)`$ colour-colour (CC)-diagram. The C/M ratio is then determined in 100 $`\times `$ 100 cells of 0.04 square degrees. In the SMC there is no clear trend with position while in the LMC the ratio appears to decrease radially. Overall the distribution is clumpy and corresponds to a spread in metallicity of 0.75 dex.
* Marigo et al. (2003; see Figure 3) fit the colour-magnitude diagram (CMD) of 2MASS objects within 4 degrees of the bar optical centre using a population synthesis code (Girardi et al. 2005), with special emphasis on the tail of carbon stars. She demonstrates the importance of including the change in molecular opacities in evolutionary calculations as stars evolve from O-rich to C-rich (Marigo 2002).
Regarding our own Milky-Way it is interesting to mention the Faint High-Latitude Carbon (FHLC) stars<sup>2</sup><sup>2</sup>2The halo usually defined as $`b>30`$ degrees in this context.. Originally one was interested in identifying carbon stars at high galactic latitude as tracers of the halo. However, quite a few turned out to have a measurable parallax or proper motion, indicating that these were not giants, but rather carbon dwarfs. Papers that use 2MASS data in connection with FHLC stars are:
* Liebert et al. (2000) looked for very red objects in the halo in the 2MASS second incremental data release and confirmed, taking optical spectra, four new halo carbon stars with $`(JK)`$ between 2.3 and 4.7. These stars are similar to two examples found earlier using IRAS (e.g. Groenewegen et al., 1997, and references therein).
* Mauron et al. (2004) did a similar study using the 2MASS but their initial selection was based on $`(HK)>0.4`$ and $`(JH)>0.95`$, hence significantly bluer than the stars found by Liebert et al., but not as blue as some carbon dwarfs. Their initial sample of about 1200 objects was trimmed by eliminating known objects using SIMBAD and stars with $`(BR)<1.5`$ from the USNO. They ended up with a list of $``$200 best candidates and optical spectra were taken of 97 and 30 were confirmed as carbon stars.
* Downes et al. (2004) follow up on the work by Margon et al. (2002) by selecting FHLC stars using 5-colour Sloan Digital Sky Survey (SDSS) photometry, and then taking optical spectra. Correlation with the 2MASS is done to provide the IR magnitudes. Margon et al. find 39, and Downes et al. 251 FHLC stars. At least half of these are dwarf carbon stars. The density of FHLC stars is about 0.06 per square degree down to $`r^{}=21`$.
### 2.2. Midcourse Space Experiment
The Midcourse Space Experiment (MSX; see Price et al. 2001) satellite mission included the spirit iii instrument that obtained mid-IR data in 6 bands designated A,B1,B2,C,D,E with central wavelengths between 8 and 21 $`\mu `$m. Band A was the most sensitive one with a limiting flux-density of order 0.1 Jy, and most of the 320 000 objects are only detected in this band.
MSX observed the Galactic plane ($`b<5^{}`$) and the LMC. Lumsden et al. (2002) and Ortiz et al. (2005) investigate the colours of AGB stars. Both papers reach similar conclusions in that the combination of MSX magnitudes with near-infrared magnitudes ($`J`$ or $`K`$) is a powerful diagnostic in discriminating between different types of reddened sources.
The MSX results on the LMC are discussed in Egan et al. (2001). 1806 sources were detected in the $`A`$ band of which 1664 had positional matches with 2MASS. This allowed them to discriminate different classes of objects in a ($`JK`$)-($`KA`$) CC-diagram.
## 3. Long Period Variables in the micro-lensing surveys
Before focusing on the results of the micro-lensing surveys per se on LPV research, it is worthwhile to mention a few variability surveys that also are in the process of providing results on variable stars. Contrary to the micro-lensing surveys the surveys mentioned below will be most useful for studies of variable stars in the solar neighbourhood and a welcome complement to the classical General Catalog of Variable Stars (GCVS, Kholopov et al., 1998) which lists 7200 Mira variables. These new surveys are:
* ASAS (All-Sky Automated Survey), see Pojmanski (2002).
This survey monitored about 15 million stars in the magnitude range $`V=514`$, south of declination +28 degrees, for 3 years every 1-3 days.
Of the 38 000 variables in the ASAS-3 release 2551 are classified as Mira variables on the web-site.
* NSVS (Northern Sky Variability Survey), see Woลบniak et al. (2004a).
This survey monitored about 14 million objects in the northern sky (and some fields with less quality down to $`38`$ in declination), in the optical magnitude range 8 to 15.5, over a 1 year baseline with 100-500 measurements per object.
Woลบniak et al. (2004b) use this dataset to create a catalog of 8600 variable AGB stars.
* TAROT (Klotz, this proceedings)
### 3.1. The Magellanic Clouds
Lets us now turn to the micro-lensing surveys. After the initial results presented by Wood et al. (1999) and Wood (2000) on a single field of MACHO data in the LMC, there has been a flood of papers on this topic. Table 1 summarises the current state of affaires. The main differences are in the location of the sky (SMC, LMC, Galactic Bulge), the variability survey used (MACHO, OGLE, MOA, EROS), use of infra-red data (publicly available 2MASS or DENIS data, or own observations), and the selection of the variables (unbiased, selection of Miras, starting from spectroscopically known AGB stars).
The main results are very similar between these papers, and one of the main results is shown in Figure 2. It shows different sequences labelled โABCDโ. Sequences $`A`$ and $`B`$ are further split-up relative to the tip of the RGB. Sequence $`B`$ is now often considered to be actually 2 sequences (e.g. Soszynski et al., 2004, Fraser et al., 2005).
Selecting on amplitude selects stars in different boxes. Selecting the largest amplitudes results in a nice Period-Luminosity relation in Box $`C`$ that can be identified with the classical Mira variables. They are almost certainly pulsating in the fundamental mode (Wood 1999, 2000, Fraser et al. 2005). Selecting progressively smaller amplitude cuts populates more and more Box $`B`$ and then $`A`$. These stars are thought to be overtone pulsators (1st, 2nd, 3rd) although theoretical $`PL`$-relation do not fit the observed sequences exactly (e.g Fraser et al. 2005).
The big surprise to come out of these observations was sequence $`D`$. Stars with a period in this box have multiple periods with the shorter period in one of the other boxes. Examples of the light curves are shown in Figure 3. The origin of these, so-called, Long Secondary Periods (LSPs) is still a mystery. Several explanations have been proposed, $`g^+`$ non-radial modes, binary companions, rotating prolate spheroids, episodic dust ejection, star spot cycles, but none is satisfactory (Olivier & Wood 2003, Wood et al. 2004).
### 3.2. The Galactic Bulge
Table 1 includes the few works that have used the results of the micro-lensing surveys to study variables in the direction of the Galactic Bulge (GB). They extend previous classical works on Bulge variable stars, like those of Lloyd Evans (1976), Glass & Feast (1982), Whitelock et al. (1991), Glass et al. (1995), Alard et al. (1996) and Glass et al. (2001).
The work by Groenewegen & Blommaert (2005, GB05) appears to be the first paper to use the full OGLE-ii database to study classical Mira variables ($`I`$-band semi-amplitude larger than 0.45 magnitudes), and some details of this study are presented here.
Figure 4 presents the $`K`$-band $`PL`$-relation for the 2691 Mira variables identified in the OGLE database after correlation with the 2MASS and DENIS. The mean $`PL`$-relation is also shown.
With this slope fixed, one can now determine the zero point (ZP) of the $`PL`$-relation for each OGLE field separately. In Figure 5 the ZP is plotted against Galactic longitude for the inner OGLE fields. There is a significant slope indicating the presence of the well-known bar.
Simulations have been carried out using the basic disk & bulge model of Binney et al. (1997), with the main parameter being the angle between the major axis of the bar and the line-of-sight towards the Galactic Centre. It turns out that for angles of 43 and 79 degrees one can account for the observed slope in Figure 5, but only the model with an angle of 43 degrees can also account for the observed number of stars.
The preferred value of $`\varphi =43^{}`$ is in agreement with the values of about $`45^{}`$ by Whitelock (1992), based on 104 IRAS detected Mira variables, and the preferred value of $`46`$ degrees by Sevenster et al. (1999), based on an analysis of OH/IR stars in the inner Galaxy.
Other values in the literature are usually much larger, between 60 and 80 degrees: Dwek et al. (1995) and Binney et al. (1997), based on COBE-DIRBE data, Stanek et al. (1997), based on bulge red clump stars, Robin et al. (2003) and Picaud & Robin, based on colour-magnitude fitting. Sevenster et al. (1999), however, argues that these values are commonly found when no velocity data is available, the longitude range is too narrow or when low latitudes are excluded. It is also possible that these studies are tracing other populations, which may be differently distributed than the Miras. Whitelock et al. and Sevenster et al. do use populations closely related to the Mira stars and find an angle of the bar close to the one we derive.
Another way to look at this dataset is shown in Figure 6 where the period distribution is shown for OGLE-fields at similar latitudes but different latitudes. To add a field even closer to the GC than surveyed by OGLE the data in Glass et al. (2001, 2002) is considered on a field centered on $`l=0.05{}_{}{}^{},b=0.05^{}`$. They present the results of a $`K`$-band survey of 24 $`\times `$ 24 arcmin<sup>2</sup> for LPVs down to $`K12.0`$.
A Kolmogorov-Smirnov test indicates that the fields at $`1.2^{}`$ and lower are statistically identical, and very different from the inner field, where a tail of longer periods is observed.
To quantify the nature of the Mira Bulge population, synthetic AGB evolutionary models have been calculated. In brief, the synthetic AGB code of Wagenhuber & Groenewegen (1998) was finetuned to reproduce the models of Vassiliadis & Wood (1993) for $`Z`$ = 0.016 and then extended to more initial masses and including mass loss on the RGB. For several initial masses the fundamental mode period distribution was calculated for stars inside the observed instability strip and when the mass loss was below a critical value to simulate the fact that they should be optically visible. The results are shown in Figure 7.
From the comparison of the observed period distribution for fields more than $`1.2^{}`$ away from the galactic centre with the theoretical ones, we deduce that the periods can be explained with a population of stars with Main Sequence masses in the range of 1.5 to 2.0 M. A possible extension to smaller masses is possible, but not necessary to explain the periods below 200 days. To explain the excess periods in the range of 350-600 days observed closer to the centre we need initial masses in the range 2.5 - 3 M. The presence of more massive stars in the inner field at $`b`$ = $`0.05^{}`$ cannot be excluded, as it turns out that for more massive stars the optically visible Mira phase is essentially absent.
The formation history of the Bulge is still a matter of debate. In several works like in Kuijken & Rich (2002) and recently in Zoccali et al. (2003), the bulge is considered to be old ($`>10`$ Gyr) and formed on a relatively short timescale ($`<1`$ Gyr) (e.g. Ferreras et al. 2003). The bulge Miras do not fit in this picture as, according to the analysis in GB05, they are considerably younger.
Their results agree more with the analysis of the infrared ISOGAL survey discussed in van Loon et al. (2003). They conclude that the bulk of the bulge population is old (more than 7 Gyr) but that a fraction of the stars is of intermediate age (1 to several Gyr). The Miras in GB05 study can thus be considered as the intermediate age population seen in their analysis.
If indeed the bulk of the bulge population is old and formed quickly and if the Miras are of intermediate age, then the Miras must be representatives of a population which was added at a later stage and it is unclear how it relates to the overall bulge. An interesting scenario suggested in Kormendy & Kennicutt (2004) is the one in which a secondary bulge or also called pseudo-bulge forms within an old bulge. Such a process would be connected to the presence of a โbarโ which would add โdiskyโ material into the old classical bulge. The Miras are indeed situated in a bar-structure as was discussed earlier.
## 4. AGB stars in the Local Group using narrow-band filter surveys
This technique uses the specific characteristic of late-type M-stars, where strong TiO bands develop, and C-stars, with C<sub>2</sub> and CN molecular bands. First introduced by Palmer & Wing (1982) and then applied by Richer et al. (1984) and Aaronson et al. (1984) the method typically uses two broad-band filters from the set $`V,R,I`$, and two narrow-band filters near 7800 and 8100 ร
, which are centred on a CN-band in carbon stars (and near-continuum in oxygen-rich stars), and a TiO band in oxygen-rich stars (and continuum in C-stars), respectively. In an \[78-81\] versus $`[VI]`$ (or $`[RI]`$) colour-colour plot, carbon stars and late-type oxygen-rich stars clearly separate redwards of $`(VI)`$ 1.6. For an illustration of this, see Cook & Aaronson (1989) or Nowotny & Kerschbaum (2002).
Originally, the applications were limited mainly because of the small format CCDs and the use of relatively small telescopes. However the last about 5 years have seen a remarkable revival through the use of these narrow-band filters using wide-field imagers on bigger telescopes (2-4m class).
At present a large fraction of LG galaxies have been surveyed, at least partially, using these narrow-band filters. For recent reviews see Azzopardi (1999) and Groenewegen (1999, 2002, 2005). The most recent works not listed in these reviews are the surveys by:
* Kerschbaum et al. (2004) who surveyed a 6.5 $`\times `$ 6.5 field covering almost all of And ii and identifying 7 C-stars.
* Demers et al. (2004) who identified 676 carbon stars in a 42 $`\times `$ 28 field centered on IC10.
* Battinelli & Demers (2005a) who surveyed a 450 arcmin<sup>2</sup> field about 40 kpc along the southern major axis of M31 finding only a handfull of C-stars, concluding they reached the edge of the M31 disk, at least as defined by an intermediate age population.
* Harbeck et al. (2005) who found one candidate carbon star in And ix.
* Rowe et al. (2005) who surveyed a 74 $`\times `$ 56 field centered on M33 finding 7936 C-stars.
A quantitative theoretical interpretation or understanding of this data is clearly lacking. At the lowest level of interpretation these data indicate a relation between the ratio of C-to-M stars and (mean) metallicity in the Galaxy, as shown in Figure 8 (data from Battinelli & Demers, 2005b).
Mouhcine & Lanรงon (2003, see Figure 9) present evolutionary population synthesis models, including chemical evolution, with special focus on intermediate age populations. Their models are the first that are able to account qualitatively for the observed trend in Figure 8 adopting โtypicalโ Star Formation Histories (SFHs) for Sa, Sb, Sc and Irr Hubble type galaxies. The AGB phase is included through a semi-analytical treatment of the third dredge-up, with efficiency parameters set to values that have been determined in other studies to fit the LMC carbon star LF and C/M ratio.
With SFH now (becoming) available for most LG galaxies (see Dolphin, this proceedings) one could investigate for individual galaxies whether these SFHs predict the observed number of M- and C-type AGB stars. Repeating this for a large set of galaxies covering a range of metallicities would possible better constrain individual SFHs for intermediate ages, as well lead to a comprehensive picture of dredge-up and mass loss efficiency as a function of metallicity in AGB stars.
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