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warning/0002/physics0002025.html | ar5iv | text | # The oscillations in the lossy medium
## 1 The theoretical part
Let’s begin from the definition: If some physical quantity $`F`$ under specified physical conditions is described periodic or almost-periodic function of time one can say that this physical quantity is in oscillatory process or in oscillations.
As is known, a function $`F(t)`$ is called periodic if $`F(t)=F(t+T)`$.<sup>1</sup><sup>1</sup>1The definition almost-periodic functions will be introduced later. At the oscillatory process the constant $`T=2\pi /\omega `$ is called an oscillation period and the constant $`\omega `$ is called an oscillation frequency (circular or cyclic). Obviously $`T`$ is a time interval by means of that the values of function $`F(t)`$ are repeated.
If the physical quantity is in oscillations described by the harmonic function of time (i.e. function $`\mathrm{sin}(\omega t)`$ or $`\mathrm{cos}(\omega t)`$) the oscillations is called harmonic.
Among all oscillatory processes the special interest is represented those which the man can observe directly without any devices. The most known oscillatory process having so remarkable property is the oscillatory motion.
According to this, the oscillatory motion of a mass point we will call any its motion that the all physical quantities describing motion are periodic (or almost-periodic) functions of time.
The major physical values describing a motion of a mass point are:
* the radius vector of a particle $`\stackrel{}{r}(t)`$, i.e. its coordinates (we shall remind the equation of a form $`\stackrel{}{r}=\stackrel{}{r}(t)`$ is called the motion equation (or law);
* and the vector of the particle acceleration $`\stackrel{}{a}(t)`$.
If we take into account that the vectors of velocity and acceleration are defined uniquely by the radius vector $`\stackrel{}{r}(t)`$ of a mass point, it is possible to formulate the following definition:
| any motion of a mass point |
| --- |
| at which the radius vector of a particle |
| is a periodic (or almost periodic) function of time |
| is called the oscillatory motion |
### 1.1 Free simple harmonic motions
The elementary oscillatory motion of a mass point is the harmonic oscillatory motion. Thus, according to the definition of harmonic oscillations, we shall called by free simple harmonic motion such oscillatory motion, at which the radius vector of a particle is harmonic function of time. It means, the equation (the law) motion of a mass point that is in a free harmonic oscillatory motion, has a form
$$\stackrel{}{r}(t)=\stackrel{}{r}_o\mathrm{sin}(\omega t+\phi _o).$$
(1)
In eq.(1) the constant $`\mathrm{\Phi }=\omega t+\phi _o`$ is called by a phase of the oscillatory motion and its value at $`t=0`$, i.e. $`\phi _o`$, is called by an epoch angle accordingly. The constant $`\stackrel{}{r}_o`$ is called by an amplitude of the oscillatory motion. From the equation (1) is obvious that the amplitude is the maximal value of radius vector the achievable at those point in time when $`\mathrm{sin}(\omega t+\phi _o)=1`$.
Let’s note one important characteristic of the oscillatory motion described by the equation (1). The vector $`\stackrel{}{r}_o`$ is a constant vector, i.e. does not change neither in magnitude nor in the direction. Therefore the vector $`\stackrel{}{r}(t)`$ can change only in magnitude (at the expense of function $`\mathrm{sin}(\mathrm{})`$), but remains parallel to the same line. It means that the harmonic oscillatory motion always has only one degree of freedom. In other words, one coordinate is enough for describing of a harmonic oscillatory motion. For example, coordinates measured along an axis $`OX`$. So the vector equation (1) can always be replaced by one equation in the coordinate form
$$x(t)=x_o\mathrm{sin}(\omega t+\phi _o),$$
(2)
where $`x_o=|\stackrel{}{r}_o|`$ is the module of the vector $`\stackrel{}{r}_o`$.
It is easy to see that the equation (1) is the solution of the differential equation
$$\frac{d^2\stackrel{}{r}}{dt^2}+\omega ^2\stackrel{}{r}=0.$$
(3)
For this reason the differential equation (3) is called by the equation of free simple harmonic motions. So, one can say that
| free harmonic oscillatory motion |
| --- |
| of a mass point is any motion described by |
| the equation of free simple harmonic motions (eq.(3)) |
Classical example of a free harmonic oscillatory motion is the particle motion with the mass $`m`$ due to action of quasi-elastic force (i.e. simulative elastic force) $`\stackrel{}{F}=k\stackrel{}{r}`$, where $`k`$ is stiffness coefficient. To be convinced of it we shall describe for such a mass point the dynamical equation (i.e. Newton’s second law)
$$m\stackrel{}{a}=k\stackrel{}{r}.$$
(4)
Taking into account that the acceleration is a second-order derivative of the particle radius vector, we shall obtain
$$\frac{d^2\stackrel{}{r}}{dt^2}+\frac{k}{m}\stackrel{}{r}=0.$$
(5)
Comparing the obtained equation with the equation of free simple harmonic motions (3), we can see that the motion a mass point due to action of quasi-elastic force is really a free harmonic oscillatory motion. And the oscillation cyclic frequency of a mass point is equal
$$\omega =\sqrt{\frac{k}{m}}.$$
(6)
### 1.2 Damped oscillations
In the previous section we have considered a free harmonic motion and were convinced that due to action of only elastic force the mass point makes just such motion.
Let’s consider now motion a mass point due to action of quasi-elastic forces $`\stackrel{}{F}=k\stackrel{}{r}`$ in medium under the action of resistance forces. Let, for example, the resistance force is proportional to a vector of the particle velocity $`\stackrel{}{F_c}=b\stackrel{}{v}`$, where $`b`$ is the resistance coefficient. Then the dynamical law (Newton’s second law) for such the mass point will have a form
$$m\stackrel{}{a}=\stackrel{}{F}+\stackrel{}{F_c}=k\stackrel{}{r}b\stackrel{}{v}.$$
(7)
Taking into account that the velocity is a first-order derivative and that the acceleration is a second-order derivative of the particle radius vector, we shall obtain
$$\frac{d^2\stackrel{}{r}}{dt^2}+\frac{b}{m}\frac{d\stackrel{}{r}}{dt}+\frac{k}{m}\stackrel{}{r}=0.$$
(8)
It is easy to be convinced that the obtained equation coincides with the equation of free simple harmonic motions only at absence of the resistance forces (i.e. at $`b=0`$ ). The solution of the equation (8) varies from the solution of the equation (3) as well. The eq.(3) is the equation of free simple harmonic motions. So, the common solution of the equation (8) will have a form
$$\stackrel{}{r}(t)=\stackrel{}{r}_oe^{\beta t}\mathrm{sin}(\omega t+\phi _o),$$
(9)
where the following notation for parameters of an oscillatory motion described by the equation (8) are conventional
damping factor $``$ $`\beta ={\displaystyle \frac{b}{2m}}`$
oscillation cyclic frequency of the free harmonic oscillatory motion (i.e. at absence of the resistance forces) $``$ $`\omega _o=\sqrt{{\displaystyle \frac{k}{m}}}`$ (10)
oscillation cyclic frequency of the studied harmonic oscillatory motion $``$ $`\omega =\sqrt{\omega _o\beta ^2}`$
Let’s note that in these notation the equation (8) will look like
$$\frac{d^2\stackrel{}{r}}{dt^2}+2\beta \frac{d\stackrel{}{r}}{dt}+\omega _o\stackrel{}{r}=0.$$
(11)
As well as in the case of free simple harmonic motions the oscillatory motion described by the equation (9) has only one degree of freedom. Hence, if to set the direction of constant vector $`\stackrel{}{r}_o`$ parallelly to axis $`OX`$ of a cartesian frame, the eq.(9) will have a form
$$x(t)=x_oe^{\beta t}\mathrm{sin}(\omega t+\phi _o),$$
(12)
where, as well as in the equation (2), x is the length of a vector $`\stackrel{}{r}_o`$.
In fig.2 the qualitative view of the solution (12) is presented. This figure demonstrate that the studied oscillatory motion represents oscillations with amplitude decreasing in time by exponential law (i.e. described by the function $`e^{\beta t}`$). Just for this reason an oscillatory motion described by the equation (11), named as the damped oscillatory motion. Accordingly, eq.(11) named as the equation of damped oscillations. So
| the damped oscillatory motion |
| --- |
| of a mass point |
| is any motion described by |
| the equation of damped oscillations (i.e. eq.(11)) |
Let’s consider more in detail properties of the damped oscillatory motion. First of all it is obvious that in contrast to the free harmonic oscillatory motion the radius vector of the mass point in damped oscillations (i.e. expression (9) or (12)) is not periodic function of time $`\stackrel{}{r}(t)\stackrel{}{r}(t+T)`$. Thus damped oscillations are not harmonic oscillations.
According to the definition by H. Bohr (Danish mathematician) the function $`f(t)`$ satisfying the requirement
$$|f(t+T)f(t)|<ϵ,$$
(13)
where $`ϵ`$ is some positive number is named an almost-periodic function. Accordingly, $`T`$ is named an almost-period such function. And the mean value of an almost-periodic function is always limitary
$$\underset{T\mathrm{}}{lim}\frac{1}{T}\underset{0}{\overset{T}{}}f(t)𝑑t<\mathrm{}.$$
(14)
It is easy to be convinced that for $`x(t)`$ from expression (12)
$$\underset{T\mathrm{}}{lim}\frac{1}{T}\underset{0}{\overset{T}{}}x(t)𝑑t=0.$$
(15)
Moreover, it always is possible to select such positive number $`ϵ`$ that the absolute value of the difference $`|x(t+T)x(t)|`$ (where $`T=2\pi /\omega `$) will be less than this number. So the requirement (13) will be satisfied.
Hence radius vector of the mass point making the damped oscillations is an almost-periodic function with almost-period $`T`$.
Let’s remind that according to (1.2) the damped oscillation cyclic frequency $`\omega `$ of mass point is equal to
$$\omega =\sqrt{\omega _{o}^{}{}_{}{}^{2}\beta ^2}.$$
(16)
Obviously the quantity $`\omega `$ has the meaning of oscillation frequency only in the case $`\omega _{o}^{}{}_{}{}^{2}<\beta ^2`$. At $`\omega _{o}^{}{}_{}{}^{2}>\beta ^2`$ the $`\omega `$ becomes imaginary and, accordingly, the trigonometrical function $`\mathrm{sin}(\omega t)`$ is transformed to the hyperbolic function $`sh(\omega t)`$. In this case the solution of the damped oscillations equation (11) becomes
$$\stackrel{}{r}(t)=\stackrel{}{r}_oe^{\beta t}sh(\omega t+\phi _o),$$
(17)
or in the coordinate notation
$$x(t)=x_oe^{\beta t}sh(\omega t+\phi _o),$$
(18)
Such a solution is neither a periodic function no an almost-periodic function. And, therefore, the motion described by the equation of damped oscillations at $`\omega _{o}^{}{}_{}{}^{2}>\beta ^2`$ is not an oscillatory motion. This process is named as aperiodic oscillations. The diagram of a function $`x(t)`$ for an aperiodic process (i.e. described by eq.(18) at $`\phi _o=0`$) is presented on fig.2.
## 2 The practices for simulation of physical processes
Before simulation initiation of physical processes it is necessary to familiarize with blanket rules of operation with the digital computer and simple set of usual activities used at operation with the Borland software menu.
| In this laboratory work there is an opportunity |
| --- |
| to get the help information on its problem (and another) |
| at any moment not quitting the program. |
To obtain the help information it is necessary to press the key $`F1`$.
The set of practices performed by the student at the study of oscillatory motions are determined by the teacher and can vary over a wide range.
Let’s consider practices the realization of which is necessary for understanding of features of the mass point oscillatory motion at presence (and absence) of resistance forces. According to the object of the laboratory work there should be two such practices.
### 2.1 Free simple harmonic motions
In this practice state problem to study an oscillatory motion of a mass point at absence of resistance force. Namely:
* to make sure that a trajectory of a mass point is the harmonic function;
* to find out how the mass point trajectory varies by change of the following parameters:
+ particle mass $`m`$ and stiffness coefficient $`k`$ in expression for elastic force (eq. (4))
+ initial kinematic parameters of a motion: the mass point coordinates of the origin $`x(0)=x_o\mathrm{sin}(\phi _o)`$ and its initial velocity $`v(0)=x_o\omega \mathrm{cos}(\phi _o)`$ (they are those parameters, you can change by changing value of the oscillations epoch angle $`\phi _o`$);
The practice consists of the following items:
#### 2.1.1
After you entered in the menu and selected necessary laboratory work (i.e. ”The oscillatory motion”) the title page of this work arises.
Press $`Enter`$ then the main menu with titles of all practices will arise. By keys $``$ and $``$ it is necessary to select practice ”Simple harmonic Motions” and press $`Enter`$.
#### 2.1.2
You will pass in the first dialog box ”Parameters of the system”. In this box you must set the particle mass and stiffness coefficient in SI units<sup>2</sup><sup>2</sup>2In the line of the context-sensitive help (bottom line of the display) range of values is indicated within the bounds of all parameters numerical values, you can change. and write down those values to table 1.
Let’s note that the ending of input in all dialog boxes is possible by two paths:
* by pressing the key $`Enter`$;
* by activation of the dialog box button Ok (with the help of the device $`Mouse`$)
#### 2.1.3
In the following dialog box (according to its title ”Epoch angle”) you should choose an epoch angle $`\phi _o`$ of a mass point oscillatory motion<sup>2</sup><sup>2</sup>footnotemark: 2. Then write down this value to table 1 and press $`Enter`$ .
#### 2.1.4
You will see the diagram representing a trajectory of a free harmonic oscillatory motion of a mass point with parameters chosen by you<sup>3</sup><sup>3</sup>3If the obtained diagrams satisfy the object of the practice (in your opinion) then sketch these diagrams in yours writing-book and press any key (according to the message in the bottom line)..
#### 2.1.5
You will pass at the next dialog box ”Change of mass”. Enter another value of the particle mass (in comparison with the value entered in item 2.1.2 i.e. in the dialog box ”Parameters of the system” ). You will see two diagrams corresponding to different values of the particle mass with unchangeable others parameters. By pressing any key (according to the message in the bottom line) you will return to the same dialog box again. Iterate the described activities for one more value of mass. As a result you will see three diagrams corresponding to three values of the particle mass is in a simple harmonic motions<sup>3</sup><sup>3</sup>footnotemark: 3.
#### 2.1.6
You will pass at the next dialog box ”Change of K”. Enter another value of the stiffness coefficient $`k`$ (in comparison with value entered in item 2.1.2 i.e. in the dialog box ”Parameters of the system” ). You will see two diagrams corresponding to different values of the stiffness coefficient with unchangeable others parameters. By pressing any key (according to the message in the bottom line) you will return to the same dialog box again. Iterate the described activities for one more stiffness coefficient. In result you will see three diagrams corresponding to three values of the stiffness coefficient of the quasi-elastic force (by due to action of this force the mass point is in a simple harmonic motions)<sup>3</sup><sup>3</sup>footnotemark: 3.
#### 2.1.7
You will pass at the next dialog box ”Change of epoch angle”. Enter another value of the epoch angle (in comparison with value entered in item 2.1.3 i.e. in the dialog box ”Epoch angle” ). You will see two diagrams corresponding to different values of the epoch angle with unchangeable others parameters. By pressing any key (according to the message in the bottom line) you will return to the same dialog box again. Iterate the described activities for one more value of epoch angle. In result you will see three diagrams corresponding to three values of the epoch angle of the simple harmonic motions<sup>3</sup><sup>3</sup>footnotemark: 3.
So the first practice is ended and you will be returned to the main menu. Let’s note that after you exit out of the first practice you can not enter there once again. Therefore, if you do not accept results of this practice and you want to iterate it you should start the program once again.
### 2.2 Damped oscillations
In this practice state problem to study an oscillatory motion of a mass point with the resistance force is proportional to a vector of velocity. Namely:
* to make sure that the trajectory of a mass point is the non harmonic function, but almost-periodic function;
* to find out how the mass point trajectory varies by change of the following parameters:
+ particle mass $`m`$, stiffness coefficient $`k`$ in expression for the elastic force (eq. (4)) and the resistance coefficient $`b`$ for the resistance force (eq. (7))
+ initial kinematic parameters of a motion: the mass point coordinates of the origin $`x(0)=x_o\mathrm{sin}(\phi _o)`$ and its initial velocity $`v(0)=x_o\omega \mathrm{cos}(\phi _o)`$ (they are those parameters you can change by changing value of the oscillations epoch angle $`\phi _o`$);
The practice consists of the following items:
#### 2.2.1
After you exited out of the first practice you will see the main menu with titles of all practices again. By keys $``$ and $``$ it is necessary to select the practice ”Damped oscillations” and press $`Enter`$.
#### 2.2.2
You will pass in the first dialog box ”Parameters of the system”. In this box you must set the particle mass, stiffness coefficient and resistance coefficient in SI units<sup>2</sup><sup>2</sup>footnotemark: 2 and write down those values to table 2.
#### 2.2.3
In the following dialog box (according to its title ”Epoch angle”) you should choose an epoch angle $`\phi _o`$ of mass point oscillatory motion<sup>2</sup><sup>2</sup>footnotemark: 2. Then write down this value to table 2 and press $`Enter`$ .
#### 2.2.4
You will see the diagram representing a trajectory of a damped oscillatory motion of a mass point with parameters chosen by you<sup>3</sup><sup>3</sup>footnotemark: 3.
#### 2.2.5
You will pass at the next dialog box ”Change of the parameters”. Here you can change values of two parameters: the particle mass $`m`$ and resistance coefficient $`b`$. In contrast to the first practice, you can have this box as much as long. Because after each new diagram (for the next pair of parameters m and b) you will be returned here. However, as well as in the first practice, at the display draw no more three diagrams. Therefore we recommend to act as follows:
* at first draw three diagrams with different values of the particle mass $`m`$ and the constant resistance coefficient $`b`$<sup>3</sup><sup>3</sup>footnotemark: 3;
* then draw three diagrams with different values of the resistance coefficient $`b`$ and the constant particle mass $`m`$<sup>3</sup><sup>3</sup>footnotemark: 3;
* then select the such underload resistance coefficient $`b_{min}`$ (with the constant particle mass $`m`$) for which the particle motion will become aperiodic; write down values $`m`$ and $`b_{min}`$ to table 2;
* iterate operations of the previous item for another two values of the particle mass $`m`$;
* then enter pairwise obtained values of mass $`m`$ and coefficient $`b_{min}`$ (from table 2) so to obtain all three aperiodic motions on one picture (i.e. in one frame) and sketch these diagrams in your writing-book;
So the second practice is ended. In order to return to the main menu (if you have the dialog box ”Change of parameters”) is necessary to press the key Esc (or make active the dialog box button Exit by the device $`Mouse`$).
## 3 Return<sup>4</sup><sup>4</sup>4Title page of the return at the laboratory work on physical processes simulation one should draw up on the same rules that the title page of the return at the laboratory work is done in chair T&EPh experimental laboratories.
### 3.1 Contents of the return
The return should include the following items:
* Object of work.
* Summary theoretical part.
* Practices:
+ Free simple harmonic motions.
T A B L E 1
Diagrams $`x=x(t)`$ with different $`m`$ for all three trajectories in one frame.
Diagrams $`x=x(t)`$ with different $`k`$ for all three trajectories in one frame.
Diagrams $`x=x(t)`$ with different $`\phi _o`$ for all three trajectories in one frame.
+ Damped oscillations.
T A B L E 2
Diagrams $`x=x(t)`$ with different $`m`$ for all three trajectories in one frame.
Diagrams $`x=x(t)`$ with different $`b`$ for all three trajectories in one frame.
Diagrams $`x=x(t)`$ with different $`b_{min}`$ for all three trajectories in one frame.
* Conclusion
### 3.2 Design of the tables.
The tables used in the return should be designed by following ways.
| T A B L E 1 | | | |
| --- | --- | --- | --- |
| | Change $`m`$ | Change $`k`$ | Change $`\phi _o`$ |
| $`m,`$ | | | |
| $`k,`$ | | | |
| $`\phi _o,`$ | | | |
| T A B L E 2 | | | |
| --- | --- | --- | --- |
| | Change $`m`$ | Change $`b`$ | Change $`b_{min}`$ |
| $`m,`$ | | | |
| $`b,`$ | | | |
| $`b_{min},`$ | | | |
| $`k,`$ | | | |
| $`\phi _o,`$ | | | |
## 4 Conclusion
Let’s mark, that this paper is written on the basis of the previous works carried out on chair T&EPh under the author leadership (or direct participation). |
warning/0002/math0002235.html | ar5iv | text | # Formal biholomorphic maps of real analytic hypersurfaces
## 1. Introduction
A formal (holomorphic) mapping $`f:(^n,p)(^n,p^{})`$, $`p,p^{}^n`$, $`n1`$, is a vector $`(f_1(z),\mathrm{},f_n(z))`$ where each $`f_j(z)[[zp]]`$, the ring of formal holomorphic power series in $`zp`$, and $`f(p)=p^{}`$. The mapping $`f`$ is called a formal biholomorphism if its Jacobian does not vanish at $`p`$. If $`M,M^{}`$ are two smooth real real-analytic hypersurfaces in $`^n`$ through $`p`$ and $`p^{}`$ respectively, we say that a formal mapping $`f`$ as above sends $`M`$ into $`M^{}`$ if $`\rho ^{}(f(z),\overline{f(z)})=a(z,\overline{z})\rho (z,\overline{z})`$, where $`\rho ,\rho ^{}`$ are local real-analytic defining functions for $`(M,p)`$ and $`(M^{},p^{})`$ respectively and $`a[[zp,\overline{z}\overline{p}]]`$. In this paper we study the convergence (and partial convergence) of formal biholomorphic mappings between germs of real analytic hypersurfaces in $`^n`$ in terms of optimal and natural geometric conditions on the source and target manifolds.
A natural geometric condition which appears in this regularity problem is the concept of holomorphic nondegeneracy. Following Stanton, a real analytic hypersurface $`M^n`$ is called holomorphically nondegenerate if, near any point in $`M`$, there is no non-trivial holomorphic vector field, with holomorphic coefficients, tangent to $`M`$ near that point . Baouendi and Rothschild recognized the importance of such a condition and used it to characterize those real algebraic hypersurfaces for which any biholomorphic self-map must be algebraic (see also ). In this paper we establish a similar statement for formal biholomorphic mappings of real analytic hypersurfaces (Theorem 2.2 below), namely that any formal biholomorphism between germs of real analytic minimal and holomorphically nondegenerate hypersurfaces is convergent. This statement will in fact be a consequence of our main result, Theorem 2.1, where a description of the analyticity properties of formal biholomorphic maps of minimal real analytic hypersurfaces is given. An application of this theorem to partial convergence of such maps is also given in §7.
The study of the convergence of formal mappings between real analytic CR submanifolds goes back to Chern and Moser , who established the convergence of formal biholomorphisms between real analytic Levi-nondegenerate hypersurfaces. More recently, Baouendi, Ebenfelt and Rothschild addressed this problem in more general situations. In particular, the following two facts follow from their work:
i) Given a generic minimal real analytic connected holomorphically nondegenerate submanifold $`M`$ in $`^N`$, $`N2`$, there exists a proper real analytic subvariety $`SM`$, such that for any point $`pMS`$, any formal biholomorphism sending $`(M,p)`$ onto another germ of a generic real analytic submanifold in $`^N`$, must be convergent.
ii) If $`M`$ is a connected generic holomorphically degenerate submanifold in $`^N`$, then, for any point $`pM`$, there exists a formal biholomorphism sending $`(M,p)`$ into itself which does not converge.
In view of these facts and other related results in the mapping problems (), it appears likely that the subvariety $`S`$ in i) could be taken to be empty. In this paper, we actually prove such a result in the case of hypersurfaces. Our proof is based on an analysis of a so-called reflection function associated to the mapping and the hypersurfaces, which was used in many other situations (cf. ). We should also mention that, recently in , we showed that the real analytic subvariety $`S`$ in i) can also be taken to be empty when the generic submanifold $`M`$ is assumed to be real algebraic (and with no such assumption on the target manifold).
## 2. Statement of main results
Let $`(M^{},p^{})^n`$, $`n2`$, be a germ at $`p^{}`$ of a smooth real real-analytic hypersurface. Let $`\rho ^{}=\rho ^{}(\zeta ,\overline{\zeta })`$ be a real analytic defining function for $`M^{}`$ near $`p^{}`$, i.e.
$$M^{}=\{\zeta (^n,p^{}):\rho ^{}(\zeta ,\overline{\zeta )}=0\}.$$
After complexification of $`\rho ^{}`$, one defines the so-called invariant Segre varieties attached to $`M^{}`$ by
$$Q_\omega ^{}=\{\zeta (^n,p^{}):\rho ^{}(\zeta ,\overline{\omega })=0\},$$
for $`\omega `$ close to $`p^{}`$. We can assume, without loss of generality, that $`p^{}=0`$ and $`\frac{\rho ^{}}{\zeta _n}(0)0`$, $`\zeta =(\zeta ^{},\zeta _n)^{n1}\times `$. Thus, the implicit function theorem exhibits any Segre variety as a graph of the form
$$Q_\omega ^{}=\{\zeta (^n,0):\zeta _n=\mathrm{\Phi }^{}(\overline{\omega },\zeta ^{})\},$$
where $`\mathrm{\Phi }^{}`$ is a holomorphic function in its arguments in a neighborhood of $`0^{2n1}`$ and such that $`\mathrm{\Phi }^{}(0)=0`$. Equivalently, this Segre variety can be defined by
$$Q_\omega ^{}=\{\zeta (^n,0):\overline{\zeta }_n=\overline{\mathrm{\Phi }}^{}(\omega ,\overline{\zeta }^{})\}.$$
Here, we have used the following notation. If $`g=g(x)`$ is some formal holomorphic power series in $`[[x]]`$, $`x=(x_1,\mathrm{},x_k)`$, $`\overline{g}`$ is the formal holomorphic power series obtained by taking the complex conjugates of the coefficients of $`g`$. (This convention of notation will be used throughout the paper.) Our main result is the following.
###### Theorem 2.1.
Let $`f:(M,0)(M^{},0)`$ be a formal biholomorphism between two germs at 0 of real-analytic hypersurfaces in $`^n`$. Assume, furthermore, that $`M`$ is minimal at $`0`$. Then, the formal holomorphic map
$$^n\times ^{n1}(z,\lambda )\overline{\mathrm{\Phi }}^{}(f(z),\lambda )$$
is convergent.
Such a result has several applications. One of its main applications lies in the following theorem mentioned in the introduction.
###### Theorem 2.2.
Any formal biholomorphic mapping between germs of minimal, holomorphically nondegenerate, real analytic hypersurfaces in $`^n`$ is convergent.
As explained in the introduction, the interest in such a result lies in the fact that, in view of , the condition of holomorphic nondegeneracy is optimal for the class of formal biholomorphisms. However, it remains an open problem to decide whether or not the condition of minimality is necessary in Theorem 2.2. Another application of Theorem 2.1 deals with partial convergence of formal biholomorphisms. For this, we refer the reader to §7.
###### Remark 1.
After this work was completed, I received a preprint by J. Merker, “Convergence of formal biholomorphic mappings between minimal holomorphically nondegenerate real analytic hypersurfaces”, in which a similar statement to Theorem 2.2 is given. In that preprint, Theorem 2.1 above is also stated for the special case where $`M^{}`$ is a rigid and polynomial hypersurface.
## 3. Definitions and notations
### 3.1. Real analytic hypersurfaces.
Let $`M`$ be a smooth real real-analytic hypersurface in $`^n`$, $`n2`$. Since the situation is purely local, we shall always work near a point $`pM`$, which will be assumed, without loss of generality, to be the origin. Let $`\rho `$ be a real analytic defining function for $`M`$ near 0 i.e.
$$M=\{z(^n,0):\rho (z,\overline{z})=0\},$$
with $`d\rho 0`$ on $`M`$. The complexification $``$ of $`M`$ is the complex hypersurface through 0 in $`^{2n}`$ given by
$$=\{(z,w)(^{2n},0):\rho (z,w)=0\}.$$
We shall assume, without loss of generality, that the coordinates $`z=(z^{},z_n)^n`$ are chosen so that $`{\displaystyle \frac{\rho }{z_n}}(0)0`$. In this case, we define the following holomorphic vector fields tangent to $``$
(3.1)
$$_j=\frac{\rho }{w_n}(z,w)\frac{}{w_j}\frac{\rho }{w_j}(z,w)\frac{}{w_n},j=1,\mathrm{},n1,$$
which are the complexifications of the usual (0,1) vector fields tangent to $`M`$. We recall that for a point $`w`$ near 0, its associated Segre surface is the complex hypersurface defined by $`Q_w=\{z(^n,0):\rho (z,\overline{w})=0\}`$. Observe that by the complex analytic implicit function theorem, each Segre variety can be described as a graph of the form
$$Q_w=\{z(^n,0):z_n=\mathrm{\Phi }(\overline{w},z^{})\},$$
$`\mathrm{\Phi }`$ denoting a convergent power series in some neighborhood of the origin in $`^{2n1}`$ satisfying the relations $`\mathrm{\Phi }(0)=0`$ and
(3.2)
$$\mathrm{\Phi }(w^{},\overline{\mathrm{\Phi }}(z,w^{}),z^{})z_n,(z,w^{})^n\times ^{n1}.$$
Equation (3.2) is a consequence of the fact that $`M`$ is a real hypersurface. The coordinates $`z`$ are said to be normal with respect to $`M`$ if the additional condition
$$\mathrm{\Phi }(w^{},w_n,0)=\mathrm{\Phi }(0,w_n,z^{})w_n,(w^{},w_n,z^{})^{n1}\times \times ^{n1},$$
holds. It is well-known that given a real-analytic hypersurface through the origin, one can always construct such coordinates . Thus, from now on and for simplicity, we will always assume that the $`z`$-coordinates are chosen to be normal for the manifold $`M`$.
The real analytic hypersurface $`M`$ is called minimal at 0 (in the sense of Trépreau and Tumanov), or, equivalently of finite type (in the sense of Kohn and Bloom-Graham) if it does not contain any complex-analytic hypersurface through 0. To use such a nondegeneracy condition, we will need the Segre set mappings associated to $`M`$ (see ) up to order 3, which, in normal coordinates, are the following three maps :
(3.3) $`(^{n1},0)z^{}v_1(z^{})`$ $`=(z^{},0)^n,`$
$`(^{2n2},0)(z^{},\xi )v_2(z^{},\xi )`$ $`=(z^{},\mathrm{\Phi }(\xi ,0,z^{}))^n,`$
$`(^{3n3},0)(z^{},\xi ,\eta )v_3(z^{},\xi ,\eta )`$ $`=(z^{},\mathrm{\Phi }(\xi ,\overline{\mathrm{\Phi }}(\eta ,0,\xi ),z^{}))^n.`$
These maps are of fundamental importance since they parametrize the so-called Segre sets (up to order three) associated to $`M`$. Moreover, the interest in such Segre sets or Segre set mappings lies in the fact that the minimality assumption is equivalent to the fact that the generic rank of $`v_2`$ (and also $`v_3`$) equals $`n`$ (see ). This will be useful for the proof of Theorem 2.1.
All the notations introduced in this section will be used for the source hypersurface $`M`$. For the target real analytic hypersurface $`M^{}`$, we shall use the notations introduced in §2, before the statement of Theorem 2.1. In particular, we denote by $`z`$ the coordinates in the source space and by $`\zeta `$ the coordinates at the target space.
### 3.2. Some commutative algebra.
We recall here, for the reader’s convenience, some basic definitions, needed in §7, about regular local rings and their ideals. All these definitions can be found for instance in .
Let $`A`$ be a Noetherian ring. If $`I`$ and $`J`$ are two ideals of $`A`$, we use the notation $`I<J`$ to mean that $`IJ`$ and $`IJ`$. Given a prime ideal $`IA`$, the height of $`I`$ is defined by the formula
$$h(I)=\mathrm{max}\{k:\{0\}<I_1<\mathrm{}<I_k=I\},$$
where $`I_1,\mathrm{},I_k`$ are prime ideals of $`A`$. If $`J`$ is any ideal of $`A`$, we define the height of $`J`$ by the formula
$$h(J)=\mathrm{inf}\{h(I):JI,I\mathrm{prime}\mathrm{ideal}\mathrm{of}A\}.$$
If $`A`$ is furthermore assumed to be a local ring, one defines the Krull dimension of $`A`$ to be the height of its maximal ideal. Observe that if $`A`$ is a Noetherian local ring and if $`I`$ is a proper ideal of $`A`$, the quotient ring $`A/I`$ is also a Noetherian local ring. This allows one to consider the Krull dimension of such a ring. Finally, a Noetherian local ring is said to be regular if its maximal ideal has $`\delta `$ generators, where $`\delta `$ is the Krull dimension of the ring $`A`$. The rings of formal holomorphic power series or convergent power series in $`p`$ indeterminates, $`p^{}`$, are regular rings of Krull dimension $`p`$ .
## 4. Two convergence results
In this section, we first state and prove a convenient lemma which will be used twice in the paper. This lemma may already be known.
###### Lemma 4.1.
Let $`(u_i(t))_{iI}`$ be a family of convergent power series in $`\{t\}`$, $`t=(t_1,\mathrm{},t_q)`$, $`q^{}`$. Let also $`(𝒦_i(\varsigma ))_{iI}`$ be a family of convergent power series in $`\{\varsigma \}`$, $`\varsigma =(\varsigma _1,\mathrm{},\varsigma _r)`$, $`r^{}`$. Assume that:
1. There exists $`R>0`$ such that the radius of convergence of any $`𝒦_i`$, $`iI`$, is at least $`R`$.
2. For all $`\varsigma ^r`$ with $`|\varsigma |<R`$, $`|𝒦_i(\varsigma )|C_i`$, with $`C_i>0`$.
3. There exists $`V(t)=(V_1(t),\mathrm{},V_r(t))([[t]])^r`$, $`V(0)=0`$, such that $`(𝒦_iV)(t)=u_i(t)`$ (in $`[[t]]`$) for all $`iI`$.
Then, there exists $`R^{}>0`$ such that the radius of convergence of any $`u_i`$, $`iI`$, is at least $`R^{}`$ and such that for all $`t^q`$ with $`|t|<R^{}`$, $`|u_i(t)|C_i`$.
###### Proof.
By Artin’s approximation theorem , one can find a convergent power series mapping $`\vartheta (t)=(\vartheta _1(t),\mathrm{},\vartheta _r(t))(\{t\})^r`$ such that $`\vartheta (0)=0`$ and for all $`iI`$, $`(𝒦_i\vartheta )(t)=u_i(t)`$ in $`\{t\}`$. Let $`R^{}>0`$ so that if $`|t|<R^{}`$ then $`|\vartheta (t)|<R`$. Since the radius of convergence of the $`𝒦_i`$, $`iI`$, is at least $`R`$, the radius of convergence of any $`𝒦_i\vartheta `$, $`iI`$, is at least $`R^{}`$. Thus, the family $`(u_i(t))_{iI}`$ has a radius of convergence at least equal to $`R^{}`$ and for all $`t^q`$ with $`|t|<R^{}`$, one has $`|u_i(t)|=|(𝒦_i\vartheta )(t)|C_i`$, $`iI`$.∎
To derive Theorem 2.2 from Theorem 2.1, we will need the following consequence of Artin’s approximation theorem, which is contained for instance in . (See also §7 for another formulation of such a result.)
###### Proposition 4.2.
Let $`R(x,y)=(R_1(x,y),\mathrm{},R_r(x,y))(\{x,y\})^r`$, $`x^q`$, $`y^r`$, $`q,r^{}`$. Let $`g(x)=(g_1(x),\mathrm{},g_r(x))([[x]])^r`$ satisfy $`R(x,g(x))=0`$. If $`\mathrm{det}\left({\displaystyle \frac{R}{y}}(x,g(x))\right)0`$ in $`[[x]]`$, then $`g(x)`$ is convergent.
###### Proof.
We reproduce here the arguments of . Write
(4.1)
$$R(x,y)R(x,z)=Q(x,y,z)(yz)$$
where $`Q`$ is an $`r\times r`$ complex-analytic matrix such that $`Q(x,y,y)={\displaystyle \frac{R}{y}}(x,y)`$; i.e. $`Q(x,y,z)={\displaystyle _0^1}{\displaystyle \frac{R}{y}}(x,ty+(1t)z)𝑑t`$. By assumption, we know that we have $`\mathrm{det}Q(x,g(x),g(x))0`$. This implies that one can find an integer $`k_g`$ such that if $`H(x)`$ is any formal power series which agrees up to order $`k_g`$ with $`g`$ then $`\mathrm{det}Q(x,g(x),H(x))0`$. For this integer $`k_g`$, according to Artin’s approximation theorem, one can find a convergent power series $`H_0(x)`$ satisfying $`R(x,H_0(x))=0`$ and agreeing with $`g(x)`$ up to order $`k_g`$. By (4.1), we get $`Q(x,g(x),H_0(x))(g(x)H_0(x))0`$ in $`[[x]]`$. Since $`\mathrm{det}Q(x,g(x),H_0(x))0`$, we obtain $`g(x)=H_0(x)`$ and thus $`g`$ is convergent.∎
## 5. The reflection principle
Let $`f:(M,0)(M^{},0)`$ be a formal biholomorphic mapping between germs at 0 of real analytic hypersurfaces in $`^n`$. We shall use the notations introduced in §2 and §3. In particular, we denote $`f=f(z)=(f_1(z),\mathrm{},f_n(z))=(f^{}(z),f_n(z))`$ in the $`\zeta `$-coordinates. In this section, we shall make no further assumptions on $`M`$ and $`M^{}`$.
As in , we define the following formal holomorphic power series
(5.1)
$$^n\times ^{n1}(z,\lambda )(z,\lambda ):=\overline{\mathrm{\Phi }}^{}(f(z),\lambda ).$$
The goal of this section is to prove the following proposition.
###### Proposition 5.1.
Let $`f:(M,0)(M^{},0)`$ be a formal biholomorphism between germs at 0 of real analytic hypersurfaces in $`^n`$, and $``$ defined by (5.1). Then for any multi-index $`\gamma ^n`$, the formal holomorphic map
$$(^{2n2},0)(z^{},\lambda )\left(_z^\gamma (z,\lambda )\right)|_{z=v_1(z^{})}$$
is convergent in some neighborhood $`V_\gamma `$ of $`0^{2n2}`$. Here, $`v_1`$ is the first Segre set mapping for $`M`$ as defined in (3.3).
Before proceeding to the proof of Proposition 5.1, we need a preliminary lemma (Lemma 5.2 below). Since $`f`$ maps formally $`M`$ into $`M^{}`$, there exists $`a(z,\overline{z})[[z,\overline{z}]]`$ such that
$$\overline{f_n(z)}\overline{\mathrm{\Phi }}^{}(f(z),\overline{f^{}(z)})=a(z,\overline{z})\rho (z,\overline{z}),\mathrm{in}[[z,\overline{z}]].$$
Equivalently, we have
(5.2)
$$\overline{f}_n(w)\overline{\mathrm{\Phi }}^{}(f(z),\overline{f}^{}(w))=a(z,w)\rho (z,w),\mathrm{in}[[z,w]].$$
We write $`\overline{\mathrm{\Phi }}_{\lambda ^\alpha }^{}(\omega ,\lambda )`$ for $`_\lambda ^\alpha \overline{\mathrm{\Phi }}^{}(\omega ,\lambda )`$. By applying the vector fields $`_j`$, $`j=1,\mathrm{},n1`$, as defined by (3.1), to (5.2) and using the fact that $`f`$ is invertible, one obtains the following known statement (see for instance).
###### Lemma 5.2.
Under the assumptions of Proposition 5.1, one has, for any multindex $`\alpha ^{n1}`$, the formal power series identity
$$\overline{\mathrm{\Phi }}_{\lambda ^\alpha }^{}(f(z),\overline{f}^{}(w))=\chi _\alpha ((^\beta \overline{f}(w))_{|\beta ||\alpha |},z,w),(z,w),$$
where each $`\chi _\alpha `$ is a convergent power series of its arguments.
###### Proof of Proposition 5.1..
We write the expansion
(5.3)
$$\overline{\mathrm{\Phi }}^{}(\omega ,\lambda )=\underset{\alpha ^{n1}}{}\varphi _\alpha ^{}(\omega )\lambda ^\alpha .$$
For the sake of clarity, we shall first give the proof of the Proposition in the case $`\gamma =0`$.
The case $`\gamma =0`$. We restrict all the identities given by Lemma 5.2 to the $`(n1)`$-dimensional subspace
$$\{(0,v_1(z^{})):z^{}(^{n1},0)\}.$$
This gives, for any multiindex $`\alpha ^{n1}`$,
(5.4)
$$\alpha !(\varphi _\alpha ^{}fv_1)(z^{})=\chi _\alpha ((^\beta \overline{f}(0))_{|\beta ||\alpha |},v_1(z^{}),0):=u_\alpha (z^{}),$$
where $`\varphi _\alpha ^{}`$ is given by (5.3). Observe that for each multiindex $`\alpha `$, $`u_\alpha (z^{})`$ is convergent.
To show that $`^{n1}\times ^{n1}(z^{},\lambda )(v_1(z^{}),\lambda )`$ is convergent, we claim that it suffices to show that there exists $`a>0`$ and $`R_0>0`$ such that the radius of convergence of each $`u_\alpha `$ is at least $`a`$, and such that the following Cauchy estimates hold:
(5.5)
$$\alpha ^{n1},z^{}^{n1},|z^{}|<a,|u_\alpha (z^{})|\alpha !R_0^{|\alpha |+1}.$$
Indeed, if (5.5) holds, then the formal holomorphic power series
$$^{n1}\times ^{n1}(z^{},\lambda )_0(z^{},\lambda ):=\underset{|\alpha |=0}{\overset{\mathrm{}}{}}\frac{u_\alpha (z^{})}{\alpha !}\lambda ^\alpha $$
defines a convergent power series in $`B_{n1}(0,a)\times B_{n1}(0,1/2R_0)`$. (Here and in what follows, for any $`c>0`$ and for any $`k^{}`$, $`B_k(0,c)`$ denotes the euclidean ball centered at 0 in $`^k`$ of radius $`c`$.) Moreover, by (5.4) and (5.1), we have for any multiindex $`\alpha ^{n1}`$,
(5.6)
$$\left[\frac{^{|\alpha |}_0}{\lambda ^\alpha }(z^{},\lambda )\right]_{\lambda =0}=\left[\frac{^{|\alpha |}}{\lambda ^\alpha }(v_1(z^{}),\lambda )\right]_{\lambda =0}$$
in $`[[z^{}]]`$ and hence
$$_0(z^{},\lambda )=(v_1(z^{}),\lambda ).$$
This proves that under the assumption (5.5), $`(v_1(z^{}),\lambda )`$ is a convergent power series in $`(z^{},\lambda )`$. It remains to prove (5.5). Since $`\overline{\mathrm{\Phi }}^{}`$ is holomorphic in a neighborhood of $`0^{2n1}`$, in view of (5.3), one can find $`\delta >0`$ and a constant $`R>0`$ such that for any multiindices $`\alpha ^{n1}`$, $`\nu ^n`$,
(5.7)
$$\omega ^n,|\omega |<\delta ,\left|\frac{^{|\nu |}\varphi _\alpha ^{}(\omega )}{\omega ^\nu }\right|\nu !R^{|\alpha |+|\nu |+1}.$$
In view of (5.4) and (5.7) (in the case $`\nu =0`$), we can apply Lemma 4.1 to conclude that there exists $`a>0`$ such that the family $`(u_\alpha (z^{}))_{\alpha ^{n1}}`$ is convergent in $`B_{n1}(0,a)`$ and such that (5.5) holds with $`R_0=R`$. This finishes the proof of Proposition 5.1 in the case $`\gamma =0`$.
The case $`|\gamma |>0`$. We proceed now to the proof of Proposition 5.1 for general $`\gamma ^n`$. For this, we need the following lemma.
###### Lemma 5.3.
Under the assumptions of Proposition 5.1, for any multiindices $`\alpha ^{n1}`$, $`\gamma ^n`$, the formal holomorphic power series
$$^{n1}z^{}_z^\gamma \left((\varphi _\alpha ^{}f)(z)\right)|_{z=v_1(z^{})}$$
is convergent.
###### Proof of Lemma 5.3..
We prove the Lemma by induction on $`|\gamma |`$ (for any multiindex $`\alpha ^{n1}`$). For $`\gamma =0`$, the statement follows from (5.4), as we previously noticed. Let $`\gamma ^n`$. For $`\alpha ^{n1}`$, using Lemma 5.2, we obtain for $`(z^{},z_n,0,z_n)(,0)`$,
$$\overline{\mathrm{\Phi }}_{\lambda ^\alpha }^{}(f(z),\overline{f}^{}(0,z_n))=\chi _\alpha ((^\beta \overline{f}(0,z_n))_{|\beta ||\alpha |},z,0,z_n).$$
If we apply $`_z^\gamma `$ to this equation, we obtain
(5.8)
$$\frac{^{|\gamma |}}{z^\gamma }\left[_{\lambda ^\alpha }(z,\overline{f}^{}(0,z_n))\right]=\frac{^{|\gamma |}}{z^\gamma }\left[\chi _\alpha ((^\beta \overline{f}(0,z_n))_{|\beta ||\alpha |},z,0,z_n)\right].$$
One can easily check that this implies that there exist a polynomial $`𝒮_\gamma `$ such that the left-hand side of (5.8) is equal to
(5.9)
$$\begin{array}{c}_{z^\gamma \lambda ^\alpha }(z,\overline{f}^{}(0,z_n))+\hfill \\ \hfill 𝒮_\gamma [\left(_{z^\nu \lambda ^\beta }(z,\overline{f}^{}(0,z_n))\right)_{\genfrac{}{}{0pt}{}{|\nu |<|\gamma |}{|\beta ||\alpha |+|\gamma |}},\left(^\mu \overline{f}(0,z_n)\right)_{|\mu ||\gamma |}],\end{array}$$
where $`\mu ,\nu ^n`$, $`\beta ^{n1}`$. Furthermore, we observe that the right-hand side of (5.8) can be written in the form
(5.10)
$$\chi _{\alpha ,\gamma }^1((^\beta \overline{f}(0,z_n))_{|\beta ||\alpha |+|\gamma |},z),$$
where $`\chi _{\alpha ,\gamma }^1`$ is a convergent power series. Restricting (5.8), (5.9) and (5.10) to $`z=v_1(z^{})`$, one obtains
(5.11)
$$\begin{array}{c}\alpha !_z^\gamma \left((\varphi _\alpha ^{}f)(z)\right)|_{z=v_1(z^{})}+\hfill \\ \hfill 𝒮_\gamma [\left(\beta !_z^\nu \left((\varphi _\beta ^{}f)(z)\right)|_{z=v_1(z^{})}\right)_{\genfrac{}{}{0pt}{}{|\nu |<|\gamma |}{|\beta ||\alpha |+|\gamma |}},(^\mu \overline{f}(0))_{|\mu ||\gamma |}]=\\ \hfill \chi _{\alpha ,\gamma }^1((^\beta \overline{f}(0))_{|\beta ||\alpha |+|\gamma |},z^{},0).\end{array}$$
The induction hypothesis tells us that for any multiindex $`\beta ^{n1}`$ and for any multiindex $`\nu ^n`$ such that $`|\nu |<|\gamma |`$, the formal holomorphic power series
$$z^{}^{n1}_z^\nu \left((\varphi _\beta ^{}f)(z)\right)|_{z=v_1(z^{})}$$
is convergent. Thus, we obtain the desired similar statement for $`^{n1}z^{}_z^\gamma \left((\varphi _\alpha ^{}f)(z)\right)|_{z=v_1(z^{})}`$, for any multiindex $`\alpha ^{n1}`$.∎
We come back to the proof of Proposition 5.1. For all multiindices $`\gamma ^n,\alpha ^{n1}`$, we put
(5.12)
$$\mathrm{\Psi }_{\alpha ,\gamma }(z^{}):=\alpha !_z^\gamma \left((\varphi _\alpha ^{}f)(z)\right)|_{z=v_1(z^{})}=_{z^\gamma \lambda ^\alpha }(v_1(z^{}),0).$$
By Lemma 5.3, the $`\mathrm{\Psi }_{\alpha ,\gamma }(z^{})`$ are convergent. We now fix $`\gamma `$, $`|\gamma |1`$. We want to prove that $`_{z^\gamma }(v_1(z^{}),\lambda )`$ is convergent in some neighborhood $`V_\gamma `$ of $`0^{2n2}`$. For this, as in the case $`\gamma =0`$, it suffices to prove that one can find $`q_\gamma >0`$ and $`R_\gamma >0`$ such that the radius of convergence of the family $`(\mathrm{\Psi }_{\alpha ,\gamma })_{\alpha ^{n1}}`$ is at least $`q_\gamma `$ and such that the following estimates hold:
(5.13)
$$\alpha ^{n1},z^{}^{n1},|z^{}|<q_\gamma ,|\mathrm{\Psi }_{\alpha ,\gamma }(z^{})|\alpha !R_\gamma ^{|\alpha |+1}.$$
We first observe that for any multiindex $`\nu ^n`$ with $`|\nu ||\gamma |`$, there exists a universal polynomial $`𝒫_{\nu ,\gamma }`$, such that
(5.14)
$$\begin{array}{c}_{z^\gamma \lambda ^\alpha }(z,\lambda )=_z^\gamma \left(\overline{\mathrm{\Phi }}_{\lambda ^\alpha }^{}(f(z),\lambda )\right)=\hfill \\ \hfill \underset{|\nu ||\gamma |}{}𝒫_{\nu ,\gamma }\left((^\beta f(z))_{1|\beta ||\gamma |}\right)\overline{\mathrm{\Phi }}_{\omega ^\nu \lambda ^\alpha }^{}(f(z),\lambda ).\end{array}$$
This means in particular that the polynomials $`𝒫_{\nu ,\gamma }`$, $`|\nu ||\gamma |`$, are independent of $`\alpha `$. Putting $`\lambda =0`$ and $`z=v_1(z^{})`$ in (5.14), we obtain
(5.15)
$$\begin{array}{c}_{z^\gamma \lambda ^\alpha }(v_1(z^{}),0)=\hfill \\ \hfill \alpha !\underset{|\nu ||\gamma |}{}𝒫_{\nu ,\gamma }\left(((^\beta f)(v_1(z^{})))_{1|\beta ||\gamma |}\right)\left(\frac{^{|\nu |}\varphi _\alpha ^{}}{\omega ^\nu }fv_1\right)(z^{}).\end{array}$$
Recall that $`\gamma `$ is fixed. For $`\alpha ^{n1}`$, consider the convergent power series of the variables $`((\mathrm{\Lambda }_\beta )_{1|\beta ||\gamma |},\omega )`$ defined by
$$h_{\alpha ,\gamma }((\mathrm{\Lambda }_\beta )_{1|\beta ||\gamma |},\omega ):=\alpha !\underset{|\nu ||\gamma |}{}𝒫_{\nu ,\gamma }\left((\mathrm{\Lambda }_\beta +^\beta f(0))_{1|\beta ||\gamma |}\right)\frac{^{|\nu |}\varphi _\alpha ^{}}{\omega ^\nu }(\omega ).$$
Let $`r(\gamma )=n\mathrm{Card}\{\nu ^n:1|\nu ||\gamma |\}`$. In view of (5.7), each $`h_{\alpha ,\gamma }`$, $`\alpha ^{n1}`$, is convergent in $`B_{r(\gamma )}(0,1)\times B_{n1}(0,\delta )`$. Moreover, for $`((\mathrm{\Lambda }_\beta )_{1|\beta ||\gamma |},\omega )B_{r(\gamma )}(0,1)\times B_{n1}(0,\delta )`$, we have the following estimates
$$|h_{\alpha ,\gamma }((\mathrm{\Lambda }_\beta )_{1|\beta ||\gamma |},\omega )|\alpha !\underset{|\nu ||\gamma |}{}|𝒫_{\nu ,\gamma }((\mathrm{\Lambda }_\beta +^\beta f(0))_{1|\beta ||\gamma |})|\nu !R^{|\alpha |+|\nu |+1}.$$
Put
(5.16)
$$\begin{array}{c}C_\gamma :=\mathrm{sup}\{|𝒫_{\nu ,\gamma }((\mathrm{\Lambda }_\beta +^\beta f(0))_{1|\beta ||\gamma |})|:\hfill \\ \hfill |\nu ||\gamma |,(\mathrm{\Lambda }_\beta ^n)_{1|\beta ||\gamma |}B_{r(\gamma )}(0,1)\}.\end{array}$$
This implies that in $`B_{r(\gamma )}(0,1)\times B_{n1}(0,\delta )`$, the following estimates hold for some suitable constant $`C_\gamma ^1`$:
$$|h_{\alpha ,\gamma }((\mathrm{\Lambda }_\beta )_{1|\beta ||\gamma |},\omega )|\alpha !C_\gamma \underset{|\nu ||\gamma |}{}\nu !R^{|\alpha |+|\nu |+1}C_\gamma ^1\alpha !R^{|\alpha |+|\gamma |+1}.$$
From this, we see that there exists $`R_\gamma >0`$ such that for $`((\mathrm{\Lambda }_\beta )_{1|\beta ||\gamma |},\omega )B_{r(\gamma )}(0,1)\times B_{n1}(0,\delta )`$,
(5.17)
$$|h_{\alpha ,\gamma }((\mathrm{\Lambda }_\beta )_{1|\beta ||\gamma |},\omega )|\alpha !R_\gamma ^{|\alpha |+1}.$$
In view of (5.12) and (5.15), we have for any multiindex $`\alpha ^{n1}`$,
$$h_{\alpha ,\gamma }(((^\beta f)(v_1(z^{}))^\beta f(0))_{1|\beta ||\gamma |},(fv_1)(z^{}))=\mathrm{\Psi }_{\alpha ,\gamma }(z^{}),$$
as formal power series in $`z^{}`$. Thus, in view of (5.17), we are in a position to apply Lemma 4.1 to conclude that there exists $`q_\gamma >0`$ such that the family $`(\mathrm{\Psi }_{\alpha ,\gamma }(z^{}))_{\alpha ^{n1}}`$ is convergent in $`B_{n1}(0,q_\gamma )`$, in which, moreover, the desired estimates (5.13) hold. This implies that $`_{z^\gamma }(v_1(z^{}),\lambda )\{z^{},\lambda \}`$. The proof of Proposition 5.1 is thus complete.∎
###### Remark 2.
It is clear that Proposition 5.1 still holds in higher codimension with the same proof. More precisely, the following holds. Let $`M,M^{}`$ be two germs through the origin in $`^n`$, $`n2`$, of smooth real real-analytic generic submanifolds of CR dimension $`N`$ and of real codimension $`d`$. Let $`f:(M,0)(M^{},0)`$ be a formal biholomorphic map. Assume that the coordinates at the target space $`\zeta =(\zeta ^{},\zeta ^{})^N\times ^d`$ are chosen so that, near the origin, $`M^{}=\{(\zeta ^{},\zeta ^{})(^n,0):\overline{\zeta }^{}=\overline{\mathrm{\Phi }}^{}(\zeta ,\overline{\zeta }^{})\}`$, for some $`^d`$-valued holomorphic map $`\overline{\mathrm{\Phi }}^{}=(\overline{\mathrm{\Phi }}_1^{},\mathrm{},\overline{\mathrm{\Phi }}_d^{})`$ near $`0^{2N+d}`$. Assume also that the coordinates at the source space are chosen to be normal coordinates for $`M`$. Then, if we define $`^n\times ^N(z,\lambda )(z,\lambda ):=\overline{\mathrm{\Phi }}^{}(f(z),\lambda )^d`$, the following holds. For any multiindex $`\gamma ^n`$, the formal holomorphic power series mapping
$$^N\times ^N(z^{},\lambda )_{z^\gamma }((z^{},0),\lambda )^d$$
is convergent.
## 6. Proofs of Theorem 2.1 and Theorem 2.2
For the proof of Theorem 2.1, we first need to prove a lemma (also used in ) which will allow us to bypass the second Segre set and to work directly on the third Segre set.
###### Lemma 6.1.
Let $`𝒯(x,u)=(𝒯_1(x,u),\mathrm{},𝒯_r(x,u))([[x,u]])^r`$, $`x^q`$, $`u^s`$, with $`𝒯(0)=0`$. Assume that $`𝒯(x,u)`$ satisfies an identity in the ring $`[[x,u,y]]`$, $`y^q`$, of the form
$$\phi (𝒯(x,u);x,u,y)=0,$$
where $`\phi [[W,x,u,y]]`$ with $`W^r`$. Assume, furthermore, that for any multi-index $`\beta ^q`$, the formal power series $`\left[{\displaystyle \frac{^{|\beta |}\phi }{y^\beta }}(W;x,u,y)\right]_{y=x}`$ is convergent, i.e. belongs to $`\{W,x,u\}`$. Then, for any given positive integer $`e`$, there exists an $`r`$-tuple of convergent power series $`𝒯^e(x,u)(\{x,u\})^r`$ such that $`\phi (𝒯^e(x,u);x,u,y)=0`$ in $`[[x,u,y]]`$ and such that $`𝒯^e(x,u)`$ agrees up to order $`e`$ (at 0) with $`𝒯(x,u)`$.
###### Proof.
First observe that $`𝒯(x,u)`$ is a formal power series solution of the analytic system in the unknown $`W`$,
(6.1)
$$\left[\frac{^{|\beta |}\phi }{y^\beta }(W;x,u,y)\right]_{y=x}0,\beta ^q.$$
Thus, an application of Artin’s approximation theorem gives, for any positive integer $`e`$, an $`r`$-tuple of convergent power series $`𝒯^e(x,u)(\{x,u\})^r`$ solution in $`W`$ of (6.1), and which agrees up to order $`e`$ with $`𝒯(x,u)`$. The Lemma follows by noticing that $`\phi (𝒯^e(x,u);x,u,y)0`$ in $`[[x,u,y]]`$ if and only if $`𝒯^e(x,u)`$ is solution of (6.1). The proof of Lemma 6.1 is complete. ∎
The following proposition will also be useful in the proof of Theorem 2.1 (see for a proof for instance).
###### Proposition 6.2.
Let $`𝒥(x)=(𝒥_1(x),\mathrm{},𝒥_r(x))(\{x\})^r`$, $`x^k`$, $`k,r1`$, $`𝒥(0)=0`$, and $`𝒱(t)[[t]]`$, $`t^r`$. If $`𝒱𝒥`$ is convergent and $`𝒥`$ is generically submersive, then $`𝒱`$ itself is convergent.
###### Proof of Theorem 2.1..
Restricting the identity (5.2) to the set
$$\{(v_3(z^{},\xi ,\eta ),\overline{v}_2(\xi ,\eta )):(z^{},\xi ,\eta )(^{3n3},0)\},$$
where $`v_j`$, $`j=2,3`$, are the Segre sets mappings defined by (3.3), we obtain
(6.2)
$$(v_3(z^{},\xi ,\eta ),(\overline{f}^{}\overline{v}_2)(\xi ,\eta ))=(\overline{f}_n\overline{v}_2)(\xi ,\eta ).$$
Here, $``$ is the formal power series defined by (5.1). We would like to apply Lemma 6.1 to the formal equation (6.2) with $`y=z^{}`$, $`x=\eta `$, $`u=\xi `$, $`𝒯(x,u)=(\overline{f}\overline{v}_2)(\xi ,\eta )`$, $`W=(\lambda ,\mu )`$, $`\lambda ^{n1}`$, $`\mu `$ and
$$\phi ((\lambda ,\mu );\eta ,\xi ,z^{}):=(v_3(z^{},\xi ,\eta ),\lambda )\mu .$$
For this, one has to check that any derivative of the formal holomorphic power series $`^{4n4}(z^{},\xi ,\eta ,\lambda )(v_3(z^{},\xi ,\eta ),\lambda )`$ with respect to $`z^{}`$ evaluated at $`z^{}=\eta `$ is convergent with respect to the variables $`\lambda `$, $`\xi `$ and $`\eta `$. All these derivatives involve derivatives of $`v_3`$ at $`z^{}=\eta `$ (which are convergent) and derivatives of the form $`\left[_{z^\gamma }(v_3(z^{},\xi ,\eta ),\lambda )\right]_{z^{}=\eta }`$, for $`\gamma ^n`$. Because of the reality condition (3.2) and the definition of $`v_1`$ and $`v_3`$ given by (3.3), we have
$$v_3(\eta ,\xi ,\eta )=v_1(\eta ).$$
This implies that for each $`\gamma ^n`$, we have
$$\left[_{z^\gamma }(v_3(z^{},\xi ,\eta ),\lambda )\right]_{z^{}=\eta }=_{z^\gamma }(v_1(\eta ),\lambda ),$$
with the right-hand side being convergent in $`(\eta ,\lambda )`$ by Proposition 5.1. Thus, by Lemma 6.1, we have, for any positive integer $`e`$, a convergent power series mapping, denoted $`𝒯^e(\xi ,\eta )=(𝒯^{}_{}{}^{}e(\xi ,\eta ),𝒯_n^e(\xi ,\eta ))^{n1}\times `$, which agree up to order $`e`$ with $`(\overline{f}\overline{v}_2)(\xi ,\eta )`$ and such that
(6.3)
$$(v_3(z^{},\xi ,\eta ),𝒯^{}_{}{}^{}e(\xi ,\eta ))=𝒯_n^e(\xi ,\eta ),\mathrm{in}[[z,\xi ,\eta ]].$$
Since $`𝒯^e(\xi ,\eta )`$ is a convergent power series mapping, in order to show that $`(z,\lambda )\{z,\lambda \}`$, it suffices to show by Proposition 6.2 that for $`e`$ large enough the generic rank of the holomorphic map
(6.4)
$$(^{3n3},0)(z^{},\xi ,\eta )(v_3(z^{},\xi ,\eta ),𝒯^{}_{}{}^{}e(\xi ,\eta ))^{2n1}$$
is $`2n1`$. For this, note that since $`M`$ is minimal at 0, the map $`\overline{v}_2`$ is of generic rank $`n`$, and thus, by the form of $`v_3`$ given in (3.3), the holomorphic map $`(^{3n3},0)(z^{},\xi ,\eta )(v_3(z^{},\xi ,\eta ),\overline{v}_2(\xi ,\eta ))`$ is of generic rank $`2n1`$ (see ). Moreover, since $`f`$ is invertible, we have $`\mathrm{det}\left(\frac{f^{}}{z^{}}(0)\right)0`$, which implies that the rank of the formal map $`(z,w)(z,\overline{f}^{}(w))`$ is $`2n1`$ (at the origin). From this, we see that the formal map $`(^{3n3},0)(z^{},\xi ,\eta )(v_3(z^{},\xi ,\eta ),(\overline{f}^{}\overline{v}_2)(\xi ,\eta ))`$ has rank $`2n1`$. (The rank of such a formal map is its rank in the quotient field of $`[[z^{},\xi ,\eta ]]`$.) Since $`𝒯^e(\xi ,\eta )`$ agrees up to order $`e`$ with $`(\overline{f}\overline{v}_2)(\xi ,\eta )`$, we obtain that for $`e`$ large enough, the mapping (6.4) is of generic rank $`2n1`$. This completes the proof of Theorem 2.1.∎
###### Remark 3.
When the target hypersurface $`M^{}`$ is given in normal coordinates i.e. $`\overline{\mathrm{\Phi }}^{}(\omega ,0)\omega _n`$, then the normal component $`f_n`$ of a formal biholomorphism $`f:(M,0)(M^{},0)`$ is convergent provided that the source hypersurface $`M`$ is minimal. Indeed, this follows by taking $`\lambda =0`$ in Theorem 2.1.
###### Proof of Theorem 2.2..
By the Taylor expansion (5.3) and by Theorem 2.1, we obtain that all the $`\varphi _\alpha ^{}f`$ are convergent in a common neighborhood $`U`$ of $`0^n`$. Since $`M^{}`$ is holomorphically nondegenerate, by , there exists $`\varphi _{\beta ^1}^{}(\omega ),\mathrm{},\varphi _{\beta ^n}^{}(\omega )`$, $`\beta ^i^{n1}`$, $`i=1,\mathrm{},n`$, such that
$$\mathrm{det}\left[\frac{\varphi _{\beta ^i}^{}}{\omega _j}(\omega )\right]_{1i,jn}0.$$
Since $`f`$ is a formal biholomorphism, this implies that
(6.5)
$$\mathrm{det}\left[\frac{\varphi _{\beta ^i}^{}}{\omega _j}(f(z))\right]_{1i,jn}0,$$
as a formal power series in $`z`$. Put $`\psi _i(z):=(\varphi _{\beta ^i}^{}f)(z)`$ and $`R_i(z,\omega ):=\varphi _{\beta ^i}^{}(\omega )\psi _i(z)`$, $`i=1,\mathrm{},n`$. Observe that since $`\psi _i(z)`$ is convergent, $`R_i(z,\omega )\{z,\omega \}`$ for $`i=1,\mathrm{},n`$. Moreover, since $`R_i(z,f(z))=0`$, $`i=1,\mathrm{},n`$, in $`[[z]]`$, by (6.5), we may apply Proposition 4.2 to conclude that $`f`$ is convergent.∎
## 7. Transcendence degree and partial convergence of formal maps
In this last section, we want to indicate how Theorem 2.1 can be viewed as a result of partial convergence for formal biholomorphic mappings of real analytic hypersurfaces. Before explaining what we mean by this, we need to recall the following. If $`M`$ is a real analytic hypersurface in $`^n`$ and $`pM`$, let $`𝕂(p)`$ be the quotient field of $`\{zp\}`$, and $`H(M,p)`$ be the vector space over $`𝕂(p)`$ consisting of the germs at $`p`$ of (1,0) vector fields, with meromorphic coefficients, tangent to $`M`$ (near $`p`$). We then define the degeneracy of $`M`$ at $`p`$, denoted $`d(M,p)`$, to be the dimension of $`H(M,p)`$ over $`𝕂(p)`$. It is shown in that the mapping $`Mpd(M,p)\{0,\mathrm{},n\}`$ is constant on any connected component of $`M`$. Thus, if $`M`$ is a connected real analytic hypersurface, one can define its degeneracy $`d(M)`$ to be the degeneracy $`d(M,q)`$ at any point $`qM`$. Observe that the germ $`(M,p)`$, $`pM`$, is holomorphically nondegenerate if and only if $`d(M)=d(M,p)=0`$.
Theorem 2.1 gives the following result of partial convergence for formal biholomorphic mappings of real analytic hypersurfaces. By this, we mean that we have the following.
###### Theorem 7.1.
Let $`f:(M,0)(M^{},0)`$ be a formal biholomorphism between two germs at 0 of smooth real real-analytic hypersurfaces in $`^n`$. Assume that $`M`$ is minimal at $`0`$ and let $`d(M^{})`$ be the degeneracy of the germ $`(M^{},0)`$. Then, there exists $`g(\omega )=(g_1(\omega ),\mathrm{},g_{nd(M^{})}(\omega ))(\{\omega \})^{nd(M^{})}`$, $`\omega ^n`$, of generic rank $`nd(M^{})`$ such that the mapping $`gf`$ is convergent.
###### Proof.
We again use the notations of §2. As in the proof of Theorem 2.2, we have, using the expansion (5.3),
$$\overline{\mathrm{\Phi }}^{}(f(z),\lambda )=\underset{\beta ^{n1}}{}(\varphi _\beta ^{}f)(z)\lambda ^\beta .$$
Thus, we know, by Theorem 2.1, that for any multi-index $`\beta ^{n1}`$, $`(\varphi _\beta ^{}f)(z)`$ is convergent in some neighborhood $`U`$ of $`0`$ in $`^n`$. We choose $`\varphi _{\beta ^1}^{}(\omega ),\mathrm{},\varphi _{\beta ^r}^{}(\omega )`$, $`r=nd(M^{})`$, of generic maximal rank equal to $`nd(M^{})`$ in a neighborhood $`U^{}`$ of $`0`$ in $`^n`$ (see ). Then, if we define $`g_j(\omega )=\varphi _{\beta ^j}^{}(\omega )`$, $`j=1,\mathrm{},nd(M^{})`$, we obtain the desired statement of the Theorem. ∎
###### Remark 4.
One should observe that the convergent power series mapping $`g`$ in Theorem 7.1 is obtained in a constructive way from the target manifold $`M^{}`$. Indeed, this is a consequence of the statement of Theorem 2.1.
Now, we want to make explicit links with the notion of transcendence degree introduced in in the $`𝒞^{\mathrm{}}`$ mapping problem. For this, we first set the corresponding definitions in the formal case.
###### Definition 7.1.
Let $`:(^N,0)(^N^{},0)`$ be a formal (holomorphic) mapping, and $`V`$ be a complex analytic set through the origin in $`^N\times ^N^{}`$. Assume that $`V`$ is given near the origin in $`^{N+N^{}}`$ by
$$V=\{(x,y)^N\times ^N^{}:b_1(x,y)=\mathrm{}=b_q(x,y)=0\},$$
$`b_i(x,y)\{x,y\}`$, $`i=1,\mathrm{},q`$. Then, the graph of $``$ is said to be formally contained in $`V`$ if $`b_1(x,(x))=\mathrm{}=b_q(x,(x))=0`$ in $`[[x]]`$.
It follows from the Nullstellensatz that this definition is independent of the choice of the defining functions $`(b_i)`$ for $`V`$.
###### Definition 7.2.
Let $`:(^N,0)(^N^{},0)`$ be a formal holomorphic mapping. Let $`V_{}`$ be the germ of the complex analytic set through the origin in $`^{N+N^{}}`$ defined as the intersection of all the complex analytic sets through the origin in $`^{N+N^{}}`$ which formally contain the graph of $``$. Then the transcendence degree of $``$ is the nonnegative integer $`\mathrm{dim}_{}V_{}N`$.
This definition is motivated by the following result.
###### Proposition 7.2.
Let $`:(^N,0)(^N^{},0)`$ be a formal holomorphic mapping. Then, the following conditions are equivalent:
i) $``$ is convergent.
ii) The transcendence degree of $``$ is zero.
###### Proof.
The implication i) $``$ ii) is clear. The other implication is equivalent to the following proposition.∎
###### Proposition 7.3.
Let $`:(^N,0)(^N^{},0)`$ be a formal (holomorphic) mapping. If there exists a germ at $`0`$ of a complex analytic set $`V^{N+N^{}}`$ which formally contains the graph of $``$ with $`\mathrm{dim}_{}V=N`$, then $``$ is convergent.
###### Proof.
Let $`V`$ be as in the proposition. We can assume that, near $`0^{N+N^{}}`$,
$$V=\{(x,y)^N\times ^N^{}:b_1(x,y)=\mathrm{}=b_p(x,y)=0\},$$
where each $`b_j(x,y)\{x,y\}`$. To this complex analytic set $`V`$, as is customary, we associate the following ideal of $`\{x,y\}`$ defined by
$$(V)=\{s\{x,y\}:s\mathrm{vanishes}\mathrm{on}V\}.$$
By the Noetherian property, we can assume that $`(V)`$ is generated by a family $`(h_i(x,y))_{i=1,\mathrm{},k}\{x,y\}`$. Furthermore, by the Nullstellensatz , for $`i=1,\mathrm{},k`$, there exists an integer $`\mu _i`$ such that $`h_i^{\mu _i}(x,y)`$ is in the ideal generated by the $`b_j(x,y)`$, $`j=1,\mathrm{},p`$. This implies that for $`i=1,\mathrm{},k`$, $`h_i(x,(x))=0`$ in $`[[x]]`$. The following result, a consequence of the Artin approximation theorem, is contained in (p.63). Assume that the height of the ideal $`(V)`$ (generated by the family $`(h_j(x,y))_{1jk}`$ in $`\{x,y\}`$) is equal to $`N^{}`$. Then, any formal solution $`𝒴(x)([[x]])^N^{}`$, $`𝒴(0)=0`$, of the system $`h_1(x,y)=\mathrm{}=h_k(x,y)=0`$ (in the unknown $`y`$) is convergent. Thus, to obtain the convergence of our original formal power series $``$, it suffices to check that the height of $`(V)`$ is equal to $`N^{}`$. Since $`\{x,y\}`$ is a local regular ring of Krull dimension $`N+N^{}`$, by Proposition 6.12, p.22 of , we have the formula
$$\mathrm{height}((V))+\mathrm{dim}\{x,y\}/(V)=N+N^{},$$
where $`\mathrm{dim}\{x,y\}/(V)`$ is the Krull dimension of the ring $`\{x,y\}/(V)`$. Since the Krull dimension of such a ring coincides with the dimension of the complex analytic set $`V`$ (cf. , p.226, Proposition 1), which is, here, equal to $`N`$, we obtain that the height of $`(V)`$ is $`N^{}`$. This completes the proof of Proposition 7.3, and hence, the proof of Proposition 7.2. ∎
With these tools at our disposal, we can now state a result which follows from Theorem 7.1.
###### Corollary 7.4.
Let $`f:(M,0)(M^{},0)`$ be a formal biholomorphism between two germs at 0 of smooth real real-analytic hypersurfaces in $`^n`$. Assume that $`M`$ is minimal at $`0`$ and denote by $`𝒟_f`$ the transcendence degree of the map $`f`$. Then, $`𝒟_fd(M^{})`$, where $`d(M^{})`$ is the degeneracy of $`M^{}`$. In other words, there exists a complex analytic set of (pure) dimension $`n+d(M^{})`$ which formally contains the graph of $`f`$.
###### Proof.
By Theorem 7.1, there exists a convergent power series mapping $`g(\omega )=(g_1(\omega ),\mathrm{},g_{nd(M^{})}(\omega ))(\{\omega \})^{nd(M^{})}`$ such that for each $`j=1,\mathrm{},r`$, $`\delta _j(z):=(g_jf)(z)`$ is convergent, $`r=nd(M^{})`$. Then, the graph of $`f`$ is formally contained in the complex analytic set
$$A=\{(z,\omega )(^{2n},0):g_1(\omega )\delta _1(z)=\mathrm{}=g_r(\omega )\delta _r(z)=0\}.$$
Let $`A=_{i=1}^k\mathrm{\Gamma }_i`$ be the decomposition of $`A`$ into irreducible components. For any positive integer $`\sigma `$, one can find, according to the Artin approximation theorem , a convergent power series mapping $`f^\sigma (z)(\{z\})^n`$ defined in some small neighborhood $`U^\sigma `$ of $`0`$ in $`^n`$, which agrees with $`f(z)`$ up to order $`\sigma `$ (at 0) and such that the graph of $`f^\sigma `$, denoted $`G(f^\sigma )`$, is contained in $`A`$. Since $`G(f^\sigma )`$ is contained in $`A`$, it must be contained in an irreducible component of $`A`$. Thus, by the pigeonhole principle, at least one subsequence of $`(f^\sigma )_\sigma ^{}`$ is contained in one of such irreducible components, say $`\mathrm{\Gamma }_1`$. There is no loss of generality in assuming that such a subsequence is $`(f^\sigma )_\sigma ^{}`$ itself. We first observe that this implies that the graph of $`f`$ is formally contained in $`\mathrm{\Gamma }_1`$. Moreover, since $`f`$ is a formal biholomorphism, the family $`(f^\sigma )_\sigma ^{}`$ is also a family of local biholomorphisms. In particular, this implies that the generic rank of the family of holomorphic functions
$$((g_if^1)(z))_{1ir},$$
is $`r`$. As a consequence, if $`z_0`$ is close enough to 0 in $`^n`$ and is chosen so that the rank of the preceding family at $`z_0`$ equals $`r`$, the implicit function theorem shows that $`A`$ is an $`n+d(M^{})`$-dimensional complex submanifold near $`(z_0,f^1(z_0))\mathrm{\Gamma }_1`$. Since $`\mathrm{\Gamma }_1`$ is irreducible, it is pure-dimensional; thus $`\mathrm{\Gamma }_1`$ is an $`n+d(M^{})`$ pure-dimensional complex analytic set formally containing the graph of $`f`$. By definition of the transcendence degree, this implies that $`𝒟_fd(M^{})`$.∎
The following example illustrates the applications of Theorem 7.1 and Corollary 7.4.
###### Example 1.
Let $`M=M^{}`$ be the minimal real algebraic hypersurface through the origin in $`^3`$ given by
$$\mathrm{Im}z_3=|z_1z_2|^2.$$
Here, $`M`$ is holomorphically degenerate and its degeneracy $`d(M)`$ is equal to 1. Consider the following formal biholomorphic self-map of $`M`$:
$$f_h:^3(z_1,z_2,z_3)(z_1e^{h(z)},z_2e^{h(z)},z_3)^3,$$
where $`h(z)=h(z_1,z_2,z_3)`$ is any non-convergent formal power series vanishing at the origin. Observe that Theorem 2.1 gives in this example that for any formal biholomorphic self-map $`f=(f_1,f_2,f_3)`$ of $`M`$, the product $`f_1f_2`$ and the third component $`f_3`$ are necessarily convergent. Observe in this example that the first two components of $`f_h`$ are not convergent, but that the transcendence degree of the map $`f_h`$ is actually $`1=d(M)`$. Indeed, the graph of $`f_h`$ is formally contained in the complex analytic set of dimension 4
$$V=\{(z,\omega )^3\times ^3:\omega _1\omega _2=z_1z_2,\omega _3=z_3\},$$
and cannot be formally contained in a complex analytic set of dimension 3, since otherwise $`f_h`$ would be convergent by Proposition 7.3.
We conclude by observing that Theorem 2.2 can be regarded as a direct consequence of Corollary 7.4. Indeed, when $`M^{}`$ is holomorphically nondegenerate, as mentioned above, $`d(M^{})=0`$ and hence $`𝒟_f=0`$ by Corollary 7.4. It then follows from Proposition 7.2 that $`f`$ is convergent in that case. We should also mention that Theorem 2.2, Theorem 7.1 and Corollary 7.4 are all consequences of our main result, namely Theorem 2.1.
## Acknowledgments
I would like to thank Makhlouf Derridj for his interest in this work and Vincent Thilliez for calling my attention to the paper . I wish also to address special thanks to Salah Baouendi and Linda Rothschild for many simplifying remarks and suggestions. |
warning/0002/hep-ph0002275.html | ar5iv | text | # Continuous Emission versus Freeze-out via HBT
## I Introduction
When describing ultra-relativistic heavy-ion collisions with hydrodynamic models, a simple picture has been extensively adopted. It is usually considered that, as the thermalized matter expands, the system gradually cools down and, when the temperature reaches a certain freeze-out value $`T_f`$, it decouples. Every observed quantity is then computed on the hypersurface $`T=T_f`$. For instance, the momentum distribution of the produced hadrons are obtained by using the Cooper-Frye integral extended over this hypersurface. Though operationally simple, such a zero-thickness freeze-out hypersurface is clearly a highly idealized concept when applied to finite-volume and finite-lifetime systems as those formed in high-energy heavy-ion collisions.
More recently, Grassi, Hama and Kodama proposed an alternative picture to the particle emission: instead of being emitted only when crossing the sharply defined freeze-out surface, they considered that the process could occur continuously. Being so, in this picture, particles could be emitted from the whole expanding volume of the system, at different temperatures, and not only from the surface with constant $`T=T_f`$. As a consequence, in the continuous emission model (CEM), the observed quantities depend on the whole history of the expanding system and not only on the instant of the freeze-out. Concretely, it has been shown that i) CEM enhances the large$`m_T`$ component of the heavy-particle ($`p,\mathrm{\Lambda },\mathrm{\Xi },\mathrm{\Omega },\mathrm{}`$) $`m_T`$ spectra, ii) it gives a concave shape for the pion $`m_T`$ spectrum even without considering transverse expansion of the fluid, iii) it can lead to the correct hyperon production ratios and spectrum shapes with conceptually reasonable choice of parameters, and iv) it reproduces the observed mass dependence of the slope parameter $`T`$ .
Naturally, we would like to further explore if the above model would present striking differences when compared to the usual sudden freeze-out picture. One expectation would be that the space-time region from which the particles were emitted would be quite different in both scenarios. In the continuous emission picture the duration of the emission processes is expected to be longer than in the freeze-out scenario, which should considerably affect the behavior of the correlation function. Previous studies have indeed shown that the influence of the emission time on the apparent transverse source dimensions were remarkably strong. It was also shown in Ref. that a prolonged freeze-out would considerably distort the two-particle correlation function. Our main object in the present work is to show the differences in two-pion correlation predicted by CEM, as compared with the results obtained under the usual assumption of sharp freeze-out. For this purpose, we will adopt the same approximations used in Ref. , namely, one dimensional Bjorken model for massless-pion gas. It turns out that, within these approximations, the HBT effect suffers a large deformation when the usual freezeout scenario is replaced by CEM, affecting substantially the conclusions achieved on the properties of the matter formed in high-energy collisions.
## II Continuous emission of particles
In CEM, it is assumed that, at each space-time point $`x^\mu `$, each particle has a certain probability of not colliding any more, due to the finite dimensions and lifetime of the thermalized matter. Then, the distribution function $`f(x,p)`$ of the expanding system has two components, one representing the portion of the fluid already free and another corresponding to the part still interacting, i.e.,
$`f(x,p)=f_{free}(x,p)+f_{int}(x,p).`$
We may write the portion of free particles as a fraction of the total distribution function, as follows
$$f_{free}(x,p)=𝒫f(x,p)=\frac{𝒫}{1𝒫}f_{int}(x,p).$$
(1)
Let us assume, as in the previous papers, that the fraction still interacting is represented by a thermal distribution function
$$f_{int}(x,p)f_{th}(x,p)=\frac{g}{(2\pi )^3}\frac{1}{\mathrm{exp}[p.u(x)/T(x)]\pm 1},$$
(2)
where $`u^\mu `$ is the fluid velocity at $`x^\mu `$ and $`T`$ is its temperature at that point. The factor $`𝒫`$ can be alternatively understood as the probability that a particle with momentum $`p^\mu `$ escapes from $`x^\mu `$ without further collisions.
If we assume that the fluid is confined to a cylinder of radius $`R_T`$ , the fraction $`𝒫`$ of free particles at each space-time point $`x^\mu `$ may be computed by using the Glauber formula
$$𝒫=\mathrm{exp}\left(_t^{t_{out}}n(x^{})\sigma v_{rel}𝑑t^{}\right),$$
(3)
where
$$t_{out}=t+(\rho \mathrm{cos}\varphi +\sqrt{R_T^2\rho ^2\mathrm{sin}^2\varphi })/(v\mathrm{sin}\theta )$$
(4)
is the time when the particle with velocity $`\stackrel{}{v}=(v\mathrm{sin}\theta \mathrm{cos}\varphi ,v\mathrm{sin}\theta \mathrm{sin}\varphi ,v\mathrm{cos}\theta )`$ reaches the surface of the fluid at $`\rho =R_T`$.
If we further consider that, initially, the energy density is approximately constant (i.e., $`ϵ=\frac{\pi ^2}{10}T_o^4`$ for all the points with $`\rho R_T`$ and zero for $`\rho >R_T`$), we can calculate the probability $`𝒫`$ analytically, resulting in
$$𝒫=(\tau /\tau _{out})^a;a3\frac{1.202}{\pi ^2}T_0^3\tau \sigma v_{rel},$$
(5)
where $`v_{rel}1`$. The previous results can be found in Ref..
## III HBT interferometry
The second-order interferometry of identical particles, also known as HBT effect is a powerful tool for probing geometrical sizes of the space-time zone from which they were emitted, as well as for testing dynamical correlations built in during the system evolution.
In its idealized version, the two-pion interferometry could be studied through the so-called two-particle correlation function
$`C_2(k_1,k_2)={\displaystyle \frac{P_2(k_1,k_2)}{P_1(k_1)P_1(k_2)}}=1+|\rho (k_1k_2)|^2,`$ (6)
where $`P_1(k_i)`$ and $`P_2(k_1,k_2)`$ are, respectively, the single-particle inclusive distribution and the joint probability for detecting two pions; $`\rho (k_1k_2)`$ is the Fourier transform of the source space-time distribution.
In realistic cases, however, it is mandatory to employ more general formalisms, as is the case of the Covariant Current Ensemble, flexible enough to include phase-space correlations resulting from the underlying dynamics. As a consequence, the HBT correlation functions would reflect a model dependent analysis. In the Covariant Current Ensemble formalism, the correlation function can be expressed as
$$C(k_1,k_2)=C(q,K)=1+\frac{|G(q,K)|^2}{G(k_1,k_1)G(k_2,k_2)},$$
(7)
where $`q^\mu =k_1^\mu k_2^\mu `$ and $`K^\mu =\frac{1}{2}(k_1^\mu +k_2^\mu )`$ and the complex amplitude, $`G(k_1,k_2)`$, can be written as
$$G(k_1,k_2)=d^4xd^4pe^{iq^\mu x_\mu }D(x,p)j_0^{}(u_f^\mu k_{1\mu })j_0(u_f^\mu k_{2\mu }),$$
(8)
where $`D(x,p)`$ is the break-up phase-space distribution and the currents, $`j_0(u_f.k_i)`$, contain information about the production dynamics. If one takes $`k_1=k_2`$ in Eq. (8), one obtains
$$G(k_i,k_i)=d^4xd^4pD(x,p)|j_0(u_f^\mu k_{i\mu })|^2,$$
(9)
which coincides with the one-particle spectrum.
As discussed in Ref., the currents $`j_0(u_f.k_i)`$ in Eqs. (8,9) can be associated to thermal models and written covariantly as $`j_0(k)\sqrt{u^\mu k_\mu }\mathrm{exp}\{u^\mu k_\mu /(2T)\}`$. However, to make the computation easier, we shall adopt throughout the paper a more convenient parametrization
$$j_0(u.k)=\mathrm{exp}\{\frac{u^\mu k_\mu }{2T_{ps}}\},$$
(10)
where, in the case of pions, the so-called pseudo temperature $`T_{ps}`$ was related with the true temperature $`T`$ by the equation
$$T_{ps}(x)=1.42T(x)12.7\text{ MeV }.$$
(11)
This mapping between $`T(x)`$ and $`T_{ps}(x)`$ was later shown to be a good approximation also in the case of kaon interferometry.
### A Bjorken model with sudden freezeout
In the ideal one dimensional Bjorken picture, using the above pseudo-thermal parameterization for the currents, an analytical form for the amplitudes can be derived
$$G(k_1,k_2)=2<\frac{dN}{dy}>\{\frac{2}{q_TR_T}J_1(q_TR_T)\}K_0(\xi ),$$
(12)
where
$`\xi ^2`$ $`=`$ $`[{\displaystyle \frac{1}{2T}}(m_{1T}+m_{2T})i\tau (m_{1T}m_{2T})]^2+`$ (14)
$`2({\displaystyle \frac{1}{4T^2}}+\tau ^2)m_{1T}m_{2T}[\mathrm{cosh}(\mathrm{\Delta }y)1],`$
$`\mathrm{\Delta }y=y_1y_2`$ and $`<>`$ indicates average over particles 1 and 2.
The single-inclusive distribution is then written as
$$G(k_i,k_i)=E\frac{d^3N}{dk_i^3}=2\frac{dN}{dy_i}K_0(\frac{m_{iT}}{T}).$$
(15)
### B Bjorken model with continuous emission
The initial expectation concerning the differences between the continuous emission versus the freeze-out scenarios were mainly focused on the different emission periods. Naturally, in the continuous emission picture the duration of the emission processes is longer than in the freeze-out scenario, which should considerably affect the behavior of the correlation function.The reason for this comes from previous studies which have shown that the influence of the emission time on the transverse source dimensions were remarkably strong.
For treating pion interferometry in the case that interests us, we consider a different but equivalent form for expressing the amplitudes in Eq. (7). The single-inclusive distribution is written as in Ref.
$$G(k_i,k_i)=d^4x𝒟_\mu \left[k_i^\mu f_{free}\right].$$
(16)
Analogously, the two-particle complex amplitude is written, instead of Eq. (8), as
$$G(k_1,k_2)=d^4xe^{iqx}\{𝒟_\mu \left[k_1^\mu f_{free}\right]\}^{\frac{1}{2}}\{𝒟_\mu \left[k_2^\mu f_{free}\right]\}^{\frac{1}{2}}.$$
(17)
In Eq.(16) and (17), $`𝒟_\mu `$ is the generalized divergence operator, which, due to the symmetry of the problem, is written in Bjorken+transverse polar coordinates.
In order to proceed further, let us recall that usually we are interested in small momentum differences $`q^\mu =k_1^\mu k_2^\mu `$, as compared with the average momentum of the pair, $`K^\mu =\frac{1}{2}(k_1^\mu +k_2^\mu )`$. If we then approximate $`k_i^\mu K^\mu `$ in Eq.(17), a substantial simplification is achieved and it could then be written as
$$G(q,K)G(k_1,k_2)=d^4xe^{iq^\nu x_\nu }𝒟_\mu \left[K^\mu f_{free}\right].$$
(18)
We should note that such a dependence on $`K^\mu `$, replacing the individual momenta $`k_1^\mu `$ and $`k_2^\mu `$ in the complex amplitude of Eq.(17), is also present in the general derivations based on the Wigner formalism.
In principle, the integral in Eq. (18) should be extended over the whole space-time with $`\tau >\tau _0`$ . However, due to the finite size and lifetime of our system, the integrand is expected to quickly vanish where the assumption embodied by Eq. (2) also breaks down. So, in computing this integral, we separated the space-time in two regions, one where $`𝒫>𝒫_{}`$ and the other with $`𝒫𝒫_{}`$, with some reasonable value of $`𝒫_{}`$. Upon partial integration, the latter is reduced to the surface contribution and the former may be estimated by using the Cooper-Frye formula on the surface $`𝒫=𝒫_{}`$, applied to the interacting component. We emphasize, however, that $`𝒫`$ is a momentum-dependent quantity, so this is not the usual Cooper-Frye integral. After some manipulation, we get for the single-inclusive distribution
$`G(k_i,k_i)`$ $`=`$ $`{\displaystyle \frac{1}{(2\pi )^3(1𝒫_{})}}{\displaystyle _0^{2\pi }}𝑑\varphi {\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑\eta `$ (19)
$`\times `$ $`\{{\displaystyle _0^{R_T}}\rho d\rho \tau _{}m_{iT}\mathrm{cosh}(y_i\eta )`$ (20)
$`+`$ $`{\displaystyle _{\tau _0}^+\mathrm{}}\tau d\tau \rho _{}k_{iT}\mathrm{cos}\varphi \}e^{m_{iT}\mathrm{cosh}(y_i\eta )/T_{ps}(x)}.`$ (21)
Analogously, instead of Eq. (12), the two-particle complex amplitude is now written as
$`G(q,K)`$ $`=`$ $`{\displaystyle \frac{1}{(2\pi )^3(1𝒫_{})}}{\displaystyle _0^{2\pi }}𝑑\varphi {\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑\eta `$ (23)
$`\times `$ $`\{{\displaystyle _0^{R_T}}\rho d\rho \tau _{}M_T\mathrm{cosh}(Y\eta )`$ (25)
$`\times e^{i[\tau _{}(q_0\mathrm{cosh}\eta q_L\mathrm{sinh}\eta )\rho q_T\mathrm{cos}(\varphi \varphi _q)]}`$
$`+`$ $`{\displaystyle _{\tau _0}^+\mathrm{}}\tau 𝑑\tau \rho _{}K_T\mathrm{cos}\varphi `$ (27)
$`\times e^{i[\tau (q_0\mathrm{cosh}\eta q_L\mathrm{sinh}\eta )\rho _{}q_T\mathrm{cos}(\varphi \varphi _q)]}\}`$
$`\times `$ $`e^{M_T\mathrm{cosh}(Y\eta )/T_{ps}(x)},`$ (28)
where
$`M_T=\sqrt{K_T^2+M^2},\stackrel{}{K}_T=\frac{1}{2}(\stackrel{}{k}_1+\stackrel{}{k}_2)_T,M^2=K_\mu K^\mu =m^2\frac{1}{4}q_\mu q^\mu `$, $`Y`$ is the rapidity corresponding to $`\stackrel{}{K}`$, $`\varphi `$ is the azimuthal angle with respect to the direction of $`\stackrel{}{K}`$, and $`\varphi _q`$ is the angle between the directions of $`\stackrel{}{q}`$ and $`\stackrel{}{K}`$.
In Eqs. (LABEL:giicontemis) and (28), $`\tau _{}`$ and $`\rho _{}`$ are the limiting values corresponding to a certain value of the escape probability $`𝒫_{}`$, i.e.,
$`\tau _{}`$ $`=`$ $`{\displaystyle \frac{\rho \mathrm{cos}\varphi +\sqrt{R_T^2\rho ^2\mathrm{sin}^2\varphi }}{(k_T/E)\mathrm{cosh}y\left[\sqrt{\mathrm{sinh}^2(\eta y)+𝒫_{}^{2/a}}\mathrm{cosh}(\eta y)\right]}}.`$ (29)
and
$`\rho _{}`$ $`=`$ $`\tau {\displaystyle \frac{k_T}{E}}\mathrm{cosh}y\mathrm{cos}\varphi \left[\sqrt{\mathrm{sinh}^2(\eta y)+𝒫_{}^{2/a}}\mathrm{cosh}(\eta y)\right]`$ (31)
$`\pm `$ $`[R_T^2\tau ^2({\displaystyle \frac{k_T}{E}})^2\mathrm{cosh}^2y\mathrm{sin}^2\varphi [\sqrt{\mathrm{sinh}^2(\eta y)+𝒫_{}^{2/a}}`$ (33)
$`\mathrm{cosh}(\eta y)]^2]^{1/2}.`$
For choosing the value of $`𝒫_{}`$, in principle, we would like to take $`𝒫_{}=1`$, corresponding to the complete integration of (16) and (18). However, we should notice that the expressions (LABEL:giicontemis) and (28) above become indeterminate in the limit $`𝒫_{}1`$. As mentioned above, the thermal assumption for $`f_{int}(x,p)`$ breaks down in the same limit. For this reason, already in Ref., it was chosen $`𝒫_{}=0.5`$ and the effect of changing this value was discussed. We shall adopt the same value here.
## IV Comparison of Results
### A Ideal Configurations
The complexity of the expressions for the amplitudes appearing in Eq. (7) in the continuous-emission scenario is evident from Eq.(LABEL:giicontemis) and (28) above. In order to get some insight regarding the differences of the correlation functions in the two scenarios under investigation, let us select some special kinematical zones, corresponding to an idealized situation in which high precision data with unlimited statistics would be available. For instance, let us fix $`y_i=0`$ ($`K_L=q_L=0`$), so that $`\theta =\pi /2`$ with respect to the collision axis. Due to the symmetry of the problem we can, without any loss of generality, choose $`\stackrel{}{K}`$ along the $`x`$-axis. We then explore the behavior of $`C(q_T,K_T)`$ for fixed $`K_T`$. For restricting even more our kinematical window, let us consider two cases. Case I (or Zone I) corresponds to considering $`\varphi _{p_1}=\varphi _{p_2}=\varphi _p`$ (the two pions are symmetrically emitted around $`\stackrel{}{K}`$), implying that $`\varphi _q=\pi /2`$; in this case, we can write the individual momenta as $`k_i^\mu =(\sqrt{m_\pi ^2+K_T^2+q_T^2/4},K_T,\pm q_T/2,\mathrm{\hspace{0.33em}0})`$, where the $`\pm `$ signs correspond to pion 1 and 2, respectively. The momentum difference $`\stackrel{}{q}_T=\stackrel{}{q}_S`$ in this situation corresponds to the so-called sidewards component introduced in Ref. For comparison, we consider that in the usual freezeout scenario, the decoupling occurs at $`T_{fo}=170`$ MeV. The other constant values assumed in the calculation that follows were:
| $`T_0`$ | $`\tau _0`$ | $`<\sigma v_{rel}>`$ | $`R_T`$ | $`m_\pi `$ |
| --- | --- | --- | --- | --- |
| (MeV) | (fm/c) | (fm<sup>2</sup>) | (fm) | (MeV) |
| 200 | 1 | 2 | 3.7 ($``$ S) | 140 |
Results corresponding to the ZONE I above are shown in Fig. 1. As expected, since we are neglecting the transverse expansion, the difference between the predictions of the two scenarios is small. Slightly broader correlation function for the CEM case, and the decreasing width with $`K_T`$, was also expected.
Case II (or Zone II) correponds to considering $`\varphi _q=0`$, $`|\stackrel{}{k}_1|>|\stackrel{}{k}_2|`$, with both $`\stackrel{}{k}_i`$ along the $`x`$-axis, i.e., $`\stackrel{}{k}_i\stackrel{}{K}_T\stackrel{}{q}_T`$; in this case, we can write the individual momenta as $`k_i^\mu =(\sqrt{m_\pi ^2+[K_T\pm q_T/2]^2},K_T\pm q_T/2,0,0)`$, where again the $`\pm `$ signs correspond to pion 1 and 2, respectively. The momentum difference $`\stackrel{}{q}_T=\stackrel{}{q}_O`$ in this situation corresponds to the so-called outwards compo-
nent introduced in Ref. Results corresponding to the ZONE II above are shown in Fig. 2. This is the case, mentioned in the Introduction, where the duration of the emission process becomes essential. Since the emission time in the usual freezeout does not depend crucially on the particle momentum, the correlation function is almost independent of $`K_T`$. On the contrary, the emission time is strongly momentum dependent in
CEM, for large-momentum particles are emitted mainly at early times, whereas small-momentum particles may be emitted also at later stage of expansion when the fluid is cooler and the system larger. So, we see in Fig. 2 a significant $`K_T`$ dependence, being the correlation narrower for smaller $`K_T`$. Also, we can see that both the CEM curves are narrower than the corresponding freezeout curves, indicating that the emission time in CEM is longer in general. Looking more carefully at the curves, we can also perceive that the tail of $`C(q_T,K_T)`$ in CEM is much flatter than in sharp freezeout. Probably this flat tail is due to the small source depth in CEM at early times. It is clear that in Bjorken model without transverse expansion, which we used in the present work, the source depth is constant and $`R_T`$ in sharp freezeout scenario.
Figure 3 represents the same situation as in Zone II but with a different freezeout temperature, $`T_{fo}=140`$ MeV, in order to show the sensitivity of the results for a lower $`T_{fo}`$ in the case of the usual freezeout. As seen, the $`T_{fo}`$ dependence has shown to be very weak, as expected.
### B Averaged Correlations
Although the selective kinematical zones could teach us interesting points concerning the differences of the behavior of the correlation functions corresponding to both scenarios, as shown in Fig. 1-3, such conditions correspond to an idealization. For putting the calculations into more realistic grounds, averages over the angles, momenta, and the unobserved projections of the momentum differences $`\stackrel{}{q}`$ should be performed. Using the azimuthal symmetry of the problem we can still select $`\stackrel{}{K}_T`$ along the $`x`$-axis, such that $`\stackrel{}{K}=(K_T,0,K_L)`$. Then, averaging over different kinematical zones or windows would correspond to integrating over $`\stackrel{}{K}`$ and $`\stackrel{}{q}`$ (except over the plotting component of $`\stackrel{}{q}`$). In order to make the analysis roughly compatible with the range covered by NA35 S+A collisions, we considered the kinematical variables in the following intervals: $`0.5y0.5`$ (or, equivalently, $`180K_L180`$ MeV); $`50K_T600`$ MeV; $`0(q_L,q_S,q_{out})30`$ MeV (corresponding to the first experimental bin). As an illustration, we show below an example about how to compute the average:
$`C(q_L)=1+`$ (34)
$`{\displaystyle \frac{_{180}^{180}𝑑K_L_{50}^{600}𝑑K_T_0^{30}𝑑q_S_0^{30}𝑑q_oC(K,q)|G(K,q)|^2}{_{180}^{180}𝑑K_L_{50}^{600}𝑑K_T_0^{30}𝑑q_S_0^{30}𝑑q_oC(K,q)G(k_1,k_1)G(k_2,k_2)}}.`$ (35)
The results are presented in the following way. First, as done in the preceding subsection, in order to stress the differences of results predicted by the two scenarios under study, we start from the same initial temperature $`T_0=200`$ MeV for both the usual freezeout and CEM. The results are shown in Figs. 4-6, respectively as functions of $`q_L`$, $`q_O`$ and $`q_S`$. One sees in Fig. 4 that, as is well known, the $`q_L`$ dependence is very sensitive to the freeze-out temperature $`T_{fo}`$ and if the same initial temperature is attained in both scenarios, the correlation function cor-
responding to the continuous emission picture is closer to the one referring to the thermal freezeout at lower $`T_{fo}`$. However, the shapes are not the same. The one related to CEM is more peaked at the small-$`q_L`$ values, becoming flatter in the tail region. This is in clear contrast to those corresponding to sharp freezeout scenario which are more similar to Gaussians. Physically, this behavior of CEM curve could be interpreted as exhibiting the history of the matter in expansion, because particles are emitted during the whole evolution in CEM. Namely, the tail of $``$C$``$ depends essentially on early times, when the size of the fluid is small and its temperature high, whereas the peak reflects later times, when the fluid has fully expanded and cooled down.
We have already seen in the preceding subsection that, when plotted as function of $`q_O`$, the correlation function in CEM is significantly narrower than the one in the sharp freezeout and has a flatter tail. These features are again seen in Fig. 5, where $`C_{\pi \pi }`$ with lower freeze-out temperature $`T_{fo}`$ is closer to the curve for CEM, as in Fig. 4. However, if the same initial temperature is attainded in both scenarios, very low $`T_{fo}`$ is necessary, in this case, in order to approximately reproduce the same correlation predicted by CEM. We can also notice that the depletion of the correlation function at small $`q_O`$ values is more dramatic in CEM. In any case, it is important to emphasize that our source is totally chaotic in both scenarios and, as is well known, $`C<2`$ at $`q_L=0`$ is originated only from the averaging processes, since Coulomb final state interactions, as well as the effect of resonances decaying into $`\pi `$’s were not considered here.
We can again notice in Fig. 6 that, also as function of
$`q_S`$, the depletion of $``$ C $`_{\pi \pi }`$ is more pronounced in CEM as compared to the curves corresponding to the sharp freezeout case. As happened in the previous cases, for the same initial temperature, $`C_{\pi \pi }`$ with lower $`T_{fo}`$ is closer to the curve for CEM, but the shape is somewhat different.
In the previous figures, we have shown and discussed the differences of CEM correlation function, confronted with the usual abrupt freezeout one, when the fluid started from the same initial conditions. Now, when analyzing the experimental data, the parameters are usually adjusted by fitting the data points as close as possible, and the conclusions are extracted from the adjusted parameters. In the present model calculations, the only parameter, besides the freezeout temperature $`T_{fo}`$, is the initial temperature $`T_0`$. For computing the results presented in Figs. 7-9, we have fixed the initial temperature for CEM as 200 MeV, and varied the initial temperature $`T_0`$ for the usual freezeout scenario, trying to get the same (or similar) result. To doing so, we have chosen the freezeout temperature as $`T_{fo}=140`$ MeV, as often done in hydrodynamic calculations and following the indications of the previous discussions.
As seen, especially in Fig. 7, the correlation functions predicted by the two different scenarios were so different in shape that it was not always possible to obtain similar curves. For example, in the case of $`q_L`$ dependence, Fig. 7, $``$C$``$ for CEM is closer to the curve with higher $`T_0`$ at small $`q_L`$, but at high $`q_L`$, it turns to be closer to the one with the lowest value of $`T_0`$. As discussed above in connection with Fig. 4, the correlation curve in CEM could be interpreted as showing the history of the hot matter in expansion, i.e., the tail of $``$C$``$ reflects essen-
tially the early times, when the size of the fluid is small and its temperature high, and the peak the later times, when the fluid has fully expanded and cooled down. In the usual freezeout picture with a fixed $`T_{fo}`$, small $`T_0`$ is enough to produce a large tail, whereas a larger expansion, so higher $`T_0`$, is required to produce a narrower $``$C$``$.
In Fig. 8, where the $`q_O`$ dependence of $`C_{\pi \pi }`$ is shown, we can again observe a distortion introduced by CEM into the shape of the correlation function. We see that $`C`$ for CEM is closer to the curve with $`T_0=230`$ MeV
at small $`q_O`$, but at high $`q_O`$, it turns to be closer to the one corresponding to $`T_0=260`$ MeV. However, differently from the previous case, Fig. 7, the shape of the freezeout correlation curves is only slightly dependent on $`T_0`$ and it becomes narrower as the temperature increases. What is clearly $`T_0`$ dependent here is the intercept at the origin, which reflects the large expansion dependence of $`C_{\pi \pi }`$ x $`q_L`$ as shown in Figs. 4 and 7.
Finally, Fig. 9 shows that in this case it is possible to find an appropriate $`T_0`$ for sharp freezeout to reproduce the CEM curve. We see that $``$ C $``$ for CEM is closer to the curve with $`T_0=230`$ MeV and the agreement is good for most of the $`q_S`$ region where the interferometric signal is present. This was expected because in the present study we neglected the transverse expansion, so the transverse size is the same in both the scenarios. The initial temperature $`T_0`$ in freezeout is higher than the one for CEM, because this is required to make the size of the fluid large enough and the correlation in the longitudinal direction sharp enough to decrease the intercept on averaging.
From the above results, mainly from the correlation curves as function of $`q_L`$, we clearly see deviations from the pure Gaussian behavior in cases where the continuous emission ansatz was assumed. This could actually reflect a signature of a continuous process in the particle emission.
## V Discussions and Concluding Remarks
As mentioned in the Introduction, treating the decoupling process in heavy-ion collisions as occurring on a sharply defined surface is an operationally simple but highly idealized description. If the consideration of a finite thickness of such a decoupling region does not bring any noticeable difference in the observable quantities, such an approximation would be unquestionable. However, previous studies have shown that several quantities, such as transverse spectra of produced particles and heavy-particle production ratios are sensitive to more involved description of the process, called continuous emission model.
In this paper, we concentrated on the two-pion interferometry, which has extensively been used as a powerful tool for extracting the space-time geometry, as well as probing the underlying dynamics of the hadronic matter formed in heavy-ion collisions, and studied the differences introduced by CEM in confront with the usual sudden freezeout. As shown in Sec. IV, also the HBT effect suffers a large deformation when the usual freezeout is replaced by CEM. This means that conclusions achieved on the properties of the matter formed in high-energy collisions may differ substantially if we adopt one or the other scenario studied here.
For the sake of conceptual clarity and, evidently, also to simplify the computation, we have adopted in this work a simplified one-dimensional Bjorken model for massless pion fluid (the pion mass has been included only to computing observable quantities), without phase transition. Nevertheless, one general result emerges, which seems to be evident especially by looking at Figs. 7-9. Namely, if we describe the same data by using CEM or sharp freeze-out (with $`T_{fo}=140`$ MeV), the initial temperature $`T_0`$ required in CEM is lower than in the usual freeze-out. If $`T_{fo}`$ is higher, the difference in $`T_0`$ becomes even larger, which is clear from Figs. 4-6. This result means that if CEM is the correct description of the decoupling process, then it is harder to reach the quark-gluon plasma phase than it appears in the usually adopted sharp freezeout scenario.
Since we have worked with a simplified model, we did not attempt to make any comparison with data. For doing this, evidently we have to do some (or all) of the following improvements. More realistic equation of state (probably including phase transition) should be used; finite longitudinal extension of the fluid should be considered; transverse expansion should be included; resonance formation should be taken into account. All these modifications require some hydrodynamic numerical code with CEM incorporated. We are now working in this direction.
Acknowledgement
This work was partially supported by FAPESP (contract nos. 1998/02249-4 and 1998/14990-0) and by CNPq (contract no. 300054/92-0). We express our gratitude to T. Kodama and T. Csörgő for stimulating discussion on the results. |
warning/0002/math0002128.html | ar5iv | text | # On Cotriangular Hopf Algebras
## 1 Introduction
In \[EG1, Theorem 2.1\] we proved that any triangular semisimple Hopf algebra over an algebraically closed field $`k`$ of characteristic $`0`$ is obtained from the group algebra $`k[G]`$ of a finite group $`G,`$ by twisting its comultiplication by a twist in the sense of Drinfeld \[Dr\]. Since semisimple Hopf algebras are finite-dimensional, dualizing yields that any cotriangular semisimple Hopf algebra over $`k`$ is obtained from $`k[G]^{},`$ the function algebra on $`G,`$ by twisting its multiplication by a Hopf $`2`$cocycle in the sense of Doi \[Do\] (see Section 2 below).
In this paper we generalize Theorem 2.1 from \[EG1\] to not necessarily finite-dimensional cotriangular Hopf algebras $`A`$ over $`k.`$ Namely, our main result (see Section 3 below) is:
Theorem A cotriangular Hopf algebra $`A`$ over $`k`$ is obtained from the function algebra $`𝒪(G)`$ of a pro-algebraic group $`G,`$ by twisting its multiplication by a Hopf $`2`$cocycle, and possibly changing its R-form by a central grouplike element of $`A^{}`$ of order $`2,`$ if and only if $`\text{tr}(S^2|_C)=dim(C)`$ for any finite-dimensional subcoalgebra $`C`$ of $`A`$ (where $`S`$ is the antipode of $`A`$).
The main challenge in the proof of this theorem (see Section 4 below) is to establish the “if” direction. The key step in the proof of the “if” part is to show that our trace condition on $`A`$ guarantees that the categorical dimensions of objects in the category of its finite-dimensional right comodules are non-negative integers (maybe after modifying the R-form). This enables us to apply the same theorem of Deligne on Tannakian categories \[De\] that we applied in the proof of Theorem 2.1 from \[EG1\].
In Section 5, we give examples of twisted function algebras. In particular, we show that in the infinite-dimensional case, the squared antipode for such an algebra may not equal the identity (see Example 5.2 below).
In Section 6, we show that in all of our examples, the operator $`S^2`$ is unipotent on $`A`$, and conjecture it to be the case for any twisted function algebra. We prove this conjecture, using the quantization theory of \[EK1-2\], in a large number of special cases.
In Section 7, we formulate a few open questions.
Throughout the paper, $`k`$ will denote an algebraically closed field of characteristic $`0.`$
Acknowledgements The first author is grateful to Ben Gurion University for its warm hospitality, and to Miriam Cohen and the Dozor Fund for making his visit possible; his work was also supported by the NSF grant DMS-9700477.
The second author is grateful to Susan Montgomery for numerous useful conversations.
The authors would like to acknowledge that this paper was inspired by the work \[BFM\].
## 2 Hopf $`2`$cocycles
Let $`A`$ be a coassociative coalgebra over $`k.`$ For $`aA,`$ we write $`\mathrm{\Delta }(a)=a_1a_2,`$ $`(I\mathrm{\Delta })\mathrm{\Delta }(a)=a_1a_2a_3`$ etc, where $`I`$ denotes the identity map of $`A.`$
Recall that $`A^{}`$ is an associative algebra with product defined by $`(fg)(a)=f(a_1)g(a_2).`$ This product is called the convolution product.
Now let $`(A,m,1,\mathrm{\Delta },\epsilon ,S)`$ be a Hopf algebra over $`k.`$
Recall \[Do\] that a linear form $`J:AAk`$ is called a Hopf $`2`$cocycle for $`A`$ if it has an inverse $`J^1`$ under the convolution product $``$ in $`\text{Hom}_k(AA,k)`$, and satisfies:
$$J(a_1b_1,c)J(a_2,b_2)=J(a,b_1c_1)J(b_2,c_2)\text{and}J(a,1)=\epsilon (a)=J(1,a)$$
(1)
for all $`a,b,cA.`$
Given a Hopf $`2`$cocycle $`J`$ for $`A,`$ one can construct a new Hopf algebra $`(A^J,m^J,1,\mathrm{\Delta },\epsilon ,S^J)`$ as follows. As a coalgebra, $`A^J=A.`$ The new multiplication is given by
$$m^J(ab)=J^1(a_1,b_1)a_2b_2J(a_3,b_3)$$
(2)
for all $`a,bA.`$ The new antipode is given by
$$S^J(a)=J^1(a_1,S(a_2))S(a_3)J(S(a_4),a_5)$$
(3)
for all $`aA.`$
Suppose $`A`$ is also co(quasi)triangular with universal R-form $`R:AAk`$ (see e.g. \[K, VIII.5.1\]). Then it is straightforward to verify that $`A^J`$ is co(quasi)triangular with universal R-form $`R^J:A^JA^Jk,`$ where $`R^J:=(J\tau )^1RJ`$ (here $`\tau :AAAA`$ is the usual flip map).
Recall \[Dr\] that a twist for a Hopf algebra $`B`$ is an invertible element $`JBB`$ which satisfies
$$(\mathrm{\Delta }I)(J)(J1)=(I\mathrm{\Delta })(J)(1J)\text{and}(\epsilon I)(J)=(I\epsilon )(J)=1.$$
(4)
It is straightforward to verify that if $`A`$ is finite-dimensional, then $`JA^{}A^{}`$ is a Hopf $`2`$cocycle for $`A`$ if and only if it is a twist for $`A^{}.`$
## 3 The Main Theorem
Let $`(A,R)`$ be a cotriangular Hopf algebra over $`k`$ (not necessarily finite-dimensional). Define the Drinfeld element of $`(A,R)`$ to be the linear form $`u:Ak`$ determined by
$$u(a)=R(a_2,S(a_1)).$$
(5)
Recall that $`uA^{}`$ is a grouplike element, and that
$$S^2(a)=(uIu^1)(a)=u(a_1)a_2u^1(a_3)$$
(6)
for all $`aA.`$
Suppose $`cA^{}`$ is a central grouplike element of order $`2,`$ and set
$$R_c:=\frac{1}{2}(\epsilon \epsilon +\epsilon c+c\epsilon cc).$$
(7)
Then it is straightforward to verify that $`(A,RR_c)`$ is cotriangular with Drinfeld element $`uc.`$
Note that by (6), $`S^2`$ preserves any subcoalgebra of $`A.`$
###### Definition 3.1
We say that $`A`$ is pseudoinvolutive if $`\text{tr}(S^2|_C)=dim(C)`$ for any finite-dimensional subcoalgebra $`C`$ of $`A.`$
###### Remark 3.2
If $`A`$ is finite-dimensional, then pseudoinvolutivity is equivalent to involutivity ($`S^2=I`$). Indeed, by \[R\], $`S^2`$ has a finite order, so its eigenvalues on $`A`$ are roots of $`1.`$ But a sum of $`dim(A)`$ roots of $`1`$ can be equal to $`dim(A)`$ only if all of them are $`1.`$ $`\mathrm{}`$
We can now state our main theorem.
###### Theorem 3.3
Let $`(A,R)`$ be a cotriangular Hopf algebra over $`k`$. Then the following two conditions are equivalent:
(i) $`A`$ is pseudoinvolutive.
(ii) $`A`$ is a twisted function algebra $`𝒪(G)^J`$ on a pro-algebraic group $`G,`$ and furthermore, there exists a central grouplike element $`c𝒪(G)^{}`$ of order $`2,`$ such that $`(A,R)`$ is isomorphic to $`(𝒪(G)^J,(J\tau )^1R_cJ)`$ as cotriangular Hopf algebras.
###### Remark 3.4
Theorem 3.3 is a generalization of Theorem 2.1 from \[EG1\]. Indeed, let $`A`$ be a finite-dimensional triangular Hopf algebra over $`k.`$ Equivalently, $`A^{}`$ is a finite-dimensional cotriangular Hopf algebra over $`k.`$ Now, by Remark 3.2, $`A^{}`$ is pseudoinvolutive if and only if $`S^2=I.`$ By \[LR\], this is equivalent to the semisimplicity of $`A`$ and $`A^{}.`$ Hence by Theorem 3.3, $`A^{}`$ is a finite-dimensional semisimple cotriangular Hopf algebra over $`k`$ if and only if it is a twisted function algebra $`𝒪(G)^J`$ on a pro-algebraic group $`G`$ (possibly, with a changed cotriangular structure). But of course, $`G`$ must be a finite group. $`\mathrm{}`$
###### Remark 3.5
The data $`(G,c,J)`$ corresponding to a pseudoinvolutive cotriangular Hopf algebra over $`k,`$ is unique up to isomorphism of such triples and gauge transformations of $`J`$ (see \[EG2\]). The proof is similar to the proof of Lemma 3.5 in \[EG2\]. $`\mathrm{}`$
###### Remark 3.6
If $`(A,R)`$ is a cocommutative cotriangular Hopf algebra over $`k`$ (hence also commutative), then Theorem 3.3 is applicable since in this situation $`S^2=I.`$ Thus, \[BFM, Theorem 3.19(i)\], which claims that in this case $`(A,R)`$ is a twisted group algebra (maybe with $`RRR_c`$), is a special case of our result. $`\mathrm{}`$
###### Remark 3.7
In \[CWZ, Theorem 2.1\] the authors prove that if $`(A,R)`$ is cotriangular and its Drinfeld element $`u`$ acts as the identity on a finite-dimensional right $`A`$-comodule $`V,`$ then the characters of the usual action of the symmetric group $`S_n`$ on $`V^n`$ and the one arising from the braiding $`R`$ are equal. We note that this result is a consequence of Theorem 3.3. Namely, one should apply the theorem to the cotriangular Hopf algebra $`H_VA`$ which is generated by the elements of the form $`(fI)(\rho __V(v)),`$ $`vV,fV^{}`$ (here $`\rho __V`$ denotes the structure map of $`V`$), and get that the symmetric category of finite-dimensional right comodules of the cotriangular Hopf algebra $`H_V`$ is equivalent to that of $`𝒪(G)`$ for some pro-algebraic group $`G`$. Since the characters of $`S_n`$ on $`V^n`$ are invariant under equivalences of symmetric categories, the result follows. We also see that Theorem 2.1 of \[CWZ\] can be strengthened by replacing the assumption $`u|_V=1`$ by the weaker assumption that $`u|_V`$ is unipotent. $`\mathrm{}`$
## 4 The Proof of Theorem 3.3
### 4.1 Changing the Cotriangular Structure
In this subsection we prove the following proposition, which is one of the key ingredients in the proof of Theorem 3.3.
###### Proposition 4.1.1
For any pseudoinvolutive cotriangular Hopf algebra $`(A,R)`$ over $`k,`$ there exists a central grouplike element $`cA^{}`$ of order $`2`$ such that for the cotriangular Hopf algebra $`(A,RR_c)`$, with Drinfeld element $`u,`$ one has $`\text{tr}(u|_V)=\text{dim}(V)`$ for any finite-dimensional right $`A`$-comodule $`V.`$
The rest of the subsection is devoted to the proof of this proposition.
Let $`(A,R)`$ be a cotriangular Hopf algebra over $`k.`$ Let $`𝒞:=\text{Comod}_{f.d}(A)`$ be the category of finite-dimensional right $`A`$comodules, and $`\text{Irr}(𝒞)`$ be the subcategory of all irreducible objects of $`𝒞.`$ It is straightforward to check that $`𝒞`$ is an abelian rigid symmetric tensor category in the sense of \[DM\]. The unit object of $`𝒞`$ is $`\mathrm{𝟏}:=k,`$ and clearly $`\text{End}(\mathrm{𝟏})=k.`$ Also, recall that for any object $`V𝒞,`$ one can define its categorical dimension \[DM\], denoted by $`dim_c(V),`$ to be the image of $`1`$ under the morphism $`kVV^{}V^{}Vk`$ (where the morphism $`VV^{}V^{}V`$ is the braiding map). It is straightforward to verify that
$$dim_c(V)=\text{tr}|_V(u),$$
(8)
where $`u`$ is regarded as the linear map $`VV`$ determined by $`v(Iu)\rho __V(v)`$ (where $`\rho __V`$ denotes the structure map of $`V`$). Observe that $`\text{tr}|_V(u)=\text{tr}|_V(u^1)`$ for any $`V𝒞.`$ Indeed, we have that $`u^1=uS`$ and $`R=R(SS),`$ hence $`u(a)=u^1(a)`$ for any cocommutative element $`aA.`$ But, $`\text{tr}|_VA`$ is a cocommutative element.
For any object $`V𝒞,`$ set
$$A_V:=\{(fI)\rho __V(v)|vV,fV^{}\}.$$
(9)
It is clear that $`A_V`$ is a finite-dimensional subcoalgebra of $`A.`$
From now on we assume that $`A`$ is pseudoinvolutive.
###### Definition 4.1.2
We say that an object $`V𝒞`$ is positive if $`dim_c(V)=dim(V),`$ and negative if $`dim_c(V)=dim(V).`$
###### Lemma 4.1.3
An object of $`\text{Irr}(𝒞)`$ is either positive or negative, and any object of $`𝒞`$ is positive (resp. negative) if and only if so are all the composition factors of its Jordan-Hölder series.
Proof: Let $`X\text{Irr}(𝒞).`$ Then $`A_X=XX^{}.`$ Since $`\text{tr}|_X(u)=\text{tr}|_X(u^1),`$ it follows from (6) that $`dim(X)^2=\text{tr}(S^2|_{A_X})=(\text{tr}|_X(u))(\text{tr}|_X(u^1))=dim_c(X)^2.`$ Thus, $`dim_c(X)=\pm dim(X)`$ as desired. Moreover, for any $`V𝒞,`$ if $`0=V_0V_1\mathrm{}V_n=V`$ is its Jordan-Hölder series, then $`dim_c(V)=_{i=1}^ndim_c(V_i/V_{i1}),`$ and the result follows.
Consider now the abelian category of finite-dimensional right bicomodules over $`A,`$ i.e. right comodules over $`AA^{cop}.`$ It is clear that irreducible objects of this category are of the form $`XY^{},`$ where $`X,Y\text{Irr}(𝒞).`$
For any $`V𝒞,`$ it is clear that $`A_V`$ is a right $`A`$bicomodule (recall that $`A_V^{}`$ is the image of $`A^{}\text{End}(V),`$ so it is a left $`A^{}A^{op}`$-module). For any $`X,Y\text{Irr}(𝒞),`$ let $`N_V(X,Y)`$ be the multiplicity of occurrence of $`XY^{}`$ as a composition factor in the Jordan-Hölder series of $`A_V`$ regarded as a right bicomodule over $`A.`$
###### Lemma 4.1.4
For any $`V𝒞`$ and $`X,Y\text{Irr}(𝒞)`$ with opposite signs, $`N_V(X,Y)=0.`$
Proof: Indeed,
$`\text{tr}(S^2|_{A_V})`$
$`=`$ $`{\displaystyle \underset{X,Y\text{Irr}(𝒞)}{}}N_V(X,Y)\text{tr}|_{XY^{}}\left(u(u^1)^{}\right)`$
$`=`$ $`{\displaystyle \underset{X,Y\text{Irr}(𝒞)}{}}N_V(X,Y)dim_c(X)dim_c(Y).`$
But by pseudoinvolutivity, this should be equal to
$$dim(A_V)=\underset{X,Y\text{Irr}(𝒞)}{}N_V(X,Y)dim(X)dim(Y).$$
This implies that if $`dim_c(X)`$ and $`dim_c(Y)`$ have opposite signs then $`N_V(X,Y)=0,`$ otherwise $`_{X,Y\text{Irr}(𝒞)}N_V(X,Y)dim_c(X)dim_c(Y)<_{X,Y\text{Irr}(𝒞)}N_V(X,Y)dim(X)dim(Y).`$
###### Lemma 4.1.5
For any $`V,W𝒞`$ with opposite signs, one has $`\mathrm{Ext}^1(V,W)=0.`$
Proof: We have to show that any exact sequence $`0VUW0`$ in $`𝒞`$ splits. Indeed, we have natural coalgebra embeddings $`A_VA_U,`$ $`A_WA_U`$ ( generated by the coactions on the subcomodule and the quotient comodule), and the sum of them is a coalgebra embedding $`A_VA_WA_U`$ (it is obvious that the sum is direct as the coalgebras $`A_V,A_W`$ do not intersect, because of the opposite signs of $`V,W`$). Consider the quotient $`A_U/(A_VA_W)`$ as an $`A`$ bicomodule. Its all composition factors are of the form $`X_+X_{}^{},`$ where $`X_+\text{Irr}(𝒞)`$ is positive and $`X_{}\text{Irr}(𝒞)`$ is negative, which is impossible by Lemma 4.1.4. Therefore, this bicomodule must be zero. Thus, $`A_VA_WA_U.`$ Let $`W^{}:=\rho _U^1(UA_W).`$ Then $`U=VW^{},`$ and we are done.
###### Lemma 4.1.6
Any object $`V𝒞`$ can be uniquely represented as a direct sum $`V_+V_{},`$ where $`V_+𝒞`$ is positive and $`V_{}𝒞`$ is negative. Furthermore, for any two objects $`V,W𝒞,`$ $`(VW)_+=(V_+W_+)(V_{}W_{}),`$ and $`(VW)_{}=(V_+W_{})(V_{}W_+).`$
Proof: We first prove the existence part of the lemma by induction in $`dim(V).`$ Let $`V𝒞,`$ and $`X\text{Irr}(𝒞)`$ be an irreducible subcomodule of $`V.`$ Let us assume $`X`$ is positive. By the induction assumption we have $`V/X=W_+W_{}.`$ Let $`V_+`$ be the preimage of $`W_+`$ under the projection $`VW_+.`$ Then $`V_+`$ is positive, and $`V/V_+=W_{}.`$ Therefore, by Lemma 4.1.5, $`V`$ is isomorphic to $`V_+W_{},`$ and the result follows. If $`X`$ is negative, the proof is similar.
We now prove the uniqueness part of the lemma. Let $`V=V_+V_{}.`$ Then it is easy to see that $`V_+`$ is the sum of all positive subcomodules of $`V,`$ and $`V_{}`$ is the sum of all negative subcomodules, which implies uniqueness.
The last statement of the lemma is obvious.
Let $`D`$ be any coalgebra over $`k,`$ let $`\text{Comod}(D)`$ be its category of right comodules and $`F:\text{Comod}(D)Vec`$ be the forgetful functor. Recall that $`\text{End}(F)D^{}`$ as algebras. Indeed, an element $`\eta \text{End}(F)`$ is by definition, a collection of linear maps $`\eta _V:VV,`$ $`V\text{Comod}(D),`$ which commute with comodule morphisms. In particular, $`\eta _D:DD`$ is a linear map between right $`D`$comodules, and it commutes with right actions by elements of $`D^{}.`$ Thus, $`\eta _D`$ comes from a left action by an element of $`D^{}`$ (namely, by $`\eta _D^{}(\epsilon )`$). Note that if moreover $`\eta _V:VV`$ is a comodule map for all $`V\text{Comod}(D),`$ then we have an endomorphism of the identity functor $`Id`$ of $`\text{Comod}(D),`$ and the resulting element of $`D^{}`$ is central. Thus, $`\text{End}(Id)\text{Center}(D^{})`$ as algebras.
For any $`V𝒞,`$ define the comodule automorphism $`c__V`$ of $`V`$ by $`c__V:=I`$ on $`V_+`$ and $`c__V:=I`$ on $`V_{},`$ where $`I`$ is the identity map of $`V.`$
###### Lemma 4.1.7
The collection $`\{c__V|V𝒞\}`$ determines a central grouplike element $`cA^{}`$ of order $`2.`$
Proof: The collection $`\{c__V|V𝒞\}`$ is an element of the algebra $`\text{End}(Id),`$ hence, by the preceding remarks, determines a central element $`cA^{}.`$ Now, it is clear from Lemma 4.1.6 that $`c`$ is a grouplike element of order $`2.`$
Now let us finally prove Proposition 4.1.1. Let $`cA^{}`$ be the central grouplike element of order $`2`$ whose existence is guaranteed by Lemma 4.1.7. Let $`R_c`$ be as in (7). Then, after changing $`R`$ to $`RR_c,`$ the new Drinfeld element is $`u^{}:=uc,`$ and we get that $`\text{tr}|_V(u^{})=dim(V)`$ for any object $`V𝒞.`$ The proposition is proved.
### 4.2 The Proof of Theorem 3.3
$`(i)(ii)`$.
Proposition 4.1.1 implies that without loss of generality, we can assume that $`\text{tr}(u|_V)=\text{dim}(V)`$ for all finite-dimensional right $`A`$comodules.
Now comes the main step of the proof, which is the usage of the following theorem of Deligne.
###### Theorem 4.2.1 (De, Theorem 7.1)
Let $`𝒞`$ be an abelian rigid symmetric tensor category over $`k`$ such that $`\text{End}(\mathrm{𝟏})=k,`$ in which categorical dimensions of objects are non-negative integers. Then there exist a pro-algebraic group $`G`$ and a $`k`$linear equivalence of abelian rigid symmetric tensor categories $`F:𝒞\text{Rep}_{f.d}(G)`$ (where $`\text{Rep}_{f.d}(G)`$ is the category of finite-dimensional algebraic $`k`$representations of $`G`$).
This theorem implies that in our situation, we have an equivalence
$$F:\text{Comod}_{f.d}(A)\text{Comod}_{f.d}(𝒪(G))$$
of rigid symmetric tensor categories. It is obvious that $`F`$ preserves dimensions.
Now we will need the following proposition, whose proof occupies the next subsection.
###### Proposition 4.2.2
Let $`A`$ and $`B`$ be two coassociative coalgebras with counit over $`k,`$ and $`F:\text{Comod}_{f.d}(A)\text{Comod}_{f.d}(B)`$ be an equivalence between the abelian categories of finite-dimensional right comodules over $`A`$ and $`B,`$ which preserves dimensions. Then there exists an isomorphism of coalgebras $`\varphi :AB`$ such that $`F`$ is isomorphic to the direct image functor $`\varphi _{}`$.
###### Remark 4.2.3
In the case when $`A`$,$`B`$ are cosemisimple, this proposition is trivial. Therefore, it was not spelled out explicitly in our previous papers, where we dealt exclusively with the semisimple case. $`\mathrm{}`$
This proposition implies that there exists an isomorphism of coalgebras $`\varphi :A𝒪(G)`$ that induces $`F.`$ Therefore, we can naturally identify the vector spaces $`V`$ and $`F(V)`$ for all $`V𝒞`$, in a functorial way.
Now, recall that $`F`$ has a tensor structure. Namely, we have a collection of right $`𝒪(G)`$comodule isomorphisms $`J_{VW}:F(V)F(W)F(VW)`$ indexed by all pairs $`V,W𝒞.`$ But for any $`U𝒞,`$ we have already identified the vector spaces $`U`$ and $`F(U).`$ Thus, we can regard the tensor structure as a collection of isomorphisms of vector spaces $`VWVW,`$ which is functorial with respect to $`V`$ and $`W.`$ This collection isomorphisms defines an element $`J(𝒪(G)𝒪(G))^{}`$ (see the preceding remarks to Lemma 4.1.7). It can be checked (similarly to Theorem 2.1 in \[EG1\]) that $`J`$ is a Hopf 2-cocycle, and that $`\varphi `$ induces an isomorphism of cotriangular Hopf algebras $`A𝒪(G)^J`$. The implication $`(i)(ii)`$ is proved.
$`(ii)(i)`$.
We may assume that $`A=𝒪(G)^J`$ (since $`S`$ does not depend on the cotriangular structure). Since the categorical dimension of any object $`V𝒞`$ does not change under twisting, we have $`dim_c(V)=dim(V)`$ in the category of finite-dimensional right comodules over $`A`$. Therefore, we have
$`\text{tr}(S^2|_{A_V})`$
$`=`$ $`{\displaystyle \underset{X,Y\text{Irr}(𝒞)}{}}N_V(X,Y)dim_c(X)dim_c(Y)`$
$`=`$ $`{\displaystyle \underset{X,Y\text{Irr}(𝒞)}{}}N_V(X,Y)dim(X)dim(Y)`$
$`=`$ $`dim(A_V).`$
Since any finite-dimensional subcoalgebra $`C`$ of $`A`$ has the form $`A_V`$ (for $`V=C`$), this completes the proof of the theorem.
### 4.3 Proof of Proposition 4.2.2
Let us first prove the proposition in the case when $`A,B`$ are finite-dimensional.
We need the following standard theorem from noncommutative algebra, which can be found for example in \[DK, Chapter 3\].
Let $`𝒜`$ be a finite-dimensional algebra. Let $`\text{Irr}(𝒜)`$ be the set of isomorphism classes of irreducible left $`𝒜`$-modules. For $`M\text{Irr}(𝒜)`$, let $`P(M)`$ be the projective cover of $`M`$.
###### Theorem 4.3.1
(i) Any finite-dimensional projective $`𝒜`$-module is a direct sum of $`P(M)^{}s.`$
(ii) For any $`M\text{Irr}(𝒜),`$ the multiplicity of $`P(M)`$ in the regular representation $`𝒜`$ is equal to $`\text{dim}(M).`$
Now we prove the proposition (in the finite-dimensional case). Let $`F_A,F_B`$ be the forgetful functors from the categories of right $`A`$-comodules and right $`B`$-comodules, respectively, to the category of vector spaces. All we need to show is that $`F_BF`$ is isomorphic to $`F_A`$. Indeed, we have $`\text{End}(F_A)=A^{}`$ and $`\text{End}(F_BF)=B^{}`$, so any isomorphism between these two functors will induce an isomorphism of coalgebras $`AB`$, which (as one can easily see) induces $`F`$.
The functor $`F_A`$ is represented by the regular representation $`A^{}`$, and $`F_BF`$ by $`F^1(B^{})`$. So it suffices to prove that $`F^1(B^{})`$ is isomorphic to $`A^{}`$ as an $`A`$-comodule.
Since $`B^{}`$ is free, it is projective, so $`F^1(B^{})`$ is also projective (as projectivity, unlike freeness, is a categorical property). Thus, by Theorem 4.3.1, we have
$$A^{}=\underset{M\text{Irr}(A^{})}{}\text{dim}(M)P(M)\text{and}F^1(B^{})=\underset{M\text{Irr}(A^{})}{}x(M)P(M),$$
where $`x(M)`$ are nonnegative integers. But since $`F`$ preserves dimensions, we have for any $`M\text{Irr}(A^{})`$:
$$\text{dim}(M)=\text{dim}(F(M))=\text{dim}\left(\text{Hom}_B^{}(B^{},F(M))\right)=\text{dim}\left(\text{Hom}_A^{}(F^1(B^{}),M)\right)=x(M).$$
This completes the proof of the proposition in the finite-dimensional case.
Now let us consider the infinite-dimensional case. For simplicity consider the case when $`A,B`$ are countably dimensional (the general case is similar). Then $`A=_{n1}A_n,`$ where $`A_n`$ are finite-dimensional coalgebras. Let $`F_n:\text{Comod}_{f.d}(A_n)\text{Comod}_{f.d}(B)`$ be the restriction of $`F,`$ and let $`B_n:=\text{End}(\text{Forget}F_n)^{},`$ where Forget is the forgetful functor on $`B`$-comodules. It is clear that $`B=_{n1}B_n.`$
It is clear from the above finite-dimensional proof that for any $`n,`$ there exists an isomorphism of coalgebras $`\varphi _n:A_nB_n,`$ such that $`\varphi _{n+1}|_{A_n}=\varphi _n\text{Ad}(a_n)^{},`$ where $`a_nA_n^{}`$ is an invertible element (this follows from the fact that $`\varphi _n`$ comes from an isomorphism of functors). Since the map of the multiplicative groups $`Gr(A^{})Gr(A_n^{})`$ is surjective (because so is the corresponding Lie algebra map), we can lift $`a_n`$ to an invertible element of $`A^{}.`$ Abusing notation, we will denote this element also by $`a_n.`$
Define $`\psi _n:=\varphi _n\text{Ad}(a_1^1\mathrm{}a_{n1}^1)^{}.`$ Then $`\psi _{n+1}|_{A_n}=\psi _n`$, so $`\{\psi _n\}`$ defines an isomorphism of coalgebras $`\psi :AB`$ which induces a functor isomorphic to $`F.`$
## 5 Examples of Twisted Function Algebras
In this section we give examples of Hopf $`2`$cocycles for $`𝒪(G)`$ for certain algebraic groups $`G`$. We will construct these cocycles in the form of linear endomorphisms of tensor products of any two finite-dimensional $`G`$-modules $`V,W`$, functorial in terms of $`V,W`$ (cf. the discussion before Lemma 4.1.7). These examples can be generalized, as usual, using the fact that if $`J`$ is a Hopf 2-cocycle for $`G`$ and $`\varphi :GG^{}`$ is a homomorphism then $`(\varphi \varphi )(J)`$ is a Hopf 2-cocycle for $`G^{}`$.
###### Example 5.1
Let $`G`$ be the group of translations of an affine space. Let $`𝔤`$ be the Lie algebra of $`G.`$ Clearly, $`𝔤`$ is abelian. Therefore, it is straightforward to verify that for any $`r\mathrm{\Lambda }^2𝔤,`$ the element $`J(h):=e^{hr/2}`$ is a Hopf 2-cocycle for any $`hk`$ (the exponential series terminates in any $`VW`$ as the components of $`r`$ are nilpotent, so we get a polynomial of $`h`$). In this case $`u=1`$ and $`S^2`$ is the identity. $`\mathrm{}`$
###### Example 5.2
Let $`G`$ be the group of affine transformations of the line. Its Lie algebra $`𝔤`$ is spanned by two elements $`X,Y`$ such that $`[X,Y]=Y`$. Define
$$J(h):=\underset{n0}{}\frac{h^n}{n!}X(X1)\mathrm{}(Xn+1)Y^n.$$
It is not difficult to check that this is a Hopf 2-cocycle. In this case $`u=1+hY+O(h^2),`$ so $`S^2`$ is not the identity. Thus $`𝒪(G)^J`$ is an example of a pseudoinvolutive but not involutive cortiangular Hopf algebra. Such a Hopf algebra can even be cosemisimple: it is enough to naturally embed $`G`$ into $`GL(2)`$ and consider the cosemisimple Hopf algebra $`𝒪(GL_2)^J`$. $`\mathrm{}`$
###### Example 5.3
For every nonnegative integer $`n,`$ let $`S_n`$ be the symmetric group of permutations of $`n`$ symbols. For any $`sS_n`$ and $`0mn,`$ and any Lie algebra $`𝔤,`$ let $`L_{s,m}:𝔤^nU(𝔤)^2`$ be the linear map determined by
$$L_{s,m}(a_1\mathrm{}a_n)=a_{s(1)}\mathrm{}a_{s(m)}a_{s(m+1)}\mathrm{}a_{s(n)}.$$
Let $`X`$ be the disjoint union of the sets $`X_n:=S_{2n}\times \{0,\mathrm{},2n\}`$ for $`n2.`$ Let $`k[X]`$ be the set of $`k`$valued functions on $`X.`$ For any $`fk[X],`$ and any Lie algebra $`𝔤,`$ define the function $`J_f:\mathrm{\Lambda }^2𝔤U(𝔤)^2[[h]],`$ by the formula
$$J_f(r,h)=1+hr/2+\underset{n2}{}h^n\underset{(s,m)X_n}{}f(s,m)L_{s,m}(r^n).$$
(10)
It is easy to show that if $`J_f(r,h)`$ is a twist for $`U(𝔤)[[h]],`$ then $`r`$ must satisfy the classical Yang-Baxter equation (CYBE):
$$[r_{12},r_{13}]+[r_{12},r_{23}]+[r_{13},r_{23}]=0.$$
(11)
Conversely, let $`r`$ be a solution of CYBE.
###### Definition 5.4
We say that $`fk[X]`$ is a quantization function for $`r`$ if the element $`J_f(r,h)`$ is a twist. We say that $`f`$ is a universal quantization function if this is the case for all $`𝔤,r`$.
It is easy to construct concrete examples of quantization functions. For instance, in Example 5.1, it is straightforward to verify that $`e^{hr/2}`$ comes from a quantization function for any $`r`$ and abelian $`𝔤.`$ The existence of universal quantization functions is not obvious, but it follows from the quantization theory of \[EK1\]. Namely, a construction of such functions can be obtained from formula (3.1) in \[EK1\].
Quantization functions allow one to construct a large family of examples of Hopf 2-cocycles. Namely, we have:
###### Theorem 5.5
Suppose that $`G`$ is an algebraic group with Lie algebra $`𝔤,`$ and let $`N`$ be its unipotent radical with Lie algebra $`𝔫.`$ Suppose that $`r`$ is an element of $`𝔤𝔫\mathrm{\Lambda }^2𝔤`$ which satisfies the CYBE. Then for any $`fk[X],`$ and any two algebraic representations $`V,W`$ of $`G,`$ $`J_f(r,h)|_{VW}`$ is a polynomial in $`h`$ (i.e. the series terminates). In particular, for any $`hk,`$ $`J_f(r,h)`$ is a well defined element in $`(𝒪(G)𝒪(G))^{}.`$ This element is a Hopf 2-cocycle for $`𝒪(G)`$ if $`f`$ is a quantization function for $`r`$.
Proof: We only need to show that the series $`J_f(r,h)`$ terminates. Let $`V,W`$ be two finite-dimensional algebraic representations of $`G`$. Let $`B_V,B_W`$ be the images of $`U(𝔤)`$ in $`\text{End}(V)`$ and $`\text{End}(W),`$ and let $`I_V,I_W`$ be the nilpotent radicals of $`B_V,B_W`$. It is clear that under the action of $`G`$ in $`V,W`$, $`𝔫`$ maps into $`I_V,I_W`$. Let $`n`$ be a positive integer such that $`I_V^n=I_W^n=0`$. Then it is clear from the definition that the series $`J_f(r,h)|_{VW}`$ is a polynomial in $`h`$ of degree $`2n2`$.
## 6 The Unipotency Conjecture
We conclude the paper with the following
###### Conjecture 6.1
For any pro-algebraic group $`G`$ and Hopf $`2`$cocycle $`J`$ for $`𝒪(G)`$ over $`k,`$ the operator $`S^2`$ is unipotent on $`𝒪(G)^J.`$
###### Remark 6.2
We know that the sum of the eigenvalues of $`S^2`$ on any finite-dimensional subcoalgebra is equal to its dimension, but the conjecture says that furthermore all of these eigenvalues are $`1.`$ $`\mathrm{}`$
###### Remark 6.3
Conjecture 6.1 is obviously satisfied in Examples 5.1 and 5.2. Moreover, it follows from the theorem below that it is satisfied in Example 5.3 as well. $`\mathrm{}`$
Let $`\mathrm{\Sigma }`$ be an irreducible affine algebraic curve with a marked smooth point $`0`$, and let $`𝒪(\mathrm{\Sigma })`$ be the ring of regular functions on $`\mathrm{\Sigma }`$. The standard example is $`\mathrm{\Sigma }:=k,`$ $`𝒪(\mathrm{\Sigma })=k[x].`$
Let $`J:𝒪(G)^2𝒪(\mathrm{\Sigma })`$ be a family of Hopf 2-cocycles for $`𝒪(G)`$ parametrized by $`a\mathrm{\Sigma },`$ with $`J(0)=1.`$ In this case, we will say that $`J(a),`$ for any $`a\mathrm{\Sigma },`$ is obtained by deformation of $`J(0).`$
###### Theorem 6.4
Let $`J`$ be obtained by deformation of $`1`$. Then Conjecture 6.1 holds for the cotriangular Hopf algebra $`𝒪(G)^J.`$
The rest of the section is devoted to the proof of the theorem.
To prove the theorem, we will choose a local parameter $`h`$ on $`\mathrm{\Sigma }`$, and write $`J`$ as a formal power series in $`h`$: $`J=1+_{n1}h^nr_n`$, $`r_n(𝒪(G)^2)^{}`$. We will say that $`J`$ is local if $`r_nU(𝔤)^2`$ for all $`n`$.
###### Lemma 6.5
The Hopf $`2`$cocycle $`J`$ is gauge equivalent to a local Hopf $`2`$cocycle. That is, there exists a “gauge transformation” $`g:=1+hg_1+h^2g_2+\mathrm{}`$, $`g_i𝒪(G)^{}`$, $`\epsilon (g_i)=0`$, such that the Hopf $`2`$cocycle $`J^g:=\mathrm{\Delta }(g)J(g^1g^1)`$ is local.
Proof: Let us prove the statement modulo $`h^{n+1}`$ by induction in $`n.`$ The base of induction ($`n=0`$) is clear. To do the inductive step, assume that $`J`$ is local modulo $`h^n`$. Observe that since $`r_n`$ satisfies the Hopf $`2`$cocycle condition it follows that
$$r_n^{12}+(\mathrm{\Delta }I)(r_n)r_n^{23}(I\mathrm{\Delta })(r_n)=f(r_1,\mathrm{},r_{n1}),$$
where $`f`$ is a polynomial. Thus, we have $`dr_nU(𝔤)^3`$, where $`d`$ is the differential in the Hochschild complex of the coalgebra $`𝒪(G)^{}`$ with trivial coefficients.
It is well known that the embedding of coalgebras $`U(𝔤)𝒪(G)^{}`$ defines an isomorphism of Hochschild cohomology of these coalgebras with trivial coefficients, and that both cohomology spaces are equal to $`\mathrm{\Lambda }𝔤`$, with the usual grading (the fact that the cohomology of $`U(𝔤)`$ is $`\mathrm{\Lambda }𝔤`$ is discussed for example in \[Dr1, p.1435\]). Therefore, there exists $`r_n^{}U(𝔤)^2`$ such that $`dr_n^{}=dr_n`$. Let $`s:=r_nr_n^{}`$. Then $`ds=0`$. Therefore, we have $`s=s_0+dz=s_0+\mathrm{\Delta }(z)z11z`$, where $`s_0\mathrm{\Lambda }^2𝔤`$ and $`z𝒪(G)^{}`$. Let us replace $`J`$ with $`J^g`$ for $`g:=1+h^n(zϵ(z))`$. Then $`J^g`$ is local modulo $`h^{n+1}`$. The lemma is proved.
Now let us continue the proof of the theorem. By Lemma 6.5, we can assume that $`J(h)`$ is local. Then $`J(h)`$ is a twist for $`U(𝔤)[[h]]`$, so one can define the triangular quantized universal enveloping (QUE) algebra $`U(𝔤)[[h]]^{J(h)}`$. It is sufficient to show that the Drinfeld element $`u`$ of this QUE algebra is unipotent on every finite-dimensional representation of $`𝔤`$.
Using the main result of \[EK2\] (in the triangular case), we conclude that $`U(𝔤)[[h]]^{J(h)}`$ is isomorphic to $`U_h(𝔤,r)`$, where $`U_h`$ is the quantization functor form \[EK2\], and $`r\mathrm{\Lambda }^2𝔤[[h]]`$ is a solution of CYBE.
The QUE algebra $`U_h(𝔤,r)`$, by definition, is obtained by twisting $`U(𝔤)`$ using a twist $`J_f(r,h)`$, where $`f`$ is a universal quantization function. This implies that the element $`u`$ in this QUE algebra has the form
$$u=1+\underset{n1}{}\underset{\sigma S_{2n}}{}c_{n,\sigma }m_{2n}(\sigma r^n),$$
where $`m_{2n}`$ is the multiplication of $`2n`$ components.
Now we show the unipotency of $`u`$ in any finite-dimensional $`𝔤`$-module by induction in the rank of $`r`$. Without loss of generality, we can assume that $`𝔤`$ is spanned by the components of $`r`$, and that $`V`$ is an irreducible faithful $`𝔤`$-module (if it is not faithful, we can go to r-matrices of smaller rank, for which the statement is known). This implies that $`𝔤`$ is reductive.
Now we will use the following theorem of Drinfeld (see e.g. \[ES, Proposition 5.2\]):
###### Theorem 6.6
Let $`𝔤`$ be a Lie algebra and $`r\mathrm{\Lambda }^2𝔤`$ be a solution of CYBE whose components span $`𝔤`$. If $`𝔤`$ is reductive, then it is abelian.
This theorem immediately implies Theorem 6.4: if $`𝔤`$ is abelian then, because of the skew-symmetry of $`r`$, we have $`u=1`$.
The rest of this section is the proof of Theorem 6.6, which we give for the reader’s convenience.
It is clear that $`r`$ defines a nondegenerate skew-symmetric bilinear form on $`𝔤^{}`$, hence one can define a skew-symmetric form $`r^1`$ on $`𝔤`$. Since $`r`$ satisfies CYBE, this form is a 2-cocycle.
Let $`𝔤=𝔤_s𝔤_a`$ be the splitting of $`𝔤`$ into the semisimple and the abelian parts. Assume that $`𝔤_s0`$. Then $`H^2(𝔤,k)=H^2(𝔤_a,k)=\mathrm{\Lambda }^2𝔤_a^{}`$, which shows that $`r^1`$ can be written as $`r^1(x,y)=f([x,y])+\rho (x,y)`$, where $`\rho \mathrm{\Lambda }^2𝔤_a^{}`$, $`f𝔤^{}`$. Thus, the decomposition $`𝔤=𝔤_s𝔤_a`$ is orthogonal under $`r^1`$, and hence the form $`f([x,y])`$ has to be nondegenerate on $`𝔤_s`$.
But $`𝔤_s`$ has a nondegenerate invariant form $`(,)`$, and we can write the functional $`f(x)`$ as $`(z,x)`$ for some $`z𝔤_s`$. Thus, the form $`(z,[x,y])`$ should be nondegenerate. However, this is impossible as $`z`$ obviously belongs to the kernel of this form (since $`(z,[z,y])=0`$). Thus, we have a contradiction and hence $`𝔤_s=0`$.
## 7 Questions
In conclusion let us discuss possible directions for future research. The main remaining problem is to obtain a classification of twisted group algebras. We believe that this should be done by generalizing the techniques of \[M\] and \[EG2\] to the infinite-dimensional case and combining them with the techniques of \[EK1,EK2\]. Let us formulate some precise questions which are related to this problem.
###### Question 7.1
Is any Hopf 2-cocycle for $`𝒪(G)`$ gauge equivalent to a deformation of a Hopf 2-cocycle of finite rank?
###### Question 7.2
Say that a cotriangular Hopf algebra is minimal if the R-form is nondegenerate. Now suppose that $`A`$ is a minimal pseudoinvolutive cotriangular Hopf algebra, and that the underlying group $`G`$ is reductive. Is it true that the connected component of the identity in $`G`$ is abelian (i.e. a torus)?
###### Remark 7.3
The answer to the classical analog of Question 7.2 is positive: if a Lie algebra spanned by the components of a skew-symmetric solution of CYBE is reductive, then it is abelian (see Theorem 6.6). $`\mathrm{}`$
###### Question 7.4
Let $`A`$ be any cotriangular Hopf algebra over $`k.`$ Is it true that the eigenvalues of the square of its antipode are all roots of unity? Are all $`\pm 1`$? Is it true at least for twisted function algebras?
Note that a positive answer to either form of Question 7.4 will imply Conjecture 6.1. |
warning/0002/math0002126.html | ar5iv | text | # Secondary Characteristic Classes and Cyclic Cohomology of Hopf Algebras
## 1 Introduction
In the paper A. Connes provided an explicit construction of the Godbillon-Vey cocycle in the cyclic cohomology. The goal of this paper is to give a similar construction for the higher secondary classes.
First, let us recall Connes’ construction. Let $`M`$ be a smooth oriented manifold and let $`\mathrm{\Gamma }\mathrm{Diff}^+(M)`$ be a discrete group of orientation-preserving diffeomorphisms of $`M`$. Let $`\omega `$ be a volume form on $`M`$. Define the following function on $`M\times \mathrm{\Gamma }`$: $`\delta (g)=\frac{\omega ^g}{\omega }`$, where the superscript denotes the group action. Then one can define one-parametric group of diffeomorphisms of the algebra $`𝒜=C_0^{\mathrm{}}(M)\mathrm{\Gamma }`$ by
$$\varphi _t(aU_g)=a\delta (g)^tU_g$$
(1.1)
This is Tomita-Takesaki group of automorphisms, associated to the weight on $`𝒜`$ given by $`\omega `$.
Consider now the transverse fundamental class – cyclic $`q`$-cocycle on $`𝒜`$ given by
$$\begin{array}{c}\tau (a_0U_{g_0},a_1U_{g_1},\mathrm{},a_qU_{g_q})=\hfill \\ \hfill \{\begin{array}{cc}\frac{1}{q!}\underset{M}{}a_0𝑑a_1^{g_0}𝑑a_2^{g_0g_1}\mathrm{}𝑑a_q^{g_0g_1\mathrm{}g_{q1}}\hfill & \text{ if }g_0g_1\mathrm{}g_q=1\hfill \\ 0\hfill & \text{ otherwise}\hfill \end{array}\end{array}$$
(1.2)
To study the behavior of this cocycle under the 1-parametric group (1.1), consider the “Lie derivative” $``$ acting on the cyclic complex by $`\xi =\frac{d}{dt}|_{t=0}\varphi _t^{}\xi `$, $`xi`$ being a cyclic cochain. It turns out that in general $`\tau `$ is not invariant under the group (1.1), and $`\tau 0`$.
However, it was noted by Connes that one always has
$$^{q+1}\tau =0$$
(1.3)
and $`^q\tau `$ is invariant under the action of the group (1.1). One deduces from this that if $`\iota _\delta `$ is the analogue of the interior derivative (see ), then $`\iota _\delta ^q\tau `$ is a cyclic cocycle.
This is Connes’ Godbillon-Vey cocycle. It can be related to the Godbillon-Vey class as follows. Let $`[GV]H^{}(M_\mathrm{\Gamma })`$ be the Godbillon-Vey class, where $`M_\mathrm{\Gamma }=M\times _\mathrm{\Gamma }\mathrm{E}\mathrm{\Gamma }`$ is the homotopy quotient. Connes defines canonical map $`\mathrm{\Phi }:H^{}(M_\mathrm{\Gamma })HP^{}(𝒜)`$. Then one has
$$\mathrm{\Phi }([GV])=[\iota _\delta ^q\tau ]$$
(1.4)
The class of this cocycle is independent of the choice of the volume form. To prove this one can use Connes’ noncommutative Radon-Nicodym theorem to conclude that if one changes the volume form, the one-parametric group $`\varphi _t`$ remains the same modulo inner automorphisms.
A natural problem then is to extend this construction to the cocycles corresponding to the other secondary characteristic classes. It was noted by Connes that if instead of 1-dimensional bundle of $`q`$-forms on $`M`$ one considers $`\mathrm{\Gamma }`$ equivariant trivial bundle of rank $`n`$, then in place of 1-parametric group (1.1) one encounters coaction of the group $`GL_n()`$ on the algebra $`𝒜`$. The difficulty is that for $`n>1`$ this group is not commutative, and one can not replace this coaction by the action of the dual group, similarly to (1.1).
In this paper we show that Connes-Moscovici theory of cyclic cohomology for Hopf algebras (cf ) provides a natural framework for the higher-dimensional situation and allows one to give construction of the secondary characteristic cocycles.
The situation we consider is the following. We have an orientation-preserving action of discrete group (or pseudogroup) $`\mathrm{\Gamma }`$ on the oriented manifold $`M`$, and a trivial bundle $`E`$ on $`M`$ equivariant with respect to this action. Well-known examples in which such a situation arises are the following (cf. , ,). Let $`V`$ be a manifold on which a discrete group (or pseudogroup) $`G`$ acts, and let be $`E_0`$ a bundle (not necessarily trivial) on $`V`$, equivariant with respect to the action of $`G`$. Let $`U_i`$, $`iI`$, be an open covering of $`V`$ such that restriction of $`F`$ on each $`U_i`$ is trivial. Put $`M=U_i`$, and let $`E`$ be the pull-back of $`E_0`$ to $`M`$ by the natural projection. Then we have an action of the following pseudogroup $`\mathrm{\Gamma }`$ on $`M`$: $`\mathrm{\Gamma }=\{g_{i,j}|gGi,jI\}id`$, where $`\mathrm{Dom}g_{i,j}=g^1\left(U_j\right)U_iU_i`$, $`\mathrm{Ran}g_{i,j}=g\left(U_i\right)U_jU_j`$, and the natural composition rules. The bundle $`E`$ is clearly equivariant with respect to this action. Our construction, described below, provides classes in the cyclic cohomology of the cross-product algebra $`C_0^{\mathrm{}}(M)\mathrm{\Gamma }`$, rather than in the cyclic cohomology of $`C_0^{\mathrm{}}(V)G`$. However, cross-product algebras $`C_0^{\mathrm{}}(M)\mathrm{\Gamma }`$ and $`C_0^{\mathrm{}}(V)G`$ are Morita equivalent, and hence have the same cyclic cohomology.
Another natural example is provided by the manifold $`V`$ with foliation $`F`$, and a bundle $`E_0`$ which is holonomy equivariant. We can always choose (possibly disconnected) complete transversal $`M`$, such that restriction $`E`$ of $`E_0`$ to $`M`$ is trivial. Let $`\mathrm{\Gamma }`$ be the holonomy pseudogroup acting on $`M`$. $`E`$ is clearly equivariant with respect to this action. In this case again the cross-product algebra $`C_0^{\mathrm{}}(M)\mathrm{\Gamma }`$ is Morita equivalent to the full algebra of the foliation $`C_0^{\mathrm{}}(V,F)`$.
We construct then a map from the cohomology of the truncated Weil algebra (cf. e.g. ) $`W(𝔤,O_n)_q`$ to the periodic cyclic cohomology $`HP^{}(𝒜)`$ of the algebra $`𝒜=C_0^{\mathrm{}}(M)\mathrm{\Gamma }`$. The construction is the following. We consider the action of the *differential graded* Hopf algebra $`\left(GL_n()\right)`$ of differential forms on the group $`GL_n()`$ on the differential graded algebra $`\mathrm{\Omega }_0^{}(M)\mathrm{\Gamma }`$, where $`\mathrm{\Omega }_0^{}(M)`$ denotes the algebra of compactly supported differential forms on $`M`$. The use of differential graded algebras allows one to conveniently encode different identities, similar to (1.3). We then show that Connes-Moscovici theory (or rather differential graded version of it) allows one to define a map from the cyclic complex of $`\left(GL_n()\right)`$ to the cyclic complex of $`\mathrm{\Omega }_0^{}(M)\mathrm{\Gamma }`$. We then relate cyclic complex of $`\left(GL_n()\right)`$ to the Weil algebra, and cyclic cohomology of $`\mathrm{\Omega }_0^{}(M)\mathrm{\Gamma }`$ to the cyclic cohomology of $`𝒜`$.
The paper is organized as follows. In the next two sections we discuss cyclic complexes for differential graded algebras and differential graded Hopf algebras respectively. In the section 4 we show that two different Hopf actions, which coincide “modulo inner automorphisms” induce the same Connes-Moscovici characteristic map in cyclic cohomology, and discuss some other properties of the characteristic map. In the section 5 we construct the action of $`\left(GL_n()\right)`$ on $`\mathrm{\Omega }_0^{}(M)\mathrm{\Gamma }`$. In the section 6 we relate cyclic complex of the Hopf algebra $`\left(GL_n()\right)`$ with the Weil algebras. Finally, in the section 7 we prove an analogue of the formula (1.4) for the cocycles we construct.
I would like to thank D. Burghelea and H. Moscovici for many helpful discussions.
## 2 Cyclic complex for differential graded algebras
In this section we collect some preliminary standard facts about cyclic cohomology of the differential graded algebras, and give cohomological version of some results of .
Recall that the cyclic module $`X^{}`$ is given by the cosimplicial module with the face maps $`\delta _i:X^{n1}X^n`$ and and degeneracy maps $`\sigma _i:X^nX^{n1}`$ $`0in`$, satisfying the usual axioms. In addition, we have for each $`n`$ an action of $`_{n+1}`$ on $`X^n`$, with the generator $`\tau _n`$ satisfying
$$\tau _n\delta _i=\delta _{i1}\tau _{n1}\text{ for }1in\text{ and }\tau _n\delta _0=\delta _n$$
(2.1)
$$\tau _n\sigma _i=\sigma _{i1}\tau _{n+1}\text{ for }1in\text{ and }\tau _n\sigma _0=\sigma _n\tau _{n+1}^2$$
(2.2)
$$\tau _n^{n+1}=id$$
(2.3)
For every cyclic object one can construct operators $`b:X^nX^{n+1}`$ and $`B:X^nX^{n1}`$, defined by the formulas
$$b=\underset{j=0}{\overset{n}{}}(1)^jd_j$$
(2.4)
$$B=\left(\underset{j=0}{\overset{n1}{}}(1)^{j(n1)}\tau _{n1}^j\right)\sigma _{n+1}(1(1)^{n1}\tau _n)$$
(2.5)
where
$$\sigma _{n+1}=\sigma _n\tau _{n+1}$$
(2.6)
These operators satisfy
$`b^2`$ $`=0`$ (2.7)
$`B^2`$ $`=0`$ (2.8)
$`bB+Bb`$ $`=0`$ (2.9)
Hence for any cyclic object $`X^{}`$ we can construct a bicomplex $`^,(X)`$ as follows: $`^{p,q}`$, $`p`$, $`q0`$ is $`X^{pq}`$, or $`0`$ if $`p<q`$, and the differential $`^{p,q}^{p+1,q}`$ (resp. $`^{p,q+1}`$) is given by $`b`$ (resp. $`B`$). By removing restriction $`p`$, $`q0`$ we obtain periodic bicomplex $`_{\text{per}}`$. Notice that it has periodicity induced by the tautological shift $`S:_{\text{per}}^{p,q}_{\text{per}}^{p+1,q+1}`$.
Let now $`\mathrm{\Omega }^{}`$ be a unital graded (DG) algebra, positively graded. We can associate with it a cyclic object as follows (the differential $`d`$ is not used in this definition).
Let $`C^k(\mathrm{\Omega }^{})`$ be the space of continuous $`k+1`$-linear functionals on $`\mathrm{\Omega }^{}`$. The face and degeneracy maps are given by
$`(\delta _j\varphi )(a_0,a_1,\mathrm{},a_{k+1})=`$ $`\varphi (a_0,\mathrm{},a_ja_{j+1},\mathrm{},a_{k+1})\text{ for }0jn1`$
$`(\delta _n\varphi )(a_0,a_1,\mathrm{},a_{k+1})=`$ $`(1)^{\mathrm{deg}a_{k+1}(\mathrm{deg}a_0+\mathrm{}+\mathrm{deg}a_k)}\varphi (a_{k+1}a_0,a_1,\mathrm{}a_k)`$ (2.10)
$$(\sigma _j\varphi )(a_0,\mathrm{},a_{k1})=\varphi (a_0,\mathrm{},a_j,1,a_{j+1},\mathrm{}a_{k1})$$
(2.11)
and the cyclic action is given by
$$(\tau _k\varphi )(a_0,\mathrm{},a_k)=(1)^{\mathrm{deg}a_k(\mathrm{deg}a_0+\mathrm{}\mathrm{deg}a_{k1})}\varphi (a_k,a_0,\mathrm{},a_{k1})$$
(2.12)
Cohomology of total complex of the bicomplex $``$ (resp. $`_{\text{per}}`$) where we consider only *finite* cochains, is the cyclic (resp. periodic cyclic) cohomology of $`\mathrm{\Omega }^{}`$, for which we use notation $`HC^{}(\mathrm{\Omega }^{})`$ (resp. $`HP^{}(\mathrm{\Omega }^{})`$).
Suppose now that $`\mathrm{\Omega }^{}`$ is a differential graded (DG) algebra with the differential of degree $`1`$. We say that $`\varphi C^k(\mathrm{\Omega }^{})`$ has weight $`m`$ if $`\varphi (a_0,a_1,\mathrm{},a_k)=0`$ unless $`\mathrm{deg}a_0+\mathrm{deg}a_1+\mathrm{}+\mathrm{deg}a_k=m`$. We denote by $`C^{k,p}(\mathrm{\Omega }^{})C^k(\mathrm{\Omega }^{})`$ set of weight $`p`$ functionals. Notice that in this case each $`C^k(\mathrm{\Omega }^{})`$ is a complex in its own right, with the grading defined above and the differential $`(1)^kd`$, where we extend $`d`$ to $`C^k(\mathrm{\Omega }^{})`$ by
$$d\varphi (a_0,a_1,\mathrm{},a_k)=\underset{j=0}{\overset{k}{}}(1)^{\mathrm{deg}a_0+\mathrm{}\mathrm{deg}a_{j1}}\varphi (a_0,\mathrm{},da_j,\mathrm{},a_k)$$
(2.13)
Then $`dbbd=0`$, $`dBBd=0`$, and hence in this situation $``$ and $`_{\text{per}}`$ become actually tricomplexes. Cyclic (resp. periodic cyclic) cohomology of the DG algebra $`(\mathrm{\Omega }^{},d)`$ is then defined as the cohomology of the total complex of the tricomplex $``$ (resp. $`_{\text{per}}`$) where we consider only *finite* cochains. Notations for the cyclic and periodic cyclic cohomologies are $`HC^{}\left((\mathrm{\Omega }^{},d)\right)`$ and $`HP^{}\left((\mathrm{\Omega }^{},d)\right)`$.
One can show that cyclic cohomology can be computed by the *normalized* complex, i.e. one where cochains satisfy
$$\varphi (a_0,a_1,\mathrm{},a_k)=0\text{ if }a_i=1,i1$$
(2.14)
We will need the following result about the cyclic cohomology.
###### Theorem 1.
Let $`𝒜=\mathrm{\Omega }^0`$ be the $`0`$-degree part of $`\mathrm{\Omega }^{}`$ (which we consider as a trivially graded algebra with the zero differential). We then have a natural map of (total) complexes $`I:_{\text{per}}(𝒜)_{\text{per}}\left((\mathrm{\Omega }^{},d)\right)`$ (extension of polilinear forms by $`0`$ to $`\mathrm{\Omega }^{}`$). Then the induced map in cohomology is an isomorphism. $`HP^{}(𝒜)HP^{}\left((\mathrm{\Omega }^{},d)\right)`$.
To prove the theorem and to write an explicit formula for the map $`R:_{\text{per}}\left((\mathrm{\Omega }^{},d)\right)_{\text{per}}(𝒜)`$, inducing the inverse isomorphism in the periodic cyclic cohomology, we need the following fact (Rinehart formula) which we use as stated in .
Let $`D`$ be a derivation of the graded algebra $`\mathrm{\Omega }^{}`$ of degree $`\mathrm{deg}D`$, i.e. a linear map $`D:\mathrm{\Omega }^{}\mathrm{\Omega }^{+\mathrm{deg}D}`$ satisfying
$$D(ab)=(Da)b+(1)^{\mathrm{deg}D\mathrm{deg}a}aD(b)$$
(2.15)
It defines an operator on the complex $`(\mathrm{\Omega }^{})`$, by
$$_D\varphi (a_0,a_1,\mathrm{},a_k)=\underset{i=0}{\overset{k}{}}(1)^{\mathrm{deg}D(a_0+\mathrm{}+a_{i1})}\varphi (a_0,\mathrm{},D(a_i),\mathrm{},a_k)$$
(2.16)
which commutes with $`b`$, $`B`$. The action of this operator on the periodic cyclic bicomplex is homotopic to zero, with the homotopy constructed as follows. Define operators $`e_D:C^{k1}(\mathrm{\Omega }^{})e_D:C^k(\mathrm{\Omega }^{})`$, $`E_D:e_D:C^{k+1}(\mathrm{\Omega }^{})e_D:C^k(\mathrm{\Omega }^{})`$ by
$`e_D\varphi (a_0,a_1,\mathrm{}a_k)`$ $`=(1)^{k+1}\varphi (D(a_k)a_0,a_1\mathrm{},a_{k1})`$ (2.17)
$`E_D\varphi (a_0,a_1,\mathrm{}a_k)`$ $`={\displaystyle \underset{1ijk}{}}(1)^{ik+1}\varphi (1,a_i,a_{i+1},\mathrm{},a_{j1},Da_j,\mathrm{},a_k,a_0,\mathrm{})`$ (2.18)
Then
$$[b+B,e_D+E_D]=_D$$
(2.19)
We now proceed with the proof of the Theorem 1
###### Proof of the Theorem 1.
Consider the derivation $`D`$ of degree $`0`$ given by $`Da=(\mathrm{deg}a)a`$. On the polilinear form of weight $`m`$ $`D`$ acts by $`m`$. Define the homotopy $`h`$ to be $`\frac{1}{m}(e_D+E_D)`$ on the forms of weight $`m>0`$ and $`0`$ on the forms of weight $`0`$. We define map of complexes $`R:_{\text{per}}\left((\mathrm{\Omega }^{},d)\right)_{\text{per}}(𝒜)`$ by
$$R\varphi =c_{k,m}(dh)^m\varphi \text{ for }\varphi C^{k,m}$$
(2.20)
where
$$c_{k,m}=(1)^{km+\frac{m^2m}{2}}$$
(2.21)
This is a map of (total) complexes. Indeed, using identities $`(b+B)h+h(b+B)=id`$ and $`(b+B)dd(b+B)=0`$ we have:
$$\begin{array}{c}(b+B)R\varphi R(b+B+(1)^kd)\varphi =\hfill \\ \hfill c_{k,m}\left(((b+B)(dh)^m(1)^m(dh)^m(b+B))+(1)^m(dh)^{m1}d\right)\varphi =\\ \hfill c_{k,m}\left((1)^m(dh)^{m1}d+(1)^m(dh)^{m1}d\right)\varphi =0\end{array}$$
(2.22)
It is clear that
$$RI=id$$
(2.23)
As for $`IR`$ we have
$$IR=id\left(H+H\right)$$
(2.24)
where $`=b+B\pm d`$ – total differential in the complex $`_{\text{per}}\left((\mathrm{\Omega }^{},d)\right)`$, $`\pm d`$ being $`(1)^kd`$ on $`C^{k,m}`$, and the homotopy $`H`$ is given by the formula
$$H\varphi =\underset{j=0}{\overset{m1}{}}c_{k,j}h(dh)^j\varphi \text{ for }\varphi C^{k,m}$$
(2.25)
This equality is also verified by direct computation.
Indeed, we have, for actions on $`C^{k,m}`$
$`H(b+B)={\displaystyle \underset{j=0}{\overset{m1}{}}}c_{k,j}((dh)^j(b+B)h(dh)^j)+{\displaystyle \underset{j=1}{\overset{m1}{}}}c_{k,j}(1)^j(hd)^j`$ (2.26)
$`(b+B)H={\displaystyle \underset{j=0}{\overset{m1}{}}}c_{k,j}(b+B)h(dh)^j`$ (2.27)
$`H(\pm d)={\displaystyle \underset{j=0}{\overset{m2}{}}}(1)^kc_{k,j}(hd)^{j+1}={\displaystyle \underset{j=1}{\overset{m1}{}}}(1)^kc_{k,j1}(hd)^j`$ (2.28)
$`(\pm d)H={\displaystyle \underset{j=0}{\overset{m1}{}}}(1)^{kj1}c_{k,j}(dh)^{j+1}={\displaystyle \underset{j=1}{\overset{m1}{}}}(1)^{kj}c_{k,j1}(dh)^j`$ (2.29)
and adding these equalities we get the desired result. ∎
## 3 Cyclic complex for differential graded Hopf algebras
In this section we reproduce Connes-Moscovici’s construction of the cyclic module of a Hopf algebra (cf. ) in the differential graded context.
Let us start with the graded Hopf algebra $`^{}`$. We need to fix a modular pair, i.e. a homomorphism $`\delta :^{}`$ and a group-like element $`\sigma ^0`$. Using the standard notations for the coproduct and antipode, define the twisted antipode $`\stackrel{~}{S}_\delta `$ by
$$\stackrel{~}{S}_\delta (h)=S(h_{(0)})\delta (h_{(1)})$$
(3.1)
Suppose that the following condition holds:
$$\left(\sigma ^1\stackrel{~}{S}_\delta \right)^2=id$$
(3.2)
Then Connes and Moscovici show that one can define a cyclic object $`\left(^{}\right)^{\mathrm{}}=\left\{\left(^{}\right)^n\right\}_{n1}`$ as follows. Face and degeneracy operators are given by
$`\delta _0(h^1\mathrm{}h^{n1})`$ $`=1h^1\mathrm{}h^{n1}`$
$`\delta _j(h^1\mathrm{}h^{n1})`$ $`=h^1\mathrm{}\mathrm{\Delta }h^j\mathrm{}h^n\text{ for }1jn1,`$
$`\delta _n(h^1\mathrm{}h^{n1})`$ $`=h^1\mathrm{}h^{n1}\sigma `$
$`\sigma _i(h^1\mathrm{}h^{n+1})`$ $`=h^1\mathrm{}\epsilon (h^{i+1})\mathrm{}h^{n+1}`$ (3.3)
The cyclic operators are given by
$$\begin{array}{c}\tau _n(h^1\mathrm{}h^{n+1})=\hfill \\ \hfill (1)^{\underset{j>i0}{}\mathrm{deg}h_i^1\mathrm{deg}h^j}\left(\stackrel{~}{S}h^1\right)_{(0)}h^2\mathrm{}\left(\stackrel{~}{S}h^1\right)_{(n2)}h^n\left(\stackrel{~}{S}h^1\right)_{(n1)}\sigma \end{array}$$
(3.4)
where
$$\left(\mathrm{\Delta }^{n1}\stackrel{~}{S}h^1\right)=\left(\stackrel{~}{S}h^1\right)_{(0)}\mathrm{}\left(\stackrel{~}{S}h^1\right)_{(n1)}$$
It is verified in that the operations above indeed define a structure of a cyclic module.
Hence we can define cyclic and periodic cyclic complexes of this cyclic module. Suppose now that our Hopf algebra $`^{}`$ is a DG Hopf algebra with the differential $`d`$ of degree $`1`$. Then complexes $``$ and $`_{\text{per}}`$ have an extra differential defined to be $`(1)^nd`$ on $`\left(^{}\right)^n`$ where we extend $`d`$ by
$$d(h^1h^2\mathrm{}h^n)=\underset{i=1}{\overset{n}{}}(1)^{\mathrm{deg}h^1+\mathrm{}\mathrm{deg}h^{i1}}h^1h^2\mathrm{}dh^i\mathrm{}h^n$$
(3.5)
We consider the total complexes of the *finite* cochains in the resulting tricomplexes, and define cyclic and peridic cyclic cohomology of DG Hopf algebra as cohomology of these complexes.
Suppose now that we are given an action $`\pi `$ of a differential graded Hopf algebra $`^{}`$ on $`\mathrm{\Omega }^{}`$, which agrees with the differential graded structures on $`^{}`$ and $`\mathrm{\Omega }^{}`$, i.e. in addition to the general properties of Hopf algebra action we have
$`\mathrm{deg}\pi (h)(a)`$ $`=\mathrm{deg}h+\mathrm{deg}a`$ (3.6)
$`d\left(\pi (h)(a)\right)`$ $`=\pi (dh)(a)+(1)^{\mathrm{deg}h}\pi (h)(da)`$ (3.7)
where $`h^{}`$, $`a\mathrm{\Omega }^{}`$. We will often omit $`\pi `$ from our notations and write just $`h(a)`$ if it is clear what action we are talking about.
Suppose that $`{\displaystyle }`$ is a closed graded $`\sigma `$-trace on $`\mathrm{\Omega }^{}`$, $`\delta `$-invariant under the action of $`^{}`$, i.e.
$`{\displaystyle }\pi (h)(a)b`$ $`={\displaystyle }a\pi \left(\stackrel{~}{Sh}\right)b`$ (3.8)
$`{\displaystyle }ab`$ $`={\displaystyle }b\pi (\sigma )(a)`$ (3.9)
Then one has a map of cyclic modules $`\chi _\pi :(^{})^{\mathrm{}}(\mathrm{\Omega }^{})^{\mathrm{}}`$, given by
$$\chi _\pi (h^1h^2\mathrm{}h^k)(a_0,a_1,\mathrm{}a_k)=\lambda a_0\pi (h^1)(a_1)\mathrm{}\pi (h^k)(a_k)$$
(3.10)
where
$$\lambda =(1)^{\underset{j>i0}{}\mathrm{deg}h^j\mathrm{deg}a_i}$$
This map also commutes with the differential $`d`$, and hence induces a characteristic map $`\chi _\pi :(^{},d)(\mathrm{\Omega }^{},d)`$, as well as corresponding maps in cohomology.
## 4 Properties of the characteristic map
We will consider now two actions of $`^{}`$ on $`\mathrm{\Omega }^{}`$ which are conjugated by the inner automorphism. We will work in the assumption that $`\mathrm{\Omega }^{}`$ is unital, indicating the changes which need to be made in the nonunital case in the Remark 7. More precisely, let $`\rho ^+`$ and $`\rho ^{}`$ be two degree-preserving linear maps from $`^{}`$ to $`\mathrm{\Omega }^{}`$, which commute with the differentials. We suppose that they are inverse to each other with respect to convolution:
$$\rho ^+(h_{(0)})\rho ^{}(h_{(1)})=\epsilon (h)1$$
(4.1)
and satisfy cocycle identities:
$`\rho ^+(hg)`$ $`={\displaystyle \rho ^+(h_{(0)})\pi (h_{(1)})(\rho ^+(g))}`$ (4.2)
$`\rho ^{}(gh)`$ $`={\displaystyle \pi (h_{(0)})(\rho ^{}(g))\rho ^{}(h_{(1)})}`$ (4.3)
$`\rho ^+(1)`$ $`=\rho ^{}(1)=\rho ^+(\sigma )=\rho ^{}(\sigma )=1`$ (4.4)
Then one can define a new action $`\pi ^{}`$ of $`^{}`$ on $`\mathrm{\Omega }^{}`$ by
$$\pi ^{}(h)(a)=(1)^{\mathrm{deg}h_{(2)}\mathrm{deg}a}\rho ^+(h_{(0)})\pi (h_{(1)})(a)\rho ^{}(h_{(2)})$$
(4.5)
###### Lemma 2.
Equation (4.5) defines an action of the DG Hopf algebra $`^{}`$ on the DG algebra $`\mathrm{\Omega }^{}`$.
###### Proof.
We check all the required properties. First
$$\begin{array}{c}\pi ^{}(h)(ab)=(1)^{\mathrm{deg}ab\mathrm{deg}h_{(2)}}\rho ^+(h_{(0)})\pi (h_{(1)})(ab)\rho ^{}(h_{(2)})=\hfill \\ \hfill (1)^{(\mathrm{deg}a+\mathrm{deg}b)\mathrm{deg}h_{(3)}}(1)^{\mathrm{deg}a\mathrm{deg}h_{(2)}}\rho ^+(h_{(0)})\pi (h_{(1)})(a)\pi (h_{(2)})(b)\rho ^{}(h_{(3)})=\\ \hfill (1)^{(\mathrm{deg}a+\mathrm{deg}b)\mathrm{deg}h_{(5)}}(1)^{\mathrm{deg}a(\mathrm{deg}h_{(2)}+\mathrm{deg}h_{(3)}+\mathrm{deg}h_{(4)})}\\ \hfill \rho ^+(h_{(0)})\pi (h_{(1)})(a)\rho ^{}(h_{(2)})\rho ^+h_{(3)}\pi (h_{(4)})(b)\rho ^{}(h_{(5)})=\\ \hfill (1)^{\mathrm{deg}a\mathrm{deg}h_{(1)}}\pi ^{}(h_{(0)})(a)\pi ^{}(h_{(1)})(b)\end{array}$$
(4.6)
Then
$$\begin{array}{c}\pi ^{}(hg)(a)=\hfill \\ \hfill (1)^{\mathrm{deg}a(\mathrm{deg}h_{(2)}+\mathrm{deg}g_{(3)})}\rho ^+(h_{(0)}g_{(0)})\pi (h_{(1)}g_{(1)})(a)\rho ^{}(h_{(2)}g_{(2)})=\\ \hfill (1)^{\mathrm{deg}a(\mathrm{deg}h_{(3)}+\mathrm{deg}g_{(2)}+\mathrm{deg}h_{(4)})}\\ \hfill \rho ^+(h_{(0)})\pi (h_{(1)})(\rho ^+(g_{(0)}))\pi (h_{(2)}g_{(1)})(a)\pi (h_{(3)})\rho ^{}(g_{(2)})\rho ^{}(h_{(4)})=\\ \hfill \pi ^{}(h)\left(\pi ^{}(g)(a)\right)\end{array}$$
(4.7)
Also
$`\pi ^{}(h)(1)=`$ $`{\displaystyle \rho ^+(h_{(0)})h_{(1)}(1)\rho ^{}(h_{(3)})}=\epsilon (h)1`$ (4.8)
$`\pi ^{}(1)(a)=`$ $`\rho ^+(1)a\rho ^{}(1)=a`$ (4.9)
and
$$\begin{array}{c}d\left(\pi ^{}(h)(a)\right)=d(1)^{\mathrm{deg}a\mathrm{deg}h_{(3)}}\rho ^+(h_{(0)})\pi (h_{(1)})(a)\rho ^{}(h_{(2)})=\hfill \\ \hfill (1)^{\mathrm{deg}a\mathrm{deg}h_{(3)}}\rho ^+(dh_{(0)})\pi (h_{(1)})(a)\rho ^{}(h_{(2)})+\\ \hfill (1)^{\mathrm{deg}a\mathrm{deg}h_{(3)}+\mathrm{deg}h_{(0)}}\\ \hfill \rho ^+(h_{(0)})\left(\pi (dh_{(1)})(a)+(1)^{\mathrm{deg}h_{(1)}}\pi (h_{(1)})(da)\right)\rho ^{}(h_{(2)})+\\ \hfill (1)^{\mathrm{deg}a\mathrm{deg}h_{(3)}+\mathrm{deg}h_{(0)}+\mathrm{deg}h_{(1)}+\mathrm{deg}a}\rho ^+(h_{(0)})\pi (h_{(1)})(a)\rho ^{}(dh_{(2)})=\\ \hfill \pi ^{}(dh)(a)+(1)^{\mathrm{deg}h}\pi ^{}(h)(da)\end{array}$$
(4.10)
Suppose now that $`{\displaystyle }`$ is the closed $`\delta `$-invariant $`\sigma `$-trace for both actions $`\pi `$ and $`\pi ^{}`$. In this case we have two characteristic maps $`\chi _\pi `$ and $`\chi _\pi ^{}`$ from $`(^{},d)`$ to $`(\mathrm{\Omega }^{},d)`$. Then we have the following
###### Proposition 3.
Let $`\pi `$ and $`\pi ^{}`$ be two actions of $`^{}`$ on $`\mathrm{\Omega }^{}`$, conjugated by inner automorphisms, and suppose that they both satisfy conditions (3.8),(3.9). Let $`\chi _\pi `$, $`\chi _\pi ^{}`$ be the corresponding characteristic maps. Then induced maps in cohomology $`HC^{}(^{},d)HC^{}\left(\mathrm{\Omega }^{}\right)`$ coincide.
###### Proof.
Let $`M_2(\mathrm{\Omega }^{})=\mathrm{\Omega }^{}M_2()`$ be the differential graded algebra of $`2\times 2`$ matrices over the algebra $`\mathrm{\Omega }^{}`$.We can define an action $`\pi _2`$ of $`^{}`$ on $`\mathrm{\Omega }^{}M_2()`$ by $`\pi _2(h)(am)=\pi (h)(a)m`$, where $`h^{}`$, $`a\mathrm{\Omega }^{}`$, $`mM_2()`$. Put now
$$\rho _2^+(h)=\left(\begin{array}{ccc}& \rho ^+(h)& 0\\ & 0& \epsilon (h)\end{array}\right)\rho _2^{}(h)=\left(\begin{array}{ccc}& \rho ^{}(h)& 0\\ & 0& \epsilon (h)\end{array}\right)$$
(4.11)
It is easy to see that $`\rho _2^+`$, $`\rho _2^{}`$ satisfy equations (4.1)-(4.4), and hence we can twist the action $`\pi _2`$ by $`\rho _2^+`$, $`\rho _2^{}`$ to define a new action $`\pi _2^{}`$, as in (4.5).
Consider now the linear functional $`{\displaystyle }_2`$ on $`M_2(\mathrm{\Omega }^{})`$ defined by
$$_2(am)=\left(a\right)\left(\mathrm{tr}m\right)$$
(4.12)
Then $`{\displaystyle }_2`$ is a closed graded $`\delta `$-invariant $`\sigma `$-trace on $`M_2(\mathrm{\Omega }^{})`$ with respect to the action $`\pi _2^{}`$. Hence we can define the characteristic map $`\chi _{\pi _2^{}}:(^{},d)(M_2(\mathrm{\Omega }^{}),d)`$
Consider now two imbeddings $`i`$, $`i^{}:\mathrm{\Omega }^{}M_2(\mathrm{\Omega }^{})`$ defined by
$$i(a)=\left(\begin{array}{ccc}& 0& 0\\ & 0& a\end{array}\right)i^{}(a)=\left(\begin{array}{ccc}& a& 0\\ & 0& 0\end{array}\right)$$
(4.13)
It is easy to see that $`i^{}\chi _{\pi _2^{}}=\chi _\pi `$ and $`(i^{})^{}\chi _{\pi _2^{}}=\chi _\pi ^{}`$. Now to finish the proof it is enough to recall the well-known fact that $`i`$ and $`i^{}`$ induce the same map in cyclic cohomology. Since we will later need an explicit homotopy between $`\chi _\pi `$ and $`\chi _\pi ^{}`$ we give the proof below.
Put $`u_t=\mathrm{exp}t\left(\begin{array}{ccc}0& 1& \\ 1& 0& \end{array}\right)=\left(\begin{array}{ccc}\mathrm{cos}t& \mathrm{sin}t& \\ \mathrm{sin}t& \mathrm{cos}t& \end{array}\right)`$ Put $`i_t(a)=u_ti(a)u_t^1`$. Notice that $`i_0=i`$, $`i_{\pi /2}=i^{}`$. Consider the family of maps $`i_t^{}:\left(M_2(\mathrm{\Omega }^{})\right)\left(\mathrm{\Omega }^{}\right)`$. Since we have $`\frac{d}{dt}i_t(a)=[g,i_t(a)]`$, where $`g=\left(\begin{array}{ccc}0& 1& \\ 1& 0& \end{array}\right)`$ these maps satisfy $`\frac{d}{dt}i_t^{}=i_t^{}L_g`$, where $`L_g:C^k(M_2(\mathrm{\Omega }^{}))C^k(M_2(\mathrm{\Omega }^{}))`$ is the operator defined by
$$L_g\varphi (x_0,\mathrm{},x_k)=\underset{j=0}{\overset{k}{}}\varphi (x_0,\mathrm{},[g,x_j],\mathrm{},x_k)$$
Define also an operator $`I_g:C^k(M_2(\mathrm{\Omega }^{}))C^{k1}(M_2(\mathrm{\Omega }^{}))`$ by
$$I_g\varphi (x_0,\mathrm{},x_{k1})=\underset{j=0}{\overset{k1}{}}\varphi (x_0,\mathrm{},x_j,g,x_{j+1},\mathrm{},x_{k1})$$
(4.14)
Then it is easy to verify that $`[b,I_g]=L_g`$, $`[B,I_g]=0`$ and $`[d,I_g]=0`$. Hence $`L_g=I_g+I_g`$, $`=\pm d+b+B`$. We conclude that $`i_1^{}i_0^{}=K+K`$, where the homotopy $`K`$ is given by $`K\varphi =_0^{\pi /2}i_t^{}I_g`$.
Hence
$$\chi _\pi ^{}\chi _\pi =H+H$$
(4.15)
where $`H=K\chi _{\pi _2^{}}`$
Now note that the complex $`\left(^{}\right)`$ has a natural weight filtration by subcomplexes $`F^l(^{},d)`$, where
$$F^l(^{},d)=\{\alpha _1\alpha _2\mathrm{}\alpha _j|\mathrm{deg}\alpha _1+\mathrm{deg}\alpha _2+\mathrm{}\mathrm{deg}\alpha _jl$$
(4.16)
Suppose now that $`{\displaystyle }`$ has weight $`q`$, i.e.
$$a=0\text{ if }\mathrm{deg}aq$$
(4.17)
Notice that in this case $`\chi _\pi `$ reduces the total degree by $`q`$ Then following then proposition is clear:
###### Proposition 4.
The characteristic map is $`0`$ on $`F^l\left(^{}\right)`$ for $`l>q`$.
Let $`(^{},d)_l`$ denote the truncated cyclic bicomplex:
$$(^{},d)_l=(^{},d)/F^{l+1}(^{},d)$$
(4.18)
Then we immediately have the following
###### Corollary 5.
The characteristic map $`\chi _\pi `$ defined in (3.10) induces the map from the complex $`(^{},d)_q`$ to the cyclic complex of the differential graded algebra $`\mathrm{\Omega }^{}`$.
This new map will also be denoted $`\chi _\pi `$. We use the notation
$$HC^{}(^{},d)_l=H^{}\left((^{},d)_l\right)$$
(4.19)
for the cohomology of the complex $`(^{},d)_l`$. The explicit form of the homotopy in the Proposition 3 now implies the following
###### Corollary 6.
Suppose in addition to the conditions of the Proposition 3 that (4.17) is satisfied. Then the two maps in cohomology induced by $`\chi _\pi ,\chi _\pi ^{}:(^{},d)_q\left(\mathrm{\Omega }^{}\right)`$ are the same.
###### Proof.
We use the notations of the proof of the Proposition 3. There we established that $`\chi _\pi ^{}\chi _\pi =H+H`$. We need only to verify that $`H`$ is well defined on the quotient complex $`(^{},d)_q`$. But since $`H=K\chi _{\pi _2^{}}`$, and $`\chi _{\pi _2^{}}`$ is easily seen to be $`0`$ on $`F^{q+1}(^{},d)`$, the result follows. ∎
###### Remark 7.
We worked above in the assumption that the DG algebra $`\mathrm{\Omega }^{}`$ is unital. If this is not the case some changes should be made. First of all, $`\rho ^+(h)`$, $`\rho ^{}(h)`$ now don’t have to be elements of the algebra, but rather multipliers, such that $`\pi ^{}`$ defined in (4.5) is a Hopf action. Moreover, we need to require that if $`m`$ is such a multiplier, then $`{\displaystyle }ma={\displaystyle }am`$ $`a\mathrm{\Omega }^{}`$. Then the the Proposition 3 remains true. Characteristic maps in this case take values in the $``$ complex of the algebra $`\mathrm{\Omega }^{}`$ with unit adjoined. The homotopy between two characteristic maps is still given by explicit formula (4.15), which continues to make sense in the nonunital situation. Indeed, it is sufficient to define $`I_g\chi _{\pi _2^{}}`$, which can be defined by the same formula as above, provided one treats $`g`$ and $`\pi _2^{}(h)(g)`$ as multipliers of the algebra $`M_2\left(\mathrm{\Omega }^{}\right)`$, with $`\pi _2^{}(g)`$ defined to be $`\left(\begin{array}{ccc}& 0& \rho ^+(h)\\ & \rho ^{}(h)& 0\end{array}\right)`$
Finally, we collect all the information we will need to use in the next sections.
###### Theorem 8.
Let $`(\mathrm{\Omega }^{},d)`$ be a differential graded algebra, and $`{\displaystyle }`$ a linear functional on $`\mathrm{\Omega }^{}`$ of weight $`q`$, and let $`𝒜=\mathrm{\Omega }^0`$ be the degree $`0`$ part of $`\mathrm{\Omega }^{}`$. Let $`\pi `$ be an action of the DG Hopf algebra $`^{}`$ act on the DGA $`\mathrm{\Omega }^{}`$. Suppose that $`{\displaystyle }`$ is a $`\delta `$-invariant $`\sigma `$-trace with respect to $`\pi `$. Then characteristic map (3.10) defines a map in cohomology $`HC^i(^{})_qHP^{iq}(𝒜)`$. Suppose now that $`\pi ^{}`$ is another action of $`^{}`$ on $`\mathrm{\Omega }^{}`$, obtained from $`\pi `$ by twisting by a cocycle (4.5). Then if $`{\displaystyle }`$ is a $`\delta `$-invariant $`\sigma `$-trace with respect to $`\pi ^{}`$, the maps in cohomology $`HC^i(^{})_qHP^{iq}(𝒜)`$ induced by $`\chi _\pi `$, $`\chi _\pi ^{}`$ are the same.
## 5 Secondary characteristic classes
Let $`M`$ be a manifold, and let $`\mathrm{\Gamma }`$ be a discrete pseudogroup of diffeomorphisms of $`M`$, acting from the right.
By this we mean a set $`\mathrm{\Gamma }`$ such that every element $`g`$ of $`\mathrm{\Gamma }`$ defines a local diffeomorphism of $`M`$, i.e. diffeomorphism $`g:\mathrm{Dom}g\mathrm{Ran}g`$, where $`\mathrm{Dom}g`$, $`\mathrm{Ran}gM`$ – open subsets of $`M`$, and that we have partially defined operations of composition and inverse such that
1. If $`g\mathrm{\Gamma }`$ then $`g^1:\mathrm{Ran}g\mathrm{Dom}g`$ is also in $`\mathrm{\Gamma }`$.
2. If $`g_1,g_2\mathrm{\Gamma }`$ then $`g_1g_2`$ with domain $`g_1^1\left(\mathrm{Dom}g_2\mathrm{Ran}g_1\right)`$ and range $`g_2\left(\mathrm{Dom}g_2\mathrm{Ran}g_1\right)`$ is in $`\mathrm{\Gamma }`$.
3. $`id:MM`$ is in $`\mathrm{\Gamma }`$.
Note that we use a wide definition of pseudogroups, and do not include any saturation axioms.
Let $`E`$ be a trivial vector bundle on $`M`$, equivariant with respect to the action of $`\mathrm{\Gamma }`$. In other words, every $`g\mathrm{\Gamma }`$ defines for every $`x\mathrm{Dom}g`$ a linear map $`E_xE_{xg}`$.
For the rest of the paper we suppose the following:
*If $`g_1`$ and $`g_2\mathrm{\Gamma }`$ are such that they induce the same diffeomorphisms and the same action on the bundle, then $`g_1=g_2`$*.
With this data one can associate the following groupoid $`𝒢`$: the objects are the points of $`M`$ and the morphisms $`xy`$, $`x`$, $`yM`$ are given by $`g\mathrm{\Gamma }`$ such that $`g(x)=y`$, with the composition given by the product in $`\mathrm{\Gamma }`$. Let $`𝒜`$ denote the convolution algebra of this groupoid, i.e. the cross-product $`C_0^{\mathrm{}}(M)\mathrm{\Gamma }`$. Let $`\mathrm{\Omega }^{}=(\mathrm{\Omega }^{}(M)\mathrm{\Gamma },d)`$ denote the differential graded algebra of forms on $`𝒢`$ with the convolution product, where the differential $`d`$ is the de Rham differential. We will use the usual cross-product notations $`\omega U_g`$ for the elements of this algebra, where $`\omega \mathrm{\Omega }^{}(M)`$, $`g\mathrm{\Gamma }`$. Since $`\mathrm{\Gamma }`$ is, in general a pseudogroup, we suppose that
$$\mathrm{supp}\omega \mathrm{Dom}g$$
(5.1)
Fix a trivialization of $`E`$. The action of $`\mathrm{\Gamma }`$ on the bundle defines then a homomorphism
$$h:𝒢GL_n()$$
(5.2)
Let $`\left(GL_n()\right)`$ denote the differential graded Hopf algebra of the forms on $`GL_n()`$, with the product given exterior multiplication, coproduct, antipode and counit induced respectively by the product $`GL_n()\times GL_n()GL_n()`$, inverse $`GL_n()GL_n()`$ and the inclusion $`1GL_n()`$. The differential is given by the de Rham differential on forms.
We now show that the map (5.2) allows one to define an action of $`\left(GL_n()\right)`$ on $`\mathrm{\Omega }^{}`$.
###### Proposition 9.
The map $`\left(GL_n()\right)\mathrm{\Omega }^{}\mathrm{\Omega }^{}`$ given by
$$\pi (\alpha )(\omega )=h^{}(\alpha )\omega $$
(5.3)
where $`\alpha \left(GL_n()\right)`$, $`\omega \mathrm{\Omega }^{}`$ defines an action of the differential graded Hopf algebra $`\left(GL_n()\right)`$ on the differential graded algebra $`\mathrm{\Omega }^{}`$.
###### Proof.
We have:
$$\pi (\alpha _1\alpha _2)(\omega )=h^{}(\alpha _1\alpha _2)\omega =h^{}(\alpha _1)h^{}(\alpha _2)\omega =\pi (\alpha _1)\left(\pi (\alpha _2)(\omega )\right)$$
(5.4)
Next, if we write $`\mathrm{\Delta }\alpha =\underset{k}{}\alpha _{(0)}\alpha _{(1)}`$ we have:
$$\begin{array}{c}\pi (\alpha )(\omega _0\omega _1)(g)=h^{}(\alpha )(g)\omega _0\omega _1(g)=h^{}(\alpha )(g)\underset{g_0g_1=g}{}\omega _0(g_0)\omega _1^{g_0}(g_1)=\hfill \\ \hfill \underset{g_0g_1=g}{}\underset{k}{}h^{}(\alpha _{(0)})(g_0)h^{}(\alpha _{(1)})^{g_0}(g_1)\omega _0(g_0)\omega _1^{g_0}(g_1)=\\ \hfill \underset{g_0g_1=g}{}\underset{k}{}(1)^{\mathrm{deg}\omega _0\mathrm{deg}\alpha _{(0)}}h^{}(\alpha _{(0)})(g_0)\omega _0(g_0)h^{}(\alpha _{(2)})(g_1)\omega _1^{g_0}(g_1)=\\ \hfill \underset{k}{}(1)^{\mathrm{deg}\omega _0\mathrm{deg}\alpha _{(0)}}\pi (\alpha _{(0)})(\omega _0)\pi (\alpha _{(1)})(\omega _1)\end{array}$$
(5.5)
Also, if $`M`$ is compact the algebra $`\mathrm{\Omega }^{}`$ has a unit given by the function
$$e(g)=\{\begin{array}{cc}1\hfill & \text{ if }g\text{ is a unit}\hfill \\ 0\hfill & \text{ otherwise}\hfill \end{array}$$
(5.6)
Then
$$\pi (\alpha )(e)=h^{}(\alpha )e=\epsilon (\alpha )e$$
(5.7)
Finally, we have
$$\begin{array}{c}d(\pi (\alpha (\omega )))=d(h^{}(\alpha )\omega ))=\hfill \\ \hfill h^{}(d\alpha )\omega +(1)^{\mathrm{deg}\alpha }h^{}(\alpha )d\omega =\pi ((d\alpha ))(\omega )+(1)^{\mathrm{deg}\alpha }\pi (\alpha )(d\omega )\end{array}$$
(5.8)
We now have a natural inclusion $`i:M𝒢`$ as the space of units. Then we define a graded trace $`{\displaystyle }`$ on $`\mathrm{\Omega }^{}`$ by
$$\omega =\underset{M}{}i^{}\omega $$
(5.9)
###### Proposition 10.
The graded trace $`{\displaystyle }`$ is closed under the de Rham differential and is invariant under the action of $``$, i.e.
$`{\displaystyle }d\omega =0`$ (5.10)
$`{\displaystyle }\alpha (\omega )=\epsilon (\alpha ){\displaystyle }\omega `$ (5.11)
###### Proof.
The first identity is clear, the second follows from the fact that $`hi:MGL_n()`$ is a constant map, taking the value 1. ∎
Hence we have a map $`_q(\left(GL_n()\right),d)(\mathrm{\Omega }^{},d)`$ where $`q=dimM`$, which also gives us a map
$$\chi :HC_q^{}(\left(GL_n()\right),d)HP^{}\left(C_0^{\mathrm{}}(M)\mathrm{\Gamma }\right)$$
(5.12)
Definition of the action of $`\left(GL_n()\right)`$ on $`\mathrm{\Omega }^{}`$, and hence definitions of the map given by (5.12) apriori depends on the choice of trivialization of $`E`$, but we will now show that this map is independent of the choice of trivialization.
###### Proposition 11.
Suppose we use another trivialization of $`E`$ to define an action of $`\left(GL_n()\right)`$ on $`\mathrm{\Omega }^{}`$. Then the two actions are conjugated by the inner automorphisms.
###### Proof.
Let us chose another trivialization of the bundle $`E`$, and let $`U(x)`$ $`xM`$ be a transition matrix between the two bases of the fiber $`E_x`$. Then we have a new map $`h^{}:𝒢GL_n()`$, related to $`h`$ by
$$h^{}(\gamma )=U(s(\gamma ))h(\gamma )U^1(r(\gamma ))$$
(5.13)
Let $`\pi ^{}`$ denote the action corresponding to the map $`h^{}`$.
Consider now the pull-back $`U^{}:\mathrm{\Omega }^{}(GL_n())\mathrm{\Omega }^{}(M)`$ as a map $`\left(GL_n()\right)\mathrm{\Omega }^{}`$, where we consider forms on $`M`$ as the form on $`𝒢`$ which is $`0`$ outside the space of units. When $`M`$ is not compact, we obtain not an element in algebra, but rather a multiplier. Hence we see that if we define
$$\rho ^+(\alpha )=U^{}(\alpha )$$
(5.14)
and
$$\rho ^{}(\alpha )=(U^1)^{}(\alpha )=\rho ^+(S\alpha )$$
(5.15)
we will have
$$\pi ^{}(\alpha )(\omega )=(1)^{\mathrm{deg}\alpha _{(2)}\mathrm{deg}\omega }\rho ^+(\alpha _{(0)})\pi (\alpha _{(1)})(\omega )\rho ^{}(\alpha _{(2)})$$
(5.16)
We can now summarize the results as follows.
###### Theorem 12.
Let $`\mathrm{\Gamma }`$ be a discrete pseudogroup acting on the manifold $`M`$ of dimension $`q`$ by orientation preserving diffeomorphisms. Let $`E`$ be a $`\mathrm{\Gamma }`$-equivariant trivial bundle of rank $`n`$ on $`M`$. Let $`\left(GL_n()\right)`$ be the DG Hopf algebra of the differential forms on the Lie group $`GL_n()`$. Then we have a map
$$\chi :HC_q^i(\left(GL_n()\right),d)HP^{iq}\left(C_0^{\mathrm{}}(M)\mathrm{\Gamma }\right)$$
(5.17)
which is independent of the trivialization of $`E`$. In conjunction with the equation (6.1) below it gives a map
$$\underset{m}{}H^{i2m}\left(W(𝔤𝔩_𝔫,O_n)_q\right)HP^{iq}\left(C_0^{\mathrm{}}(M)\mathrm{\Gamma }\right)$$
(5.18)
## 6 Relation with Weil algebras
In this section we use methods of , and to identify the cyclic cohomology $`HC^{}(\left(GL_n()\right),d)_q`$. It turns out that
$$HC^i(\left(GL_n()\right),d)_q=\underset{m0}{}H^{i2m}\left(W(𝔤𝔩_𝔫,O_n)_q\right).$$
(6.1)
where $`H^{}\left(W(𝔤𝔩_𝔫,O_n)_q\right)`$ is the cohomology of truncated Weil algebra (cf. ). As a matter of fact, one can work with the DG Hopf algebra $`(G)`$ of differential forms on any almost connected Lie group, and the result in this case is
$$HC^i((G),d)_q=\underset{m0}{}H^{i2m}\left(W(𝔤,K)_q\right).$$
(6.2)
where $`K`$ is the maximal compact subgroup of $`G`$. Computation of the Hochschild cohomology of this Hopf algebra is essentially contained in , and with little care using ideas from one recovers the cyclic cohomology.
Let $`G`$ be a Lie group with finitely many connected components. Similarly to the previous section we can define a differential graded Hopf algebra $`(G)`$. Let $`K`$ be the maximal compact subgroup of $`G`$. We will now construct the map of complexes from the truncated relative Weil algebra $`W(𝔤,K)_q`$ to the complex $`_q\left((G)\right)`$.
Let $`NG`$ denote the simplicial manifold with $`NG_p=\underset{p}{\underset{}{G\times G\times \mathrm{}G}}`$ The simplicial structure is given by the face maps
$$_i(g_1,g_2,\mathrm{},g_k)=\{\begin{array}{cc}(g_2,\mathrm{},g_k)\hfill & \text{ if }i=0\hfill \\ (g_1,g_2,\mathrm{}g_ig_{i+1},\mathrm{},g_k)\hfill & \text{ if }1ik1\hfill \\ (g_1,\mathrm{},g_{k1})\hfill & \text{ if }i=k\hfill \end{array}$$
(6.3)
and degeneracy maps
$$\sigma _i(g_1,g_2,\mathrm{},g_k)=(g_1,\mathrm{},g_{i1},1,g_i,\mathrm{},g_k)$$
(6.4)
The geometric realization of this simplicial manifold is the classifying space $`\mathrm{B}G`$. It is a union of manifolds $`\mathrm{\Delta }^p\times NG_p`$ with the modulo the equivalence relation (cf. ).
We will also consider simplicial manifold $`\overline{N}G`$, with $`\overline{N}G_p=\underset{p+1}{\underset{}{G\times G\times \mathrm{}\times G}}`$. The face and degeneracy maps are given by
$$_i(g_0,g_1,g_2,\mathrm{},g_k)=(g_0,\mathrm{}\widehat{g}_i,\mathrm{},g_k)$$
(6.5)
$$\sigma _i(g_0,g_1,g_2,\mathrm{},g_k)=(g_0,\mathrm{},g_{i1},g_i,g_i,g_{i+1},g_k)$$
(6.6)
The geometric realization of this simplicial manifold is $`\mathrm{EG}`$. The map $`pr:\overline{N}GNG`$ given by
$$pr(g_0,g_1,\mathrm{},g_p)=(g_0g_1^1,g_1g_2^1,\mathrm{},g_{p1}g_p^1)$$
(6.7)
defines a simplicial principal $`G`$-bundle $`\mathrm{EG}\mathrm{BG}`$. Simplicial manifolds $`\overline{N}G`$ and $`NG`$ moreover have a cyclic structure, i.e. an action of the cyclic groups $`_{p+1}`$ on the $`p`$-th component, which satisfy all the necessary relations with the face and degeneracy maps. The actions are given on $`\overline{N}G`$ by
$$\tau _p(g_0,g_1,\mathrm{},g_p)=(g_1,g_2,\mathrm{},g_p,g_0)$$
(6.8)
Since the maps $`\tau _p`$ are $`G`$-equivariant, they induce corresponding actions on $`NG`$:
$$\tau _p(g_1,g_2,\mathrm{},g_p)=(g_2,g_3,\mathrm{},g_p,(g_1g_2\mathrm{}g_p)^1)$$
(6.9)
We will identify the $`p`$-cochains for the Hopf algebra $`(G)`$ with the forms on $`NG_p`$. Under this identification the simplicial structure on the Hopf cochains corresponds to the one induced by the simplicial structure on $`NG`$, the de Rham differential on the Hopf cochains corresponds to the de Rham differential on $`NG`$, and cyclic structure on the Hopf cochains is induced by the cyclic structure on $`NG`$. Filtration by the form degree on the Hopf algebra cochains corresponds to the filtration by the form degree on the manifold $`NG`$.
We will now construct the map $`\mu `$ from the complex $`W(𝔤,K)`$ to the simplicial-de Rham complex of forms on $`NG`$, which preserves filtration on these complexes. We do it by constructing the map from $`W(𝔤,K)`$ to the complex of simplicial forms on $`\mathrm{B}G`$, and then applying the integration map. The complex of simplicial forms on $`\mathrm{B}G`$ has a natural bigrading. Let $`\theta `$ be the Maurer-Cartan form on $`G`$. Let $`p_i:\overline{N}G=G^{p+1}G`$ be the projection on $`i`$-th component. Consider on $`\mathrm{EG}_\mathrm{p}`$ the $`𝔤`$ valued differential form $`\omega t_i\theta _i`$, where $`\theta _i=p_i^{}\theta `$. It defines a simplicial connection in the bundle $`\mathrm{EG}\mathrm{BG}`$. The standard construction defines a differential graded algebra homomorphism $`\psi `$ from $`W(𝔤,K)`$ to the complex of $`K`$-basic simplicial forms on $`\mathrm{EG}`$, which we identify with forms on the space $`\mathrm{EG}/\mathrm{K}`$. This space is a bundle over $`\mathrm{B}G`$ with the fiber $`G/K`$. This bundle has a section, which can be explicitly described as follows. Since $`G/K`$ has a natural structure of a manifold of constant negative curvature, for any finite set of points $`x_0`$, $`x_1`$, …, $`x_kG/K`$ one can construct a canonical simplex in $`G/K`$, i.e. a map $`\sigma (x_0,x_1,\mathrm{},x_k):\mathrm{\Delta }^kG/K`$, with vertices $`x_0`$, $`x_1`$, …, $`x_k`$, and this construction agrees with taking faces of a simplex, and is $`G`$-equivariant:
$$\sigma (gx_0,gx_1,\mathrm{},gx_k)(t_0,t_1,\mathrm{},t_k)=g\sigma (x_0,x_1,\mathrm{},x_k)(t_0,t_1,\mathrm{},t_k)$$
(6.10)
Denote by $`\pi `$ the canonical projection $`GG/K`$. Then the section $`s`$ is given by the simplicial map defined by the following formula, where we write just $`\sigma `$ for $`\sigma (\pi (1),\pi (g_1),\pi (g_1g_2),\mathrm{}\pi (g_1\mathrm{}g_k)(t_0,t_1,\mathrm{},t_k)`$:
$$\begin{array}{c}s(g_1,g_2,\mathrm{},g_k;t_0,t_1,\mathrm{}t_k)=\hfill \\ \hfill (\sigma ^1,g_1\sigma ^1,g_1g_2\sigma ^1,\mathrm{}g_1\mathrm{}g_k\sigma ^1;t_0,t_1,\mathrm{}t_k)\end{array}$$
(6.11)
Notice that this section intertwines the actions of the cyclic group on the spaces $`\mathrm{B}G`$ and $`\mathrm{EG}/\mathrm{K}`$.
Consider now the map $`s^{}\psi `$. It is clearly a homomorphism from the differential graded algebra $`W(𝔤,K)`$ to the differential graded algebra of simplicial forms on $`\mathrm{B}G`$. We will show now that it preserves filtration’s on both algebras. First we need the following statement:
###### Lemma 13.
If $`\xi `$ is a horizontal form on $`\mathrm{EG}`$ of the type $`(k,l)`$,then $`s^{}\xi `$ is also of the type $`(k,l)`$.
###### Proof.
Since $`\xi `$ is horizontal, it can be written as a sum of the expressions of the form $`fpr^{}\zeta `$, where $`f`$ is a function, $`pr:\mathrm{EG}\mathrm{BG}`$ – projection, and $`\zeta `$ is a form on $`\mathrm{B}G`$ of the type $`(k,l)`$. Then $`s^{}\left(fpr^{}\zeta \right)=(s^{}f)\zeta `$ is also of the type $`(k,l)`$. ∎
###### Proposition 14.
The homomorphism $`s^{}\psi `$ agrees with filtrations on the Weil algebra and on the forms on $`\mathrm{B}G`$.
###### Proof.
The curvature of the connection $`\omega `$ is a horizontal form $`\mathrm{\Omega }`$ on $`\mathrm{B}G`$, given by
$$\mathrm{\Omega }=dt_i\theta _i+t_id\theta _i+\underset{i<j}{}t_it_j[\theta _i,\theta _j]$$
(6.12)
Hence $`\mathrm{\Omega }`$ has only components of the type $`(1,1)`$ and $`(0,2)`$. The statement of the lemma will then follow from the fact that $`s^{}\mathrm{\Omega }`$ also has only components of the type $`(1,1)`$ and $`(0,2)`$. But this follows from the Lemma 13. ∎
We can now apply the integration map and obtain the map $`\mu `$ from the Weil algebra to the simplicial-de Rham complex of $`NG`$. Since the integration map respects filtrations, the resulting map $`\mu `$ also respects filtrations. We identify the Hochschild complex of $`(G)`$ with the simplicial-de Rham complex of $`NG`$. Results of , imply that this is actually an isomorphism, i.e.
$$HH^i((G),d)_q=H^i\left(W(𝔤,K)_q\right)$$
(6.13)
But since the connection $`\omega `$ and the section $`s`$ are invariant under the cyclic action, the resulting Hochschild cochains are actually cyclic. This implies that Connes’ long exact sequence is equivalent to the collection of short exact sequnces
$$\begin{array}{c}0HC^{i2}((G),d)_q\stackrel{𝑆}{}HC^i((G),d)_q\stackrel{𝐼}{}\hfill \\ \hfill \stackrel{𝐼}{}HH^i((G),d)_q0\end{array}$$
(6.14)
and the map $`I`$ splits. Hence
$$\begin{array}{c}HC^i((G),d)_q=\hfill \\ \hfill \underset{m0}{}HH^{i2m}((G),d)_q=\underset{m0}{}H^{i2m}\left(W(𝔤,K)_q\right).\end{array}$$
(6.15)
Explicitly, maps $`H^{i2m}\left(W(𝔤,K)_q\right)HC^i((G),d)_q`$ are given by $`S^m\mu `$, where we consider $`\mu `$ as a map into cyclic complex.
## 7 Relation with other constructions
Suppose, as before, that we have an orientation-preserving action of discrete group $`\mathrm{\Gamma }`$ on an oriented manifold $`M`$, and an equivariant trivial bundle $`E`$ over $`M`$. Then results of previous sections provide us a map $`H^{}(W(𝔤,O_n))HP^{}(𝒜)`$, where $`𝒜=C_0^{\mathrm{}}(M)\mathrm{\Gamma }`$. We also have a construction of the map $`H^{}(W(𝔤,O_n))H^{}(M_\mathrm{\Gamma })`$ (see e.g. ) where $`M_\mathrm{\Gamma }=M\times _\mathrm{\Gamma }\mathrm{E}\mathrm{\Gamma }`$ is the homotopy quotient. In this section we prove that these constructions are compatible, i.e. that the following diagram is commutative
(7.1)
where $`\mathrm{\Phi }`$ is the canonical map given by Connes .
The proof goes as follows. We construct a map $`\mathrm{\Psi }`$ from the complex computing $`H^{}(M_\mathrm{\Gamma })`$ to the cyclic complex $`(\mathrm{\Omega }^{},d)`$, where $`\mathrm{\Omega }^{}=\mathrm{\Omega }_0^{}(M)\mathrm{\Gamma }`$, which has the following properties. First, it agrees with the map $`\mathrm{\Phi }`$, in the sense that the following diagram is commutative:
(7.2)
where the map $`R`$ is defined by (2.20). Then it is clear from the definitions that the diagram similar to (7.1) is valid with the map $`\mathrm{\Psi }`$ already on the level of cochains, not just cohomology.
The definition of the map $`\mathrm{\Psi }`$ is the following. Recall that the cohomology of $`M_\mathrm{\Gamma }`$ can be computed by the following bicomplex $`𝒞^,`$. $`𝒞^{k,l}`$ denotes the set of totally antisymmetric functions $`\tau `$ on $`\underset{k+1}{\underset{}{\mathrm{\Gamma }\times \mathrm{\Gamma }\mathrm{}\times \mathrm{\Gamma }}}`$ with values in $`l`$-currents on $`\mathrm{Dom}g_0\mathrm{Dom}g_1\mathrm{}\mathrm{Dom}g_k`$, which satisfy the invariance condition
$$\tau (gg_0,gg_1,\mathrm{},gg_k)=\tau (g_0,g_1,\mathrm{},g_k)^{g^1}.$$
(7.3)
The two differentials of this complex are given by the the group cohomology complex differential given on $`𝒞^{k,l}`$ by
$$(d_1\tau )(g_0,g_1,\mathrm{},g_k,g_{k+1})=(1)^l\underset{j=0}{\overset{k+1}{}}(1)^j\tau (g_0,g_1,\mathrm{},\widehat{g}_j,\mathrm{},g_{k+1})$$
(7.4)
and the de Rham differential $`d`$ given by
$$(d_2\tau )(g_0,g_1\mathrm{},g_k)=d\left(\tau (g_0,g_1,\mathrm{},g_k)\right)$$
(7.5)
We now define the map $`\mathrm{\Psi }`$ from the complex $`𝒞^,`$ to the cyclic complex $`(\mathrm{\Omega }^{},d)`$, where $`\mathrm{\Omega }^{}=\mathrm{\Omega }_0^{}(M)\mathrm{\Gamma }`$, by the following formula.
$$\begin{array}{c}\mathrm{\Psi }(\tau )(\omega _0U_{g_0},\omega _1U_{g_1},\mathrm{},\omega _kU_{g_k})=\hfill \\ \hfill \{\begin{array}{cc}(1)^{kl}\tau (1,g_0,g_0g_1,\mathrm{},g_0\mathrm{}g_{k1}),\omega _0\omega _1^{g_0}\mathrm{}\omega _k^{g_0\mathrm{}g_{k1}}\hfill & \text{ if }g_0\mathrm{}g_k=1\hfill \\ 0\hfill & \text{ otherwise}\hfill \end{array}\end{array}$$
(7.6)
Map $`\mathrm{\Psi }`$ satisfies the following identities
$`b\mathrm{\Psi }(\tau )`$ $`=\mathrm{\Psi }(d_1\tau )`$ (7.7)
$`d\mathrm{\Psi }(\tau )`$ $`=\mathrm{\Psi }(d_2\tau )`$ (7.8)
$`B\mathrm{\Psi }(\tau )`$ $`=0`$ (7.9)
and hence it is indeed a map of complexes. It is clear from the definition of the map $`\mathrm{\Psi }`$ that the diagram obtained from the diagram (7.1) by replacing $`\mathrm{\Phi }`$ by $`\mathrm{\Psi }`$ commutes, even on the level of complexes. It remains to prove that the map $`\mathrm{\Psi }`$ induces the same map in cohomology as the map $`\mathrm{\Phi }`$.
To do this we note that Connes’ map $`\mathrm{\Phi }`$ is characterized uniquely by the following properties (cf. ):
1. $`\mathrm{\Phi }`$ takes values in the part of cyclic cohomology supported at the identity of the group.
2. Let $`M`$ be an oriented manifold on which discrete group $`\mathrm{\Gamma }`$ acts freely and properly preserving orientation. Then any class in $`cH^{}(M_\mathrm{\Gamma })`$ can be represented by a $`\mathrm{\Gamma }`$-invariant current $`C`$ on $`M`$. Then $`\mathrm{\Phi }(c)`$ is the same as the class of the cyclic cocycle defined by
$$\begin{array}{c}\varphi (a_0U_{g_0},a_1U_{g_1},\mathrm{},a_kU_{g_k})=\hfill \\ \hfill \{\begin{array}{cc}C,a_0da_1^{g_0}\mathrm{}da_k^{g_0g_1\mathrm{}g_{k1}}\hfill & \text{ if }g_0g_1\mathrm{}g_k=1\hfill \\ 0\hfill & \text{ otherwise}\hfill \end{array}\end{array}$$
(7.10)
3. Let $`X`$ be another oriented $`\mathrm{\Gamma }`$-manifold. Then one has the following commutative diagram
(7.11)
The left vertical arrow here is induced by the natural map $`(M\times X)_\mathrm{\Gamma }M_\mathrm{\Gamma }`$ and the right one is induced by the product with the transverse fundamental class of $`X`$.
It is easy to see that the map $`\mathrm{\Psi }`$ satisfies the same properties, and hence gives the same map in cohomology. |
warning/0002/hep-ph0002186.html | ar5iv | text | # On the Size of the Dark Side of the Solar Neutrino Parameter Space
\[
## Abstract
We present an analysis of the MSW neutrino oscillation solutions of the solar neutrino problem in the framework of two–neutrino mixing in the enlarged parameter space $`(\mathrm{\Delta }m^2,\mathrm{tan}^2\theta )`$ with $`\theta (0,\frac{\pi }{2})`$. Recently, it was pointed out that the allowed region of parameters from a fit to the measured total rates can extend to values $`\theta \frac{\pi }{4}`$ (the so called “dark side”) when higher confidence levels are allowed. The purpose of this letter is to reanalize the problem including all the solar neutrino data available, to discuss the dependence on the statistical criteria in the determination of the CL of the “dark side” and to extract the corresponding limits on the largest mixing allowed by the data. Our results show that when the Super–Kamiokande data on the zenith angle distribution of events and the spectrum information is included, the regions extend more into the second octant.
\]
In a recent work, de Gouvea et al. have pointed out that the study of two–active neutrino oscillations in the framework of the MSW solutions to the solar neutrino deficit as done traditionally on the $`(\mathrm{\Delta }m^2,\mathrm{sin}^2(2\theta )`$ parameter space is incomplete since it covers only the range $`0\theta \frac{\pi }{4}`$ (which they denote as “light side”). By fitting the data on the total rates meassured at Chlorine , Gallium and Super–Kamiokande experiments, they claim that the allowed region of parameters extends to the region with $`\theta \frac{\pi }{4}`$ (denoted as the “dark side”) at some reasonable confidence level (CL). In fact, to our knowledge, the need of extending the mixing parameter space was first discussed in Ref. where the mixing variable $`\mathrm{tan}^2\theta `$ was introduced to chart the full mixing range $`0\theta \frac{\pi }{2}`$.
In this note we revisit the problem of the extension of the “dark side”, that we will denote simply as second octant, of the solar neutrino parameter space after including the effect of all other solar neutrino observables. In particular in this analysis we use all measured total event rates as well as the 825-day Super–Kamiokande data on the zenith angle dependence and the recoil electron energy spectrum of the events. We also discuss the dependence on the statistical criteria used in the construction of the allowed regions. Our results are summarized in Tables I and II where we show the maximum allowed values of $`\mathrm{tan}^2\theta `$ at the 90 and 99% CL when the different observables are included as well as the CL for the second octant of the parameter space. We show that when adding the zenith angle and the spectrum information in the analysis the regions extend more into the second octant. In particular for the LMA solution, this behaviour is mainly driven by the data on the zenith angular dependence since the best fit point for the zenith angle distribution is in the second octant as pointed out in Ref. . For details on the data and statistical analysis employed in this study we refer to our detailed work on Ref. and references therein.
Let’s first recall that in the framework of two massive neutrinos, the weak eigenstates ($`\nu _e`$ and $`\nu _x`$ for the solar neutrino problem) can be writen as a linear combination of the mass eigenstates $`\nu _1`$ and $`\nu _2`$
$`\nu _e`$ $`=`$ $`\mathrm{cos}\theta \nu _1+\mathrm{sin}\theta \nu _2`$ (1)
$`\nu _x`$ $`=`$ $`\mathrm{sin}\theta \nu _1+\mathrm{cos}\theta \nu _2,`$ (2)
where $`\theta `$ is refered to as the mixing angle in vacuum. The mass–squared difference is defined as $`\mathrm{\Delta }m^2=m_2^2m_1^2`$. For the solar neutrino problem $`\nu _x`$ can label either an active neutrino $`x=\mu ,\tau `$ or an sterile neutrino. In what follows we restrict our discussion to oscillations into active neutrinos. For oscillations into sterile neutrinos large mixing solutions are not allowed .
By inspection of the symmetry properties of Eq. (2) one sees that the full parameter space can be exhausted by using the mass-squared difference $`\mathrm{\Delta }m^2`$ always positive and the mixing angle in the interval $`0\theta \frac{\pi }{2}`$. In the case of vacuum oscillations, moreover, the transition probabilities as directly derived from Eq. (2) can be writen in terms of $`\mathrm{sin}^2(2\theta )`$ and therefore they are symmetric under the change $`\theta \theta \frac{\pi }{2}`$ so each allowed value of $`\mathrm{sin}^2(2\theta )`$ corresponds to two allowed values of $`\theta `$.
On the other hand, in the case of the MSW solutions the transition probability takes the form
$$P_{ee}=P_{e1}^{Sun}P_{1e}^{Earth}+P_{e2}^{Sun}P_{2e}^{Earth}$$
(3)
where $`P_{e1}^{Sun}`$ is the probability that a solar neutrino, that is created as $`\nu _e`$, leaves the Sun as a mass eigenstate $`\nu _1`$, and $`P_{1e}^{Earth}`$ is the probability that a neutrino which enters the Earth as $`\nu _1`$ arrives at the detector as $`\nu _e`$ . Similar definitions apply to $`P_{e2}^{Sun}`$ and $`P_{2e}^{Earth}`$. For $`P_{ie}^{Earth}`$ we integrate numerically the evolution equation in matter using the Earth density profile given in the Preliminary Reference Earth Model (PREM) .
The quantity $`P_{e1}^{Sun}`$ is given, after discarding the oscillation terms, as
$$P_{e1}^{Sun}=1P_{e2}^{Sun}=\frac{1}{2}+(\frac{1}{2}P_{LZ})\mathrm{cos}(2\theta _{m,0})$$
(4)
where $`P_{LZ}`$ denotes the standard Landau-Zener probability and $`\theta _{m,0}`$ is the mixing angle in matter at the neutrino production point:
$`\mathrm{cos}(2\theta _{m,0})`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Delta }m^2\mathrm{cos}(2\theta )A_0}{\sqrt{(\mathrm{\Delta }m^2\mathrm{cos}(2\theta )A_0)^2+(\mathrm{\Delta }m^2\mathrm{sin}(2\theta ))^2}}}`$ (5)
$`P_{LZ}`$ $`=`$ $`{\displaystyle \frac{\mathrm{exp}[\gamma \mathrm{sin}^2\theta ]\mathrm{exp}[\gamma ]}{1\mathrm{exp}[\gamma ]}}`$ (6)
$`\gamma `$ $`=`$ $`\pi {\displaystyle \frac{\mathrm{\Delta }m^2}{E}}\left({\displaystyle \frac{d\mathrm{ln}N_e(r)}{dr}}|_{r=r_{res}}\right)^1`$ (7)
with $`A_0=2\sqrt{2}G_FEN_e(r_0)`$ where $`N_e(r_0)`$ is the number density of electrons in the production point, $`E`$ is the neutrino energy,$`G_F`$ is the Fermi constant and $`r_{res}`$ is the resonant point ($`\mathrm{\Delta }m^2\mathrm{cos}2\theta =A_{res}`$). This probability is clearly not invariant under the change $`\theta \frac{\pi }{2}\theta `$ as resonant transitions are only possible for values of $`\theta `$ smaller than $`\frac{\pi }{4}`$ as seen in Eq. (6) what considerably supresses the transitions for $`\theta >\frac{\pi }{4}`$. For this reason most of the earlier papers on the MSW effect that considered the two-neutrino mixing case used ($`\mathrm{\Delta }m^2`$ , $`\mathrm{sin}^22\theta `$) as parameters in the fitting procedure with $`\theta [0,\frac{\pi }{4}]`$. In some cases the parameter space was represented as $`\frac{\mathrm{sin}^22\theta }{\mathrm{cos}2\theta }`$ but still assuming $`\theta `$ to be in the first octant.
Due to the fact that this choice does not exhaust the full parameter space once matter effects are included, other representations have been used to show the enlarged space particularly in the framework of three–neutrino and four–neutrino oscillations . In this way two suggestions to parametrize the mixing angle have been made in the literature: $`\mathrm{sin}^2\theta `$ and $`\mathrm{tan}^2\theta `$ with $`\theta [0,\frac{\pi }{2}]`$. For the sake of comparison with the analysis in Ref. we choose to show our results in terms of $`\mathrm{tan}^2\theta `$.
We now turn to our results. We first determine the allowed range of oscillation parameters using only the total event rates of the Chlorine, Gallium and Super–Kamiokande experiments. We have not included in our analysis the Kamiokande data as it is well in agreement with the results from the Super–Kamiokande experiment and the precision of this last one is much higher. For the Gallium experiments we have used the weighted average of the results from GALLEX and SAGE detectors. Our results are shown in Fig. 1 where we show the allowed regions. We choose to show our allowed regions for the “conventional” 90 and 99 % CL. It is obvious that increasing the allowed CL would lead to larger regions. The CL for the second octant and the maximum values of $`\mathrm{tan}^2\theta `$ for which the LMA and LOW solutions are allowed at those CL can be found in Tables I and II. Furthermore, in Fig. 2 we show the value of $`\chi _{min}^2(\mathrm{tan}^2\theta )\chi _{min}^2`$ in the LMA and LOW regions (where $`\chi _{min}^2(\mathrm{tan}^2\theta `$) is minimized in $`\mathrm{\Delta }m^2`$) as a fuction of $`\mathrm{tan}^2\theta `$. From the figure it is possible to extract the maximum allowed values of $`\mathrm{tan}^2\theta `$ at any other CL. For instance, taking the global analysis in the LMA region from the figure we see that $`\mathrm{tan}^2\theta =3`$ ($`\theta =60^{}`$) is possible only for $`\mathrm{\Delta }\chi ^2>16.6`$ which corresponds to 99.98 for 2 dof ( 3.65$`\sigma `$ ).
We present our results according to two different statistical criteria used in the literature in the defintion of the allowed paramers. The use of each criterion depends on the physics scenario to which the result of the analysis is to be applied.
* Criterion 1 (C1): The regions at certain CL are defined in terms of shifts of the $`\chi ^2`$ function for 2-d.o.f, $`\mathrm{\Delta }\chi ^2`$=4.6 (9.2) at 90 (99) % CL, with respect to the global minimum in the full plane. This criterion is used, for instance, in Refs. . It is applicable when no region of the parameter space SMA, LMA, or LOW is a priory assumed to be the right one.
* Criterion 2 (C2): The regions at certain CL are defined in terms of shifts of the $`\chi ^2`$ funcion for 2-d.o.f with respect to the local minimum in the corresponding region. This criterion holds when assuming that a given solution, SMA, LMA, or LOW, is the valid one. It clearly yields less restrictive limits. This criterion is used, for instance, in Refs. .
In the figures we mark with a star the location of the minima used to define the contours in each case. For instance, in Fig. 1 when using criterion 1, the contours are defined with respect to minimum $`\chi _{min}^2=0.37/1`$ dof which is obtained in the small mixing angle solution. For criterion 2 we define each region with respect to its local minimum whose value and the corresponding goodness of the fit is given in Tables I and II respectively. Notice that the relatively bad goodness of the fit for the LMA and LOW regions from the analysis of the rates only is the reason why for example in Fig. 1 the LMA and LOW regions are quite different using C1 and C2.
As seen in Fig. 1 and Tables I and II when using criterion 1, we do not find any solution in the second octant at 99 % CL from the analysis of the rates only. One must increase the CL to 99.6 (99.6) for the LMA (LOW) region to extend into values $`\mathrm{tan}^2\theta >1`$. This is in agreement, for instance, with the results of Ref. . In Ref. they find some small allowed region in the second octant for the LOW solution at the 99 % CL. We have traced the origin of this small discrepancy to their use of the exponential approximation for the electron number density in the sun. In our calculation we use the solar neutrino fluxes from Ref. and the new numerical parametrization of the sun density as given by Bahcall . We have explicitely verified that if using the exponential approximation in our calculation the allowed LOW region extends into the second octant at 99 % CL.
The situation is changed when the regions are defined according to criterion 2, as seen in Fig. 1 and Tables I and II. In particular both the LMA and LOW regions overlap at 99% CL because the value of $`\chi ^2`$ in between the two regions is below the 99 % CL defined respect to the local LOW miminum (for this reason the LOW region in the figure is only shown at the 90 % CL). Notice also that in our representation we have chosen to cut the parameter space at $`\mathrm{\Delta }m^2>10^8`$ eV<sup>2</sup>. Recently, in Ref. it has been pointed out that matter effects may be relevant for lower values of $`\mathrm{\Delta }m^2`$. One must notice, however that for such lower mass values the simple analytic expresions in Eqs.(4) and (6) start loosing validity .
In Figs. 3 and 4 we show the allowed regions when either the data on the zenith angular dependence or the recoil electron energy spectrum are combined with the results from the total rates. The corresponding values of the absolute minimum of the $`\chi ^2`$ fuction for the combination of rates plus zenith angular dependence data (rates plus recoil electron energy spectrum) are $`\chi _{min}^2=5.9/3`$ dof (22.1/15 dof) which are obtained for the SMA (LMA) solutions and are used in the construction of the allowed regions for criterion 1. In the figures we also show the corresponding excluded regions at 99 %CL by the new observables.
As seen in Fig. 3 and in the tables the inclusion of the data on the zenith angular dependence of the Super–Kamiokande events leads, in general, to an increase in the maximum allowed values of $`\mathrm{tan}^2\theta `$ and a better CL for the second octant for both LMA and LOW regions. This is due to the fact that the best fit point for the angular distribution is obtained in the second octant ($`\mathrm{\Delta }m_{21}^2=3.7\times 10^6`$ eV<sup>2</sup>, $`\mathrm{tan}^2(\theta )=5.9`$ with $`\chi _{min,zen}^2=1.5`$) as shown in Fig. 3 .
This fact has been obviated in past analysis in the two neutrino oscillations and it must be taken into account to do properly the $`\chi ^2`$ analysis when the data on zenith dependence is included. In this respect, one must notice, for instance, that in their preliminary analysis on the zenith angle dependence SuperKamiokande obtains the minimum $`\chi ^2`$ at the boundary of their “cut” mixing parameter space $`\mathrm{sin}^2(2\theta )=1`$. In this way, their $`\chi _{min}^2`$ is higher than the “true” minimum which is missed. Although, at present, the difference in the excluded region defined with respect the “true” or with respect to the “cut” minimum is small, this can be of further importance when more data is acummulated.
We have also explicitely verified that the results on the zenith angle excluded region as well as the best fit point position are very mildly dependent on the exact profile of the Earth. Very similar results can be obtained by using the analytical expressions valid for the two–step Earth density profile .
The effect of the inclusion of the recoil electron energy spectrum is shown in Fig. 4. In this case when using criterion 1 we also find an increase in the maximum allowed values of $`\mathrm{tan}^2\theta `$ and a better CL for the second octant for both LMA and LOW regions as compared to the results for the analysis of the rates only. However this is not the case when using criterion 2 as can be seen in the tables. The allowed regions from the global analysis are displayed in Fig. 5. We see from the figure that at 99 % CL both the LMA and the LOW regions extend into the second octant when using any of the two statistical criteria. In Fig. 2 we show the value of $`\chi _{min}^2(\mathrm{tan}^2\theta )\chi _{min}^2`$ in the LMA and LOW regions (where $`\chi _{min}^2(\mathrm{tan}^2\theta `$) is minimized in $`\mathrm{\Delta }m^2`$) as a fuction of $`\mathrm{tan}^2\theta `$ coming from the fit to the total rates (for which we recall that $`\chi _{min}^2=0.37`$) and to the global data set (for which $`\chi _{min}^2=\chi _{min,LMA}^2=27.0`$). As seen in the figure the inclusion of the new observables leads to a small shift in the position of the local best fit points towards slightly larger values of $`\mathrm{tan}^2\theta `$ . However, as discussed before this not always translates into an increase on the maximum allowed value of $`\mathrm{tan}^2\theta `$ as the minimum also becomes deeper.
To summarize, in this paper we have studied the extension of the second octant of the solar neutrino parameter space after including in the analysis all measured total event rates as well as all the 825-day Super–Kamiokande data on the zenith angle dependence and the recoil electron energy spectrum of the events. We also have discuss the dependence of the results on the statistical criterion used in the construction of the allowed regions. Our results are summarized in Tables I and II where we show the maximum allowed values of the $`\mathrm{tan}\theta `$ at 90 and 99% CL when the different observables are included as well as the CL for the second octant of the parameter space. We have shown that when the zenith angle and the spectrum information is included the regions extend slightly more into the second octant. For the LMA this behaviour is mainly driven by the data on the zenith angular dependence since the best fit point for the zenith angle distribution is in the second octant.
Finally just to comment that, the existence of solutions to the solar neutrino problem for $`\theta >\frac{\pi }{4}`$ is not only of academic interest to the extent that in general they are perfectly allowed by the models of neutrino masses . A particular interesting recent example of a predictive model can be found in Ref. where, despite the predictivity of the model, solutions in both octants are equally possible. Furthermore, in the context of three neutrino mixing, the determination of the sign of $`\mathrm{\Delta }m^2`$ of the LMA solution can be of relevance in the exact determination of the $`CP`$-violating phase .
###### Acknowledgements.
We thank Carlo Giunti, H. Murayama and J. W. F. Valle for discussions. We are specially endebted to J. Bahcall for providing us with the latest parametrization of the sun matter density. This work was supported by DGICYT under grants PB98-0693 and PB97-1261, and by the TMR network grant ERBFMRXCT960090 of the European Union. |
warning/0002/hep-ph0002016.html | ar5iv | text | # 1 Introduction
## 1 Introduction
A suitable method for accelerating the convergence of power series is based on conformal mappings. As is known, a power series converges inside the circle passing through the nearest singularity of the function to be approximated. Some time ago, in Ref. , it was shown that, if the position of the singularities of the expanded function is known, one can reach the fastest convergence rate by expanding in powers of the function that conformally maps the whole holomorphy domain onto a unit disk. In addition, the convergence region extends over the whole holomorphy domain. In a recent paper we applied the technique proposed in to the Borel transform of the Green functions in perturbative QCD. As discussed in , the Borel plane is very suitable for applying the method, since some information about the singularities of the Borel transform is available from the study of certain classes of Feynman digrams and from nonperturbative arguments. By the technique of conformal mapping, this additional information can to a certain extent be incorporated even into the lowest-order terms. In this way, the convergence of the perturbative expansion is improved, allowing one in particular to approximately predict the next-order perturbative terms from the calculated low-order ones .
In Ref. the expansion in powers of the optimal conformal mapping variable<sup>1</sup><sup>1</sup>1The conformal mapping that maps the whole holomorphy domain of the expanded function onto the unit disk will be called optimal. In this case, the singularities are mapped onto the boundary circle, and the requirement of holomorphy implies convergence of the power series at every point of the disk, which is the map of the holomorphy domain. was also used to calculate the Borel-Laplace integral, which is supposed to give, with a certain prescription of treating the infrared renormalons, the Borel summation of the large orders in the Green functions. The numerical results on mathematical models discussed in indicate that the power expansion in the optimal variable makes also the calculation of the Borel integral convergent, in addition with a very high convergence rate. However, only qualitative arguments explaining the results were given, and the problem whether the improved expansion of the integral is convergent, or signs of divergence might appear at large orders, remained open. In the present paper we address this problem and investigate the convergence of the expansion of the Borel integral in perturbative QCD, improved by the use of conformal mapping.
## 2 Optimal expansion of the Laplace-Borel integral
We study the following integral
$$I(a)=\underset{0}{\overset{\mathrm{}}{}}\mathrm{e}^{\frac{u}{a}}B(u)du,$$
(1)
where $`B(u)`$ is assumed to be analytic near $`u=0`$, where it can be expanded as a Taylor series
$$B(u)=\underset{n=0}{\overset{\mathrm{}}{}}b_nu^n$$
(2)
converging inside a circle of non-vanishing radius. The function $`I(a)`$ is of interest for the Borel summation of the Green functions in perturbative QCD; note, however, that the integral on the right hand side of (1) is ill-defined if $`B(u)`$ has singularities along the positive real semiaxis, which is the case of QCD (infrared renormalons). First, the singularities of $`B(u)`$ (renormalons of either kind) make the expansion (2) badly divergent along the integration path. Second, the function $`B(u)`$ itself is, because of infrared renormalons, not uniquely defined along the integration path. We shall discuss both these problems below in this paper.
We shall consider for illustration the Adler function $`D(s)`$ of the massless QCD vacuum polarization, which can be expressed formally as -
$$D(s)=1+\frac{1}{\pi \beta _0}I(a),$$
(3)
with $`a=\beta _0\alpha _s(s)`$, where $`\alpha _s(s)`$ is the running coupling, and $`\beta _0=(332n_f)/12\pi `$ is the first coefficient of the $`\beta `$ function. The expression (3) formally reproduces the renormalization-group-improved expansion of the Adler function
$$D(s)=1+\underset{n=1}{\overset{\mathrm{}}{}}D_n\left(\frac{\alpha _s(s)}{\pi }\right)^n,$$
(4)
by taking the coefficients $`b_n`$ in the expansion (2) of the form
$$b_n=\frac{1}{n!}\frac{D_{n+1}}{(\pi \beta _0)^n}.$$
(5)
We consider also minkowskian quantities, like the hadronic decay rate of the $`\tau `$ lepton, $`R_\tau `$, which can be expressed formally as
$$R_\tau =\mathrm{\hspace{0.17em}3}(1+\delta _{\mathrm{EW}})\left[1+\frac{1}{\pi \beta _0}_0^{\mathrm{}}du\mathrm{exp}\left(\frac{u}{\beta _0\alpha _s(m_\tau ^2)}\right)B(u)F(u)\right].$$
(6)
Here $`\delta _{EW}`$ is an electroweak correction, $`B(u)`$ is the Borel transform of the Adler function and
$$F(u)=\frac{12\mathrm{sin}(\pi u)}{\pi u(u1)(u3)(u4)}.$$
(7)
The extra factor $`\mathrm{sin}(\pi u)`$ in the Laplace-Borel integral is generic for minkowskian quantities. Note that strictly speaking the expressions (3) and (6) are not equivalent to the Borel summation method, which requires an analytic continuation of $`B(u)`$ from the convergence disk to an infinite strip of non-vanishing width, bisected by the real positive semiaxis. This condition is not fulfilled in QCD because of infrared renormalons, which produce cuts of $`B(u)`$ located along the real positive semiaxis.
We shall be concerned with the evaluation of the integral (1) for complex $`a`$ of the general form $`a=|a|\mathrm{e}^{i\psi }`$, where $`\psi =\mathrm{arg}a`$ is the phase of $`a`$. In the case of the Adler function, with the running coupling at one loop in the $`V`$ scheme $`\alpha _s^{(V)}(s)=1/[\beta _0\mathrm{ln}(s/\mathrm{\Lambda }_V^2)]`$, and writing $`s=|s|e^{i(\varphi \pi )}`$, we have
$$\frac{1}{a}=\mathrm{ln}\frac{|s|}{\mathrm{\Lambda }_V^2}i(\pi \varphi ).$$
(8)
Outside the Landau region, i.e. for $`|s|>\mathrm{\Lambda }_V^2`$, we have $`\mathrm{cos}\psi >0`$, so that
$$|\psi |<\frac{\pi }{2},$$
(9)
and $`\psi `$ is related to the momentum plane variable $`s`$ by
$$\psi =\mathrm{Arctg}\left[(\pi \varphi )/\mathrm{ln}\frac{|s|}{\mathrm{\Lambda }_V^2}\right].$$
(10)
The phase $`\psi `$ is positive for $`s`$ in the upper half of the $`s`$\- plane, where $`0<\varphi <\pi `$, negative for $`s`$ in the lower half-plane, where $`\pi <\varphi <2\pi `$, and $`0`$ along the euclidean axis.
For the minkowskian quantity (6) we combine the additional factors $`\mathrm{exp}(\pm i\pi u)`$ due to the sinus with the exponential, which amounts to taking $`a`$ complex with
$$\psi =\pm \mathrm{Arctg}[\pi \overline{a}],$$
(11)
where $`\overline{a}=\beta _0\alpha _s(m_\tau ^2)`$.
As already mentioned, the Borel transform $`B(u)`$ has singularities in the complex plane, correlated to the factorial increase of the perturbative coefficients of the Green functions at large orders , . The precise form of the singularities is not known for the exact theory, but the position and the nature of the first renormalons can be inferred from general principles. In the case of the Adler function the first ultraviolet (UV) renormalon is situated at $`u=1`$ and the first infrared (IR) one at $`u=2`$, and they are branch points of the type $`(1+u)^{\gamma _1}`$ and $`(2u)^{\gamma _2}`$ respectively, with $`\gamma _1`$ computed in Ref. and $`\gamma _2`$ in . We mention also that the summation of the one-renormalon chains in massless QCD in the large $`\beta _0`$ limit gives , :
$$B(u)=\frac{32\mathrm{e}^{Cu}}{3(2u)}\underset{k=2}{\overset{\mathrm{}}{}}\frac{(1)^kk}{\left[k^2(1u)^2\right]^2},$$
(12)
i.e. all the singularities are poles ($`C`$ is a scheme-dependent constant, with $`C=5/3`$ in the $`\overline{\mathrm{MS}}`$ scheme, and $`C=0`$ in the V scheme) , .
The series (2) converges only inside the circle $`|u|<R`$ passing through the nearest singularity ($`R=1`$ for the Adler function). Since the integration range in (1) extends far outside this region, by inserting (2) in (1) and integrating term by term one obtains a divergent expansion. By the technique of conformal mappings one extends the domain of convergence of a series beyond the limit imposed by the first singularity. In we used the optimal conformal mapping
$$w=w(u)=\frac{\sqrt{1+u}\sqrt{1u/2}}{\sqrt{1+u}+\sqrt{1u/2}},$$
(13)
with the inverse $`u(w)=8w/(32w+3w^2)`$. The transformation (13) preserves the origin and maps the complex $`u`$ plane, cut along the real axis for $`u>2`$ and for $`u<1`$, onto the interior of the circle $`|w|<\mathrm{\hspace{0.17em}1}`$, all the singularities of the Borel transform being mapped onto the boundary $`|w|=\mathrm{\hspace{0.17em}1}`$. The expansion in powers of $`w`$ ,
$$B(u)=\underset{n=0}{\overset{\mathrm{}}{}}c_nw^n,$$
(14)
is called optimal because it converges inside the circle $`|w|<1`$, i.e. in the whole domain of holomorphy of $`B(u)`$ (which is the doubly cut complex $`u`$-plane in our case), up to points close to the branch cuts produced by renormalons. As was already pointed out above, this power expansion yields, when compared with other conformal mappings, the fastest large-order convergence rate (see a proof in ). In practice, as discussed in , the expansion (14) is obtained by suitably reorganizing the summation of the original series (2). More precisely, consider the expansion of each $`u^n`$ in powers of $`w`$, truncated at a finite order $`N`$. In particular, in our case this expansion has the general form
$$u_N^n=\underset{j=n}{\overset{N}{}}c_{nj}w^j,$$
(15)
with the coefficients $`c_{nj}`$ obtained by expressing $`u`$ in terms of $`w`$ (in our case $`u(w)`$ is given explicitly after formula (13)). Starting now with the expansion (2) truncated at finite order $`N`$, and replacing each $`u^n`$ by its approximant $`u_N^n`$, one obtains a truncated expansion of the function $`B`$ in powers of $`w`$, which in the limit $`N\mathrm{}`$ gives (14).
By inserting the optimal expansion (14) into the integral (1) we obtain the formal development
$$I(a)=\underset{n=0}{\overset{\mathrm{}}{}}c_nI_n(a),$$
(16)
with
$$I_n(a)=\underset{0}{\overset{\mathrm{}}{}}\mathrm{e}^{\frac{u}{a}}w^ndu.$$
(17)
In the present work we shall adopt (16) as the optimal expansion of the Laplace-Borel integral. We point out that in the physical case this seems to be a natural definition. Indeed, when attempting to make the Borel summation of a perturbation expansion in QCD, one starts with a finite sum of the form
$$I_N(a)=\underset{n=0}{\overset{N}{}}b_n\underset{0}{\overset{\mathrm{}}{}}\mathrm{e}^{\frac{u}{a}}u^ndu.$$
(18)
By replacing here the powers $`u^n`$ with the approximations (15), we replace $`I_N(a)`$ by an expansion of the form
$$\underset{n=0}{\overset{N}{}}c_nI_n(a),$$
(19)
which in the limit $`N\mathrm{}`$ leads to (16).
Actually, as mentioned above, the conditions for the Borel summation are, because of the infrared renormalons, not fulfilled. Therefore, Eq. (16) can be considered as a definition of $`I(a)`$, provided that (i) the integration path in expressions like (1) or (17) is consistently defined, and (ii) the series (16) is convergent. Let us devote a brief discussion to these conditions.
(i) As concerns the integration contour, let us notice that the expansion (16) has not a precise mathematical sense with the $`I_n(a)`$ defined by (17), because the integration path runs along the positive real semiaxis, where the $`w^n`$ have cuts. We shall adopt, as in , the generalized principal value (PV) prescription, defining the $`I_n^{PV}(a)`$ as
$$I_n^{PV}(a)=\frac{1}{2}\underset{𝒞_+}{}\mathrm{e}^{\frac{u}{a}}(w(u))^ndu+\frac{1}{2}\underset{𝒞_{}}{}\mathrm{e}^{\frac{u}{a}}(w(u))^ndu$$
(20)
for $`n=0,1,2,\mathrm{}`$, where $`𝒞_+`$ ($`𝒞_{}`$) are lines parallel to the real positive axis, slightly above (below) it. While the PV prescription does not always give the expected results , in QCD it has the advantage that it reproduces, to a larger extent than other choices, the momentum plane analyticity properties of the Green functions derived from the general principles of field theory. In particular, as discussed in , the Adler function calculated with the PV prescription has no unphysical singularities in the region $`|s|>\mathrm{\Lambda }^2`$. The functions $`I_n^{PV}(a)`$, $`n=1,2\mathrm{}`$ are chosen so as to share some of the known properties with the unknown $`I^{PV}(a)`$. This makes them suited for the definition of $`I^{PV}(a)`$ by means of the expansion
$$I^{PV}(a)=\underset{n=0}{\overset{\mathrm{}}{}}c_nI_n^{PV}(a).$$
(21)
(ii) The convergence of the series (21) for complex $`a`$ is not a priori obvious. Indeed, the expansion (14) converges at points $`|w|<1`$, therefore in the neigbourhood of the integration axis, but not necessarily on the boundary. One might therefore expect that the boundary singularities could manifest in a dramatic way for very large orders $`N`$, making the series (21) divergent, like in the case of the original expansion (2). In we investigated mathematical models with $`B(u)`$ having a few number of isolated branch point singularities, and real values of $`a`$. The numerical results confirm that the expansion (2) in powers of $`u`$ gives results which deviate dramatically from the exact value for large $`N`$, which is typical for a divergent expansion. On the other hand, the improved series (14) in powers of the optimal variable led to results improving continuously with increasing $`N`$, and no signs of divergence appeared even at very high $`N`$. In the next Section we shall discuss the convergence of the optimal expansion of the Laplace-Borel integral, bringing analytic arguments which explain the numerical results obtained in .
## 3 Convergence of the optimal expansion
We investigate the expansion (21) with the functions $`I_n^{PV}(a)`$ defined by means of the PV prescription (20). We consider in our discussion analytic functions $`B`$ of real type, i.e. which satisfy $`B^{}(u)=B(u^{})`$, where $`u^{}`$ is the complex conjugate of $`u`$. Therefore, the coefficients $`b_n`$ in the expansion (2), as well as the coefficients $`c_n`$ in the expansion (14) are real.
The contribution to (20) of integral along the contour $`𝒞_+`$ can be written as
$$I_n^+(a)=\underset{𝒞_+}{}\mathrm{e}^{F_n(u)}du,$$
(22)
where
$$F_n(u)=\frac{u}{a}n\mathrm{ln}w(u).$$
(23)
We evaluate the integral (22) for large $`n`$ by applying the method of steepest descent , . The saddle points are given by the equation
$$\frac{w^{}(u)}{w(u)}=\frac{1}{an},$$
(24)
which has four solutions, having at large $`n`$ the form
$$\frac{1+i}{2^{1/4}}\sqrt{an},\frac{1i}{2^{1/4}}\sqrt{an},\frac{1+i}{2^{1/4}}\sqrt{an},\frac{1i}{2^{1/4}}\sqrt{an}.$$
(25)
Of interest for the evaluation of (22) is the point
$$u_0=2^{1/4}(1+i)\sqrt{an}=|u_0|\mathrm{e}^{i\alpha }$$
(26)
with
$$|u_0|=2^{1/4}\sqrt{|a|n},\alpha =\frac{\pi }{4}+\frac{\psi }{2},$$
(27)
which is situated in the first quadrant of the $`u`$-plane. Indeed, since the phase $`\psi `$ of the parameter $`a`$ satisfies the condition (9), then $`\mathrm{Re}u_0>0`$ and $`\mathrm{Im}u_0>0`$.
Near the saddle point $`F_n(u)`$ can be expanded as
$$F_n(u)=F_n(u_0)+\frac{1}{2}F_n^{\prime \prime }(u_0)(uu_0)^2+\mathrm{}.$$
(28)
By using the expansion of $`w(u)`$ for large $`u`$ in the upper half plane ( $`w\zeta (1i\sqrt{2}/u)`$, where $`\zeta =(\sqrt{2}+i)/(\sqrt{2}i)`$), we obtain after a straightfoward calculation
$$\mathrm{e}^{F_n(u_0)}\zeta ^n\left(1\frac{2^{3/4}i}{(1+i)\sqrt{an}}\right)^n\mathrm{e}^{2^{1/4}(1+i)\sqrt{\frac{n}{a}}}\zeta ^n\mathrm{e}^{2^{3/4}(1+i)\sqrt{\frac{n}{a}}}$$
(29)
and
$$F_n^{\prime \prime }(u_0)\frac{2^{1/4}(1i)}{\sqrt{na^3}}=|F_n^{\prime \prime }(u_0)|\mathrm{e}^{i\beta },$$
(30)
where
$$|F_n^{\prime \prime }(u_0)|=\frac{2^{3/4}}{\sqrt{n|a|^3}},\beta =\frac{\pi }{4}\frac{3\psi }{2}.$$
Therefore (22) becomes
$$I_n^+(a)\zeta ^n\mathrm{e}^{2^{3/4}(1+i)\sqrt{\frac{n}{a}}}\underset{𝒞_+}{}\mathrm{e}^{\frac{|F_n^{\prime \prime }(u_0)|}{2}\mathrm{e}^{i\beta }(uu_0)^2}du.$$
(31)
In order to evaluate the integral we first rotate the contour $`𝒞_+`$ in the trigonometric direction in the upper half-plane, until it becomes a line passing through the origin and the saddle point $`u_0`$. The rotation is possible since $`B(u)`$ (and therefore also $`w`$) has no singularities outside the real axis, and the arc of the circle at infinity gives a vanishing contribution, as can be easily verified. Along the rotated line $`u=\mathrm{e}^{i\alpha }t`$, where $`\alpha `$ is the phase of $`u_0`$ defined in (27) and $`t`$ is real, so the integral in (31) becomes
$$\mathrm{e}^{i\alpha }\underset{0}{\overset{\mathrm{}}{}}\mathrm{e}^{\left[\frac{|F_n^{\prime \prime }(u_0)|}{2}\mathrm{e}^{i(2\alpha +\beta )}(t|u_0|)^2\right]}dt.$$
(32)
Since $`\mathrm{cos}(2\alpha +\beta )>0`$ for $`\psi `$ satisfying the condition (9), the integration axis lies in the two valleys near the saddle point $`u_0`$. Therefore it can be deformed into the path of steepest descent through $`u_0`$, without passing outside the valleys. We take the integral along the path going to infinity
$$uu_0\sqrt{2/|F_n^{\prime \prime }(u_0)|}\mathrm{e}^{i\beta /2}\rho $$
(33)
with real $`\rho `$. The phase of $`(uu_0)^2`$ exactly compensates the phase of $`F_n^{\prime \prime }(u_0)`$, making the exponent of the integrand in (31) real. The integrand can be written as $`\mathrm{e}^{\rho ^2}`$ and the integral done explicitly gives
$$I_n^+(a)\zeta ^n\mathrm{e}^{2^{3/4}(1+i)\sqrt{\frac{n}{a}}}\frac{\mathrm{e}^{i\beta /2}}{\sqrt{|F_n^{\prime \prime }(u_0)|/2}}\frac{\sqrt{\pi }}{2},$$
(34)
i.e. up to a constant independent of $`n`$
$$I_n^+(a)n^{\frac{1}{4}}\zeta ^n\mathrm{e}^{2^{3/4}(1+i)\sqrt{\frac{n}{a}}}.$$
(35)
It is important to note that the path of steepest descent must not cross the real axis, where $`B(u)`$ has singularities. From (33) this implies $`\beta /2=\pi /8+3\psi /4>0`$ which writes as
$$\psi >\frac{\pi }{6}.$$
(36)
The evaluation of the integral along the contour $`𝒞_{}`$ in (20) proceeds in a similar way. The saddle point of interest is
$$u_0^{}=2^{1/4}(1i)\sqrt{an}=\mathrm{\hspace{0.17em}2}^{1/4}\sqrt{|a|n}\mathrm{e}^{i(\frac{\pi }{4}\frac{\psi }{2})},$$
(37)
which satisfies $`\mathrm{Re}u_0>0`$ and $`\mathrm{Im}u_0<0`$ for $`\psi `$ in the range given in (9). Instead of (31) we have
$$I_n^{}(a)(\zeta ^{})^n\mathrm{e}^{2^{3/4}(1i)\sqrt{\frac{n}{a}}}\underset{𝒞_{}}{}\mathrm{e}^{\frac{|F_n^{\prime \prime }(u_0^{})|}{2}\mathrm{e}^{i\beta ^{}}(uu_0^{})^2}du,$$
(38)
where $`\beta ^{}=\pi /43\psi /2`$. We rotate the contour $`𝒞_{}`$ in the lower half-plane up to a line passing through the point $`u_0^{}`$, and then deform it into the steepest descent path. One can easily verify that this path does not cross the real axis for $`\pi /8+3\psi /4<0`$ which writes as
$$\psi <\frac{\pi }{6}.$$
(39)
Collecting the terms we obtain the coefficients $`I_n(a)`$ in the PV prescription (20) as
$$I_n(a)n^{\frac{1}{4}}\zeta ^n\mathrm{e}^{2^{3/4}(1+i)\sqrt{\frac{n}{a}}}+n^{\frac{1}{4}}(\zeta ^{})^n\mathrm{e}^{2^{3/4}(1i)\sqrt{\frac{n}{a}}}.$$
(40)
In order to examine the convergence of the expansion (21), we consider the ratio
$$\left|\frac{c_nI_n(a)}{c_{n1}I_{n1}(a)}\right|,$$
(41)
for large $`n`$. If the coefficients $`c_n`$ do not grow too rapidly, i.e.
$$|c_n|<C\mathrm{e}^{ϵn^{1/2}},$$
(42)
for all $`ϵ>0`$, then the expansion (21) converges for $`a`$ complex in the domain
$$\mathrm{Re}[(1\pm i)a^{1/2}]>0.$$
(43)
As we already discussed these conditions are equivalent to (9). If the coefficients $`c_n`$ behave at large $`n`$ like
$$|c_n|\mathrm{e}^{cn^{1/2}}$$
(44)
for some positive $`c`$, then the expansion (21) converges in the domain
$$\mathrm{Re}[(1\pm i)a^{1/2}+c]>0,$$
(45)
while for coefficients $`c_n`$ which grow faster than $`\mathrm{exp}(cn^{1/2})`$ the new series (21) is also divergent. We mention that such a behaviour is not excluded in general for series of the form (14) with a radius of convergence equal to 1 .
We recall however that the expression (40) is valid only for $`\psi `$ which satisfy the conditions (36) and (39), i.e.
$$|\psi |<\frac{\pi }{6},$$
(46)
which define a sector in the $`a`$\- complex plane (we recall that $`\psi `$ is the phase of $`a`$). This inequality is a condition of applicability of the steepest descent method used by us. We found therefore that the expansion (21), improved by the optimal conformal mapping of the Borel plane, is convergent if the Taylor coefficients $`c_n`$ of the expansion (14) satisfy the condition (42), at least inside the sector (46) of the complex plane of $`a`$, or, if they behave like (44), in the smallest of the domains (45) and (46). For the Adler function in massless QCD, using (10) we write the condition (46) in the form
$$|\pi \varphi |<\frac{1}{\sqrt{3}}\mathrm{ln}\frac{|s|}{\mathrm{\Lambda }_V^2},$$
(47)
where $`\varphi `$ is the phase of $`s`$ and we have $`|s|>\mathrm{\Lambda }_V^2`$. For the minkowskian quantities, from (11) we obtain
$$\pi \overline{a}<\frac{1}{\sqrt{3}},$$
(48)
which means in particular that for the $`\tau `$\- hadronic decay rate the expansion defined as in (21) is convergent for $`\alpha _s(m_\tau ^2)<4/(9\sqrt{3})0.257`$.
The behaviour of the coefficients $`c_n`$ depends on the singularities of $`B(w)`$. By the conformal mapping (13) all the renormalons are situated on the circle $`|w|=1`$, appearing in conjugate pairs since $`B(u)`$ is of real type. Assuming that all the singularities are poles or branch points, $`c_n`$ has the generic form
$$c_n\frac{1}{n!}\mathrm{Re}\underset{j}{}r_jp_j(p_j+1)(p_j+2)\mathrm{}(p_j+n)\mathrm{e}^{i\beta _j\gamma _j(p_j+n)},$$
(49)
where $`\mathrm{exp}(\pm i\beta _j)`$ denote the position of the renormalon in the $`w`$-plane, $`r_j`$ the residue, and $`p_j`$ the exponent of the singularity. In Ref. we investigated simple models with a finite number of singularities, and real values of the parameter $`a`$, for which the conditions of convergence are satisfied. In the physical case, one knows only that for the first UV renormalon $`\alpha _1=\pi `$ and $`p_1=2\gamma _1`$, and for the first IR renormalon $`\alpha _2=0`$ and $`p_2=2\gamma _2`$. In the large $`\beta _0`$ case, as seen from (12), all the singularities are poles, $`p_j`$ in (49) is independent of $`j`$, and $`r_j`$ are known. In this case the condition of convergence (42) is satisfied. Therefore, the optimal expansion on the Laplace-Borel integral, in the PV prescription, for the summation of one renormalon chains in the large $`\beta _0`$ limit, is convergent, at least in the sector of the complex $`a`$ plane defined by the condition (46).
In conclusion, we investigated the expansion of the Laplace- Borel integral in perturbative QCD, improved by the analytic continuation of the Borel transform outside the perturbative convergence disk (and, simultaneously, by reaching the fastest convergence rate) by means of the optimal conformal mapping . The convergence properties of the new expansion depend on the strength of the singularities of the Borel transform, reflected in the behaviour of the Taylor coefficients of the expansion (14). If the Taylor coefficients satisfy the condition (42), the new expansion of the Laplace- Borel integral converges in the sector of the complex plane of the coupling $`a`$ defined by (46). The conditions are satisfied in the case of the resummation of one-loop renormalons in the large- $`\beta _0`$ limit. We mention that in the region where the series converges the function $`I(a)`$ must be analytic. For the Adler function in the complex momentum plane this corresponds to the region described by Eq. (47), where $`D(s)`$ is analytic.
Acknowledgements: One of us (I.C.) thanks Prof. S. Randjbar-Daemi for his kind hospitality at the High Energy Section of the Abdus Salam International Centre of Theoretical Physics, Trieste. The other author (J.F.) is indebted to Prof. A. De Rújula for hospitality at the CERN Theory Division. |
warning/0002/nlin0002030.html | ar5iv | text | # On the stability of long-range sound propagation through a structured ocean
## I Introduction
There is a great deal of experimental and theoretical interest in long-range, low-frequency acoustic pulse propagation through the deep ocean’s sound channel. It has been investigated as a problem of wave propagation in random media (WPRM) , and as a basis for tomography . Recent results from the Acoustic Engineering Test (AET) as part of the Acoustic Thermometry of Ocean Climate (ATOC) project can be found in Colosi et al. and Worcester et al. . One of the main challenges in analyzing and understanding long range acoustic propagation is in dealing with difficulties arising from the ocean environment’s tendency to generate multiple, weak, small-angle (forward) scattering . At sufficiently long ranges of propagation, the multiple scattering should effectively randomize an acoustic pulse so that it is very difficult to deduce much information. However, several long range experiments have found a great deal of stability in the earlier portions of the received wave fronts in spite of the fluctuations inherent in the ocean environment . In addition, it has been found that wave field intensity fluctuations at long range are consistent with a lognormal density which would be reminiscent of earlier work in optics on WPRM , except that this earlier work was for the short range (weak focusing) regime.
In the past 10-15 years, simplified models inspired by the ocean environment have been shown to possess chaotic ray limits . Essentially simultaneously, there has been enormous progress in the understanding of chaotic systems . Some of the most familiar emerging concepts are simpler for bounded systems and are not easily applicable to open, scattering systems as we have here. However, there is an important tool which does straightforwardly generalize for our purposes, the stability analysis of the rays. Stability matrices can be constructed as a function of range for each ray. Their properties, such as the stability exponents, reveal the basic character of the rays, and are at the foundation of the findings reported in this paper.
There are several intriguing questions that arise from comparing the theoretical results to date regarding chaotic acoustic ray dynamics in the ocean and the high amount of stability observed in the data. The most general question concerns how an acoustic pulse – which at multi-megameter ranges extends to nearly 10 seconds in time and 2 km in depth – loses it’s coherence from multiple forward scattering through interaction with internal waves and mesoscale energetics. Because refraction is adequate to explain the scattering physics , the ray limit should suffice for understanding long-range propagation. Some manifestations of the underlying chaotic dynamics should be observed.
It has been suggested that there exists a “predictability horizon” at the range of propagation defined by the scale over which chaotic dynamics develops . Beyond this range, the wave fields should appear as random superpositions of many plane waves which would imply that acoustic field intensity fluctuations are Rayleigh distributed . Several problems crop up beyond the predictability horizon. It becomes increasingly difficult to get numerically calculated rays to converge to true rays of the system. Worse, semiclassical approximations (i.e. wave front reconstructions) from the rays might fail for fundamental reasons related to the breakdown of stationary phase approximations, but one should recognize that more optimistic viewpoints exist on this issue . Whether or not this is true, it is currently not known to what extent tomographic inversions fail for any system beyond its predictability horizon where eigenrays are proliferating exponentially fast with increasing range. In order to begin addressing these and related issues, we focus on the ‘forward propagation problem’ by performing a statistical analysis that should be much less sensitive to the difficulties engendered by the predictability horizon.
In fact, justifications for statistical laws derived by invoking stochastic or ergodic postulates are often ultimately founded on the presence of fully developed chaos; see for example Ref. . Systems that once were approached by stochastic methods have more recently begun to be regarded from the perspective of dynamical systems. The two approaches mostly give consistent results, but there are important distinctions. Stochastic ray modeling is the traditional approach to the geometric limit of the problem of WPRM . This nondeterministic treatment leads one to pessimistic conclusions regarding the overall stability expected in an ocean acoustic pulse at sufficiently long range . By carefully defining the Lyapunov exponent, it turns out to be roughly half the value reported in Ref. . The scales relevant to the ocean are such that this factor two increase in an important length scale might prove to be significant. Also, for this problem, the validity of the stochastic or ergodic assumptions deserves to be critically examined. It is not obvious that a dynamical systems perspective would lead to similar pessimistic conclusions as does the stochastic ray theory. For example, the predictability horizon concept that has grown out of the chaotic dynamics point of view does not necessarily lead to a sudden transition — regular behavior at short ranges, completely stochastic just beyond — and remnants of stability that violate assumptions of stochasticity could persist well into the horizon’s initial onset. We anticipate several features of deterministic dynamics playing an interconnected role in this regard, but we focus on the importance of only one, intermittent-like dynamics. Intermittency is a common feature for nonintegrable dynamical systems . For the ray acoustics problem, intermittent-like behavior is evident through the appearance of rays which persist in remaining relatively insensitive to their initial conditions (also environment) for remarkably long ranges, as measured on the inverse scale of the mean Lyapunov exponent. It might then be expected that the existence of intermittent-like dynamics might allow linear based tomographic inversions based on acoustic ray models to be suitable to greater ranges than previously anticipated.
One objective of this article is to illustrate the existence of intermittent-like dynamics in the generic long-range ocean acoustics problem. This is a direct consequence of the wide variability in the eigenvalues of the stability matrix which is defined in Sect. IIIA. We demonstrate that the magnitude of its largest eigenvalue follows a lognormal distribution, and that the stability exponent follows a Gaussian distribution. Importantly, there is preliminary evidence suggesting that these distributions are robust, i.e. that they would be found in much more realistic, sophisticated ocean models . To be more explicit, if one knows the probability density of the stability exponents, then one can determine the expected measure of intermittent-like rays that will persist out to the reception range. It follows that these rays will not require extremely precise numerical interpolation schemes for quantities such as the gradient of the potential.
The model upon which we rely in this paper is admittedly extremely simplified. However, it is not the model that is of concern, it is whether or not general features of simplified WPRM models carry over to the ocean itself. If we are careful enough, the simplifications that we accept remove non-essential complications for uncovering the general physical features of interest, and no more. A follow-up study to this one is underway which uses a more realistic ocean sound speed model. It is important for confirming the applicability of our results to long-range ocean acoustics experiments.
The organization is as follows: in Sect. II, we introduce and motivate a simple model leading to a one-degree-of-freedom, non-autonomous Hamiltonian dynamical system for the rays. This is followed by a discussion of the analysis methods which are most critical for our study. They are based on the stability matrix and its well known properties. Sect. IV examines the fluctuation behavior of the stability exponents giving their densities as a function of range. The proportion of intermittent-like rays is deduced and compared with the numerical results of the model. We finish with a discussion and conclusions.
## II From wave equation to ray model
We briefly outline the assumptions and approximations leading to the highly idealized ray model used in this paper. The primary physics we are concerned about involve refraction of acoustic energy due to volume inhomogeneities in the ocean sound speed. We assume that interactions of the acoustic energy with both the surface and sub-bottom are negligible. For multi-megameter ranges of propagation in mid-latitude, deep ocean environments, a significant amount of acoustic energy is received that satisfies this assumption . As alluded to above, the necessary assumptions leading to the primary results are that: i) the linear, one-way Helmholtz wave equation is valid (the important point here is that backscattering is negligible), and ii) the spatial scales of the sound speed field are long compared to the acoustic wavelength so that ray theory is justified. A detailed derivation is readily available . We point out a priori that the coordinate system is three-dimensional Cartesian $`𝐱=(r,y,z)`$, with $`r`$ the range from the source, $`y`$ the transverse or cross-range coordinate, and $`z`$ the depth from the surface. Thus, Earth curvature effects are neglected.
The fundamental starting point is the linear acoustic wave equation :
$$\frac{1}{c^2}\frac{^2\psi (𝐱;t)}{t^2}=^2\psi (𝐱;t),$$
(1)
where $`\psi (𝐱;t)`$ is the complex scalar wave function whose real part denotes the acoustic pressure. The sound speed field $`c`$ can be taken as a function of space $`𝐱`$ only whereby it has been assumed that the time scales for the propagation of the acoustic wave function are small compared to the time scale associated with the evolution of the sound speed field. For non-dispersive sources, the acoustic group and phase speeds are equivalent, and one can linearly transform Eq. (1) from time to frequency, arriving at a Helmholtz equation with the magnitude of the wave vector defined as $`k=2\pi f/c`$, where $`f`$ is the continuous wave (CW) source frequency. Attenuation effects can of course be incorporated by modifying $`k`$ to be a complex quantity, but since we are interested in: 1) acoustic energy that interacts negligibly with the ocean bottom, and 2) typical sources operating at frequencies with minimal volume attenuation (with center frequencies of about 100 Hz), ignoring attenuation effects seems reasonable. Also, one can similarly derive a reduced wave equation which includes variations in density, but we ignore this effect because it is known to be important predominantly with acoustic energy that interacts with the ocean sub-bottom, which is not considered herein.
The next assumption (which is quite a reasonable one) is that the strength of the sound speed fluctuations, whatever the physical process that produces them, are small. This allows one to neglect backscattered acoustic energy, and admits the one-way Helmholtz wave equation, whereby one assumes a primary direction of propagation along the range. The so-called “ standard parabolic approximation” is invoked next. This allows one to derive a linear partial differential wave equation of parabolic type for the complex envelope of $`\psi `$. The principle assumption is that this envelope wave function evolves slowly on the scale of the acoustic wavelength. There many flavors of parabolic approximations that have varying degrees of phase errors in the complex wave function $`\psi `$ as compared to the one-way Helmholtz equation , but we choose to use the standard parabolic approximation, which takes the form
$$\frac{i}{k_0}\frac{\varphi (y,z;r)}{r}=\frac{1}{k_0^2}_{}^2\varphi (y,z;r)+V(y,z;r)\varphi (y,z;r),$$
(2)
where the transverse Laplacian is represented by $`_{}^2=_y^2+_z^2`$, and the variable $`r`$ is the range (propagation variable), but plays an exact analogous role to time in the Schrödinger equation of quantum mechanics. The parameter $`k_0=2\pi f/c_0`$ represents the reference wave number, and depends on the choice of a reference sound speed $`c_0`$, which we take to be 1.5 km/s. The potential, $`V(y,z;r)`$, is related to the sound speed fluctuations as
$$V(y,z;r)=\frac{1}{2}\left[\left(\frac{c_0}{c(y,z;r)}\right)^21\right]\frac{\delta c(y,z;r)}{c_0},$$
(3)
where the sound speed variations away from an average profile has been expressed as $`c(y,z;r)=c_0+\delta c(y,z;r)`$. The sound speed fluctuations refract the rays and lead to chaos in a deterministic, mathematically defined sense. Under the parabolic approximation, the basic problem maps precisely onto problems of quantum chaos . The fields of long range acoustic propagation in the ocean and quantum chaos thus have the opportunity of cross-fertilization.
Because the instability does not critically depend on having multiple degrees of freedom, we make a significant, practical simplification in the model of ignoring the depth degree of freedom ($`z`$); see Ref. for a more detailed discussion of the model presented here. The system could be thought of as lying in the plane of the sound channel axis, but this is really just the generic problem of WPRM (see, for example, ). The gain in simplicity more than compensates for the loss of realism at this point as long as the main physical phenomena carry over to more realistic models. As was mentioned in the Introduction, preliminary evidence for our main results have been found in recent calculations incorporating a much more realistic model .
The magnitude of the wave vector $`k`$ is large enough that for the purposes of this study, we can focus on the ray limit. The rays can be generated by a system of Hamilton’s equations
$`{\displaystyle \frac{dy}{dr}}`$ $`=`$ $`{\displaystyle \frac{H(y,p;r)}{p}},`$ (4)
$`{\displaystyle \frac{dp}{dr}}`$ $`=`$ $`{\displaystyle \frac{H(y,p;r)}{y}},`$ (5)
where $`y`$ and $`p`$ are the phase space variables cross-range (position) and horizontal slowness (momentum) respectively. The independent variable $`r`$ denotes range. Correspondence with Eq. (2) necessitates that the Hamiltonian is explicitly
$$H=\frac{p^2}{2}+V(y;r).$$
(6)
The physical meaning of momentum is $`p=\mathrm{tan}\theta `$, where $`\theta `$ represents the angle a ray subtends in cross-range about the range axis.
The state of the ocean is constantly changing, and its exact state is unknown. A statistical ansatz is thus fruitful for making assertions concerning its “typical” state. Assuming isotropy in the sound speed fluctuations in range and cross-range, the potential is taken to be a realization of a zero-mean, stationary, random function. Thus a single correlation length scale $`L`$ exists. The standard deviation is denoted by $`ϵ=c_0^1\delta c^2^{1/2}`$, where $`\delta c^2^{1/2}`$ is the root-mean-square fluctuation of the sound speed about $`c_0`$. Typical values for underwater acoustics are $`ϵ=O(10^3)`$, and $`L=O(100)\text{km}`$, but both $`ϵ`$ and $`L`$ vary plus or minus an order of magnitude depending on what ocean structure is considered and the geographic location. For purposes of studying a fully defined, deterministic dynamical system, we complete the description of $`V`$ by defining its correlation function to be Gaussian,
$$V(y;r)V(y+\delta y;r+\delta r)=ϵ^2\mathrm{exp}[(\delta y^2+\delta r^2)/L^2].$$
(7)
We exploit this single scale throughout the rest of this article by transforming space variables as $`rr/L`$ and $`yy/L`$, so the physical dimensions will always be in units of $`L`$. One should envision the potential as being deterministic, even though it is selected from an ensemble of realizations. This implies that the potential is to be considered a highly complicated (albeit smooth and fixed) function of both $`y`$ and $`r`$. To provide some idea of the character of this potential, contours of sound speed fluctuations based on a typical region of $`V(y;r)`$ is shown contoured in Fig. 1. The boundary conditions are taken as open in $`y`$, but numerically $`y`$ is treated as periodic, with the ray coordinate unfolded a posteriori to simulate the open boundary condition. A variety of initial conditions are possible with the restriction that the initial momentum is always kept small enough that the parabolic approximation is valid all along the rays. The rays deriving from two such initial conditions are plotted in Fig. 2 which shows their phase space portraits (position, momentum). In the absence of a varying potential, the solutions to the equations of motion are $`p(t)=p_0`$, $`q(t)=p_0t+q_0`$. In this figure, rays would trace out vertical lines except in the case, $`p_0=0`$, in which rays would show up as points, $`(q(t),p(t))=(q_0,0)`$. With the potential included, the rays trace out a random-walk-like motion with some drift as they move further away from the $`p=0`$ line.
## III Analysis Methods
The standard analysis of ray stability in the theory of dynamical systems begins with the stability matrix. From here, it is possible to calculate whether a ray is stable or unstable, what its Lyapunov exponent is, and for the unstable ray, determine the orientations of the associated stable and unstable manifolds that characterize the exponential sensitivity to initial conditions. All of our results and conclusions are based on the behavior of the stability matrices of the rays in the model introduced in the previous section. The stability matrix is a strictly local analysis in range about some particular reference trajectory. It may be stable at one range, yet for a greater range be unstable. There is no restriction that various portions of its full history cannot have completely different stability properties. In fact, one expects the portions to be almost entirely uncorrelated .
Often research done in chaotic dynamics uses either time (range) independent or periodic Hamiltonians, and the stability matrix is investigated about periodic orbits. The Hamiltonian of Eq. (6) is aperiodic, and as many others have done before, we slightly generalize those treatments by considering arbitrary, aperiodic rays.
### A Stability Matrix
The stability matrix for a ray describes the behavior of other rays that remain within its infinitesimal neighborhood, $`\{\delta y,\delta p\}`$, for all ranges. It is derived by linearizing the dynamics locally; see Ref. for more details. At the range $`r`$, one has
$$\left(\begin{array}{c}\delta p_r\\ \delta y_r\end{array}\right)=M\left(\begin{array}{c}\delta p_0\\ \delta y_0\end{array}\right),$$
(8)
with the stability matrix being given by the partial derivatives
$$M=\left(\begin{array}{cc}m_{11}& m_{12}\\ m_{21}& m_{22}\end{array}\right)=\left(\begin{array}{cc}\frac{p_r}{p_0}|_{y_0}& \frac{p_r}{y_0}|_{p_0}\\ \frac{y_r}{p_0}|_{y_0}& \frac{y_r}{y_0}|_{p_0}\end{array}\right).$$
(9)
The multi-dimensional generalizations are immediate. The $`m_{21}`$ matrix element is well known for its appearance in the prefactor of the standard time (range) Green’s function of the parabolic equation; it therefore gives directly information on wave amplitudes.
The evolution of $`M`$ is governed by
$$\frac{d}{dr}M=KM,$$
(10)
with the initial condition $`M(r=0)`$ being the identity matrix, and
$$K=\left(\begin{array}{cc}\frac{^2H}{yp}& \frac{^2H}{y^2}\\ \frac{^2H}{p^2}& \frac{^2H}{yp}\end{array}\right)\left(\begin{array}{cc}0& \frac{^2V}{y^2}\\ 1& 0\end{array}\right).$$
(11)
The latter form is the simplification relevant for Hamiltonians of the so-called mechanical type as in Eq. (6). Since Eq. (10) represents linear, coupled, first-order differential equations, the elements of $`M`$ can be numerically calculated as a function of range simultaneously with the calculation of its reference ray using identical numerical techniques, e.g. variable step, fourth-order Runge-Kutta.
### B Stability and Lyapunov exponents
The stability matrix has several important properties. It can be viewed as generating a linear, canonical transformation, and therefore its determinant is equal to unity. It is diagonalized by a linear, similarity transformation
$$\mathrm{\Lambda }=LML^1\left(\begin{array}{cc}\lambda & 0\\ 0& \lambda ^1\end{array}\right),$$
(12)
where the last form applies specifically to the case of a single degree of freedom. Here, the second eigenvalue must be the inverse of the first in order for $`\text{det}[M]=1`$. The diagonalizing similarity transformation leaves the sum of the diagonal elements (trace), $`\text{Tr}(M)`$, invariant. It is then clear that $`\text{Tr}(M)`$ is real, and three distinct cases may arise. The first is $`|\text{Tr}(M)|<2`$ which is linked to stable motion, and it is then customary to denote $`\lambda =\mathrm{exp}(i\theta r)`$. The second case is $`|\text{Tr}(M)|=2`$, and it is often called marginally stable because it is the boundary case between stable and unstable motion. The third case represents unstable motion, and is characterized by $`|\text{Tr}(M)|>2`$. In Fig. 2, a typically stable ray is represented by the solid line, and a typically highly unstable ray is represented by the dashed line. Their distinctions are not immediately obvious.
The evolution of neighboring rays about a ray that has $`|\text{Tr}(M)|<2`$ satisfied from the source to the reception range will undergo only rotations in phase space, and subsets of phase space of finite measure where this behavior dominates the dynamics is precisely where intermittent-like rays reside. Fig. 3 illustrates this characteristic behavior by showing a group of stable rays winding about each other as they propagate. The dashed ray in the group is the stable ray of Fig. 2. They perform their “random walks”, yet remain winding about each other. For the purposes of this paper, we make a slightly generalized definition of intermittent-like rays as being all those for which $`|\text{Tr}(M)|`$ remains sufficiently small over the range of propagation, i.e. not far from two.
For unstable motion, it is customary to denote $`\lambda =\pm \mathrm{exp}(\nu r)`$ where $`\nu `$ is positive and real. The neighboring rays move hyperbolically relative to each other. We add a collection of unstable (chaotic) rays onto Fig. 3 to illustrate the distinction between neighboring groups of stable and unstable rays. The dashed ray is the highly unstable ray from Fig. 2. The rays were selected to span the same size initial neighborhood as the stable group, yet they fan out and become completely independent. For the unstable case one can introduce a definition for the Lyapunov exponent as
$$\nu _L\underset{r\mathrm{}}{lim}\frac{\mathrm{ln}\left(|\text{Tr}(M)|\right)}{r}.$$
(13)
Note that there is no ensemble averaging implied in the definition of $`\nu _L`$. None of the theory presented thus far prevents it from taking on a distinct value for each ray for each realization of the random potential. In Sect. IV, the value of $`\nu _L`$ will be shown to be independent of the particular ray, and the particular realization of the potential as well. It thus defines a unique length scale, $`\nu _L^1`$, which is used from Fig. 3 onward wherever the range variable is involved.
For unstable motion, $`|\lambda |`$ tends to be very large leaving $`\lambda ^1`$ negligible. With little inaccuracy, $`\text{Tr}(M)=\lambda +\lambda ^1\lambda `$ , even for finite ranges. One then deduces a stability exponent, $`\nu `$ from $`\text{Tr}(M)`$ as
$$\nu =\frac{\mathrm{ln}|\text{Tr}(M)|}{r}.$$
(14)
Thus $`\nu `$ depends on the particular ray and varies with range whereas the Lyapunov exponent has no range dependence by definition.
For any fixed range, an ensemble of $`\nu `$ can be created by considering various initial conditions (by exploiting the isotropic and stationary properties of $`V(y;r)`$), and different realizations of $`V(y;r)`$. The resulting statistical densities of $`|\text{Tr}(M)|`$, $`\rho _{|\text{Tr}(M)|}(x)`$, and similarly $`\nu `$, $`\rho _\nu (x)`$, are the main objects of concern; the two densities are directly tied to each other. The cumulative probability distribution is given as
$$F_\nu (x)=_{\mathrm{}}^xdx^{}\rho _\nu (x^{})$$
(15)
which provides a useful tool for numerically studying the behavior of $`\rho _\nu (x)`$. It also has the utility of directly giving the proportion of nearly stable rays up to some argument set to a maximum instability criterium, $`\nu =x`$.
We denote the mean and variance respectively as
$`\nu _0`$ $`=`$ $`\nu ={\displaystyle dxx\rho _\nu (x)},`$ (16)
$`\sigma _\nu ^2`$ $`=`$ $`\left(\nu \nu _0\right)^2={\displaystyle dx\left(x\nu _0\right)^2\rho _\nu (x)},`$ (17)
where the brackets $``$ denote ensemble averaging. For any real $`\gamma `$, ensemble averages of powers of $`|\text{Tr}(m)|`$ are expressed as
$$\left|\text{Tr}(M)\right|^\gamma =\mathrm{exp}(\gamma \nu r)=dx\mathrm{exp}(\gamma xr)\rho _\nu (x).$$
(18)
Note that the case of $`\gamma =1/2`$ relates to wave amplitude statistics resulting from a semiclassical reconstruction of the wave field, and will be discussed in a future work.
To continue the theoretical development, it is useful to introduce a slightly modified stability exponent, $`\overline{\nu }`$:
$$\overline{\nu }\frac{\mathrm{ln}|\text{Tr}(M)|^2}{2r}.$$
(19)
Clearly $`\overline{\nu }`$ is necessarily greater than $`\nu _0`$ because of the important distinction of ensemble averaging before taking the natural logarithm as opposed to the inverse order and the root mean square fluctuation contributions. It is shown in the next section that the Lyapunov exponent becomes $`\nu _L=lim_r\mathrm{}\nu _0`$, and not $`lim_r\mathrm{}\overline{\nu }`$ which surprisingly remains greater than $`\nu _L`$. Near a parameter regime motivated by the ocean, we find numerically that analytical estimates of $`\nu _L`$ as being equal to $`\overline{\nu }`$ are roughly double their actual values.
### C Stochastic analysis results
An analytic estimate of $`\overline{\nu }`$ can be derived from previous analytic results based on a stochastic analysis which involves a strong Markovian assumption . It was verified in Ref. that, in the context of the present acoustic ray model, the stochastic analysis predictions for $`\overline{\nu }`$ (actually $`\nu ^{}`$, see text ahead) matched to a high degree of precision with numerical tests. In fact, no statistically significant deviations were observed. Thus, although the stochastic system is not strictly mathematically equivalent to the deterministic dynamics, we accept the applicability of those specific results at sufficiently long ranges (defining this range scale is admittedly not as trivial to determine for the general ocean acoustics scenario as it is for the idealized problem). We begin with
$$[\text{Tr}(M)]^2=m_{11}^2+m_{22}^2+2m_{11}m_{22}.$$
(20)
By appealing to stochastic integration techniques , it has been shown that in the small-$`ϵ`$, large-$`r`$ limit that
$$m_{22}^2=\frac{1}{3}\mathrm{exp}(2\nu ^{}r),$$
(21)
where
$`\nu ^{}`$ $``$ $`({\displaystyle \frac{1}{2}}{\displaystyle _0^{\mathrm{}}}d\xi {\displaystyle \frac{^2V(y;r\xi )}{y^2}}|_{\genfrac{}{}{0pt}{}{y=y_0}{p=p_0}}`$ (23)
$`{\displaystyle \frac{^2V(y;r)}{y^2}}|_{\genfrac{}{}{0pt}{}{y=y_0}{p=p_0}})^{1/3}`$
$`=`$ $`(3\sqrt{\pi })^{1/3}ϵ^{2/3}`$ (24)
(in dimensional units $`\nu ^{}=(3\sqrt{\pi })^{1/3}ϵ^{2/3}/L`$). The last result of Eq. (23) is for the specific example of a Gaussian single scale potential of Eq. (7). The first result of Eq. (23) is more general, but requires numerical confirmation for models with greater realism, and will also depend on the ray’s initial conditions for models with a nonuniform background sound speed field.
By symmetry considerations of the stochastic equations, $`m_{11}^2=m_{22}^2`$. It is also deduced that $`m_{11}m_{22}`$ can, at most, grow on the same scale. Defining a correlation coefficient,
$$\mu =\frac{m_{11}m_{22}}{m_{22}^2},$$
(25)
where $`|\mu |1`$, it follows that
$$[\text{Tr}(M)]^2=\frac{2}{3}(1+\mu )\mathrm{exp}(2\nu ^{}r).$$
(26)
Then, by using the definition of Eq. (19), one obtains
$$\overline{\nu }=\nu ^{}+\frac{1}{2r}\mathrm{ln}\left[\frac{2}{3}(1+\mu )\right].$$
(27)
Note that the second term disappears if $`\mu `$ equals $`1/2`$; i.e. $`m_{11}m_{22}=m_{22}^2/2`$. We give its value numerically in the next section. In that case, $`\overline{\nu }=\nu ^{}`$ at finite range, and we have an analytic estimate for $`\overline{\nu }`$ (which has not been derived previously to our knowledge). It is also worth remarking that $`\nu ^{}`$ is not the Lyapunov exponent itself (just as $`\overline{\nu }`$ is not), but rather only an upper bound. By analogy with the behavior of $`\overline{\nu }`$ stated at the end of the last subsection, $`\nu ^{}`$ will turn out numerically to be about double the actual $`\nu _L`$.
## IV Fluctuations
The ocean is not infinite in extent, and so the distribution of the stability exponents, $`\nu `$ (or $`|\text{Tr}(M)|`$) at a specified range $`r`$, is more directly relevant to the ocean acoustics problem than the Lyapunov exponent, $`\nu _L`$ (or $`\mathrm{exp}(\nu _Lr)`$). In order to visualize the magnitude of the fluctuations we are discussing, Fig. 4 displays the $`\mathrm{ln}|\text{Tr}(M)|`$ for eight of the rays from Fig. 3 as a function of range out to $`7.5\nu _L^1`$. By the right end of the figure, for any fixed range one ray might have a $`|\mathrm{Tr}(M)|=\mathrm{e}^{13}`$, and another one might have $`|\mathrm{Tr}(M)|=\mathrm{e}^0`$. At $`7.5\nu _L^1`$, there exist fluctuations in the stabilities of at least six orders of magnitude which is characteristic of broad tailed densities.
To characterize the fluctuations more quantitatively, we consider the cumulative densities for $`\nu `$ and $`|\text{Tr}(M)|`$. An initial working hypothesis might be to check whether at some long, fixed range, a diagonal element of $`M`$, say $`m_{ii}`$, is distributed as a Gaussian random variable across the ensemble of $`V(y;r)`$ and derive the implied cumulative densities from there. However, there ought to be an identifiable mechanism for a central limit theorem (CLT) to be operating with respect to $`m_{ii}`$. From Eq. (10), one can deduce that $`M`$ can be decomposed into a product of shorter range stability matrices. For very long $`r`$, consider a range $`\mathrm{\Delta }r`$ which is short compared to the final range r, yet long compared with $`\nu _L^1`$. Let $`N\mathrm{\Delta }r=r`$ where $`N`$ is large. Then it follows that the stability matrix is given by the left-ordered product
$$M=\underset{l=1}{\overset{N}{}}M_l,$$
(28)
where $`M_l`$ is the stability matrix for the range $`l\mathrm{\Delta }r`$ to range $`(l1)\mathrm{\Delta }r`$. To a great degree of accuracy the set of $`M_l`$ should behave independently with the only correlations being amongst the matrix elements necessary for maintaining unit determinant. The stability matrix should have the statistical properties of an ensemble of products of uncorrelated, random matrices.
If there exists a limiting form for a distribution at long range $`r`$, one would expect the same form (with different parameters, i.e. mean, variance) at $`r/2`$. In other words, the limiting form would have to be invariant under the matrix multiplication process. Denoting $`m_{l,ij}`$ as the matrix elements of $`M_l`$, for the $`N=2`$ case, we have
$$m_{11}=m_{2,11}m_{1,11}+m_{2,12}m_{1,21}$$
(29)
If the $`m_{l,ij}`$ behave as independent, random Gaussian variables, then $`m_{11}`$ could not be Gaussian because of the product form. The applicability of a CLT results from an additive process involving random variables. Instead, we anticipate something closer to a lognormal density because the log of a product of random variables acts like a sum of random variables. It should be mentioned here that this concept has been in use in many problems involving statistical physics .
To test whether $`|\text{Tr}(M)|`$ is lognormally distributed, we calculate $`100,000`$ rays through $`5`$ realizations of $`V(y;r)`$ to $`7.5\nu _L^1`$ ($`320L`$) (a reasonable upper bound for global acoustic propagation) for values of $`ϵ=2\times 10^3,3\times 10^3,\text{and}\mathrm{\hspace{0.33em}5}\times 10^3`$. If $`|\text{Tr}(M)|`$ is distributed lognormally, then $`\nu `$ is distributed in a Gaussian manner by definition. In Fig. 5, we plot the cumulative density for $`\nu `$. The corresponding analytic Gaussian form is superposed. It is impossible to distinguish the numerical results from the Gaussian form from this plot. A similar plot for $`|\mathrm{Tr}(M)|`$ carries little new information, and is not pictured in this paper, though we have verified its excellent consistency with a lognormal density as well. By plotting the differences between the numerical and analytical curves for three different ranges, we see in Fig. 6 that the consistency with a Gaussian density is excellent, and that as range increases the consistency of $`\rho _\nu (x)`$ with a Gaussian density improves. Note the small scale of the deviations. We have verified that they are roughly of the order of expected sample size errors for the curve at maximum range.
There are several relationships implied by the lognormal density that are straightforward to test. First, if we denote the variance of $`\nu `$ as $`\sigma _\nu ^2`$, then a relationship between $`\overline{\nu }`$ and $`\nu _0`$ can be derived. With
$`\rho _\nu (x)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2\pi \sigma _\nu ^2}}}\mathrm{exp}\left[{\displaystyle \frac{(x\nu _0)^2}{2\sigma _\nu ^2}}\right],`$ (30)
$`\overline{\nu }`$ $`=`$ $`{\displaystyle \frac{1}{2r}}\mathrm{ln}e^{2\nu r}`$ (31)
$`=`$ $`{\displaystyle \frac{1}{2r}}\mathrm{ln}({\displaystyle \frac{1}{\sqrt{2\pi \sigma _\nu ^2}}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}\mathrm{d}x\mathrm{exp}(2xr)`$ (33)
$`\mathrm{exp}[{\displaystyle \frac{(x\nu _0)^2}{2\sigma _\nu ^2}}])`$
$`=`$ $`r\sigma _\nu ^2+\nu _0.`$ (34)
Inverting this last relation for $`\sigma _\nu ^2`$, one obtains
$$\sigma _\nu ^2=\frac{\overline{\nu }\nu _0}{r}.$$
(35)
Both exponents $`\overline{\nu },\nu _0`$ were defined (see Eq. (14,16,19)) to be independent of $`r`$ to leading order; see the upper panel of Fig. 7 where $`\overline{\nu },\nu _0`$ are plotted as a function of range. The stochastic approximation for $`\nu ^{}`$ also given matches precisely the value of $`\overline{\nu }`$ implying that $`\mu =1/2`$. From numerical simulations, it turns out that $`\mu `$ is 0.466, but this number is poorly determined due to sample size errors. There is no discernible $`r`$-dependence in either $`\overline{\nu }`$ or $`\nu _0`$ beyond the scale at which the stochastic approximation begins to work for $`\overline{\nu }`$. They maintain a rather constant ratio of $`2.20`$ amongst themselves. The lognormal $`|\text{Tr}(M)|`$ density thus implies that the standard deviation of $`\rho _\nu (x)`$ approaches zero as $`r^{1/2}`$. Again there is excellent consistency; see the lower panel of Fig. 7 where $`\sigma _\nu `$ is plotted versus $`[(\overline{\nu }\nu _0)/r]^{1/2}`$. Thus, in the limit of $`r\mathrm{}`$, $`\rho _\nu (x)`$ goes to a $`\delta `$-density; all $`\nu `$ converge to the single value $`\nu _0`$. This value would also have to be the Lyapunov exponent from the definition in Eq. (13), and the Lyapunov exponent would be a constant for all trajectories independent of the specific realization of the potential or the initial conditions. It appears that the approach of $`\nu _0`$ to $`\nu _L`$ is so rapid as to warrant replacing $`\nu _0`$ with $`\nu _L`$ in all the formulae of this section. In fact, the lower panel of Fig. 7 actually incorporates our best value for $`\nu _L`$ and not $`\nu _0`$ as a function of range.
A consequence of the lognormal density for $`|\text{Tr}(M)|`$ is that the density of $`|\text{Tr}(M)|^\gamma `$ for any real $`\gamma `$ must also be lognormally distributed. This follows from the fact that $`\gamma \nu `$ would be Gaussian distributed with mean $`\gamma \nu _0`$ and variance $`\gamma ^2\sigma _\nu ^2`$, and $`\gamma \nu =\mathrm{ln}|\text{Tr}(M)|^\gamma /r`$. Thus $`\gamma `$ enters as a linear scale factor in the parameters that define the lognormal density. It is given by
$`\rho _{|\text{Tr}(M)|^\gamma }(x)`$ $`=`$ $`\sqrt{{\displaystyle \frac{1}{2\pi r(\overline{\nu }\nu _0)}}}{\displaystyle \frac{1}{|\gamma |x}}`$ (37)
$`\mathrm{exp}\left[{\displaystyle \frac{\left(\mathrm{ln}(x)/\gamma \nu _0r\right)^2}{2r(\overline{\nu }\nu _0)}}\right],x0.`$
Straightforward integration gives
$$|\text{Tr}(M)|^\gamma =\mathrm{exp}\left(\left[\gamma \nu _0+\gamma ^2(\overline{\nu }\nu _0)/2\right]r\right)$$
(38)
for its ensemble averaged value. Note that the $`\gamma =2`$ case for which the stochastic theory was worked out is the only one independent of $`\nu _0`$, and thus also $`\nu _L`$ in the large-$`r`$ limit. Using Eq. (38), a variety of estimates for the Lyapunov exponent can be constructed. For example, for $`r`$ large enough
$$\nu _L\frac{2}{r}\mathrm{ln}|\mathrm{Tr}(M)|\frac{1}{2r}\mathrm{ln}|\mathrm{Tr}(M)|^2$$
(39)
as given in . Another example would be
$$\nu _L\frac{1}{r}\mathrm{ln}|\mathrm{Tr}(M)|^2\frac{1}{4r}\mathrm{ln}|\mathrm{Tr}(M)|^4$$
(40)
etc.
Another interesting, rather curious consequence of the constant ratio of $`\overline{\nu }`$ to $`\nu _0`$ is that $`\overline{\nu }`$ does not approach $`\nu _L`$ in the $`r\mathrm{}`$ limit even though $`\rho _\nu (x)\delta (x\nu _0)`$. Care must be taken to perform the non-commuting operations of taking the infinite range limit and ensemble averaging in the correct order. Furthermore, the variation of $`|\text{Tr}(M)|`$ grows without bound as a function of range $`r`$, in spite of the fact that all the trajectories possess equal stability exponents in the limit $`r\mathrm{}`$. From Eq. (38), it follows that
$$\sigma _{|\text{Tr}(M)|}^2=\mathrm{e}^{\overline{\nu }r}\left(\mathrm{e}^{\overline{\nu }r}\mathrm{e}^{\nu _0r}\right)\mathrm{e}^{2\overline{\nu }r}$$
(41)
where the last form applies in the large-$`r`$ limit, even though $`\sigma _\nu ^2`$ is approaching zero.
Finally, we point out that a lognormal density has long tails and, as already noted, allows for many orders of magnitude fluctuations in $`|\text{Tr}(M)|`$. To return to the issue of intermittent-like rays, at any range, all rays whose corresponding $`|\text{Tr}(M)|`$ are less than some O(1) constant can be considered as intermittent-like. Values of $`e`$ or $`e^2`$ could be taken as criteria, for example. The equivalent criteria expressed for the maximum of $`\nu `$ would be $`O(r^1)`$. In the present model the proportion of intermittent-like rays approaches zero as $`r\mathrm{}`$, but for finite range the proportion of intermittent rays up to some maximum value $`|\text{Tr}(M)|=x`$ is given by the cumulative density
$`F_{|\text{Tr}(M)|}(x)`$ $`=`$ $`F_\nu \left({\displaystyle \frac{\mathrm{ln}(x)}{r}}\right)`$ (42)
$`=`$ $`{\displaystyle \frac{1}{2}}\left(1+\mathrm{erf}\left[{\displaystyle \frac{\mathrm{ln}(x)\nu _0r}{\sqrt{2r(\overline{\nu }\nu _0)}}}\right]\right),`$ (43)
where $`\mathrm{erf}(z)`$ is the error function of argument $`z`$. With the replacement of $`\nu _0`$ by $`\nu _L`$, this gives a very interesting, nontrivial connection between the Lyapunov exponent ($`\nu _L`$), $`\overline{\nu }`$, and the proportion of intermittent rays as a function of range. The validity of this expression is verified in Fig. 8. The behavior is just as predicted. The small deviations seen may be indicative of some slight non-lognormal behavior in the lower tail, or it may just be due to our not using best fit values of $`\overline{\nu }`$ and $`\nu _0`$. For long ranges, the proportion of intermittent-like rays decreases exponentially as $`a_0r^{1/2}\mathrm{exp}(b_0r)`$ where $`a_0`$ and $`b_0`$ can be deduced from the asymptotic properties of the error function, and this behavior is independent of the precise criterium used for the maximum desired $`|\text{Tr}(M)|`$. As can be seen in Fig. 8, $`10\%`$ of the initial ray density remains stable or nearly stable out to ranges of order $`5\nu _L^1`$. This $`10\%`$ of the initial acoustic energy is then only linearly sensitive to the fluctuating sound speed field, and since energy remains in coherent bundles (see Fig. 3), they will be expected to have a longer time coherence over repeated experiments as the environment evolves. Performing repeated experiments and applying coherent averaging as a filter to pick up this energy, one can imagine being able to use this apportionment of the initial acoustic pulse for acoustic tomography.
## V Discussion and Summary
Long-range, low-frequency sound propagation in the ocean has been previously investigated both as a problem of wave propagation through a random medium, and as a basis for tomography. Several outstanding quandaries remain that our results only begin to address: i) in the early arriving portion of a wave front, there seems to be more coherence and stability than would be expected from an analysis based on stochastic ray techniques common to the subject of WPRM; ii) one expects that as one moves from the weak focusing to strong focusing regimes (roughly speaking, from short range to long range), there should be a transition from lognormally distributed wave field intensities to Rayleigh distributed ones. Data analyses suggest that the lognormal densities extend well beyond the weak focusing regime, and the cross-over is not understood theoretically; and iii) related to the first item, given the presence of more stability than seems consistent with theoretical modeling, how valid are the underlying assumptions of tomographic inversions performed at long ranges?
Complete solutions to these problems are well beyond the scope of this paper. Nevertheless, we believe that our results form one cornerstone for their eventual resolution. As the ocean acoustic problem mainly involves refraction, and is in a wavelength regime that should be extremely well-suited to semiclassical analysis, we have focused our attention exclusively on a simple ray model inspired by the ocean. Our approach is from a dynamical systems perspective as opposed to a stochastic ray method. It has the advantage of being a more fundamental starting point in the sense that a system’s dynamics may determine where a stochastic ray method is appropriate, but a stochastic ray method just presupposes a certain randomness that may or may not actually exist in the system’s dynamics.
Our main concern is the ray stability properties that govern wave field amplitudes in semiclassical approximations. A follow-up study is underway to address the phases (classical actions and geometric indices), correlations amongst ray properties, and robustness, i.e. the generality and applicability to the ocean of our results. The stability matrix is our key analysis tool because it contains all the necessary information about how stable or unstable each ray is. The distribution of stabilities reflects on statistical properties arising in the study of WPRM whereas the existence or lack thereof of stable rays impacts tomographic inversion. We also note that studying the stability matrix has the utility of providing additional strong checks on one’s numerical integration techniques. Its determinant must remain unity.
We have carefully introduced several stability exponents depending upon whether ensemble averaging is taken before or after the logarithm (or at all), and whether the range is finite or the infinite range limit is taken. We have related them to the absolute value of the trace of the stability matrix which we have found to fluctuate to a high degree of consistency with a lognormal density; note that this also applies to the absolute value of individual matrix elements of the stability matrix. We have given a heuristic argument for this distribution, and are not aware of any known analytic derivations of this result.
An important consequence follows from the appearance of lognormally distributed stabilities, or equivalently Gaussian densities in the stability exponents (the logarithmic variables). As shown in Sect. IV \[see Eq. (37)\], any power of the stability matrix trace, or individual matrix elements, is also distributed lognormally. Thus, each individual contribution of an eigenray to the semiclassical approximation of the Green’s function has a magnitude fluctuating as a lognormal density. Further study is underway to determine theoretically the cross-over from lognormal wave field intensity distributions characteristic of the weak focusing limit to Rayleigh densities in the strong focusing limit. It is tempting to extrapolate our results to compute statistically relevant quantities such as the scintillation index (normalized variance of intensity) by using Eq. (38). Although one can immediately deduce that the normalized variance of intensity due to a single ray contribution grows exponentially with range with a e-folding scale of ($`\overline{\nu }\nu _0`$), one cannot infer anything about the scintillation index in the region where multipathing is important since the phase of each contributing ray must be incorporated into the calculation. Also, our work assumes one is at or beyond the regime of strong focusing. This is confirmed in the upper panel of Fig. 7, where $`\nu _0`$ is seen to converge at the range $`O(\nu _L^1)`$. Since the strong Markov assumption is valid for this problem, this range can be shown to be of the order where strong focusing occurs. Thus it is erroneous to compute the scintillation index from our work in the weak focusing regime ($`r\nu _L^1`$) where only one ray contributes to the intensity distribution.
All rays in the model possess identical Lyapunov exponents, and the finite-range, mean stability exponent, $`\nu _0`$, converges rapidly to it (see upper panel of Fig. 7). This follows from the finite range stability exponent acting as a Gaussian random variable with a standard deviation shrinking with range as $`r^{1/2}`$. This is also a consequence of the lognormal density, and not the single scale nature of the sound speed fluctuations per se. This behavior may be rather general (as general as the lognormal behavior), and it would be interesting to verify it in more realistic models. It is quite unlike the $`ϵ^{2/3}`$-scaling law for the stability exponents which should only apply to a model with a single correlation scale in range for the sound speed fluctuations.
The lognormal distribution has very broad tails. One typically observes stability matrix traces that fluctuate many orders of magnitude at a given range. Long after the appearance of highly unstable rays as a function of range, some stable or nearly stable rays will still be present. Their proportion decays essentially exponentially with range where the parameters are uniquely fixed by the Lyapunov exponent and the related stability exponent $`\overline{\nu }`$. However, they may be tomographically invertible, and relatively speaking, more important than their proportion would suggest. We have pointed out the distinctiveness of their behavior relative to unstable rays such as the way they twist about each other, and hang together as they propagate. Their collective properties appear to be highly correlated.
The $`r^{1/2}`$ behavior of the standard deviation of the stability exponent leads to a paradoxical situation in which all the rays possess a unique Lyapunov exponent, yet the exponentiated quantity, the matrix elements or trace of the stability matrix possess a divergent variance as the range approaches infinity. This illustrates dramatically the differences arising when non-commuting operations, i.e. ensemble averaging, taking the logarithm, and taking the infinite range limit, are interchanged. It is the stability matrix elements which are relevant to semiclassical approximations. So the individual terms in a summation over eigenrays will vary infinitely in their relative importance.
Acknowledgments
We gratefully acknowledge helpful discussions with M. G. Brown, F. D. Tappert, J. A. Colosi, and financial support from the Office of Naval Research. |
warning/0002/cond-mat0002101.html | ar5iv | text | # Gap deformation and classical wave localization in disordered two-dimensional photonic band gap materials
## I Introduction
Electromagnetic waves traveling in periodic dielectric structures will undergo multiple scattering. For the proper structural parameters and wave frequencies, all waves may backscatter coherently; the result is total inhibition of propagation inside the structure. Such structures are called photonic band gap (PBG) materials or photonic crystals, and the corresponding frequency ranges, for which propagation is not allowed, photonic band gaps or stop bands. PBG materials can be artificially made in one, two, or three dimensions. For example, a periodic lattice of dielectric spheres embedded in a different dielectric medium would work as a three-dimensional PBG material, for the proper choice of lattice symmetry, dielectric contrast, and sphere volume filling ratio. In two dimensions, a periodic array of parallel, infinitely long, dielectric cylinders could work as a two-dimensional PBG material, prohibiting propagation in a direction perpendicular to the cylinders’ axis for some frequency range(s). The absence of optical modes in a photonic band gap is often considered as analogous to the absence of electronic energy eigenstates in the semiconductor energy gap. The ability of PBG materials to modulate electromagnetic wave propagation, in a similar way semiconductors modulate the electric current flow, can have a profound impact in many areas in pure and applied physics. It is then of fundamental importance to study the effects of disorder on the transmission properties of such materials.
Besides the non-resonant, macroscopic Bragg-like multiple scattering, there is also a second, resonant mechanism, that contributes to the formation of the spectral gaps. This is the excitation of single scatterer Mie resonances . In a previous publication it was shown that for two-dimensional PBG materials, for the $`E`$ polarization scalar wave case (electric field parallel to the cylinders’ axis), these Mie resonances are analogous to the electronic orbitals in semiconductors. The idea of the linear combination of atomic orbitals (LCAO) method was extended to the classical wave case as a linear combination of Mie resonances (LCMR), leading to a successful tight-binding (TB) parameterization for photonic band gap materials. This moves the picture for the photon states, from a one analogous to the nearly free electron model, to the one analogous to the strongly localized electron whose transport is achieved only by hopping (tunneling) from atom to atom. Depending then on which mechanism is dominant for the formation of the photonic gaps, we expect different changes to the system’s properties when disorder is introduced. If the Bragg-like multiple scattering mechanism is the dominant one, the photonic gaps should close quickly with increasing disorder, while if it is the excitation of Mie resonances, the photonic gaps should survive large amounts of disorder, in a similar way the electronic energy gap survives in amorphous silicon.
In this paper we will use two ab initio numerical methods to study the effects of disorder on photonic gap formation and wave localization in two-dimensional PBG materials. The first is the finite difference time domain (FDTD) spectral method , from which we obtain the photonic density of states for an infinite, disordered PBG material, and the second is the transfer matrix technique , from which we obtain the transmission coefficient for a wave incident onto a finite slab of the disordered PBG material. From the transmission coefficient we can obtain the localization length for the photonic states of the disordered material . The study will be on both PBG material realizations (solid high dielectric cylinders in air and cylindrical air holes in high dielectric), for both wave polarizations, and it will incorporate three different disorder realizations: disorder in position, radius, and dielectric constant (these systems, though, will still be periodic on the average). We will find that only the case of solid dielectric cylinders in air with the wave $`E_z`$-polarized exhibits the behavior expected from the strongly localized photon picture, while for all other cases, the nearly free photon picture seems to be the dominant one.
## II Numerical methods
Electromagnetic wave propagation in lossless composite dielectric media is described by Maxwell’s equations
$`\mu {\displaystyle \frac{\stackrel{}{H}}{t}}=\stackrel{}{}\times \stackrel{}{E},`$ $`ϵ(\stackrel{}{r}){\displaystyle \frac{\stackrel{}{E}}{t}}=\stackrel{}{}\times \stackrel{}{H},`$ (1)
where the dielectric constant $`ϵ(\stackrel{}{r})`$ is a function of position. In two dimensions, the two independent wave polarizations are decoupled. We assume the variation of the dielectric constant, as well as the propagation direction, along the $`xy`$ plane, and so, the cylinders along the $`z`$ axis. One of the polarizations is with the electric field parallel to the $`z`$ axis and the magnetic field on the $`xy`$ plane ($`E_z`$ or TM polarized) and obeys a scalar wave equation. The other one with the magnetic field parallel to the $`z`$ axis and the electric field on the $`xy`$ plane ($`H_z`$ or TE polarized) and obeys a vector wave equation.
The first method we will use, to study disordered PBG materials, is the FDTD spectral method . In our FDTD scheme, we first discretize the $`xy`$ plane into a fine uniform grid. Each grid point is centered in a unit cell which is further discretized into a 10$`\times `$10 subgrid, on which an arithmetic average of the dielectric constant is performed. In our problem we will assume dispersionless and lossless materials. For the $`E_z`$ polarization case we define the electric field on this grid and the magnetic field on two additional grids, one tilted by $`(d/2,0)`$, on which $`H_y`$ is defined, and one tilted by $`(0,d/2)`$, on which we define $`H_x`$. $`d`$ is the side of the grid cell. The corresponding finite-difference equations for the space derivatives that are used in the curl operators are then central-difference in nature and second-order accurate. The electric and magnetic fields are also displaced in time by a half time step $`\mathrm{\Delta }t/2`$, resulting in a “leapfrog” arrangement and central-difference equations for the time derivatives as well. If one initialize the electric and magnetic fields at $`t=t_0`$ and $`t=t_0+\mathrm{\Delta }t/2`$ respectively, then updating the values of the electric field for each grid point $`(i,j)`$ at $`t=t_0+\mathrm{\Delta }t`$ is done by
$`E_z|_{i,j}^{t_0+\mathrm{\Delta }t}=E_z|_{i,j}^{t_0}`$ $`+`$ $`{\displaystyle \frac{\mathrm{\Delta }t}{dϵ_{i,j}}}(H_y|_{i+1/2,j}^{t_0+\mathrm{\Delta }t/2}H_y|_{i1/2,j}^{t_0+\mathrm{\Delta }t/2}`$ (2)
$``$ $`H_x|_{i,j+1/2}^{t_0+\mathrm{\Delta }t/2}+H_x|_{i,j1/2}^{t_0+\mathrm{\Delta }t/2})`$ (3)
where $`ϵ_{i,j}`$ is the averaged dielectric constant for the grid point $`(i,j)`$. Similar equations follow for updating the magnetic field components at $`t=t_0+3\mathrm{\Delta }t/2`$, then again Eq. (2) for $`E_z`$ at $`t=t_0+2\mathrm{\Delta }t`$ etc. This way the time evolution of the system can be recorded. For numerical stability and good convergence the number of grid points per wavelength $`\lambda /d`$ must be at least 20, and also $`\mathrm{\Delta }td/\sqrt{2}c`$, where $`c`$ the speed of light in vacuum. Similar equations, with the roles of the electric and magnetic fields interchanged, apply for the $`H_z`$ polarization case.
In order to find the eigenmodes of a particular periodic (or disordered) system, we first initialize the electric and magnetic fields in the unit cell (or a suitable supercell) using periodic boundary conditions: $`\stackrel{}{E}(\stackrel{}{r}+\stackrel{}{a})=e^{i\stackrel{}{k}\stackrel{}{a}}\stackrel{}{E}(\stackrel{}{r})`$ and similarly for $`\stackrel{}{H}(\stackrel{}{r})`$, where $`\stackrel{}{k}`$ is the corresponding Bloch wave vector and $`\stackrel{}{a}`$ the lattice vector. These fields must have nonzero projections with the modes in search. We choose a superposition of Bloch waves for the magnetic field and set zero the electric field:
$`\stackrel{}{H}(\stackrel{}{r})={\displaystyle \underset{\stackrel{}{g}}{}}\widehat{v}_\stackrel{}{g}e^{i(\stackrel{}{k}+\stackrel{}{g})\stackrel{}{r}+i\varphi _\stackrel{}{g}},`$ $`\stackrel{}{E}(\stackrel{}{r})=0,`$ (4)
where $`\varphi _\stackrel{}{g}`$ is just a random phase and the unit vector $`\widehat{v}`$ is perpendicular to both $`\stackrel{}{E}`$ and $`(\stackrel{}{k}+\stackrel{}{g})`$, ensuring that $`\stackrel{}{H}`$ is transverse and that $`\stackrel{}{}\stackrel{}{H}=0`$. Once the initial fields are defined, we can evolve them in time using the “leapfrog” difference equations, while recording the field values as a time series for some sampling points. As the electric fields “builds” up, some particular modes dominate while most are depressed, reflecting the underline lattice symmetries. Here we record only the $`E_z`$ field for the $`E_z`$ polarization case, and the $`H_z`$ field for the $`H_z`$ polarization. At the end of the simulation, the time series are Fourier transformed back into frequency space, and the eigenmodes $`\omega (\stackrel{}{k})`$ of the system appear as sharp peaks. The length of the simulation determines the frequency resolution while the time difference between successive recordings determines the maximum frequency considered. This method scales linearly with size: a larger system will still need the same number of time steps for the same frequency resolution, thus sometimes referred to also as an “order-N” method .
Here we will use this method to obtain the system’s density of states (DOS). If one chooses a large supercell instead of the unit cell, then for each $`\stackrel{}{k}`$ point inside it’s first Brillouin zone, the Fourier transformed time series will consist of a number of peaks. Adding all contributions from all $`\stackrel{}{k}`$’s will result to a smooth function for the DOS. This is in contrast to older methods that where using random fields as initial boundary conditions . Random initial fields will ensure the condition for nonzero projections to all of the system’s eigenmodes, but in order to get coupled with them during “built” up, a large simulation time is required. Furthermore, the produced DOS is not a smooth function of frequency, still consisting of a large collection of peaks, and thus being useful only as an indication for the existence of spectral gaps. In our method, the underline symmetries of the modes are already in the initial fields and so they couple easier with them. Also, the larger the supercell, the smaller is its first Brillouin zone, and so the smaller the frequencies we initialize through the various $`\stackrel{}{k}`$. This is why we can get smooth results even for very low frequencies. In Figs. 1 and 2, we show the calculated density of states for the case of solid dielectric cylinders in air and cylindrical air holes in dielectric respectively, both for a square lattice arrangement, and for both polarizations. Along with them we also plot the corresponding band structure as obtained with the plane wave expansion method. Our study is going to be based on these two photonic structures.
The second method we will use is the transfer matrix technique in order to obtain the transmission coefficient for a wave incident along the $`xy`$ plane on a slab (or a slice) of the photonic material. The slice is assumed uniform along the $`z`$ axis, and periodic along the $`x`$ direction through application of periodic boundary conditions, while in the $`y`$ direction it has a finite width $`L`$. In this method one first constructs the transmitted waves at one side of the slice and then integrates numerically the time-independent Maxwell’s equations to the other side. There, the waves are projected into incident and reflected waves, and so a value for the transmission coefficient $`T`$ can be obtained. Here, we are interested in the wave localization in disordered photonic band gap materials, and in particularly in the localization length $`\mathrm{}2L/\mathrm{ln}T`$.
A few remarks about the results of this method are in order. Waves with different incidence directions will have different reflection and transmission coefficients, so if one is looking for an average transmission, all directions should be included. It is shown, however, that there is also a large dependence on the surface plane along which the structure is cut. More specifically, a wave normally incident on a (1,0) surface will have different transmission characteristics than a wave incident with 45<sup>o</sup> on a (1,1) surface. This is because certain modes can not always get coupled with the incident wave. One should then also average for the two different surface cuts, otherwise it will not be a true average. This is shown in Figs. 3 and 4 where we plot the (1,0) and the (1,1) cuts, each with both incidence directions (normal and 45<sup>o</sup> with respect to the surface) averaged. We see that taken individually, none of them corresponds to the true gaps as shown in Figs. 1 and 2, but rather, to wider and generally displaced gaps. For example, in the $`E_z`$ polarization case in the first spectral gap (Figs. 3a and 4a), with the (1,0) cut, the incident waves fail to couple with the the M modes of the first band, while with the (1,1) cut, the incident waves fail to couple with the X modes of the second band.
This is expected to be lifted once disorder is introduced into our system, since the sense of direction will be somehow lost. Disorder can be introduced as a random displacement, a random change in the radius, or, a random change in the dielectric constant of the cylinders. It is not clear however what amount of disorder would be needed for this. We repeated the calculations for small enough amounts of disorder so that the spectral gaps, as found from the FDTD method, remain almost unchanged, for all three different disorder mechanisms. As seen in Figs. 3 and 4, indeed, in some cases the coupling is achieved. For example, for the first gap in the $`E_z`$-polarization case, with the (1,0) cut, the M modes of the first band are now coupled with the incident waves and appear in the transmission diagram. These could be easily mistaken for disorder-induced localized states entering the gap, but they are not, since for the values of disorder used, the first gap is virtually unchanged. On the other hand, with the (1,1) cut, the coupling to X modes of the second band is not yet achieved, still yielding a wrong picture for the gap. Increasing the disorder further will eventually destroy any sense of direction and there will be no distinction between the two cases. Figs. 3 and 4 will be useful as a guide of which results can be trusted and which can not, if one uses only one surface cut and small values of disorder. As a general rule, we can deduce that the (1,0) cut should be used for the $`E_z`$ polarization case, while the (1,1) cut would be better for the $`H_z`$ polarization case.
## III Results and discussion
We first looked into the spectral gaps’ dependence on disorder using the FDTD spectral method. Our system consisted of a 8$`\times `$8 supercell, each cell discretized into a 32$`\times `$32 grid. We studied two systems: a square lattice array of solid cylinders, with dielectric constant $`ϵ_a`$=10, in air ($`ϵ_b`$=1) with a filling ratio $`f`$=0.28%, and a square lattice array of air cylinders ($`ϵ_a`$=1) in dielectric material $`ϵ_b`$=10, with air filling ratio $`f`$=0.71%, as described in Figs. 1 and 2. We divided the supercell’s first Brillouin zone into 10$`\times `$10 grid, which for the irreducible part yields 66 different $`\stackrel{}{k}`$ points. For each particular disorder realization (i.e. disorder type) and disorder strength, we run the simulation for all these 66 $`\stackrel{}{k}`$’s. At each $`\stackrel{}{k}`$ however we use a different disordered configuration (i.e. a different seed in the random number generator), and so a large statistical sample is automatically included in our result. In each case the effective disorder is measured by the rms error of the dielectric constant $`<ϵ>`$, which is defined as
$`ϵ^2={\displaystyle \frac{1}{N}}{\displaystyle \underset{i=1}{\overset{N}{}}}(ϵ_i^dϵ_i^p)^2,`$ (5)
where the sum goes over all $`N=8\times 8\times 32\times 32=65536`$ grid points, $`ϵ_i^d`$ and $`ϵ_i^p`$ are the dielectric constants at site $`i`$ in the disordered and periodic case respectively, and $`<\mathrm{}>`$ means the average over different configurations (different $`\stackrel{}{k}`$’s in our case). In both settings (dielectric cylinders in air and vise versa) the filling ratio of the high dielectric material is similar, and so $`<ϵ>`$ is expected to have the same meaning and weight.
Four different disorder realizations are studied: 1) disorder in position, without though allowing any cylinders to overlap with each other, 2) disorder in position allowing cylinder overlapping to occur, 3) disorder in radius, and 4) disorder in dielectric constant (the last one only in the solid cylinder case). For each different realization we consider various disorder strengths, and thus different effective disorders $`<ϵ>`$, for which we record the upper and lower gap edges for the first two photonic band gaps (if they exist). Results are summarized in Figs. 5 and 6, for the solid and air cylinder cases respectively. We note that the $`E_z`$ polarization case for the solid cylinders is very different from all other cases: the gaps survive very large amounts of positional disorder, especially if no overlaps are allowed. In fact, once the disorder becomes large enough for overlaps to be possible, the gap quickly closes, as shown in Fig. 5. The actual DOS graphs for the two different realizations of the positional disorder are shown in Fig. 7, for three different values of the effective disorder. On the other hand, if the disorder is of the third or fourth kind, the gaps close very quickly, even for modest values of the effective disorder.
The picture is very different in the other cases, as seen in Fig. 6. The effect of the positional disorder is the same, independent of whether overlaps are allowed or not. This is most clearly seen in Fig. 8, where the actual DOS graphs are plotted for the air cylinder case for both polarizations, and for both positional disorder realizations. Allowing the air cylinders to overlap, though, means that the connectivity of the background material will break. Our results, thus, indicate that there is no connection between the connectivity of the background material and the formation of the spectral gaps in this 2D case. Most importantly, however, we note that the disorder in radius has a similar effect with that of the positional disorder in closing the gaps. In fact, it is also similar to the effect of the disorder in radius for the $`E_z`$-solid-cylinder case. So, in the case of air cylinders, the type of the disorder that is introduced into the system does not play a significant role, but rather, it is only the effective disorder (measured through the dielectric constant’s error function) that determines the effect on the spectral gaps. On the other hand, for the $`E_z`$-solid-cylinder case, the type of disorder plays a profound role: if the “shape” of the individual scatterer is preserved, the gaps can sustain large amounts of disorder, while if it is not preserved, the gaps collapse in a manner similar to the air cylinder case.
We next go over the localization length results, which were obtained with the transfer matrix technique. Here, our system consisted of a 3$`\times `$7 supercell (3 along the $`x`$ axis), with each cell discretized into a 18$`\times `$18 grid (a small supercell was used in order to ease the computation burden). In the $`x`$ direction we applied periodic boundary conditions, while in the $`y`$ direction the supercell was repeated 4 times, to provide a total length $`L`$ for the slab of $`L`$=28 unit cells. The structures studied are exactly the same as described before. The lattice was cut along one only symmetry direction, the (1,0), since for large disorders we expect all “hidden” modes to be coupled with the incident wave (in any case, we know from Figs. 3 and 4 which results can be completely trusted and which can not). For each disorder realization and strength, we used 11 different $`\stackrel{}{k}`$ values uniformly distributed between normal and 45<sup>o</sup> angle incidence, and for each $`\stackrel{}{k}`$ we used a different disordered configuration, so these will constitute our statistical sample. For each $`\stackrel{}{k}`$ we find the minimum transmission coefficient inside each gap, from which we find the minimum localization length, and then average over all $`\stackrel{}{k}`$’s, ie. $`\mathrm{}2L/<\mathrm{ln}T>`$ (in the periodic case we first averaged over $`T`$ in order to correctly account for different propagation directions, but in the highly disordered case it is not so much important any more, and so we just average over the localization lengths).
Our results are shown in Figs. 9 and 10 (because of the small statistical sample and the small supercell used, the data points appear very “noisy”, especially for large disorders). We note here, as well, the distinct difference between the $`E_z`$-solid-cylinder case for positional disorder and all other cases. Especially for the first spectral gap, the localization length not only remains unaffected by the disorder, but it even decreases (this is not an artifact of the averaging procedure). The first conclusion from this, is that the mechanisms responsible for the gap formation in this case are unaffected by the presence of positional disorder, and so they are definitely not macroscopic (long-range) in nature. The fact that the localization length decreases, is attributed to the coupling of more symmetry modes with the incident wave as the disorder increases (they provide a smaller $`\mathrm{}`$ to the average, as seen in Fig. 3a). This decrease should not be mistaken for additional localization induced by the disorder (the classical analog of Anderson localization in electrons), since the latter is macroscopic in nature, and does not apply for strongly localized waves. The decrease in the localization length continues until a fairly large disorder value, and then it increases to a saturation value (the dielectric error function can reach only up to some value for positional disorder). This saturation value is higher for the case where overlaps are allowed, but still is very small compared to other cases, so waves remain strongly localized.
All other cases, on the other hand, show a common pattern of behavior: photon states become quickly de-localized with increasing disorder. The localization length is increased, until the point where the localization induced by the disorder becomes dominant. After this it starts decreasing, until finally it reaches some saturation point. Note also that there is an almost quantitative agreement between some cases that was not really expected, eg. for the disorder in radius in the first gap with the wave $`E_z`$-polarized, for both lattice settings, as seen in Figs. 9a and 10a. Only the case of disorder in the dielectric constant seems to deviate, having very quickly a very large effect, with the localization length directly saturating to some constant value. So, for air cylinders in dielectric with any type of disorder, and for the $`E_z`$-solid-cylinder case with disorder that does not preserve the scatterer’s “shape”, the behavior under disorder is similar.
All these results can be understood if we adopt two different “pictures” for the photon states, depending on which is the dominant mechanism that is responsible for the formation of the spectral gaps in each case. The first is the “nearly free” photon picture, in which the gap forming mechanism is the non-resonant macroscopic Bragg-like multiple scattering, while the second is the “strongly localized” photon picture, in which the gap forming mechanism is the microscopic (short-range) excitation of single scatterer Mie resonances.
Sharp Mie resonances appear only for the solid cylinder case, and they can be thought as analogous to the atomic orbitals in semiconductors. Using this analogy, a tight-binding model, based on a linear combination of Mie resonances, was recently developed for the photonic states in the $`E_z`$-solid-cylinder case . But if a tight-binding model can give a satisfactory description of the photonic states, then it is expected that certain behavioral patterns found in semiconductors should apply in our case too. So, positional disorder should have a small effect on the gaps, in a similar way the energy gap survives in amorphous silicon. Also, changing the scatterer should have a similar effect as changing the atoms in the semiconductor, yielding a large amount of impurity modes that quickly destroys the gap. This pattern is definitely confirmed here for the $`E_z`$-solid-cylinder case. In this case, multiple scattering and interference can only help to make the gaps wider, but are definitely not decisive on the existence of a gap.
For the macroscopic Bragg-like multiple scattering mechanism, the lattice periodicity is a very important factor for the existence of a spectral gap. If it is destroyed, then coherence in the backscattered waves will be destroyed, and so will the spectral gaps. It is of small consequence the exact way that the periodicity is destroyed, and so different disorder realizations will have similar effects. Also, since the gaps close more easily, it will be easier to observe the localization induced on the waves by the disorder itself, ie. the classical analog of Anderson localization in electrons. All these are recognized in the case of air cylinders in dielectric.
Finally, in the $`H_z`$-solid-cylinder case, there were no gaps to begin with, and so we can have no results about it. However, sharp Mie resonances appear for this case as well, and if their excitation was the dominant scattering mechanism, a gap would be expected here as well. The difference with the $`E_z`$ is that the former is described by a vector wave equation, while the latter by a scalar one (and thus closer to the electronic case). The form of the wave equation must, then, be an important factor in determining the relative strength of the two gap forming mechanisms.
## IV Conclusions
We have shown that several results in periodic and random photonic band gap materials can be understood in terms of two distinct photonic states: (a) The “local” states, based on a single scatterer Mie resonance, with the multiple scattering playing a minor role; these states are more conveniently described in terms of an LCAO-type of approach and are the analog of the $`d`$-states in transition metals. “Local” photonic states appear in the case of high dielectric cylinders surrounded by a low-dielectric host and for $`E`$-polarized waves. (b) The “nearly free” photonic states, where Bragg-like multiple scattering is the dominant mechanism responsible for their appearance; these states are more conveniently described in terms of a pseudopotential-type of approach and are the analog of $`s`$ (or $`p`$) states in simple metals.
Each type of photonic states responds differently to the presence of disorder: For the “local” states case, the gap is robust as the periodicity is destroyed, and it is hardly affected by the disorder as long as the identity of each individual scatterer is preserved; however, if the shape, or other characteristics influencing the scattering cross section of each individual scatterer, is altered by disorder, the gap tends to disappear. On the other hand, for the “nearly free” states case, the gap is very sensitive and tends to disappear easily as the periodicity is destroyed.
## V Acknowledgments
Ames Laboratory is operated for the U. S. Department of Energy by Iowa State University under contract No. W-7405-ENG-82. This work was supported by the Director of Energy Research office of Basic Energy Science and Advanced Energy Projects. It was also supported by the Army Research Office, a E.U. grant, and a NATO grant. |
warning/0002/hep-th0002228.html | ar5iv | text | # 1 Introduction
## 1 Introduction
In the last few years there has been a renewed intense interest in gauged supergravity theories. The work on AdS/CFT (anti-de Sitter/conformal field theory) dualities in recent years has reaffirmed the importance of gauged supergravity theories in various dimensions to the understanding of the dynamics of M/superstring-theory . The best studied example of this duality is between the IIB superstring theory on the background manifold $`AdS_5\times S^5`$ with $`N`$ units of five-form flux through the five-sphere and $`4d`$, $`𝒩=4`$ super Yang Mills theory with gauge group $`SU(N)`$, which is a conformally invariant quantum field theory. In the limit of small string coupling and large $`N`$, the classical (i.e. tree level) IIB supergravity approximation becomes valid. The lowest lying Kaluza Klein modes of IIB supergravity on $`AdS_5\times S^5`$ are believed to form a consistent nonlinear truncation<sup>4</sup><sup>4</sup>4The consistency of the nonlinear truncation for a subsector of the scalar manifold has been shown recently . which is described by five-dimensional $`𝒩=8`$ gauged supergravity. Many aspects of the AdS/CFT correspondence, like eg. RG flows , can therefore be studied entirely within the framework of $`5d`$ gauged supergravity due to the lack of interference with the higher Kaluza-Klein modes.
On the other hand, five-dimensional, $`𝒩=2`$ gauged supergravity theories naturally occur as effective field theories in certain brane world scenarios based on heterotic M-theory compactifications . Since gauged supergravity theories typically also allow for AdS ground states, they have recently been discussed as a potential framework for embedding the Randall/Sundrum scenario into M/string theory.
Several attempts in this direction have been made (see eg. .) Many of them focused on what we will later call $`𝒩=2`$ “gauged Maxwell/Einstein theories” . It was found, however, that the scalar potentials of these theories are not of the right form to admit a supersymmetric embedding of Randall/Sundrum-type models . The question whether this is a generic feature of all gauged supergravity theories provides one of the motivations to study the potentials of more general gauged supergravity theories in five dimensions.
Recently, we have constructed the general gaugings of $`5d`$, $`𝒩=2`$ supergravity coupled to vector as well as tensor multiplets . This was an extension and generalization of earlier work on the gaugings of $`𝒩=2`$ supergravity coupled to vector multiplets only .
Starting point of our construction were the ungauged Maxwell/Einstein supergravity theories (MESGT’s) of ref. , which describe the coupling of Abelian vector multiplets to supergravity. These theories have a global symmetry group of the form $`SU(2)_R\times G`$, where $`G`$ is the subgroup of the isometry group of the scalar field target manifold that extends to a global symmetry group of the full Lagrangian, and $`SU(2)_R`$ denotes the R-symmetry group of the $`𝒩=2`$ supersymmetry algebra. In general, there are various ways to turn a subgroup of $`SU(2)_R\times G`$ into a local gauge group. We will use different names for these different possibilities :
We refer to theories in which $`U(1)_RSU(2)_R`$ is gauged as “gauged Maxwell/Einstein supergravity theories”.
In order to gauge a subgroup $`K`$ of $`G`$, a subset of the vector fields of the ungauged theory has to transform in the adjoint representation of $`K`$. If such a group $`K`$ exists, there are two possibilities:
(i) There are additional vector fields outside the adjoint of $`K`$ which transform nontrivially under $`K`$. These vector fields have to be dualized to “self-dual” antisymmetric tensor fields in order to perform the gauging of $`K`$ in a supersymmetric way . <sup>5</sup><sup>5</sup>5 We should note that the gauging of $`𝒩=8`$ Poincaré supergravity in $`5d`$ requires the dualization of twelve of the vector fields of the $`𝒩=8`$ Poincaré supermultiplet to self-dual tensor fields for completely analogous reasons.
(ii) If there are no vector fields outside the adjoint of $`K`$, or if the additional vectors are all singlets under $`K`$ (“spectator vector fields”), the gauging of $`K`$ proceeds in a straightforward way, and no tensor fields have to be introduced .
In order to distinguish between gaugings of $`U(1)_R`$ and $`K`$, we will refer to theories in which $`K`$ is gauged as “Yang-Mills/Einstein supergravity theories” (“with or without tensor fields”, depending on which of the possibilities (i) or (ii) is realized)<sup>6</sup><sup>6</sup>6We will use the term “Yang-Mills” also when $`K`$ is Abelian (as is the case for our examples in Section 4).
The most general gauging in this framework is then obviously a simultaneous gauging of $`U(1)_R`$ and $`K`$. For consistency with our terminology, we will sometimes use the term “gauged Yang-Mills/Einstein supergravity theories (with or without tensor multiplets)” for this type of gauging.
As for the scalar potentials that are introduced by these different types of gaugings, one makes the following observation:
(i) The gauging of $`U(1)_R`$ introduces a scalar potential, which in all known cases
a) either has a maximum that corresponds to an anti-de Sitter space or
b) vanishes identically or
c) has no critical points at all.
(ii) The gauging of $`K`$ introduces no potential when no vector fields have to be dualized to tensor fields.
(iii) If tensor fields have to be introduced, the gauging of $`K`$ introduces a scalar potential which is positive semidefinite and can therefore not lead to AdS vacua.
(iv) The simultaneous gauging of $`U(1)_R`$ and $`K`$ leads to a scalar potential which is simply the sum of the potentials that would result from the gaugings of $`U(1)_R`$ and $`K`$ alone. The critical points of this combined potential have not yet been fully investigated.
The purpose of this paper is to give an explicit example of a gauged Yang-Mills/Einstein supergravity theory with tensor fields which is simple enough to admit a complete analysis of its scalar potential. The model we discuss describes the coupling of one vector multiplet and one self-dual tensor multiplet (which contains two real tensor fields) to supergravity. The three scalar fields from the vector/tensor multiplets parametrize the space $`=SO(1,1)\times SO(2,1)/SO(2)`$, and the possible gauge groups are $`U(1)_R\times SO(2)`$ and $`U(1)_R\times SO(1,1)`$. We will find that the structure of the resulting scalar potentials is much richer than for gaugings without tensor fields.
The organization of the paper is as follows. Section 2 briefly summarizes the most general form of a gauged Yang-Mills/Einstein supergravity theory with tensor fields. Section 3 discusses some general properties of the scalar potentials of these theories. The ungauged MESGT with scalar manifold $`=SO(1,1)\times SO(2,1)/SO(2)`$, its $`U(1)_R\times SO(2)`$ and $`U(1)_R\times SO(1,1)`$ gaugings and the resulting scalar potentials are analyzed in section 4, which represents the main part of this paper. Section 5 discusses the generalization to the scalar manifolds $`SO(1,1)\times SO(n1,1)/SO(n1)`$, and Section 6 finally ends with some conclusions. An appendix summarizes the “very special geometry” of the ungauged $`=SO(1,1)\times SO(2,1)/SO(2)`$ theory.
## 2 Gauged Yang-Mills/Einstein supergravity with tensor fields
In this section, we briefly review the most relevant features of $`𝒩=2`$ gauged Yang-Mills/Einstein supergravity theories coupled to tensor multiplets . Unless otherwise stated, our conventions will coincide with those of ref. , where further details can be found. In particular, we will use the metric signature $`(++++)`$ and impose the ‘symplectic’ Majorana condition on all fermionic quantities.
The fields of the $`𝒩=2`$ supergravity multiplet are the fünfbein $`e_\mu ^m`$, two gravitini $`\mathrm{\Psi }_\mu ^i`$ ($`i=1,2`$) and a vector field $`A_\mu `$. An $`𝒩=2`$ vector multiplet contains a vector field $`A_\mu `$, two spin-$`1/2`$ fermions $`\lambda ^i`$ and one real scalar field $`\phi `$. The fermions of each of these multiplets transform as doublets under the $`USp(2)_RSU(2)_R`$ R-symmetry group of the $`𝒩=2`$ Poincaré superalgebra; all other fields are $`SU(2)_R`$-inert. A tensor field satisfying a 5-dimensional “self-duality” condition must necessarily be complex . We choose to work with the real and imaginary parts of the complex tensors. A self-dual $`𝒩=2`$ tensor multiplet contains such a pair of tensor fields, four spin-$`1/2`$ fermions (i.e. two $`SU(2)_R`$ doublets) and two scalars.
The general coupling of $`m`$ self-dual tensor multiplets to $`𝒩=2`$ gauged Yang-Mills/Einstein supergravity was given in . The field content of these theories is
$$\{e_\mu ^m,\mathrm{\Psi }_\mu ^i,A_\mu ^I,B_{\mu \nu }^M,\lambda ^{i\stackrel{~}{a}},\phi ^{\stackrel{~}{x}}\}$$
(2.1)
where
$`I`$ $`=`$ $`0,1,\mathrm{}n`$
$`M`$ $`=`$ $`1,2,\mathrm{}2m`$
$`\stackrel{~}{a}`$ $`=`$ $`1,\mathrm{},\stackrel{~}{n}`$
$`\stackrel{~}{x}`$ $`=`$ $`1,\mathrm{},\stackrel{~}{n},`$
with $`\stackrel{~}{n}=n+2m`$. Note that we have combined the ‘graviphoton’ with the $`n`$ vector fields of the $`n`$ vector multiplets into a single $`(n+1)`$-plet of vector fields $`A_\mu ^I`$ labelled by the index $`I`$. Also, the spinor and scalar fields of the vector and tensor multiplets are combined into $`\stackrel{~}{n}`$-tupels of spinor and scalar fields. The indices $`\stackrel{~}{a},\stackrel{~}{b},\mathrm{}`$ and $`\stackrel{~}{x},\stackrel{~}{y},\mathrm{}`$ are the flat and curved indices, respectively, of the $`\stackrel{~}{n}`$-dimensional target manifold $``$ of the scalar fields. The metric, vielbein and spin connection on $``$ will be denoted by $`g_{\stackrel{~}{x}\stackrel{~}{y}}`$, $`f_{\stackrel{~}{x}}^{\stackrel{~}{a}}`$ and $`\mathrm{\Omega }_{\stackrel{~}{x}}^{\stackrel{~}{a}\stackrel{~}{b}}`$, respectively. The $`SU(2)_R`$ index $`i`$ is raised and lowered with the antisymmetric metric $`\epsilon _{12}=\epsilon ^{12}=1`$ according to
$$X^i=\epsilon ^{ij}X_j,X_i=X^j\epsilon _{ji}.$$
The fermions $`\mathrm{\Psi }_\mu ^i`$ and $`\lambda ^{i\stackrel{~}{a}}`$ are $`U(1)_R`$-charged, whereas the fields $`\phi ^{\stackrel{~}{x}}`$, $`\lambda ^{i\stackrel{~}{a}}`$ and $`B_{\mu \nu }^M`$ carry charge under $`K`$.
Denoting the $`U(1)_R`$ and $`K`$ coupling constants by $`g_R`$ and $`g`$, respectively, the $`(U(1)_R\times K)`$ gauge covariant derivatives of these fields are as follows ($``$ denotes the ordinary spacetime covariant derivative)
$`𝔇_\mu \mathrm{\Psi }_\nu ^i`$ $``$ $`_\mu \mathrm{\Psi }_\nu ^i+g_RV_IA_\mu ^I\delta ^{ij}\mathrm{\Psi }_{\nu j}`$
$`𝔇_\mu \lambda ^{i\stackrel{~}{a}}`$ $``$ $`_\mu \lambda ^{i\stackrel{~}{a}}+g_RV_IA_\mu ^I\delta ^{ij}\lambda _j^{\stackrel{~}{a}}+gA_\mu ^IL_I^{\stackrel{~}{a}\stackrel{~}{b}}\lambda ^{i\stackrel{~}{b}}`$
$`𝔇_\mu \phi ^{\stackrel{~}{x}}`$ $``$ $`_\mu \phi ^{\stackrel{~}{x}}+gA_\mu ^IK_I^{\stackrel{~}{x}}`$
$`𝔇_\mu B_{\nu \rho }^M`$ $``$ $`_\mu B_{\nu \rho }^M+gA_\mu ^I\mathrm{\Lambda }_{IN}^MB_{\nu \rho }^N.`$ (2.2)
Here, $`K_I^{\stackrel{~}{x}}`$ are the Killing vector fields on $``$ that generate the subgroup $`K`$ of its isometry group. The $`\phi `$-dependent matrices $`L_I^{\stackrel{~}{a}\stackrel{~}{b}}`$ and the *constant* matrices $`\mathrm{\Lambda }_{IN}^M`$ are the $`K`$-transformation matrices of $`\lambda ^{i\stackrel{~}{a}}`$ and $`B_{\mu \nu }^M`$, respectively. The $`V_I`$ are some constants that define the linear combination of the vector fields $`A_\mu ^I`$ that is used as the $`U(1)_R`$-gauge field
$$A_\mu [U(1)_R]=V_IA_\mu ^I.$$
(2.3)
They have to be constrained by
$$V_If_{JK}^I=0,$$
(2.4)
with $`f_{JK}^I`$ being the structure constants of $`K`$ <sup>7</sup><sup>7</sup>7If there are spectator vector fields among the $`A_\mu ^I`$, the corresponding $`f_{IJ}^K`$ are just zero.
We denote the curls of the vector fields $`A_\mu ^I`$ by $`F_{\mu \nu }^I`$. The non-Abelian field strengths $`_{\mu \nu }^IF_{\mu \nu }^I+gf_{JK}^IA_\mu ^JA_\nu ^K`$ ($`I=0,1,\mathrm{}n`$) of the gauge group $`K`$ and the self-dual tensor fields $`B_{\mu \nu }^M`$ ($`M=1,2\mathrm{},2m`$) are grouped together to define the tensorial quantity $`_{\mu \nu }^{\stackrel{~}{I}}=(_{\mu \nu }^I,B_{\mu \nu }^M)`$ with $`\stackrel{~}{I}=0,1,\mathrm{},n+2m`$.
The Lagrangian is then given by (up to 4-fermion terms)
$`e^1`$ $`=`$ $`{\displaystyle \frac{1}{2}}R(\omega ){\displaystyle \frac{1}{2}}\overline{\mathrm{\Psi }}_\mu ^i\mathrm{\Gamma }^{\mu \nu \rho }𝔇_\nu \mathrm{\Psi }_{\rho i}{\displaystyle \frac{1}{4}}\stackrel{o}{a}_{\stackrel{~}{I}\stackrel{~}{J}}_{\mu \nu }^{\stackrel{~}{I}}^{\stackrel{~}{J}\mu \nu }`$ (2.5)
$`{\displaystyle \frac{1}{2}}\overline{\lambda }^{i\stackrel{~}{a}}\left(\mathrm{\Gamma }^\mu 𝔇_\mu \delta ^{\stackrel{~}{a}\stackrel{~}{b}}+\mathrm{\Omega }_{\stackrel{~}{x}}^{\stackrel{~}{a}\stackrel{~}{b}}\mathrm{\Gamma }^\mu 𝔇_\mu \phi ^{\stackrel{~}{x}}\right)\lambda _i^{\stackrel{~}{b}}{\displaystyle \frac{1}{2}}g_{\stackrel{~}{x}\stackrel{~}{y}}(𝔇_\mu \phi ^{\stackrel{~}{x}})(𝔇^\mu \phi ^{\stackrel{~}{y}})`$
$`{\displaystyle \frac{i}{2}}\overline{\lambda }^{i\stackrel{~}{a}}\mathrm{\Gamma }^\mu \mathrm{\Gamma }^\nu \mathrm{\Psi }_{\mu i}f_{\stackrel{~}{x}}^{\stackrel{~}{a}}𝔇_\nu \phi ^{\stackrel{~}{x}}+{\displaystyle \frac{1}{4}}h_{\stackrel{~}{I}}^{\stackrel{~}{a}}\overline{\lambda }^{i\stackrel{~}{a}}\mathrm{\Gamma }^\mu \mathrm{\Gamma }^{\lambda \rho }\mathrm{\Psi }_{\mu i}_{\lambda \rho }^{\stackrel{~}{I}}`$
$`+{\displaystyle \frac{i}{2\sqrt{6}}}\left({\displaystyle \frac{1}{4}}\delta _{\stackrel{~}{a}\stackrel{~}{b}}h_{\stackrel{~}{I}}+T_{\stackrel{~}{a}\stackrel{~}{b}\stackrel{~}{c}}h_{\stackrel{~}{I}}^{\stackrel{~}{c}}\right)\overline{\lambda }^{i\stackrel{~}{a}}\mathrm{\Gamma }^{\mu \nu }\lambda _i^{\stackrel{~}{b}}_{\mu \nu }^{\stackrel{~}{I}}`$
$`{\displaystyle \frac{3i}{8\sqrt{6}}}h_{\stackrel{~}{I}}\left[\overline{\mathrm{\Psi }}_\mu ^i\mathrm{\Gamma }^{\mu \nu \rho \sigma }\mathrm{\Psi }_{\nu i}_{\rho \sigma }^{\stackrel{~}{I}}+2\overline{\mathrm{\Psi }}^{\mu i}\mathrm{\Psi }_i^\nu _{\mu \nu }^{\stackrel{~}{I}}\right]`$
$`+{\displaystyle \frac{e^1}{6\sqrt{6}}}C_{IJK}\epsilon ^{\mu \nu \rho \sigma \lambda }\{F_{\mu \nu }^IF_{\rho \sigma }^JA_\lambda ^K+{\displaystyle \frac{3}{2}}gF_{\mu \nu }^IA_\rho ^J(f_{LF}^KA_\sigma ^LA_\lambda ^F)`$
$`+{\displaystyle \frac{3}{5}}g^2(f_{GH}^JA_\nu ^GA_\rho ^H)(f_{LF}^KA_\sigma ^LA_\lambda ^F)A_\mu ^I\}`$
$`+{\displaystyle \frac{e^1}{4g}}\epsilon ^{\mu \nu \rho \sigma \lambda }\mathrm{\Omega }_{MN}B_{\mu \nu }^M𝔇_\rho B_{\sigma \lambda }^N`$
$`+g\overline{\lambda }^{i\stackrel{~}{a}}\mathrm{\Gamma }^\mu \mathrm{\Psi }_{\mu i}W_{\stackrel{~}{a}}+g\overline{\lambda }^{i\stackrel{~}{a}}\lambda _i^{\stackrel{~}{b}}W_{\stackrel{~}{a}\stackrel{~}{b}}g^2P`$
$`{\displaystyle \frac{i\sqrt{6}}{8}}g_R\overline{\mathrm{\Psi }}_\mu ^i\mathrm{\Gamma }^{\mu \nu }\mathrm{\Psi }_\nu ^j\delta _{ij}P_0{\displaystyle \frac{1}{\sqrt{2}}}g_R\overline{\lambda }^{i\stackrel{~}{a}}\mathrm{\Gamma }^\mu \mathrm{\Psi }_\mu ^j\delta _{ij}P_{\stackrel{~}{a}}`$
$`+{\displaystyle \frac{i}{2\sqrt{6}}}g_R\overline{\lambda }^{i\stackrel{~}{a}}\lambda ^{j\stackrel{~}{b}}\delta _{ij}P_{\stackrel{~}{a}\stackrel{~}{b}}g_R^2P^{(R)}.`$
The transformation laws are (to leading order in fermion fields)
$`\delta e_\mu ^m`$ $`=`$ $`{\displaystyle \frac{1}{2}}\overline{\epsilon }^i\mathrm{\Gamma }^m\mathrm{\Psi }_{\mu i}`$
$`\delta \mathrm{\Psi }_\mu ^i`$ $`=`$ $`𝔇_\mu \epsilon ^i+{\displaystyle \frac{i}{4\sqrt{6}}}h_{\stackrel{~}{I}}(\mathrm{\Gamma }_\mu ^{\nu \rho }4\delta _\mu ^\nu \mathrm{\Gamma }^\rho )_{\nu \rho }^{\stackrel{~}{I}}\epsilon ^i+{\displaystyle \frac{i}{2\sqrt{6}}}g_RP_0\mathrm{\Gamma }_\mu \delta ^{ij}\epsilon _j`$
$`\delta A_\mu ^I`$ $`=`$ $`\vartheta _\mu ^I`$
$`\delta B_{\mu \nu }^M`$ $`=`$ $`2𝔇_{[\mu }\vartheta _{\nu ]}^M+{\displaystyle \frac{\sqrt{6}g}{4}}\mathrm{\Omega }^{MN}h_N\overline{\mathrm{\Psi }}_{[\mu }^i\mathrm{\Gamma }_{\nu ]}\epsilon _i+{\displaystyle \frac{ig}{4}}\mathrm{\Omega }^{MN}h_{N\stackrel{~}{a}}\overline{\lambda }^{i\stackrel{~}{a}}\mathrm{\Gamma }_{\mu \nu }\epsilon _i`$
$`\delta \lambda ^{i\stackrel{~}{a}}`$ $`=`$ $`{\displaystyle \frac{i}{2}}f_{\stackrel{~}{x}}^{\stackrel{~}{a}}\mathrm{\Gamma }^\mu (𝔇_\mu \phi ^{\stackrel{~}{x}})\epsilon ^i+{\displaystyle \frac{1}{4}}h_{\stackrel{~}{I}}^{\stackrel{~}{a}}\mathrm{\Gamma }^{\mu \nu }\epsilon ^i_{\mu \nu }^{\stackrel{~}{I}}+gW^{\stackrel{~}{a}}\epsilon ^i+{\displaystyle \frac{1}{\sqrt{2}}}g_RP^{\stackrel{~}{a}}\delta ^{ij}\epsilon _j`$
$`\delta \phi ^{\stackrel{~}{x}}`$ $`=`$ $`{\displaystyle \frac{i}{2}}f_{\stackrel{~}{a}}^{\stackrel{~}{x}}\overline{\epsilon }^i\lambda _i^{\stackrel{~}{a}}`$ (2.6)
with
$$\vartheta _\mu ^{\stackrel{~}{I}}\frac{1}{2}h_{\stackrel{~}{a}}^{\stackrel{~}{I}}\overline{\epsilon }^i\mathrm{\Gamma }_\mu \lambda _i^{\stackrel{~}{a}}+\frac{i\sqrt{6}}{4}h^{\stackrel{~}{I}}\overline{\mathrm{\Psi }}_\mu ^i\epsilon _i.$$
(2.7)
The various scalar field dependent quantities $`\underset{\stackrel{~}{I}\stackrel{~}{J}}{\overset{o}{a}}`$, $`h_{\stackrel{~}{I}}`$, $`h^{\stackrel{~}{I}}`$, $`h_{\stackrel{~}{I}}^{\stackrel{~}{a}}`$, $`h^{\stackrel{~}{I}\stackrel{~}{a}}`$ and $`T_{\stackrel{~}{a}\stackrel{~}{b}\stackrel{~}{c}}`$ that contract the different types of indices are already present in the corresponding *ungauged* MESGT’s and describe the “very special”geometry of the scalar manifold $``$ (see for details). The ungauged MESGT’s also contain a constant symmetric tensor $`C_{\stackrel{~}{I}\stackrel{~}{J}\stackrel{~}{K}}`$. If the gauging of $`K`$ involves the introduction of tensor fields, the coefficients of the type $`C_{MNP}`$ and $`C_{IJM}`$ have to vanish . The only components that survive such a gauging are thus $`C_{IJK}`$, which appear in the Chern-Simons-like term of (2.5), and $`C_{IMN}`$, which are related to the transformation matrices of the tensor fields by
$$\mathrm{\Lambda }_{IN}^M=\frac{2}{\sqrt{6}}\mathrm{\Omega }^{MP}C_{IPN}.$$
Here $`\mathrm{\Omega }^{MN}`$ is the inverse of $`\mathrm{\Omega }_{MN}`$, which is a (constant) invariant antisymmetric tensor of the gauge group $`K`$:
$$\mathrm{\Omega }_{MN}=\mathrm{\Omega }_{NM},\mathrm{\Omega }_{MN}\mathrm{\Omega }^{NP}=\delta _M^P.$$
(2.8)
The terms proportional to
$`W^{\stackrel{~}{a}}(\phi )`$ $`=`$ $`{\displaystyle \frac{\sqrt{6}}{8}}h_M^{\stackrel{~}{a}}\mathrm{\Omega }^{MN}h_N`$
$`W^{\stackrel{~}{a}\stackrel{~}{b}}(\phi )`$ $`=`$ $`W^{\stackrel{~}{b}\stackrel{~}{a}}(\phi )=ih_{}^{J[\stackrel{~}{a}}K_J^{\stackrel{~}{b}]}+{\displaystyle \frac{i\sqrt{6}}{4}}h^JK_J^{[\stackrel{~}{a};\stackrel{~}{b}]}`$ (2.9)
(the semicolon denotes covariant differentiation on the target space $``$) and the potential term
$$P(\phi )=2W_{\stackrel{~}{a}}W^{\stackrel{~}{a}}$$
(2.10)
are due to the presence of the tensor fields.
The supersymmetric gauging of the $`U(1)_R`$-factor, on the other hand, introduces the terms proportional to
$`P^{\stackrel{~}{a}}(\phi )`$ $`=`$ $`\sqrt{2}h^{\stackrel{~}{a}I}V_I`$ (2.11)
$`P_0(\phi )`$ $`=`$ $`2h^IV_I`$ (2.12)
$`P_{\stackrel{~}{a}\stackrel{~}{b}}(\phi )`$ $`=`$ $`{\displaystyle \frac{1}{2}}\delta _{\stackrel{~}{a}\stackrel{~}{b}}P_0+2\sqrt{2}T_{\stackrel{~}{a}\stackrel{~}{b}\stackrel{~}{c}}P^{\stackrel{~}{c}}`$ (2.13)
in (2.5) and (2) and leads to the scalar potential contribution
$$P^{(R)}(\phi )=(P_0)^2+P_{\stackrel{~}{a}}P^{\stackrel{~}{a}}.$$
(2.14)
## 3 Some general properties of the scalar potential
As summarized in the previous section, the simultaneous gauging of $`U(1)_RSU(2)_R`$ and a subgroup $`KG`$ of the isometry group $`G`$ of the vector/tensor multiplets moduli space $``$ leads to a scalar potential of the form
$$e^1_{pot}=g^2Pg_R^2P^{(R)},$$
(3.1)
where $`P^{(R)}`$ arises from the gauging of $`U(1)_R`$, whereas $`P`$ is nonzero if and only if some $`K`$-charged vector fields $`A_\mu ^M`$ had to be dualized to tensor fields $`B_{\mu \nu }^M`$ in order to perform the gauging of $`K`$ in a supersymmetric way. In the remainder we will write
$$P_{tot}:=P+\lambda P^{(R)},\text{ with }\lambda :=\frac{g_R^2}{g^2}$$
(3.2)
so that
$$e^1_{pot}=g^2P_{tot}.$$
(3.3)
The potentials $`P`$ and $`P^{(R)}`$ are given by
$`P`$ $`=`$ $`2W_{\stackrel{~}{a}}W^{\stackrel{~}{a}}`$ (3.4)
$`P^{(R)}`$ $`=`$ $`(P_0)^2+P_{\stackrel{~}{a}}P^{\stackrel{~}{a}}.`$
Using $`h_{\stackrel{~}{a}}^{\stackrel{~}{I}}h_{\stackrel{~}{J}}^{\stackrel{~}{a}}=\delta _{\stackrel{~}{J}}^{\stackrel{~}{I}}h^{\stackrel{~}{I}}h_{\stackrel{~}{J}}`$ , it is easy to verify that $`W^{\stackrel{~}{a}}`$ and $`P^{\stackrel{~}{a}}`$ are orthogonal:
$$W_{\stackrel{~}{a}}P^{\stackrel{~}{a}}=0.$$
Contracting $`\delta \lambda ^{i\stackrel{~}{a}}=0`$ with $`W^{\stackrel{~}{a}}`$ and $`P^{\stackrel{~}{a}}`$ then shows that an $`𝒩=2`$ supersymmetric ground state requires
$$W^{\stackrel{~}{a}}=P^{\stackrel{~}{a}}=0.$$
(3.5)
This implies, in particular, that the cosmological constant of an $`𝒩=2`$ supersymmetric vacuum is given by $`P^{(R)}(\phi _c)`$ alone, i.e. $`P(\phi _{c,susy})=0`$, as has also been pointed out in . Nevertheless, $`P`$ can still have a non-trivial effect on the *form* of a supersymmetric critical point, i.e. it can change it from a maximum to a saddle point. In addition, there might be critical points which do not preserve the full $`𝒩=2`$ supersymmetry and therefore *can* have $`P(\phi _c)0`$. We will see examples for all this in the next section.
Using
$$C_{\stackrel{~}{I}\stackrel{~}{J}\stackrel{~}{K}}h^{\stackrel{~}{K}}=h_{\stackrel{~}{I}}h_{\stackrel{~}{J}}\frac{1}{2}h_{\stackrel{~}{I}\stackrel{~}{a}}h_{\stackrel{~}{J}}^{\stackrel{~}{a}},$$
$`P`$ and $`P^{(R)}`$ can be expressed in a more compact form which will facilitate the analysis of the critical points:
$`P`$ $`=`$ $`{\displaystyle \frac{3\sqrt{6}}{16}}h^I\mathrm{\Lambda }_I^{MN}h_Mh_N`$ (3.6)
$`P^{(R)}`$ $`=`$ $`4C^{IJ\stackrel{~}{K}}V_IV_Jh_{\stackrel{~}{K}},`$ (3.7)
where we have defined
$`\mathrm{\Lambda }_I^{MN}`$ $``$ $`\mathrm{\Lambda }_{IP}^M\mathrm{\Omega }^{PN}={\displaystyle \frac{2}{\sqrt{6}}}\mathrm{\Omega }^{MR}C_{IRP}\mathrm{\Omega }^{PN}`$ (3.8)
$`C^{\stackrel{~}{I}\stackrel{~}{J}\stackrel{~}{K}}`$ $``$ $`\stackrel{}{a}^{\stackrel{~}{I}\stackrel{~}{I}^{}}\stackrel{}{a}^{\stackrel{~}{J}\stackrel{~}{J}^{}}\stackrel{}{a}^{\stackrel{~}{K}\stackrel{~}{K}^{}}C_{\stackrel{~}{I}^{}\stackrel{~}{J}^{}\stackrel{~}{K}^{}}`$ (3.9)
with $`\stackrel{}{a}^{\stackrel{~}{I}\stackrel{~}{J}}`$ being the inverse of $`\stackrel{}{a}_{\stackrel{~}{I}\stackrel{~}{J}}`$.
If $``$ is associated with a Jordan algebra <sup>8</sup><sup>8</sup>8We recall that the MESGT’s associated with Jordan algebras are those for which the cubic form defined by the symmetric tensor $`C_{\stackrel{~}{I}\stackrel{~}{J}\stackrel{~}{K}}`$ can be identified with the norm form of a Jordan algebra of degree three. , one has (componentwise)
$$C^{\stackrel{~}{I}\stackrel{~}{J}\stackrel{~}{K}}=C_{\stackrel{~}{I}\stackrel{~}{J}\stackrel{~}{K}}=\text{const.}$$
In this case, because of $`C^{IJM}=C_{IJM}=0`$, $`P^{(R)}`$ simplifies to
$$P^{(R)}=4C^{IJK}V_IV_Jh_K\text{(for the Jordan family)}$$
(3.10)
with *constant* $`C^{IJK}=C_{IJK}`$ and summation over $`K`$ instead of $`\stackrel{~}{K}`$.
The critical points of $`P^{(R)}`$ have been analyzed in for the purely $`U(1)_R`$-gauged MESGTs of the Jordan type. It was found that they are characterized by the “dual” element
$$V^{\mathrm{\#}\stackrel{~}{I}}\sqrt{\frac{2}{3}}C^{\stackrel{~}{I}\stackrel{~}{J}\stackrel{~}{K}}V_{\stackrel{~}{J}}V_{\stackrel{~}{K}}$$
(3.11)
of $`V_{\stackrel{~}{I}}`$. Three cases could be distinguished:
1. $`V^{\mathrm{\#}\stackrel{~}{I}}=0`$. In this case, the scalar potential $`P^{(R)}`$ vanishes identically, leading to Minkowski ground states with broken supersymmetry.
2. $`V^{\mathrm{\#}\stackrel{~}{I}}`$ is in the “domain of positivity” of the corresponding Jordan algebra $`J`$. In this case, there exists precisely one critical point, which sits at the unique global maximum of the scalar potential $`P^{(R)}`$ and corresponds to an Anti-de Sitter ground state with unbroken $`𝒩=2`$ supersymmetry and unbroken global Aut(J)-invariance, where Aut(J) denotes the automorphism group of the Jordan algebra $`J`$.
3. $`V^{\mathrm{\#}\stackrel{~}{I}}`$ is non-zero and not in the domain of positivity of $`J`$. In this case, the scalar potential $`P^{(R)}`$ has no critical points at all.
In order to get a better understanding as to whether and how the presence of the tensor field related potential $`P`$ changes this picture, we will analyze the simplest non-trivial example of a gauged Yang-Mills/Einstein supergravity theory with tensor multiplets in full detail in the next section.
## 4 The simplest nontrivial example: $`=SO(1,1)\times SO(2,1)/SO(2)`$
### 4.1 The ungauged theory
The *ungauged* MESGT with the scalar manifold $`=SO(1,1)\times SO(2,1)/SO(2)`$ allows the construction of two of the simplest non-trivial examples of a gauged Yang-Mills/Einstein supergravity theory with tensor multiplets. Let us consider this *ungauged* theory first. It belongs to the generic Jordan family<sup>9</sup><sup>9</sup>9The “generic Jordan family” consists of the MESGT’s with scalar manifolds of the form $`=SO(1,1)\times SO(n1,1)/SO(n1)`$. and describes the coupling of three Abelian vector multiplets to supergravity. Consequently, the field content is
$$\{e_\mu ^m,\mathrm{\Psi }_\mu ^i,A_\mu ^{\stackrel{~}{I}},\lambda ^{i\stackrel{~}{a}},\phi ^{\stackrel{~}{x}}\}$$
(4.1)
with
$`i`$ $`=`$ $`1,2`$
$`\stackrel{~}{I}`$ $`=`$ $`0,1,\mathrm{},3`$
$`\stackrel{~}{a}`$ $`=`$ $`1,\mathrm{},3`$
$`\stackrel{~}{x}`$ $`=`$ $`1,\mathrm{},3,`$
where the three scalar fields $`\phi ^{\stackrel{~}{x}}`$ parametrize the target space $`=SO(1,1)\times SO(2,1)/SO(2)`$. The latter can be described as the hypersurface
$$N(\xi )=\left(\frac{2}{3}\right)^{\frac{3}{2}}C_{\stackrel{~}{I}\stackrel{~}{J}\stackrel{~}{K}}\xi ^{\stackrel{~}{I}}\xi ^{\stackrel{~}{J}}\xi ^{\stackrel{~}{K}}=1$$
in a four-dimensional ambient vector space parametrized by coordinates $`\xi ^{\stackrel{~}{I}}`$. In the case at hand, this vector space can be identified with the Jordan algebra
$$J=\text{I}\text{R}\mathrm{\Sigma }_3,$$
where $`\mathrm{\Sigma }_3`$ is the Jordan algebra of degree two corresponding to a quadratic form $`Q`$ with signature $`(+,,)`$ . In the most natural basis of this Jordan algebra, $`N(\xi )`$ takes on the following form
$$N(\xi )=\sqrt{2}\xi ^0\left[(\xi ^1)^2(\xi ^2)^2(\xi ^3)^2\right],$$
where the normalization factor $`\sqrt{2}`$ ensures that the unique selfdual point $`\xi ^{\stackrel{~}{I}}=\xi ^{\mathrm{\#}\stackrel{~}{I}}`$ (i.e. the “basepoint” $`c^{\stackrel{~}{I}}`$ of the Jordan algebra ) really lies on the hypersurface $`N(\xi )=1`$, or equivalently, that there is a point on $``$ where $`\stackrel{}{a}_{\stackrel{~}{I}\stackrel{~}{J}}=\delta _{\stackrel{~}{I}\stackrel{~}{J}}`$ .<sup>10</sup><sup>10</sup>10In our parametrization, $`c^{\stackrel{~}{I}}=(\frac{1}{\sqrt{2}},1,0,0)`$, which corresponds to $`\phi ^{\stackrel{~}{x}}=(1,0,0)`$ (cf. the Appendix).
Hence, the nonvanishing $`C_{\stackrel{~}{I}\stackrel{~}{J}\stackrel{~}{K}}`$ are
$`C_{011}`$ $`=`$ $`{\displaystyle \frac{\sqrt{3}}{2}}`$
$`C_{022}`$ $`=`$ $`C_{033}={\displaystyle \frac{\sqrt{3}}{2}}.`$ (4.2)
The constraint $`N=1`$ can be solved by
$`\xi ^0`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}\phi ^2}}`$ (4.3)
$`\xi ^1`$ $`=`$ $`\phi ^1`$ (4.4)
$`\xi ^2`$ $`=`$ $`\phi ^2`$ (4.5)
$`\xi ^3`$ $`=`$ $`\phi ^3`$ (4.6)
with
$$\phi ^2(\phi ^1)^2(\phi ^2)^2(\phi ^3)^2.$$
Obviously, the hypersurface $`N=1`$ decomposes into three disconnected components:
(i) $`\phi ^2>0`$ and $`\phi ^1>0`$
(ii) $`\phi ^2<0`$
(iii) $`\phi ^2>0`$ and $`\phi ^1<0`$.
In the following, we will consider the “positive timelike” region (i) only, since in region (ii), $`g_{\stackrel{~}{x}\stackrel{~}{y}}`$ and $`\stackrel{}{a}_{\stackrel{~}{I}\stackrel{~}{J}}`$ are not positive definite (see the Appendix), and region (iii) is isomorphic to region (i).
All the scalar field dependent quantities in the Lagrangian and the supersymmetry transformation laws can be derived from $`N(\xi )`$, and they are listed in the Appendix.
### 4.2 The $`U(1)_R\times SO(2)`$ gauging
We will now turn the above ungauged $`SO(1,1)\times SO(2,1)/SO(2)`$ model into a gauged Yang-Mills/Einstein supergravity theory with tensor fields.
The isometry group of the scalar manifold $``$ is $`G=SO(2,1)\times SO(1,1)`$, which is simply the invariance group of $`N(\xi )`$. There are now two different ways to construct a Yang-Mills/Einstein supergravity theory with tensor multiplets: Either one gauges the compact subgroup $`SO(2)SO(2,1)`$, or one gauges the noncompact subgroup $`SO(1,1)SO(2,1)`$. We will focus on the compact gauging first and discuss the noncompact $`SO(1,1)`$ gauging in the next subsection. The $`SO(2)`$ subgroup of $`SO(2,1)`$ rotates $`\xi ^2`$ and $`\xi ^3`$ into each other and therefore acts nontrivially on the vector fields $`A_\mu ^2`$ and $`A_\mu ^3`$. Hence, gauging this $`SO(2)`$ requires the dualization of $`A_\mu ^2`$ and $`A_\mu ^3`$ to antisymmetric tensor fields. Accordingly, we decompose the index $`\stackrel{~}{I}`$ as follows
$$\stackrel{~}{I}=(I,M)$$
with $`I,J,K,\mathrm{}=0,1`$ and $`M,N,P,\mathrm{}=2,3`$.
It is easy to verify that our $`C_{\stackrel{~}{I}\stackrel{~}{J}\stackrel{~}{K}}`$ in eqs. (4.1) are consistent with the requirements $`C_{IJM}=C_{MNP}=0`$ for this type of gauging. Having a closer look at the $`C_{\stackrel{~}{I}\stackrel{~}{J}\stackrel{~}{K}}`$ of the type $`C_{IMN}`$ we also see that $`C_{1MN}`$ is zero, whereas $`C_{0MN}`$ is non-vanishing. This means, because of $`\mathrm{\Lambda }_{IN}^M\mathrm{\Omega }^{MP}C_{IPN}`$, that the vector field $`A_\mu ^0`$ plays the rôle of the $`SO(2)`$-gauge field, whereas $`A_\mu ^1`$ is just a “spectator vector field” with respect to the $`SO(2)`$-gauging.
In addition to this $`SO(2)`$-gauging, one can now use a linear combination $`A_\mu [U(1)_R]=A_\mu ^IV_I`$ of the vector fields $`A_\mu ^0`$ and $`A_\mu ^1`$ as the $`U(1)_R`$-gauge field, and simultaneously gauge $`U(1)_R`$ and $`SO(2)`$. The result is a *gauged* Yang-Mills/Einstein supergravity theory with tensor fields with the full gauge group $`U(1)_R\times SO(2)`$.
In our parametrization, the resulting potentials $`P`$ and $`P^{(R)}`$ (cf. eqs. (3.6), (3.10) and the Appendix) are found to be<sup>11</sup><sup>11</sup>11We are choosing $`\mathrm{\Omega }^{23}=\mathrm{\Omega }^{32}=1`$.
$`P`$ $`=`$ $`{\displaystyle \frac{1}{8}}{\displaystyle \frac{\left[(\phi ^2)^2+(\phi ^3)^2\right]}{\phi ^6}}`$ (4.7)
$`P^{(R)}`$ $`=`$ $`2\left[2\sqrt{2}{\displaystyle \frac{\phi ^1}{\phi ^2}}V_0V_1+\phi ^2(V_1)^2\right].`$ (4.8)
For the functions $`W_{\stackrel{~}{x}}`$, $`P_{\stackrel{~}{x}}`$ and $`P_0`$ that enter the supersymmetry transformation laws of the fermions, one obtains
$`W_1`$ $`=`$ $`0`$ (4.9)
$`W_2`$ $`=`$ $`{\displaystyle \frac{1}{4}}{\displaystyle \frac{\phi ^3}{\phi ^4}}`$ (4.10)
$`W_3`$ $`=`$ $`{\displaystyle \frac{1}{4}}{\displaystyle \frac{\phi ^2}{\phi ^4}},`$ (4.11)
respectively,
$`P_1`$ $`=`$ $`\sqrt{2}\left(\sqrt{2}{\displaystyle \frac{\phi ^1}{\phi ^4}}V_0V_1\right)`$ (4.12)
$`P_2`$ $`=`$ $`2{\displaystyle \frac{\phi ^2}{\phi ^4}}V_0`$ (4.13)
$`P_3`$ $`=`$ $`2{\displaystyle \frac{\phi ^3}{\phi ^4}}V_0`$ (4.14)
and
$$P_0=\frac{2}{\sqrt{3}}\left(\frac{V_0}{\phi ^2}+\sqrt{2}\phi ^1V_1\right).$$
(4.15)
This shows that the necessary condition for an $`𝒩=2`$ supersymmetric critical point, $`W_{\stackrel{~}{x}}(\phi _c)=P_{\stackrel{~}{x}}(\phi _c)=0`$, is equivalent to
$`\phi ^2`$ $`=`$ $`\phi ^3=0`$ (4.16)
$`\phi ^1^3V_1`$ $`=`$ $`\sqrt{2}V_0.`$ (4.17)
Let us now analyze the critical points of the above scalar potentials. We will first investigate the critical points of $`P`$ and $`P^{(R)}`$ separately and then consider the combined potential $`P_{tot}=P+\lambda P^{(R)}`$.
The critical points of $`𝐏`$:
Taking the deriative of $`P(\phi )`$ with respect to $`\phi ^{\stackrel{~}{x}}`$, one finds
$`P_{,1}`$ $`=`$ $`{\displaystyle \frac{3}{4}}{\displaystyle \frac{\left[(\phi ^2)^2+(\phi ^3)^2\right]}{\phi ^8}}\phi ^1`$ (4.18)
$`=`$ $`A\phi ^1+{\displaystyle \frac{\phi ^1}{4\phi ^6}}`$ (4.19)
$`P_{,2}`$ $`=`$ $`A\phi ^2`$ (4.20)
$`P_{,3}`$ $`=`$ $`A\phi ^3,`$ (4.21)
where
$$A\frac{3}{4}\frac{\left[(\phi ^2)^2+(\phi ^3)^2\right]}{\phi ^8}+\frac{1}{4\phi ^6}$$
has been introduced. There are now two possibilities:
Case 1: $`A0`$
Then $`P_{,2}=P_{,3}=0`$ implies $`\phi _c^2=\phi _c^3=0`$ (which then also implies $`P_{,1}=0`$). But then $`P(\phi _c)=0`$, and we have a Minkowski ground state, which, because of $`W_{\stackrel{~}{x}}(\phi _c)=0`$, preserves the full $`𝒩=2`$ supersymmetry (as long as the $`U(1)_R`$ gauging is turned off).
Case 2: $`A=0`$
Then $`P_{,1}=0`$ implies $`\frac{\phi ^1}{4\phi ^6}=0`$, which is inconsistent with $`\phi ^1>0`$ and $`\phi ^2>0`$.
Summary for $`P`$:
There exists a one parameter family of $`𝒩=2`$ supersymmetric Minkowski ground states, given by $`\phi ^2=\phi ^3=0`$ and arbitrary $`\phi ^1>0`$. These vacua also preserve the $`SO(2)`$-gauge invariance. There are no other critical points.
The critical points of $`𝐏^{(𝐑)}`$:
The gradient of $`P^{(R)}`$ is
$`P_{,1}^{(R)}`$ $`=`$ $`B\phi ^1{\displaystyle \frac{4\sqrt{2}V_0V_1}{\phi ^2}}`$ (4.22)
$`P_{,2}^{(R)}`$ $`=`$ $`B\phi ^2`$ (4.23)
$`P_{,3}^{(R)}`$ $`=`$ $`B\phi ^3,`$ (4.24)
where
$$B8\sqrt{2}\frac{\phi ^1}{\phi ^4}V_0V_1+4(V_1)^2.$$
There are now two possibilities:
Case 1: $`B=0`$
$`P_{,1}^{(R)}=0`$ then requires $`V_0V_1=0`$. Thus either $`V_0`$ or $`V_1`$ (or both of them) have to be zero. If $`V_0=0`$, $`B=0`$ implies $`V_1=0`$. Thus, $`B=0`$ automatically implies $`V_1=0`$, and the potential $`P^{(R)}`$ vanishes identically (cf. eq. (4.8)) resulting in a Minkowski vacuum. The $`U(1)_R`$-gauging is non-trivial only when at least one $`V_I`$ is non-zero. Since $`V_1=0`$ in the case at hand, a non-trivial $`U(1)_R`$ gauging requires $`V_00`$, implying $`P_10`$, ie. broken supersymmetry.
Case 2: $`B0`$
The vanishing of $`P_{,2}^{(R)}`$ and $`P_{,3}^{(R)}`$ then requires $`\phi ^2=\phi ^3=0`$, i.e. $`\phi ^2=\phi ^1^2`$. Because $`P_{,1}^{(R)}`$ has to vanish, this implies $`\phi ^1^3(V_1)^2=\sqrt{2}V_0V_1`$. Thus there are two possibilities: Either $`V_1=0`$, or $`V_0`$ and $`V_1`$ are both non-vanishing. The former case leads us back to the case of identically vanishing potential $`P^{(R)}0`$. The second possibility leads to a critical point with $`\phi ^2=\phi ^3=0`$ and
$$\phi ^1^3=\sqrt{2}\frac{V_0}{V_1}$$
(4.25)
whenever $`V_0V_1>0`$ (since $`\phi ^1>0`$). It is easy to see that this critical point satisfies the necessary conditions (4.16)-(4.17) for $`𝒩=2`$ supersymmetry. The value of the potential $`P^{(R)}`$ at this critical point is
$$P^{(R)}(\phi _c)=6(\phi _c^1)^2(V_1)^2<0,$$
(4.26)
i.e. it corresponds to an Anti-de Sitter ground state.
Summary for $`P^{(R)}`$:
There are three possibilities:
a) $`V_1=0`$.
This implies a flat potential $`P^{(R)}0`$ and Minkowski ground states with broken supersymmetry (supersymmetry is broken as long as the $`U(1)_R`$-gauging is non-trivial, ie. if $`V_00`$).
b) $`V_0V_1>0`$.
In this case, there exists exactly one critical point. It is given by $`\phi ^2=\phi ^3=0`$ and $`\phi ^1^3=\sqrt{2}V_0/V_1`$ and corresponds to an $`𝒩=2`$ supersymmetric AdS ground state whose cosmological constant can be read off from (4.26). This vacuum breaks the global symmetry group $`SO(1,1)\times SO(2,1)`$ down to its maximal compact subgroup $`SO(2)`$.
c) $`V_0V_1<0`$.
No critical points exist in this case.
It is instructive to recover the characterization of the critical points in terms of the dual element $`V^{\mathrm{\#}\stackrel{~}{I}}`$ mentioned in section 3. Using (3.11), one finds
$$V^{\mathrm{\#}\stackrel{~}{I}}=((V_1)^2/\sqrt{2},\sqrt{2}V_0V_1,0,0).$$
This shows that $`V^{\mathrm{\#}\stackrel{~}{I}}=0`$ is equivalent to $`V_1=0`$ and that $`V^{\mathrm{\#}\stackrel{~}{I}}`$ is in the domain of positivity iff $`V_0V_1>0`$ so that our cases a), b), c) are equivalent to the cases (i), (ii), (iii), respectively, listed in section 3.
The critical points of the combined potential $`𝐏_{\mathrm{𝐭𝐨𝐭}}=𝐏+\lambda 𝐏^{(𝐑)}`$:
The gradient of $`P_{tot}`$ is given by
$`P_{tot,1}`$ $`=`$ $`(A+\lambda B)\phi ^1+{\displaystyle \frac{\phi ^1}{4\phi ^6}}\lambda 4\sqrt{2}{\displaystyle \frac{V_0V_1}{\phi ^2}}`$ (4.27)
$`P_{tot,2}`$ $`=`$ $`(A+\lambda B)\phi ^2`$ (4.28)
$`P_{tot,3}`$ $`=`$ $`(A+\lambda B)\phi ^3,`$ (4.29)
where $`A`$ and $`B`$ are again as defined above.
There are now two possibilities:
Case 1: $`\phi ^2=\phi ^3=0`$.
In this case, $`P_{,\stackrel{~}{x}}(\phi _c)`$ vanishes automatically (see the discussion of $`P`$ above). This implies that $`P_{,\stackrel{~}{x}}^{(R)}(\phi _c)`$ also has to vanish separately, i.e. we are dealing with critical points that are just simultaneous critical points of the individual potentials $`P`$ and $`P^{(R)}`$. These have already been discussed above.
Case 2: $`\phi ^2^2+\phi ^3^2>0`$.
This case involves a nontrivial interplay of the two potentials $`P`$ and $`P^{(R)}`$. For $`P_{tot,2}`$ and $`P_{tot,3}`$ to vanish, one obviously needs $`A+\lambda B=0`$. $`P_{tot,1}=0`$ then implies
$$\frac{\phi ^1}{\phi ^4}=16\sqrt{2}\lambda V_0V_1.$$
(4.30)
This implies (remembering $`\lambda >0`$ and $`\phi ^1>0`$)
$$V_0V_1>0.$$
(4.31)
Inserting (4.30) into $`A+\lambda B=0`$, and reexpressing $`(\phi ^2)^2+(\phi ^3)^2`$ in terms of $`\phi ^2`$ and $`(\phi ^1)^2`$, one derives the additional condition
$$\frac{1}{\phi ^6}=\frac{1}{2}(16\sqrt{2}\lambda V_0V_1)^2+8\lambda (V_1)^2.$$
(4.32)
Now, by assumption, $`\phi ^2^2+\phi ^3^2>0`$. Hence
$$\frac{(\phi ^1)^2}{\phi ^8}>\frac{1}{\phi ^6},$$
so that in order for (4.30) and (4.32) to be consistent, one needs
$$32\lambda (V_0)^2>1.$$
(4.33)
Thus, if $`V_0`$ is big enough such that (4.33) is fulfilled and if $`V_1V_0>0`$ (cf. (4.31)), new non-trivial critical points exist. Eq. (4.32) fixes $`\phi _c^2`$ so that eq. (4.30) fixes $`\phi _c^1`$. This in turn fixes $`((\phi _c^2)^2+(\phi _c^3)^2)`$, but not $`\phi _c^2`$ and $`\phi _c^3`$ individually. Hence, we obtain a one-parameter family of critical points, which, because of $`((\phi _c^2)^2+(\phi _c^3)^2)>0`$, do not preserve the full $`𝒩=2`$ supersymmetry (cf. (4.16)) and spontaneously break the $`SO(2)`$-gauge invariance. Using (4.30) and (4.32), one finds for the value of $`P_{tot}`$ at these critical points
$$P_{tot}(\phi _c)=\frac{3}{8}\frac{1}{\phi ^4}<0,$$
(4.34)
which again corresponds to an Anti-de Sitter solution. Putting everything together, we arrive at the following
Summary for $`P_{tot}`$:
Depending on the values of the $`V_I`$, the total potential $`P_{tot}=P+\lambda P^{(R)}`$ admits the following types of critical points:
a) $`V_1=0`$.
In this case, $`P^{(R)}`$ vanishes identically, and one has a one-parameter family of $`SO(2)`$ gauge invariant Minkowski ground states. They are given by $`\phi _c^2=\phi _c^3=0`$ and an arbitrary $`\phi _c^1>0`$. If $`V_00`$ (i.e. if the $`U(1)_R`$-gauging is non-trivial), these ground states break the $`𝒩=2`$ supersymmetry. If $`V_0=0`$, the $`U(1)_R`$-gauging is switched off, and supersymmetry is unbroken, corresponding to case 1 in the discussion of $`P`$.
b1) $`V_0V_1>0`$, and $`32\lambda (V_0)^21`$.
In this case, there is precisely one ground state. It preserves the full $`𝒩=2`$ supersymmetry and the $`SO(2)`$ gauge invariance. It corresponds to an Anti-de Sitter solution, and is given by $`\phi _c^2=\phi _c^3=0`$ and $`(\phi _c^1)^3=\sqrt{2}V_0/V_1`$ with $`P_{tot}(\phi _c)=\lambda P^{(R)}(\phi _c)=6\lambda (\phi _c^1)^2(V_1)^2`$. Although the potential $`P`$ due to the tensor fields does not contribute to this cosmological constant, it does have an effect on the form of the extremum of the total potential: It is now a saddle point, as opposed to the case of pure $`U(1)_R`$ gauging, where the supersymmetric critical point is always a maximum.
b2) $`V_0V_1>0`$, and in addition $`32\lambda (V_0)^2>1`$
In this case, there are two types of critical points. The first one is an isolated supersymmetric critical point which has exactly the same properties as the one described in b1) above, with one exception: it is now a local *maximum* of the total scalar potential. Apart from this point, there is an additional one-parameter family of critical points. They are given by eqs. (4.30) and (4.32), which fix $`\phi _c^1`$, and $`[(\phi _c^2)^2+(\phi _c^3)^2]`$. They break the $`𝒩=2`$ supersymmetry and the $`SO(2)`$-gauge invariance and correspond to an Anti-de Sitter solution with $`P_{tot}(\phi _c)=3/(8\phi _c^4)`$. These critical points are saddle points of the total potential.
c) $`V_0V_1<0`$.
In this case, there are no critical points.
### 4.3 The $`U(1)_R\times SO(1,1)`$ gauging
We now come to the noncompact version of the above theory. Since the analysis is very similar to the compact case, our presentation can be less detailed.
We choose the $`SO(1,1)`$ subgroup of $`SO(2,1)`$ to rotate the components $`\xi ^1`$ and $`\xi ^2`$ into each other. Consequently, this $`SO(1,1)`$ acts nontrivially on the vector fields $`A_\mu ^1`$ and $`A_\mu ^2`$, and its gauging requires the dualization of $`A_\mu ^1`$ and $`A_\mu ^2`$ to antisymmetric tensor fields. Accordingly, we decompose the index $`\stackrel{~}{I}`$ as follows
$$\stackrel{~}{I}=(I,M)$$
with $`I,J,K,\mathrm{}=0,3`$ and $`M,N,P,\mathrm{}=1,2`$.
Since $`C_{0MN}0`$ and $`C_{3MN}=0`$, $`A_\mu ^0`$ plays the rôle of the $`SO(1,1)`$-gauge field, whereas $`A_\mu ^3`$ is a “spectator vector field” with respect to the $`SO(1,1)`$-gauging.
Using a linear combination $`A_\mu [U(1)_R]=A_\mu ^IV_I`$ of the vector fields $`A_\mu ^0`$ and $`A_\mu ^3`$ as the $`U(1)_R`$-gauge field, one can then simultaneously gauge $`U(1)_R`$ and $`SO(1,1)`$, and obtains the ($`U(1)_R\times SO(1,1)`$)-gauged analog of the ($`U(1)_R\times SO(2)`$)-theory discussed in the previous subsection.
The scalar potentials $`P`$ and $`P^{(R)}`$ are now (we use $`\mathrm{\Omega }^{12}=\mathrm{\Omega }^{21}=1`$)
$`P`$ $`=`$ $`{\displaystyle \frac{1}{8}}{\displaystyle \frac{\left[(\phi ^1)^2(\phi ^2)^2\right]}{\phi ^6}}`$ (4.35)
$`P^{(R)}`$ $`=`$ $`2\left[2\sqrt{2}{\displaystyle \frac{\phi ^3}{\phi ^2}}V_0V_3\phi ^2(V_3)^2\right].`$ (4.36)
For the functions $`W_{\stackrel{~}{x}}`$, $`P_{\stackrel{~}{x}}`$ and $`P_0`$ that enter the supersymmetry transformation laws of the fermions, one obtains
$`W_1`$ $`=`$ $`{\displaystyle \frac{1}{4}}{\displaystyle \frac{\phi ^2}{\phi ^4}}`$ (4.37)
$`W_2`$ $`=`$ $`{\displaystyle \frac{1}{4}}{\displaystyle \frac{\phi ^1}{\phi ^4}}`$ (4.38)
$`W_3`$ $`=`$ $`0,`$ (4.39)
respectively,
$`P_1`$ $`=`$ $`2{\displaystyle \frac{\phi ^1}{\phi ^4}}V_0`$ (4.40)
$`P_2`$ $`=`$ $`2{\displaystyle \frac{\phi ^2}{\phi ^4}}V_0`$ (4.41)
$`P_3`$ $`=`$ $`\sqrt{2}\left(\sqrt{2}{\displaystyle \frac{\phi ^3}{\phi ^4}}V_0+V_3\right)`$ (4.42)
and
$$P_0=\frac{2}{\sqrt{3}}\left(\frac{V_0}{\phi ^2}+\sqrt{2}\phi ^3V_3\right).$$
(4.43)
This already shows that there can be no $`𝒩=2`$ supersymmetric critical point, because $`W_2`$ can never vanish.
Let us now come to the critical points of the scalar potentials. We will again first analyse the critical points of $`P`$ and $`P^{(R)}`$ separately and then consider the combined potential $`P_{tot}=P+\lambda P^{(R)}`$.
The critical points of $`𝐏`$:
For the gradient of $`P(\phi )`$ with respect to $`\phi ^{\stackrel{~}{x}}`$, one obtains
$`P_{,1}`$ $`=`$ $`\stackrel{~}{A}\phi ^1`$ (4.44)
$`P_{,2}`$ $`=`$ $`\stackrel{~}{A}\phi ^2`$ (4.45)
$`P_{,3}`$ $`=`$ $`\stackrel{~}{A}\phi ^3+{\displaystyle \frac{\phi ^3}{4\phi ^6}}`$ (4.46)
with
$$\stackrel{~}{A}\frac{3}{4}\frac{\left[(\phi ^1)^2(\phi ^2)^2\right]}{\phi ^8}+\frac{1}{4\phi ^6}.$$
Since $`\phi ^1`$ cannot vanish, $`P_{,1}=0`$ requires $`\stackrel{~}{A}=0`$
But then $`P_{,3}=0`$ implies $`\phi ^3=0`$. The assumption $`\stackrel{~}{A}=0`$ then leads to the contradiction $`1=3`$, and is therefore inconsistent.
Summary for $`P`$:
$`P`$ alone has no critical points at all.
The critical points of $`𝐏^{(𝐑)}`$:
The gradient of $`P^{(R)}`$ is
$`P_{,1}^{(R)}`$ $`=`$ $`\stackrel{~}{B}\phi ^1`$ (4.47)
$`P_{,2}^{(R)}`$ $`=`$ $`\stackrel{~}{B}\phi ^2`$ (4.48)
$`P_{,3}^{(R)}`$ $`=`$ $`\stackrel{~}{B}\phi ^3{\displaystyle \frac{4\sqrt{2}V_0V_3}{\phi ^2}},`$ (4.49)
where
$$\stackrel{~}{B}8\sqrt{2}\frac{\phi ^3}{\phi ^4}V_0V_3+4(V_3)^2.$$
Since $`\phi ^1`$ cannot vanish, $`P_{,1}^{(R)}=0`$ implies $`\stackrel{~}{B}=0`$. The condition $`P_{,3}^{(R)}=0`$ then implies $`V_0V_3=0`$. Assume $`V_30`$. Then $`V_0=0`$ would imply $`V_3=0`$ by virtue of $`\stackrel{~}{B}=0`$. Thus, $`V_3`$ has to vanish in any case if a critical point of $`P^{(R)}`$ is assumed to exist. However, $`P^{(R)}`$ then vanishes identically.
Summary for $`P^{(R)}`$:
A critical point of $`P^{(R)}`$ exists if and only if $`P^{(R)}`$ vanishes identically (which is equivalent to $`V_3=0`$).
It is easy to recover the characterization of the critical points of $`P^{(R)}`$ in terms of the dual element $`V^{\mathrm{\#}\stackrel{~}{I}}`$ mentioned in section 3. In the case at hand, one finds
$$V^{\mathrm{\#}\stackrel{~}{I}}=((V_3)^2/\sqrt{2},0,0,\sqrt{2}V_0V_3).$$
This shows that $`V^{\mathrm{\#}\stackrel{~}{I}}=0`$ is equivalent to $`V_3=0`$ and that $`V^{\mathrm{\#}\stackrel{~}{I}}`$ can never be in the domain of positivity if $`V_30`$. Thus, our results are consistent with the discussion given in section 3.
The critical points of the combined potential $`𝐏_{\mathrm{𝐭𝐨𝐭}}=𝐏+\lambda 𝐏^{(𝐑)}`$:
The gradient of $`P_{tot}`$ is given by
$`P_{tot,1}`$ $`=`$ $`(\stackrel{~}{A}+\lambda \stackrel{~}{B})\phi ^1`$ (4.50)
$`P_{tot,2}`$ $`=`$ $`(\stackrel{~}{A}+\lambda \stackrel{~}{B})\phi ^2`$ (4.51)
$`P_{tot,3}`$ $`=`$ $`(\stackrel{~}{A}+\lambda \stackrel{~}{B})\phi ^3+{\displaystyle \frac{\phi ^3}{4\phi ^6}}\lambda 4\sqrt{2}{\displaystyle \frac{V_0V_3}{\phi ^2}},`$ (4.52)
where $`\stackrel{~}{A}`$ and $`\stackrel{~}{B}`$ are again as defined above.
Since $`\phi ^1`$ cannot vanish, $`P_{tot,1}=0`$ requires $`(\stackrel{~}{A}+\lambda \stackrel{~}{B})=0`$. $`P_{tot,3}=0`$ then implies
$$\frac{\phi ^3}{\phi ^4}=16\sqrt{2}\lambda V_0V_3.$$
(4.53)
The analogous equation (4.30) in the compact gauging implied $`V_0V_1>0`$ (cf. eq. (4.31)). In the case at hand, however, eq. (4.53) does not imply any constraint for $`V_0V_3`$, because $`\phi ^3/\phi ^4`$ does not have to be positive, as opposed to $`\phi ^1/\phi ^4`$, which is always greater than zero.
Inserting (4.53) into $`\stackrel{~}{A}+\lambda \stackrel{~}{B}=0`$, one derives the additional condition (i.e. the analog of (4.32))
$$\frac{1}{\phi ^6}=\frac{1}{2}(16\sqrt{2}\lambda V_0V_3)^2+8\lambda (V_3)^2.$$
(4.54)
Since $`1/\phi ^6>0`$, the last equation implies the consistency conditions
$`V_3`$ $``$ $`0`$
$`32\lambda (V_0)^2`$ $`<`$ $`1.`$ (4.55)
(The analogous equation (4.33) in the compact gauging arose as a consistency condition of (4.30) and (4.32). However, it is easy to see that (4.53) and (4.54) do not imply any additional constraints on $`V_0`$ or $`V_3`$, so that eqs. (4.55) remain the only constraints on the $`V_I`$.)
For a given set of $`V_I`$ and $`\lambda `$ subject to (4.55), $`\phi ^2`$ is fixed by (4.54). This in turn fixes $`\phi ^3`$ and $`((\phi ^1)^2(\phi ^2)^2)`$ by virtue of eq. (4.53), but leaves the $`\phi ^1`$ and $`\phi ^2`$ otherwise undetermined. We thus obtain a one parameter family of critical points which can be viewed as the noncompact analog of the nontrivial non-supersymmetric critical points found for the compact gauging (i.e. the ones mentioned in case b2) in the discussion of $`P_{tot}`$). However, for the non-compact gauging, these critical points have very different physical properties. In particular, the total scalar potential becomes
$$P_{tot}(\phi _c)=3\lambda \phi ^2(V_3)^2[132\lambda (V_0)^2],$$
(4.56)
which is positive because of the condition (4.55) and therefore corresponds to a *de Sitter* rather than an Anti-de Sitter spacetime.
Summary for $`P_{tot}`$:
If $`V_30`$ and $`32\lambda (V_0)^2<1`$, there exists a one parameter family of critical points given by (4.53) and (4.54). They correspond to a de Sitter spacetime with $`P_{tot}(\phi _c)=3\lambda \phi ^2(V_3)^2[132\lambda (V_0)^2]>0`$ and break the $`𝒩=2`$ supersymmetry and the $`SO(1,1)`$ gauge invariance. There are no other critical points of the combined potential. In particular, neither the analog of the $`𝒩=2`$ supersymmetric critical point mentioned in case b1) and b2), nor the analogs of the Minkowski ground states mentioned in case a) in the summary for $`P_{tot}`$ in section 4.2, exist for the non-compact gauging.
## 5 The generic Jordan family of $`𝒩=2`$ gauged Yang-Mills/Einstein supergravity theories coupled to tensor multiplets
In the previous section we studied in detail the critical points of the potentials of the simplest non-trivial gauged Yang-Mills Einstein supergravity theories with tensor multiplets. The corresponding $`𝒩=2`$ MESGT belongs to the generic Jordan family and has the scalar manifold $`SO(1,1)\times SO(2,1)/SO(2)`$. The MESGT’s of the generic Jordan family have the scalar manifold $`SO(1,1)\times SO(n1,1)/SO(n1)`$. From the results of and the arguments given in the previous section it follows that any gaugeable subgroup $`K`$ of the isometry group with $`K`$-charged vectors dualized to tensor fields must be Abelian. Since the vector field $`A_\mu ^0`$ must be the gauge field it follows that one can only gauge $`SO(2)`$ or $`SO(1,1)`$ and have some $`K`$-charged tensor fields under them. We should also note that the gaugeable $`SO(1,1)`$ must be a subgroup of $`SO(n1,1)`$ and can not be the $`SO(1,1)`$ factor in the isometry group since all the vector fields are charged under the latter $`SO(1,1)`$. The $`SO(2)`$ gauge group is some diagonal subgroup of the maximal Abelian subgroup $`SO(2)_1\times SO(2)_2\times \mathrm{}\times SO(2)_p`$ of $`SO(n1,1)`$ (for $`n=2p+1`$ or $`n=2p+2`$). The gaugeable $`SO(1,1)`$ subgroup is unique modulo some $`SO(n1)`$ rotation.
After the gauging of the Abelian subgroup of the isometry group with the charged vectors dualized to tensor fields, the remaining vector fields can be used to gauge some non-Abelian subgroup $`S`$ of the full isometry group so long as they decompose as the adjoint plus some singlets of $`S`$. This non-Abelian gauging does not introduce any additional potential . A linear combination of the remaining $`S`$ singlet vector fields can then be used to gauge the $`U(1)_R`$ subgroup of the $`R`$-symmetry group $`SU(2)_R`$. The full potential of the $`K\times U(1)_R\times S`$ gauged Yang-Mills Einstein supergravity with $`K`$-charged tensor fields must have novel critical points of the type we discussed in the previous section since these theories can be truncated to the the simplest non-trivial model consistently.
There exist an infinite family of non-Jordan MESGT’s with the scalar manifold $`SO(n,1)/SO(n)`$ . For this family only the parabolic subgroup $`SO(n1)\times SO(1,1)T_{(n1)}`$, which is simply an “internal Euclidean group” in $`(n1)`$ dimensions times a dilatation factor, extends to a symmetry of the full action . The analysis of the possible gauge groups $`K`$ that involve a dualization of $`K`$-charged vectors to tensor fields is very similar to the generic Jordan case . In this case too one finds that only a one dimensional Abelian subgroup $`K`$ can be gauged with nontrivial tensor fields carrying charge under $`K`$. However, there is one crucial difference between the Jordan family and the non-Jordan family. For the non-Jordan family the tensor $`C_{\stackrel{~}{I}\stackrel{~}{J}\stackrel{~}{K}}`$ is not an invariant tensor of the full isometry group $`SO(n,1)`$ of the scalar manifold. As a consequence one finds that
$$C_{\stackrel{~}{I}\stackrel{~}{J}\stackrel{~}{K}}C^{\stackrel{~}{I}\stackrel{~}{J}\stackrel{~}{K}},$$
(5.57)
and the $`C^{\stackrel{~}{I}\stackrel{~}{J}\stackrel{~}{K}}`$ are no longer constant tensors but depend on the scalars.
## 6 Conclusions
In this paper we have analyzed the scalar potentials of the simplest examples of a gauged Yang-Mills/Einstein supergravity theory coupled to tensor multiplets.
Although not all the results we have derived for these examples may carry over to the most general gauged Yang-Mills/Einstein supergravity theory with tensor fields, they show that the scalar potentials of these theories can exhibit a much richer structure than the purely $`U(1)_R`$-gauged supergravity theories or the gauged Yang-Mills/Einstein supergravity theories *without* tensor fields. Our analysis revealed that even though the total potential is just a sum of the potentials that appear in the separate gaugings of $`K`$ and $`U(1)_R`$, there can be critical points of the total potential which would not be critical points of the individual potentials. In particular, we found that, for a certain parameter range (case b2) in Section 4.2), the ($`U(1)_R\times SO(2))`$ gauging leads to a new one-parameter family of non-supersymmetric critical points, which are saddle points of the total potential. These are accompanied by an isolated $`𝒩=2`$ supersymmetric *maximum*, which is already present in the purely $`U(1)_R`$ gauged theory without tensor fields. In another parameter range (case b1)), the novel non-supersymmetric one-parameter family of critical points disappears and the $`𝒩=2`$ supersymmetric critical point becomes a *saddle point* (and remains supersymmetric). In yet another parameter range (case a)), the theory has a one-parameter family of Minkowski ground states which break the $`𝒩=2`$ supersymmetry as long as the $`U(1)_R`$ gauging is nontrivial. If the $`U(1)_R`$ gauging is switched off, these critical points become supersymmetric.
The possible types of critical points are much more restricted for the non-compact $`U(1)_R\times SO(1,1)`$ gauging, which can have at most a one-parameter family of non-supersymmetric *de Sitter* ground states (which are presumably unstable). This is consistent with the experience from compact and non-compact gaugings of the $`𝒩=8`$ theory where a non-supersymmetric de Sitter critical point was found in the $`SO(3,3)`$-gauged version of the $`𝒩=8`$ theory.
In this paper we have not studied the critical points of the potential when one gauges a *non-Abelian* subgroup $`K`$ of the isometry group of the scalar manifold with tensor multiplets transforming in a nontrivial representation of $`K`$. Such gauge groups are possible for the magical Jordan $`𝒩=2`$ theories as well as for the infinite family of theories with $`SU(n)`$ isometries discussed in . The study of the critical points of these theories as well as those of the non-Jordan family will be the subject of a future investigation.
## Appendix A The “very special geometry” of the $`SO(1,1)\times SO(2,1)/SO(2)`$-model
This Appendix contains a list of the basic scalar field dependent quantities that enter the Lagrangian and the transformation laws of the ungauged and gauged $`SO(1,1)\times SO(2,1)/SO(2)`$-theory.
In our parametrization, the $`h^{\stackrel{~}{I}}=\sqrt{\frac{2}{3}}\xi ^{\stackrel{~}{I}}|_{N=1}`$ are
$$h^0=\frac{1}{\sqrt{3}\phi ^2},h^1=\sqrt{\frac{2}{3}}\phi ^1,h^2=\sqrt{\frac{2}{3}}\phi ^2,h^3=\sqrt{\frac{2}{3}}\phi ^3.$$
For the $`h_{\stackrel{~}{I}}=\frac{1}{\sqrt{6}}\frac{}{\xi ^{\stackrel{~}{I}}}N|_{N=1}`$ one obtains
$$h_0=\frac{1}{\sqrt{3}}\phi ^2,h_1=\frac{2}{\sqrt{6}}\frac{\phi ^1}{\phi ^2},h_2=\frac{2}{\sqrt{6}}\frac{\phi ^2}{\phi ^2},h_3=\frac{2}{\sqrt{6}}\frac{\phi ^3}{\phi ^2}.$$
The vector/tensor field metric $`\stackrel{}{a}_{\stackrel{~}{I}\stackrel{~}{J}}=\frac{1}{2}\frac{}{\xi ^{\stackrel{~}{I}}}\frac{}{\xi ^{\stackrel{~}{J}}}\mathrm{ln}N(\xi )|_{N=1}`$ turns out to be
$$\stackrel{}{a}_{\stackrel{~}{I}\stackrel{~}{J}}=\left(\begin{array}{cccc}\phi ^4& 0& 0& 0\\ 0& 2(\phi ^1)^2\phi ^4\phi ^2& 2\phi ^1\phi ^2\phi ^4& 2\phi ^1\phi ^3\phi ^4\\ 0& 2\phi ^1\phi ^2\phi ^4& 2(\phi ^2)^2\phi ^4+\phi ^2& 2\phi ^2\phi ^3\phi ^4\\ 0& 2\phi ^1\phi ^3\phi ^4& 2\phi ^2\phi ^3\phi ^4& 2(\phi ^3)^2\phi ^4+\phi ^2\end{array}\right).$$
This shows that the unique point with $`\stackrel{}{a}_{\stackrel{~}{I}\stackrel{~}{J}}=\delta _{\stackrel{~}{I}\stackrel{~}{J}}`$ corresponds to $`\phi ^{\stackrel{~}{x}}=(1,0,0)`$, as has been mentioned earlier.<sup>12</sup><sup>12</sup>12If we had chosen another normalization for $`N`$, ie. $`N(\xi )=a\xi ^0\left[(\xi ^1)^2(\xi ^2)^2(\xi ^3)^2\right]`$ for some $`a\text{I}\text{R}`$, $`\stackrel{}{a}_{00}`$ would have been $`a^2\phi ^4/2`$ with the other components unchanged. It is easy to see that only $`a=\sqrt{2}`$ can lead to a point where $`\stackrel{}{a}_{\stackrel{~}{I}\stackrel{~}{J}}=\delta _{\stackrel{~}{I}\stackrel{~}{J}}`$.
Finally, the metric $`g_{\stackrel{~}{x}\stackrel{~}{y}}`$ on $``$ reads
$$g_{\stackrel{~}{x}\stackrel{~}{y}}=\left(\begin{array}{ccc}4(\phi ^1)^2\phi ^4\phi ^2& 4\phi ^1\phi ^2\phi ^4& 4\phi ^1\phi ^3\phi ^4\\ 4\phi ^1\phi ^2\phi ^4& 4(\phi ^2)^2\phi ^4+\phi ^2& 4\phi ^2\phi ^3\phi ^4\\ 4\phi ^1\phi ^3\phi ^4& 4\phi ^2\phi ^3\phi ^4& 4(\phi ^3)^2\phi ^4+\phi ^2\end{array}\right).$$
For the determinants of $`\stackrel{}{a}_{\stackrel{~}{I}\stackrel{~}{J}}`$ and $`g_{\stackrel{~}{x}\stackrel{~}{y}}`$, one finds
$`det\stackrel{}{a}_{\stackrel{~}{I}\stackrel{~}{J}}`$ $`=`$ $`\phi ^2`$ (1.58)
$`detg_{\stackrel{~}{x}\stackrel{~}{y}}`$ $`=`$ $`3\phi ^6,`$ (1.59)
which shows that $`\stackrel{}{a}_{\stackrel{~}{I}\stackrel{~}{J}}`$ and $`g_{\stackrel{~}{x}\stackrel{~}{y}}`$ are positive definite and well-behaved throughout the entire “positive timelike” region (i) and that both are not positive definite in region (ii), where $`\phi ^2<0`$.
Acknowledgements: We would like to thank Eric Bergshoeff, Renata Kallosh, Andrei Linde and Toine van Proeyen for fruitful discussions. |
warning/0002/hep-th0002085.html | ar5iv | text | # UCTP101.00 Supersymmetric Wilson Lines and Loops, and Super Non-Abelian Stokes Theorem
## 1 Introduction
The notion of Wilson loop provides a systematic method of obtaining gauge invariant observables. Its standard applications range from particle phenomenology and lattice field theories to strings and topological gauge theories. More recently, in the context of the AdS/CFT correspondence , an interesting connection between Wilson loops in supersymmetric gauge theories and membranes in supergravity theories has been suggested . In view of this and other important developments in supersymmetric gauge theories, it is natural to ask whether the notions of Wilson line and Wilson loop permit a supersymmetric generalization. Some formal work in this direction was carried out early in the development of supersymmetric gauge theories . There are also some recent suggestions in the $`N=4`$ case . In contrast to these attempts, our aim is to construct supersymmetric Wilson lines and Wilson loops in terms of supersymmetric product integrals. For non-supersymmetric gauge theories, it has been shown recently that standard product integrals provide a natural framework for describing Wilson lines and Wilson loops. This is because they have a built-in feature for keeping track of the order of matrices in path ordered quantities. The main purpose of the present work is to extend these results to theories which involve supersymmetric matrices. Thus, our construction of supersymmetric Wilson lines and loops is the natural supersymmetric extension of the definition of their non-supersymmetric counter parts. This will permit us to give, among other things, an unambiguous proof of the supersymmetric version of the non-Abelian Stokes theorem.
To provide a supersymmetric generalization of the notions of Wilson line and Wilson loop in terms of product integrals, we must address a number of questions. The first among these has to do with the fact that in supersymmetric gauge theories, the superfields have values in a Grassmann algebra. To be able to explore the properties of these theories in terms of product integrals, we must first ensure that Grassmann valued product integrals exist. We address this question in Section 2, where we construct supersymmetric product integrals and explore their properties. In Section 3, we use this formalism to define supersymmetric Wilson lines and Wilson loops. In Section 4, we construct a surface integral representation for the supersymmetric Wilson loop, thus establishing the supersymmetric version of the non-Abelian Stokes theorem. As a further confirmation of this theorem, in Section 5, we show the gauge covariance of the surface integral representation of the super Wilson loop operator. Section 6 is devoted to concluding remarks.
## 2 Supersymmetric Product Integrals
Comprehensive accounts of ordinary product integrals and their applications exist in the literature . Here we mention in passing that the justification for the word “product” lies in the property that the product integral is to the product what the ordinary integral is to the sum and that one of their most common applications is in solving systems of linear differential equations of the form
$$𝐲^{}(𝐬)=A(s)𝐲(𝐱),𝐲(𝐬_\mathrm{𝟎})=𝐲_\mathrm{𝟎}.$$
(1)
The solution of this system can be constructed in terms of the limit of the finite ordered product : $`\mathrm{\Pi }_p(A)=_{k=1}^ne^{A(s_k)\mathrm{\Delta }s_k}`$. In this expression, $`\mathrm{\Delta }s_k=s_ks_{k1}`$ for $`k=1,\mathrm{},n`$, where $`\{s_0,s_1,\mathrm{}.,s_n\}`$ is a partition of the real interval $`[a,b]`$. In the limit of large $`n`$ and under suitable conditions, this ordered product leads to the definition of the product integral.
The properties of standard product integrals rest heavily on the Banach algebra structure of matrix valued functions . In supersymmetric theories, the corresponding matrices take values in a Grassmann algebra. Since product integrals are products of the exponentials of the matrix valued functions, and in a supersymmetric theory the exponents must necessarily belong to the even part of the Grassmann algebra, we expect intuitively that all the properties of standard product integrals can be extended to supersymmetric product integrals. To put this on firm mathematical foundation, we must specify a suitable norm on the Grassmann algebra, with respect to which supersymmetric matrices also acquire a Banach algebra structure.
The Banach algebra structure of the Grassmann algebra is well known . Consider for definiteness the finite dimensional Grassmann algebra generated by the anticommuting quantities $`\theta ^1,\theta ^2,\mathrm{},\theta ^p`$. In this case, a generic element of the algebra can be written as a linear combination of the products $`\theta ^{i_1}\theta ^{i_2}\mathrm{}\theta ^{i_k}`$, $`k=0,\mathrm{},p`$, with complex coefficients $`a_{i_1i_2\mathrm{}i_k}`$. As a complex vector space the Grassmann algebra is $`2^p`$ dimensional. A norm on the above vector space (more precisely a valuation of the algebra) can be defined as the sum of the moduli of the coefficients. For example in the Grassmann algebra generated by a single element $`\theta `$, the norm of the generic element $`x=a+b\theta `$, $`a,b𝐂`$, is $`x=|a|+|b|`$, with $`|a|,|b|`$ the complex moduli. From this definition, one can show that the norm of the product of any two elements $`x`$ and $`y`$ of the Grassmann algebra satisfies the inequality: $`xyxy`$. This result is true not only for the above simple example but for the general Grassmann algebra generated by $`\theta ^1,\theta ^2,\mathrm{},\theta ^p`$. It is also straightforward to show that this norm is complete. In other words, with respect to this norm, the Grassmann algebra becomes a Banach algebra. As we will see below, this allows us to extend to supersymmetric product integrals most of the theorems which apply to ordinary product integrals .
Having specified a suitable norm on the Grassmann algebra, we turn to the construction of supersymmetric product integrals and to the study of some of their basic properties.
###### Definition 1
Let $`\mathrm{\Gamma }:[a,b]𝐂_{n\times n}^{1|p}`$ be an $`n\times n`$ matrix valued function with entries in the complex superspace $`𝐂^{1|p}`$. Let $`P=\{s_0,s_1,\mathrm{},s_n\}`$ be a partition of the interval $`[a,b]`$, with $`\mathrm{\Delta }s_k=s_ks_{k1}`$ for all $`k=1,\mathrm{},n`$.
$`\mathrm{\Gamma }`$ is called a step function iff there is a partition $`P`$ such that $`\mathrm{\Gamma }`$ is constant on each open subinterval $`(s_{k1},s_k)`$, for all $`k=1,\mathrm{},n`$.
The point value approximant $`\mathrm{\Gamma }_P`$ corresponding to the function $`\mathrm{\Gamma }`$ and partition P is the step function taking the value $`\mathrm{\Gamma }(s_k)`$ on the interval $`(s_{k1},s_k]`$ for all $`k=1,\mathrm{},n`$.
If $`\mathrm{\Gamma }`$ is a step function, then we define the function $`E_\mathrm{\Gamma }:[a,b]𝐂_{n\times n}^{1|p}`$ by $`E_\mathrm{\Gamma }(x):=e^{\mathrm{\Gamma }(s_k)(xs_{k1})}\mathrm{}e^{\mathrm{\Gamma }(s_2)\mathrm{\Delta }s_2}e^{\mathrm{\Gamma }(s_1)\mathrm{\Delta }s_1}`$ for any $`x(s_{k1},s_k]`$, for all $`k=1,\mathrm{},n`$, and $`E_\mathrm{\Gamma }(a):=I`$.
Based on the product integral formalism developed for ordinary matrices , we want the functions $`E_\mathrm{\Gamma }`$ to converge to the product integral as the partition of $`[a,b]`$ is refined. For the proof of the existence of the supersymmetric product integral we need some preliminary results. We start with estimating the norm of $`E_\mathrm{\Gamma }`$. This requires one more ingredient, namely the norm of a Grassmann algebra valued matrix. This will be defined in analogy with that of ordinary matrices: for any $`n\times n`$ matrix $`\mathrm{\Gamma }`$ as above, we define
$$\mathrm{\Gamma }_M=\underset{x𝐂^{n|p}x_n1}{sup}\frac{\mathrm{\Gamma }x_n}{x_n},$$
(2)
where $`x_n`$ refers to the norm of $`x`$ as an element in $`(𝐂^{1|p})^n𝐂^{n|p}`$. For $`x=(x_1,\mathrm{},x_n)𝐂^{n|p}`$ we can define $`x_n`$ for example by $`x_n=_{i=1}^nx_i`$. Some clarifications are necessary at this point. The Grassmann algebra $`𝐂^{1|p}`$ is a $`𝐂`$ vector space, but it is not a field. As a result $`𝐂^{n|p}`$ is not a vector space over $`𝐂^{1|p}`$, but only a rank $`n`$ module. Though it is a $`𝐂`$ vector space, it has no canonical norm on it. With all these preparations we have:
| $`E_\mathrm{\Gamma }(x)_M`$ | $`=e^{\mathrm{\Gamma }(s_k)(xs_{k1})}\mathrm{}e^{\mathrm{\Gamma }(s_2)\mathrm{\Delta }s_2}e^{\mathrm{\Gamma }(s_1)\mathrm{\Delta }s_1}_M`$ |
| --- | --- |
| | $`e^{\mathrm{\Gamma }(s_k)(xs_{k1})}_M\mathrm{}e^{\mathrm{\Gamma }(s_1)\mathrm{\Delta }s_1}_Me^{_a^x\mathrm{\Gamma }(s)_M𝑑s}.`$ |
(3)
In summary, we have obtained the following result:
$$E_\mathrm{\Gamma }(x)_Me^{_a^x𝑑s\mathrm{\Gamma }(s)_M}.$$
(4)
As a final preparation, we prove the following lemma: Let $`\mathrm{\Gamma }_1,\mathrm{\Gamma }_2:[a,b]𝐂_{n\times n}^{1|p}`$ be step-functions. Then,
$$E_{\mathrm{\Gamma }_1}(x)E_{\mathrm{\Gamma }_2}(x)=E_{\mathrm{\Gamma }_2}(x)_a^x𝑑sE_{\mathrm{\Gamma }_2}^1(s)[\mathrm{\Gamma }_1(s)\mathrm{\Gamma }_2(s)]E_{\mathrm{\Gamma }_1}(s).$$
(5)
To prove this, we define $`G(x)=E_{\mathrm{\Gamma }_2}^1(x)E_{\mathrm{\Gamma }_1}(x)`$. It follows immediately that $`G(a)=I`$ and $`G(x)`$ is continuous, and differentiable except for the division points of the partitions associated to $`\mathrm{\Gamma }_1`$ and $`\mathrm{\Gamma }_2`$. As a result, except for the division points, we have
$$G^{}(x)=E_{\mathrm{\Gamma }_2}^1(x)[\mathrm{\Gamma }_1(x)\mathrm{\Gamma }_2(x)]E_{\mathrm{\Gamma }_1}(x).$$
(6)
The quantity $`G(x)`$ is continuous and is continuously differentiable on each open division subinterval. Then, using the fundamental theorem of calculus on the subintervals and piecing the results together, we get:
$$G(x)=I+_a^x𝑑sE_{\mathrm{\Gamma }_2}^1(s)[\mathrm{\Gamma }_1(s)\mathrm{\Gamma }_2(s)]E_{\mathrm{\Gamma }_1}(s).$$
(7)
Multiplication from the left by $`E_{\mathrm{\Gamma }_2}(x)`$ leads to Eq. (5).
We are now in a position to define the supersymmetric product integral, and prove its existence:
###### Definition-Theorem 1
Given a continuous function $`\mathrm{\Gamma }:[a,b]𝐂_{n\times n}^{1|p}`$ and a sequence of step functions $`\{\mathrm{\Gamma }_n\}`$, which converges to $`\mathrm{\Gamma }`$ in the $`L^1([a,b])`$ sense, then the sequence $`\{E_{\mathrm{\Gamma }_n}(x)\}`$ converges uniformly on $`[a,b]`$ to a matrix called the supersymmetric product integral of $`\mathrm{\Gamma }`$ over $`[a,b]`$.
To prove the existence of super product integrals, we must demonstrate the convergence of the sequence $`\{E_{\mathrm{\Gamma }_n}(x)\}`$. By the lemma given above, we have
$$E_{\mathrm{\Gamma }_n}(x)E_{\mathrm{\Gamma }_m}(x)=E_{\mathrm{\Gamma }_m}^1(x)_a^x𝑑sE_{\mathrm{\Gamma }_m}(s)[\mathrm{\Gamma }_n(s)\mathrm{\Gamma }_m(s)]E_{\mathrm{\Gamma }_n}(s).$$
(8)
We can estimate the norm of the left-hand-side (lhs) as follows:
$$E_{\mathrm{\Gamma }_n}(x)E_{\mathrm{\Gamma }_m}(x)_ME_{\mathrm{\Gamma }_m}^1(x)_M_a^x𝑑sE_{\mathrm{\Gamma }_m}(s)_M\mathrm{\Gamma }_n(s)\mathrm{\Gamma }_m(s)_ME_{\mathrm{\Gamma }_n}(s)_M.$$
(9)
Using Eq. (4), we can estimate the difference of the norms as
$$E_{\mathrm{\Gamma }_n}(x)E_{\mathrm{\Gamma }_m}(x)_Me^{2_a^b𝑑s\mathrm{\Gamma }_m(s)_M}e^{_a^b𝑑s\mathrm{\Gamma }_n(s)_M}_a^x𝑑s\mathrm{\Gamma }_n(s)\mathrm{\Gamma }_m(s)_M.$$
(10)
Since $`\{\mathrm{\Gamma }_n\}`$ converges to $`\mathrm{\Gamma }`$ in the $`L^1([a,b])`$ sense, the first two terms on the rhs are bounded, and the sequence $`\{\mathrm{\Gamma }_n\}`$ is Cauchy in the $`L^1([a,b])`$ sense. Accordingly, the rhs goes to zero as $`m,n\mathrm{}`$. Since the rhs is independent of $`x`$, $`\{E_{\mathrm{\Gamma }_n}\}`$ is uniformly Cauchy, hence uniformly convergent. This establishes the existence the supersymmetric product integral. To prove the uniqueness of the limit, we estimate the difference $`E_{B_n}(x)E_{C_n}(x)_M`$ for two sequences $`\{B_n\}`$ and $`\{C_n\}`$, converging to $`\mathrm{\Gamma }`$ in the $`L^1`$ sense. Proceeding as we did above, it is immediate that $`\{E_{B_n}\}`$ and $`\{E_{C_n}\}`$ have the same limit. This concludes the proof of the existence and uniqueness of the supersymmetric product integral.
The structure of the supersymmetric product integrals described above permits the generalization of some of the well-known theorems of product integration to the supersymmetric case. Here we give a summary of the results which are relevant to the proof of the supersymmetric non-Abelian Stokes’ theorem. The proofs and further discussion of these results will be given elsewhere .
We will follow the notation and the conventions of as much as possible. Let $`\mathrm{\Gamma }:[a,b]𝐂_{n\times n}^{1|p}`$ be a continuous Grassmann valued function. For any $`x[a,b]`$, we express the supersymmetric product integral from $`a`$ to $`x`$ as
$$F(x,a):=\underset{a}{\overset{x}{}}e^{\mathrm{\Gamma }(s)ds}.$$
(11)
Then, $`F`$ satisfies the integral equation:
$$F(x,a)=1+_a^x𝑑s\mathrm{\Gamma }(s)F(s,a).$$
(12)
It is also a solution of the initial value problem:
$$\frac{dF}{dx}(x,a)=\mathrm{\Gamma }(x)F(x,a),F(a,a)=I.$$
(13)
The determinant of a supersymmetric product integral is given by
$$det\left(\underset{a}{\overset{x}{}}e^{\mathrm{\Gamma }(s)ds}\right)=e^{_a^x\mathrm{Str}\mathrm{\Gamma }(\mathrm{s})\mathrm{ds}},$$
(14)
where Str stands for supertrace. The intuitive composition rule holds:
$$\underset{a}{\overset{b}{}}e^{\mathrm{\Gamma }(s)ds}=\underset{c}{\overset{b}{}}e^{\mathrm{\Gamma }(s)ds}\underset{a}{\overset{c}{}}e^{\mathrm{\Gamma }(s)ds}.$$
(15)
It is possible to differentiate with respect to the endpoints:
$$\frac{}{x}\left(\underset{y}{\overset{x}{}}e^{\mathrm{\Gamma }(s)ds}\right)=\mathrm{\Gamma }(x)\underset{y}{\overset{x}{}}e^{\mathrm{\Gamma }(s)ds},\frac{}{y}\left(\underset{y}{\overset{x}{}}e^{\mathrm{\Gamma }(s)ds}\right)=\underset{y}{\overset{x}{}}e^{\mathrm{\Gamma }(s)ds}\mathrm{\Gamma }(y).$$
(16)
The L-derivative of ordinary product integrals can be extended to super product integrals: for a non-singular differentiable Grassmann valued function $`\mathrm{\Gamma }:[a,b]𝐂_{n\times n}^{1|p}`$, we define
$$L\mathrm{\Gamma }(x):=\mathrm{\Gamma }^{}(x)\mathrm{\Gamma }^1(x),$$
(17)
where prime indicates differentiation with respect to $`x`$. Defining
$$P(x)=\underset{a}{\overset{x}{}}e^{\mathrm{\Gamma }(s)ds},$$
(18)
and using Eq. (13), we can extend the analog of the fundamental theorem of calculus to super product integrals:
$$\underset{a}{\overset{x}{}}e^{(LP)(s)ds}=P(x)P^1(a).$$
(19)
The proof of the super non-Abelian Stokes’ theorem given below will rely heavily on the contents of the next three theorems. The first one is the sum rule. With $`P(x)=_a^xe^{\mathrm{\Gamma }_1(s)ds}`$, we have
$$\underset{a}{\overset{x}{}}e^{[\mathrm{\Gamma }_1(s)+\mathrm{\Gamma }_2(s)]ds}=P(x)\underset{a}{\overset{x}{}}e^{P^1(s)\mathrm{\Gamma }_2(s)P(s)ds}.$$
(20)
The second one is the similarity rule:
$$P(x)\left(\underset{a}{\overset{x}{}}e^{\mathrm{\Gamma }_2(s)ds}\right)P^1(a)=\underset{a}{\overset{x}{}}e^{[LP(s)+P(s)\mathrm{\Gamma }_2(s)P^1(s)]ds}.$$
(21)
Finally, the third one is differentiation with respect to a parameter. Given a Grassmann valued function $`\mathrm{\Gamma }:[a,b]\times [c,d]𝐂_{n\times n}^{1|p}`$ satisfying proper differentiability conditions, and given $`P(x,y;\lambda )=_y^xe^{\mathrm{\Gamma }(s;\lambda )ds}`$, we have:
$$\frac{}{\lambda }P(x,y;\lambda )=_y^x𝑑sP(x,s;\lambda )\frac{\mathrm{\Gamma }}{\lambda }(s;\lambda )P(s,y;\lambda ).$$
(22)
## 3 Supersymmetric Wilson Lines and Loops
Our results for supersymmetric product integrals are fairly general. In this section, we will use them as a basis to provide a natural and mathematically sound definition of supersymmetric Wilson lines and loops. To this end, we introduce our notations in a manner which naturally arises in supersymmetric gauge theories. We focus on the supersymmetric Wilson loop first. Consider an oriented manifold $`M`$ and a closed path $`C`$ in $`M`$. For simplicity, we assume that the target space is a simply connected manifold $`M`$, i.e. $`\pi _1(M)=0`$. This insures that the loop may be taken to be the boundary of an orientable two dimensional surface $`\mathrm{\Sigma }`$ in $`M`$. It will be convenient to describe the properties of such a 2-surface in terms of local coordinates $`\sigma ^0=\tau `$ and $`\sigma ^1=\sigma `$. So, for the points of the manifold $`M`$, which lie on $`\mathrm{\Sigma }`$, we have $`x=x(\sigma ,\tau )`$.
Let, in standard two component spinor notation , the local coordinates of a superspace be given by $`z^M=(x^{\alpha \dot{\alpha }},\theta ^\alpha ,\theta ^{\dot{\alpha }})`$. Also let the components of a supersymmetric connection $`\mathrm{\Gamma }`$ be given by $`\mathrm{\Gamma }_M`$. In terms of local coordinates, the connection $`\mathrm{\Gamma }`$ is a Lie superalgebra valued superform, which can be expressed as $`\mathrm{\Gamma }=dz^M\mathrm{\Gamma }_M`$. From the point of view of covariance under supersymmetry transformations, it is more convenient to express $`\mathrm{\Gamma }`$ in a basis in which the exterior derivative operator $`d=dz^M_M`$ maps superfields to superfields . So, we shall work, instead, in the basis where $`d=e^AD_A`$, with $`D_A`$ the supersymmetric covariant derivative, and $`e^A(z)=dz^Me_M^A(z)`$. In this expression, $`e_M^A(z)`$ are the well-known super-beins. Thus, we have
$$\mathrm{\Gamma }(z)=dz^M\mathrm{\Gamma }_M(z)=e^A(z)\mathrm{\Gamma }_A(z).$$
(23)
To describe Wilson lines and Wilson loops, we need the pull-back of this quantity on the path $`C`$ in $`M`$, described by an intrinsic parameter $`s`$: $`x^{\alpha \dot{\alpha }}=x^{\alpha \dot{\alpha }}(s)`$, $`\theta ^\alpha =\theta ^\alpha (s)`$, and $`\theta ^{\dot{\alpha }}=\theta ^{\dot{\alpha }}(s)`$. In terms of the embedding map $`i:CM`$ we have:
$$\mathrm{\Gamma }(s)=i^{}\mathrm{\Gamma }(z)=_sz^M(s)\mathrm{\Gamma }_M(z(s)).$$
(24)
Similarly, to obtain the pull-back of $`\mathrm{\Gamma }`$ on the 2-surface, we use the supersymmetric vielbeins:
$$\mathrm{\Gamma }_a=v_a^A\mathrm{\Gamma }_A;v_a^M=_az^M;v_a^A=v_a^Me_M^A(z).$$
(25)
It is the quantity $`\mathrm{\Gamma }=\mathrm{\Gamma }(s)ds`$ or $`\mathrm{\Gamma }=\mathrm{\Gamma }_ad\sigma ^a`$ that we will identify with the matrix valued functions of the supersymmetric product integral formalism described above. The corresponding pull-backs of the components of the supersymmetric covariant derivative on the line and on the 2- surface are given, respectively, by $`\frac{}{s}`$ and
$$D_a=v_a^AD_A=\frac{}{\sigma ^a}=_a.$$
(26)
The components of the supersymmetric field strength $`F_{ab}`$ on the 2-surface can be computed in two different ways. The first method is the obvious pull-back of the target space supersymmetric field strength:
$$F_{ab}=v_a^Av_b^BF_{BA}=v_a^Mv_b^NF_{NM}.$$
(27)
The second method is to make use of the pulled-back connection $`\mathrm{\Gamma }_a`$ given above:
$$F_{ab}=_a\mathrm{\Gamma }_b_b\mathrm{\Gamma }_a+[\mathrm{\Gamma }_a,\mathrm{\Gamma }_b].$$
(28)
To show the consistency of the above two expresions, multiply both (27) and (28) with the wedge product of differential forms $`\frac{1}{2}d\sigma ^ad\sigma ^b`$ to get the corresponding field strength two-forms on the two-surface. Then, the consistecy amounts to showing that the two field strength expressions are equal. Since on the two-surface $`d\sigma ^av_a^M=dz^M`$, Eq. (27) becomes $`\frac{1}{2}dz^Mdz^NF_{NM}`$. Moreover, $`d\sigma ^a_a=d`$ on the two-surface, so that Eq. (28) becomes $`d\mathrm{\Gamma }\mathrm{\Gamma }^2`$. But this expression is equal to the previous one by definition .
Consider now the continuous map $`\mathrm{\Gamma }:[a,b]𝐑_{n\times n}^{1|4}`$, where the latter is an $`n`$ by $`n`$ matrix valued function, with entries in the superspace $`𝐑^{1|4}`$, corresponding to the pull-back on the path C. Then, we define the supersymmetric Wilson line in terms of a super product integral as follows:
$$𝒫e^{_a^b\mathrm{\Gamma }(s)𝑑s}\underset{a}{\overset{b}{}}e^{\mathrm{\Gamma }(s)ds},$$
(29)
where $`𝒫`$ indicates path ordering as defined by the super product integral on the right-hand-side. Anticipating that we will identify the closed path $`C`$ over which the Wilson loop is defined with the boundary of a 2-surface, it is convenient to work from the beginning with Wilson lines depending on a parameter. Define $`\mathrm{\Gamma }_a:[\sigma _0,\sigma _1]\times [\tau _0,\tau _1]𝐑_{n\times n}^{1|4}`$, where $`[\sigma _0,\sigma _1]`$ and $`[\tau _0,\tau _1]`$ are the range of the local coordinates on the two surface $`\mathrm{\Sigma }`$. For later convenience, we also define the following elementary supersymmetric Wilson lines:
$$P(\sigma ,\sigma _0;\tau )=\underset{\sigma _0}{\overset{\sigma }{}}e^{v_1^A\mathrm{\Gamma }_A(\sigma ^{};\tau )d\sigma ^{}},Q(\sigma ;\tau ,\tau _0)=\underset{\tau _0}{\overset{\tau }{}}e^{v_0^A\mathrm{\Gamma }_A(\sigma ;\tau ^{})d\tau ^{}}.$$
(30)
To prove the supersymmetric version of the non-Abelian Stokes theorem, we want to make use of super product integration techniques to express the super Wilson loop operator as an integral over a two dimensional surface bounded by the corresponding loop. For this purpose, we define the super Wilson loop operator as
$$W_s[C]=𝒫\mathrm{exp}(_C\mathrm{\Gamma }(\tau )𝑑\tau )e^{_Ci^{}(dz^M\mathrm{\Gamma }_M)}.$$
(31)
In this expression, as in Eq. (24), $`i^{}`$ denotes the pull-back of the embedding $`i:CM`$. We have written this expression in a notation familiar from the physics literature. It is to be understood, however, that the right-hand-side is to be composed of the super product integrals as given in Eq. (29) above. The expression for the supersymmetric Wilson loop depends on the homotopy class of the loop $`C`$ in $`M`$. We can, therefore, parameterize $`C`$ in any convenient manner consistent with its homotopy class. In particular, we can break up the closed path into piecewise continuous segments, along which either $`\sigma `$ or $`\tau `$ remains constant. The composition rule for super product integrals given by Eq. (15) ensures that this break up of the super Wilson loop into super Wilson lines does not depend on the intermediate points chosen on the closed path. Inspired by the typical paths which are used in the actual computations of of both ordinary and supersymmetric Wilson loops (see e.g. ), we break up the super Wilson loop into a product of four super Wilson lines. Using the same notation as in the non-supersymmetric case , we write
$$W_s[C]=W_4W_3W_2W_1.$$
(32)
In this expression, $`W_k`$, $`k=1,..,4`$, are super Wilson lines such that $`\tau =const.`$ along $`W_1`$ and $`W_3`$, and $`\sigma =const.`$ along $`W_2`$ and $`W_4`$. We emphasize that $`\sigma =const.`$ and $`\tau =const.`$ are arbitrary curves.
To see the advantage of parameterizing the closed path in this manner, consider the exponent of Eq. (31). Along each segment, only one of the terms is non-vanishing. For example, along the segment $`[\sigma _0,\sigma ]`$, we have $`\tau ^{}=\tau _0=const`$. As a result, we obtain:
$$W_1=P(\sigma ,\sigma _0;\tau _0),W_2=Q(\sigma ;\tau ,\tau _0),W_3=P^1(\sigma ,\sigma _0;\tau ),W_4=Q^1(\sigma _0;\tau ,\tau _0).$$
(33)
Using these expressions, the supersymmetric Wilson loop can be expressed as
$$W_s[C]=Q(\sigma _0;\tau ,\tau _0)^1P(\sigma ,\sigma _0;\tau )^1Q(\sigma ;\tau ,\tau _0)P(\sigma ,\sigma _0;\tau ).$$
(34)
For definiteness, in the rest of the paper we will confine ourselves to the case in which the 2-surface, $`\mathrm{\Sigma }`$, can be covered by a single coordinate patch. If $`\mathrm{\Sigma }`$ requires more than one patch to be covered, then using partition of unity and the product integral composition rule, Eq. (15), it is straightforward to extend our upcoming reasonings.
## 4 Super Non-Abelian Stokes Theorem
As an application of the supersymmetric product integral formalism, we prove the supersymmetric version of the non-Abelian Stokes theorem . The proof makes essential use of the generalized theorems listed in the previous paragraphs, and is the supersymmetric version of one of the proofs given for the non-supersymmetric case in reference . The other proof give in this reference can also be extended to the supersymmetric case and will be given in a subsequent work .
We start with the form of $`W_s[C]`$ given in Eq. (34) and take its derivatives with respect to the parameter $`\tau `$:
$`{\displaystyle \frac{W_s[C]}{\tau }}`$ $`=`$ $`_\tau Q^1(\sigma _0;\tau ,\tau _0)P^1(\sigma ,\sigma _0;\tau )Q(\sigma ;\tau ,\tau _0)P(\sigma ,\sigma _0;\tau _0)+`$ (35)
$`+Q^1(\sigma _0;\tau ,\tau _0)_\tau P^1(\sigma ,\sigma _0;\tau )Q(\sigma ;\tau ,\tau _0)P(\sigma ,\sigma _0;\tau _0)+`$
$`+Q^1(\sigma _0;\tau ,\tau _0)P^1(\sigma ,\sigma _0;\tau )_\tau Q(\sigma ;\tau ,\tau _0)P(\sigma ,\sigma _0;\tau _0).`$
Here, we have made use of the fact that $`P(\sigma ,\sigma _0;\tau _0)`$ is independent of $`\tau `$. As a preparation for using Eq. (19), we start with Eq. (17) for $`W_s[C]`$, and make use of Eq. (16) to get
$`L_\tau W_s[C]={\displaystyle \frac{W_s[C]}{\tau }}W_s[C]^1=`$ $`T^1(\sigma ;\tau )[\mathrm{\Gamma }_0(\sigma ;\tau )P(\sigma ,\sigma _0;\tau )\mathrm{\Gamma }_0(\sigma _0;\tau )P^1(\sigma ,\sigma _0;\tau )`$ (36)
$`_\tau P(\sigma ,\sigma _0;\tau )P^1(\sigma ,\sigma _0;\tau )]T(\sigma ;\tau ),`$
In this expression, $`T(\sigma ;\tau )=P(\sigma ,\sigma _0;\tau )Q(\sigma _0;\tau ,\tau _0)`$. Next, by means of differentiation with respect to a parameter given by Eq. (22), we evaluate the derivative of the super product integral $`P(\sigma ,\sigma _0;\tau )`$ with respect to the parameter $`\tau `$:
$$_\tau P(\sigma ,\sigma _0;\tau )=_{\sigma _0}^\sigma 𝑑\sigma ^{}P(\sigma ,\sigma ^{};\tau )_\tau \mathrm{\Gamma }_1(\sigma ^{};\tau )P(\sigma ^{},\sigma _0;\tau ).$$
(37)
Then, after some simple manipulations using the defining equations for the various terms in Eq. (36), we get:
$$T^1(\sigma ;\tau )_\tau P(\tau )P^1(\tau )T(\sigma ;\tau )=_{\sigma _0}^\sigma 𝑑\sigma ^{}T^1(\sigma ^{};\tau )_\tau \mathrm{\Gamma }_1(\sigma ^{};\tau )T(\sigma ^{};\tau ).$$
(38)
Using Eq. (16) and the fact that $`P(\sigma _0,\sigma _0;\tau )=1`$, we can rewrite the rest of Eq. (36) also as an integral:
$`T^1(\sigma ;\tau )[\mathrm{\Gamma }_0(\sigma ;\tau )P(\sigma ,\sigma _0;\tau )\mathrm{\Gamma }_0(\sigma _0;\tau )P^1(\sigma ,\sigma _0;\tau )]T(\sigma ;\tau )=`$
$`={\displaystyle _{\sigma _0}^\sigma }𝑑\sigma ^{}P^1(\sigma ^{},\sigma _0;\tau )(_\tau \mathrm{\Gamma }_0(\sigma ^{},\tau )+[\mathrm{\Gamma }_0(\sigma ^{},\tau ),\mathrm{\Gamma }_1(\sigma ^{},\tau )])P(\sigma ^{},\sigma _0;\tau ).`$ (39)
Combining Eqs. (36), (38), and (4), we obtain:
$$L_\tau W_s[C]=_{\sigma _0}^\sigma 𝑑\sigma ^{}T^1(\sigma ^{},\tau )F_{01}(\sigma ^{},\tau )T(\sigma ^{},\tau ).$$
(40)
Here $`F_{01}`$ is the field strength component as defined in Eq. (28), but based on the discussion in that paragraph, we know that it also equals the pull-back of the supersymmetric field strength to the surface. Using Eq. (19), we are immediately led to the supersymmetric version of the non-Abelian Stokes theorem:
$$W_s[C]=\underset{\tau _0}{\overset{\tau }{}}e^{_{\sigma _0}^\sigma T^1(\sigma ^{};\tau ^{})F_{01}(\sigma ^{};\tau ^{})T(\sigma ^{};\tau ^{})𝑑\sigma ^{}𝑑\tau ^{}}.$$
(41)
Recalling the antisymmetry of the components of the field strength, we can rewrite this expression in a more familiar reparameterization invariant form
$$W_s[C]=𝒫_\tau e^{\scriptscriptstyle \mathrm{\Gamma }}=\underset{\tau _0}{\overset{\tau }{}}e^{\frac{1}{2}_\mathrm{\Sigma }𝑑\sigma ^{ab}T^1(\sigma ;\tau )F_{ab}(\sigma ;\tau )T(\sigma ;\tau )},$$
(42)
where $`d\sigma ^{ab}`$ is the area element of the 2-surface. Despite appearances, it must be remembered that $`\sigma `$ and $`\tau `$ play very different roles in this expression.
The above result also applies to the special case in which the gauge group is Abelian. In that case, however, since the corresponding matrices commute, the machinery of the super product integrals is not needed, and one can establish the super Stokes theorem directly .
## 5 Gauge Covariance of the Super Loop Operator
To demonstrate the gauge covariance of the supersymmetric Wilson loop operator and its 2-surface representation, we must show how the supersymmetric Wilson line transforms under gauge transformations. For this, we need to know, in turn, how the pull-back of the connection $`\mathrm{\Gamma }_A(z)`$ transforms. The transformation properties of the connection itself follows from that of the vector superfield : $`e^V^{}=e^{i\mathrm{\Lambda }^{}}e^Ve^\mathrm{\Lambda }.`$ More specifically, we have
$$\mathrm{\Gamma }^{}(z)=g(z)\mathrm{\Gamma }(z)g(z)^1g(z)dg(z)^1,$$
(43)
where, $`g(z)=e^{i\mathrm{\Lambda }(z)}`$ and $`d=e^A(z)D_A`$. As we have seen, the pull-back of this quantity on the line is given by $`d=ds_s`$. Thus, we get for the transformation of the supersymmetric connection on the line:
$$\mathrm{\Gamma }^{}(s)=g(s)\mathrm{\Gamma }(s)g^1(s)g(s)_sg^1(s).$$
(44)
This is formally identical to that for the plain Yang-Mills theory . As a result, under a gauge transformation we obtain:
$$\underset{a}{\overset{b}{}}e^{ds\mathrm{\Gamma }(s)}\underset{a}{\overset{b}{}}e^{[g(s)\mathrm{\Gamma }(s)g^1(s)g(s)_sg^1(s)]ds}.$$
(45)
By Eq. (17), we have $`g(s)_sg^1(s)=L_sg(s).`$ Thus, for the gauge transformed super Wilson line we have
$$\underset{a}{\overset{b}{}}e^{[g(s)\mathrm{\Gamma }(s)g^1(s)+L_sg(s)]ds}.$$
(46)
Moreover, using Eq. (21) and recalling from Eq. (19) that $`_a^be^{L_sg(s)ds}=g(b)g^1(a)`$, the gauge transformed expression takes the form
$$g(b)g^1(a)\underset{a}{\overset{b}{}}e^{g(a)\mathrm{\Gamma }(s)g^1(a)}.$$
(47)
Finally, using the same argument as in reference , the constant terms in the exponents can be factored out from the super product integral. Thus, we get for the gauge transformed super Wilson line
$$\underset{a}{\overset{b}{}}e^{ds\mathrm{\Gamma }(s)}g(b)\left(\underset{a}{\overset{b}{}}e^{ds\mathrm{\Gamma }(s)}\right)g^1(a).$$
(48)
We can use this result to determine the gauge transforms of operators which are products of simple super Wilson lines. Consider, e.g., the operator $`T(\sigma ;\tau )`$ which is the product of two super Wilson lines. Applying Eq. (48) to each factor, we obtain:
$$T(\sigma ;\tau )g(\sigma ;\tau )T(\sigma ;\tau )g^1(\sigma _0;\tau _0).$$
(49)
From this, we can easily obtain the transformation properties of the super Wilson loop operator which is also a composite of super Wilson lines. The transformation has the same form as Eq. (48) with $`a=b`$.
Finally, let us consider how the surface integral representation of super Wilson loop operator given by Eq. (41) transforms under gauge transformation. From knowing how each factor in the exponent transforms, it follows that
$$W_s[C]\underset{\tau _0}{\overset{\tau }{}}e^{g(\sigma _0;\tau _0)\left(_{\sigma _0}^\sigma T^1(\sigma ^{};\tau ^{})F_{01}(\sigma ^{};\tau ^{})T(\sigma ^{};\tau ^{})𝑑t^{}\right)g^1(\sigma _0;\tau _0)}.$$
(50)
Just as for super Wilson line, the constant terms in the exponent factorize, so that under gauge transformations the surface integral representation of the super Wilson loop transforms covariantly:
$$W_s[C]g(\sigma _0;\tau _0)\underset{\tau _0}{\overset{\tau }{}}e^{_{\sigma _0}^\sigma T^1(\sigma ^{};\tau ^{})F_{01}(\sigma ^{};\tau ^{})T(\sigma ^{};\tau ^{})𝑑t^{}}g^1(\sigma _0;\tau _0).$$
(51)
## 6 Concluding Remarks
In this work, we have presented a supersymmetric generalization of ordinary product integral formalism. Given that Wilson lines and Wilson loops can be expressed in terms of ordinary product integrals, we have constructed the supersymmetric extensions of these notions for supersymmetric gauge theories in terms of supersymmetric product integrals. These constructions are natural in the sense that the supersymmetric representations given in this paper reduce to the ordinary product integral representations of standard Wilson lines and Wilson loops.
It is hoped that this formalism provides a reliable non-perturbative means of extracting information from supersymmetric gauge theories. In this respect, we note that the construction of the supersymmetric Wilson lines and Wilson loops as well as the proof of the super non-Abelian Stokes theorem given in the previous sections are independent of any specific physical applications. To apply these concepts to supersymmetric gauge theories, it is necessary to clarify the physical content of the operators such as the connection and the field strength which appear in the relevant expressions . It is well known that in supersymmetric gauge theories the superfield strength $`F_{AB}`$ contains more degrees of freedom than is required by supersymmetry and gauge invariance . As a result, it is necessary to impose constraints on the components of the field strength to eliminate the unphysical degrees of freedom. This means that in the expressions for supersymmetric Wilson lines and loops, $`\mathrm{\Gamma }`$ and $`F`$ must be expressed in terms of unconstrained superfields, just as in the abelian case . Such a description in terms of unconstrained superfields already exist in the literature and can be adapted to specific applications.
This work was supported in part by the Department of Energy under the contract number DOE-FGO2-84ER40153. We are grateful to M. Awada for valuable input at the initial stages of this work. We would also like to thank R. Grimm, A. Kiss and R.L. Mkrtchian for helpful communications. |
warning/0002/hep-th0002075.html | ar5iv | text | # 1 The UV/IR Connection in Non-commutative Field Theory
## 1 The UV/IR Connection in Non-commutative Field Theory
Field theories formulated on non-commutative spaces are interesting in both their own right as well as for their applications to string and matrix theories . These theories are characterized by a non-commutativity parameter $`\theta `$ with dimensions of length squared. Classically and in the tree–level approximation the behavior of the theory for momenta much less than $`\theta ^{1/2}`$ is the same as for the corresponding commutative theory. However this is not necessarily the case for the quantum theory. The non-commutativity can lead to unfamiliar effects of the ultraviolet modes on the infrared behavior which have no analog in conventional quantum field theory .
The origin of the strange mixing of IR and UV effects in non-commutative field theory can be understood in a simple way . The field quanta in such a theory can be thought of as pairs of opposite charges moving in a strong magnetic field . The spatial locations of the two charges are defined by a center of mass position $`x_{cm}^i`$ and a relative coordinate $`\mathrm{\Delta }^m`$. The relative coordinate is related to the spatial momentum $`p`$ by
$$\mathrm{\Delta }^i=\theta ^{ij}p_j$$
(1.1)
where $`\theta `$ is an antisymmetric matrix with components in the spatial directions. In this paper we will consider the case of 3 dimensional space. Without loss of generality $`\theta `$ can be taken to lie in the $`(1,2)`$ plane
$`\theta ^{1,2}`$ $`=`$ $`\theta ^{2,1}\theta `$ (1.2)
$`\theta ^{1,3}`$ $`=`$ $`\theta ^{3,2}=0.`$ (1.3)
The momentum in the $`(1,2)`$ plane will be called $`P`$. We will also use the notation $`\stackrel{~}{P}_i=\theta _{ij}P^j.`$ Thus a particle moving with momentum $`P`$ along the $`X^1`$ axis has a spatial extension of size $`|\theta P|`$ in the $`X^2`$ direction. The growth of the size of a particle with its momentum has interesting consequences. For example, when a quantum of momentum $`P`$ scatters off a target at rest, the scattering amplitude will spread in impact parameter space over a distance $`|\theta P|`$.
There are also important and somewhat bizarre consequences for Feynman loop integrations. Roughly speaking, when a particle of momentum $`P`$ circulates in a loop it can induce an effect at distance $`|\theta P|`$. The high momentum end of the integrals can give rise to power law long range forces which are absent entirely in the classical theory. We will call such effects “anomalies” although we emphasize that they do not signal an inconsistency in the theory but rather a violation of naive expectations.
In the case of planar diagrams the effect of non-commutativity is simple. Every diagram gets multiplied by a phase factor that depends only on the momenta of the external lines. Thus the Feynman integrals are exactly as in the commutative theory. In the non-planar case the situation is more interesting. The phase factors now involve products of external momenta $`p`$ and internal momenta $`l`$ in the form
$$e^{ip\theta l}=e^{ip\stackrel{~}{l}}.$$
(1.4)
If the diagram in question is divergent in the commutative theory, the effect of the oscillating phases is typically to regulate the diagram and render it finite. But as $`P0`$ the phases become ineffective and the diagram diverges at $`P=0`$. This is the mechanism described in detail in . We will begin by reviewing a simple example from non-commutative $`\varphi ^4`$ theory.
The diagram in question, fig(1), is the lowest order mass renormalization correction to the propagator. We are interested in the non-planar contribution which in the commutative theory has the form
$$d^4l\frac{1}{l^2}.$$
(1.5)
The diagram is quadratically divergent and is renormalized by a mass counter term.
In the non-commutative case the integrand has an additional factor $`\mathrm{exp}(i\stackrel{~}{p}l)`$ where $`p`$ and $`l`$ are the external <sup>1</sup><sup>1</sup>1In the rest of the paper, we denote by small p the 4–momentum vector with $`p_{1,2}P_{1,2}`$ in the non–commutativity plane. and loop momenta. The integral has the form
$$d^4l\frac{1}{l^2}e^{i\stackrel{~}{p}l}\frac{1}{\theta ^2P^2}.$$
(1.6)
As emphasized in there are some very striking features of this result. The first is that the pole at $`P=0`$ arises from the high momentum region of integration. Although we evaluated it for the massless theory, the pole itself is independent of mass. Furthermore this contribution to the self energy has a huge effect on the propagation of long wavelength particles. The on-shell condition or dispersion relation becomes
$$p_0^2=p_3^2+P^2+c\frac{1}{\theta ^2P^2}$$
(1.7)
where $`c`$ is proportional to the coupling constant. Thus, as discussed in , the behavior of the non-commutative theory below the non-commutativity scale seems to be nothing like the commutative theory. In this case the low momentum end of the spectrum is completely removed from the low energy theory.
Commutative gauge theories are better behaved in the UV than commutative scalar theories. The worst divergences in pure Yang mills theory or Yang Mills theory with fermions are logarithmic. This naively suggests that in their non-commutative versions the worst anomalous effects will be logarithmic in $`P`$. As an example consider the vacuum polarization correction to the gauge boson propagator. The divergences have the form
$$\mathrm{\Pi }g^2p^2\mathrm{log}\kappa $$
(1.8)
where $`p^2`$ is the squared four-momentum of the gauge boson. Note in particular that the mass correction vanishes since $`\mathrm{\Pi }`$ vanishes at $`p^2=0`$. This situation suggests that in the noncommutative theory the worst anomalous effect in the propagator has the form
$$\mathrm{\Pi }g^2p^2\mathrm{log}\stackrel{~}{P}^2$$
(1.9)
If this were so, the dispersion relation of a low energy gauge boson would be unaffected by the non-commutativity. As we will see in the next section this is generally incorrect.
## 2 U(1) Non–Commutative Yang–Mills
In this section we analyze $`U(1)`$ Yang–Mills theory on a non-commutative space. The classical action is given by
$$S=\frac{1}{4}d^4xF^2,$$
(2.1)
with the field strength F given by
$$F_{\mu \nu }=_\mu A_\nu _\nu A_\mu ig[A_\mu ,A_\nu ]$$
(2.2)
and
$$[A_\mu ,A_\nu ]=A_\mu A_\nu A_\nu A_\mu .$$
(2.3)
The $``$–product between two functions $`\varphi _1(x)`$ and $`\varphi _2(x)`$ is defined by
$$\varphi _1\varphi _2(x)=e^{\frac{i}{2}\theta ^{\mu \nu }_\mu ^y_\nu ^z}\varphi _1(y)\varphi _2(z)|_{y=z=x}.$$
(2.4)
The theory is invariant under non–commutative gauge transformations
$$\delta _\lambda A_\mu =_\mu \lambda gi(A_\mu \lambda \lambda A_\mu ).$$
(2.5)
We may add matter fields in the theory as well. The scalar and fermionic parts of the action that involve interactions of the matter fields with the gauge field are given by
$$S_{matter}=d^4xi\overline{\psi }\gamma ^\mu D_\mu \psi +\frac{1}{2}(D_\mu \mathrm{\Phi })^2.$$
(2.6)
The covariant derivative acts on the fields as follows
$$D_\mu X=_\mu Xig[A_\mu ,X].$$
(2.7)
The commutator is defined through the $``$–product as before. The matter fields are covariant under the family of gauge transformations given in Eq(2.5).
The Feynman rules for the theory have been worked out in references . The vertices look similar to those of a commutative non–abelian gauge theory with all matter fields in the adjoint of the gauge group. The structure constants are replaced by sines of external momenta as shown in detail in the Appendix. The Feynman rules for the ghosts are also included. The vertices vanish when $`\theta `$ is taken to be zero as expected. Our computations are done in the Feynman gauge. In this gauge the gauge field propagator is given by
$$i\frac{g_{\mu \nu }}{p^2}.$$
(2.8)
We refer the reader in for a derivation of the perturbative Feynman rules and gauge fixing of the theory.
## 3 Photon self–energy correction
In non-commutative gauge theory the most serious anomalous effects, namely those which exhibit inverse powers of $`\stackrel{~}{P}`$, are associated with the 2 and 3 point functions. We will begin with the computation of the 2-point photon self energy diagrams. We will consider the contributions from loops involving fermions, scalars and gauge bosons (including ghosts). It should be noted that the individual contributions are gauge invariant. Similar calculations have been done by Hayakawa . The relevant diagrams are shown in Figures 2 and 3. We will illustrate the procedure with the fermion loop (line (1) in Fig. 3).
Using the Feynman rules from the Appendix we find
$$i\mathrm{\Pi }_f^{\mu \nu }(p)=4g^2N_f\frac{d^4l}{(2\pi )^4}\frac{\mathrm{Tr}\left[\gamma ^\mu (/l/p)\gamma ^\nu /l\right]}{(lp)^2l^2}\mathrm{sin}^2\left(\frac{1}{2}\stackrel{~}{p}l\right),$$
(3.1)
where $`N_f`$ is the number of Majorana fermions, each of which counts as two fermion species.
We are interested in the contribution to the integral coming from very high loop momentum. We therefore drop the sub-leading dependence in the integrand and replace eq.(3.1) by
$$i\mathrm{\Pi }_f^{\mu \nu }(p)=4g^2N_f\frac{d^4l}{(2\pi )^4}\frac{\mathrm{Tr}\left[\gamma ^\mu /l\gamma ^\nu /l\right]}{(l)^4}\mathrm{sin}^2\left(\frac{1}{2}\stackrel{~}{p}l\right),$$
(3.2)
Using
$$\mathrm{sin}^2\left(\frac{1}{2}\stackrel{~}{p}l\right)=\frac{1}{2}[1\mathrm{cos}\left(\stackrel{~}{p}l\right)],$$
(3.3)
we can isolate the planar and non-planar contributions. The non-planar contribution is obtained by dropping the first term and keeping only the $`\mathrm{cos}`$ term. Working out the trace, we get
$$i\mathrm{\Pi }_f^{\mu \nu }(p)=4g^2N_f\frac{d^4l}{(2\pi )^4}\frac{\left(2l^\mu l^\nu g^{\mu \nu }l^2\right)}{l^4}e^{i\stackrel{~}{p}l}.$$
(3.4)
Since $`d^4l/(2\pi )^4\frac{1}{l^4}e^{i\stackrel{~}{p}l}=i\alpha (\mathrm{log}\mathrm{\Lambda }\mathrm{log}|\stackrel{~}{p}|)`$ <sup>2</sup><sup>2</sup>2In Euclidean space the integral is of course real, $`d^4l/(2\pi )^4\frac{1}{l^4}e^{i\stackrel{~}{p}l}=\alpha (\mathrm{log}\mathrm{\Lambda }\mathrm{log}|\stackrel{~}{p}|),\alpha >0.`$, where $`\mathrm{\Lambda }`$ denotes a short–momentum cut–off and $`\alpha `$ is a positive real constant, we can rewrite (3.4) as
$$i\mathrm{\Pi }_f^{\mu \nu }(p)=4ig^2N_f\alpha \left(2^\mu ^\nu g^{\mu \nu }^2\right)\mathrm{log}|\stackrel{~}{p}|.$$
(3.5)
The integral (3.4) is finite, and the cut–off dependence vanishes after differentiating. We note that the term $`g^{\mu \nu }/\stackrel{~}{p}^2`$ cancels in this expression leaving us with
$$i\mathrm{\Pi }_f^{\mu \nu }(p)=8ig^2(2N_f)\alpha \frac{\stackrel{~}{p}^\mu \stackrel{~}{p}^\nu }{\stackrel{~}{p}^4}.$$
(3.6)
The answer is somewhat surprising. If the factor $`\mathrm{sin}^2\left(\frac{1}{2}\stackrel{~}{p}l\right)`$ were not present in eq(3.1) the expression would be the conventional self energy diagram of the commutative theory. Gauge invariance would be invoked to say that any quadratic divergence is absent. Alternatively the diagram can be Pauli Villars regulated eliminating the quadratic divergence. The integral with the trigonometric factor is finite and well defined. However it quadratically diverges as $`\stackrel{~}{p}0`$. Thus we see that there is an anomalous effect of order $`\theta ^2`$ arising out of a diagram which in the commutative theory is quadratically divergent by power counting but for which the divergence vanishes as a consequence of symmetry. As noted in in the context of scalar theories this type of behavior is proportional to inverse powers of $`\theta `$. Evidently the limit in which $`\theta 0`$ does not smoothly tend to the commutative theory.
The physical interpretation of terms like eq(3.6) is very interesting. For small non-commutative momentum, the one-loop inverse propagator is given by
$$\mathrm{\Gamma }_{\mu \nu }=i\left[(p_0^2p_3^2P^2)g_{\mu \nu }g^2c\frac{\stackrel{~}{p}_\mu \stackrel{~}{p}_\nu }{\stackrel{~}{p}^4}\right],$$
(3.7)
where $`P`$ represents the projection of the spatial momentum on the $`(1,2)`$ plane. From this matrix we can read the dispersion relation for the two physical, transversely polarized photons. Suppose $`P`$ is along the 2–direction so that $`\stackrel{~}{P}`$ is in the 1-direction. Then the photon polarized in the direction perpendicular to $`\stackrel{~}{P}`$ satisfies the same dispersion relation as a photon would in the commutative theory
$$p_0^2=p_3^2+P^2.$$
(3.8)
However, the photon polarized along the 1-direction, parallel to $`\stackrel{~}{P}`$, satisfies a different dispersion relation given by
$$p_0^2=p_3^2+P^2+cg^2\frac{1}{\theta ^2P^2}.$$
(3.9)
This splitting of the polarization states of the gauge boson is perfectly consistent with gauge invariance. Indeed the vacuum polarization tensor in eq.(3.6) is purely transverse which follows from the identity $`p\stackrel{~}{p}=0`$. This effect would not be possible without the breaking of Lorentz invariance caused by $`\theta `$.
We remark at this point that since the contributions from the scalar loops and from the gauge sector are gauge–invariant by themselves, they individually give combinations of the form (3.4). The computations involving gauge bosons (Fig. 2) and scalars (line (2) on Fig. 3) in the loop are very similar, and without giving explicit details, we add all three sectors together obtaining
$$i\mathrm{\Pi }^{\mu \nu }(p)=8ig^2(N_s+22N_f)\alpha \frac{\stackrel{~}{p}^\mu \stackrel{~}{p}^\nu }{\stackrel{~}{p}^4}.$$
(3.10)
This shows that in the supersymmetric theories with an equal number of bosons and fermions, quadratic divergences in the photon self–energy do not appear at one loop.
As we have mentioned before, the coefficient $`c=8(N_s+22N_f)`$ of the quadratic anomalous term is proportional to the number of bosons minus the number of fermions in the theory. In particular it is zero in any supersymmetric model. This had to be the case since the splitting of the two photon states would be inconsistent with supersymmetry. In particular, the fermions in the theory remain massless with dispersion relation
$$/p=0.$$
(3.11)
Note also that in the $`𝒩=4`$ case the SO(6) R–symmetry of the theory insures that the scalars do not split. So it would be impossible to achieve bose–fermi degeneracy with multiplets of $`𝒩=4`$ supersymmetry had $`c`$ been non–zero.
We also note that the coefficient $`c`$ will vanish in softly broken supersymmetric theories as well. In particular, the effect we find arises from high momenta circulating the loop. Therefore, it is independent of the mass of the loop particle. As long as the number of fermions equals the number of bosons in the theory c will cancel. Simple dimensional analysis shows that the only effect of adding a mass term in the Lagrangian is to modify the logarithmic singularities as $`\stackrel{~}{p}0`$.
Finally, we note that $`c`$ becomes negative if the number of fermions is bigger than the number of bosons. This is also the case in the $`\varphi ^4`$ theory when we include fermions coupled to the scalar field through Yukawa couplings. In both cases, the theory becomes unstable at low energies.
## 4 Vertex Corrections
Next we analyze 1–loop corrections to the three–photon vertex. In Figures 4 and 5, we have drawn all 1PI Feynman diagrams that contribute corrections to the vertex up to 1–loop order in perturbation theory. As before, we shall see that there are unfamiliar long distance effects arising from the region of integration of high loop momenta.
Let us consider in detail the scalar graph involving cubic vertices only (Fig. 5, 1st graph on line (1)). Applying the Feynman rules, we find
$$8ig^3\frac{d^4l}{(2\pi )^4}\frac{(2l+p_1)^{\mu _1}(2lp_2)^{\mu _2}(2l+p_1p_2)^{\mu _3}}{l^2(lp_2)^2(l+p_1)^2}\mathrm{sin}\left(\frac{\stackrel{~}{p}_1l}{2}\right)\mathrm{sin}\left(\frac{\stackrel{~}{p}_2l}{2}\right)\mathrm{sin}\left(\frac{\stackrel{~}{p}_3(l+p_1)}{2}\right).$$
(4.1)
We must also impose overall momentum conservation so that
$$p_1+p_2+p_3=0.$$
(4.2)
We have amputated the external propagators for compactness. Eq(4.1) has to be multiplied by the number of scalars in the theory. Ignoring the phases for a moment, we note that the integrand is linear in $`l`$ at high loop momentum
$$\frac{d^4l}{(2\pi )^4}\frac{l^{\mu _1}l^{\mu _2}l^{\mu _3}}{l^6}.$$
(4.3)
In the commutative non–abelian theory, however, no linear divergence arises because the integral is zero by symmetry. Thus it is at most logarithmic in the cutoff. In the non–commutative case the oscillating phases spoil the rotational symmetry. They make the integral finite but they induce low momentum poles of the form
$$\frac{\stackrel{~}{p}^{\mu _1}\stackrel{~}{p}^{\mu _2}\stackrel{~}{p}^{\mu _3}}{\stackrel{~}{p}^4}.$$
(4.4)
Unlike the previous effect which was of order $`\theta ^2`$, this is a $`\theta ^1`$ effect.
To compute the precise coefficient in front of such anomalous terms, it is enough to consider high momentum running in the loop ignoring external momenta in the denominators. Then the integral becomes
$$8ig^3\mathrm{cos}(\frac{\stackrel{~}{p}_3p_1}{2})\frac{d^4l}{(2\pi )^4}\frac{8l^{\mu _1}l^{\mu _2}l^{\mu _3}}{l^6}\mathrm{sin}\left(\frac{\stackrel{~}{p}_1l}{2}\right)\mathrm{sin}\left(\frac{\stackrel{~}{p}_2l}{2}\right)\mathrm{sin}\left(\frac{\stackrel{~}{p}_3l}{2}\right).$$
(4.5)
To carry out this integral, it is useful to express the product of sines as a sum of exponentials. Using Eq(4.2), we can write the product of sines as
$$\frac{1}{4}\left[\mathrm{sin}\left(\stackrel{~}{p}_1l\right)+\mathrm{sin}\left(\stackrel{~}{p}_2l\right)+\mathrm{sin}\left(\stackrel{~}{p}_3l\right)\right].$$
(4.6)
We are left with a sum of three simpler integrals, and, in addition, we can replace each sine by an exponential in the integral. In this form, it is easy to see that the contributions arise solely from the six non–planar graphs in the double line notation.
To these, we must add the contributions from the scalar graphs that involve a quartic vertex (Second graph and its two permutations on line (1) in Fig. 5). The graphs produce a different tensor structure but the phase factors are the same. There are three such graphs from permuting the external particles among the external lines. Each graph has to be multiplied by a symmetry factor of $`1/2`$. Adding all four graphs together yields the following integrals
$$4g^3N_s\mathrm{cos}(\frac{\stackrel{~}{p}_3p_1}{2})\frac{d^4l}{(2\pi )^4}\frac{1}{l^6}\left[4l^{\mu _1}l^{\mu _2}l^{\mu _3}l^2(l^{\mu _1}g^{\mu _2\mu _3}+\mathrm{perms})\right](e^{i\stackrel{~}{p}_1l}+e^{i\stackrel{~}{p}_2l}+e^{i\stackrel{~}{p}_3l}).$$
(4.7)
The pure gauge sector graphs and the fermionic graphs can be computed in the same way. One has to remember to include a combinatorics factor of 2 for the ghosts and fermion graphs arising from the two different cyclic orderings of the external particles on the loop. The gauge boson quartic graph has a symmetry factor of a $`1/2`$. Including these has the effect of changing the coefficient in front of the integrals to
$$N_s+22N_f,$$
(4.8)
where $`N_f`$ is the number of Majorana fermions in the theory. We see that in any supersymmetric theory the linear poles are absent.
We will now explicitly find the structure of the linear poles. To compute the integral in (4.7),
$$I^{\mu _1\mu _2\mu _3}=\frac{d^4l}{(2\pi )^4}\frac{\left(4l^{\mu _1}l^{\mu _2}l^{\mu _3}l^2(l^{\mu _1}g^{\mu _2\mu _3}+l^{\mu _2}g^{\mu _1\mu _3}+l^{\mu _3}g^{\mu _1\mu _2})\right)}{l^6}e^{i\stackrel{~}{p}l},$$
(4.9)
we again use $`d^4l/(2\pi )^4\frac{1}{l^4}e^{i\stackrel{~}{p}l}=i\alpha (\mathrm{log}\mathrm{\Lambda }\mathrm{log}|\stackrel{~}{p}|)`$. Notice that (4.9) is convergent, and the $`\mathrm{\Lambda }`$–dependent part can thus be dropped. We can write
$$I^{\mu _1\mu _2\mu _3}=i\left(4^{\mu _1}^{\mu _2}^{\mu _3}\mathrm{}(g^{\mu _2\mu _3}^{\mu _1}+g^{\mu _1\mu _3}^{\mu _2}+g^{\mu _1\mu _2}^{\mu _3})\right)J(\stackrel{~}{p}),$$
(4.10)
where $`J(\stackrel{~}{p})=d^4l/(2\pi )^4\frac{1}{l^6}e^{i\stackrel{~}{p}l}=iJ_E(\stackrel{~}{p})`$. The integral over Euclidean space, $`J_E(\stackrel{~}{p})`$ is obtained from $`J(\stackrel{~}{p})`$ after performing a Wick rotation. It satisfies
$$\mathrm{}J_E(\stackrel{~}{p})=\frac{1}{\stackrel{~}{p}^3}\frac{d}{d\stackrel{~}{p}}\stackrel{~}{p}^3\frac{d}{d\stackrel{~}{p}}J_E(\stackrel{~}{p})=\alpha \mathrm{log}|\stackrel{~}{p}|.$$
(4.11)
Solving for $`J`$ and substituting
$$J(\stackrel{~}{p})=iJ_E(\stackrel{~}{p})=\frac{i\alpha }{32}\stackrel{~}{p}^2(4\mathrm{log}\stackrel{~}{p}3)$$
(4.12)
in (4.10) we find
$$I^{\mu _1\mu _2\mu _3}=2\alpha \frac{\stackrel{~}{p}^{\mu _1}\stackrel{~}{p}^{\mu _2}\stackrel{~}{p}^{\mu _3}}{\stackrel{~}{p}^4}.$$
(4.13)
Note that the second possible tensor structure, $`(\stackrel{~}{p}^{\mu _1}g^{\mu _2\mu _3}+\stackrel{~}{p}^{\mu _2}g^{\mu _1\mu _3}+\stackrel{~}{p}^{\mu _3}g^{\mu _1\mu _2})/P^2,`$ cancelled completely, just like the photon mass term in the photon self–energy correction (3.6). This fact will turn out to be essential in showing the gauge invariance of the S–matrix in the next chapter.
We can now summarize the computations in this section by giving the linearly divergent terms in the correction to the 3–point photon vertex:
$$\mathrm{\Gamma }^{\mu _1\mu _2\mu _3}(p_1,p_2,p_3)=8\alpha g^3(N_s+22N_f)\mathrm{cos}\left(\frac{1}{2}\stackrel{~}{p}_3p_1\right)\underset{i=1}{\overset{3}{}}\frac{\stackrel{~}{p}_i^{\mu _1}\stackrel{~}{p}_i^{\mu _2}\stackrel{~}{p}_i^{\mu _3}}{\stackrel{~}{p}_i^4}+\mathrm{},$$
(4.14)
where the terms denoted by $`\mathrm{}`$ are at most logarithmic in $`\stackrel{~}{p}_1,\stackrel{~}{p}_2,\stackrel{~}{p}_3`$.
The anomalous effects found in this paper and in are highly nonlocal but in a particular way. The matrix elements in eqs(3.10) and (4.14) depend depend only on the components of momentum in the $`x^1,x^2`$ plane and are independent of $`x^3`$ and $`x^4`$. Thus while nonlocal in the noncommutative directions they are completely local in the commutative directions.
In supersymmetric theories the nonlocal $`\theta ^2`$ and $`\theta ^1`$ effects are absent, at least at one loop. However there are still logarithmic dependences on $`P`$ which are proportional to the one-loop coefficient of the $`\beta `$ function<sup>3</sup><sup>3</sup>3The $`\beta `$ function is the one controlling the running of the t’Hooft coupling in the large $`N`$ limit. in the corresponding commutative non-abelian gauge theory. Once again the corresponding effects are nonlocal only in the noncommutative directions. Theories such as $`N=4`$ super Yang Mills theory with vanishing $`\beta `$ function seem to be free of the non-analytic dependence on $`\theta `$.
## 5 Gauge Invariant S–Matrix?
In this section we show that the anomalous terms in the 2–point and 3–point functions we computed are consistent with gauge invariance. To check gauge invariance, we study a case involving scattering of two gauge bosons into two fermions. We consider 1–loop diagrams only and choose the kinematic variables so that the anomalous terms we found dominate.
All particles must be put on shell using the appropriate, corrected dispersion relations up to 1–loop order in perturbation theory. We choose one of the gauge bosons to be transversely polarized. We denote its polarization by $`ϵ`$ and its momentum by p; then
$$ϵ_\mu p^\mu =0.$$
(5.1)
In order to test the gauge invariance we set the polarization of the other gauge boson, $`\theta `$, equal to its momentum $`q`$. The momenta of the fermions are denoted by $`(l_1,l_2)`$. Gauge invariance requires this scattering amplitude to be zero. The tree level diagrams contributing to the process are shown in Fig. 6.
Now we choose a specific kinematic limit. We choose $`l_1+l_2=l`$ so that $`\stackrel{~}{l}`$ is small. Furthermore, we let $`l^2`$ be small so that only the s–channel diagrams are important. Two types of one loop diagrams are important (Fig. 7), the diagrams involving corrections to the intermediate gauge boson’s propagator and the diagrams involving corrections to the three gauge boson vertex. Corrections to the two fermion–gauge boson vertex are at most logarithmic in $`\stackrel{~}{l}`$ and, therefore, sub-leading. We can think of the intermediate gauge boson being coupled to some current $`j_\mu `$. The exact form of $`j_\mu `$ is not important for our purposes.
First consider 1–loop diagrams involving the gauge boson self–energy. In the limit when $`\stackrel{~}{l}`$ is small the contribution to the amplitude becomes
$$i_1=16i\alpha g^4\mathrm{sin}(\frac{q\stackrel{~}{l}}{2})ϵ_\nu q_\mu \left[g^{\mu \nu }(qp)^\rho +g^{\nu \rho }(pl)^\mu +g^{\rho \mu }(lq)^\nu \right]\frac{\stackrel{~}{l}_\rho (\stackrel{~}{l}j)}{l^4\stackrel{~}{l}^4}.$$
(5.2)
The first factor is a phase factor coming from the 3–photon vertex. Now, using momentum conservation $`l=p+q`$, $`ep=0`$, $`l\stackrel{~}{l}=0`$ and $`q^2=0`$, we see that this piece becomes
$$8i\alpha \frac{(ϵ\stackrel{~}{l})(q\stackrel{~}{l})(\stackrel{~}{l}j)}{l^2\stackrel{~}{l}^4}.$$
(5.3)
For small $`\stackrel{~}{l}`$, $`\mathrm{sin}\left(q\stackrel{~}{l}/2\right)`$ is just $`q\stackrel{~}{l}/2`$. We note that the dispersion relation of a longitudinally polarized photon remains $`q^2=0`$ since the self–energy corrections we found are transverse.
Next we turn to diagrams involving one loop corrections to the vertex. In the limit $`\stackrel{~}{l}0`$, the contribution to the amplitude becomes
$$i_2=8i\alpha g^4\frac{(ϵ\stackrel{~}{l})(q\stackrel{~}{l})(\stackrel{~}{l}j)}{l^2\stackrel{~}{l}^4}.$$
(5.4)
This is because in this limit the vertex is dominated by anomalous terms of the form $`\stackrel{~}{l}^\mu \stackrel{~}{l}^\nu \stackrel{~}{l}^\rho /\stackrel{~}{l}^4`$. Adding the two contributions together, we see that the amplitude becomes zero as required by gauge invariance.
Next we study the case when $`\stackrel{~}{p}0`$. Again, we study this particular case because we can isolate the possible singular terms. The most dangerous 1–loop diagram contributing to the scattering amplitude is the s–channel diagram involving corrections to the 3–gauge boson vertex. All other diagrams are sub–leading. The 1–loop contribution to the amplitude is then given by
$$i=8i\alpha g^4\frac{(ϵ\stackrel{~}{p})(q\stackrel{~}{p})(\stackrel{~}{p}j)}{l^2\stackrel{~}{p}^4}.$$
(5.5)
We now distinguish between two cases. First we note that if $`ϵ`$ is perpendicular to $`\stackrel{~}{p}`$, the amplitude is zero. Recall also from the previous section that this is the photon with un-corrected dispersion relation. If on the other hand $`ϵ`$ is along $`\stackrel{~}{p}`$, the amplitude is not zero but given by
$$i_2=8i\alpha g^4\frac{(q\stackrel{~}{p})(\stackrel{~}{p}j)}{l^2\stackrel{~}{p}^2}.$$
(5.6)
This contribution is cancelled by the tree level graph once we use the correct dispersion relation for the gauge boson. The tree level graph is given by
$$i_3=2i\alpha g^2\mathrm{sin}\left(\frac{q\stackrel{~}{p}}{2}\right)\frac{(ϵj)(2pq)}{l^2}.$$
(5.7)
The first factor is a phase factor. We also used the following relations: $`jl=0`$, $`ϵp=0`$ and $`q^2=0`$. Using momentum conservation, we also find that
$$2pq=l^2p^2.$$
(5.8)
Now, when the polarization of the gauge boson is along the $`\stackrel{~}{p}`$ direction, the dispersion relation gets corrected as follows
$$p^2=8\alpha g^2\frac{1}{\stackrel{~}{p}^2}.$$
(5.9)
Therefore, there is order a $`g^4`$ contribution from the tree–level graph given by
$$i_3=8i\alpha g^4(q\stackrel{~}{p})\frac{(\stackrel{~}{p}j)}{l^2\stackrel{~}{p}^2}.$$
(5.10)
Adding the two together we see that to order $`g^4`$ the amplitude is zero as it is required by gauge invariance.
The case involving scattering of scalars is more subtle because the two scalar – gauge boson vertex also contains linear divergences in $`\theta `$. We defer this case for an upcoming paper .
## 6 Conclusions
The most naive expectation about the non–commutative field theories is that they become commutative when the non–commutativity parameter $`\theta ^{ij}0`$ limit is taken. In other words, star–products are replaced by ordinary products of the fields in the Lagrangian in this limit, and one may expect that the theory becomes commutative also at the quantum level. As it was shown in this is generally not true due to the appearance of the new divergences at low non–commutative momenta. In the commutative theories, these divergences appear at high momenta in the superficially divergent loop integrals, but they can be eliminated by an appropriate choice of the regularization scheme. In the non–commutative theories, some of these divergences simply do not occur due to oscillating phases associated with star–products in the vertices which make the integrals finite. The non–commutative momenta thus play the role of the regulator. The dependence on these non–commutative momenta, however, does not disappear and manifests itself in the form of the new infrared divergences at small values of non–commutative momenta. These effects can also be characterized as non-analytic behavior in $`\theta `$.
A less naive expectation would be that a non–commutative gauge theory must be free of quadratic and linear poles at low non–commutative momenta since the corresponding commutative non–abelian gauge theory contains at most logarithmic divergences. In this paper we have shown that even this expectation does not hold, and both quadratic and linear poles appear in a generic gauge theory. The structure of these new poles is consistent with gauge invariance but not Lorentz invariance. These effects are local in the direction perpendicular to the non–commutativity plane and completely non–local in the non–commutative directions.
In supersymmetric gauge theories these poles cancel between the bosons and the fermions at the one loop level but even these theories typically contain logarithmic divergences at small values of non–commutative momenta. We expect that in $`𝒩=4`$ SYM theory even the logarithmic divergences do not occur, and thus this theory is completely free of anomalous effects in the small non–commutative momentum limit. This is the only theory that we know that reduces to its commutative counterpart in the limit $`\theta ^{ij}0.`$
## Acknowledgements
L.S. would like to thank Nathan Seiberg for numerous discussions. We would also like to thank Steve Shenker for useful conversations. This work was supported in part by the NSF grant 9870115.
## Appendix |
warning/0002/cond-mat0002273.html | ar5iv | text | # Geometrical phase effects in biaxial nanomagnetic particles
## I INTRODUCTION
For last decade, many physicists have explored the nanoscale magnetic particles to study the extrapolation of quantum mechanics to macroscopic realm. The magnetizations of these particles show the macroscopic quantum tunneling (MQT) from a metastable state to a stable one, or macroscopic quantum coherence(MQC) between two degenerate ground states separated by potential energy barrier. In the latter case, the geometrical phase, known as the Wess-Zumino phase, plays a crucial role in tunnel splitting. In particular, the biaxial nanomagnetic particles present some exotic phenomena related to this topological phase. In the absence of external magnetic field, the ground state tunneling vanishes only for half-integer spins, which is called as spin parity effect and related to well-known Kramer degeneracy. When the magnetic field is applied along the hard anisotropy axis, however, the tunnel splitting oscillates with magnetic field regardless of spin values. These spin parity effect and oscillations of tunnel splittings are interpreted as the interferences of geometrical phases between two opposite tunneling paths in the instantonic approach.
Most theoretical treatments so far about the oscillations of level splittings have been restricted to the instanton method for the mapped effective one-dimensional potential by performing the cos $`\theta `$ integration in spin coherent state representation or by direct particle mapping method. In these approaches, although the ground state tunneling can be obtained well, it is not possible to treat the tunneling between excited energy levels exactly. In fact, the result of numerical diagonalization for excited states is that the number of oscillations decreases one by one as the quantum number increases. In addition, due to the mapping onto one-dimension, the dominant tunneling path between two wells in $`(\theta ,\varphi )`$ space is not clear. This makes it difficult to perform the direct geometrical analysis of Wess-Zumino phase. Recently, Garg derived the splitting oscillations in high energy levels using the discrete WKB approximation. Another semiclassical approach is the complex periodic orbit theory, the extension of the trace formula in periodic orbit theory to tunneling systems. We obtained successfully the energy levels for not only the ground state but also all the higher excited states by applying this method to spin system for the first time, and showed the complete agreement with numerical diagonalization. Since this theory also uses the effective mapping potential obtained by cos $`\theta `$-integration of original Hamiltonian, the dominant tunneling path in phase space remains to be unspecified yet. The level splitting oscillations were confirmed by experiment very recently in the Fe<sub>8</sub> molecular clusters with spin $`S=10`$. They showed that, although the oscillations occurrs as predicted by theory, in quantitative aspect the oscillation periods is about 1.6 times larger than previous theoretical results, and this discrepancy can be resolved introducing the fourth-order terms in spin Hamiltonian.
In this paper, we find the dominant tunneling paths in phase space, and analyze geometrically the Wess-Zumino phase without any mapping in order to obtain the ground and excited level splitting oscillations with magnetic field along the hard axis. Furthermore, the experimentally discovered parity effect which is very similar to spin parity effect in the asymmetric system is derived theoretically using this geometrical analysis and complex periodic orbit theory. We also discuss the possibility of improving the discrepancy with experiments by adding the fourth-order terms in spin Hamiltonian. In section II, after finding the dominant tunneling path, we obtain the splitting oscillations for the tunneling between both the ground and the excited states in symmetric system by investigating the relation between the geometrical phase and the tunneling path in phase space without mapping onto a particle. In Section III, the parity effect occuring in the asymmetric system is derived using the arguement of Section II, and furthermore is also certified from the complex periodic orbit theory. We discuss the possibility of improving the quantitative discrepancy with experiment and give conclusions in Section IV.
## II LEVEL SPLITTING OSCILLATIONS IN THE SYMMETRIC SYSTEM
In this section, we calculate the both ground and excited level splitting oscillations by analyzing the topological Wess-Zumino phase geometrically in phase space. The tunneling rate in spin system is
$$\mathrm{\Gamma }=𝒟[\mathrm{cos}\theta ]𝒟[\varphi ]\mathrm{exp}(S_E)$$
(1)
in the path integral formalism, where $`𝒟`$ means the integration over all paths and $`S_E`$ is the Euclidean action which is given as
$$S_E=iS_{WZ}+𝑑\tau (\theta ,\varphi ).$$
(2)
Here, $`\tau `$ is the imaginary time, $`(\theta ,\varphi )`$ is the spin Hamiltonian in the $`(\theta ,\varphi )`$ representation, and
$`S_{WZ}`$ $``$ $`S{\displaystyle 𝑑\tau \dot{\varphi }\left[1\mathrm{cos}\theta (\tau )\right]}`$ (3)
$``$ $`{\displaystyle _0^\pi }𝑑\varphi \left\{Sp[\varphi (\tau )]\right\}`$ (4)
is the Wess-Zumino action with total spin number $`S`$. Here $`pS\mathrm{cos}\theta `$ and corresponds to the conjugate momentum to coordinate $`\varphi `$, since $`\varphi `$ and $`p`$ satisfy the Hamilton’s equations of motion in the spin system. When we consider the tunneling from $`\varphi =0`$ to $`\pi `$, Eq. (4) means the area surrounded by $`p=S`$, $`\varphi =0,\pi `$ lines and the tunneling path $`p(\varphi )`$ in $`(\varphi ,p)`$ phase space. \[See the Fig. 1 (a).\] Since the energy splitting due to tunneling, $`\mathrm{\Delta }E`$, is determined by Wess-Zumino action (the area mentioned above), i.e.,
$$\mathrm{\Delta }E\mathrm{cos}(S_{WZ}),$$
(5)
whenever this area becomes $`(n+1/2)\pi `$ with integer $`n`$ the tunnel splitting vanishes. In order to perform the geometrical evaluation of Wess-Zumino phase without mapping, it is essential to find the dominant tunneling path $`p(\varphi )`$ in the phase space. It will be shown below that in the case of $`k`$th state, the Wess-Zumino phase continuously decreases from $`(Sk)\pi `$ to zero as the external field $`h`$ increases, so that the moment of area being the half-integral multiple of $`\pi `$, at which the tunnel splitting quenches, occurs $`(Sk)`$ times.
Let us consider the Hamiltonian describing the biaxial spin system which is given by
$$=DS_x^2+E(S_z^2S_y^2)+g\mu _B𝐇𝐒,$$
(6)
where $`D`$ and $`E`$ are longitudinal and transverse anisotropy constants, respectively, $`g`$ is the gyromagnetic ratio and $`\mu _B`$ is the Bohr magneton. This Hamiltonian represents that the system has an easy axis in the $`x`$-direction, and hard axis in the $`z`$-direction. When the magnetic field $`𝐇`$ is applied along the $`z`$-direction, the system is symmetric and the reduced Hamiltonian is given in the spin coherent state representation by
$`^{}`$ $`=`$ $`[2(1\lambda )+(2\lambda 1)\mathrm{cos}^2\varphi ]p^22Shp`$ (8)
$`(1\lambda )S^2(2\lambda 1)S^2\mathrm{cos}^2\varphi ,`$
where $`^{}/(D+E)`$, $`\lambda =D/(D+E)`$ and $`hH/H_c=g\mu _BH/2(D+E)S`$ with coersive field $`H_c`$. In $`(p,\varphi )`$ phase space, $`^{}`$ has two degenerate minima when $`p=Sh`$ and $`\varphi =0,\pi `$. At $`\varphi =\frac{\pi }{2}`$, however, $`p=Sh/h_1`$ with $`h_12(1\lambda )`$ gives the saddle point for $`0<h<h_1`$, local minimum for $`h_1<h<h_2`$ with $`h_2\sqrt{2(1\lambda )}`$ and global minimum for $`h_2<h<1`$, respectively, if $`p`$ can be also defined beyond the range $`SpS`$. These pictures of various external fields are shown in Fig. 1 as the energy contour plots.
In order to get an information about the dominant tunneling path, we revisit, for a moment, the mapping onto a one-dimensional particle problem by integration over $`\mathrm{cos}\theta `$ in Eq. (1). The resulting Euclidean action is given by
$$S_E(\varphi )=iA(\varphi )+S\sqrt{2\lambda 1}𝑑\tau \left[\frac{1}{2}M(\varphi )\dot{\varphi }^2+V(\varphi )\right],$$
(9)
where the dot means the Euclidean time derivative and the imaginary part of Euclidean action which is responsible for the phase effect is
$$A(\varphi )=S𝑑\varphi \left[1\frac{h}{1(2\lambda 1)\mathrm{sin}^2\varphi }\right],$$
(10)
the $`\varphi `$-dependent effective mass is
$$M(\varphi )=\frac{1}{1(2\lambda 1)\mathrm{sin}^2\varphi }$$
(11)
and the effective potential is
$$V(\varphi )=\frac{1}{2}\mathrm{sin}^2\varphi \left[1\frac{h^2}{1(2\lambda 1)\mathrm{sin}^2\varphi }\right].$$
(12)
For the ground state tunneling, if we compare Eq. (4) and Eq. (10), the tunneling path $`p(\varphi )`$ is found to be the second
It is emphasized that the tunneling path $`p(\varphi )`$ in the phase space $`(\varphi ,p)`$ is just the path satisfying $`H^{}(\varphi ,p)/p=0`$. This means that the tunneling path follows the minimum energy positions at all $`\varphi `$ values. These tunneling paths for vaious fields are shown in Fig. 1 as thick solid lines. This simple condition for dominant tunneling path in the phase space $`(\varphi ,p)`$ works well only when the Gaussian integration is possible as in the present case.
Now, we calculate the Wess-Zumino action (the area) mentioned above as $`h`$ increases from 0 to 1. When $`h=0`$, the phase is just $`S\pi `$ (whole area of phase space considered) because the dominant tunneling path is $`p(\varphi )=0`$. It is important to note that, although the quantum mechanical ground state energy is not the well minimum, the path starts and ends at well minima, because the Euclidean action with the energy of the minimum gives the information about the ground state splitting in semiclassical theory. When $`0<h<h_1`$ \[Fig. 1 (a)\], the path is confined within a range $`0<p<S`$. Therefore, we can easily calculate the area enclosed by the path and $`p=S`$ line. Until $`h`$ becomes $`h_1`$, where the north-pole becomes the saddle point of energy, the area continuously decreases from $`S\pi `$. For $`h_1<h<h_2`$ \[Fig. 1 (b)\], however, the local minimum exists outside the $`p=S`$ line. In this case the tunneling path runs beyond the $`p=S`$ line and thus, the Wess-Zumino action becomes the area of (region B + region C - region A). As $`h`$ increases the action continues to decrease and becomes zero at $`h=h_2`$. As a result, within a range $`0<h<h_2`$, the quenching take places $`S`$ times for the ground state by Eq.(5). The range $`h_2<h<1`$ is more or less subtle \[Fig. 1 (c)\]. The local minimum turns into the global one, i.e., the point $`(Sh/h_1,\pi /2)`$ becomes lower than the well minimum. In this case, the path giving the Wess-Zumino action should be determined by somewhat different way, because there exists the contour of $`E_0`$ (well minimum energy). If the path obtained by the previous method encounters the contour of $`E_0`$, it should follow the contour in a real time due to energy conservation. Therefore, the action is equivalent to the area of \[region B + region C - region A \] of Fig. 1(c) which always vanishes in the range of $`h_2<h<1`$. This vanishing of the action is closely related to the cancellation between resonant tunneling and quenching phenomena and also between the Euclidean action of real time motion and Wess-Zumino action.
The extention to the excited states is straightforward. The $`k`$th energy level of each well can be calculated from the well-known EBK (Einstein-Brillouin-Keller) quantization rule,
$$p(\varphi )𝑑\varphi =2\pi \left(k+\frac{1}{2}\right),k=0,1,2,\mathrm{},$$
(13)
In this quantization, whenever the action for closed orbit in one well differs by $`2\pi `$, the energy level is given. It will be shown that, by this fact, the number of splitting oscillations decreases one by one with increasing quantum number. When $`h=0`$, there exist the contours of $`E_k`$ centered at $`p=0`$ and $`\varphi =0,\pi `$ whose areas are $`2k\pi `$ for $`k`$th excited state (see also Fig. 1). If we use the analogy for the case $`h_2<h<1`$ of the ground state tunneling, the area of $`k\pi `$ is excluded, due to the energy conservation, for the $`k`$th excited state ($`k\pi /2`$ for each well) in the phase space considered, i.e., $`(Sk)\pi `$. As $`h`$ increases the continuous reduction to zero is same as the ground state case. This leads to the $`(Sk)`$ oscillations for $`k`$th energy level. The calculated periods of oscillations for both ground and excited state levels completely agree with the numerical diagonalizations. Comparing with the previous semiclassical theories, in this theory the dominant tunneling path in $`(\varphi ,p)`$ phase space can be found to make possible to perform the geometrical analysis of the effect of topological phase, and in addition to the gound state case the structures of splitting oscillations for the excited states can be obtained easily on the basis of EBK quantization scheme.
## III PARITY EFFECT IN THE ASYMMETRIC SYSTEM
Unlike the symmetric case in Sec. II, in this section we treat the asymmetric system having the longitudinal field as well as the transverse field. In such system, we derive the parity effect which was discovered in recent experiment, by using the simple geometrical argument developed in section II, and also lead to the same result in the semiclassical complex periodic orbit theory. This parity effect is very similar to the spin parity effect proposed by Loss et. al. and von Delft and Henley. This is observed when one measures the transtion between the ground state in upper well and the $`k`$th excited state in lower well, by applying the external field $`h^{}`$ along the easy axis (longitudinal field) which makes the system asymmetric in addtion to $`h`$ in the hard direction (tansverse field). For the case of the tunneling from the ground state of the upper well to even $`k`$th state of lower well, the phase of the oscillation is same as that of the symmetric case, while for the case of tunneling to odd $`k`$th state of the lower well the phase is shifted by $`\pi /2`$.
This interesting observation can naturally arises in our above analysis. As shown in Fig. 2, the energy structure of spin system is asymmetric due to $`h^{}`$. At certain values of $`h^{}`$ the ground state energy level $`E_0^l`$ in upper (left) well exactly coincides with the $`k`$th excited level $`E_k^r`$ in lower (right) well, so that one can consider the tunnel splittings. The tunneling path starts from the minimum at upper well, and when it exits the barrier it follows the contour of $`k`$th excited state in the lower well due to energy conservation. Therefore, by EBK quantization relation and analogous arguments with the excited states in Section II, the area $`k\pi /2`$ in the $`(\varphi ,p)`$ phase space is excluded when $`h=0`$ case. As the transverse field $`h`$ increases, total area surrounded by this path continuously decreases from $`(S\pi k\pi /2)`$ to zero in the same manner with the symmetric case except the $`k\pi /2`$ shift. This leads to the parity effect in the asymmetric system.
Next, let us consider this effect in complex periodic orbit theory. In quantum mechanics the energy spectrum of the system can be derived from the singularities of the trace of the energy-dependent Green’s function which is given by
$$g(E)=\mathrm{Tr}\left[(E\widehat{H})^1\right]=\underset{n}{}\frac{d_n}{EE_n},$$
(14)
where Tr means the trace and $`d_n`$ is the degeneracy factor. The semiclassical approximation of Eq. (14) was given by Gutzwiller in the form
$$g_{\mathrm{sc}}(E)=\frac{1}{i}\underset{j}{}T_j\underset{r=1}{\overset{\mathrm{}}{}}e^{ir(S_j\mu _j\pi /2)},$$
(15)
where $`T_j`$, $`S_j`$ and $`\mu _j`$ are the period, the action and the Maslov index of the $`j`$th primitive classical periodic orbit, and index $`r`$ corresponds to the repetition of primitive periodic orbit. In evaluating this sum, we restrict on the one-dimensional effective particle problem. The system with $`h^{}=0`$ which are given as Eq. (9) - Eq. (12) is recently treated in detail within the complex periodic orbit theory in Ref. . When the longitudinal magnetic field $`h^{}`$ is applied, the effective potential of Eq. (12) becomes asymmetric as shown in Fig. 3. This figure presents the case when the ground state energy of upper well coincides with the $`k`$th excited level of lower well. The classical orbits in two wells and tunneling orbits within the barrier are represented as solid and dotted lines, respectively. The directions denoted by arrows are given by the rules of Ref. .
Now we should consider all possible periodic orbits in order to perform the summation in Eq. (15). Let us take the classical segments in the respective wells as the half of whole orbit having the actions $`W_1S_1/2`$ and $`W_2S_2/2`$, and similarly for the tunneling segments within the barrier having an action $`\mathrm{\Theta }S_c/2`$. Let us consider the case where the ground state energy in the upper well is equal to the $`k`$th excited state in the lower well. Note, then, that the relation between two actions $`W_1`$ and $`W_2`$ is
$$W_2=W_1+k\pi ,k=0,1,2,\mathrm{},$$
(16)
due to EBK quantization rule. The possible periodic orbits can be parametrized as follows: i) the paths that start and end at $`\varphi _A`$ or $`\varphi _A^{}`$ without complete rotations, ii) the paths that start at $`\varphi _A(\varphi _A^{})`$ and end at $`\varphi _A^{}(\varphi _A)`$, i.e., having complete rotations (See Fig. 3. $`\varphi _A`$ and $`\varphi _A^{}`$ are the equivalent position in $`\varphi `$ coordinate). The Wess-Zumino phase contributes only to the paths belonging to category ii) as
$$\alpha =2\pi S\left(1\frac{h}{h_2}\right).$$
(17)
For the paths that start at $`\varphi _A`$ and end at $`\varphi _A^{}`$, i.e., proceed from left to right, we take the Wess-Zumino phase as positive, whereas for the paths from right ($`\varphi _A^{}`$) to left ($`\varphi _A`$) we take as negative. The Maslov index is easily evaluated following the method for double well potential in Ref. , i.e., each classical segment $`W_1`$, $`W_2`$ accompanies the additional phase $`+\pi /2`$, while the tunneling segment $`\pi /2`$. The all possible primitive periodic orbits are drawn graphically in Fig. 4, where the solid line corresponds to $`\mathrm{exp}[i(W_1+\pi /2)]`$, thick solid line to $`\mathrm{exp}[i(W_2+\pi /2)]`$ and dotted line to $`\mathrm{exp}[\mathrm{\Theta }i\pi /2]`$. Therefore, the semiclassical trace of Green’s function can be expressed as the summation over all repetitions of these orbits, and its poles are found to be when
$`1e^{i(2W_1+\pi )}+e^{2\mathrm{\Theta }}+{\displaystyle \frac{e^{2\mathrm{\Theta }}}{1e^{i(2W_2+\pi )}+e^{2\mathrm{\Theta }}}}`$ (18)
$`\times [e^{i(W_1+W_2+\pi )}(e^{i\alpha }+e^{i\alpha })`$ (19)
$`e^{i(2W_1+\pi )}+e^{i(2W_2+\pi )}]=0`$ (20)
If we neglect, for the moment, the tunneling probability $`e^{2\mathrm{\Theta }}`$, then the poles give the EBK energies $`E_n`$ for upper well where
$$S_1(E_n)=2W_1(E_n)=2\pi \left(n+\frac{1}{2}\right),$$
(21)
with $`n=0`$ for the ground state. Therefore, we can expand the actions around the EBK energy like
$`2W_1(E)`$ $`=`$ $`\pi +T_1(EE_0)+\mathrm{},`$ (22)
$`2W_2(E)`$ $`=`$ $`\pi +2k\pi +T_2(EE_0)+\mathrm{}`$ (23)
near the ground state of upper well, where $`T_1`$ and $`T_2`$ are given by $`dS_1/dE`$ and $`dS_2/dE`$, respectively. Here, we used the fact that the difference of two actions at EBK energy is $`2\pi k`$ for the transition between ground state in upper well and $`k`$th excited state in lower well. Expanding Eq. (20) up to the first-order in $`e^{2\mathrm{\Theta }}`$,
$$T_1T_2(EE_0)^2e^{2\mathrm{\Theta }}\left[(1)^k\left(e^{i\alpha }+e^{i\alpha }\right)+2\right]=0$$
(24)
The final energy splitting is obtained as
$`E`$ $`=`$ $`E_0\pm {\displaystyle \frac{2}{\sqrt{T_1T_2}}}e^\mathrm{\Theta }\mathrm{cos}{\displaystyle \frac{\alpha }{2}},\mathrm{even}k,`$ (25)
$`E`$ $`=`$ $`E_0\pm {\displaystyle \frac{2}{\sqrt{T_1T_2}}}e^\mathrm{\Theta }\mathrm{sin}{\displaystyle \frac{\alpha }{2}},\mathrm{odd}k,`$ (26)
which is just the parity effect in the asymmetric system of Ref. .
## IV DISCUSSIONS AND CONCLUSIONS
Until now, we have analyzed the Wess-Zumino phase of biaxial spin system in $`(p,\varphi )`$ phase space, and obtained not only the same results with those in the mapping formalism for ground state, but also the excited energy level splitting oscillations. Furthermore, the parity effect in the asymmetric system shown in recent experiment is derived by both the geometrical analysis and the complex periodic orbit theory. However, the quantitative discrepancy with experiment in the period of oscillation exists, i.e., about 1.6 times larger oscillation period than expected in the theory. In order to resolve these discrepancies the authors of Ref. introduced the fourth-order term $`C(S_+^4+S_{}^4)`$ in spin Hamiltonian (Eq. (8)), where $`C`$ is the adjustable parameter which is shown to be $`2.9\times 10^5`$K in Kelvin unit through numerical diagonalization, when anisotropy constants are given by $`D=0.292`$K and $`E=0.046`$K. Then the fourth-order Hamiltonian is expressed in the $`S_x`$-representation as
$$_1=^{}+2C(S_z^4+S_y^46S_z^2S_y^2)$$
(27)
within the semiclassical approximation. In fact, if we consider the classical commutation relation satisfied by angular momenta, the cubic and quadratic terms must be added in this Hamiltonian. But we have ignored them since their contributions are negligible. It is not possible to get an anaytical effective potential through integration, since Eq. (27) is beyond the Gaussian approximation. However, the qualitative analysis of this Hamiltonian can be performed within our theory on the basis of the energy structure of $`_1`$ and the dominant tunneling paths in $`(\varphi ,p)`$ phase space.
The structure of energy barrier for $`_1`$ is different from that for $`^{}`$ as shown in Fig. 5. Since $`C`$ is negative, the effect of $`S_z^2S_y^2`$ in Eq. (27) is to raise the barrier around $`\varphi =\pi /2`$ and the mid-value of $`0<p<S`$, whereas $`S_z^4(S_y^4)`$ to lower the barrier at $`z(y)`$-axis. Therefore, as the field $`h`$ increases, the saddle point moves more slowly to the north-pole, compared with the $`C=0`$ case. If we use the same analogy about finding the dominant tunneling path, this means that the slower decrease of the area surrounded by tunneling path, and thus the larger oscillation period. Although it is possible to get a qualitative correction to $`C=0`$ case in the desirable direction, in order to be consistent with experiment, the adjustable parameter $`C`$ must be about 3 times greater than the present value. The quantitative disagreement implies that the tunneling paths which follow all energy minima are no longer the dominant paths in the quartic case. In order to treat the quartic case exactly, new theoretical approach in order to obtain the dominant tunneling path in phase space is needed.
In conclusion, we derive both the ground and excited tunnel splitting oscillations in the biaxial nanomagnetic particle with the magnetic field along the hard anisotropy axis by finding the dominant tunneling path and geometrically analyzing the topological Wess-Zumino phase in the phase space. All the results of analysis are in agreement with the numerical diagonalization. Furthermore, the interesting parity effect in the asymmetric case is naturally clear in this analysis, and also certified within the complex periodic orbit theory. We also discussed the possibility of improving the discrepancies with experiment in the period, by introducing the quartic terms in spin variables into Hamiltonian. In order to resolve the discrepancy quantitatively, the new approach to find the dominant tunneling path in quartic case is required.
ACKNOWLEDGMENTS
We thank to Dr. Cheol-Hong Kim for very helpful discussions on mathematics. |
warning/0002/nucl-ex0002006.html | ar5iv | text | # Statistical signatures of critical behavior in small systems
## I Introduction
Beginning in the 1970’s significant advances in the understanding of nuclear multifragmentation were made possible with the advent of high statistics inclusive experiments. Typically, only one intermediate mass fragment ( $`3Z_f30`$ ) was detected per event. From these inclusive studies came the first evidence that intermediate mass fragments (IMFs) were associated with a simultaneous multi-body breakup of a system which had undergone expansion. A study of the fragment mass yield distribution obtained in an inclusive gas jet experiment conducted at Fermilab contained the first indication that nuclear multifragmentation might be related to critical phenomena normally observed in macroscopic systems . The Purdue Group was the first to make the suggestion that the observed power law in the fragment yield distribution might result from a system whose excitation energy was comparable to its total binding energy . The exponent of the power law was $`2\tau 3`$, within the range expected for a system near its critical point. The presence of the power law and the value of the exponent, coupled with the strong similarity of the nuclear and van der Waals potentials, led the Purdue group to suggest that multifragmentation of nuclei might be analogous to a fluid undergoing a continuous phase transition from a liquid to a gas. Furthermore, the Fisher Droplet Model (FMD) -, used to describe condensation in a fluid system near its critical point, after modification for nuclear physics effects, was capable of describing the isotopic yields of 50 fragments with one set of parameters , . The temperature of the system was determined to be about 5 MeV , a reasonable value considering that the average binding energy per nucleon in a nucleus is approximately 8 MeV. The success of this approach reinforced the notion that multifragmentation was both a thermal process and that it was related to critical phenomena.
With the advent of exclusive experiments capable of detecting all of the charged reaction products, the possibility of studying multifragmentation on an event-by-event basis became a reality. High statistics exclusive experiments in which the fragmenting system is characterized according to its nucleon number and excitation energy permit both the correlation of dynamical and statistical information and the study of fluctuations in experimental observables. Fluctuations are central to all critical phenomena, and indeed, such fluctuations are apparent in exclusive multifragmentation data. In this paper, the focus will be on the statistical signals of multifragmentation data observed in the EOS experiment -. Comparisons will be made with two other systems, one of which exhibits critical behavior and one of which does not.
Much of the pioneering work in understanding the statistical aspects of multifragmentation has been performed by Campi - and Mekjian -. Both efforts have compared multifragmentation data to model systems in order to gain some insight into the nuclear breakup process. In this paper, many of the ideas suggested by these authors are followed and applied to both the EOS data and the model systems in order to demonstrate which of the many suggested signals are useful for the identification of critical behavior. A major goal of this paper is to present a comprehensive review of several methods proposed for detecting signals of critical phenomena in multifragmentation.
It is tempting to compare the experimental data to dynamical models that attempt to describe nuclear multifragmentation. However, the task of modeling multifragmentation from the initial collision phase of the reaction to freeze-out has proven to be a daunting task. Models that adequately describe the initial stage of the reaction - do not satisfactorily describe the fragment formation stage, in either statistical or dynamical aspects. Likewise, the most successful models in describing the statistical properties of nuclear multifragmentation -, assume thermodynamic equilibrium, yet fail to adequately match the dynamical features of the data.
Molecular dynamical approaches, which have enjoyed considerable success in describing critical behavior in classical systems -, have not been conclusive in describing nuclear multifragmentation and at times have yielded contradictory results , . Later studies suggested flaws in the application of molecular dynamical models to nuclear multifragmentation, therefore calling into question the conclusions drawn from the earlier studies .
The most striking of the early theoretical efforts came from Campi’s analysis of a few hundred completely reconstructed emulsion multifragmentation events and the comparison of these data to clusters generated from a percolation calculation , . In this series of papers it was shown that the fragment distributions from multifragmentation bore a striking similarity to the cluster distributions from percolation lattices. This analysis provided strong evidence that multifragmentation was a statistical process which appeared to be related to critical phenomena. In this analysis another estimate of the exponent $`\tau `$ was made which agreed with the first measurements from the Purdue Group and several later analyses of various fragment distributions.
In the early 1990’s the ALADIN Group from GSI performed several multifragmentation experiments -. Of particular importance was the “rise and fall” of multifragmentation. In one analysis the ALADIN group plotted the “rise and fall” curve of the production of IMFs versus an observable related to the excitation energy of the reaction for several multifragmenting systems. With the appropriate scaling the data collapsed to a single curve suggesting that the multifragmenting systems retained no memory of the reaction entrance channel. This is expected for an equilibrated system.
The results of some statistical analyses of multifragmentation data could be interpreted to suggest that multifragmentation is a sequential decay in contrast to the phase transition picture. The same sort of statistical analysis has also been applied to explicitly simultaneous models and produced results that were similar to those of multifragmentation data. Thus those signals could be interpreted as evidence for either sequential or simultaneous multifragmentation .
This last effort puts into focus the main question in this work: what type of analysis of the statistical aspects of a cluster distribution can provide the most insight into the nature of the mechanism which created the clusters? Specifically, can those systems which contain critical behavior be distinguished from those which do not? It will be argued that this question has two answers. Analysis of the insensitive features of the cluster distribution cannot make the above mentioned distinction . However, an analysis of the sensitive features of the cluster distribution will be shown to provide deeper insight into the cluster production mechanism. This type of analysis has been previously reported for clusters resulting from nuclear multifragmentation -. Note that the more generic term cluster will be used to refer to any composite of constituents, whether these be molecules of a fluid, nuclear fragments or percolation clusters.
The method employed to address the above question is as follows. The same analysis is performed on the cluster distributions produced by three different systems. In one case, clusters are generated by randomly partitioning an integer. Such one-dimensional partitioning does not posses critical behavior indicative of a continuous phase transition. In the second case, three-dimensional bond building percolation is used to produce clusters. Percolation is well-known mathematical construct that possesses a continuous phase transition, i.e. a critical point. Finally, the cluster distributions resulting from the multifragmentation of gold nuclei are analyzed. Although it is not known, a priori, whether the nuclear multifragmentation bears any relation to critical phenomena, it will be seen that the analysis presented in this work yields suggestive results.
This paper is organized as follows. In section II a description of each system is presented. In section III the Fisher Droplet Model is reviewed. In section IV-A the insensitive signatures of the cluster distributions for all systems are examined. In section IV-B the sensitive signatures are examined. Sections V and VI present possible corrections to the analysis of the multifragmentation data. Finally, Section VII discusses the conclusions reached upon the completion of the analyses in sections IV and V. Throughout this paper the term continuous phase transition will be used instead of second order phase transition, the latter from the outdated Ehrenfest theory of phase transitions .
## II Description of systems under study
### A 1.0 A GeV Au $`+`$ C multifragmentation
Approximately 40,000 fully reconstructed events ($`76Z_{observed}82`$) were collected with the EOS experimental apparatus discussed in ref. . In the collision of the projectile gold nucleus ($`197`$, $`79`$) and the target carbon nucleus, so-called prompt nucleons are knocked out of the gold nucleus by quasi-elastic and inelastic collisions between projectile and target nucleons . Immdeiately following the collision, he gold projectile remnant is in an excited state with fewer than 197 nucleons. The excited remnant cools and expands, evolving to the neighborhood of the critical point in the temperature-density plane , where clusters condense from a high temperature low density vapor of nucleons. The charge and mass of the projectile remnant, $`Z_0`$ and $`A_0`$, were determined for each event by subtracting the charge and mass of the prompt particles from the charge and mass of the gold nucleus . Prompt particles have $`Z_f=0`$, $`1`$ and $`2`$ and are removed from the cluster distributions analyzed in this work. Only clusters created from the excited gold projectile are considered in the ensuing analysis. For events with the lowest charged particle multiplicities, $`m`$, the remnant had $`Z_076`$, $`A_0194`$ and $`E^{}/A_02`$ MeV$`/`$nucleon, while for events with the highest multiplicities the remnant had $`Z_039`$, $`A_092`$ and $`E^{}/A_016`$ MeV$`/`$nucleon .
Clusters of a given charge, $`Z_f`$, were counted on an event by event basis to determine the cluster charge distribution, $`N_{Z_f}`$. In this analysis, although the mass number of the clusters is of interest, a cluster’s charge will be used as an index. Mass numbers for clusters of charge one and two were measured in the EOS time projection chamber. Clusters with $`Z_f3`$ were assigned a mass number, $`A_f`$, by multiplying the cluster charge by the mass to charge ratio of the excited gold projectile remnant; for low $`m`$ events $`A_0/Z_02.55`$ and for high $`m`$ events $`A_0/Z_02.36`$. This procedure provided an estimate of a cluster’s mass number prior to any secondary decay effects. It is assumed that on average $`N_{A_f}=N_{Z_f}`$. Finally, it is the normalized cluster distribution, $`n_{A_f}=N_{A_f}/A_0(m)`$, that is used in the analysis presented in this paper.
### B Percolation
Bond building percolation calculations were performed on three dimensional simple cubic lattices of 216 sites. Cluster distributions for $`100,000`$ lattice realizations were generated in the standard fashion by forming bonds between sites. Bonds were either active (on) or inactive (off) according to the following algorithm.
The control parameter (e.g. temperature in thermodynamic systems) for percolation is the lattice probability, $`p_l`$. A single value of $`p_l`$ was chosen for the entire lattice. All probabilities were between 0 and 1. Next, a bond probability, $`p_{b_i}`$, was randomly chosen from a uniform distribution on (0,1) for the $`i^{th}`$ bond. If $`p_{b_i}`$ was less than $`p_l`$, then the $`i^{th}`$ bond was active and two sites were joined into a cluster. This process was performed for each bond in the lattice.
At low values of $`p_l`$, few bonds were formed resulting in a high multiplicity, $`m`$, of small clusters, a distribution analogous to the gaseous phase of a fluid. At high values of $`p_l`$, many bonds were formed resulting in a low multiplicity of mostly large clusters, analogous to the liquid phase of a fluid. In an infinite lattice the phase transition occurs at a unique value of the lattice probability, $`p_c`$, when the probability of forming a percolating cluster changes from zero to unity.
To examine the behavior of the average cluster distribution, the number of clusters of size $`A_f`$ per lattice site was calculated by histogramming the $`100,000`$ lattice realizations into 100 bins from $`0`$ to $`1`$. The use of $`m`$ as a control parameter and the ensuing effects on signatures of continuous phase transition were investigated by calculating the average number of clusters of size $`A_f`$ with the $`100,000`$ lattice realizations histogrammed in units of $`m`$.
### C Random partitions
Random partitions were generated from 79 total system constituents, chosen to approximate the number of charges in the gold multifragmentation system. The algorithm is as follows. First a random choice of $`m`$ was made from a uniform distribution on (1,79). Next the maximum size of a cluster, $`A_{max}^1`$, for an event with $`m`$ was determined; this depended on the constraints of the system size, $`A_0=79`$ and the choice of $`m`$. The size of the first cluster, $`A_1`$, was then randomly chosen from a uniform distribution on $`(1,A_{max}^1)`$. There were then $`m1`$ clusters to be generated from $`79A_1`$ constituents. The maximum size of a cluster for an $`m1`$ event from a $`791`$ constituent system was determined: $`A_{max}^2`$. The size of the second cluster, $`A_2`$, was then randomly chosen form a uniform distribution on $`(1,A_{max}^2)`$. This process was repeated until all constituents belonged to a cluster. $`100,000`$ partitions were generated in this manner. This particular weighting results in a power law cluster distribution .
## III Review of the Fisher Droplet Model
The focus of most studies of phase transitions is on standard thermodynamical variables such as a system’s temperature, density, compressibility, etc. These quantities are difficult or impossible to measure directly in present nuclear multifragmentation experiments. Thus a theory which addresses quantities accessible to MF experiments is needed. To that end Fisher’s gas-to-liquid phase transition model, based on Mayer’s condensation theory, is followed , , .
Fisher begins his model, called the Fisher Droplet Model (FDM) hereafter, by writing the free energy for the formation of clusters of size $`A_f`$ as:
$`\mathrm{\Delta }G_{A_f}`$ $`=`$ $`k_bTA_f\mathrm{ln}(g(\mu ,T))`$ (2)
$`k_bT\mathrm{ln}(f(A_f,T))+k_bT\tau \mathrm{ln}(A_f)+\mathrm{}`$
Where $`k_b`$ is the Boltzmann constant and the $`g`$-term is the bulk formation energy, or volume term and:
$$g(\mu ,T)=exp[(\mu \mu _{coex})/k_bT],$$
(3)
where $`\mu `$ is the chemical potential and $`\mu _{coex}`$ is the chemical potential along the coexistence curve.
The $`f`$-term is related to the surface free energy of cluster formation. It’s a form given by Fisher is:
$$f(A_f,T)=exp[a_0\omega A_f^\sigma ϵT_c/k_bT],$$
(4)
where $`\sigma `$ is a critical exponent and is related to the ratio of the dimensionality of the surface to the dimensionality of the volume, $`a_0`$ is a constant of proportionality relating the average surface area of a droplet to its number of constituents and $`\omega `$ is the surface entropy density; $`ϵ`$ is a measure of the distance from the critical point. For usual thermodynamic systems $`ϵ=(T_cT)/T_c`$, in the percolation treatment $`ϵ=(p_lp_c)/p_c`$ and for multifragmentation $`ϵ=(m_cm)/m_c`$ will be used. All formulations of $`ϵ`$ are such that $`ϵ>0`$ ($`ϵ<0`$) corresponds to the liquid (gas) region. This form of the surface free energy is applicable on only one side of the critical point, the single phase side. A more general form suggested by efforts from percolation theory - that can be applied on both sides of the critical point and leads to a power law which describes the behavior of the order parameter is:
$$f(z)=Aexp[(zB)^2/C],$$
(5)
where the scaling variable, $`z`$, is
$$z=A_f^\sigma ϵ.$$
(6)
The physical interpretation of the parameters $`A`$, $`B`$ and $`C`$ is an open question.
Finally $`\tau `$ is another critical exponent depending principally on the dimensionality of the system and has its origins in considerations of a three dimensional random walk of a surface closing in on itself. For three dimensions $`2\tau 3`$ . In eq. (2), $`q_0`$ is a normalization constant which will be shown to depend solely on the value of $`\tau `$ .
From the free energy of cluster formation the average cluster distribution normalized to the size of the system is:
$$n_{A_f}(ϵ)=\mathrm{exp}(\mathrm{\Delta }G_{A_f}/k_bT)=q_0A_f^\tau f(z)g(\mu ,T)^{A_f}.$$
(7)
At the critical point, $`ϵ=0`$, both $`f`$ and $`g`$ are unity and the cluster distribution is given by a pure power law:
$$n_{A_f}(ϵ)=q_0A_f^\tau .$$
(8)
If the first moment of the normalized cluster distribution is considered at the critical point then :
$$M_1(ϵ=0)=\underset{A_f}{}n_{A_f}(ϵ)A_f=q_0\underset{A_f}{}A_f^{1\tau }=1.0$$
(9)
when the sum runs over all clusters. From eq. (9) it is obvious that the value of the overall cluster distribution normalization constant, $`q_0`$, is dependent on $`\tau `$ via a Riemann $`\zeta `$-function:
$$q_0=1.0/\underset{A_f}{}A_f^{1\tau }.$$
(10)
The above is true only if the scaling assumptions in the FDM apply to all clusters. For finite size systems even at the critical point this is only approximately true. However, it will be seen that eq. (10) holds reasonably well at the critical point for systems with a continuous phase transition over some range in cluster size.
In the FDM it is assumed that all clusters of size $`A_f`$ can be treated as an ideal gas, so that the total pressure of the entire cluster distribution can be determined by summing all of the partial pressures:
$`P/(k_bT)`$ $`=`$ $`{\displaystyle \underset{A_f}{}}n_{A_f}(ϵ)`$ (11)
$`=`$ $`q_0{\displaystyle \underset{A_f}{}}A_f^\tau f(z)g(\mu ,T)^{A_f}=M_0(ϵ)`$ (12)
Is is clear from eq. (12) that the pressure of the system is related to the zeroth moment of the cluster distribution.
The density is then:
$`\rho ={\displaystyle \frac{P}{\mu }}`$ $`=`$ $`q_0{\displaystyle \underset{A_f}{}}A_f^{1\tau }f(z)g(\mu ,T)^{A_f}`$ (14)
$`={\displaystyle \underset{A_f}{}}n_{A_f}(ϵ)A_f=M_1(ϵ).`$
The density is given by the first moment of the cluster distribution.
It is now a simple matter to derive the power law which describes the divergence of the isothermal compressibility, $`\kappa _T`$. By definition:
$$\kappa _T=\frac{1}{V}(\frac{V}{P})_T=\frac{1}{\rho }(\frac{\rho }{P})_T.$$
(15)
Noting that $`k_bT\rho =g(\mu ,T)(P/g(\mu ,T))`$, eq. (15) can be rewritten as:
$`\kappa _T`$ $`=`$ $`{\displaystyle \frac{1}{\rho ^2}}\times `$ (17)
$`\left(g(\mu ,T){\displaystyle \frac{P}{g(\mu ,T)}}+g(\mu ,T)^2{\displaystyle \frac{^2P}{g(\mu ,T)^2}}\right)_T,`$
which leads to:
$`\kappa _T`$ $`=`$ $`(\rho k_bT)^1+(\rho ^2k_bT)^1{\displaystyle \underset{A_f}{}}n_{A_f}(ϵ)A_f^2`$ (18)
$`=`$ $`(\rho k_bT)^1+(\rho ^2k_bT)^1M_2(ϵ).`$ (19)
The sum in the second term illustrates the relation of the second moment of the cluster distribution, $`M_2(ϵ)`$, to the isothermal compressibility. The sums in eq. (12), (14) and (19) run over all clusters in the gas and exclude the bulk liquid drop. In percolation and multifragmentation the largest cluster on the liquid side of the critical point will be considered as the liquid drop and will thus be excluded from the sum. On the gas side of the critical point, the sum runs over all clusters as there is no longer a liquid drop.
In the thermodynamic limit, large $`A_f`$ dominate the sum so that it may be treated as an integral giving:
$$\kappa _T=(\rho k_bT)^1+(\rho ^2k_bT)^1_0^{\mathrm{}}n_{A_f}(ϵ)A_f^2𝑑A_f.$$
(20)
Working along the liquid-gas coexistence curve so that $`g(\mu ,T)=1`$ eq. (20) reduces to:
$$\kappa _T=(\rho k_bT)^1+(\rho ^2k_bT)^1_0^{\mathrm{}}A_f^{2\tau }f(z)𝑑A_f.$$
(21)
A change of variables from $`A_f`$ to $`z`$ shows that near the critical point:
$`\kappa _T`$ $``$ $`(\rho ^2k_bT)^1\left|{\displaystyle \frac{q_0}{\sigma }}{\displaystyle _0^\pm \mathrm{}}𝑑zf(z)|z|^{\frac{3\tau \sigma }{\sigma }}\right||ϵ|^{\frac{\tau 3}{\sigma }}`$ (22)
$`=`$ $`(\rho ^2k_bT)^1\mathrm{\Gamma }_\pm |ϵ|^\gamma `$ (23)
This is the so-called $`\gamma `$-power law which describes the divergence of the isothermal compressibility and the second moment of the cluster distribution near the critical point. The scaling relation between the exponents $`\gamma `$, $`\sigma `$ and $`\tau `$ is:
$$\gamma =\frac{3\tau }{\sigma }.$$
(24)
The absolute normalization constants of the $`M_2(ϵ)`$ power law depend on the scaling function, $`f(z)`$, the exponent $`\sigma `$ and the overall normalization of the cluster distribution, $`q_0`$, which in turn depends on the exponent $`\tau `$:
$$\mathrm{\Gamma }_\pm =\left|\frac{q_0}{\sigma }_0^\pm \mathrm{}𝑑zf(z)|z|^{\frac{3\tau \sigma }{\sigma }}\right|.$$
(25)
The second moment is related to the isothermal compressibility by the temperature and density of the system.
The derivation of the $`\gamma `$-power law demonstrates one way to arrive at the scaling relations between the critical exponents. In addition it illustrates the existence of only two independent exponents and shows the relation of the moments of the cluster distribution to familiar thermodynamic quantities. Fisher’s framework here illustrated and tempered by percolation theory will be used in the analysis of the cluster distributions of the three systems discussed above. It will be seen that in the case of systems which exhibit a continuous phase transition, the framework of Fisher is well followed, while for systems with no such phase transition, the framework fails, as it should.
## IV Phase transition signatures in cluster distributions
### A Insensitive signatures
In this section the insensitive features of the cluster distribution for each system are examined. It will be demonstrated that on this level of analysis each system exhibits behaviors that are consistent with systems which undergo a continuous phase transition. The conclusion is inescapable that this sort of analysis can yield necessary, but not sufficient, signals and no further insight to the mechanism behind multifragmentation. A deeper analysis will be necessary to distinguish those systems which undergo such a phase transition from those which do not.
#### 1 Fluctuations
One of the most striking characteristics of systems undergoing continuous phase transitions is the occurrence of fluctuations that exist on all length scales in a small range of the control parameter. In fluid systems this was observed as critical opalescence, first noted by Andrews in the latter half of the $`19^{th}`$ century . Fluctuations in cluster size and the density of the system arise because of the disappearance of the latent heat at the critical point. This is illustrated in the FDM when the isothermal compressibility diverges at the critical point and small changes in pressure gives rise to great changes in the density. In the FDM as the volume and surface contribution to the free energy of cluster formation vanishes the power law dominates and clusters of all length scales are observed .
In a cluster distribution the most readily observed fluctuations are those in the size of the largest cluster. For each system the root mean square (RMS) fluctuations in the size of the largest cluster normalized to the size of the system, $`\mathrm{\Delta }\left(A_{max}/A_0\right)`$, have been calculated as a function of the system’s control parameter. This measure of the fluctuations in the cluster distribution was first studied by Campi for gold multifragmentation and percolation . Those results are replicated here for those two systems.
Figure 1a shows $`\mathrm{\Delta }\left(A_{max}/A_0\right)`$, as a function of $`p_l`$ for percolation. As expected for a system known to exhibit a continuous phase transition, the RMS fluctuations peak over a narrow range in the control parameter. The location of this peak provides a first estimate of the critical point; $`p_c=0.33\pm 0.01`$. See Table I.
Next the percolation lattice is examined using the multiplicity of clusters, $`m`$, as an estimate of the control parameter. This is done because in the case of nuclear multifragmentation $`m`$ is experimentally measurable. Figure 1b shows much the same qualitative behavior as Figure 1a. The fluctuations peak over some narrow range of $`m`$ and suggest the value of the multiplicity at the critical point, the critical multiplicity, to be $`m_c=55\pm 5`$.
For random partitions a peaking behavior in the fluctuations of the size of the largest cluster as a function of $`m`$ was observed, see Figure 1c. These fluctuations can be understood as follows. At $`m=1`$ there can be no fluctuations in the size of the largest cluster because of the dual constraints of event cluster multiplicity and the fixed number of constituents. As the multiplicity increases from one, the constraints ease and fluctuations in the size of the largest cluster grow. At the maximum possible multiplicity, i.e. when $`m`$ is equal to the total number of constituents, the size of the largest cluster is constrained to be equal to unity. Thus, the fluctuations show a peak, but for a reason that has nothing to do with a continuous phase transition. Therefore it must be concluded that the observation of a maximum in the fluctuations in the size of the largest cluster is not sufficient to distinguish systems with and without critical behavior. On the other hand, the absence of a peak in fluctuations would indicate that the clusters of the system were not produced near a critical point. If the system’s phase space has been fully explored, then the stronger statement that the system does not possess a critical point could be made. At this level of analysis the critical multiplicity of this system can be estimated to be $`m_c=5\pm 2`$.
Finally, Figure 1d shows the Au $`+`$ C multifragmentation data with the cluster distribution normalized to the size of the system, $`A_0(m)`$. The fluctuations in the mass of the largest cluster exhibit a peak when plotted as a function of the event total charged particle multiplicity, $`m`$. This behavior is consistent with what is expected for a critical phenomenon. However, as illustrated above, it is far from conclusive. At this level of analysis the estimate for the critical multiplicity is $`m_c=31\pm 6`$.
It is also possible to study the fluctuations in the average size of a cluster. From the example of critical opalescence it is clear that the greatest fluctuations in cluster size should occur at the critical point. To that end the quantity known as $`\gamma _2`$ is constructed again following the work of Campi \- . The variance in the mean cluster size, $`A_f`$, is defined as:
$$\sigma ^2=\underset{N\mathrm{}}{lim}(\frac{1}{N}A_f^2)A_f^2.$$
(26)
The average cluster size is given by the ratio of the first moment to the zeroth moment:
$$A_f=n_{A_f}A_f/n_{A_f}=M_1/M_0.$$
(27)
The first term in eq. (26) is just the ratio of the second moment to the zeroth moment. Therefore, the variance in the average cluster size can be written in terms of the $`k^{th}`$-moments:
$$\sigma ^2=\frac{M_2}{M_0}(\frac{M_1}{M_0})^2.$$
(28)
This quantity is directly related to Campi’s $`\gamma _2`$ via:
$$\gamma _2=\frac{\sigma ^2}{A^2}+1=\frac{M_2M_0}{M_1^2},$$
(29)
which is easily measured and was coined by Campi as the reduced variance .
In a later paper, , Campi discussed the differences in methods to measure $`\gamma _2`$. Specifically, the manner in which the $`k^{th}`$-moments are computed from the observed cluster distribution. One method is to measure the $`k^{th}`$-moments on an event by event basis and then compute an average based on the control parameter, e.g.:
$$M_k(ϵ)=\frac{1}{N}\underset{i=1}{\overset{N}{}}M_k^i(ϵ)=\frac{1}{N}\underset{i=1}{\overset{N}{}}(\underset{A_f}{}n_{A_f}^i(ϵ)A_f^k),$$
(30)
where $`N`$ is the number of events at a control parameter value of $`ϵ`$, and $`i`$ denotes the $`i^{th}`$ event. This method of calculation of the $`k^{th}`$-moments will be termed averaging the sums and will yield: $`\gamma _2`$.
The alternate method involves calculating an average cluster distribution at each value of the control parameter and then calculating the $`k^{th}`$-moments from the resulting average cluster distribution:
$$\overline{M}_k=n_{A_f}(ϵ)A_f^k=\underset{A_f}{}(\frac{1}{N}\underset{i=1}{\overset{N}{}}n_{A_f}^i(ϵ))A_f^k.$$
(31)
This method of calculation will be termed summing the averages and will give: $`\overline{\gamma }_2`$.
For quantities linear in $`n_{A_f}`$ there is no difference in the two methods so that $`M_k(ϵ)=\overline{M}_k(ϵ)`$. However, due to the dependence of $`\gamma _2`$ on the square of the first moment, there will be a difference in the two methods of calculation. Results for both methods for each system are shown in Figure 2.
Of primary significance is the presence of a peak in both measurements of $`\gamma _2`$ for all systems. For an infinite system exhibiting critical phenomena, the location of the peak in $`\gamma _2`$ will coincide with the location of the critical point. For the percolation system Figures 2a and 2b show that both the location and magnitude of the peak in $`\gamma _2`$ is dependent on the choice of calculation method. Solid lines indicate this measure of the critical point. For the random partitions, Figure 2c, the location of the peak in $`\gamma _2`$ shows no dependence on the method of calculation while the magnitude of the peak does. The gold multifragmentation data exhibit a dependence on the method of calculation both in the magnitude and location of a peak in $`\gamma _2`$.
Having noted the peaking behavior of $`\gamma _2`$, the significance of the amplitude of the peak is now addressed. It has been suggested that the height of the peak can be used to differentiate between the presence of a power law and that of an exponential: for a power law $`\gamma _2>2`$ while for an exponential $`\gamma _2<2`$. This is not definitive proof of the existence of a continuous phase transition as other systems show power laws in the absence of such a phase transition. All of the percolation figures show peaks above two, as do the multifragmentation data plots and the random partitions. However, the value of $`\gamma _2`$ depends on the size of the system in question . For a percolation system with 64 sites, peaks in $`\gamma _2`$ under two are observed, see Figure 3a and 3c. Therefore, the criterion $`\gamma _2>2`$ is not sufficient to discriminate between those finite systems which do and those which do not posses a power law cluster distribution.
Finally the question of the difference between the alternative methods of calculating $`\gamma _2`$ is examined via: $`\mathrm{\Delta }\gamma _2=\gamma _2\overline{\gamma }_2`$. It has been suggested that a peak in the difference could indicate critical phenomena and the location of the critical point . Unfortunately, the cause of this peak is not well understood and vanishes at the limits of the system size: $`(0,\mathrm{})`$. Figures 4a and 4b do show peaks in $`\mathrm{\Delta }\gamma _2`$ at some intermediate value of the control parameter for this percolation lattice of 216 sites. However, as the size of the percolation lattice increases this signal vanishes . For a percolation lattice with 64 sites Figures 3c and 3d, respectively, look like a cross between the percolation ($`L=6`$, $`m`$) results, Figures 2b and 4b, and the random partition results shown in Figures 2c and 4c. This is believed to be due to the twin constraints of the multiplicity and the conservation of constituents imposed upon the system at the extremes in cluster multiplicity. Similar behavior is observed in the gold multifragmentation data in Figures 2d and 4d.
Neither the $`\gamma _2`$ measure of fluctuations nor the observation of fluctuations in the size of the largest cluster provide definitive insight into the nature of the cluster producing mechanism. For both random partitions and percolation $`\gamma _2`$ peaks at nearly the same value of the control parameter regardless of the method of averaging used. For the percolation system the value of $`p_l`$ at the peak in $`\gamma _2`$ is close to the value of $`p_l`$ where $`\mathrm{\Delta }(A_{max}/A_0)`$ is a maximum. This coincidence does not hold for random partitions; compare Figure 1c and 4c. For both percolation ($`m`$) and multifragmentation, there is better agreement on the critical point from fluctuations and than from $`\gamma _2`$ is computed via eq. (30).
#### 2 Divergences
Another signature previously used to infer the existence of a continuous phase transition from cluster distributions is the observance of a peak in the second moment , . It has been pointed out that models with no phase transition can exhibit a peaking behavior in the second moment . Figure 5 shows the behavior of the second moment for each of the systems examined in this work. In this figure, for the sake of illustration, the largest cluster has been excluded from the sum at all values of the control parameter. Each system shows a peak at some intermediate value of its control parameter. Table I lists the location of the second moment peaks. It is clear from the peak observed for the random partitions that it is possible to observe a peak in the second moment for a non-critical cluster distribution. Thus this quantity cannot be used to distinguish between critical and non-critical systems.
An issue with the use of the second moment’s peaking behavior is the exclusion of the largest cluster from the sum in eq. (19). Again, in the FDM formalism the sum runs over all clusters in the gas. On the liquid side of the critical point a gas exists in addition to a liquid drop. Thus, the largest cluster represents the bulk liquid. On the gas side of the critical point there is no liquid drop and the largest cluster is merely the largest gas particle. With this understanding it is clear that the largest cluster should be omitted from the summation in the second moment only in the liquid region, whereas the summation should run over all clusters in the gas region. For a proper construction of the second moment, knowledge of the location of the critical point is required. In the thermodynamic limit of infinite system size, exclusion of the largest cluster makes little difference. However in small systems the proper construction of the second moment is crucial if critical behavior is to be observed in ref. . See ref. for an example of the improper construction of the second moment.
#### 3 Campi plots
Plots of the natural log of the normalized size of the largest cluster, $`\mathrm{ln}(A_{max}/A_0)`$, versus the natural log of the second moment, $`\mathrm{ln}(M_2)`$, were first presented by Campi in a comparison of gold multifragmentation and percolation . Figure 6 shows the resulting plots for each of the systems discussed in this paper. In each plot there is a liquid leg for the largest $`A_{max}`$ and small $`M_2`$ and a gas leg for smaller $`A_{max}`$ and mid-range values of $`M_2`$. That similar behavior is observed for all systems is in indication that this is a necessary, but not sufficient, signal for critical behavior.
#### 4 Rise and fall of intermediate mass fragments
In many nuclear multifragmentation studies the term intermediate mass fragment, IMF, has been defined as a cluster which has a charge between $`3Z_f30`$. For the percolation system presented here a charge has been assigned to each cluster by multiplying the number of constituents in the cluster by the charge to mass ratio of a gold nucleus. For the random partitions the number of constituents is used as the charge. Since the definition of an IMF is arbitrary it makes little qualitative difference what range in some measure of the cluster size is used.
Aside from the equilibrium arguments made by the ALIDIN group -, little insight towards the presence or absence of a continuous phase transition is gained from a simple plot of the average number of IMF’s, $`M_{imf}`$ versus the control parameter. Figure 7 shows the results for the systems discussed in this work. Each system shows a peak in $`M_{imf}`$ at some intermediate value of the control parameter. Comparing the peak position in Figure 7 to the values listed in Table I shows that there is little correspondence between the numerous proposed methods for locating the critical point. The arbitrary nature of the definition of an IMF makes it unlikely that the peak in $`M_{imf}`$ occurs at the critical point. To some degree the rise and fall feature is due to the constraint of a fixed number of constituents. It it obvious that a the extreme values of the control parameter, the number of IMFs must diminish, while at intermediate values, it must be at least as great. Thus, the occurrence of a peak at some intermediate value of the control parameter is expected.
#### 5 $`\tau _{eff}`$-minimum
With the first observation of a power law in the nuclear multifragmentation yield distribution , it became a common analysis tool to fit cluster distributions to a power laws and extract exponent values. In an effort to make this a more quantitative analysis the value of the extracted exponent, $`\tau _{eff}`$, was examined as a function of some control parameter that was experimentally or numerically accessible. It was assumed that at the critical point the value of $`\tau _{eff}`$ should attain a lower value than fits which were performed away from the critical point -. The logic of this assumption was based upon the idea that at low temperatures a system has few small clusters, so the power law should be steep, leading to a high $`\tau _{eff}`$ value. At high temperatures there are many small clusters and little else, which is reflected in a high value of $`\tau _{eff}`$ and a steep power law. At the critical point clusters on all length scales appear and the power law is shallow with a lower value of $`\tau _{eff}`$. In this analysis the largest cluster was generally omitted from the fitting procedure and both the constant of proportionality and $`\tau _{eff}`$ were allowed to vary independently. Many investigations of nuclear multifragmentation, both theoretical and experimental, employed this method of analysis -.
There are two flaws in this analysis method. The first is the use of a two parameter fit for the power law. Allowing both the overall normalization of the power law and the exponent to vary independently is in conflict to the scaling assumptions underlying the FDM as shown in eq.’s (9) and (10). A proper fit for a power law within the context of the FDM should be based on single parameter. Furthermore, the cluster distribution must be normalized to the size of the system as was outlined in III. Without this normalization, which requires knowledge of the system’s size, power law fits lose much of their ability to contribute useful information to the presence of critical phenomena.
Leaving aside for a moment that the execution of the $`\tau _{eff}`$-minimum analysis violates the scaling assumptions of the FDM, the signal of a minimum in the cluster yield power law will be examined. A two parameter fit for $`\tau _{eff}`$ searches for the minimum in an effective exponent which is defined as -:
$$\tau _{eff}=\frac{\mathrm{ln}n_{A_f}(ϵ)}{\mathrm{ln}A_f}.$$
(32)
If it is assumed that the system under study follows a power law in the cluster yield at the critical point, and away from the critical point the cluster yield is affected by a scaling function such as in eq. (8), then:
$$\tau _{eff}=\tau A_f\frac{\mathrm{ln}f}{A_f}.$$
(33)
The minimum in $`\tau _{eff}`$ can be found by differentiating eq. (33):
$$\frac{d\tau _{eff}}{dϵ}=A_f\frac{}{ϵ}\frac{\mathrm{ln}f}{A_f}=0.$$
(34)
This indicates that the location of the minimum in $`\tau _{eff}`$ is dependent on the form of the scaling function, $`f`$. Assuming the scaling function has the form of eq.(5) then the minimum in $`\tau _{eff}`$ will be at $`ϵ=B/2A_f^\sigma `$ , and not at the critical point $`ϵ_c=0`$.
Despite the flaws in the $`\tau _{eff}`$-minimum analysis it is of interest to examine the results for the systems discussed in this paper. Figures 8 through 11 show the results for a two parameter fit to the cluster distribution for percolation (probability and multiplicity), random partitions and gold multifragmentation, respectively. For all systems, the cluster distributions were fit at each value of the control parameter. Only clusters with $`0.02(A_f/A_0)0.22`$ were included in the fits. The first three systems weighted $`\chi _\nu ^2`$ with errors associated with $`n_{A_f}(ϵ)`$ while the $`\chi _\nu ^2`$ for the gold multifragmentation cluster distributions were weighted with errors on both $`n_{A_f}(ϵ)`$ and $`A_f`$.
For the percolation ($`p_l`$) a minimum in $`\tau _{eff}`$ was observed at $`p_l=0.3`$ with $`\chi _\nu ^2=2.3`$, $`q_0=0.214\pm 0.005`$ and $`\tau =2.19\pm 0.01`$; shown in Figure 8a, b and c by the dotted lines. However, at $`p_l=0.33`$ the $`\chi _\nu ^2=1.02`$ , $`q_0=0.181\pm 0.003`$ and $`\tau =2.27\pm 0.01`$; shown in Figure 8a, b and c with the dashed lines. Based on a goodness of fit comparison, the latter value of $`p_l`$ is a better choice for the critical point as the cluster distribution is better fit by a power law. This result is in agreement with the analytic discussion of $`\tau _{eff}`$ above, namely that a minimum in $`\tau _{eff}`$ is a poor indicator of the critical point. If the results for $`p_l=0.33`$ are compared to the center of the $`\tau _{eff}`$, $`p_l0.28`$, the differences in the $`\chi _\nu ^2`$ and $`q_0`$ results are even more striking.
Similar results were seen for percolation ($`m`$), see Figure 9. Here the minimum in the $`\tau _{eff}`$-well yielded worse results for both $`\chi _\nu ^2`$ and $`q_0`$ than does the choice of the critical point based on a choice from the $`\chi _\nu ^21`$ region where there is good agreement between the fitted $`q_0`$ and the value computed using eq. (10) and the canonical $`\tau `$ value for three dimensional percolation.
Significant differences between percolation and random partitions are observed in this analysis as seen in Figure 10. The solid lines show the $`\tau _{eff}`$ and $`q_0`$ values for systems in the three dimensional Ising universality class, while the dashed line shows the $`\tau _{eff}`$ and $`q_0`$ for three dimensional percolation. The first noticeable difference is a lack of a valley shape in the plot of $`\tau _{eff}`$ versus control parameter, see Figure 10b. The value of $`\tau _{eff}`$ is below $`2.2`$ for all but $`m>60`$. Next is the lack of a region in $`m`$ where $`\chi _\nu ^21`$ (other than at $`m=2`$), Figure 10a. All fits yield large $`\chi _\nu ^2`$ values indicating poor fits to the cluster distribution by a power law for the range of clusters examined.
The gold multifragmentation data show results similar to those of percolation ($`m`$). Here the cluster size is measured in terms of the nucleon number and the cluster distribution is normalized to the mass of the gold projectile remnant. Figure 11b shows a valley in $`\tau _{eff}`$ as a function of $`m`$, albeit one with a shallow and questionable upwards slope at high $`m`$. Figure 11c shows fitted values of $`q_0`$ that coincide with canonical values. Figure 11a shows a region of low $`\chi _\nu ^2`$ values followed by steadily increasing values. If no knowledge of the $`q_0`$ and $`\tau `$ values is assumed, then this analysis shows no definitive signals. The $`\tau _{eff}`$-valley shows a broad minimum in $`\chi _\nu ^2`$ thus no one value of $`m`$ can be selected for the critical point based on goodness of fit arguments. At best one could argue for the neighborhood of the critical point and a value of $`q_0`$ and $`\tau `$ in some broad range.
#### 6 Conclusion
The analysis presented in this section shows that many proposed indicators of critical behavior are inconclusive. All of the considered systems show similar signals which are qualitative in nature and open to interpretation. It is therefore impossible, based solely on this level of analysis, to make a definitive conclusion as to the presence of a continuous phase transition in any of these systems. What is needed is an analysis or set of analyses that more clearly differentiates between systems with and without critical behavior.
### B Sensitive signatures
#### 1 The Fisher $`\tau `$-power law and the critical point
In this section the cluster yields are fit to a power law in a manner consistent with the FDM formalism. As with the two parameter fits the yields for clusters with $`0.02(A_f/A_0)0.22`$ were fit at each value of the control parameter. However, only a single parameter, the value of $`\tau `$, was allowed to vary to minimize the $`\chi _\nu ^2`$ of the fit. The value of the normalization, $`q_0`$, was determined via the Riemann $`\zeta `$-function in eq. (10). As suggested by Fisher , the value of $`\tau `$ was constrained to be between 2 and 3 so that the sum in the $`\zeta `$-function converges.
If the cluster distribution is well described by the FDM, then at the critical point the fit to a single parameter power law should show a minimum in $`\chi _\nu ^2`$. Away from the critical point the power law is modified by a scaling function with a form similar to that given in eq. (5). Therefore, fits to a single parameter power law should become worse as the modification from the scaling function increases away from the critical point.
Figure 12 shows the results for the percolation system with $`p_l`$ as the control parameter. In Figure 12a, a minimum in $`\chi _\nu ^2`$ is observed for fits in the mid $`p_l`$ range. This minimum indicates the location in $`p_l`$ of the cluster yield distribution which is best fit by a single parameter power law as suggested by the FDM formalism. By this estimation the critical point for this 216 site percolation lattice is at $`p_c=0.31\pm 0.05`$ with $`\tau =2.2\pm 0.1`$, $`q_0=0.20\pm 0.01`$ and a $`\chi _\nu ^2=1.62`$. The canonical values of $`\tau `$ and $`q_0`$ are not extracted due in part to unavoidable finite size effects, and to the binning of cluster yields together over a range of 0.01 in $`p_l`$, which causes the true cluster distribution at the critical point to be contaminated by distributions at other values of the control parameter. In spite of these difficulties, the signature of the critical behavior suggested by the FDM formalism is unmistakable. The location of the critical point determined here is consistent, at the 10$`\%`$ level, with the insensitive signatures presented in the previous section. See Table I. Figure 16a shows the best fit power law. For the percolation system clusters consisting of a single site are excluded from the fitting procedure. It is accepted that those clusters reflect the effects of the finite size of the system to a higher degree than larger clusters. Clusters with $`A_f53`$ were included in the fit. The largest cluster from each event was excluded from consideration when generating the average cluster distribution in keeping with the FDM formalism. Figure 16a shows the data for the entire cluster distribution in open circles. It is clear from this figure that the majority of the cluster distribution was used in the power law fit and further, that the exclusion or inclusion of the larger clusters has almost no effect on the results of this procedure. In short, the extracted parameters, namely $`\tau `$, $`q_0`$ and $`p_c`$, do not depend on the fit range.
Figure 13 shows the results of the single parameter fit analysis when applied to the same percolation system but with the cluster multiplicity used as a measure of the control parameter. Again there is a minimum in the $`\chi _\nu ^2`$ values at some intermediate value of the control parameter which indicates that $`m_c=57\pm 3`$ with $`\tau =2.2\pm 0.1`$, $`q_0=0.20\pm 0.01`$ and a $`\chi _\nu ^2=0.72`$. Note the consistency between these values of $`q_0`$ and $`\tau `$ and those obtained with $`p_l`$ following this method. The location of the critical point determined here is consistent, at the 10$`\%`$ level, with the insensitive signatures presented in the previous section. The lower $`\chi _\nu ^2`$ value is due to the finer bins over which the cluster distributions were grouped. Figure 16b shows the best fit power law. Here only clusters of size $`A=1`$ and size $`A=A_{max}`$ were excluded from the fitting procedure.
From Figures 12 and 13 it could be argued, based on the best agreement between the fitted $`\tau `$ and the accepted three dimensional percolation value, that there are better choices for the critical point than those quoted above. However, those arguments assume knowledge of the value of $`\tau `$ as an input. The use of the location of the best fit to a single parameter power law as an indicator of the critical point makes no assumption regarding the value of $`\tau `$ and is a test of the FDM formalism in which only the range of $`\tau `$ is suggested: $`2<\tau <3`$. The values of $`\tau `$ and $`q_0`$ are outputs rather than inputs of this analysis. Much of the following analysis presented in this paper follows the same philosophy. That is, the analysis is designed to test the cluster distribution in question for behavior consistent with the FDM formalism. The values of quantities, such as critical exponents, are results of the analysis method and are in no way selected for on the basis of their particular values. Agreement between exponent values determined by this procedure and the canonical values in various universality classes is then significant because the values of the exponents are determined solely by the behavior of the cluster distributions so analyzed.
The results of the single parameter power law fits for the random partitions are presented in Figure 14. There is a minimum in the $`\chi _\nu ^2`$ value at $`m=59`$. However $`\chi _\nu ^2=10.83`$, which is an order of magnitude above the percolation results, should not be used as an indication of a good fit of the cluster distribution by a single parameter power law. The location of the $`\chi _\nu ^2`$ minimum is also in disagreement with the insensitive signatures presented in the last section. Here only clusters of size $`A_f=1`$ and size $`A_f=A_{max}`$ were excluded from the fitting procedure.
Figure 15 shows the results of this analysis applied to the gold multifragmentation data. As with the percolation results, the $`\chi _\nu ^2`$ shows a minimum that drops nearly two orders of magnitude from the peaks for high and low $`m`$ to the valley at a mid range value of $`m`$, see Figure 15a. In the context of the FDM analysis this result suggests that the critical point is located at $`m_c=22\pm 1`$ with $`\tau =2.2\pm 0.1`$ and $`q_o=0.18\pm 0.01`$ and $`\chi _\nu ^2=2.70`$. The best fit power law is show in Figure 16d. An uncertainty of one unit of multiplicity is assigned to $`m_c`$ to take into account the relatively low values of $`\chi _\nu ^2`$ of the neighboring fits.
For the above fits to the gold multifragmentation data the $`\chi _\nu ^2`$ is weighted by the errors in both $`n_{A_f}`$ and $`A_f`$. The fitting procedure has also been performed with no error weighting on $`\chi _\nu ^2`$ and with errors only in $`n_{A_f}`$ for weighting. Both analyses shows results that were not significantly different from those quoted here. As mentioned previously, clusters with $`Z_f=2`$ are created in both the prompt first stage and in the multifragmenting of the gold nuclear remnant. The prompt $`Z_f=2`$ clusters have been excluded from the filtered gold multifragmentation analysis. However, as a further test of the single parameter power law fit, only clusters with $`3Z_f16`$, i.e. clusters with no contamination from the prompt first stage, were included in a repeat of this analysis. Again the results show practically the same behavior as those shown here. As yet another test, clusters with $`2Z_f<Z_{max}`$ were included in the fitting procedure, and again the results showed no difference from those presented here. Finally clusters with $`3Z_f<Z_{max}`$ were included in the fitting procedure, and again the results showed no difference from those presented here. The data always showed a deep valley in the $`\chi _\nu ^2`$ versus $`m`$ plot which indicated that the location of the critical point was $`m_c22`$ and that $`\tau 2.2`$, $`q_00.18`$ and $`1<\chi _\nu ^2<4`$. This analysis shows that the value of $`\tau `$ and the location of $`m_c`$ are not sensitive to the fit region. The behavior of the data show this clearly. See open circles in Figure 16d.
The single parameter power law analysis of the cluster distributions of the various systems produced the first result which can differentiate between systems that follow the FDM formalism and systems that do not. The differences between Figures 12a, 13a, 15a and Figure 14a are clear and distinct. For both percolation and gold multifragmentation the behavior of $`\chi _\nu ^2`$ is exactly what is predicted by the FDM formalism for continuous phase transitions. Far from the critical point the cluster distribution is fit poorly by a single parameter power law due to the influence of a scaling function where volume and surface effects overwhelm the underlying power law. At the critical point, where the influence of the scaling function vanishes, the cluster distributions are well described by a single parameter power law with an exponent value, $`\tau 2.2`$ and thus $`q_00.2`$, in keeping with what is expected for many universality classes. This fitting procedure does not merely search out a cluster distribution which is well fit by a power law, but finds the cluster distribution which is well fit by the FDM formalism. This is achieved via the coupling between the exponent $`\tau `$ and the normalization factor $`q_0`$. See eq. (10). The random partitions fail to produce such signals. This is expected as that system does not obey the FDM formalism and thus should not show the same signals as systems that are known to follow the FDM such as percolation. This analysis of the cluster yield of gold multifragmentation yields a signal that is suggestive of critical phenomena.
#### 2 The critical exponent $`\sigma `$
In section III it was shown that in the context of the FDM the surface of a cluster makes a contribution to the free energy of cluster formation, via the scaling function $`f(z)`$, that depends on the number of constituents of the cluster raised to the power $`\sigma `$. See eqn’s (2), (5) and (6). The behavior of the order parameter suggests that the scaling function $`f(z)`$ has a maximum . At the maximum of the scaling function, $`f_{max}(z_{max})`$, the production of $`A_f`$ sized clusters is greatest:
$$n_{A_f}^{max}(ϵ_{max})=q_0A_f^\tau f(z_{max}).$$
(35)
The argument of $`f_{max}`$ is:
$$z_{max}=A_f^\sigma ϵ_{max},$$
(36)
where the value of $`z_{max}`$ depends on the specific details of the system in question . Rearranging eq. (36) yields:
$$ϵ_{max}=z_{max}A_f^\sigma .$$
(37)
Thus $`ϵ_{max}`$, the value of the control parameter at which the greatest number of clusters of size $`A_f`$ are produced, is related to the cluster size through a simple power law with exponent $`\sigma `$. The exponent $`\sigma `$ can then be determined from knowledge of the location of the critical point and the value of the control parameter at the greatest production of clusters of size $`A_f`$.
The location of the critical point was determined in the search for the Fisher $`\tau `$-power law in section IV-B-1 and will be used here to determine the $`\sigma `$. The value of the control parameter which yields the greatest production of each $`A_f`$ cluster size was determined from the peak location in a plot of $`n_{A_f}(ϵ)`$ versus the system’s control parameter. See Figure 17.
For each system at each cluster size plots such as those shown in Figure 17 were used to determine the location of the peak of $`n_{A_f}(ϵ)`$. For example, in percolation ($`p_l`$), $`ϵ_{max}=(p_cp_{max})/p_c`$, pairs of points $`(n_{A_f}(ϵ),p_l)`$ for a particular $`A_f`$ were fed into a SPLINE routine . Input pairs were then smeared by assigning $`\delta n_{A_f}(ϵ)`$ as the standard deviation of a gaussian centered on $`n_{A_f}(ϵ)`$. Output of the SPLINE routine was used to interpolate the behavior of a smooth curve between the pairs of input points. Stepping along the interpolations in increments much smaller than the separation of the input $`p_l`$, a maximum of $`n_{A_f}(ϵ)`$ was determined and $`p_{max}`$ was recorded. This process was repeated thousands of times for each cluster size and lead to an estimate of $`p_{max}\pm \delta p_{max}`$ as a function of $`A_f`$.
Using $`p_{max}(A_f)\pm \delta p_{max}(A_f)`$ and the value of $`p_c\pm \delta p_c`$, from the Fisher $`\tau `$-power law determination process, the value of the exponent $`\sigma `$ was determined by taking the slope of $`\mathrm{ln}(ϵ_{max})`$ versus $`\mathrm{ln}(A_f)`$. The value of $`z_{max}`$ was determined by exponentiating the offset. The value of $`p_c`$ was varied uniformly throughout the range suggested by $`\delta p_c`$ and tens of fits were made with varying starting and ending points in $`A_f`$ of the fitting region. The final value of $`\sigma \pm \delta \sigma `$ and $`z_{max}\pm \delta z_{max}`$ are the average and RMS values resulting from all the fits.
Results of this analysis performed on percolation ($`p_l`$) are shown in Figure 18a. Here the value of the control parameter that coincides with the maximum in production of clusters of size $`A_f`$, $`ϵ_{max}`$, is plotted against the cluster size. Results of the average power law fits to eq. (37) are plotted as a solid line. The agreement between the values returned by the procedure discussed above, $`\sigma =0.52\pm 0.02`$ and $`z_{max}=0.89\pm 0.03`$, and the accepted values for three dimensional percolation, $`\sigma =0.45`$ and $`z_{max}=0.8`$ , establishes the reliability of this exponent extraction method. The analysis differs in method from previous efforts on percolation lattices , but not in result.
The next test of this analysis is to extract a value of $`\sigma `$ from percolation ($`m`$). In order for this procedure to be useful on multifragmentation data it must be shown that the exponent $`\sigma `$ can be determined using cluster multiplicity as the control parameter. To that end the multiplicity at which the production of each cluster size is maximal, $`m_{max}`$ was determined via the procedure described previously. Using the value of $`m_c\pm \delta m_c`$ determined via searching for the Fisher $`\tau `$-power law and $`m_{max}`$ the exponent $`\sigma `$ was determined by taking the slope of $`\mathrm{ln}(ϵ_{max})`$ versus $`\mathrm{ln}(A_f)`$. The value of $`m_c`$ was varied uniformly throughout the range suggested by $`\delta m_c`$ and several plots were made with varying starting and ending points in $`A_f`$ of the fitting region. The value of $`z_{max}`$ was determined by exponentiating the offset. Results of the average power law fits to eq. (37) are plotted as a solid line in Figure 18b. The agreement between the values returned by this procedure, $`\sigma =0.52\pm 0.02`$, the value for $`\sigma `$ quoted in the above paragraph and the accepted values for three dimensional percolation again establishes the reliability of this exponent extraction method and shows that the use of $`m`$ as a control parameter is acceptable.
The value of $`z_{max}=2.4\pm 0.1`$ extracted for percolation as a function of multiplicity is different from the value quoted above, $`z_{max}=0.89\pm 0.03`$, for the percolation system as a function of probability. This is a result of changing the measure of the control parameter from probability to multiplicity. This difference was observed in previous percolation efforts and explained therein. A plot of $`ϵ(p_l)`$ against $`ϵ(m)`$ show that $`z_{max}(p_l)`$ and $`z_{max}(m)`$ map to each other. See Figure 9 of ref. .
Clusters from the random distribution were also subjected to this analysis. Due to the failure of the search for the Fisher $`\tau `$-power law in the random partitions, the value of $`m_c`$ determined in the analysis of the gold multifragmentation was used, $`m_c=22\pm 1`$. The value of the cluster multiplicity for maximum production of $`A_f`$ sized clusters was determined in the same manner as with the percolation system. The flatness of the $`n_{A_f}(ϵ)`$ versus $`m`$ curves, see Figure 17c, makes finding a unique value of $`m_{max}`$ impossible. This is reflected in the large error bars on $`ϵ_{max}\pm \delta ϵ_{max}`$ seen when plotted against $`A_f`$ in Figure 18c. The value of $`m_{max}`$ reported by the peak finding procedure employed here reflects, approximately, the mid point of the multiplicity range of $`n_{A_f}(ϵ)`$ for a particular $`A_f`$. Coupling the $`m_c`$ from the filtered gold multifragmentation data with the $`m_{max}`$ and fitting $`\mathrm{ln}(ϵ_{max})`$ versus $`\mathrm{ln}(A_f)`$ gave $`\sigma =0.4\pm 0.2`$ and $`z_{max}=2.0\pm 0.8`$. However, it is clear when comparing the resulting average fit for the random partitions shown in Figure 18c with either of the percolation results shown in Figure 18a and 18b that the $`\sigma `$ resulting for the random partitions cluster distribution is meaningless. This is to be expected as the framework of the FDM, used in the extraction of the exponent $`\sigma `$, is meaningful only when applied to systems which undergo a continuous phase transition. The failure of this analysis on this system is expected based on the basis of the failure of the analysis in the preceding section that aimed to find the Fisher $`\tau `$-power law and the critical point.
Results for the extraction of $`\sigma `$ from the gold multifragmentation data have been published previously , . In those analyses the largest cluster was excluded from consideration at every value of the control parameter. This is at odds with the formalism of the FDM where the sums excludes the largest cluster for $`ϵ>0`$ and include the largest cluster for $`ϵ<0`$.
The previous analyses yielded values of $`\sigma =0.68\pm 0.05`$ and $`0.65\pm 0.06`$ for work with the un-normalized charge distribution and normalized mass distribution respectively. When this analysis was redone using formalism of the FDM, i.e. the largest cluster excluded on one side of the critical point (liquid) and included on the other side (gas), the values of $`\sigma `$ were reduced by approximately $`50`$%: $`\sigma =0.32\pm 0.02`$. In the case of percolation the difference introduced in the value of $`\sigma `$ when following the FDM formalism (as was done above) or not (as was the case in ref. ) is on the order of a few percent. This is the first notable difference observed in the qualitative behaviors of percolation cluster distributions and gold multifragmentation cluster distributions.
One source of this differing behavior is the changing mass of the system. For gold multifragmentation, from $`A_0194`$ at low $`m`$ to $`A_092`$ at high $`m`$ , while the system size is constant for percolation. For gold multifragmentation effects of the finite size of the system are felt more at high multiplicities than low. Since the percolation system size is constant, finite size effects are felt more evenly.
It is the higher values of $`m`$ where cluster production peaks in multifragmentation. The finite size of the system limits the size to which a cluster can grow. Thus the number of clusters of size $`A_f`$, $`n_{A_f}`$, as a function of $`m`$ is contaminated when the largest cluster, $`A_{max}`$, is included in a plot of $`n_{A_f}`$ versus $`A_f`$ because $`A_{max}`$ would like to be larger, but finite size effects limit the size $`A_{max}`$ can attain. Therefore, one method to account for this effect is to exclude $`A_{max}`$ from the cluster distribution at large $`m`$ values where these effects are largest This was done for the gold multifragmentation data.
The multiplicity at which the production of each cluster size is maximal, $`m_{max}`$ was determined via the procedure described previously. The value of $`m_c`$ determined in the Fisher $`\tau `$-power law analysis was used, $`m_c=22\pm 1`$. The value of $`m_c`$ was varied uniformly throughout the range suggested by $`\delta m_c`$ and several plots were made with varying starting and ending points in $`A_f`$ of the fitting region. The exponent $`\sigma `$ was determined by taking the slope of $`\mathrm{ln}(ϵ_{max})`$ versus $`\mathrm{ln}(A_f)`$ and the value of $`z_{max}`$ was determined by exponentiating the offset. The results were $`\sigma =0.64\pm 0.05`$ and $`z_{max}=6.0\pm 0.8`$, see Table II and III. The average power law fits are shown in Figure 18d.
#### 3 The scaling function $`f(z)`$
With the critical point ($`p_c`$ or $`m_c`$), $`\tau `$, $`q_0`$ and $`\sigma `$ determined and assuming coexistence, $`g=1`$, it is possible to find the scaling function by rewriting eq. (7) as
$$n_{A_f}(ϵ)/q_0A_f^\tau =f(z).$$
(38)
Doing this has the effect of appropriately scaling $`n_{A_f}(ϵ)`$ and collapsing the data onto a single curve. Figure 19 shows the results of this sort of scaling.
In Figure 19a, b and d, the data from percolation ($`p_l`$ and $`m`$) and multifragmentation, respectively, show collapse onto a single curve for a wide range in cluster size and over nearly the full range in control parameter. Random partitions shows no such collapse, see Figure 19c.
As a demonstration of this type of scaling the same data has been scaled in the same fashion, but with a different choice of the critical point. Figure 20 shows the systems using a critical point with a value of half of the critical point determined via the Fisher $`\tau `$-power law, while Figure 21 shows the same analysis with a value of twice the critical point determined via the Fisher $`\tau `$-power law. A visual inspection of Figures 19, 20 and 21 reveals the greatest data collapse occurs when the choice of the Fisher $`\tau `$-power law critical point is used, at least for the percolation ($`p_l`$ and $`m`$) and multifragmentation systems. Random partitions show no such collapse. Using different values of $`\tau `$ and $`\sigma `$ in this scaling analysis of random partitions does not significantly alter the data collapse. In Figures 19, 20 and 21 error bars on the data points are not shown for the sake of clarity. The size of the error bars reflect the scatter of the data and are larger for larger negative values of $`z`$ since there were lower statistics for higher multiplicity events .
Figure 22 shows a quantitative measure of the data collapse from this scaling analysis. For a number of different choices of control parameter scaling plots, as in Figures 19 through 21, were made. Each plot was binned along the abscissa and the RMS fluctuations for each bin were calculated. The RMS fluctuations in all bins were then summed and plotted as a function of the choice of critical point. See Figure 22. In the percolation ($`p_l`$ and $`m`$) and multifragmentation systems the data shows the most collapse in the neighborhood of the Fisher $`\tau `$-power law critical point. No such behavior is observed in the random partition system. This analysis serves as another, albeit crude, estimate of the location of the critical point. Table I lists the results.
The scaled data were used to determine the functional form of the scaling function by fitting the data with an empirical parameterization consisting of two gaussians instead of the single gaussian in eq (5):
$`f(z)`$ $`=`$ $`a_1\mathrm{exp}[{\displaystyle \frac{1}{2}}({\displaystyle \frac{zb_1}{c_1}})^2]+`$ (40)
$`a_2\mathrm{exp}[{\displaystyle \frac{1}{2}}({\displaystyle \frac{zb_2}{c_2}})^2].`$
This was suggested by the asymmetry of the percolation ($`p_l`$) data, Figure 19a, and is consistent with a simplified version of corrections to scaling as discussed in section V. Figure 19 shows the resulting fits for all systems. Fit parameter values can be found in Table IV. Errors on the parameters of the fits, e.g. $`a_1`$ etc., reflect the change in those parameters when the range of clusters included was changed, e.g. clusters with $`Z_f=2`$ were included or excluded and so on, and the weighting on the fit was changed, e.g. $`\chi _\nu ^2`$ is unweighted, weighted with errors on $`n_{A_f}(ϵ)/q_0A_f^\tau `$ or with errors on $`n_{A_f}(ϵ)/q_0A_f^\tau `$ and $`ϵ`$.
The scaling function for percolation ($`p_l`$ and $`m`$) determined here is the scaling function for percolation in three dimensions, i.e. it is universal for three dimensional percolation independent of size. The scaling functions for $`p_l`$ and $`m`$ determined above agree well with the scaled cluster distributions of different size lattices, see Figure 23, and can be used to predict the behavior of the second moment for any size lattice . In the same spirit, the scaling function determined here for gold multifragmentation is the scaling function for charged nuclear matter which describes the cluster distributions produced in the multifragmentation of any nucleus, not just the excited gold remnant discussed in this work. With the knowledge of the form of the scaling function various other quantities can be determined as illustrated in section III and shown below.
The cluster distributions for the random partitions is fit, by eye, with the same empirical parameterization as in eq. (40) see Figure 19c. The random partitions cannot be described by eq. (40). The solid curve in Figure 19c will be used in the following section to demonstrate the failure of the scaling analysis, as is also seen here, when applied to a system where a continuous phase transition is absent.
Finally, a consistency check in this analysis is the agreement between the location of the peak in the scaling functions and the values of $`z_{max}`$ determined in the $`\sigma `$ analysis, see Table III.
#### 4 $`\gamma `$-power law from the scaling function
The behavior of $`\kappa _T`$ or $`M_2`$ can be derived from the functional form of the scaling function and the critical parameters via eq. (23). Performing the integration in eq. (23) using the functional form of the scaling function determined above yields a direct calculation of the critical amplitudes, $`\mathrm{\Gamma }_\pm `$ via eq. (25). The critical exponent $`\gamma `$ is calculated from the values of $`\tau `$ and $`\sigma `$ via a scaling relation in eq. (24). Combining these two, $`\mathrm{\Gamma }_\pm `$ and $`\gamma `$, it is possible to calculate the $`\gamma `$-power law that describes the behavior of the second moment. This calculated $`\gamma `$-power law can then be compared to the behavior of $`M_2`$ as measured from the cluster distribution. Figure 24a, b and d shows the agreement between the measured $`M_2`$ data (largest cluster omitted in the liquid region) and the calculated $`\gamma `$-power law curves for percolation ($`p_l`$ and $`m`$) and gold multifragmentation, respectively and Tables II and III list the results.
The values of $`\gamma `$ determined via the scaling relation in eq. (24) for percolation ($`p_l`$ and $`m`$) show approximate agreement with the accepted value of 1.8. The high value of $`\sigma `$ extracted above leads to a low value of $`\gamma `$ here. Figures 24a and b also show the behavior of the second moment of a $`250,047`$ site lattice. The power law predicted using the scaling function determined with a 216 site lattice shows rough agreement with the measured $`M_2`$ of the larger lattice in both the amplitude ($`\mathrm{\Gamma }_\pm `$) and exponent ($`\gamma `$). There is approximate agreement between the predicted power law and the measured $`M_2`$ of the smaller lattice over some region in $`ϵ`$ that is neither too near to, nor too far from the critical point, $`ϵ=0`$. It is this region that will be determined, independently, in the following section.
For the percolation ($`p_l`$ and $`m`$) system, the disagreement between the measured $`M_2`$ data and the calculated curves is due to two well known reasons: far from the critical point, the assumptions of scaling are no longer valid and the analytic background overwhelms the singular behavior. Near the critical point finite size effects dominate $`M_2`$, limiting the sizes of the large clusters which make the most significant contribution. In contrast, the $`\tau `$-power law was observed at the critical point because it is determined by smaller clusters which suffer the least from the finite size effects.
Figure 24c shows the results when this analysis was applied to random partitions. The power law predicted from the scaling function analysis applied to the cluster distribution of the random partitions fails to reflect the behavior of the measured second moment. This is not surprising as the random partitions presented here are not the result of a system undergoing a continuous phase transition. The disagreement observed in Figure 24c then serves as an indication of how this particular analysis probes for the presence of a continuous phase transition. This figure shows the results of this analysis for a system with no phase transition, while Figures 24a and b show the results of this analysis on a system where such a phase transition is present.
The results of this analysis when applied to nuclear multifragmentation are shown in Figure 24d. In this case, the comparison to the predicted $`\gamma `$-power law is neither as good as that for percolation nor as poor as that for the random partitions. It shall be shown in section VI that considerable improvement can be achieved if account is taken of the changing system size, $`A_0(m)`$, and finite size scaling effects.
The approximate agreement between the predicted $`\gamma `$-power law and the measured $`M_2`$ behavior is in keeping with the behavior expected for small systems undergoing a continuous phase transition, e.g. the percolation system. The multifragmentation results are clearly different that then results of a system without a continuous phase transition, e.g. random partitions.
#### 5 $`\gamma `$-matching
In the previous works the procedure for determining critical exponent values and the location of the critical point from the cluster distribution was based on a method of matching exponent values on both sides of the critical point , . The idea was to find the region on either side of the critical point where the power law behavior predicted by the scaling function holds. As is seen in Figure 24 there is some intermediate $`ϵ`$ region where the second moment data are described by a power law, a region where the $`M_2`$ behavior is dominated by the $`\gamma `$-power law and all other effects are small in comparison. In earlier percolation studies general guidelines based on the correlation length and size of the fluctuations were used to find the boundaries in $`ϵ`$ of the regions to be fit. In nuclear multifragmentation analyses it was impossible to use such guidelines. Instead a method was developed that searched for regions best fit by power laws and determined the location of the critical point and exponent values simultaneously. The values of the critical exponents and the normalizations associated with power laws were obtained from the best fit power laws in those regions. As with the previous analyses presented in this paper, this method of exponent matching does not select a particular value of a critical exponent or the critical point. Instead the values found are the outcome of an unbiased procedure.
The method is as follows. A choice of the critical point, $`p_c`$ or $`m_c`$ was made. From this choice plots such as those shown in Figure 24 were made. Then fitting boundaries in $`ϵ`$ were chosen. The fitting range was defined by $`ϵ_\pm ^{far}`$ and $`ϵ_\pm ^{near}`$. For example, on the gas side of the critical point a fit of $`ln(M_2)`$ versus $`ln(|ϵ|)`$ was made for all data with $`|ϵ_+^{near}||ϵ||ϵ_+^{far}|`$. The slope of the resulting linear fit was recorded as $`\gamma _+`$, the offset as $`ln(\mathrm{\Gamma }_+)`$ and the goodness of fit as $`\chi _{\nu +}^2`$. The same procedure was applied to the liquid side of the chosen critical point, recording $`\gamma _{}`$, $`ln(\mathrm{\Gamma }_{})`$ and $`\chi _\nu ^2`$. For each choice of the critical point, several choices of fitting regions, $`ϵ_\pm ^{far}`$ and $`ϵ_\pm ^{near}`$ were made and results recorded. Five parameters were chosen for each set of power law regions examined: $`ϵ_\pm ^{far}`$, $`ϵ_\pm ^{near}`$ and $`p_c`$ or $`m_c`$.
The $`\gamma `$-power law fit regions and critical point locations were evaluated by demanding that: (1) they yield $`\gamma _+`$ and $`\gamma _{}`$ values that matched each other to within the error bars on those values returned by the fitting routine and (2) that the $`\chi _\nu ^2`$ of the fits were in the lowest quarter of the distribution resulting from all the fits which satisfy condition (1). The results from the power law fit regions that passed these two criteria were then histogrammed and average values for all quantities concerned were determined. The results are summarized in Tables I, II and III and shown in Figure 25.
The lines plotted in Figure 25 do not result from any single fit, but display the average results for $`\gamma _\pm `$ and $`\mathrm{\Gamma }_\pm `$ that have satisfied conditions (1) and (2). The points in Figure 25 are the measured second moment for the particular cluster distribution in questions plotted against $`ϵ`$, which depends on the average value of $`p_c`$ or $`m_c`$ that satisfies conditions (1) and (2). Therefore the lines in Figure 25 should not be interpreted as a fit to the data points shown in the same figure, but as the average results from the $`\gamma `$-matching procedure. Full circles in Figure 25 show the average fitting regions that satisfy conditions (1) and (2).
For percolation $`p_l`$ the value of $`\gamma `$ determined in this manner is within a few percent of the value determined in and the infinite lattice value. The ratio of $`\mathrm{\Gamma }_+/\mathrm{\Gamma }_{}`$ determined by this method, a ratio that depends on the universality class of the system in question, is also in agreement with the infinite lattice value and the $`\mathrm{\Gamma }_\pm `$ values predicted by the scaling function, see Table III. The value of $`p_c`$ determined here is within 15$`\%`$ of the value determined in a previous analysis of the $`L=6`$ lattice and the value determined above in the Fisher $`\tau `$-power law analysis, see Table II and Figure 25a.
The results for the analysis of percolation with $`m`$ as a measure of the control parameter are worse that the results when the natural control parameter $`p_l`$ is used, the difference in $`\gamma _+`$ and $`\gamma _{}`$ was: $`\mathrm{\Delta }\gamma ^m=0.06\pm 0.1`$ compared to $`\mathrm{\Delta }\gamma ^{p_l}=0.0\pm 0.3`$. This is to be expected because for each value of $`p_l`$ there is some spread in the resulting values of $`m`$, so that binning in $`m`$ groups together events with different values of $`p_l`$. There is also a non-linear relation between the average values of $`p_l`$ and $`m`$ . In spite of these two effects the results of the analysis in section IV-B-4 suggests that vestiges of the signature of a phase transition are still present even when $`m`$ is used as the control parameter. That is also the case in the present analysis. Table II shows that the $`\gamma `$ value agrees, within error bars, with the infinite lattice value. The values of the critical amplitudes, $`\mathrm{\Gamma }_\pm `$, do not yield a ratio that agrees with the infinite lattice value. This is due to the non-linear mapping of $`p_l`$ onto $`m`$ and is discussed in .
When the $`\gamma `$-matching procedure was applied to the cluster distributions from random partitions a very limited amount of trial fits passed the combined tests of (1) and (2). The results compared poorly to the percolation results. At best the values of $`\gamma _+`$ and $`\gamma _{}`$ match to within 20$`\%`$ of the average value of $`\gamma `$, compared to perfect matching for percolation $`p`$ and matching within 5$`\%`$ for percolation $`m`$. The value of the critical point, $`m_c`$, returned from this analysis also compared poorly to other outcome of previous analyses, see Table I. Finally, while fit regions for all systems were limited, the fit regions are the smallest and the poorest of quality for the random partitions.
The results of the $`\gamma `$-matching analysis applied to multifragmentation data has been published in ref. . In that work the data were contaminated by the inclusion of prompt nucleons; prompt nucleons are excluded from consideration in this work. In that work the second moment of the cluster distribution was determined based on the charge of a cluster rather than its mass as is done in this work. Previously, the second moment was generated from a cluster distribution that was not normalized to the changing size of the system as is done here. Furthermore the prior analysis consisted of only one quarter of the total number of events used in the present analysis. Thus the current analysis has higher statistics, has been freed of prompt nucleons, has a second moment that has been constructed with the masses from the cluster distribution and a cluster distribution that has been normalized to the changing system size. The exclusion of prompt nucleons and normalization to the changing system size are an effort to address the criticisms raised in and rebutted in . When the $`\gamma `$-matching procedure was applied to the data presented in this paper essentially the same results as presented in ref. were recovered. See Table II and Figure 25d. One difference observed is in the value of the critical point returned, $`m_c^{94}=26\pm 1`$ reported in ref. and $`m_c^{99}=21\pm 2`$ reported in this work. The difference is not as great as it appears to be. The origin of the published value of $`m_c^{94}`$ lies in picking the peak of the distribution of $`m_c`$ values that satisfied conditions (1) and (2) as the location of the critical point. The value was estimated based on the height of the peak and the error based on the width of the peak. The mean and RMS of the $`m_c`$ distribution in ref. suggest a value of the critical point of $`m_c^{94}=25\pm 3`$. This value agrees, to within error bars, with the value of $`m_c^{99}`$ presented here. The relatively small shift in $`m_c`$ can then be understood to arise from the differences in the data sets. Noting this it is clear that the present $`\gamma `$-matching analysis is in agreement with the previous work.
The results of the present work are, again, in keeping with the expected results of a small system undergoing a continuous phase transition. There is some region where matching $`\gamma `$ values can be obtained, some regions in $`ϵ`$ where th $`\gamma `$-power law overwhelms all other effects. The fits in Figure 25d are of quality than those for random partitions in Figure 25c and cover a greater range. When compared to the percolation $`m`$ results the multifragmentation data compare favorably in terms of overall goodness of fits, width of fit region and matching of $`\gamma _\pm `$. See Table II. The location of the critical point returned by this analysis also compares well with the location from other analyses. See Table I.
## V Corrections to Scaling
In the last section it was seen that the $`\gamma `$-power law and the data for the second moment in all systems have agreed over only a limited area. To some degree this is to be expected. Near the critical point, assumptions valid for thermodynamic systems are invalid for the finite systems discussed in this work. For that reason, finite size effects dominate at the critical point and the second moment merely peaks instead of diverging. Far from the critical point other effects come into play. The scaling assumptions inherent in the FDM are valid only in the neighborhood of the critical point. The size of this neighborhood is somewhat ill defined and seems to depend on many factors, e.g. the quantity in question, the nature of the system, the size of the system and so on. Scaling behavior in physical systems can be observed over a wide range in temperatures and densities. This is most elegantly illustrated in the Guggenheim Plot of scaled temperature ($`T/T_c`$) as function of scaled density ($`\rho /\rho _c`$) for several different gases (Ne, Ar, Kr, Xe, N<sub>2</sub>, O<sub>2</sub>, CO and CH<sub>4</sub>). In that plot the data collapse onto a curve that is well described by a power law with an exponent of $`\beta =1/3`$. The range in validity of this agreement between data and power law is shown on the Guggenheim Plot to be over a range of $`\mathrm{\Delta }T0.5T_c`$ and $`\mathrm{\Delta }\rho 2.5\rho _c`$. However, another system, the combination of isobutyric acid and water, shows the Guggenheim type of scaling only very near the critical point . Already when the range considered is $`\mathrm{\Delta }T0.04T_c`$ and $`\mathrm{\Delta }\rho 0.01\rho _c`$ corrections to scaling can be observed. To that end, higher order corrections to scaling are now examined in order to determine if fits such as those shown in previous sections can be improved. However, any improvement comes at the expense of more fit parameters and assumptions.
To fully explore corrections to scaling in the context of the present systems where the cluster distributions serve as the main observable the FDM is revisited in a fashion employed in references and . Assuming coexistence eq. (7) is then re-written as
$$n_{A_f}(ϵ)=q_0A_f^\tau \left(f_0(z)+A_f^\mathrm{\Omega }f_1(z)+\mathrm{}\right),$$
(41)
where $`f_1(z)`$ is the correction-to-scaling function and $`\mathrm{\Omega }`$ is the correction-to-scaling exponent. The form of eq. (41) anticipates the presence of a second function of $`z`$. In section IV-3 it was found empirically that both the scaled percolation and multifragmentation cluster distributions ($`n_{A_f}(ϵ)/q_0A_f^\tau `$) could be reasonably well described by the sum of two gaussians, eq. (40). In that treatment, the amplitude of each gaussian was a constant, $`a_1`$ and $`a_2`$. If $`A_f`$ is restricted to a single value, the prescription give by eq. (41) is equivalent to that of eq. (40). Eq. (41) predicts that there should be an ordering to the scaled cluster distributions, i.e. smaller cluster sizes should lie above the larger clusters due to the correction term. This can be observed in Figures 19a and b in the neighborhood of the maximum of the scaled data for the percolation systems. In the tails of the distribution, either large or small cluster production is suppressed. In the case of multifragmentation data, Figure 19d, the ordering is generally observed where the statistics are adequate, namely, near $`z=0`$. The ordering of the random partitions implies that $`\mathrm{\Omega }<0`$.
From eq. (41) it possible to derive the corrected isothermal compressibility (second moment) power law. Following the method in section III leads to:
$`\kappa _T`$ $``$ $`(\rho ^2k_bT)^1\times `$ (45)
$`\left(\right|{\displaystyle \frac{q_0}{\sigma }}{\displaystyle _0^\pm \mathrm{}}dzf_0(z)|z|^{\frac{3\tau \sigma }{\sigma }}|+`$
$`|{\displaystyle \frac{q_0}{\sigma }}{\displaystyle _0^\pm \mathrm{}}dzf_1(z)|z|^{\frac{3\tau \sigma \mathrm{\Omega }}{\sigma }}||ϵ|^{\frac{\mathrm{\Omega }}{\sigma }})`$
$`\times |ϵ|^{\frac{\tau 3}{\sigma }}`$
which is usually simplified and written as
$$\kappa _T\mathrm{\Gamma }_\pm |ϵ|^\gamma \left(1+a_\pm |ϵ|^\mathrm{\Delta }\right).$$
(46)
Now the overall amplitude, $`\mathrm{\Gamma }_\pm `$ is given by the first integral, and the correction-to-scaling amplitude is given by the second integral divided by the first. The correction-to-scaling exponent is $`\mathrm{\Delta }=\mathrm{\Omega }/\sigma `$.
Using eq. (46) to fit the second moment distribution would lead to determining four fit parameters: two amplitudes and two exponents. To explore the effects of corrections to scaling an assumption was made as to the universality class of the system in question and thus the choice of exponent values. For the three dimensional percolation universality class $`\mathrm{\Delta }=1.22`$ , and for the three dimensional Ising universality class $`\mathrm{\Delta }=0.56`$ -. The amplitudes were left as free parameters and the second moment of the cluster distributions were fit. The value of the critical point determined from the $`\gamma `$-matching analysis of section IV-B-5 was used for each system. Figure 26 shows the results.
For the percolation system with the corrections-to-scaling a better fit to the second moment data was possible over a greater range in $`ϵ`$, up to twice the range of the average fitted region in the $`\gamma `$-matching analysis, see Figures 26a and b. The fits still failed to reproduce the behavior of $`M_2`$ near the critical point where finite size effects dominate the system. Table III lists the results for the critical amplitudes, $`\mathrm{\Gamma }_\pm `$. The agreement in the critical amplitudes determined in this analysis and the amplitudes from the $`\gamma `$-matching analysis is due to the agreement of the behavior of eq. (46) and the $`\gamma `$-power law from the $`\gamma `$-matching analysis over the region in $`ϵ`$ determined by $`\gamma `$-matching. Thus the $`\gamma `$-matching analysis finds regions that are the least affected by higher-order corrections to scaling.
For the random partitions an improvement is only observed for the high multiplicity region where a better fit over a larger range was obtained for both choices of universality class, see Figure 26c. The low multiplicity events showed no such improvement partly due to the limited range in $`ϵ`$ available. Both the three dimensional percolation and three dimensional Ising exponents were used in this analysis for the random partitions and the multifragmentation data. Both choices of universality classes showed similar results. The lack of effect of corrections to scaling is to be expected in a system that does not follow FDM like scaling laws.
The multifragmentation data also showed improvement resulting in a better agreement between the fits and the $`M_2`$ data points over a larger range in $`ϵ`$, see Figure 26d. The improvement was observed for both choices of universality classes thus indicating this analysis is insensitive to the differences , , though the goodness of fit was better for the choice of the three dimensional Ising exponents over the same fit regions (3D Ising: $`\chi _{\nu +}^2=0.4`$ and $`\chi _\nu ^2=0.7`$; 3D percolation: $`\chi _{\nu +}^2=0.7`$ and $`\chi _\nu ^2=1.1`$). At this level of analysis it appears that corrections to scaling improves the fits for the $`\gamma `$-power law. Whether this is due to the presence of a continuous phase transition in nuclear multifragmentation, or merely extra terms in the fitting function remains an open question.
## VI Changing System Size and Finite Size Corrections
It has been pointed out that the previous statistical analysis of gold multifragmentation ignored the changing size of the system . To first order this may have been a reasonable procedure as many statistical signatures of a continuous phase transition have been observed both before and after the scaling to account for the changing system size has been performed; e.g. the $`\gamma `$ power law shown here and in previous works agree well. However, the data collapse in Figure 19d is qualitatively not as great as that shown by the percolation ($`m`$) system in Figure 19b and the agreement between the calculated and measured $`M_2`$ behavior in gold multifragmentation in Figure 24d is qualitatively not as good as that shown by percolation ($`m`$) in Figure 24b. In this section, the effect of the changing size of the system is explored and accounted for.
The size of the multifragmenting system is shown in Figure 27a after ref. . An approximately linear relation between the system size, $`A_0`$, and $`m`$ was found: $`A_0=1991.6m`$, see Figure 27a. The functional form of $`A_0(m)`$ was used in the following analysis to account for the changing system size in an average way, i.e. not on an event-to-event basis.
If the multifragmenting system is assumed to be a system undergoing a phase transition, then the theory of finite size scaling of the critical point - asserts that the effective critical temperature, $`T_c(A_0)`$, changes as a function of the system size. Coupling this with the changing size of the system indicates that at each value of $`m`$ the value of $`T_c^{eff}(A_0)`$ is different.
The value of $`T_c^{eff}(A_0)`$ can be determined in the following manner. First a relation between the multifragmenting system’s temperature, $`T`$, and $`m`$ must be determined. Again from ref. a relation can be found, see Figure 27b. Two estimates of $`T`$ were made, one from a Fermi gas (uncorrected for the effect of expansion energy), $`T_i`$, and the other from an isotopic yield ratio thermometer, $`T_f`$, . These temperatures give an approximate indication of the initial and final temperatures of the system. Figure 27b shows that $`T_i`$ is well fit by a quadratic function, while $`T_f`$ is well fit by a linear function. Another linear function reproduces the average of $`T_i`$ and $`T_f`$: $`T=3.0+0.14m`$, this was used for the following analysis. The critical temperature of infinite nuclear matter was assumed to be $`T_c^{\mathrm{}}=17\pm 1`$ MeV .
From the Fisher $`\tau `$-power law analysis the value of the multiplicity at the critical point was determined to be $`m=m_c=22\pm 1`$. The system size at that point is then $`A_0(m=22)=164\pm 2`$ and the temperature is $`T=6\pm 2`$ MeV. This indicates that the critical temperature for a nuclear system with 164 nucleons, $`T_c^{eff}(A_0)`$, is approximately $`6\pm 2`$ MeV.
According to theory, to first order the critical point scales with system size as:
$$(T_c^{eff}(A_0)T_c^{\mathrm{}})/T_c^{\mathrm{}}=bA_0^{1/d\nu },$$
(47)
where $`d`$ is the Euclidian dimension of the system and $`\nu `$ is the so-called hyperscaling exponent. At the smallest of system sizes higher order correction terms may play an important role in the scaling of the critical point . It is assumed in this procedure that the above mentioned finite size scaling dominates all other effects, including those which arise due to the charge of the protons present in the excited remnant. This is in keeping with the general philosophy of the theory of critical phenomena where near the critical point the precise details of the system are irrelevant.
If it is assumed that gold multifragmentation is the result of a continuous phase transition and that transition occurs in three dimensions, $`d=3`$, then using the hyperscaling relation
$$\nu =\frac{\tau 1}{d\sigma },$$
(48)
with the extracted values of $`\sigma `$ and $`\tau `$, gives $`\nu =0.63\pm 0.07`$. The coefficient $`b`$ in eq. (47) can be determined using $`T_c^{eff}(A_0(m_c))`$, $`T_c^{\mathrm{}}`$, $`d`$ and $`\nu `$; resulting in $`b=9\pm 2`$. Note that this value of $`b`$ suggests that $`T_c^{eff}(A_0)=0`$ for $`A_0=70\pm 40`$. This is a result of the form of eq. (47) and the notion that the critical temperature lowers as the size of the system decreases. Presumably higher order effects not taken into account in eq. (47) will affect the location in system size, $`A_0`$, where the effective critical temperature vanishes. For the form of finite size scaling corrections shown in eq. (47), only $`b=1`$ yields an effective critical temperature that vanishes at $`A_0=1`$.
Now eq. (47) can be used to solve for $`T_c^{eff}(A_0(m))`$ and thus determine the scaled control parameter:
$$ϵ^{scaled}=(T_c^{eff}(A_0(m))T)/T_c^{eff}(A_0(m)).$$
(49)
The analysis to extract the exponent $`\sigma `$ was performed with this $`ϵ^{scaled}`$ by finding the peak in $`A_f`$ sized cluster production as a function of $`n_{A_f}(ϵ^{scaled})`$ versus $`ϵ^{scaled}`$. Previously, it was argued that the largest cluster should be excluded for all values of $`ϵ`$ to account for finite size effects in the analysis to extract $`\sigma `$. The analysis in this section directly accounts for finite size effects, thus the standard FDM formalism with respect to the largest cluster is followed. This results in a value of $`\sigma =0.65\pm 0.07`$ and $`z_{max}^{scaled}=11\pm 2`$, Figure 28a shows the resulting power law.
The scaling function for gold multifragmentation was then plotted using the above corrections for the changing system size and finite size scaling, see Figure 28b. The data collapse is qualitatively better than in Figure 19d. The two gaussian parameterization of $`f(z^{scaled})`$ was fit to the scaled scaling function and is shown in Figure 28b with parameters listed in Table IV.
Using the fitted parameterization of the scaling function and other quantities, the $`\gamma `$-power law can be determined as before. Figures 28b and c show the measured $`M_2`$ of gold multifragmentation plotted as a function of $`ϵ^{scaled}`$. A $`\gamma `$-power law was plotted on, not fitted, Figures 28b and c with $`\gamma =1.3\pm 0.2`$ (from $`\tau `$, $`\sigma `$ and eq. (24)) and offsets determined via the scaled scaling function. This $`\gamma `$-power law agrees with the measured $`M_2`$ over nearly all of the $`ϵ^{scaled}`$-range, the exception being near $`ϵ^{scaled}=0`$ where finite size of the system limits the maximum of $`M_2`$.
As stated above, this analysis makes no attempt to account for the Coulomb energy of the system in the formulation of finite size scaling. The effect of the Coulomb energy may evidence itself in the degree of collapse of the data. The fact that, to a large extent, the data do collapse without any explicit adjustment to the theory (i.e. formalation of the scaling variable, $`z^{scaled}`$) may then be an indication that the Coulomb energy is not a major perturbation. Hence, it is perhaps justifiable to adhere to the Coulomb-free theory when scaling $`T_c`$.
## VII Discussion and conclusions
The focus in the present paper has been on cluster distributions and the types of analyses which can shed light on their creation mechanism. In particular, attempts were made to identify procedures that can distinguish those distributions which are related to critical behavior from those which are not. While this question is easily answered for systems containing large numbers of constituents, it is much more difficult to address the case of interest here, namely systems with at most only a few hundred particles. In such systems, finite size effects play a large role. Therefore, two different computational models, bond building percolation and a random partitions, have been used. It is well known that in the macroscopic limit, the former system possesses a continuous phase transition characterized by a unique scaling function and set of critical exponents while the latter system does not. In addition, data arising from the multifragmentation of gold nuclei has been studied using the same procedures. For this system, it is not known, a priori, whether a critical point is present.
Many cluster properties have been proposed as being suitable measures of critical behavior. Among these are the fluctuations in the size of largest fragment (Figure 1), peaking in the quantity $`\gamma _2`$ (Figure 3), peaking behavior in $`M_2`$ (Figure 5), Campi plots (Figure 6) and simple power law behavior in the cluster mass distribution for a particular value of the appropriate control parameter. It was seen that none of these measures, taken alone or together, was sufficient to distinguish a system possessing critical behavior from one which does not.
The first procedure which produced different results for critical and noncritical systems was the single parameter power law fit to the cluster mass distribution (Figures 12 through 15). For the percolation systems and for the multifragmentation data, it was shown that the one parameter power law fit describes the data well only over a very limited range of the control parameter. The value of the control parameter where the power law fit is best is very close to where the abovementioned peaks occur. However, for the random partitions, this was not the case.
If a system possesses a critical point, it is also expected to possess a scaling function that describes its behavior away from the critical point. Phase transition theory specifies how the argument of this function depends on cluster mass and distance from the critical point. If such a function exists, the theory permits the determination of the critical exponent $`\sigma `$. This determination was done for the percolation system yielding satisfactory agreement with its known value. The same procedure was applied to the random breakup model and to the multifragmentation data. It was clear from this analysis that percolation and multifragmentation were very similar in many features, while the random system was distinctly different. See Figures 17 and 18.
Again, if a system possesses a critical point, it is expected to possess a scaling function that describes its behavior away from the critical point. Therefore, when the data is properly scaled, it should collapse onto a single curve. Figure 19 shows the amount of data collapse for the systems discussed here. The quality of the data collapse (Figures 19-21) reinforces the notion that the random breakup system is different from the others. Although the precise form of the scaling function is not dictated by phase transition theory, both the percolation system and multifragmentation data were satisfactorily described by a sum of two gaussians. The random breakup model was not.
The issue of finite size scaling was discussed. Unlike the other systems examined here, which had a fixed number of constituents, the nuclear multifragmentation data originated from systems whose size varied monotonically with observed charged particle multiplicity. See Figure 27a. Phase transition theory makes a prediction, eq. (47), as to how the value of the control parameter at the critical point changes as a function of system size. Applying eq. (47) produced an improvement in the quality of the data collapse for the scaling function and yielded a better prediction for the behavior of $`M_2`$. See Figure 28.
Critical exponent values have been determined in an unbiased manner for each system. For both sets of analyses on the percolation clusters, the standard percolation exponents were recovered to within error bars. For the random partitioning, exponents could be extracted, but none that fulfilled well known scaling laws. The exponent values determined from the gold multifragmentation cluster distributions fulfill the scaling laws, to within error bars, and fall near the three dimensional Ising universality class.
The effect of secondary decays from hot initial fragments on the critical exponents has not been explicitly considered in this paper. In the SMM - such effects become significant above $`E^{}/A_0=7`$ MeV/nucleon. Thus $`\tau `$ and $`\gamma `$, which are determined at lower excitation energies, will be unaffected in the SMM’s fragment distributions. The exponent $`\sigma `$ is determined by the multiplicities at which individual light fragment yields attain their peak values. As shown in Figure 17d, the lightest fragments peak at large multiplicities, corresponding to excitation energies for which secondary decay are important in the SMM. An SMM calculation indicates that the value of $`\sigma `$ from the SMM’s fragment distribution was increased by about $`70`$% due to this effect. However, it is unclear from that calculation that the effects of secondary decay are as great in the experimental data as they are in the SMM. The SMM calculation over predicts the yield of light fragments which could indicate that the SMM estimates of secondary decay are too severe. Corrections to the model independent quantities determined in this paper based on the SMM calculations are premature. See Appendix C for discussions of another set of model based interpretations.
Although the multifragmentation data possess many of the gross and detailed characteristics that the percolation system does, it is not at all obvious why this should be the case. After all, real nuclei obey quantum mechanics, have varying binding energies per particle, and, most significantly, are charged. On the other hand, it is well known that near a critical point the details of the interaction become unimportant and only the dimensionality of the system and the dimension of the order parameter are important. As noted in the introduction, the attractive nuclear force bears a similarity to a van der Waals force. However, the Coulomb force is a long range force and imposes a natural limit to the size of stable nuclei. Thus, it is not clear to what extent a finite charged system can exhibit critical behavior when the macroscopic system cannot exist. The exact role of the Coulomb force in physical systems undergoing a change of phase is currently of great interest , and is, at this point, an open question. The philosophy of this paper has been to make use of phase transition theory as it applies to uncharged systems. What results for the analysis of the gold multifragmentation data bears great similarity to the results of the same procedures applied to a system known to possess a critical point. It is tempting then to conclude that multifragmentation is related to critical behavior occurring in a finite nuclear system.
## A Discussion of power law in random partitions
As a demonstration of the power of the sort of scaling analysis presented above it is shown that the random partitions follow a simple power law of $`N_{A_f}A_f^1`$. Figure A1 shows the scaled cluster distribution as a function of $`m`$ from clusters with $`A_f3`$. The data nearly collapses to unity along the horizontal axis over the multiplicity range for $`m>5`$. The deviations are due to the constraints of $`m`$ and finite size. With a simple scaling analysis the underlying power law describing the cluster distribution becomes clear.
## B Riemann $`\zeta `$ function summation
The value of $`q_0`$ used in this work based on the $`\zeta `$-function was generated with VAX FORTRAN code using double precision and letting the sum run from $`1`$ to $`10^6`$. The sum was terminated at this point in order to keep computing times within reason. For a value of $`\tau =2.18`$ summing to $`10^6`$ gives a value of $`q_0`$ that is within 10$`\%`$ of the value when the sum is terminated at $`10^{10}`$, see Figure A2. Increasing the upper limit of the summation in the $`\zeta `$-function causes no significant changes in the analysis presented in this work.
## C Lattice gas model interpretations
In a recent paper , Gulminelli and Chomaz (GC) report on calculations made using a lattice gas model. Their results can be divided into two different categories: (1) an analysis of the thermodynamic quantities of the system; (2) an analysis of the statistical properties of the droplet or cluster distribution.
In their paper GC performed canonical calculations on a three dimensional lattice of $`n`$ sites ($`n=216`$ or $`512`$) with periodic boundary conditions at a given temperature, $`T`$, and density, $`\rho `$. The density was varied by changing the number of particles in the system, $`A_0`$, with $`\rho =A_0/n`$. For a choice of $`T`$, $`\rho `$ was varied and the system’s chemical potential, $`\mu `$, was determined. At some values of $`T`$ the $`\mu `$ vs. $`\rho `$ ($`\mu `$ vs. $`A_0`$) isotherm exhibited a back bend. As $`T`$ increased the back bend disappeared. This was taken as evidence that a first order phase transition (back bending $`\mu `$ vs. $`A_0`$) had culminated in a continuous phase transition (flat $`\mu `$ vs. $`A_0`$). The critical point, ($`T_c`$, $`\rho _c`$), was thus defined. Via a Maxwell construction, GC also determined the boundary of the coexistence zone.
The theoretical ideas behind the procedure followed by GC are sound (with the exception of the questionable definition of volume in a system with periodic boundary conditions), however the execution is flawed. The error in execution lies in the method by which $`\rho `$ was varied. GC varied $`\rho `$ by holding $`n`$, the volume of the system, fixed and varying $`A_0`$, the number of constituents or “size” of the system. Therefore, at each value of $`\rho `$ GC have a system of different size. It is well known that the effective $`T_c`$ of a system varies with its system’s size -. Thus each time GC varied $`A_0`$ to vary $`\rho `$ they were dealing with a system with a different $`T_c`$. With this understanding it is clear that a plot of $`\mu `$ versus $`A_0`$ as shown by GC in their Figure. 1a, which serves as the basis for their determination of the critical point, is difficult to interpret. $`T_c`$ cannot be deduced in the manner followed by GC. Similarly, it is impossible to determine $`\rho _c`$ for a system of $`A_0=100`$ using information from a system of $`A_0=200`$ without accounting for the finite size scaling of the critical point. This makes the thermodynamically extracted values of the critical point suspect and calls into question the validity of the coexistence line shown by GC in the lower part of their Figure 3.
Via an analysis of the droplet distributions GC determined the value of the critical exponent $`\tau =\tau _{max}`$ based on the peak in the production of different size droplets. The value of $`T_c`$ was determined coupling the $`\tau _{max}`$ results with a two parameter power law fit the droplet distribution at various temperatures, $`\tau (T)`$, at $`T_c`$ $`\tau _{max}=\tau (T_c)`$. The use of two parameter fits is generally an improper method to determine $`\tau `$ -. In this analysis GC determine $`T_c`$ on a system-by-system basis, avoiding the error of grouping systems of different size together. GC also extracted the exponent $`\sigma `$ and then used the above results to collapse the droplet data onto a single curve illustrating the scaling function that modifies the power law away from the critical point. While the data collapses onto a single curve, the value of the curve at $`T_c`$ is problematic. From GC’s work the value of the scaling function at $`T_c`$ is $`f(T_c)0.5`$ which contradicts the $`\tau `$ values extracted by GC according to the Riemann $`\zeta `$-function relation between $`\tau `$ and $`f(T_c)`$ ; e.g. $`\tau 2.2`$ gives $`f(T_c)0.2`$, while GC’s $`f(T_c)0.5`$ gives $`\tau 2.7`$.
Finally GC combine the results of their analyses and assert that the Kértesz line of $`\tau `$ power laws continues into the coexistence region of the lattice gas and that this behavior is observed due to the finite size of the system and is not observed in systems of much greater size. However, due to the problems in GC’s method of varying $`\rho `$ to find the coexistence curve, such a claim is premature. Instead, their estimates of $`T_c`$ from the droplet analysis, which are shown as a string of points in their Figure 3, may reflect the finite size scaling of the critical point and not the Kértesz line extending into the coexistence zone. That each different size system appears to be at $`\rho _c`$ as well as $`T_c`$, illustrated by the data collapse, may be explained by the relation $`\mu \mu _{coex}(\rho \rho _c)^\delta `$. For a 3D Ising system $`\delta 4.8`$. Thus, systems of finite size near their critical density are nearly on the coexistence curve and show the sort of data collapse illustrated in GC’s Figure 2 .
This work was supported in part by the U.S. Department of Energy Contracts or Grants No. DE-ACO3-76F00098, DE-FG02-89ER-40513, DE-FG02-88ER-40408, DE-FG02-88ER40412, DE-FG05-ER40437 and by the U.S. National Science Foundation under Grant No. PHY-91-23301. |
warning/0002/math0002206.html | ar5iv | text | # Chapter 1 Fiber with Intrinsic Action on a 1+1 Dimensional Spacetime
## Chapter 1 Fiber with Intrinsic Action on a $`1+1`$ Dimensional Spacetime
Robert W. Johnson
footnotetext: This work is an outgrowth of a paper presented at the 5th International Conference on Clifford Algebras and their Applications in Mathematical Physics, Ixtapa, Mexico, June 27 - July 4, 1999.
AMS Subject Classification: 83C10, 15A66, 53B30.
### 1 Introduction
In a recent paper I described a construction for a vector space with metric. In this construction one forms elements $`xx`$ in a direct product algebra where $`x`$ is an element in an underlying finite group algebra . One uses a particular decomposition of the direct product algebra to obtain a direct sum of subspaces. One then observes that vectors in the various subspaces are interrelated. In particular examples that I considered there existed a $`1d`$ subspace whose measure was determined by the components of a second vector in a higher dimensional subspace. I proposed that the measure in the $`1d`$ subspace be identified with the norm of the vector in the higher dimensional subspace. In this way both a vector space and its norm are viewed as residing in particular subspaces of a given direct product algebra $`xx.`$ In this construction the signature of the metric arises intrinsically from the particular underlying finite group algebra.
The realization of Clifford algebras in terms of underlying finite group algebras is described by Salingaros . The present work differs from the approaches of, e.g., Hestenes, Lounesto, and Greider by considering vector spaces obtained through the decomposition of a direct product algebra (having elements of form $`xx)`$ rather than through a decomposition of a Clifford algebra. In particular cases that I consider, however, the underlying finite group algebra (containing elements $`x)`$ does correspond to a Clifford algebra.
This work is motivated by the analogous construction in quantum mechanics where one forms observable vector spaces in terms of bilinear functions of an underlying state vector. The quantum mechanical $`2`$-state problem provides one instructive example . For this problem the underlying finite group is $`C_4H,`$ the direct product of the cyclic group of order 4 and the quaternion group. The quantum wave function $`\psi `$ is an element in a left ideal of the $`C_4H`$ group algebra over the real numbers. <sup>1</sup><sup>1</sup>1This algebra is also termed the complex quaternion algebra. Elements in this algebra are linear combinations of the $`C_4H`$ group elements with real coefficients. The $`C_4H`$ group multiplication is used along with the distributive law to induce the algebra product rule. The polarization vector and its norm, the total probability, reside in particular subspaces of the direct product algebra whose elements $`\psi \psi `$ are formed from the product of two copies of the quantum mechanical state $`\psi .`$
Vector spaces with metric that are constructed using this procedure are also of interest as models for typical fiber vector spaces residing at each location of a configurational manifold such as arises in the context of classical Lagrangian mechanics. In this paper I develop an algebraic model for a typical fiber at a location $`P`$ of a configurational manifold $`M`$ having 1 space and 1 time dimension. The fiber that I construct contains both the tangent and momentum vector spaces as subspaces. In addition, a subspace that can be associated with the flow of action at $`P`$ is included and arises in an intrinsic way. The construction for a configuration space with 2 space dimensions follows in a completely analogous way. The details for these two 2-dimensional cases are transparent. The extension of this construction to higher dimensional cases is also straightforward; however, the multiplicity of subspaces in the corresponding direct product algebras makes the interpretation of the overall structure more challenging, and it has not been fully addressed by this author.
In section 2, I review the construction of a vector space with metric signature $`(p,q)=(1,1)`$ that corresponds to the tangent space at a point of a $`1+1`$ dimensional configurational manifold . This construction uses the very simple $`C_2`$ group algebra. In section 3 I use the $`D_2=C_2C_2`$ group algebra to obtain the vector space that is the primary focus of this paper. In section 4 I summarize these results and indicate work that still remains to be done.
### 2 Algebraic Model for Tangent Space at a Point of a $`1+1`$ Dimension Spacetime
As a starting point for this paper let us review my earlier construction of an algebraic model for the vector space with signature $`(p,q)=(1,1).`$ This vector space corresponds to the tangent space at each location of a $`1+1`$ dimensional spacetime $`M.`$ The tangent space at an arbitrary location $`(t_0,q_0)M`$ consists of tangent vectors to curves passing through $`(t_0,q_0).`$ Here $`t_0`$ and $`q_0`$ denote, respectively, time and spatial location. For a curve $`\gamma =\gamma (\lambda ),\lambda ,`$ the tangent vector is defined by
$$\frac{d\gamma }{d\lambda }=\underset{\lambda 0}{lim}\frac{\gamma (\lambda )\gamma (0)}{\lambda }=(\frac{dt}{d\lambda },\frac{dq}{d\lambda })$$
where $`\gamma (0)=(t_0,q_0)`$ and $`\gamma (\lambda )M.`$
We begin the construction with the $`C_2`$ group which contains two elements $`C_2=\{\mathrm{𝟏},𝐞\}`$ with $`𝐞^_2=\mathrm{𝟏},`$ and then form the vector space $`V_{C_2}`$ whose elements consist of formal sums of the elements of $`C_2`$ with real coefficients. An arbitrary element $`xV_{C_2}`$ can be written $`x=x_0\mathrm{𝟏}+x_1𝐞.`$ The product rule, that is induced by the group multiplication, is used to form the group algebra. We then consider the direct product of two copies of a vector $`x`$ in the algebra $`(\mathrm{𝟏}\mathrm{𝟏}\mathrm{𝟏})`$:
$`xx`$ $`=(x_0\mathrm{𝟏}+x_1𝐞)(x_0\mathrm{𝟏}+x_1𝐞)`$
$`=x_0^2(\mathrm{𝟏}\mathrm{𝟏})+x_0x_1(\mathrm{𝟏}𝐞)+x_0x_1(𝐞\mathrm{𝟏})+x_1^2(𝐞𝐞)`$
$`=(x_0\mathrm{𝟏}+x_1\mathrm{𝐄𝐞})(x_0\mathrm{𝟏}+x_1𝐞)`$
where the second line follows from bilinearity of the direct product and in the third line we introduce the notation $`𝐄=𝐞𝐞`$ and $`𝐞=\mathrm{𝟏}𝐞.`$ $`\mathrm{𝐄𝐞}=𝐞\mathrm{𝟏}`$ follows from the product rule in the direct product algebra. Acting on $`xx`$ with the projection operators $`P_\pm =\frac{1}{2}(\mathrm{𝟏}\mathrm{𝟏}\pm 𝐄)`$ we obtain
$$xx=[P_+(\frac{dt}{d\lambda }\mathrm{𝟏}+\frac{dq}{d\lambda }𝐞)+P_{}\frac{ds}{d\lambda }](\mathrm{𝟏}\mathrm{𝟏}),$$
where
$`{\displaystyle \frac{dt}{d\lambda }}`$ $`=x_0^2+x_1^2`$
$`{\displaystyle \frac{dq}{d\lambda }}`$ $`=2x_0x_1`$
$`{\displaystyle \frac{ds}{d\lambda }}`$ $`=x_0^2x_1^2=({\displaystyle \frac{dt}{d\lambda }}^2{\displaystyle \frac{dq}{d\lambda }}^2)^{\frac{1}{2}}`$
The measure $`\frac{ds}{d\lambda }`$ of the $`1d`$ $`P_{}xx`$ subspace is determined up to sign by the two components $`(\frac{dt}{d\lambda },\frac{dq}{d\lambda })`$ of the $`2d`$ $`P_+xx`$ subspace and can be interpreted as their norm. We note that the measure $`x_0^2+x_1^2`$ associated with the $`\frac{dt}{d\lambda }`$ increment is positive definite and so has an intrinsic directionality. Also, since $`(2x_0x_1)^2(x_0^2+x_1^2)^2,`$ we have $`|\frac{dq}{dt}|1`$ so that there is a maximum speed.
In this way the product $`xx`$ provides a model for an element in the tangent space of a $`1+1`$ dimensional spacetime. The set of all such elements $`xx`$ can be identified with the tangent space itself.
Continuing further we find that rotations of the $`2d`$ vector space $`P_+xx`$ are induced by acting on $`x`$ with an element $`u=u_0\mathrm{𝟏}+u_1𝐞`$ and forming the product $`(xx)(uu)=xuxu.`$ For $`u`$ such that $`u_0^2u_1^2=1,`$ $`P_{}xx`$ is unchanged while the $`P_+xx`$ vector undergoes a proper orthochronous rotation.
A completely analogous treatment for the $`2d`$ Euclidean case is obtained by substituting the $`C_4`$ group algebra for the $`C_2`$ group algebra. This approach also extends to higher dimensional vector spaces, though in these cases we encounter a multiplicity of subspaces in the direct product algebra which make the interpretation of the overall structure more involved .
In the following section we extend this treatment to the case of a typical fiber at a location $`P`$ of a configurational manifold $`M`$ that contains both the tangent and momentum vector spaces.
### 3 Algebraic Model for Fiber on $`1+1`$ Spacetime with an Intrinsic Action
Let us briefly review the classical mechanics that motivate this construction. The action function
$$S_{t_0,q_0}(t,q)=_\gamma L𝑑t$$
is the integral of the Lagrangian $`L=L(\stackrel{}{q},q,t)`$ along an extremal path $`\gamma (\lambda ),`$ $`\lambda `$ real numbers, connecting an initial point $`(t_0,q_0)`$ with $`(t,q)`$ . The action function for a free particle with mass m in a locally Minkowski coordinate system can be written as
$$S_{t_0,q_0}(t,q)=_\gamma m\frac{ds}{dt}dt$$
where $`ds`$ is the proper time with measure $`ds=(dt^2dq^2)^{\frac{1}{2}}`$ and $`\gamma (\lambda )`$ is a locally straight line . This action corresponds to the Lagrangian $`L=m\frac{ds}{dt}`$ . The rate of change of the action along the path $`\gamma `$ for a fixed initial point is
$$\frac{dS}{d\lambda }=p\frac{dq}{d\lambda }H\frac{dt}{d\lambda }$$
(3.1)
where
$$p=\frac{L}{\stackrel{}{q}}=m\stackrel{}{q}/(1\stackrel{}{q}^2)^{\frac{1}{2}}$$
and
$$H=p\stackrel{}{q}L=m/(1\stackrel{}{q}^2)^{\frac{1}{2}}.$$
We now construct a model for the local tangent and momentum vectors of a trajectory $`\gamma (\lambda )`$ on $`1+1`$ dimensional manifold for which such an action function arises intrinsically. For this construction we use the abelian group $`(𝐞_{\mathrm{𝟏𝟐}}:=𝐞_\mathrm{𝟏}𝐞_\mathrm{𝟐})`$
$$D_2=C_2C_2=\{\mathrm{𝟏},𝐞_\mathrm{𝟏},𝐞_\mathrm{𝟐},𝐞_{\mathrm{𝟏𝟐}}\}$$
where $`𝐞_1^2=𝐞_2^2=𝐞_{12}^2=\mathrm{𝟏}.`$ A general element of the $`D_2`$ group algebra can be written
$$x=x_0\mathrm{𝟏}+x_1𝐞_\mathrm{𝟏}+x_2𝐞_2+x_3𝐞_{12}.$$
We decompose $`x`$ into a sum of two left ideals obtained by acting on the right with the projection operators $`P_{\pm 2}=\frac{1}{2}(\mathrm{𝟏}\pm 𝐞_2).`$ We have
$$x=x(P_{+2}+P_2)$$
where
$`xP_{+2}`$ $`=[(x_0+x_2)\mathrm{𝟏}+(x_1+x_3)𝐞_\mathrm{𝟏}]P_{+2}`$
$`xP_2`$ $`=[(x_0x_2)\mathrm{𝟏}+(x_1x_3)𝐞_\mathrm{𝟏}]P_2.`$
We now form the tensor product of two copies of $`x,`$
$$xx=x(P_{+2}+P_2)x(P_{+2}+P_2)$$
Expanding out this expression and acting on the left side with the projection operator $`P_{\pm 1}=\frac{1}{2}(\mathrm{𝟏}\mathrm{𝟏}\pm 𝐄_1),`$ where $`𝐄_1`$ = $`𝐞_\mathrm{𝟏}𝐞_\mathrm{𝟏},`$ we obtain, after a change of variables,
$`xx=`$ $`[P_{+1}({\displaystyle \frac{dt}{d\lambda }}\mathrm{𝟏}+{\displaystyle \frac{dq}{d\lambda }}𝐞_1)+P_1{\displaystyle \frac{ds}{d\lambda }}](P_{+2}P_{+2})`$
$`+[P_{+1}([{\displaystyle \frac{1}{2}}(H{\displaystyle \frac{dt}{d\lambda }}+p{\displaystyle \frac{dq}{d\lambda }}+m{\displaystyle \frac{ds}{d\lambda }})]^{\frac{1}{2}}\mathrm{𝟏}`$
$`\text{ }+[{\displaystyle \frac{1}{2}}(H{\displaystyle \frac{dt}{d\lambda }}+p{\displaystyle \frac{dq}{d\lambda }}m{\displaystyle \frac{ds}{d\lambda }})]^{\frac{1}{2}}𝐞_1)`$
$`+P_1([{\displaystyle \frac{1}{2}}(H{\displaystyle \frac{dt}{d\lambda }}p{\displaystyle \frac{dq}{d\lambda }}+m{\displaystyle \frac{ds}{d\lambda }})]^{\frac{1}{2}}\mathrm{𝟏}`$
$`\text{ }[{\displaystyle \frac{1}{2}}(H{\displaystyle \frac{dt}{d\lambda }}p{\displaystyle \frac{dq}{d\lambda }}m{\displaystyle \frac{ds}{d\lambda }})]^{\frac{1}{2}}𝐞_{12})](P_{+2}P_2+P_2P_{+2})`$
$`+[P_{+1}(H\mathrm{𝟏}+p𝐞_1)+P_1m](P_2P_2)`$
where $`𝐞_1=\mathrm{𝟏}𝐞_1`$ and $`𝐞_{12}=\mathrm{𝟏}𝐞_{12}.`$ In this expression
$`{\displaystyle \frac{dt}{d\lambda }}`$ $`=(x_0+x_2)^2+(x_1+x_3)^2`$
$`{\displaystyle \frac{dq}{d\lambda }}`$ $`=2(x_0+x_2)(x_1+x_3)`$
$`{\displaystyle \frac{ds}{d\lambda }}`$ $`=(x_0+x_2)^2(x_1+x_3)^2=[{\displaystyle \frac{dt}{d\lambda }}^2{\displaystyle \frac{dq}{d\lambda }}^2]^{\frac{1}{2}}`$
while for the momenta variables we have,
$`H`$ $`=(x_0x_2)^2+(x_1x_3)^2`$
$`p`$ $`=2(x_0x_2)(x_1x_3)`$
$`m`$ $`=(x_0x_2)^2(x_1x_3)^2=[H^2p^2]^{\frac{1}{2}}`$
In analogy to the $`C_2`$ case discussed above, we associate the $`xP_{+2}xP_{+2}`$ subspace with the tangent to a curve $`\gamma (\lambda )`$ on the configurational manifold. The $`P_{+1}(xP_{+2}xP_{+2})`$ portion is identified with the tangent vector $`(\frac{dt}{d\lambda },\frac{dq}{d\lambda })`$ while the $`P_1(xP_{+2}xP_{+2})`$ portion with measure $`\frac{ds}{d\lambda }=[\frac{dt}{d\lambda }^2\frac{dq}{d\lambda }^2]^{\frac{1}{2}}`$ is associated with the norm of the tangent vector.
Similarly, we identify the $`xP_2xP_2`$ subspace with the momentum of the trajectory $`\gamma (\lambda )`$ at $`\lambda .`$ The $`2d`$ $`P_{+1}(xP_2xP_2)`$ projection is identified with the momentum vector $`(H,p)`$ while the $`1d`$ $`P_1(xP_2xP_2)`$ projection with measure $`m=[H^2p^2]^{\frac{1}{2}}`$ is associated with the norm of the momentum vector.
So far in this development the velocity tangent vector $`(\frac{dt}{d\lambda },\frac{dq}{d\lambda })`$ is completely independent of the momentum vector $`(H,p).`$ These two vectors contain the 4 degrees of freedom of the original vector $`x`$ in the $`D_2`$ algebra.
Let us now consider the $`(xx)(P_{+2}P_2+P_2P_{+2})`$ subspace of this algebra. Both the $`P_{+1}`$and $`P_1`$ projections on this subspace are 2-dimensional. Taking the difference of the squares of the component measures for each of these two 2-vectors we find the resultant $`m\frac{ds}{d\lambda }.`$ The measure squared of the $`\mathrm{𝟏}`$ component of the $`P_1`$ subspace
$$\frac{1}{2}(H\frac{dt}{d\lambda }p\frac{dq}{d\lambda }+m\frac{ds}{d\lambda })$$
motivates its association with $`\frac{dS}{d\lambda }`$ of Eq. (1). Minimizing this measure while keeping $`m`$ and $`\frac{ds}{d\lambda }`$ constant and their product $`m\frac{ds}{d\lambda }`$ greater than zero requires setting the $`𝐞_{12}`$ component of the $`P_1`$ subspace to zero, since the difference of the squares of the two component measures is fixed. In this way we obtain the condition
$$m\frac{ds}{d\lambda }=p\frac{dq}{d\lambda }H\frac{dt}{d\lambda }$$
that corresponds to eq. 1. This equation links the tangent vector $`(\frac{dt}{d\lambda },\frac{dq}{d\lambda })`$ and the momentum vector $`(H,p).`$ It is equivalent to the condition $`\frac{dq}{dt}=\frac{p}{E}`$ that holds for the trajectory of a free particle. Interestingly, this condition translates to the requirement
$$x=\frac{1}{x_0}(x_0+x_1𝐞_1)(x_0+x_2𝐞_2)\text{ for }x_00$$
for the original vector $`x`$ in the $`D_2`$ group algebra; or, stated differently, it ensures that $`x`$ can be written as the tensor product of two vectors in the $`C_2`$ group algebra.
Transformations that preserve the norms of the tangent and momentum vectors and the action differential $`\frac{dS}{d\lambda }=m\frac{ds}{d\lambda }`$ can be induced by acting on $`x`$ with an element $`u=u_0+u_1𝐞_1+u_2𝐞_2+u_3𝐞_3`$ in the $`D_2`$ algebra and forming the product $`(xx)(uu)=xuxu.`$ For $`u`$ such that $`u=u_0+u_1𝐞_1`$ and $`u_0^2u_1^2=1,`$ the $`P_{}xx`$ subspace that contains the three norms is unchanged. The tangent and momentum vectors undergo a proper orthochronous transformation under this rule.
Finally, we also note that if instead of using $`D_2=C_2C_2`$ in the construction above, we use
$$C_2C_4=\{\mathrm{𝟏},𝐞_2\}\{\mathrm{𝟏},𝐞_1\}=\{\mathrm{𝟏},𝐞_1,𝐞_2,𝐞_{12}\}$$
where $`𝐞_{12}=𝐞_1𝐞_2,`$ $`𝐞_{21}=𝐞_{12},`$ $`𝐞_2^2=+1`$ and $`𝐞_1^2=𝐞_{12}^2=1`$ then we obtain the $`2d`$ Euclidean counterpart to the above $`1+1`$ spacetime case.
### 4 <br>Summary
We associate the vector $`xx`$ in the $`C_2C_2`$ group algebra with an element in a typical fiber residing at a point in a $`1+1`$ dimensional configurational manifold:
$`xx=`$ $`\{P_{+1}({\displaystyle \frac{dt}{d\lambda }}\mathrm{𝟏}+{\displaystyle \frac{dq}{d\lambda }}𝐞_1)+P_1{\displaystyle \frac{ds}{d\lambda }}\}(P_{+2}P_{+2})`$
$`+\{P_{+1}((H{\displaystyle \frac{dt}{d\lambda }})^{\frac{1}{2}}\mathrm{𝟏}+(p{\displaystyle \frac{dq}{d\lambda }})^{\frac{1}{2}}𝐞_1)`$
$`+P_1(m{\displaystyle \frac{ds}{d\lambda }})^{\frac{1}{2}}\}(P_{+2}P_2+P_2P_{+2})`$
$`+\{P_{+1}(H\mathrm{𝟏}+p𝐞_1)+P_1m\}(P_2P_2).`$
The collection of all such vectors $`xx`$ comprise the typical fiber. The $`(xx)(P_{+2}P_{+2})`$ portion of $`xx`$ is identified with the tangent vector and its norm. The $`(xx)(P_2P_2)`$ portion is identified with the momentum vector and its norm. The
$$(xx)(P_{+2}P_2+P_2P_{+2})$$
subspace is associated with the flow of action and its norm. The condition of minimum action translates into the condition that $`x`$ has the form
$$x=\frac{1}{x_0}(x_0+x_1𝐞_1)(x_0+x_2𝐞_2)\text{for}x_00.$$
It remains to determine how the fibers at different locations are connected. The description of this connection should also lead to a description of extended motions on this manifold.
Robert W. Johnson
878 Sunnyhills Road
Oakland, CA 94610 |
warning/0002/cond-mat0002323.html | ar5iv | text | # Computer simulation of uniformly heated granular fluids
## 1 Introduction
Most of the recent studies of rapid granular flow C90 are based on the Enskog equation for the velocity distribution function $`f(𝐫,𝐯,t)`$ of an assembly of inelastic hard spheres BDS97 . In the special case of a spatially uniform state, the Enskog equation reads
$$\frac{}{t}f(𝐯_1,t)+f(𝐯_1,t)=\chi I[𝐯_1|f(t),f(t)],$$
(1)
where
$`I[𝐯_1|f(t),f(t)]`$ $``$ $`\sigma ^{d1}{\displaystyle 𝑑𝐯_2𝑑\widehat{𝝈}\mathrm{\Theta }(𝐯_{12}\widehat{𝝈})(𝐯_{12}\widehat{𝝈})}`$
$`\times \left[\alpha ^2f(𝐯_1^{\prime \prime },t)f(𝐯_2^{\prime \prime },t)f(𝐯_1,t)f(𝐯_2,t)\right].`$
In Eq. (1) $``$ is an operator representing the effect of an external force (if it exists), $`\chi `$ is the pair correlation function at contact and $`I`$ is the collision operator. In Eq. (1), $`d`$ is the dimensionality of the system, $`\sigma `$ is the diameter of the spheres, $`𝐯_{12}𝐯_1𝐯_2`$ is the relative velocity of the colliding particles, $`\widehat{𝝈}`$ is a unit vector directed along the line of centers from the sphere $`1`$ to the sphere $`2`$, $`\mathrm{\Theta }`$ is the Heaviside step function and $`\alpha <1`$ is the coefficient of normal restitution, here assumed to be constant. In addition, $`(𝐯_1^{\prime \prime },𝐯_2^{\prime \prime })`$ are the precollisional velocities yielding $`(𝐯_1,𝐯_2)`$ as the postcollisional ones, i.e. $`𝐯_{1,2}^{\prime \prime }=𝐯_{1,2}\frac{1}{2}(1+\alpha ^1)(𝐯_{12}\widehat{𝝈})\widehat{𝝈}`$. Except for the presence of the factor $`\chi `$, which accounts for the increase of the collision frequency due to excluded volume effects, the Enskog equation for uniform states, Eq. (1), becomes identical with the Boltzmann equation.
In the case of elastic particles ($`\alpha =1`$) and in the absence of external forcing ($`=0`$), it is well known that the long-time solution of Eq. (1) is the Maxwell-Boltzmann equilibrium distribution function, $`f(𝐯,t)nv_0^d\varphi (𝐯/v_0)`$, $`\varphi (𝐜)\pi ^{d/2}e^{c^2}`$, where $`n`$ is the number density, $`v_0=(2v^2/d)^{1/2}`$ is the thermal velocity and $`𝐜=𝐯/v_0`$ is the reduced velocity. On the other hand, if the particles are inelastic ($`\alpha <1`$) and $`=0`$, a steady state is not possible in uniform situations since, due to the dissipation of energy through collisions, the thermal velocity $`v_0(t)`$ decreases monotonically with time. Regardless of the initial uniform state, the solution of Eq. (1) tends to the so-called homogeneous cooling state GS95 ; NE96 ; NE98 ; BMC96 , characterized by the fact that the time dependence occurs only through the thermal velocity $`v_0(t)`$: $`f(𝐯,t)nv_0^d(t)\stackrel{~}{f}(𝐯/v_0(t))`$. In addition, $`\stackrel{~}{f}(𝐜)`$ deviates from a Maxwellian, $`\stackrel{~}{f}(𝐜)\varphi (𝐜)`$, as measured by the fourth cumulant
$$a_2\frac{d}{d+2}\frac{v^4}{v^2^2}1=\frac{4}{d(d+2)}c^41,$$
(3)
where
$$c^p𝑑𝐜c^p\stackrel{~}{f}(𝐜).$$
(4)
By expanding $`\stackrel{~}{f}(𝐜)/\varphi (𝐜)`$ in a set of Sonine polynomials $`\{S_p(c^2)\}`$ and neglecting the terms beyond $`p=2`$, van Noije and Ernst NE96 ; NE98 have estimated the value of $`a_2`$:
$$a_2(\alpha )\frac{16(1\alpha )(12\alpha ^2)}{9+24d\alpha (418d)+30(1\alpha )\alpha ^2}.$$
(5)
The above expression corrects an algebraic error in a previous calculation of $`a_2`$ in the three-dimensional case GS95 . According to Eq. (5), $`a_2`$ changes sign at $`\alpha =1/\sqrt{2}0.71`$. By using the same method, Garzó and Dufty GD99 have recently extended the evaluation of $`a_2`$ to a binary mixture of hard spheres. The accuracy of Eq. (5) has been quantitatively confirmed by Brey et al. BMC96 from Monte Carlo simulations of the Boltzmann equation for hard spheres ($`d=3`$) in the range $`0.7\alpha 1`$. As a complementary measure of the departure of $`\stackrel{~}{f}(𝐜)`$ from $`\varphi (𝐜)`$, Esipov and Pöschel EP97 and van Noije and Ernst NE98 have analyzed the high energy tail of the distribution function and have found an asymptotic behavior of the form
$$\mathrm{log}\stackrel{~}{f}(𝐜)c,$$
(6)
in contrast to $`\mathrm{log}\varphi (𝐜)c^2`$. The high energy tail (6) has been confirmed by simulations in the case of hard disks ($`d=2`$) BCR99 .
In order to reach a steady state, energy injection is needed to compensate for the energy dissipated through collisions. This can be achieved by vibration of vessels ER89 or in fluidized beds IH95 . The same effect can be obtained by means of external driving forces acting locally on each particle WM96 . Borrowing a terminology frequently used in nonequilibrium molecular dynamics of elastic particles EM90 , we will call this type of external forces “thermostats”. In general, the equation of motion for a particle $`i`$ is then
$$m\dot{𝐯}_i=𝐅_i^{\text{coll}}+𝐅_i^{\text{th}},$$
(7)
where $`m`$ is the mass of a particle, $`𝐅_i^{\text{coll}}`$ is the force due to collisions and $`𝐅_i^{\text{th}}`$ is the thermostat force. Williams and MacKintosh WM96 introduced a stochastic force assumed to have the form of a Gaussian white noise:
$$𝐅_i^{\text{th}}(t)=\mathrm{𝟎},𝐅_i^{\text{th}}(t)𝐅_j^{\text{th}}(t^{})=𝖨m^2\xi _0^2\delta _{ij}\delta (tt^{}),$$
(8)
where $`𝖨`$ is the $`d\times d`$ unit matrix and $`\xi _0^2`$ represents the strength of the correlation. The corresponding operator $``$ appearing in Eq. (1) has a Fokker-Planck form NE98 :
$$f(𝐯_1)=\frac{\xi _0^2}{2}\left(\frac{}{𝐯_1}\right)^2f(𝐯_1).$$
(9)
Van Noije and Ernst NE98 have studied the stationary solution of the uniform equation (1) with the thermostat (9). They have found for the coefficient $`a_2`$ defined by Eq. (3) the value
$$a_2(\alpha )\frac{16(1\alpha )(12\alpha ^2)}{73+56d3\alpha (35+8d)+30(1\alpha )\alpha ^2}.$$
(10)
The high energy tail is NE98
$$\mathrm{log}\stackrel{~}{f}(𝐜)c^{3/2}.$$
(11)
Of course, deterministic thermostats can also be used. For instance, the use of Gauss’s principle of least constraint leads to the thermostat force EM90
$$𝐅_i^{\text{th}}=m\zeta 𝐯_i,$$
(12)
where $`\zeta `$ is a positive constant. In this case,
$$f(𝐯_1)=\zeta \frac{}{𝐯_1}\left[𝐯_1f(𝐯_1)\right].$$
(13)
It is interesting pointing out that the Enskog-Boltzmann equation (1) for the above Gaussian thermostat force is formally identical with the equation for the homogeneous cooling state (i.e. with $`=0`$) when both equations are expressed in terms of the reduced distribution $`\stackrel{~}{f}(𝐜)`$ (see Sect. 2). As a consequence, the results (5) and (6) apply to this thermostatted case as well.
The differences between Eqs. (5) and (10) and between (6) and (11) illustrate the influence of the thermostat force on the departure of the steady-state distribution function from the Maxwell-Boltzmann distribution. In the case of the stochastic force, Eq. (9), $`\stackrel{~}{f}(𝐜)`$ is closer to $`\varphi (𝐜)`$ than in the case of the Gaussian force, Eq. (13), since in the former case $`a_2`$ and the high energy overpopulation are smaller than in the latter case. Of course, other types of thermostats are also possible. For example, a different choice for a deterministic thermostat is
$$𝐅_i^{\text{th}}=mg\widehat{𝐯}_i,$$
(14)
where $`\widehat{𝐯}_i𝐯_i/v_i`$. While the Gaussian force, Eq. (12), is proportional to the velocity of the particle, Eq. (14) corresponds to a force that is parallel to the direction of motion but constant in magnitude. The corresponding operator $``$ is
$$f(𝐯_1)=\frac{g}{v_1}\left\{\frac{}{𝐯_1}\left[𝐯_1f(𝐯_1)\right]f(𝐯_1)\right\}.$$
(15)
The aim of this paper is to present direct Monte Carlo simulations of Eq. (1) with the three choices for the thermostat, Eqs. (9), (13) and (15). In the cases of the stochastic and the Gaussian thermostats, we will confirm the tails (11) and (6) and will check the accuracy of the estimates (10) and (5). In the latter case, however, we will see that a better agreement with simulation results for $`\alpha <0.5`$ is obtained if an estimate slightly different from (5) is used. The simulation results corresponding to the non-Gaussian thermostat (15) show that, in contrast to what happens in the two previous cases, $`a_2`$ remains negative for all $`\alpha `$. This feature is qualitatively captured by an estimate derived from a Sonine approximation. In this problem, however, the Sonine polynomials do not constitute a good set for the expansion of $`\stackrel{~}{f}(𝐜)/\varphi (𝐜)`$ and, consequently, the estimate is not quantitatively good. Besides, the high energy tail is of the form $`\mathrm{log}\stackrel{~}{f}(𝐜)c^2`$, but with a coefficient different from that of the Maxwell-Boltzmann distribution.
The organization of this paper is as follows. The theoretical analysis is reviewed in Sect. 2. The computer simulation method employed to solve numerically the uniform Enskog-Boltzmann equation is described in Sect. 3. The results are presented and compared with the theoretical predictions in Sect. 4. The paper ends with a summary and discussion in Sect. 5.
## 2 Theoretical predictions
In the steady state, the Enskog-Boltzmann equation (1) can be expressed in terms of the reduced velocity distribution function $`\stackrel{~}{f}(𝐜)`$ as
$$\stackrel{~}{}\stackrel{~}{f}(𝐜_1)=\chi \stackrel{~}{I}[𝐜_1|\stackrel{~}{f},\stackrel{~}{f}],$$
(16)
where
$`\stackrel{~}{I}[𝐜_1|\stackrel{~}{f},\stackrel{~}{f}]`$ $``$ $`{\displaystyle 𝑑𝐜_2𝑑\widehat{𝝈}\mathrm{\Theta }(𝐜_{12}\widehat{𝝈})(𝐜_{12}\widehat{𝝈})}`$ (17)
$`\times \left[\alpha ^2\stackrel{~}{f}(𝐜_1^{\prime \prime })\stackrel{~}{f}(𝐜_2^{\prime \prime })\stackrel{~}{f}(𝐜_1)\stackrel{~}{f}(𝐜_2)\right].`$
The reduced operator $`\stackrel{~}{}`$ for the stochastic \[Eq. (9)\], Gaussian \[Eq. (13)\] and non-Gaussian \[Eq. (15)\] thermostats, is
$$\stackrel{~}{}\stackrel{~}{f}(𝐜_1)=\frac{\xi _0^2}{2v_0^3n\sigma ^{d1}}c_1^{(d1)}\frac{}{c_1}\left[c_1^{d1}\frac{\stackrel{~}{f}(𝐜_1)}{c_1}\right],$$
(18)
$$\stackrel{~}{}\stackrel{~}{f}(𝐜_1)=\frac{\zeta }{v_0n\sigma ^{d1}}c_1^{(d1)}\frac{}{c_1}\left[c_1^d\stackrel{~}{f}(𝐜_1)\right],$$
(19)
$$\stackrel{~}{}\stackrel{~}{f}(𝐜_1)=\frac{g}{v_0^2n\sigma ^{d1}}c_1^{(d1)}\frac{}{c_1}\left[c_1^{d1}\stackrel{~}{f}(𝐜_1)\right],$$
(20)
respectively. In Eqs. (18)–(20) we have already taken into account that the distribution function must be isotropic in the steady state. Equation (16) with the term (19) is fully equivalent to Eq. (10) of Ref. NE98 , the latter being derived in the context of the homogeneous cooling state. This formal equivalence between the free evolving state and the one controlled by a Gaussian external force is also present in the case of elastic particles interacting via arbitrary power-law potentials in homogeneous situations GSB90 or via the Maxwell potential in the uniform shear flow DSBR86 .
### 2.1 Stochastic thermostat
For the sake of completeness, we summarize now some of the results obtained in Ref. NE98 . In order to characterize the deviation of $`\stackrel{~}{f}(𝐜)`$ from $`\varphi (𝐜)`$ by means of the cumulant (3), it is useful to consider the hierarchy of moment equations. Multiplying both sides of Eq. (16) by $`c_1^p`$ and integrating over $`𝐜_1`$, we get
$$\mu _p=\mu _2\frac{p(p+d2)}{2d}c^{p2}$$
(21)
for the stochastic thermostat, where we have defined
$$\mu _p𝑑𝐜c^p\stackrel{~}{I}[𝐜|\stackrel{~}{f},\stackrel{~}{f}].$$
(22)
In Eq. (21) we have taken into account the normalization condition $`c^0=1`$, so that $`\mu _2=d\xi _0^2/v_0^3\chi n\sigma ^{d1}`$. In the special case of $`p=4`$, Eq. (21) becomes
$$\mu _4=(d+2)\mu _2,$$
(23)
where we have used the fact that, by definition, $`c^2=d/2`$. Equations (22) and (23) are still exact. To get an approximate expression for $`a_2(\alpha )`$, three steps will be taken NE98 . First, we assume that $`\stackrel{~}{f}`$ can be well described by the simplest Sonine approximation, at least for the velocities relevant to the evaluation of $`a_2`$. Thus,
$$\stackrel{~}{f}(𝐜)\varphi (𝐜)\left[1+a_2S_2(c^2)\right],$$
(24)
where
$$S_2(x)=\frac{1}{2}x^2\frac{d+2}{2}x+\frac{d(d+2)}{8}.$$
(25)
The approximation (24) is justified by the fact that $`a_2`$ is expected to be small. The second step consists of inserting Eq. (24) into Eq. (22) and neglecting terms nonlinear in $`a_2`$. For $`\mu _2`$ and $`\mu _4`$ the results are NE98
$$\mu _p\mu _p^{(0)}+\mu _p^{(1)}a_2,$$
(26)
with
$$\mu _2^{(0)}\frac{\pi ^{(d1)/2}}{\sqrt{2}\mathrm{\Gamma }(d/2)}(1\alpha ^2),$$
(27)
$$\mu _2^{(1)}\frac{3}{16}\mu _2^{(0)},$$
(28)
$$\mu _4^{(0)}\left(d+\frac{3}{2}+\alpha ^2\right)\mu _2^{(0)},$$
(29)
$$\mu _4^{(1)}\left[\frac{3}{32}\left(10d+39+10\alpha ^2\right)+\frac{d1}{1\alpha }\right]\mu _2^{(0)}.$$
(30)
In the third step, the approximations (26) with $`p=2\text{ and }4`$ are inserted into the exact equation (23) and $`a_2`$ is obtained from the resulting linear equation:
$$a_2\frac{\mu _4^{(0)}(d+2)\mu _2^{(0)}}{\mu _4^{(1)}(d+2)\mu _2^{(1)}}.$$
(31)
This is the result derived by van Noije and Ernst NE98 , Eq. (10). It must be pointed out that a certain degree of ambiguity is present in this last step. For instance, if Eq. (23) were written as $`\mu _4/\mu _2=d+2`$, we could expand the ratio $`\mu _4/\mu _2`$ in powers of $`a_2`$ and neglect nonlinear terms to find
$`a_2`$ $``$ $`{\displaystyle \frac{\mu _4^{(0)}(d+2)\mu _2^{(0)}}{\mu _4^{(1)}\mu _2^{(1)}\mu _4^{(0)}/\mu _2^{(0)}}}`$ (32)
$`=`$ $`{\displaystyle \frac{4(1\alpha )(12\alpha ^2)}{19+14d3\alpha (9+2d)+6(1\alpha )\alpha ^2}}.`$
However, since $`a_2`$ is indeed small ($`|a_2|<0.1`$), Eqs. (10) and (32) give practically identical results, the maximum deviation being less than about 0.001.
Now we consider the high energy tail. In general, the collision integral can be decomposed into a gain and a loss term: $`\stackrel{~}{I}[𝐜_1|\stackrel{~}{f},\stackrel{~}{f}]=\stackrel{~}{I}_g[𝐜_1|\stackrel{~}{f},\stackrel{~}{f}]\stackrel{~}{I}_l[𝐜_1|\stackrel{~}{f},\stackrel{~}{f}]`$. For large $`c_1`$ the loss term can be approximated as
$`\stackrel{~}{I}_l[𝐜_1|\stackrel{~}{f},\stackrel{~}{f}]`$ $`=`$ $`\beta _1{\displaystyle 𝑑𝐜_2c_{12}\stackrel{~}{f}(𝐜_1)\stackrel{~}{f}(𝐜_2)}`$ (33)
$``$ $`\beta _1c_1\stackrel{~}{f}(𝐜_1),`$
where $`\beta _1`$ is defined by Eq. (84). Let us assume that for large velocities the gain term is negligible versus the loss term, i.e.
$$\underset{c_1\mathrm{}}{lim}\frac{\stackrel{~}{I}_g[𝐜_1|\stackrel{~}{f},\stackrel{~}{f}]}{\stackrel{~}{I}_l[𝐜_1|\stackrel{~}{f},\stackrel{~}{f}]}=0.$$
(34)
In that case, the Enskog-Boltzmann equation for the stochastic thermostat becomes
$$\frac{\mu _2}{2d}c^{(d1)}\frac{}{c}\left[c^{d1}\frac{\stackrel{~}{f}(𝐜)}{c}\right]\beta _1c\stackrel{~}{f}(𝐜).$$
(35)
The solution of this equation for large $`c`$ is
$$\stackrel{~}{f}(𝐜)K\mathrm{exp}\left(Ac^{3/2}\right),A\frac{2}{3}\left(\frac{2d\beta _1}{\mu _2}\right)^{1/2},$$
(36)
where $`K`$ is an undetermined constant. By arguments given in Ref. NE98 , it can be seen that the result (36) is indeed consistent with the assumption (34). Equation (36) shows an overpopulation with respect to the Maxwell-Boltzmann tail. On the other hand, as $`\alpha 1`$, the amplitude $`A`$ diverges as $`(1\alpha )^{1/2}`$, thus indicating that the overpopulation effect is restricted to larger and larger energies in the limit $`\alpha 1`$.
### 2.2 Gaussian thermostat
In the case of the deterministic Gaussian thermostat, Eq. (19), the moment equation is
$$\mu _p=\mu _2\frac{p}{d}c^p,$$
(37)
where now $`\mu _2=d\zeta /v_0\chi n\sigma ^{d1}`$. If we set $`p=4`$,
$$\mu _4=(d+2)(1+a_2)\mu _2,$$
(38)
where we have made use of Eq. (3). Substituting the approximation (26) and neglecting terms nonlinear in $`a_2`$, we get
$$a_2\frac{\mu _4^{(0)}(d+2)\mu _2^{(0)}}{\mu _4^{(1)}(d+2)(\mu _2^{(1)}+\mu _2^{(0)})},$$
(39)
which is the same as Eq. (5). There exists again some arbitrariness about the use of the exact equation (38) in connection with the approximation (26). If we rewrite (38) as $`\mu _4/\mu _2=(d+2)(1+a_2)`$ and neglect nonlinear terms, the resulting $`a_2`$ is fairly close to Eq. (39). On the other hand, if we start from $`\mu _4/(1+a_2)=(d+2)\mu _2`$, the result is
$`a_2`$ $``$ $`{\displaystyle \frac{\mu _4^{(0)}(d+2)\mu _2^{(0)}}{\mu _4^{(1)}\mu _4^{(0)}(d+2)\mu _2^{(1)}}}`$ (40)
$`=`$ $`{\displaystyle \frac{16(1\alpha )(12\alpha ^2)}{25+24d\alpha (578d)2(1\alpha )\alpha ^2}}.`$
The estimates (5) and (40) practically coincide in the region $`0.5<\alpha 1`$. However, they visibly separate for larger dissipation. In the interval $`0<\alpha <0.3`$, the values given by Eq. (5) are 12%–20% ($`d=3`$) or 18%–28% ($`d=2`$) larger than those given by Eq. (40). As we will see later, the simulation results indicate that Eq. (40) is a better estimate than Eq. (5).
For large $`c`$ the Enskog-Boltzmann equation becomes
$$\frac{\mu _2}{d}c^{(d1)}\frac{}{c}\left[c^d\stackrel{~}{f}(𝐜)\right]\beta _1c\stackrel{~}{f}(𝐜),$$
(41)
where we have used Eqs. (33) and (34). Its solution is
$$\stackrel{~}{f}(𝐜)K\mathrm{exp}\left(Ac\right),A\frac{d\beta _1}{\mu _2}.$$
(42)
Again, this result is seen to be consistent with (34) NE98 . Equation (42) indicates an overpopulation effect even larger than with the stochastic thermostat.
### 2.3 Non-Gaussian thermostat
Now we consider the deterministic non-Gaussian thermostat (14), represented by the operator (20). To the best of our knowledge, this external force has not been analyzed before. The corresponding moment equation is
$$\mu _p=\mu _2\frac{p}{2}\frac{c^{p1}}{c},$$
(43)
where $`\mu _2=2gc/v_0^2\chi n\sigma ^{d1}`$. In particular,
$$\mu _4=2\mu _2\frac{c^3}{c}.$$
(44)
In contrast to the two previous cases, now the even collisional moments $`\mu _p`$ are coupled to the odd moments $`c^{p1}`$, and vice versa. In terms of the energy variable $`ϵ=c^2`$, this means that the integer collisional moments are coupled to the half-integers energy moments. This is related to the fact that the force (14) is singular at $`ϵ=0`$. As a consequence, while $`\stackrel{~}{f}(𝐜)`$ is expected to be close to the Maxwellian $`\varphi (𝐜)`$, the ratio $`\stackrel{~}{f}(𝐜)/\varphi (𝐜)`$ is singular at $`ϵ=0`$ and thus it is not well represented by an expansion in $`\{S_p(ϵ)\}`$. To be more precise, let us define the function $`\mathrm{\Delta }(c)`$ by the equation
$$\stackrel{~}{f}(𝐜)=\varphi (𝐜)\left[1+a_2\mathrm{\Delta }(c)\right].$$
(45)
Therefore,
$$𝑑𝐜\varphi (𝐜)\mathrm{\Delta }(c)=𝑑𝐜\varphi (𝐜)c^2\mathrm{\Delta }(c)=0,$$
(46)
$$𝑑𝐜\varphi (𝐜)c^4\mathrm{\Delta }(c)=\frac{d(d+2)}{4}.$$
(47)
The polynomial $`S_2(c^2)`$ verifies the above equalities. As a matter of fact, $`\mathrm{\Delta }(c)S_2(c^2)`$ in the cases of the thermostats (18) and (19). This is not so, however, in the case of (20), even in the limit of low dissipation. As we will see in Sect. 4, $`\mathrm{\Delta }(c)/c|_{c=0}0`$, what indicates that $`\mathrm{\Delta }(c)`$ is essentially different from a polynomial in $`c^2`$. All of this complicates the evaluation of $`a_2`$. Nevertheless, since $`\mathrm{\Delta }(c)`$ and $`S_2(c^2)`$ share the moments of degrees 0, 2 and 4 \[cf. Eqs. (46) and (47)\], we can expect to obtain a crude estimate of $`a_2`$ by assuming that in the calculation of $`c`$, $`c^3`$, $`\mu _2`$ and $`\mu _4`$ we can replace $`\mathrm{\Delta }(c)`$ by $`S_2(c^2)`$. If that were the case, $`\mu _2`$ and $`\mu _4`$ would be given by Eqs. (26)–(30) and
$$c\frac{\mathrm{\Gamma }((d+1)/2)}{\mathrm{\Gamma }(d/2)}\left(1\frac{1}{8}a_2\right),$$
(48)
$$c^3\frac{\mathrm{\Gamma }((d+3)/2)}{\mathrm{\Gamma }(d/2)}\left(1+\frac{3}{8}a_2\right).$$
(49)
Inserting this into Eq. (44) and neglecting nonlinear terms, we get
$`a_2`$ $``$ $`{\displaystyle \frac{\mu _4^{(0)}(d+1)\mu _2^{(0)}}{\mu _4^{(1)}(d+1)(\mu _2^{(1)}+\mu _2^{(0)}/2)}}`$ (50)
$`=`$ $`{\displaystyle \frac{16(1\alpha )(1+2\alpha ^2)}{63+40d\alpha (95+8d)+30(1\alpha )\alpha ^2}}.`$
While in the cases of the stochastic thermostat, Eqs. (10) or (32), and the Gaussian thermostat, Eqs. (5) or (40), the cumulant $`a_2`$ changes from negative to positive values at $`\alpha 0.71`$, Eq. (50) indicates that $`a_2`$ remains negative in the case of the non-Gaussian thermostat. We will see in Sect. 4 that our computer simulations confirm this feature. At a quantitative level, however, the estimate (50) is about 20% too small in magnitude.
To analyze the high energy tail, let us assume for the moment the validity of (34), so that the Enskog-Boltzmann equation can be replaced by
$$\frac{\mu _2}{2c}c^{(d1)}\frac{}{c}\left[c^{d1}\stackrel{~}{f}(𝐜)\right]\beta _1c\stackrel{~}{f}(𝐜),$$
(51)
whose solution for large $`c`$ is
$$\stackrel{~}{f}(𝐜)K\mathrm{exp}\left(A^{}c^2\right),A^{}\frac{\beta _1c}{\mu _2}.$$
(52)
According to (52), $`\stackrel{~}{f}(𝐜)`$ has a Maxwellian tail that is underpopulated with respect to the Maxwell-Boltzmann distribution $`\varphi (𝐜)`$, since the amplitude $`A^{}\sqrt{2}/(1\alpha ^2)`$ is larger than 1. But now we get an unphysical result: the underpopulation effect increases as one approaches the elastic limit, since $`A^{}\mathrm{}`$ as $`\alpha 1`$. The solution to this paradox lies in the fact that the assumption (34) is not justified in this case. Let us assume instead that the gain and loss term are comparable, namely
$$\stackrel{~}{I}_g[𝐜_1|\stackrel{~}{f},\stackrel{~}{f}]\gamma \beta _1c_1\stackrel{~}{f}(𝐜_1),$$
(53)
where $`\gamma <1`$ is an unknown function of $`\alpha `$. According to this, Eq. (52) is replaced by
$$\stackrel{~}{f}(𝐜)K\mathrm{exp}\left(Ac^2\right),AA^{}(1\gamma ).$$
(54)
On physical grounds we expect that $`A1`$ when $`\alpha 1`$, which implies that $`\gamma 1\sqrt{2}(1\alpha )`$ in that limit. As will be shown in Sect. 4, comparison with simulation results confirms a behavior of the form (54).
## 3 Direct Simulation Monte Carlo method
The Direct Simulation Monte Carlo (DSMC) method devised by Bird Bird has proven to be a very efficient tool to solve numerically the Boltzmann equation. The DSMC method has been recently extended to the Enskog equation MS96 and its application to inelastic particles is straightforward BMC96 ; MGSB99 . Here we briefly describe the specific method we have used to solve the uniform Enskog-Boltzmann equation (1) in the case of a three-dimensional system ($`d=3`$).
The velocity distribution function is represented by the velocities $`\{𝐯_i\}`$ of $`N`$ “simulated” particles:
$$f(𝐯,t)n\frac{1}{N}\underset{i=1}{\overset{N}{}}\delta \left(𝐯_i(t)𝐯\right).$$
(55)
At the initial state the particles are assigned velocities drawn from a Maxwell-Boltzmann probability distribution:
$$n^1f(𝐯,0)=\pi ^{3/2}v_0^3(0)e^{v^2/v_0^2(0)},$$
(56)
where $`v_0(0)`$ is an arbitrary initial thermal velocity. To enforce a vanishing initial total momentum, the velocity of every particle is subsequently subtracted by the amount $`N^1_i𝐯_i(0)`$.
The velocities are updated from time $`t`$ to time $`t+h`$, where the time step $`h`$ is much smaller than the mean free time, by following two successive stages: collisions and free streaming. In the collision stage, a sample of $`\frac{1}{2}N\omega _{\text{max}}h`$ pairs is chosen at random with equiprobability, where $`\omega _{\text{max}}`$ is an upper bound estimate of the probability that a particle collides per unit of time. For each pair $`ij`$ belonging to this sample, the following steps are taken: (1) a given direction $`\widehat{𝝈}_{ij}`$ is chosen at random with equiprobability; (2) the collision between particles $`i`$ and $`j`$ is accepted with a probability equal to $`\mathrm{\Theta }(𝐯_{ij}\widehat{𝝈}_{ij})\omega _{ij}/\omega _{\text{max}}`$, where $`\omega _{ij}=(4\pi \sigma ^2\chi n)|𝐯_{ij}\widehat{𝝈}_{ij}|`$; if the collision is accepted, postcollisional velocities are assigned to both particles: $`𝐯_{i,j}𝐯_{i,j}\frac{1}{2}(1+\alpha )(𝐯_{ij}\widehat{𝝈}_{ij})\widehat{𝝈}_{ij}`$. In the case that in one of the collisions $`\omega _{ij}>\omega _{\text{max}}`$, the estimate of $`\omega _{\text{max}}`$ is updated as $`\omega _{\text{max}}=\omega _{ij}`$.
In the free streaming stage the velocity of every particle is changed according to the thermostat force under consideration:
$$𝐯_i𝐯_i+𝐰_i,$$
(57)
where
$$𝐰_i=\frac{1}{m}_t^{t+h}𝑑t^{}𝐅_i^{\text{th}}(t^{}).$$
(58)
In the case of the stochastic thermostat, Eq. (8), one has
$$𝐰_i=\mathrm{𝟎},𝐰_i𝐰_j=𝖨\xi _0^2h\delta _{ij}.$$
(59)
Consequently, each vector $`𝐰_i`$ is randomly drawn from the Gaussian probability distribution
$$P(𝐰)=\left(2\pi \xi _0^2h\right)^{3/2}e^{w^2/2\xi _0^2h}.$$
(60)
In the case of deterministic external forces the velocity increment $`𝐰_i`$ is assigned in a more direct way. If the thermostat is the Gaussian one, Eq. (12),
$$𝐰_i=\left(e^{\zeta h}1\right)𝐯_i.$$
(61)
In the case of the non-Gaussian thermostat defined by Eq. (14),
$$𝐰_i=gh\left(\widehat{𝐯}_i𝐤\right),$$
(62)
where the vector $`𝐤N^1_i\widehat{𝐯}_i`$ is introduced to preserve the detailed conservation of momentum, i.e. $`_i𝐰_i=\mathrm{𝟎}`$.
The moments of the distribution are simply obtained as
$$v^p=\frac{1}{N}\underset{i=1}{\overset{N}{}}v_i^p,c^p=v^p/v_0^p,$$
(63)
where $`v_0=\left(2v^2/3\right)^{1/2}`$. The evaluation of the collisional moments $`\mu _p`$, $`p=2\text{ and }4`$, is more complicated. In the Appendix it is shown that
$$\mu _p=n^2v_0^{p1}𝑑𝐯_1𝑑𝐯_2f(𝐯_1)f(𝐯_2)\mathrm{\Phi }_p(𝐯_1,𝐯_2),$$
(64)
where
$$\mathrm{\Phi }_2(𝐯_1,𝐯_2)=\frac{\pi (1\alpha ^2)}{8}v_{12}^3,$$
(65)
$`\mathrm{\Phi }_4(𝐯_1,𝐯_2)`$ $`=`$ $`{\displaystyle \frac{\pi }{4}}v_{12}\{{\displaystyle \frac{5(1\alpha ^2)}{3}}v_{12}^2V_{12}^2`$
$`+{\displaystyle \frac{(1\alpha ^2)(2+\alpha ^2)}{12}}v_{12}^4`$
$`+(3\alpha )(1+\alpha )[(𝐯_{12}𝐕_{12})^2{\displaystyle \frac{1}{3}}v_{12}^2V_{12}^2]\}.`$
In the above equations, $`𝐯_{12}𝐯_1𝐯_2`$, $`𝐕_{12}\frac{1}{2}(𝐯_1+𝐯_2)`$. Starting from the exact expression (64) and using (55), we arrive at the following formula for the numerical computation of $`\mu _p`$:
$$\mu _p=v_0^{p1}\frac{1}{N^{}}\underset{ij}{}^{}\mathrm{\Phi }_p(𝐯_i,𝐯_j).$$
(67)
The prime in the summation means that we restrict ourselves to $`N^{}`$ pairs $`ij`$ randomly chosen out of the total number $`N(N1)/2`$ of pairs in the system. This allows us to compute $`c^p`$ and $`\mu _p`$ with similar accuracy within reasonable computer times. Once the steady state is reached, the relevant quantities are subsequently averaged over $`M`$ independent instantaneous values.
In our simulations we have typically taken $`N=2\times 10^5`$, $`N^{}=10^7`$ and $`M=10^3`$. Since the thermal velocity is not constant in the transient regime, we have taken a time-dependent time step $`h=0.01\lambda /v_0(t)`$, where $`\lambda =(\sqrt{2}\pi \chi n\sigma ^2)^1`$ is the mean free path.
## 4 Results
By using the numerical method described in the previous section, we have computed the steady-state values of the first few moments $`c^p`$ and $`\mu _p`$. We have also evaluated the reduced velocity distribution function $`\stackrel{~}{f}(𝐜)`$. As a test of the accuracy of the simulations and also to check that the steady state has been reached, we compare in Table 1 the values of $`\mu _4`$ obtained directly from Eq. (67) with those given by Eqs. (23), (38) or (44). The values corresponding to a Maxwell-Boltzmann distribution, $`\mu _4^{(0)}`$, are also included in the table. We can observe that the direct and indirect routes to the computation of $`\mu _4`$ disagree less than $`0.1\%`$ in all the cases. The difference between $`\mu _4`$ and $`\mu _4^{(0)}`$ is a measure of the departure of $`\stackrel{~}{f}(𝐜)`$ from $`\varphi (𝐜)`$.
Now we present the results separately for each one of the three thermostats considered.
### 4.1 Stochastic thermostat
The basic quantity measuring the deviation of the distribution function from the Maxwell-Boltzmann distribution is the cumulant $`a_2`$, Eq. (3). Figure 1 shows the $`\alpha `$-dependence of the simulation values of $`a_2`$, $`(\mu _2\mu _2^{(0)})/\mu _2^{(1)}`$, $`(\mu _4\mu _4^{(0)})/\mu _4^{(1)}`$ and the theoretical estimate (10), first derived in Ref. NE98 . As said in Sect. 2, the estimate (32) gives practically the same results as (10) and therefore it is not plotted. The agreement between the simulation data and the theoretical prediction is excellent, thus indicating that the approximation (26) was justified.
The above agreement indicates that the distribution function $`\stackrel{~}{f}(𝐜)`$ for thermal velocities is well represented by Eq. (24). To confirm this, the function $`\mathrm{\Delta }(c)`$ defined by Eq. (45) is plotted in Fig. 2 for $`\alpha =0.5`$. The simulation curve agrees very well with the Sonine polynomial $`S_2(c^2)`$.
It is worth noting that this deviation from the Maxwell-Boltzmann distribution in the case of the stochastic thermostat could not be observed in recent two-dimensional molecular dynamics simulations PO98 because the statistical accuracy was not high enough.
The theoretical prediction for the asymptotic high energy tail, Eq. (36), is much harder to confirm in the simulations since it involves a very small fraction of particles. Equation (36) implies that
$$\underset{c\mathrm{}}{lim}G(c)=K=\text{const},$$
(68)
where
$$G(c)e^{Ac^{3/2}}\stackrel{~}{f}(𝐜).$$
(69)
The function $`G(c)`$ is plotted (in logarithmic scale) in Fig. 3 for $`\alpha =0.4`$ and $`\alpha =0.5`$. In both cases the values of $`A`$ have been obtained from (36) by using the simulation values of $`\mu _2`$, which yields $`A1.99`$ ($`\alpha =0.4`$) and $`A2.11`$ ($`\alpha =0.5`$). The figure is convincingly consistent with Eq. (68), where $`K1.3`$ and $`K2.2`$ for $`\alpha =0.4`$ and $`\alpha =0.5`$, respectively. Figure 3 also shows the corresponding functions $`G(c)`$ obtained from Eq. (69) by replacing $`\stackrel{~}{f}(𝐜)`$ by the Maxwell-Boltzmann distribution $`\varphi (𝐜)`$. The overpopulation phenomenon for $`c>2`$ is quite apparent. At $`c=4`$, for instance, $`\stackrel{~}{f}/\varphi 8`$ for $`\alpha =0.4`$ and $`\stackrel{~}{f}/\varphi 5`$ for $`\alpha =0.5`$.
### 4.2 Gaussian thermostat
Now we carry out a parallel analysis in the case of the deterministic Gaussian thermostat. The $`\alpha `$-dependence of the simulation values of $`a_2`$, $`(\mu _2\mu _2^{(0)})/\mu _2^{(1)}`$ and $`(\mu _4\mu _4^{(0)})/\mu _4^{(1)}`$ are shown in Figure 4. The values of $`a_2`$ are in this case generally larger than in the previous case. In addition, Eq. (26) tends to overestimate $`\mu _4`$ and underestimate $`\mu _2`$ for small $`\alpha `$. As a consequence, the theoretical estimate (5) gives values larger than the simulation data for $`\alpha <0.5`$, while the estimate (40) is fairly good in that region.
For values of the coefficient of restitution for which the fourth cumulant $`a_2`$ is not small enough (say $`a_20.1`$), we may expect a non-negligible deviation from (24). This is confirmed in Fig. 5, where $`\mathrm{\Delta }(c)`$ is plotted for $`\alpha =0.4`$. Here the contributions associated with higher-order Sonine polynomials are relatively important.
As a quantitative measure of the difference between $`\mathrm{\Delta }(c)`$ and $`S_2(c^2)`$, we have obtained preliminary simulation results for the sixth cumulant $`a_3`$ defined as
$`a_3`$ $``$ $`{\displaystyle \frac{48}{d(d+2)(d+4)}}{\displaystyle 𝑑𝐜S_3(c^2)\stackrel{~}{f}(𝐜)}`$ (70)
$`=`$ $`{\displaystyle \frac{8}{d(d+2)(d+4)}}c^6+1+3a_2.`$
This quantity is plotted in Fig. 6 for $`d=3`$.
For $`\alpha 0.6`$ $`|a_3|`$ remains small, but for larger dissipation the values of $`|a_3|`$ increase rapidly.
The high energy tail predicted by Eq. (42) NE98 ; EP97 is tested in Fig. 7, where
$$G(c)e^{Ac}\stackrel{~}{f}(𝐜)$$
(71)
is plotted for $`\alpha =0.2`$ and $`\alpha =0.4`$. The corresponding values of $`A`$ are $`A3.82`$ and $`A4.41`$, respectively. The agreement with Eq. (68) is excellent; from the simulation data we can estimate $`K7`$ for $`\alpha =0.2`$ and $`K31`$ for $`\alpha =0.4`$.
In this case of a Gaussian thermostat, the overpopulation effect is much more important than in the previous case. At $`c=4`$, $`\stackrel{~}{f}/\varphi 80`$ for $`\alpha =0.2`$ and $`\stackrel{~}{f}/\varphi 34`$ for $`\alpha =0.4`$. The results reported here for inelastic hard spheres complement those obtained by Brey et al. BCR99 , where the asymptotic behavior (42) was verified for inelastic hard disks.
Recently, Sela and Goldhirsch SG98 have obtained numerically the function $`\mathrm{\Delta }(c)`$ in the low dissipation limit. In their notation, $`lim_{\alpha 1}\mathrm{\Delta }(c)8\widehat{\mathrm{\Phi }}_ϵ(c)`$. From simulation results presented in Ref. BMC96 for $`\alpha =0.99`$ it follows that the function $`\widehat{\mathrm{\Phi }}_ϵ(c)`$ is well represented by the Sonine polynomial $`S_2(c^2)/8`$ in the range $`0c1`$. However, this agrees only qualitatively with the function $`\widehat{\mathrm{\Phi }}_ϵ(c)`$ obtained numerically by Sela and Goldhirsch SG98 . For instance, from Fig. 3 of Ref. SG98 one gets $`\widehat{\mathrm{\Phi }}_ϵ(0)0.35`$, while $`S_2(0)/8=15/640.23`$. Moreover, it is claimed in Ref. SG98 that $`\widehat{\mathrm{\Phi }}_ϵ(c)c^2\mathrm{log}c`$ for large $`c`$, which differs from the behavior (42) that has been confirmed here and in Ref. BCR99 . It is possible that the high energy tail obtained from the perturbative approach presented in Ref. SG98 only holds for $`1c(1\alpha ^2)^1`$ and thus it is not representative of the general asymptotic behavior for arbitrary $`\alpha `$.
### 4.3 Non-Gaussian thermostat
In contrast to the two previous cases, the Sonine polynomials $`\{S_p(c^2)\}`$ are not expected to constitute a good set for the expansion of the ratio $`\stackrel{~}{f}(𝐜)/\varphi (𝐜)`$ in the case of the non-Gaussian thermostat (15) since the latter is singular at $`𝐜=\mathrm{𝟎}`$. Consequently, we do not expect the estimate (50) to be quantitatively accurate. This is confirmed in Fig. 8, where we observe that Eq. (50) gives values that are about 20% smaller in magnitude than the simulation ones.
Also, the approximation (26) with $`\mu _{2,4}^{(1)}`$ given by Eqs. (28) and (30) is rather poor. It is reasonable to expect that a better approximation would be obtained if $`c`$, $`c^3`$, $`\mu _2`$ and $`\mu _4`$ were computed from the unknown function $`\mathrm{\Delta }(c)`$ rather than from $`S_2(c^2)`$. When plotting the simulation data of $`(\mu _2\mu _2^{(0)})/\mu _2^{(1)}`$, $`(\mu _4\mu _4^{(0)})/\mu _4^{(1)}`$, $`c`$ and $`c^3`$ versus $`a_2`$ we have observed that the points fit well in straight lines, as predicted by Eqs. (26), (48) and (49), but with different slopes. More specifically, our simulation results indicate that, instead of Eqs. (26), (48) and (49), one should have (for $`d=3`$)
$$\mu _2\mu _2^{(0)}+\frac{21}{20}\mu _2^{(1)}a_2,$$
(72)
$$\mu _4\mu _4^{(0)}+\frac{13}{15}\mu _4^{(1)}a_2,$$
(73)
$$c\frac{2}{\pi ^{1/2}}\left(1\frac{3}{2}\frac{1}{8}a_2\right),$$
(74)
$$c^3\frac{2\sqrt{2}}{\pi ^{1/2}}\left(1+\frac{23}{20}\frac{3}{8}a_2\right).$$
(75)
If we insert the above expressions into Eq. (44) and neglect terms nonlinear in $`a_2`$, we get
$$a_2\frac{240(1\alpha )(1+2\alpha ^2)}{19571125\alpha +390(1\alpha )\alpha ^2}.$$
(76)
This semi-empirical estimate exhibits a fairly good agreement with the simulation data, as shown in Fig. 8.
The limitations of a Sonine description in the case of the non-Gaussian thermostat are quite apparent in Fig. 9, where $`\mathrm{\Delta }(c)`$ is plotted for $`\alpha =0.4`$, 0.6 and 0.95. The curves corresponding to $`\alpha =0.4`$ and $`\alpha =0.6`$ practically coincide, while the curve corresponding to $`\alpha =0.95`$ clearly deviates in the region of very small velocities. As a matter of fact, $`\mathrm{\Delta }(0)`$ is roughly equal to $`a_2^1`$, which indicates an almost vanishing population of rest particles, i.e. $`\stackrel{~}{f}(\mathrm{𝟎})0`$, even at $`\alpha =0.95`$.
A key feature of Fig. 9 is the existence of a non-zero initial slope, $`\mathrm{\Delta }(c)/c|_{c=0}0`$, that cannot be described by any polynomial in $`c^2`$.
From the analysis made at the end of Subsect. 2.3, we expect an underpopulated high energy tail of the form (54), where the coefficient $`A`$ is unknown. By a fitting of the simulation results we have estimated $`A1.48`$ for $`\alpha =0.3`$ and $`A1.51`$ for $`\alpha =0.4`$. Figure 10 shows the function
$$G(c)e^{Ac^2}\stackrel{~}{f}(𝐜)$$
(77)
for $`\alpha =0.3`$ and $`\alpha =0.4`$. In both cases the value of $`K`$ is $`K1.7`$.
The regions of small and large velocities are highly underpopulated with respect to the Maxwell-Boltzmann distribution. At $`c=3`$, for instance, $`\stackrel{~}{f}/\varphi 0.12`$ for $`\alpha =0.3`$ and $`\stackrel{~}{f}/\varphi 0.10`$ for $`\alpha =0.4`$.
## 5 Summary and discussion
In this paper we have performed direct Monte Carlo simulations of the Enskog-Boltzmann equation for a fluid of smooth inelastic spheres in spatially uniform states. Upon describing the velocity distribution of the granular fluid by the Enskog-Boltzmann equation (1) it has been implicitly assumed the validity of the “molecular chaos” hypothesis of uncorrelated binary collisions. However, molecular dynamics simulations of hard disks have shown a non-uniform distribution of impact parameters for high enough dissipation ($`\alpha <0.8`$) Luding . In addition, there exist long range spatial correlations in density and flow fields which cannot be understood on the basis of the Enskog-Boltzmann equation NEBO97 . These two effects are associated with the appearance of the so-called cluster instability BM90 for systems sufficiently large. Since we have simulated directly the spatially uniform equation (1), such an instability is precluded in the simulations.
To compensate for cooling effects associated with the inelasticity of collisions, three types of “thermostatting” external driving forces have been considered. We have analyzed the deviation of the steady-state velocity distribution function from the Maxwell-Boltzmann distribution, as measured by the fourth cumulant $`a_2`$ and by the high energy tail.
A simple mechanism for thermostatting the system is to assume that the particles are subjected to random kicks WM96 , what mimics the effects of shaking or vibrating the vessel ER89 . If this stochastic force has the properties of a white noise \[cf. Eq. (8)\], it gives rise to a Fokker-Planck diffusion term in the Enskog-Boltzmann equation NE98 . By making a first Sonine approximation, van Noije and Ernst NE98 have obtained an approximate expression for $`a_2`$ as a function of the coefficient of normal restitution $`\alpha `$. Our simulation results confirm the accuracy of that expression even for large dissipation ($`\alpha =0.2`$). We have also confirmed a high energy tail of the form $`f(𝐯)\mathrm{exp}[A(v/v_0)^{3/2}]`$ (where $`v_0`$ is the thermal velocity) derived in Ref. NE98 . Moreover, that asymptotic behavior (which represents an overpopulation with respect to the Maxwell-Boltzmann distribution) is already practically reached for $`v>4v_0`$, at least for $`\alpha =0.4\text{ and }0.5`$.
In the absence of any external forcing, the freely evolving granular fluid reaches a homogeneous cooling state in which all the time dependence of the velocity distribution occurs through the thermal velocity $`v_0(t)`$, so that the distribution of the reduced velocity $`𝐜=𝐯/v_0`$ is stationary. When the Enskog-Boltzmann equation is written in terms of this reduced velocity, the operator $`/t`$ gives rise to an operator that coincides with the one representing the action of an external force proportional to the particle velocity \[cf. Eq. (12)\]. This type of “anti-drag” force can also be justified by Gauss’s principle of least constraint EM90 and has been widely used in nonequilibrium molecular dynamics simulations of molecular fluids. Thus, the homogeneous cooling state is equivalent to the steady state reached under a Gaussian thermostat. In their simulations, Brey et al. BMC96 ; BCR99 used the former point of view, while in this paper we have used the latter. Our simulations complement those of Ref. BMC96 also in that we have considered a wide range $`0.2\alpha 1`$, while Brey et al. BMC96 analyzed in detail the region $`0.7\alpha 1`$. They obtained an excellent agreement with the estimate (5) based on a Sonine approximation, first derived in Ref. NE96 . However, as $`\alpha `$ decreases and $`a_2`$ becomes larger, we have seen in this paper that Eq. (5) overestimates $`a_2`$. This discrepancy can be traced back to contributions associated with higher order Sonine polynomials as well as to the ambiguity involved in the approximate determination of $`a_2`$ by neglecting nonlinear terms in the exact equation (38). If Eq. (38) is rewritten in another equivalent form (for example, by transferring a quantity from one side to the other), the same method yields a different approximation for $`a_2`$. As long as $`a_2`$ remains small (say $`|a_2|<0.05`$), all the approximations give practically undistinguishable results. On the other hand, for larger values of $`a_2`$ (i.e., for $`\alpha <0.5`$) the result is relatively dependent of the route followed. By starting from Eq. (38) rewritten as $`\mu _4/(1+a_2)=(d+2)\mu _2`$, we have obtained the estimate (40), which is seen to agree fairly well with the simulation results for the whole range of coefficients of restitution considered. The asymptotic analysis of the kinetic equation predicts a high energy tail of the form NE98 ; EP97 $`f(𝐯)\mathrm{exp}[A(v/v_0)]`$, what represents an overpopulation phenomenon stronger than in the previous case. This behavior was already confirmed in Ref. BCR99 for $`d=2`$ and has now been confirmed by our simulation results for $`d=3`$.
In the case of the Gaussian thermostat, the heating force points in the motion direction and its magnitude is proportional to that of the particle velocity. This is a very efficient thermostat because it gives more energy to fast particles, which are the ones colliding more frequently. In contrast, the stochastic thermostat adds a velocity increment per unit of time that is random both in direction and in magnitude. This is why the high energy population is larger with the Gaussian thermostat than with the stochastic thermostat. Nevertheless, in both cases such a population is larger than in the case of elastic particles at equilibrium. One could be tempted to expect that this overpopulation is a common feature of heated granular fluids, regardless of the mechanism of heating. Our third choice of thermostat, Eq. (14), proves that this is not the case. Like in the case of the stochastic thermostat, the force is independent of the magnitude of the particle velocity; like in the case of the Gaussian thermostat, the force is deterministic and points in the motion direction. The action of this third thermostat can be graphically described by saying that, between two successive collisions, a particle feels a “pseudo-gravity” field that makes it to “fall” along its motion direction. With this choice of a non-Gaussian deterministic thermostat, the Sonine polynomials $`\{S_p(c^2)\}`$ are not a good set to represent the ratio $`\stackrel{~}{f}(𝐜)/\varphi (𝐜)`$, even for low dissipation. As a consequence, the theoretical estimate of $`a_2`$ derived by assuming that $`[\stackrel{~}{f}(𝐜)/\varphi (𝐜)1]/a_2\mathrm{\Delta }(c)S_2(c^2)`$, while being qualitatively correct, is not quantitatively accurate. We have not been able to get the functional form of $`\mathrm{\Delta }(c)`$ in the limit of low dissipation. However, we have estimated its contributions to $`c`$, $`c^3`$, $`\mu _2`$ and $`\mu _4`$ from the simulation data. This has allowed us to obtain an approximate expression for $`a_2`$ that fits well the simulation results. An interesting feature of the velocity distribution function in this case is that it is highly underpopulated with respect to the Maxwell-Boltzmann distribution both for small and large velocities. Between two successive collisions, every particle experiences a constant tangential acceleration $`g`$. The total work done by this force is exactly compensated by the total loss of energy through collisions, which are much more frequent for fast particles than for slow ones. Therefore, the population of slow particles decreases because of the action of the external force, while that of fast particles decreases because of the effect of collisions. The high energy tail of the distribution function is of the form $`\stackrel{~}{f}(𝐜)\mathrm{exp}(Ac^2)`$ with $`A>1`$. In this case the gain and loss terms of the collision integral are comparable, so that the dependence of $`A`$ on $`\alpha `$ is an open problem.
###### Acknowledgements.
Partial support from the DGES (Spain) through grant No. PB97-1501 and from the Junta de Extremadura–Fondo Social Europeo through grant No. IPR98C019 is gratefully acknowledged.
## Appendix A Collisional moments
In this Appendix we derive the expressions (64)–(3). Starting from Eq. (22) and by a standard change of variables, it is easy to get NE98
$$\mu _p=𝑑𝐜_1𝑑𝐜_2\stackrel{~}{f}(𝐜_1)\stackrel{~}{f}(𝐜_2)\stackrel{~}{\mathrm{\Phi }}_p(𝐜_1,𝐜_2),$$
(78)
where
$$\stackrel{~}{\mathrm{\Phi }}_p(𝐜_1,𝐜_2)=\frac{1}{2}𝑑\widehat{𝝈}\mathrm{\Theta }(𝐜_{12}\widehat{𝝈})(𝐜_{12}\widehat{𝝈})\left[c_1^p+c_2^pc_{1}^{}{}_{}{}^{p}c_{2}^{}{}_{}{}^{p}\right],$$
(79)
with $`𝐜_{1,2}^{}=𝐜_{1,2}\frac{1}{2}(1+\alpha )(𝐜_{12}\widehat{𝝈})\widehat{𝝈}`$. In the cases $`p=2`$ and $`p=4`$ we have
$$c_1^2+c_2^2c_{1}^{}{}_{}{}^{2}c_{2}^{}{}_{}{}^{2}=\frac{1\alpha ^2}{2}(𝐜_{12}\widehat{𝝈})^2,$$
(80)
$`c_1^4+c_2^4c_{1}^{}{}_{}{}^{4}c_{2}^{}{}_{}{}^{4}`$ $`=`$ $`(1\alpha ^2)(𝐜_{12}\widehat{𝝈})^2\left({\displaystyle \frac{c_{12}^2}{4}}+C_{12}^2\right)`$ (81)
$`{\displaystyle \frac{(1\alpha ^2)^2}{8}}(𝐜_{12}\widehat{𝝈})^4`$
$`2(1+\alpha )^2(𝐜_{12}\widehat{𝝈})^2(𝐂_{12}\widehat{𝝈})^2`$
$`+4(1+\alpha )(𝐜_{12}\widehat{𝝈})(𝐂_{12}\widehat{𝝈})`$
$`\times (𝐜_{12}𝐂_{12}),`$
where $`𝐜_{12}=𝐜_1𝐜_2`$, $`𝐂_{12}=\frac{1}{2}(𝐜_1+𝐜_2)`$. Consequently,
$$\stackrel{~}{\mathrm{\Phi }}_2(𝐜_1,𝐜_2)=\frac{\beta _3(1\alpha ^2)}{4}c_{12}^3,$$
(82)
$`\stackrel{~}{\mathrm{\Phi }}_4(𝐜_1,𝐜_2)`$ $`=`$ $`\beta _3c_{12}\{{\displaystyle \frac{(d+2)(1\alpha ^2)}{2d}}c_{12}^2C_{12}^2`$ (83)
$`+{\displaystyle \frac{(1\alpha ^2)(d+1+2\alpha ^2)}{8(d+3)}}c_{12}^4`$
$`+{\displaystyle \frac{(2d+33\alpha )(1+\alpha )}{d+3}}`$
$`\times [(𝐜_{12}𝐂_{12})^2{\displaystyle \frac{1}{d}}c_{12}^2C_{12}^2]\},`$
where we have taken into account that
$$𝑑\widehat{𝝈}\mathrm{\Theta }(\widehat{𝐜}\widehat{𝝈})(\widehat{𝐜}\widehat{𝝈})^n=\pi ^{(d1)/2}\frac{\mathrm{\Gamma }((n+1)/2)}{\mathrm{\Gamma }((n+d)/2)}\beta _n,$$
(84)
$$𝑑\widehat{𝝈}\mathrm{\Theta }(\widehat{𝐜}\widehat{𝝈})(\widehat{𝐜}\widehat{𝝈})^n\widehat{𝝈}=\beta _{n+1}\widehat{𝒄},$$
(85)
$$𝑑\widehat{𝝈}\mathrm{\Theta }(\widehat{𝐜}\widehat{𝝈})(\widehat{𝐜}\widehat{𝝈})^n\widehat{𝝈}\widehat{𝝈}=\frac{\beta _n}{n+d}\left(n\widehat{𝒄}\widehat{𝒄}+𝖨\right).$$
(86)
In the three-dimensional case, Eqs. (78), (82) and (83) yield Eqs. (64)–(3). |
warning/0002/hep-th0002134.html | ar5iv | text | # 1 Introduction
## 1 Introduction
The Born-Infeld (BI) action in flat spacetime, <sup>2</sup><sup>2</sup>2We use $`\eta _{\mu \nu }=\mathrm{diag}(+,,,)`$ and $`\mathrm{}=c=1`$.
$$S_{\mathrm{BI}}=\frac{1}{b^2}d^4x\left\{1\sqrt{det(\eta _{\mu \nu }+bF_{\mu \nu })}\right\},$$
$`(1.1)`$
is the particular non-linear generalization of Maxwell theory, $`F_{\mu \nu }=_\mu A_\nu _\nu A_\mu `$. The action (1.1) was initially introduced to regularize both the electric field and the self-energy of a point-like charge in electrodynamics . Much later, the BI action was recognized as the leading contribution to the effective action of open strings in an abelian background with constant field strength $`F`$ , and as the essential part of the D3-brane action as well , with $`b=2\pi \alpha ^{}`$. The action (1.1) has many remarkable properties, e.g., causal propagation and electric-magnetic duality .
The BI Lagrangian can be rewritten to the form
$$L=\frac{1}{2}p^{\mu \nu }F_{\mu \nu }+H(P,Q),$$
$`(1.2)`$
where the auxiliary antisymmetric tensor $`p_{\mu \nu }`$ and the BI structure function
$$H(P,Q)=\frac{1}{b^2}\left(1\sqrt{12b^2P+b^4Q^2}\right),$$
$`(1.3)`$
as well as the definitions
$$P=\frac{1}{4}p_{\mu \nu }p^{\mu \nu },Q=\frac{i}{4}p_{\mu \nu }\stackrel{~}{p}^{\mu \nu },\stackrel{~}{p}^{\mu \nu }=\frac{1}{2}\epsilon ^{\mu \nu \rho \sigma }p_{\rho \sigma },$$
$`(1.4)`$
have been introduced. Eliminating $`p_{\mu \nu }`$ from eq. (1.2) results in the equivalent Lagrangian
$$L=\frac{1}{b^2}\left[1\sqrt{1+\frac{b^2}{2}F^2\frac{b^4}{16}(F\stackrel{~}{F})^2}\right],$$
$`(1.5)`$
where we have defined $`F^2=F^{\mu \nu }F_{\mu \nu }`$, $`\stackrel{~}{F}^{\mu \nu }=\frac{1}{2}\epsilon ^{\mu \nu \rho \sigma }F_{\rho \sigma }`$ and $`F\stackrel{~}{F}=F^{\mu \nu }\stackrel{~}{F}_{\mu \nu }`$.
Supersymmetric generalizations of the BI action are of particular interest in connection to superstring theory (see ref. for a recent review). The super-BI actions describing D-branes can be naturally interpreted as the Goldstone-type actions associated with partial supersymmetry breaking, while they can still be duality invariant too. The manifestly N=1 supersymmetric generalization of the four-dimensional BI action in N=1 superspace was discovered long time ago (see also ref. ), while its manifestly N=2 supersymmetric generalization in N=2 superspace was found only recently (see ref. too). To our knowledge, the higher (N=3 or N=4) manifestly supersymmetric generalizations of the four-dimensional bosonic BI action (1.1) are not known in any form.
Supersymmetry apparently prefers the parametrization of the BI action in terms of the Maxwell term $`L_2=\frac{1}{4}F^2`$ and the Maxwell stress-energy tensor squared ,
$$L_4=\frac{1}{32}\left\{(F^2)^2+(F\stackrel{~}{F})^2\right\}=\frac{1}{8}(F^+)^2(F^{})^2,F_{\mu \nu }^\pm =\frac{1}{2}\left(F_{\mu \nu }\pm i\stackrel{~}{F}_{\mu \nu }\right).$$
$`(1.6)`$
This term is known as the Euler-Heisenberg (EH) Lagrangian . The EH action also appears as the bosonic part of the one-loop effective action in N=1 supersymmetric scalar electrodynamics with the parameter $`b^1=2\sqrt{6}\pi m^2/e^2`$. One easily finds that
$$L_{\mathrm{BI}}=\frac{1}{b^2}\left\{1\sqrt{(1b^2L_2)^22b^4L_4}\right\}=L_2+b^2L_4+O(F^6).$$
$`(1.7)`$
A manifestly N=4 supersymmetric generalization of the BI action is known to be the formidable problem, though it is highly desirable, e.g., for an investigation of quantum properties of D3-branes and their comparison to supergravity . Even the $`N>2`$ supersymmetrization of the EH-term $`L_4`$, representing the four-derivative terms $`(F^4)`$, is non-trivial. The additional terms with four derivatives in the N=4 BI action were determined in ref. in N=1 superspace, by imposing the $`SU(3)`$ internal symmetry on three N=1 chiral multiplets extending an N=1 (abelian) vector multiplet to an N=4 vector multiplet, with manifest (linearly realised) N=1 off-shell supersymmetry. The manifestly N=2 supersymmetric form of the N=4 EH action was derived in ref. in N=2 projective superspace, while its equations of motion can also be written in terms of on-shell N=4 superfields in harmonic superspace . It is the purpose of this Letter to write down an off-shell, manifestly N=3 supersymmetric formulation of the N=4 EH action in N=3 light-cone superspace.
Our paper is organized as follows. In sect. 2 we introduce a light-cone gauge and rewrite the EH Lagrangian in terms of physical (transverse) degrees of freedom up to the relevant order. In sect. 3 we introduce N=3 light-cone superspace and deduce an N=3 supersymmetric generalization of the EH action in terms of a single N=3 light-cone superfield. The obstructions encountered in our efforts to find a similar, manifestly N=4 supersymmetric EH action in N=4 light-cone superspace are discussed in Conclusion (sect. 4).
## 2 EH action in light-cone gauge
The light-cone formulation of a gauge theory (in light-cone gauge) keeps only physical (transverse) degrees of freedom in the field theory by giving up its manifest Lorentz invariance. The light-cone formulation is, therefore, very suitable for an off-shell formulation of N-extended supersymmetric gauge field theories with manifest supersymmetry in N-extended light-cone superspace .
We define light-cone coordinates in Minkowski spacetime as
| $`x^+={\displaystyle \frac{1}{\sqrt{2}}}\left(x^0+x^3\right),`$ | $`x^{}={\displaystyle \frac{1}{\sqrt{2}}}\left(x^0x^3\right),`$ |
| --- | --- |
| $`x={\displaystyle \frac{1}{\sqrt{2}}}\left(x^1+ix^2\right),`$ | $`\overline{x}={\displaystyle \frac{1}{\sqrt{2}}}\left(x^1ix^2\right),`$ |
$`(2.1)`$
and similarly for the gauge vector field, $`A_\mu (A^+,A^{},A,\overline{A})`$. The real coordinate $`x^+`$ is going to be considered as ‘light-cone time’. The linear transformation (2.1) of spacetime coordinates is obviously non-singular (with the Jacobian equal to $`i`$), while it does not preserve the Minkowski metric (i.e. it is not a Lorentz-transformation).
The light-cone gauge reads
$$A^+=0.$$
$`(2.2)`$
In this (physical) gauge the $`A^{}`$ component of the gauge field $`A_\mu `$ is supposed to be eliminated via its (non-dynamical) equation of motion, whereas the transverse components $`(A,\overline{A})`$ are supposed to represent the physical propagating fields.
It is easy to solve the equation of motion for $`A^{}`$ in the Maxwell theory, where it takes the form of a linear equation in the light-cone gauge (cf. refs. ). It becomes, however, a highly non-trivial problem in the BI or EH theory, where it takes the form of a non-linear partial differential equation. The equations of motion amount to the conservation law for the $`p`$-tensor,
$$^\mu p_{\mu \nu }=0,$$
$`(2.3)`$
while the $`p_{\mu \nu }`$ in the BI theory is given by
$$p_{\mu \nu }=\frac{b^2F_{\mu \nu }\frac{b^4}{4}(F\stackrel{~}{F})\stackrel{~}{F}_{\mu \nu }}{\sqrt{1+\frac{b^2}{2}F^2\frac{b^4}{16}(F\stackrel{~}{F})^2}}.$$
$`(2.4)`$
By the use of the Bianchi identity, $`^\mu \stackrel{~}{F}_{\mu \nu }=0`$, we find the following equation for $`A^{}`$:
| $`^\mu F_\mu =`$ | $`b^2\left\{\frac{1}{2}^\mu F_\mu F^2+\frac{1}{4}\stackrel{~}{F}_\mu ^\mu (F\stackrel{~}{F})+\frac{1}{4}F_\mu ^\mu F^2\right\}`$ |
| --- | --- |
| | $`+b^4\{\frac{1}{16}(F\stackrel{~}{F})\stackrel{~}{F}_\mu ^\mu F^2+\frac{1}{16}^\mu F_\mu (F\stackrel{~}{F})^2`$ |
| | $`+\frac{1}{16}\stackrel{~}{F}_\mu ^\mu (F\stackrel{~}{F})F^2\frac{1}{32}F_\mu ^\mu (F\stackrel{~}{F})^2\}.`$ |
$`(2.5)`$
We use a perturbative Ansatz, in powers of the small parameter $`b^2`$, for a solution to eq. (2.5),
$$A^{}(x)=\underset{n=0}{\overset{\mathrm{}}{}}b^{2n}A_{(2n)}^{}(x).$$
$`(2.6)`$
As regards the leading and sub-leading terms, we find
| $`A_{(0)}^{}=`$ | $`{\displaystyle \frac{1}{^+}}\left(\overline{}A+\overline{A}\right),`$ |
| --- | --- |
| $`A_{(2)}^{}=`$ | $`{\displaystyle \frac{1}{(^+)^2}}\left[\frac{1}{2}^\mu F_\mu F^2+\frac{1}{4}\stackrel{~}{F}_\mu ^\mu (F\stackrel{~}{F})+\frac{1}{4}F_\mu ^\mu F^2\right]|_{A^{}=A_{(0)}^{}},`$ |
$`(2.7)`$
where we have used the notation $`^+=/x^{}`$. The multiple factors $`(^+)^1`$ in our actions are harmless after rewriting them to momentum space. The first line of eq. (2.7) coincides with the exact solution in the Maxwell theory.
According to eq. (1.7), the EH term $`L_4`$ is the leading $`b^2`$-correction to the Maxwell term $`L_2`$ in the BI theory. The light-cone formulation of the BI Lagrangian in the same approximation is thus given by the terms written down on the right-hand-side of eq. (1.7) after a substitution of eq. (2.2) and the first line of eq. (2.8). After some algebra and partial integration we find
| $`L[A,\overline{A}]`$ | $`={\displaystyle \frac{1}{4}}F^2+{\displaystyle \frac{b^2}{8}}(F^+)^2(F^{})^2`$ |
| --- | --- |
| | $`=A\mathrm{}\overline{A}+2b^2\left|(\overline{A})^2+^+\overline{A}{\displaystyle \frac{\mathrm{}}{2^+}}A^+\overline{A}{\displaystyle \frac{^2}{^+}}\overline{A}\right|^2+O(b^4),`$ |
$`(2.8)`$
where we have used the notation $`=/x`$ and $`\overline{}=/\overline{x}`$. Eq. (2.8) can be thought of as the light-cone EH Lagrangian. Its N=3 supersymmetrization is discussed in the next sect. 3.
## 3 N=3 light-cone superspace action
The light-cone N=3 supersymmetry algebra reads
$$\{Q^m,\overline{Q}_n\}=\sqrt{2}\delta _n^mP^+,m,n=1,2,3,$$
$`(3.1)`$
where the supersymmetry charges $`Q^r`$ transform in the fundamental representation of $`SU(3)`$. A natural representation of the algebra (3.1) in N=3 light-cone superspace $`Z=(x^\mu ,\theta ^m,\overline{\theta }_n)`$ is given by
| $`Q^m=`$ | $`{\displaystyle \frac{}{\overline{\theta }_m}}+{\displaystyle \frac{i}{\sqrt{2}}}\theta ^m^+,`$ |
| --- | --- |
| $`\overline{Q}_n=`$ | $`{\displaystyle \frac{}{\theta ^n}}{\displaystyle \frac{i}{\sqrt{2}}}\overline{\theta }_n^+.`$ |
$`(3.2)`$
The covariant derivatives in N=3 light-cone superspace are
| $`D^m=`$ | $`{\displaystyle \frac{}{\overline{\theta }_m}}{\displaystyle \frac{i}{\sqrt{2}}}\theta ^m^+,`$ |
| --- | --- |
| $`\overline{D}_n=`$ | $`{\displaystyle \frac{}{\theta ^n}}+{\displaystyle \frac{i}{\sqrt{2}}}\overline{\theta }_n^+.`$ |
$`(3.3)`$
They anticommute with the supersymmetry charges (3.2) and obey the same algebra (3.1). The irreducible off-shell representations of N=3 light-cone supersymmetry are easily obtained by imposing the covariant chirality condition on N=3 light-cone superfields $`\varphi (Z)`$,
$$D^m\varphi (Z)=0.$$
$`(3.4)`$
A solution to eq. (3.4) in components is just given by an arbitrary complex function $`\varphi (x^+,x^{}+\frac{i}{\sqrt{2}}\theta ^m\overline{\theta }_m,x,\overline{x};\theta ^n)\varphi (y;\theta )`$. Its expansion in the chiral superspace reads
$$\varphi (y;\theta )=\frac{1}{^+}A(y)+\frac{i}{^+}\theta ^m\overline{\chi }_m(y)+\frac{i}{2}\theta ^m\theta ^n\epsilon _{mnp}C^p(y)+\frac{1}{3!}\epsilon _{mnp}\theta ^m\theta ^n\theta ^p\psi (y).$$
$`(3.5)`$
The light-cone N=3 supersymmetry transformation laws for the components are
| $`\delta A=`$ | $`i\epsilon ^n\overline{\chi }_n,`$ |
| --- | --- |
| $`\delta \overline{\chi }_m=`$ | $`\sqrt{2}\overline{\epsilon }_m^+A+\epsilon _{mnp}\epsilon ^n^+C^p,`$ |
| $`\delta C^p=`$ | $`i\sqrt{2}\epsilon ^{pqr}\overline{\epsilon }_q\overline{\chi }_ri\epsilon ^p\psi ,`$ |
| $`\delta \psi =`$ | $`\sqrt{2}\overline{\epsilon }_n^+C^n,`$ |
$`(3.6)`$
where $`(\epsilon ^n,\overline{\epsilon }_m)`$ are the infinitesimal anticommuting parameters.
All our field components have canonical dimensions. The complex field $`A`$ can be identified with the physical (translational) vector field components, the spinors $`\overline{\chi }_m`$ in the fundamental representation $`\mathrm{𝟑}`$ of $`SU(3)`$ with a triplet of photinos, the singlet spinor $`\psi `$ with extra photino, and the complex triplet $`C^m`$ with Higgs fields in $`\mathrm{𝟑}`$ of $`SU(3)`$. The physical content thus coincides with that of the N=4 supersymmetric abelian vector multiplet having a single photon field, photinos in the fundamental representation $`\mathrm{𝟒}`$ of $`SU(4)`$ and Higgs fields in real $`\mathrm{𝟔}`$ of $`SU(4)`$, after their decomposition with respect to the $`SU(3)`$ subgroup of the internal symmetry $`SU(4)`$. This is the manifestation of the well-known fact that N=3 and N=4 supersymmetric vector multiplets are physically equivalent.
It is now straightforward (though very tedious) to find the N=3 supersymmetric generalization of the bosonic EH light-cone action (2.8) in N=3 light-cone superspace,
$$S=d^4xd^3\theta d^3\overline{\theta }(\varphi ,\overline{\varphi })=d^4x(D)^3(\overline{D})^3(\varphi ,\overline{\varphi }),$$
$`(3.7)`$
where $`(D)^3=\epsilon _{mnp}D^mD^nD^p`$ and similarly for $`(\overline{D})^3`$. After some trials and errors, we find
| $`36(i\sqrt{2})^3(\varphi ,\overline{\varphi })=`$ | $`\varphi {\displaystyle \frac{\mathrm{}}{^+}}\overline{\varphi }+2b^2\{{\displaystyle \frac{1}{^{+3}}}\left(\overline{}^+\varphi \overline{}^+\varphi \right)(^+\overline{\varphi })^2`$ |
| --- | --- |
| | $`+{\displaystyle \frac{1}{^{+3}}}\left(^{+2}\varphi \overline{}^2\varphi \right)^{+2}\overline{\varphi }^2\overline{\varphi }+{\displaystyle \frac{1}{2^+}}(\varphi )(^+\overline{\varphi })^2\mathrm{}\overline{\varphi }`$ |
| | $`+{\displaystyle \frac{1}{4^{+3}}}\left(^{+2}\varphi \mathrm{}\varphi \right)^{+2}\overline{\varphi }\mathrm{}\overline{\varphi }{\displaystyle \frac{1}{^{+3}}}\left(^{+2}\varphi \overline{}^2\varphi \right)(^+\overline{\varphi })^2`$ |
| | $`{\displaystyle \frac{1}{2^{+3}}}\left(^{+2}\varphi \mathrm{}\varphi \overline{}^2\varphi \right)^{+2}\overline{\varphi }{\displaystyle \frac{1}{^{+3}}}\left(\overline{}^+\varphi \overline{}^+\varphi \right)^{+2}\overline{\varphi }^2\overline{\varphi }`$ |
| | $`+{\displaystyle \frac{1}{2^{+3}}}\left(\overline{}^+\varphi \overline{}^+\varphi \mathrm{}\varphi \right)^{+2}\overline{\varphi }{\displaystyle \frac{1}{2^+}}(\varphi )^{+2}\overline{\varphi }\mathrm{}\overline{\varphi }^2\overline{\varphi }\}.`$ |
$`(3.8)`$
The bosonic part of this action is given by
| $`_{\mathrm{bos}.}=`$ | $`A\mathrm{}\overline{A}+2b^2\left|(\overline{A})^2+^+\overline{A}{\displaystyle \frac{\mathrm{}}{2^+}}A^+\overline{A}{\displaystyle \frac{^2}{^+}}\overline{A}\right|^2`$ |
| --- | --- |
| | $`{\displaystyle \frac{1}{2}}C^p\mathrm{}\overline{C}_p2b^2\{{\displaystyle \frac{2}{^{+2}}}\left(\overline{}^+C^m\overline{}^+A\right)(^+\overline{C}_m^+\overline{A})`$ |
| | $`+{\displaystyle \frac{1}{2^{+2}}}\left(^{+2}C^p\overline{}^2A+\overline{}^2C^p^{+2}A\right)\left(^{+2}\overline{C}_p^2\overline{A}+^2\overline{C}_p^{+2}\overline{A}\right)`$ |
| | $`+{\displaystyle \frac{1}{8^{+2}}}\left(^{+2}C^p\mathrm{}A+\mathrm{}C^p^{+2}A\right)\left(^{+2}\overline{C}_p\mathrm{}\overline{A}+\mathrm{}\overline{C}_p^{+2}\overline{A}\right)`$ |
| | $`+[{\displaystyle \frac{1}{4}}C^p(2^+\overline{C}_p^+\overline{A}\mathrm{}\overline{A}+\mathrm{}\overline{C}_p^+\overline{A}^+\overline{A})(3.9)`$ |
| | $`{\displaystyle \frac{1}{^{+2}}}\left(^{+2}C^p\overline{}^2A+\overline{}^2C^p^{+2}A\right)^+\overline{C}_p^+\overline{A}`$ |
| | $`{\displaystyle \frac{1}{4^{+2}}}(^{+2}C^p\mathrm{}A\overline{}^2A+\overline{}^2C^p^{+2}A\mathrm{}A+\mathrm{}C^p\overline{}^2A^{+2}A)^{+2}\overline{C}_p+\mathrm{h}.\mathrm{c}.]\}.`$ |
One of the obvious features of both eqs. (3.8) and (3.9) is the apparent presence of higher derivatives, as may have been expected from the experience with the manifestly N=2 supersymmetric generalization of the BI action in the covariant N=2 superspace . The expected correspondence to the component D3-brane effective action having non-linearly realized extended supersymmetry and no higher derivatives implies the existence of a field redefinition that would eliminate the higher-derivative terms in our action and make its N=3 supersymmetry to be non-linearly realised (i.e non-manifest) . We also note the absence of quartic $`(C^4)`$ scalar terms and the on-shell ($`\mathrm{}A=\mathrm{}C=0`$) invariance of our action under constant shifts, $`C_p(x)C_p(x)+c_p`$, which are supposed to be related to the possible interpretation of the $`C_p`$ fields as the Goldstone scalars associated with spontaneoulsy broken translations in the full N=3 BI action.
## 4 Conclusion
Our main results are given by eqs. (2.8), (3.8) and (3.9). Our initial motivation was to construct an N=4 supersymmetric generalization of the EH action in the light-cone gauge. The N=4 light-cone supersymmetry algebra is given by eq. (3.1), where the indices now take four values. Equations (3.2), (3.3) and (3.4) are still valid in N=4 light-cone superspace, where they have to supplemented by an extra (generalized reality) condition ,
$$D^mD^n\overline{\varphi }=\frac{1}{2}\epsilon ^{mnpq}\overline{D}^p\overline{D}_q\varphi ,\mathrm{or},\mathrm{equivalently},\overline{\varphi }=\frac{1}{48^{+2}}\epsilon ^{mnpq}\overline{D}_m\overline{D}_n\overline{D}_p\overline{D}_q\varphi .$$
$`(4.1)`$
The restricted chiral N=4 light-cone superfield $`\varphi `$ is equivalent to the chiral N=3 superfield in eq. (3.5). Our efforts to construct an N=4 generalization of eq. (2.8) along the similar lines (sect. 3) unexpectedly failed, while eq. (4.1) was the main obstruction. We conclude that even a manifestly N=4 supersymmetric generalization of the EH action in the light-cone gauge seems to be highly non-trivial, if any, not to mention an even more ambitious (manifest) N=4 supersymmetrization of the BI action.
## Acknowledgements
We are grateful to Norbert Dragon, Gordon Chalmers, Olaf Lechtenfeld and Daniela Zanon for useful discussions. |
warning/0002/quant-ph0002087.html | ar5iv | text | # A TENTATIVE EXPRESSION OF THE KÁROLYHÁZY UNCERTAINTY OF THE SPACE-TIME STRUCTURE THROUGH VACUUM SPREADS IN QUANTUM GRAVITY
## 1 Introduction
In his pioneering work on the stochastic modification of the Schrödinger time evolution, known as the K model (for a recent review, see ), Károlyházy relates the loss of coherence and the breakdown of the superposition principle to the uncertainty of the Einsteinian space-time structure, caused by the quantum mechanical uncertainties of the position and of the momentum of material objects. The K model leads to sound results concerning the occurrence (or the absence) of the breakdown of the superposition principle for various physical systems. It should be noted that in obtaining those results there is no room for maneuvers, because there are no free parameters in the K model. However, the theoretical reliability of the results is weakened by the fact that in the calculations the quantum mechanical uncertainty of the space-time structure has been replaced by a classical spread. In (see also ) the single space-time $`S`$, with a quantum mechanical uncertainty in its structure, has been mimicked by a stochastic set $`\{S_\beta \}`$ of classical space-times $`S_\beta `$, with appropriately chosen metric tensors $`(g_{\mu \nu }(𝐱,t))_\beta `$. Instead of propagating on the single but “hazy” space-time $`S`$, the quantum mechanical wave function had to propagate on the concise, but different from each other space-times $`S_\beta `$, and the spread of the relative phases of the wave function over the set $`\{S_\beta \}`$ was supposed to be equal to the amount of the quantum mechanical uncertainty of the relative phase, induced by the uncertain structure of the single space-time $`S`$. In an other stochastic space-time model, leading to the same result as the original model, has been used. In this second model $`g_{\mu \nu }`$ remains concisely Minkowskian, but the moments of time are randomized via the introduction of an appropriate random set $`\{t_\beta (𝐱,t)\}`$ of moments of time $`t_\beta (𝐱,t)`$.
Károlyházy was well aware of the shortcomings of such space-time models. In he wrote: “To avoid possible misunderstanding, we would like to stress that … no physical significance should be attached to the individual members of the family $`\{(g_{\mu \nu })_\beta \}`$. The only role of the family is to provide us with a mathematical model of a single physical space-time with smeared metric…”
Until recently, I thought that because of the lack of a unified theory of general relativity and quantum theory, the use of a classically stochastic space-time model is inevitable. It turns out that this is not so. As shown in the present paper, with the help of two entities of the future unified theory of general relativity and quantum theory, one can obtain the known results of the K model without relying on such a space-time model. In addition, the calculations are simpler than those carried out with the help of the $`(g_{\mu \nu })_\beta `$’s in (outlined also in ). There, many expressions contain $`𝐤`$-space Fourier sums or integrals. In the formulas for the physical quantities these sums or integrals are convergent, but in some intermediary expressions they diverge. In the present work the calculations are carried out in $`𝐱`$-space, $`𝐤`$-space integrals do not appear at all, and it is easy to see that the claims that the K model needs a cutoff parameter to make divergent $`𝐤`$-space integrals finite are not justified.
In Sections 2 and 3 two relations exhibiting the uncertainty of the space-time structure are given and discussed. In Section 4 these relations are expressed with the help of entities of the future unified theory of general relativity and quantum theory. In Sections 5 and 6 a new derivation of the formulas for two basic quantities of the K model — the quantum mechanical spread $`\mathrm{\Delta }_\mathrm{\Phi }`$ of the relative phases and the cell length $`a_c`$ — is presented. The characterisation of quantum and of classical behavior, and the transition between them are also discussed in Section 6. Section 7 is devoted to the presentation of the stochastically modified Schrödinger evolution. In Section 8 a few concluding remarks are made.
## 2 Uncertainty in the Structure of the Einsteinian Space-Time and in the Relative Phases of the Wave Functions
### 2.1 The Lower Bound of the Uncertainty of the Length of Time Intervals
Investigating the uncertainties of the Einsteinian space-time structure induced by the quantum mechanical uncertainties in the position and in the momentum of various quantum objects, Károlyházy discovered that the uncertainty of the length $`T`$ of a time interval has a lower bound $`\mathrm{\Delta }_LT`$ . The relation between $`T`$ and $`\mathrm{\Delta }_LT`$ is simple:
$$\mathrm{\Delta }_LT\text{ }T_P^{2/3}T^{1/3},$$
(1)
where
$$T_P=\frac{\mathrm{\Lambda }}{c}=\sqrt{\frac{G\mathrm{}}{c^5}}=5.3\times 10^{44}\text{sec}$$
(2)
is the Planck time,
$$\mathrm{\Lambda }=1.6\times 10^{33}$$
(3)
being the Planck length and $`G`$ the constant of gravitation.
The approximate equality sign $``$ in relation (1) and in further equations takes into account that because of the absence of a unified theory of general relativity and quantum theory, numerical factors, unlike the $`1/2`$ in Heisenberg’s relation $`\mathrm{\Delta }x\mathrm{\Delta }p\mathrm{}/2`$, could not be fixed. In the context of the present paper such factors are unimportant. Indeed, as we shall see, the value of the relevant parameter — of the cell length — characterizing the coherence changes by tens of orders of magnitude while going from quantum behavior to classical behavior. Compared to this change, a shift by a factor between $`10^1`$ and $`10`$, or even between $`10^2`$ and $`10^2`$, is irrelevant. However, in a prospective experimental search for the anomalous Brownian motion predicted by the K model , the said loose factors may cause problems.
Taking advantage of the regrettable looseness in the basic relation (1), we shall often lump into the symbol $``$ known but neglected factors, e.g. when rounding numerical values or dropping $`\pi `$’s.
Two restrictions should be made concerning the applicability of relation (1).
(i) In agreement with the fact that the K model is a model for non-relativistic quantum mechanics, in relation (1) $`T`$ refers to a time interval along a world line of a body slowly moving (or standing) in a reference frame in which the 2.7 K background radiation is isotropic.
(ii) When
$$TT_P,$$
(4)
the very concept of space-time becomes questionable, and relation (1) may lose its physical meaning. Therefore, this relation should be applied only to time intervals for which
$$TT_P.$$
(5)
This restriction means that besides being a fundamental parameter of the K model, $`T_P`$ is also a physical cutoff parameter. It is the only inherent cutoff parameter in the model. It is needed in order to keep out from those very small space-time domains inside which the physical laws are not known, and not in order to make divergent mathematical expressions finite. Of course, similarly to regular non-relativistic quantum mechanics, the predictions of the K model become unreliable (but not divergent) when the realm of high energy particle physics is reached, that is already for time intervals of the order of $`10^{24}`$ sec and for spatial distances of the order of $`10^{13}`$ cm, much larger than $`T_P`$ and $`\mathrm{\Lambda }`$, respectively.
Notice that for $`T=1`$ sec, $`\mathrm{\Delta }_LT`$ is only of the order of $`10^{29}`$ sec. On the other hand, $`\mathrm{\Delta }_LT`$ is an absolute, inescapable lower bound. For a given $`T`$, it cannot be diminished at the expense of some other quantity, like $`\mathrm{\Delta }x`$ at the expense of $`\mathrm{\Delta }p`$ in Heisenberg’s uncertainty relation.
The extreme smallness of $`\mathrm{\Delta }_LT`$ is due to the fact that when deriving relation (1), only the basic laws of general relativity and of quantum mechanics should be respected, all other theoretical, as well as all practical limitations should be ignored. As a result, the minimal time uncertainty of actual physical processes of duration $`T`$ is much larger than the lower bound $`\mathrm{\Delta }_LT`$ for that value of $`T`$. $`\mathrm{\Delta }_LT`$ could be reached only in processes with bodies of irrealistically high density, although this density is still much-much smaller than the Planck density $`\varrho _P=m_P/\mathrm{\Lambda }^310^{94}\text{ g/cm}\text{3}`$. Thus, the existence of such irrealistic bodies would not contradict the basic laws of general relativity and of quantum mechanics.
Only the existence and the order of magnitude of the lower bound $`\mathrm{\Delta }_LT`$ is exploited in the K model, the possibility of reaching it is not necessary.
### 2.2 The Uncertainty of the Space-Time Structure and the Breakdown of the Superposition Principle
Attention should be payed to the rather remarkable fact that all the physical quantities referring to a particular body (like its mass, its velocity, etc.) dropped out from relation (1), the relation between $`T`$ and $`\mathrm{\Delta }_LT`$ involves only the universal constant of nature $`T_P=\sqrt{G\mathrm{}/c^5}`$. Now, a space-time relation independent of any particular property of matter can be, perhaps even must be attributed to space-time itself, therefore also to the empty space-time. Accordingly, Károlyházy proposed to regard relation (1) as an expression of the uncertainty of the structure of space-time. $`\mathrm{\Delta }_LT`$ gives then the measure of the limitation of the sharpness of the Einsteinian space-time structure, imposed by quantum mechanics. As shown in , this tiny uncertainty of the space-time structure induces uncertainties in the relative phases of the wave function of any isolated system, and thereby limits, in return, the sharpness of the phase relations. The amount of the uncertainty of the relative phases turns out to be negligible in the case of microsystems, but it is large enough to destroy the coherence of the wave function of a macrosystem, in agreement with the observed breakdown of the superposition principle.
## 3 The Structural Uncertainty of Synchronization
The structural uncertainty $`\mathrm{\Delta }_LT`$ of the length of the time intervals along the $`|𝐯|c`$ worldlines produces a structural uncertainty of the same order of magnitude in the synchronization of the times between two such worldlines . In order to see this, let us consider, first in Minkowskian space-time, two $`|𝐯|=0`$ worldlines $`W_1`$ and $`W_2`$ at a distance $`r`$ from each other. Let the times along these worldlines be synchronized. A light signal emitted on $`W_1`$ arrives back to $`W_1`$ from $`W_2`$ after a time $`2T`$, where
$$T=\frac{r}{c}.$$
(6)
In the space-time of Károlyházy, the time interval of length $`2T`$ along $`W_1`$ has the structural uncertainty $`\mathrm{\Delta }_LT`$ given by relation (1). (The uncertainties of $`2T`$ and of $`T`$ are of the same order of magnitude.) Consequently, the moment of arrival of a signal to $`W_2`$ suffers from the same uncertainty relative to the time along $`W_1`$. This means that the uncertainty of the synchronization of the times along two $`|𝐯|=0`$ worldlines at a distance $`r=cT`$ from each other has a structural lower bound
$$(\mathrm{\Delta }_LT)_{\mathrm{syn}}\text{ }T_P^{3/2}\left(\frac{r}{c}\right)^{1/3}=\frac{\mathrm{\Lambda }^{2/3}r^{1/3}}{c}.$$
(7)
This lower bound is by many orders of magnitude smaller than the uncertainty in the synchronization carried out by realistic quantum clocks. But again, we shall rely only on the existence of relation (7), the impossibility of actually reaching the lower bound $`(\mathrm{\Delta }_LT)_{\mathrm{syn}}`$ is not important.
## 4 The Expression of the Uncertainties $`\mathrm{\Delta }_LT`$ and $`(\mathrm{\Delta }_LT)_{\mathrm{syn}}`$ in Terms of Entities of the Future Unified Theory of General Relativity and Quantum Theory
Since in relations (1) and (7) $`\mathrm{\Delta }_LT`$ and $`(\mathrm{\Delta }_LT)_{\mathrm{syn}}`$ are uncertainties of the empty space-time, in the future unified theory they will presumably take the form of vacuum spreads of appropriate operators. The vacuum state $`|V`$ should “know” about general relativity, i.e. about gravitons, and if the K model is sound, then in the non-relativistic approximation $`|V`$ should represent the empty space-time of Károlyházy, instead of the empty Minkowskian space-time. As far as the appropriate operators are concerned, they should refer to time because relations (1) and (7) are time uncertainty relations, and since these relations do not contain quantities describing particular objects, the operators should not refer to particular objects either. An effective local time operator $`\widehat{t}(𝐱,t)`$ meets the above requirements. We call this operator effective, because it may well be that similarly to non-relativistic quantum mechanics, there will be no time operator at the basic level of the unified theory. The reason for considering the operator $`\widehat{t}(𝐱,t)`$ is that, as we shall see presently, the uncertainties $`\mathrm{\Delta }_LT`$ and $`(\mathrm{\Delta }_LT)_{\mathrm{syn}}`$ can be expressed in a simple way with the help of $`\widehat{t}(𝐱,t)`$ and $`|V`$. Also, $`\widehat{t}(𝐱,t)`$ exhibits two important features of the K model. It states that time is not a global, but a local quantity, in agreement with the involvement of general relativity in the K model, and says that the values of the moments of time have an uncertainty, corresponding to the uncertainty of the space-time structure.
We assume that the vacuum expectation value (the “vev”) of $`\widehat{t}(𝐱,t)`$ is independent of $`𝐱`$ and equals $`t`$,
$$\widehat{t}(𝐱,t)_V=t,$$
(8)
and we write $`\widehat{t}`$ in the convenient form
$$\widehat{t}(𝐱,t)=t+\widehat{\tau }(𝐱,t),$$
(9)
where, due to (8),
$$\widehat{\tau }(𝐱,t)_V=0.$$
(10)
Let us now look at relation (1). In this relation $`T`$ stands for the length of a time interval belonging to a segment $`[(𝐱,t),(𝐱,t^{})]`$ of a $`|𝐯|=0`$ worldline in the empty space-time of Károlyházy. It is therefore reasonable to identify $`T`$ with the vev of the operator difference
$$\widehat{t}(𝐱,t^{})\widehat{t}(𝐱,t).$$
(11)
From (8) and (9) we see that $`T`$ is equal to the Minkowskian length $`t^{}t`$ of our time interval (we take $`t^{}>t`$):
$$T:=\widehat{t}(𝐱,t^{})\widehat{t}(𝐱,t)_V=t^{}t.$$
(12)
Since $`\mathrm{\Delta }_LT`$ is the uncertainty of $`T`$ in the empty Károlyházy space-time, it should be given by the vacuum spread of the operator difference in (11):
$$(\mathrm{\Delta }_LT\text{ })^2:=(\widehat{t}(𝐱,t^{})\widehat{t}(𝐱,t)T)^2_V.$$
(13)
With equations (9) and (12) one finds that
$$(\mathrm{\Delta }_LT\text{ })^2=(\widehat{\tau }(𝐱,t^{})\widehat{\tau }(𝐱,t))^2_V.$$
(14)
Relation (1) can be written now in the desired form
$$(\widehat{\tau }(𝐱,t^{})\widehat{\tau }(𝐱,t))^2_VT_P^{4/3}T^{2/3}.$$
(15)
We turn now to relation (7). There $`(\mathrm{\Delta }_LT)_{\mathrm{syn}}`$ refers to the relative time uncertainty at two world points on a $`t=`$ constant hyperplane. $`(\mathrm{\Delta }_LT)_{\mathrm{syn}}`$ should therefore be the vacuum spread of the operator difference
$$\widehat{t}(𝐱^{},t)\widehat{t}(𝐱,t)=\widehat{\tau }(𝐱^{},t)\widehat{\tau }(𝐱,t).$$
(16)
With Equation (8) one finds that
$$\widehat{t}(𝐱^{},t)\widehat{t}(𝐱,t)_V=0,$$
(17)
therefore
$$(\mathrm{\Delta }_LT)_{\mathrm{syn}}^2=(\widehat{\tau }(𝐱^{},t)\widehat{\tau }(𝐱,t))^2_V,$$
(18)
and relation (7) takes the form
$$(\widehat{\tau }(𝐱^{},t)\widehat{\tau }(𝐱,t))^2_V\frac{T_P^{4/3}r^{2/3}}{c^2},$$
(19)
where
$$r=|𝐱^{}𝐱|.$$
(20)
A mathematical remark should be made here. The left-hand side of relation (19) can be written in the form
$$\widehat{\tau }^2(𝐱^{},t)+\widehat{\tau }^2(𝐱,t)\widehat{\tau }(𝐱^{},t)\widehat{\tau }(𝐱,t)\widehat{\tau }(𝐱,t)\widehat{\tau }(𝐱^{},t)_V,$$
(21)
involving vev’s of bilinear products of $`\widehat{\tau }`$’s taken at equal time. It is well known from quantum field theory that such vev’s of local field operators are, as a rule, divergent, and become finite only after renormalization. In the absence of a detailed theory, one cannot evaluate the individual vev’s in (21). However, if relation (7) of the K model, leading to (19), is correct, then the divergences in (21) should cancel and the finite part should give the right-hand side of (19). A similar remark applies to relation (14).
## 5 Uncertainties in the Phases of the Quantum <br>States and the Spread of the Relative Phases
In regular non-relativistic quantum mechanics, the pure state of an isolated physical system, constituted by $`N`$ microparticles with masses $`M_1`$, $`M_2`$, …, $`M_N`$, is usually represented by a Schrödinger wave function
$$\mathrm{\Psi }(x,t)=\mathrm{exp}\left(\frac{i}{\mathrm{}}\widehat{H}t\right)\mathrm{\Psi }(x,0),$$
(22)
where $`\widehat{H}`$ stands for the Hamiltonian of the system, and the evolution of $`\mathrm{\Psi }(x,t)`$ takes place on the Minkowskian space-time. Since the Schrödinger evolution is deterministic, $`\mathrm{\Psi }`$, and consequently also its relative phases between any pairs of points
$$x=(𝐱_1,\mathrm{},𝐱_N)$$
(23)
and $`x^{}`$ of the configuration space<sup>1</sup><sup>1</sup>1Spin variables are omitted, because they do not play a role in the K model. are sharply determined. We shall call such a wave function “perfectly coherent”.
In the K model, the quantum state has to propagate on the Károlyházy space-time having the discussed uncertainty in its structure. As we have seen, this uncertainty can be taken into account by substituting the effective time operator $`\widehat{t}(𝐱,t)`$ for the global Minkowskian time $`t`$:
$$t\widehat{t}(𝐱,t)=t+\widehat{\tau }(𝐱,t).$$
(24)
At this point we have to recall that on both sides of Equation (22) the rest energy phase factor
$$\mathrm{exp}(i\mathrm{\Phi }(t)):=\mathrm{exp}\left(\frac{i}{\mathrm{}}\underset{\mathrm{}=1}{\overset{N}{}}M_{\mathrm{}}c^2t\right)$$
(25)
is omitted, because being independent of $`x`$ and of $`p=i\mathrm{}(_1,\mathrm{},_N)`$, it drops out from all the observables. However, under the substitution (24), one has to put in $`\mathrm{\Phi }(t)`$
$$M_{\mathrm{}}c^2tM_{\mathrm{}}c^2\widehat{t}(𝐱_{\mathrm{}},t),$$
(26)
because the coordinate of the $`\mathrm{}`$-th particle is $`𝐱_{\mathrm{}}`$. This implies that
$$\mathrm{\Phi }(t)\mathrm{\Phi }(t)+\widehat{\mathrm{\Phi }}(x,t),$$
(27)
where
$$\widehat{\mathrm{\Phi }}(x,t)=\frac{c^2}{\mathrm{}}\underset{\mathrm{}=1}{\overset{N}{}}M_{\mathrm{}}\widehat{\tau }(𝐱_{\mathrm{}},t).$$
(28)
The rest energy phase $`\mathrm{\Phi }(t)`$ can be omitted again, but the $`x`$-dependent rest energy phase operator $`\widehat{\mathrm{\Phi }}(x,t)`$ should be kept.
The substitution $`t\widehat{t}(𝐱,t)`$ should have been carried out in the Schrödinger wave function (22), too. However, the non-relativistic matrix elements of $`\widehat{H}`$ are much smaller than the rest energies of the particles. For solid bodies their contribution has been estimated in in the framework of the $`(g_{\mu \nu })_\beta `$ model. In the present paper we shall neglect it.
Keeping only the contribution of the rest energy phase, one realizes that in the K model the Schrödinger wave function acquires an operator phase factor $`\mathrm{exp}(i\widehat{\mathrm{\Phi }}(x,t))`$. In other words, with any Schrödinger wave function $`\psi (x,t)`$ one has to associate a Károlyházy state
$$\widehat{\mathrm{\Psi }}_K(x,t)=\mathrm{exp}(i\widehat{\mathrm{\Phi }}(x,t))\mathrm{\Psi }(x,t).$$
(29)
Unlike $`\mathrm{\Psi }(x,t)`$, the K state $`\widehat{\mathrm{\Psi }}_K(x,t)`$ is not perfectly coherent. Through the operator phase factor it feels the uncertainty of the space-time structure. Its departure from perfect coherence between two points $`x,x^{}`$ of the configureation space can be assessed by the amount of the uncertainty of its relative phase, i.e. by the vacuum spread $`\mathrm{\Delta }_\mathrm{\Phi }(x,x^{},t)`$ of the relative phase operator
$$\widehat{\mathrm{\Phi }}_R(x,x^{},t)=\widehat{\mathrm{\Phi }}(x^{},t)+\phi (x^{},t)\widehat{\mathrm{\Phi }}(x,t)\phi (x,t)$$
(30)
of the K state between those points. Here
$$\phi (x,t)=\mathrm{arg}\mathrm{\Psi }(x,t)$$
(31)
denotes the phase of $`\mathrm{\Psi }(x,t)`$.
The $`x,x^{}`$ dependence of $`\mathrm{\Delta }_\mathrm{\Phi }`$ can be found with the help of relation (19). From Equations (28) and (10) it follows that the vev of $`\widehat{\mathrm{\Phi }}_R`$ is equal to the Schrödingerian relative phase,
$$\widehat{\mathrm{\Phi }}_R(x,x^{},t)_V=\phi (x^{},t)\phi (x,t),$$
(32)
which drops then out from the vacuum spread of $`\widehat{\mathrm{\Phi }}_R`$:
$`\mathrm{\Delta }_\mathrm{\Phi }^2(x,x^{},t)`$ $`=`$ $`(\widehat{\mathrm{\Phi }}_R(x,x^{},t)\phi (x^{},t)\phi (x,t))^2_V`$ (33)
$`=`$ $`(\widehat{\mathrm{\Phi }}(x^{},t)\widehat{\mathrm{\Phi }}(x,t))^2_V.`$
With Equation (28) one finds
$$\mathrm{\Delta }_\mathrm{\Phi }^2(x,x^{},t)=\frac{c^4}{\mathrm{}^2}\underset{i,\mathrm{}=1}{\overset{N}{}}M_iM_{\mathrm{}}(\widehat{\tau }(𝐱_i^{},t)\widehat{\tau }(𝐱_i,t))(\widehat{\tau }(𝐱_{\mathrm{}}^{},t)\widehat{\tau }(x_{\mathrm{}},t)_V.$$
(34)
The vev in the last expression is identically equal to
$`{\displaystyle \frac{1}{2}}(\widehat{\tau }(𝐱_i^{},t)\widehat{\tau }(𝐱_{\mathrm{}},t))^2+(\widehat{\tau }(𝐱_i,t)\widehat{\tau }(𝐱_{\mathrm{}}^{},t))^2`$
$`(\widehat{\tau }(𝐱_i^{},t)\widehat{\tau }(𝐱_{\mathrm{}}^{},t))^2(\widehat{\tau }(𝐱_i,t)\widehat{\tau }(𝐱_{\mathrm{}},t))^2_V,`$ (35)
and with relation (19) one obtains for the spread of the relative phase the formula
$$\mathrm{\Delta }_\mathrm{\Phi }^2(x,x^{})\mathrm{\Lambda }^{4/3}\frac{c^2}{\mathrm{}^2}\underset{i,\mathrm{}=1}{\overset{N}{}}M_iM_{\mathrm{}}(|𝐱_i^{}𝐱_{\mathrm{}}|^{2/3}\frac{1}{2}|𝐱_i𝐱_{\mathrm{}}|^{2/3}\frac{1}{2}|𝐱_i^{}𝐱_{\mathrm{}}^{}|^{2/3}).$$
(36)
The time argument of $`\mathrm{\Delta }_\mathrm{\Phi }`$ has been omitted since $`\mathrm{\Delta }_\mathrm{\Phi }`$ turned out to be time independent, although $`\widehat{\mathrm{\Phi }}(x,t)`$ may depend on $`t`$.
According to Equation (36), $`\mathrm{\Delta }_\mathrm{\Phi }`$ increases with the masses and with the number of the microparticles constituting the system, and for a given system it increases with the distances $`|𝐱_i^{}𝐱_{\mathrm{}}|`$, that is with the separation between the points $`x,x^{}`$ in the configuration space. These are encouraging features concerning the expected loss of coherence between “macroscopically distinct” components of the quantum state of a macroscopic body.
With the notations of the present paper, the formula for $`\mathrm{\Delta }_\mathrm{\Phi }`$ derived in (see also ) with the help of the $`(g_{\mu \nu })_\beta `$’s reads
$$\mathrm{\Delta }_\mathrm{\Phi }^2(x,x^{})\mathrm{\Lambda }^{4/3}\frac{c^2}{\mathrm{}^2}\frac{d^3k}{k^{11/3}}|\mu _𝐤(x^{})\mu _𝐤(x)|^2,$$
(37)
where
$$\mu _𝐤(x)=\underset{\mathrm{}}{}M_{\mathrm{}}e^{i\mathrm{𝐤𝐱}_{\mathrm{}}}$$
(38)
is the Fourier transform of the mass distribution
$$\varrho (𝐱)=\underset{\mathrm{}=1}{\overset{N}{}}M_{\mathrm{}}\delta (𝐱𝐱_{\mathrm{}})$$
(39)
of $`N`$ pointlike particles of masses $`M_1,\mathrm{},M_N`$ in the configuration $`x=[𝐱_1,\mathrm{},𝐱_N]`$. From Equation (37) one sees that the uncertainty of the relative phase increases with the difference (of the absolute values of the Fourier transform) of the mass distribution of the $`N`$ particles in the configurations $`x`$ and $`x^{}`$.
The formula (36) for $`\mathrm{\Delta }_\mathrm{\Phi }`$ has been obtained previously with the help of the $`\{t_\beta \}`$ model by the present author , who realized then that the Fourier integral in the original formula (37) can be calculated analytically and is equal to the sum in formula (36). So, if the expressions (13) and (18) for $`\mathrm{\Delta }_LT`$ and $`(\mathrm{\Delta }_LT)_{\mathrm{syn}}`$ are reliable, then the influence of the single, quantum mechanically uncertain space-time on the relative phases has been correctly mimicked by both classical space-time models.
## 6 Coherence Properties of the K states, Coherence Cells and Cell Lengths
### 6.1 K States with Nearly Perfect and with Destroyed Coherence. The Coherence Cell
One sees from formula (37) that for any pairs of points $`x^{}x`$, $`\mathrm{\Delta }_\mathrm{\Phi }>0`$. This means that a K state $`\widehat{\mathrm{\Psi }}_K`$, normalized to $`1`$, is never perfectly coherent.
Concerning the norm of $`\widehat{\mathrm{\Psi }}_K`$, from (29) one gets
$$\widehat{\mathrm{\Psi }}_K^+(x,t)\widehat{\mathrm{\Psi }}_K(x,t)=|\mathrm{\Psi }(x,t)|^2,$$
(40)
because the unitary operator phase factor drops out from the product. Therefore, the weight $`w_\mathrm{\Omega }`$ of a K state in a domain $`\mathrm{\Omega }`$ of the configuration space is equal to the weight of the Schrödinger wave function associated with $`\widehat{\mathrm{\Psi }}_K`$,
$$w_\mathrm{\Omega }=_\mathrm{\Omega }𝑑x|\mathrm{\Psi }(x,t)|^2,$$
(41)
and $`\widehat{\mathrm{\Psi }}_K`$ is normalized to $`1`$ together with $`\mathrm{\Psi }`$. (Pedantically, in (40) one should consider the vev of $`\widehat{\mathrm{\Psi }}_K^+(x,t)\widehat{\mathrm{\Psi }}_K(x,t)`$, but it is obviously equal to the product itself.)
Let us now see how the coherence of the K states can be characterized.
(1) If $`\widehat{\mathrm{\Psi }}_K`$ occupies<sup>2</sup><sup>2</sup>2As a rule, $`\mathrm{\Psi }`$ is different from zero almost everywhere. Therefore, strictly speaking, $`\mathrm{\Psi }`$, and then $`\widehat{\mathrm{\Psi }}_K`$ too, occupy the whole configuration space. In our terminology the “domain occupied by $`\mathrm{\Psi }`$” is the smallest domain in which the weight $`w`$ of $`\mathrm{\Psi }`$ is close to the maximal weight $`1`$ (e.g. $`w=110^4`$). The expansion and the shrinking (or the localization) of $`\mathrm{\Psi }`$ means that this domain expands or shrinks. a domain $`\mathrm{\Omega }`$ of the configuration space such that
$$\mathrm{\Delta }_\mathrm{\Phi }(x,x^{})\pi \text{ for all }x,x^{}\mathrm{\Omega },$$
(42)
then in good approximation the uncertainties of the relative phases of $`\widehat{\mathrm{\Psi }}_K`$ can be neglected, and the relative phases of $`\widehat{\mathrm{\Psi }}_K`$ are practically equal to those of $`\mathrm{\Psi }`$. In such a case we shall say that the coherence of $`\widehat{\mathrm{\Psi }}_K`$ is nearly perfect.
(2) If $`\widehat{\mathrm{\Psi }}_K`$ occupies a domain containing non-overlapping subdomains $`\mathrm{\Omega }`$, $`\mathrm{\Omega }^{}`$ such that
$$\mathrm{\Delta }_\mathrm{\Phi }(x,x^{})\pi \text{ for all }x\mathrm{\Omega }\text{ and all }x^{}\mathrm{\Omega }^{},$$
(43)
then the relative phases of $`\widehat{\mathrm{\Psi }}_K`$ between these subdomains are completely uncertain, the coherence between the components of $`\widehat{\mathrm{\Psi }}_K`$ belonging to these subdomains is destroyed. Notice, however, that within smaller subdomains inside which $`\mathrm{\Delta }_\mathrm{\Phi }<\pi `$, a certain degree of coherence persists, and inside sufficiently small subdomains even $`\mathrm{\Delta }_\mathrm{\Phi }\pi `$ holds, so that within such a small domain $`\widehat{\mathrm{\Psi }}_K`$ is near to perfect coherence.
The maximal domains $`\mathrm{\Omega }_c`$ of the configuration space such that
$$\mathrm{\Delta }_\mathrm{\Phi }(x,x^{})\pi \text{ for all }x,x^{}\mathrm{\Omega }_c,$$
(44)
have been called “coherence cells” in . A K state occupying a single cell is still not incoherent, but if it occupies non-overlapping cells, then it is incoherent (it has incoherent components of non-negligible weights).
The size and the shape of the coherence cell depends strongly on the composition of the physical system considered. Below we shall look at the cells of microparticles and of a class of solid objects. As we shall see, the cells of these systems are spherical, and can therefore be characterized by a single parameter, the diameter of the cell.
### 6.2 The Coherence Cell and the Cell Length of a Single Microparticle
For a single microparticle of mass $`M=M_1`$, one easily finds from formula (36) that
$$\mathrm{\Delta }_\mathrm{\Phi }(a)\frac{\mathrm{\Lambda }^{2/3}}{L}a^{1/3},$$
(45)
where (with $`𝐱=𝐱_1`$)
$$a=|𝐱^{}𝐱|,$$
(46)
and
$$L=\frac{\mathrm{}}{Mc}$$
(47)
is the Compton wavelength of the particle. Notice that $`\mathrm{\Delta }_\mathrm{\Phi }`$ increases monotonically with the distance $`a`$. The coherence cell is therefore a sphere of diameter $`a_c`$ in the configuration space $`x=(𝐱)`$, where $`a_c`$ is equal to the value of $`a`$ at which $`\mathrm{\Delta }_\mathrm{\Phi }`$ reaches the value $`\pi `$:
$$\mathrm{\Delta }_\mathrm{\Phi }(a_c)=\pi 1.$$
(48)
From (45) one finds that
$$a_c\left(\frac{L}{\mathrm{\Lambda }}\right)^2L.$$
(49)
So, for a microparticle the coherence cell is characterized by a single parameter, the cell length $`a_c`$. In the term “cell diameter” has been used. However, there are physical systems the coherence cell of which is not spherical, and then $`a_c`$ is not a diameter.
For the electron $`L10^{11}`$ cm, and
$$a_c10^{33}\mathrm{cm}.$$
(50)
The cell length of the electron, and also of the other microparticles, has a supraastronomical value. Therefore, any realistic Schrödinger wave function $`\mathrm{\Psi }`$ of an isolated microparticle occupies only a tiny part of a coherence cell, a part inside which $`\mathrm{\Delta }_\mathrm{\Phi }\pi `$, and the K state associated with $`\mathrm{\Psi }`$ is always practically perfectly coherent. Due to the very small masses of the microparticles, the uncertainty of the space-time structure has no appreciable effect on a single microparticle. For the same reason this is also true for any isolated microsystem (for a system consisting of a few microparticles, free or bound).
### 6.3 The Coherence Cell and the Cell Length of Spherical, Homogeneous Solid Objects
In a homogeneous object there are $`N1`$ identical microscopic constituents (e.g. molecules). Formula (36) for $`\mathrm{\Delta }_\mathrm{\Phi }`$ then becomes
$$\mathrm{\Delta }_\mathrm{\Phi }^2(x,x^{})\frac{\mathrm{\Lambda }^{4/3}}{L_{\mathrm{micro}}^2}\underset{i,\mathrm{}=1}{\overset{N}{}}\left(|𝐱_i^{}𝐱_{\mathrm{}}|^{2/3}\frac{1}{2}|𝐱_i𝐱_{\mathrm{}}|^{2/3}\frac{1}{2}|𝐱_i^{}𝐱_{\mathrm{}}^{}|^{2/3}\right),$$
(51)
where
$$L_{\mathrm{micro}}=\frac{\mathrm{}}{M_{\mathrm{micro}}c}$$
(52)
is the Compton wavelength of a constituent of mass $`M_{\mathrm{micro}}`$.
In a solid object the constituents are at, or very close to, their equilibrium positions. Consequently, the Schrödinger wave function of an isolated solid object is practically zero everywhere, except in such points $`x=(𝐱_1,\mathrm{},𝐱_N)`$ of the configuration space, in which the constituents are at their equilibrium positions. Therefore, in the case of a homogeneous, spherical solid object of radius $`R`$, in $`\mathrm{\Delta }_\mathrm{\Phi }(x,x^{})`$ the $`𝐱_i`$ coordinates belonging to $`x`$ are distributed uniformly in the volume of a sphere of radius $`R`$, and the $`𝐱^{}`$ coordinates belonging to $`x^{}`$ fill uniformly an other such sphere in the three dimensional $`XYZ`$ space.
For a solid object of arbitrary shape, two equilibrium configurations $`x,x^{}`$ differ from each other by a translation of their center of mass (c.m.) and by a rotation leaving the c.m. fixed. For a spherical homogeneous object, $`\mathrm{\Delta }_\mathrm{\Phi }(x,x^{})`$ does not change appreciably under a rotation, because $`\mathrm{\Delta }_\mathrm{\Phi }`$ is invariant under any permutation of the $`𝐱`$ coordinates among themselves and of the $`𝐱^{}`$ coordinates among themselves. Therefore, we have to consider only such configurations which differ from each other by a translation
$$𝐱_i^{}=𝐱_i+𝒂,$$
(53)
where
$$𝒂=𝐱_{\mathrm{c}.\mathrm{m}.}^{}𝐱_{\mathrm{c}.\mathrm{m}.}$$
(54)
is the vector joining the centers of the spheres, which are also the centers of mass of the objects in the two configurations.
With (53) the sum $``$ in (51) becomes
$$=\underset{i,\mathrm{}=1}{\overset{N}{}}\left(|𝐱_i𝐱_{\mathrm{}}+𝒂|^{2/3}|𝐱_i𝐱_{\mathrm{}}|^{2/3}\right).$$
(55)
The $`𝐱_i`$ and the $`𝐱_{\mathrm{}}`$ coordinates fill a sphere uniformly, $`N1`$, and the expression under the sum is a continuous, slowly varying function of the coordinates. Therefore, the double sum in (55) can be replaced, in a good approximation, by a double integral. With $`𝐱_i𝐫`$, $`𝐱_{\mathrm{}}𝐫^{}`$, one gets
$$=\frac{N^2}{V^2}_Vd^3r_Vd^3r^{}\left(|𝐫𝐫^{}+𝒂|^{2/3}|𝐫𝐫^{}|^{2/3}\right),$$
(56)
where the integrals have to be taken over the volume of the same sphere of radius $`R`$. $``$ can be calculated for any value of $`𝒂`$, but it is more enlightening to evaluate it in two extreme situations, namely when $`|𝒂|aR`$ and when $`aR`$.
a) The case $`aR`$
In this case one has also $`a|𝐫𝐫^{}|`$, because in (56) $`|𝐫𝐫^{}|2R`$. Therefore, in good approximation,
$$|𝐫𝐫^{}+𝒂|^{2/3}|𝐫𝐫^{}|^{2/3}=a^{2/3},$$
(57)
and
$$=N^2a^{2/3}.$$
(58)
Noticing that
$$\frac{N}{L_{\mathrm{micro}}}=\frac{1}{L},$$
(59)
where $`L`$ is the Compton wavelength correspondig to the mass $`M=NM_{\mathrm{micro}}`$ of the object considered, Equation (51) becomes
$$\mathrm{\Delta }_\mathrm{\Phi }(x,x^{})\mathrm{\Lambda }^{2/3}\frac{a^{1/3}}{L},aR.$$
(60)
This formula is formally identical with the one obtained for a single microparticle, but $`a`$ stands now for the distance between the centers of mass of the object in the two configurations $`x,x^{}`$. The coherence cell is again a sphere of diameter
$$a_c\left(\frac{L}{\mathrm{\Lambda }}\right)^2L,a_cR,$$
(61)
but now in the center of mass coordinate subspace of the configuration space.
b) The case $`aR`$
In this case in (56) $`a|𝐫𝐫^{}|`$, except for a small subdomain of the integration domain. The detailed calculation shows that one can forget about the violation of the condition $`a|𝐫𝐫^{}|`$. Expanding $`|𝐫𝐫^{}+𝒂|^{2/3}`$ in powers of $`a`$, one finds that the leading contribution to $``$ is of the order $`a^2`$. Since the dimension of $``$ is $`cm^{3/2}`$, and apart from $`a`$ the only length parameter in (56) is $`R`$, one finds that
$$N^2\frac{a^2}{R^{4/3}},aR.$$
(62)
Equation (51) takes now the form
$$\mathrm{\Delta }_\mathrm{\Phi }(x,x^{})\left(\frac{\mathrm{\Lambda }}{R}\right)^{2/3}\frac{a}{L},aR.$$
(63)
$`\mathrm{\Delta }_\mathrm{\Phi }`$ increases again monotonically with $`a`$, and reaches the value $`\pi 1`$ when $`a`$ is equal to
$$a_c\left(\frac{R}{\mathrm{\Lambda }}\right)^{2/3}L,a_cR.$$
(64)
Formulas (61) and (64) for the cell length of spherical, homogeneous objects have been presented in as “the most important results” of the model. Their physical meaning has been discussed in due detail in . Here we recall only that $`a_cR`$ is the region where quantum behavior, $`a_cR`$ — the region where classical behavior dominates. Indeed, if $`a_cR`$, then the Schrödinger wave function $`\mathrm{\Psi }(x,t)`$ may have an uncertainty of the order of $`a_c`$ in the position of the c.m., much larger than the geometrical size $`R`$ of the object, without making the K state associated with $`\mathrm{\Psi }`$ incoherent. Large coherent uncertainties correspond to quantum behavior. On the contrary, when $`a_cR`$, the $`K`$ state becomes incoherent when the uncertainty of the position of the c.m. is still much smaller than the geometrical size $`R`$ of the object. In other words, when $`a_cR`$, the coherent, quantum mechanical uncertainty of the c.m. is much smaller than $`R`$. Small quantum mechanical positional uncertainty is a characteristic of classical behavior.
### 6.4 Tiny Grains, Macroscopic Bodies and the Transition Region between Quantum and Classical Behavior
It can be shown that formulas (61) and (64) hold, with small corrections absorbable in the symbol $``$, in the whole region $`a_cR`$ and $`a_cR`$, respectively. In particular, they hold also in the region $`a_cR`$. It follows from the discussion in the preceding subsection that the latter region is, for spherical homogeneous objects, the transition region between quantum and classical behavior. In this region the maximal quantum mechanical uncertainty of the position of the c.m. is of the same order of magnitude as the geometrical size of the object. For usual terrestrial densities $`\varrho 1\text{ g/cm}\text{3}`$, one easily finds from Equation (61) (as well as from (64), of course) that
$`a_c^{\mathrm{tr}}`$ $``$ $`R^{\mathrm{tr}}10^5\mathrm{cm},`$ (65)
$`M^{\mathrm{tr}}`$ $`=`$ $`{\displaystyle \frac{4\pi }{3}}\varrho (R^{\mathrm{tr}})^310^{14}\mathrm{g}.`$ (66)
Thus, the transition mass region is the region of the colloidal grains and dust particles. The region $`a_cR`$ of quantum behavior corresponds to tiny grains with masses much smaller than $`M^{\mathrm{tr}}`$. To give an example, for a grain of $`R10^6`$ cm and of mass $`M\varrho R^310^{18}`$ g, one finds<sup>3</sup><sup>3</sup>3The estimate $`a_c10`$ km, quoted in , is erroneous. from (61) that $`a_c10^5`$ km, indeed much larger than $`R`$. If the wave function of such an isolated grain expanded, say, over $`10^3`$ km in the c.m. coordinate subspace, its associated K state would remain still very coherent. On the contrary, for a ball of $`M1`$ g and radius $`R1`$ cm, one finds from (64) that
$$a_c10^{16}\mathrm{cm},$$
(67)
a value much smaller indeed than $`R`$. The K model states that two positions of this ball, with separation $`2a_c210^{16}`$ cm between the centers of mass, correspond already to an incoherent K state, these positions are already “macroscopically distinct”. We are in the region of predominantly classical behavior, with $`a_cR`$, $`MM^{\mathrm{tr}}`$. It should be noted that formula (64) for $`a_c`$ is applicable only to bodies of moderate size. Above $`R=1`$ m the elastic vibrations of the body are not negligible .
## 7 The Stochastic Modification of the Schrödinger Evolution of Isolated Systems
The mode of the combination of a stochastic process with the Schrödinger evolution is suggested by the behavior of the K state
$$\widehat{\mathrm{\Psi }}_K(x,t)=\mathrm{exp}(i\widehat{\mathrm{\Phi }}(x,t))\mathrm{\Psi }(x,t)$$
(68)
during the Schrödinger evolution of its associated $`\mathrm{\Psi }`$ function. As it is well known, during that evolution $`\mathrm{\Psi }`$ expands<sup>4</sup><sup>4</sup>4For appropriately awkward initial wave functions, a transient shrinking precedes the steady expansion., at least in the c.m. subspace of the configuration space. According to Equations (68) and (40), $`\widehat{\mathrm{\Psi }}_K`$ expands (or shrinks) exactly in the same way as $`\mathrm{\Psi }`$. However, while $`\mathrm{\Psi }`$ remains perfectly coherent, the coherence of $`\widehat{\mathrm{\Psi }}_K`$ deteriorates during the expansion, and when it occupies a domain larger than a coherence cell, it has already incoherent components.
The basic idea for the introduction of the stochastic element is that the K state $`\widehat{\mathrm{\Psi }}_K`$ counterbalances its loss of coherence by stochastically and instantaneously localizing itself, of course together with $`\mathrm{\Psi }`$, to one of the coherence cells lying in the domain to which the state expanded. After a localization, $`\mathrm{\Psi }`$ expands again under the Schrödinger equation, until the situation gets ripe for a new localization, and so on.
There are many, essentially equivalent ways to specify these expansion–localization cycles. In , for an isolated system the coherence cell of which is characterized by a single cell diameter $`a_c`$, an instantaneous random localization occurs whenever the quantum state occupies a domain of diameter $`2a_c`$. The localizations occur then at discrete moments of time, at intervals
$$\tau _c\frac{Ma_c^2}{\mathrm{}},$$
(69)
the time needed for the Schrödinger wave function of an isolated system to expand from size $`a_c`$ to $`2a_c`$.
It has been shown, among others by Ghirardi, Rimini and Weber and by Caves and Milburn , that the stochastic localizations from size $`2a_c`$ to $`a_c`$ occurring at equally spaced discrete moments of time, can be replaced by infinitesimal stochastic localizations from size $`a_c+da_c`$ to $`a_c`$ occurring continuously in time, blended with the continuous Schrödinger evolution. The application of the infinitesimal GRW localizations to the K model, in the case of a single cell length $`a_c`$, has been carried out in .
Let us now look at the expansion–localization cycles of various physical systems. Since for the electron $`a_c10^{33}`$ cm, a realistic wave function of an electron occupies only a tiny part of a coherence cell. Therefore, the Schrödinger evolution is not interrupted by localizations. However, if a ball of 1 g could be isolated, then as soon as its $`\mathrm{\Psi }`$ function would expand to a domain of diameter $`2a_c210^{16}`$ cm (in the c.m. subspace), a stochastic localization onto a single cell of diameter $`a_c`$ would take place. So, according to the K model, the localization of the c.m. of the ball would remain practically pointlike even if the ball were isolated. Notice that from $`\mathrm{\Delta }x\mathrm{\Delta }v\mathrm{}/M`$ one finds $`\mathrm{\Delta }v10^{11}`$ cm/sec. The uncertainty of the velocity would be very small, too.
Of course, a macroscopic ball cannot be perfectly isolated. The interplay of the Károlyházy law of time evolution with the effects caused by the surroundings has been discussed in .
Károlyházy suceeded in applying the basic idea of his time evolution law to systems of many quasi independent degrees of freedom. In particular, he worked out the process of the decay of the superposition of tracks in a cloud chamber. The detailed discussion is given only in . For a short outline, see .
## 8 Concluding Remarks
The expression of the structural uncertainties $`\mathrm{\Delta }_LT`$ and $`(\mathrm{\Delta }_LT)_{\mathrm{syn}}`$ through $`\widehat{\tau }`$ and $`|V`$ and the derivation of the formula for the spread $`\mathrm{\Delta }_\mathrm{\Phi }`$ of the relative phases without the use of a space-time model, is a tentative step towards turning the K model into a theory. A further step could concern the law of time evolution described in the preceding section. There, the phase operator $`\widehat{\mathrm{\Phi }}`$ producing the spread of the relative phases acts as a “hop-master”. It tells when $`\mathrm{\Psi }`$, and with it $`\widehat{\mathrm{\Psi }}_K`$, should jump in order to prevent the loss of coherence of $`\widehat{\mathrm{\Psi }}_K`$. It would be much better if one could set up a general stochastic differential equation for $`\widehat{\mathrm{\Psi }}_K`$, presumably of Itô type, and let the equation work. However, this does not seem to be possible without knowing more about the phase operator $`\widehat{\mathrm{\Phi }}`$, that is about the local time operator $`\widehat{\tau }(𝐱,t)`$ entering into $`\widehat{\mathrm{\Phi }}`$. The basic relations (1) and (7) made possible only to evaluate the vacuum spreads (15) and (19), involving $`\widehat{\tau }`$’s. This was sufficient to derive the formulas for $`\mathrm{\Delta }_\mathrm{\Phi }`$ and then for $`a_c`$, and to formulate with their help the hop-master’s rule. Of course, it is possible, after having calculated the cell length $`a_c`$ (or the cell lengths $`a_c^{(1)},a_c^{(2)},\mathrm{})`$ of a given system, to set up an Itô equation containing these $`a_c`$’s, which would smoothly reproduce the results of the bumpy hop-master’s rule. What we have in mind is a general equation for $`\widehat{\mathrm{\Psi }}_K`$, not containing the parameters $`a_c^{(i)}`$ of the particular system to which the equation is applied. It is probable that such an equation cannot be found without knowing more about the unified theory of general relativity and quantum theory.
Another open probelm is the relativistic generalization of the K model. This is a common open problem of the existing models with stochastic modification of the Schrödinger evolution. In the case of the K model, a specific task arises. One should find, first of all, a covariant description for the structural uncertainty of the Károlyházy space-time. Again, there is little hope for progress without the unified theory.
A direction into which progress can be made is the derivation of the basic relation (1) between $`T`$ and $`\mathrm{\Delta }_LT`$ . This relation has been deduced by Károlyházy in three different ways ; ; . On the one hand, it is reassuring that various approaches lead to the same result. On the other hand, each of these approaches, taken separately, has loose ends. For instance, in and relation (1) has been obtained from the study of the quantum behavior of the hand of a clock, and one can legitimately ask what would happen if the dial also entered the game. The answer is that the dial would not upset relation (1). This will be shown in a forthcoming paper, where the derivation(s) of relation (1) will be scrutinized and a comparison of Károlyházy’s clock with a Wigner–Salecker clock will also be made.
Relations (1) and (7) have been recently rediscovered by Ng and Van Dam . These authors obtain also the formula for the spread of the relative phases for a single particle (our formula (45), but they do not consider composite systems. Having only $`\mathrm{\Lambda }`$, $`a`$ and $`L`$ at their disposal, they had to choose $`(\mathrm{\Delta }_\mathrm{\Phi }1`$, $`aL)`$ as the condition for classical behavior, which leads then to the condition $`mm_P10^5`$ g, $`m_P`$ being the Planck mass. The authors remark that this condition is only a sufficient condition for classical behavior. Indeed, objects with masses much smaller than $`10^5`$ g are known to behave classically. As we have seen, in the K model the transition mass depends not only on $`m_P`$ (i.e. on $`\mathrm{\Lambda }`$), and for homogeneous spherical solid objects the sufficient and necessary condition for classical behavior is $`a_cR`$. It should be noted that in there are many interesting relations, not overlapping with the K model.
The author is indebted to F. Károlyházy for countless discussions scattered over the last three decades. This work was partially supported by the Hungarian Research Fund, under grants OTKA T016246 and OTKA T030374. |
warning/0002/math0002228.html | ar5iv | text | # Connections on locally trivial quantum principal fibre bundles
## 1 Introduction
Since the appearance of quantum groups there has been a hope that it should be possible to use them instead of the classical symmetry groups of physical theories, in particular for quantum field theories. It was expected that the greater variety of group-like structures should lead, perhaps, to greater flexibility in the formulation of physical theories, thereby paving the way to a better understanding of fundamental problems of quantum theory and gravitation.
In (Lagrangian) quantum field theory, symmetry groups can be considered to appear in a very natural geometrical scheme: They are structure groups of principal fibre bundles. Moreover, on the classical level, all fields are geometrical objects living on the principal bundle or on associated fibre bundles. Thus, it is natural to ask for a generalization of the notion of principal bundle to a noncommutative situation. Thereby, in order to avoid unnecessary restrictions, one should replace not only the structure group by a quantum group, but also the base manifold (space-time) by a noncommutative space, which may even be necessary for physical reasons (see , , and ).
In recent years, there have been several attempts to define such quantum principal bundles and the usual geometric objects that are needed to formulate gauge field theories on them, see , , , , , , and . Roughly following the same idea (“reversing the arrows”), the approaches differ in the details of the definitions. Closest to the classical idea that a locally trivial bundle should be imagined as being glued together from trivial pieces is the definition given in . There, one starts with the notion of a covering of a quantum space. Being in the context of C\*-algebras, a covering is defined to be a (finite) family of closed ideals with zero intersection, which is easily seen to correspond to finite coverings by closed sets in the commutative case. C\*-algebras which have such a covering can be reconstructed from their “restriction” to the elements of the covering by a gluing procedure. Such a reconstruction is not always possible for general (not C\*-)algebras, as was noticed in . The aim of was to introduce differential calculi over algebras with covering. Leaving the C\*-category, one is confronted with the above difficulty, called “noncompleteness of a covering”. Nevertheless, making use of “covering completions”, if necessary, a general scheme for differential calculi on quantum spaces with covering was developed, and the example of the gluing of two quantum discs, being homeomorphic to the quantum sphere $`S_{\mu c}^2`$, $`c>0`$, including the gluing of suitable differential calculi on the discs, was described in detail.
In , a locally trivial quantum principal fibre bundle having as base $`B`$ such a quantum space with covering, and as fibre a compact quantum group $`H`$, is defined as a right $`H`$-comodule algebra with a covering adapted to the covering of the base. “Adapted” means that the ideals defining the covering appear as kernels of “locally trivializing” homomorphisms such that the intersections of these kernels with the embedded base are just the embeddings of the ideals defining the covering of $`B`$. Given such a locally trivial principal fibre bundle, one can define analogues of the classical transition functions which have the usual cocycle properties. Reversely, given such a cocycle one can reconstruct the bundle. The transition functions are algebra homomorphisms $`HB_{ij}`$, where $`B_{ij}`$ is the algebra corresponding to the “overlap” of two elements of the covering of $`B`$. It turns out that they must have values in the center of $`B_{ij}`$, which is related to the fact that principal bundles with structure group $`H`$ are determined by bundles which have as structure group the classical subgroup of $`H`$, see .
The aim of the present paper is to introduce notions of differential geometry on locally trivial bundles in the sense of in such a way that all objects can be glued together from local pieces.
Let us describe the contents of the paper: In Section 2, locally trivial principal bundles are defined slightly different from . Not assuming C\*-algebras, we add to the definition of the assumption that the “base” algebra is embedded as the algebra of right invariants into the “total space” algebra. This assumption has to be made in order to come back to the usual notion in the classical case, as is shown by an example. We prove a technical proposition about the kernels of the local trivializing homomorphisms which in turn makes it possible to prove a reconstruction theorem for locally trivial principal bundles in terms of transition functions in the context of general algebras.
The aim of Section 3 is to introduce differential calculi on locally trivial quantum principal bundles. They are defined in such a way that they are uniquely determined by giving differential calculi on the “local pieces” of the base and a right covariant differential calculus on the Hopf algebra (assuming that the calculi on the trivializations are graded tensor products). Uniqueness follows from the assumption that the local trivializing homomorphisms should be differentiable and that the kernels of their differential extensions should form a covering of the differential calculus on the total space, i. e. the differential calculus is “adapted” in the sense of . This covering need not be complete. Thus, in order to have reconstructability, one has to use the covering completion, which in general is only a differential algebra.
Section 4 is the central part of the paper. Whereas in the classical situation there is a canonically given vertical part in the tangent space of a bundle, in the dual algebraic situation there is a canonically given horizontal subbimodule in the bimodule of forms of first degree on the bundle space. We start with the definition of left (right) covariant derivatives, which involves a Leibniz rule, a covariance condition, invariance of the submodule of horizontal forms, and a locality condition. Covariant derivatives can be characterized by families of linear maps $`A_i:H\mathrm{\Gamma }(B_i)`$ satisfying $`A_i(1)=0`$ and a compatibility condition being analogous to the classical relation between local connection forms. At this point a bigger differential algebra on the basis $`B`$ appears, which is maximal among all the (LC) differential algebras being embeddable into the differential structure of the total space. Next we define left (right) connections as a choice of a projection of the left (right) $`𝒫`$-module of one-forms onto the submodule of horizontal forms being covariant under the right coaction and satisfying a locality condition. This is equivalent to the choice of a vertical complement to the submodule of horizontal forms. Left and right connections are equivalent. With this definition it is possible to reconstruct a connection from connections on the local pieces of the bundle. The corresponding linear maps $`A_i:H\mathrm{\Gamma }(B_i)`$ satisfy the conditions for the $`A_i`$ of covariant derivatives, and in addition $`RkerA_i`$ ($`S^1(R)kerA_i`$), where $`R`$ is the right ideal in $`H`$ defining the right covariant differential calculus there. Thus, connections are special cases of covariant derivatives. There is a corresponding notion of connection form as well as a corresponding notion of an exterior covariant derivative. The curvature can be defined as the square of the exterior covariant derivative, and is nicely related to a curvature form being defined by analogues of the structure equation. The local components of the curvature are related to the local connection forms in a nice way, and they are related among themselves by a homogeneous formula analogous to the classical one.
Finally, in Section 5, we give an example of a locally trivial principal bundle with a connection. The basis of the bundle, constructed in , is a C\*-algebra glued together from two copies of a quantum disc. The structure group is the classical group $`U(1)`$, and the bundle is defined by giving one transition function, which is sufficient because the covering of the basis has only two elements. Since all other coverings appearing in the example then have also two elements, there are no problems with noncomplete coverings. The differential calculus on the total space is determined by differential ideals in the universal differential calculi over the two copies of the quantum disc and the structure group. For the group, the ideal is chosen in nonclassical way. Then, a connection is defined by giving explicitely the two local connection forms on the generator of $`P(U(1))`$ and extending them using the properties a local connection form should have. The curvature of this connection is nonzero.
In the appendix, the relevant facts about coverings and gluings of algebras and differential algebras are collected, for the convenience of the reader. Details can be found in . Moreover, we recall there some well-known facts about covariant differential calculi on quantum groups.
In the following, algebras are always assumed to be over $``$, associative and unital. Ideals are assumed to be two-sided, up to some occasions, where their properties are explicitely specified.
## 2 Locally trivial quantum principal fibre bundles
Following the ideas of we introduce in this section the definition of a locally trivial quantum principal fibre bundle and prove propositions about the existence of trivial subbundles and about the reconstruction of the bundle. Essentially, this is contained in , up to some modifications: We do not assume C\*-algebras, and we add to the axioms the condition that the embedded base algebra coincides with the subalgebra of coinvariants. As structure group we take a general Hopf algebra.
In the sequel we use the results of , see in the appendix. We recall here that for an algebra $`B`$ with a covering $`(J_i)_{iI}`$, there are canonical mappings $`\pi _i:BB_i:=B/J_i`$, $`\pi _j^i:B_iB_{ij}:=B/(J_i+J_j)`$, $`\pi _{ij}:BB_{ij}`$.
###### Definition 1
A locally trivial quantum principal fibre bundle (QPFB) is a tupel
$$(𝒫,\mathrm{\Delta }_𝒫,H,B,\iota ,(\chi _i,J_i)_{iI})$$
(1)
where $`B`$ is an algebra, $`H`$ is a Hopf algebra, $`𝒫`$ is a right $`H`$ comodule algebra with coaction $`\mathrm{\Delta }_𝒫`$, $`(J_i)_{iI}`$ is a complete covering of $`B`$, and $`\chi _i`$ and $`\iota `$ are homomorphisms with the following properties:
$`\chi _i:𝒫`$ $``$ $`B_iH\text{surjective},`$
$`\iota :B`$ $``$ $`𝒫\text{injective},`$
$`(id\mathrm{\Delta })\chi _i`$ $`=`$ $`(\chi _iid)\mathrm{\Delta }_𝒫,`$
$`\chi _i\iota (a)`$ $`=`$ $`\pi _i(a)1aB,`$
$`(ker\chi _i)_{iI}`$ $`\text{complete covering of }𝒫,`$
$`\iota (B)`$ $`=`$ $`\{f𝒫|\mathrm{\Delta }_𝒫(f)=fI\}.`$
Such a tupel we often denote simply by $`𝒫`$. Occasionally, $`𝒫`$, $`B`$ and $`H`$ are called total space, base space and structure group of the bundle.
The last assumption in Definition 1 does not appear in the definition of QPFB given in . It is however used by other authors ( , , ). Already in the classical case this condition is needed to guarantee the transitive action of the structure group on the fibres, as shows the following example.
Example: Let $`M`$ be a compact topological space covered by two closed subsets $`U_1`$ and $`U_2`$ being the closure of two open subsets covering $`M`$. Define $`M_0=U_1\dot{}U_2`$ (disjoint union). $`M`$ is obtained from $`M_0`$ identifying all corresponding points of $`U_1`$ and $`U_2`$. There is a natural projection $`M_0M`$. Let us consider the algebras of continuous functions $`C(M)`$ and $`C(M_0)`$ over $`M`$ and $`M_0`$ respectively. There exists an injective homomorphism $`\kappa :C(M)C(M_0)`$ being the pull back of the natural projection $`M_0M`$. Suppose we have constructed a principal fibre bundle $`P`$ over $`M_0`$ with structure group $`G`$, which is trivial on each of the disjoint components. Then we have an injective homomorphism $`\iota _o:C(M_0)C(P)`$ and two trivialisations $`\chi _{1,2}:C(P)C(U_{1,2})C(G)`$ with the properties assumed in Definition 1. The injective homomorphismus $`\iota :C(M)C(P)`$, $`\iota :=\iota _0\kappa `$, fullfills all the assumptions in Definition 1 up to the last one, and one obtains a fibration $`P`$ over the base manifold $`M`$ which is not a principal fibre bundle.
###### Proposition 1
Let $`𝒫_c`$ be the covering completion of $`𝒫`$ with respect to the complete covering $`(ker\chi _i)_{iI}`$. Let $`K:𝒫𝒫_c`$ be the corresponding isomorphism. The tupel
$$(𝒫_c,\mathrm{\Delta }_{𝒫_c},H,B,\iota _c,(\chi _{i_c},J_i)_{iI}),$$
where
$`\mathrm{\Delta }_{𝒫_c}`$ $`=`$ $`(Kid)\mathrm{\Delta }_𝒫K^1,`$
$`\chi _{i_c}`$ $`=`$ $`\chi _iK^1,`$
$`\iota _c`$ $`=`$ $`K\iota ,`$
is a locally trivial QPFB.
The proof ist trivial (transport of the structure using $`K`$).
###### Definition 2
A locally trivial QPFB $`𝒫`$ is called trivial if there exists an isomorphism $`\chi :𝒫BH`$ such that
$`\chi \iota `$ $`=`$ $`id1,`$
$`(\chi id)\mathrm{\Delta }_𝒫`$ $`=`$ $`(id\mathrm{\Delta })\chi .`$
Remark: A locally trivial QPFB with $`cardI=1`$, i.e. with trivial covering of $`B`$, is trivial.
Triviality of the covering means that it consists of only one ideal $`J=0`$. Moreover, there is only one trivializing epimorphism $`\chi :𝒫BH`$ which necessarily fulfills $`ker\chi =0`$.
There are several trivial QPFB related to a locally trivial QPFB. Define $`𝒫_i:=𝒫/ker\chi _i`$. Then $`\stackrel{~}{\chi _i}:𝒫_iB_iH`$ defined by
$$\stackrel{~}{\chi }_i(f+ker\chi _i):=\chi _i(f)$$
(2)
is a well defined isomorphism. $`\iota _i:B_i𝒫_i`$ defined by
$$\iota _i(b):=\stackrel{~}{\chi }_i^1(b1)$$
is injective and fulfills $`\stackrel{~}{\chi }_i\iota _i=id1`$. Moreover $`\mathrm{\Delta }_{𝒫_i}:𝒫_i𝒫_iH`$ is well defined by
$$\mathrm{\Delta }_{𝒫_i}(f+ker\chi _i):=\mathrm{\Delta }_𝒫(f)+ker\chi _iH,$$
because from $`(id\mathrm{\Delta })\chi _i=(\chi _iid)\mathrm{\Delta }_𝒫`$ follows $`\mathrm{\Delta }_𝒫(ker\chi _i)ker\chi _iH`$. Obviously, $`\mathrm{\Delta }_{𝒫_i}`$ is a right coaction. Moreover, $`(\stackrel{~}{\chi }_iid)\mathrm{\Delta }_{𝒫_i}=(id\mathrm{\Delta })\stackrel{~}{\chi }_i`$, and $`\iota _i(B_i)=\{f𝒫_i|\mathrm{\Delta }_{𝒫_i}(f)=f1\}`$. Thus $`(𝒫_i,\mathrm{\Delta }_{𝒫_i},H,B_i,\iota _i,(\stackrel{~}{\chi }_i,0))`$ is a trivial QPFB.
Let $`𝒫_{ij}:=𝒫/(ker\chi _i+ker\chi _j)`$. Then there is an isomorphism $`\stackrel{~}{\chi }_{ij}^i:𝒫_{ij}(B_iH)/\chi _i(ker\chi _j)`$ given by
$$\stackrel{~}{\chi }_{ij}^i(f+ker\chi _i+ker\chi _j):=\chi _i(f)+\chi _i(ker\chi _j).$$
(3)
It is natural to expect that $`P_{ij}`$ should be a trivial bundle isomorphic to $`B_{ij}H`$. In fact, we will show that there is a natural isomorphism $`(B_iH)/\chi _i(ker\chi _j)B_{ij}H`$, leading to trivialization maps $`\chi _{ij}^i:𝒫_{ij}B_{ij}H`$. Let us introduce the natural projections $`\pi _{i_𝒫}:𝒫𝒫_i`$, $`\pi _{ij_𝒫}:𝒫𝒫_{ij}`$ and $`\pi _{j_𝒫}^i:𝒫_i𝒫_{ij}`$. Obviously, $`\stackrel{~}{\chi }_i\pi _{i_𝒫}=\chi _i`$, $`\pi _{i_𝒫}=\stackrel{~}{\chi }_i^1\chi _i`$ and $`\pi _{ij_𝒫}=\pi _{j_𝒫}^i\pi _{i_𝒫}`$. We will need the following lemma.
###### Lemma 1
Let $`B`$ be an algebra and $`H`$ be a Hopf algebra. Let $`JBH`$ be an ideal with the property
$$(id\mathrm{\Delta })JJH.$$
Then there exists an ideal $`IB`$ such that $`J=IH`$. This ideal is uniquely determined and equals $`(id\epsilon )(J)`$.
Proof: Let $`m_H:HHH`$ be the algebra product in $`H`$. It is not difficult to verify that $`I:=(id\epsilon )(J)`$ is an ideal in $`B`$. We will show $`J=IH`$ First we prove $`JIH`$. Because of $`(id\epsilon id)(id\mathrm{\Delta })=id`$ and $`(id\mathrm{\Delta })JJH`$ we have $`(id\epsilon id)(id\mathrm{\Delta })J=JIH`$. $`IHJ`$ is a consequence of $`I1J`$, which is proved as follows: A general element of $`I`$ has the form $`_ka_k\epsilon (h_k)`$ where $`_ka_kh_kJ`$. Because of
$$\underset{k}{}a_k\epsilon (h_k)1=\underset{k}{}(a_kh_{k_1})(1S(h_{k_2}))$$
and
$$(id\mathrm{\Delta })(\underset{k}{}a_kh_k)=\underset{k}{}a_kh_{k_1}h_{k_2}JH,$$
$`_ka_k\epsilon (h_k)1`$ is an element of $`J`$. $`\mathrm{}`$
###### Proposition 2
$`𝒫_{ij}`$ is a trivial QPFB, i.e. there exist
$`\chi _{ij}^i:𝒫_{ij}`$ $``$ $`B_{ij}H,`$
$`\mathrm{\Delta }_{𝒫_{ij}}:𝒫_{ij}`$ $``$ $`𝒫_{ij}H,`$
$`\iota _{ij}:B_{ij}`$ $``$ $`𝒫_{ij},`$
such that the conditions of Definitions 1 and 2 are satisfied.
Remark: $`𝒫_{ij}`$ is a trivial QPFB in two ways by choosing $`\chi _{ij}^i`$ or $`\chi _{ij}^j`$. The composition of these maps just gives the transition functions.
Proof: Applying $`\chi _iid`$ to $`\mathrm{\Delta }_𝒫(ker\chi _i)ker\chi _iH`$ and using $`(id\mathrm{\Delta })\chi _i=(\chi _iid)\mathrm{\Delta }_𝒫`$ it follows that $`(id\mathrm{\Delta })\chi _i(ker\chi _j)\chi _i(ker\chi _j)H`$. By Lemma 1, there exist ideals $`\stackrel{~}{K}_j^iB_i`$ such that $`\chi _i(ker\chi _j)=\stackrel{~}{K}_j^iH`$. $`K_j^i:=\pi _j^i(\stackrel{~}{K}_j^i)`$ is an ideal in $`B_{ij}`$.
We show now $`\pi _i(J_j)\stackrel{~}{K}_j^i`$: According to Lemma 1, we have $`\stackrel{~}{K}_j^i=(id\epsilon )(\chi _i(ker\chi _j))`$. We have to show that for a $`bJ_j`$ there exists $`\stackrel{~}{b}ker\chi _j`$ with $`(id\epsilon )\chi _i(\stackrel{~}{b})=\pi _i(b)`$. It is obvious that we can take $`\stackrel{~}{b}=\iota (b)`$.
Using this inclusion, one finds that there is a canonical isomorphism $`(B_iH)/\chi _i(ker\chi _j)(B_{ij}/K_j^i)H`$ given by $`bh+\chi _i(ker\chi _j)(\pi _i(b)+K_j^i)h`$. Composing with $`\stackrel{~}{\chi }_{ij}^i`$ (see (3)), there results an isomorphism $`\chi _{ij}^i:𝒫_{ij}B_{ij}/K_j^iH`$ given by
$$\chi _{ij}^i(f+ker\chi _i+ker\chi _j):=(\pi _j^iid)\chi _i(f)+K_j^iH.$$
Our goal is now to show $`K_j^i=\pi _j^i(\stackrel{~}{K}_j^i)=0`$, so that $`\chi _{ij}^i`$ will become the isomorphism $`𝒫_{ij}B_{ij}H`$ wanted in the proposition. As a first step we will prove $`K_j^i=K_i^j`$. To this end, we note that
$$\stackrel{~}{\varphi }_{ji}:=\chi _{ij}^j\pi _{j_𝒫}^i\stackrel{~}{\chi _i}^1$$
is a homomorphism $`\stackrel{~}{\varphi }_{ji}:B_iHB_{ij}/K_i^jH`$ with $`ker\stackrel{~}{\varphi }_{ji}=\stackrel{~}{K}_j^iH`$. In terms of this homomorphism we define a homomorphism $`\psi _{ji}:B_{ij}B_{ij}/K_i^j`$ by
$$\psi _{ji}(a+J_i+J_j):=(id\epsilon )\stackrel{~}{\varphi }_{ji}((a+J_i)1).$$
It is easy to show that $`ker\psi _{ji}=K_j^i`$. On the other hand one shows that $`\psi _{ji}:B_{ij}B_{ij}/K_i^j`$ is the natural projection, and therefore $`K_j^i=K_i^j`$:
We calculate
$`\psi _{ji}(a+J_i+J_j)`$ $`=`$ $`(id\epsilon )\stackrel{~}{\varphi }_{ji}((a+J_i)1)`$
$`=`$ $`(id\epsilon )\chi _{ij}^j\pi _{j_𝒫}^i\stackrel{~}{\chi }_i^1((a+J_i)1)`$
$`=`$ $`(id\epsilon )\chi _{ij}^j\pi _{ij_𝒫}(\iota (a))`$
$`=`$ $`(id\epsilon )\chi _{ij}^j(\iota (a)+ker\chi _i+ker\chi _j)`$
$`=`$ $`(id\epsilon )((\pi _i^jid)\chi _j(\iota (a)+K_i^jH`$
$`=`$ $`(id\epsilon )((\pi _i^jid)(\pi _j(a)1)+K_i^jH)`$
$`=`$ $`\pi _{ij}(a)+K_i^j.`$
For showing $`K_j^i=0`$ we use the completeness of the covering $`(ker\chi _i)_{iI}`$. The covering completion of $`𝒫`$ is by definition
$$𝒫_c=\{(f_i)_{iI}\underset{iI}{}𝒫/ker\chi _i|\pi _{j_𝒫}^i(f_i)=\pi _{i_𝒫}^j(f_j)\}.$$
We introduce a locally trivial QPFB $`\stackrel{˘}{𝒫}𝒫_c`$ such that a comparison of $`\stackrel{˘}{𝒫}^{coH}=\{f\stackrel{˘}{𝒫}|\mathrm{\Delta }_{\stackrel{˘}{𝒫}}(f)=f1\}`$ with $`BB_c`$ allows to read off $`ker\psi _{ij}=K_j^i=0`$. Let $`\varphi _{ij}:B_{ij}/K_j^iHB_{ij}/K_j^iH`$ be the isomorphisms defined by
$$\varphi _{ij}:=\chi _{ij}^i\chi _{ij}^{j}{}_{}{}^{1}.$$
Using the identities
$$\chi _{ij}^i\pi _{j_𝒫}^i\stackrel{~}{\chi }_i^1=(\psi _{ij}id)(\pi _j^iid)$$
it is easy to verify that the algebra $`𝒫_c`$ is isomorphic to the algebra
$$\stackrel{˘}{𝒫}=\{(g_i)_{iI}\underset{iI}{}(B_iH)|(\psi _{ji}id)(\pi _j^iid)(g_i)=\varphi _{ij}(\psi _{ji}id)(\pi _i^jid)(g_j)\}$$
(cf. Lemma 1 in ), and the corresponding isomorphism $`\chi :𝒫_c\stackrel{˘}{𝒫}`$ is defined by $`\chi ((f_i)_{iI}):=(\stackrel{~}{\chi }_i(f_i))_{iI}`$. Transporting the homomorphisms $`\mathrm{\Delta }_{𝒫_c}`$, $`\chi _{i_c}`$ and $`\iota _c`$ to $`\mathrm{\Delta }_{\stackrel{˘}{𝒫}}:=(\chi id)\mathrm{\Delta }_{𝒫_c}\chi ^1`$, $`\stackrel{˘}{\chi }_i:=\chi _{i_c}\chi ^1`$ and $`\stackrel{˘}{\iota }:=\chi \iota _c`$ respectively, one obtains a locally trivial QPFB again. Explicitly,
$`\mathrm{\Delta }_{\stackrel{˘}{𝒫}}((g_i)_{iI})`$ $`=`$ $`((id\mathrm{\Delta })(g_i))_{iI},`$
$`\stackrel{˘}{\chi }_i((g_k)_{kI})`$ $`=`$ $`g_i,`$
$`\stackrel{˘}{\iota }(a)`$ $`=`$ $`(\pi _i(a)1)_{iI}.`$
We note that the isomorphisms $`\varphi _{ij}`$ fulfill
$`(id\mathrm{\Delta })\varphi _{ij}`$ $`=`$ $`(\varphi _{ij}id)(id\mathrm{\Delta }),`$ (4)
$`\varphi _{ij}(a1)`$ $`=`$ $`a1,aB_{ij}/K_j^i.`$ (5)
Using (4) and (5) it follows that the subalgebra $`\stackrel{˘}{𝒫}^{coH}=\{f\stackrel{˘}{𝒫}|\mathrm{\Delta }_{\stackrel{˘}{𝒫}}(f)=f1\}`$ is isomorphic to
$$\stackrel{˘}{𝒫}^{coH}=\{\{(a_i1)\}_{iI}\underset{iI}{}B_i1|\psi _{ji}\pi _j^i(a_i)1=\psi _{ji}\pi _i^j(a_j)1\}.$$
This algebra is by Definition 1 isomorphic to
$$BB_c=\{(a_i)_{iI}\underset{iI}{}B_i|\pi _j^i(a_i)=\pi _i^j(a_j)\}$$
(see and the appendix). It follows that the $`\psi _{ij}`$ have to be isomorphisms, i.e. $`ker\psi _{ij}=K_j^i=0`$, which means in fact $`\psi _{ij}=id`$. Thus, the $`\chi _{ij}^i:𝒫_{ij}B_{ij}H`$ are isomorphisms.
Further define $`\mathrm{\Delta }_{𝒫_{ij}}:𝒫_{ij}𝒫_{ij}H`$ by
$$\mathrm{\Delta }_{𝒫_{ij}}(f+ker\chi _i+ker\chi _j):=\mathrm{\Delta }_𝒫(f)+(ker\chi _i+ker\chi _j)H$$
and $`\iota _{ij}:B_{ij}𝒫_{ij}`$ by
$$\iota _{ij}(\pi _{ij}(a)):=\iota (a)+ker\chi _i+ker\chi _j.$$
It is easy to verify that all the conditions of Definition 2 are satisfied. $`\mathrm{}`$
Notice that due to $`K_j^i=0`$, we have $`\stackrel{~}{K}_j^i=\pi _i(J_j)`$. This means
$$\chi _i(ker\chi _j)=\pi _i(J_j)H.$$
(6)
The isomorphisms $`\chi _{ij}^i`$ satisfy
$$\chi _{ij}^i\pi _{ij_𝒫}=(\pi _j^iid)\chi _i,$$
(7)
and the $`\varphi _{ij}`$ defined above are isomorphisms $`B_{ij}HB_{ij}H`$ fulfilling (4), (5) and $`\varphi _{ij}\varphi _{ji}=id`$.
###### Proposition 3
(cf. ) Locally trivial QPFB’s over a basis $`B`$ with complete covering $`(J_i)_{iI}`$ and with structure group H are in one-to-one correspondence with families of homomorphisms
$$\tau _{ij}:HB_{ij},$$
called transition functions, satisfying the conditions
$`\tau _{ii}(h)`$ $`=`$ $`1\epsilon (h)hH,`$
$`\tau _{ji}(S(h))`$ $`=`$ $`\tau _{ij}(h)hH,`$
$`\tau _{ij}(h)a`$ $`=`$ $`a\tau _{ij}(h)aB_{ij}hH,`$
$`\pi _k^{ij}\tau _{ij}(h)`$ $`=`$ $`m_{B_{ijk}}(\pi _j^{ik}\tau _{ik}\pi _i^{jk}\tau _{kj})\mathrm{\Delta }(h)hH.`$
Proof: Let a bundle $`𝒫`$ be given and let the $`\varphi _{ij}:B_{ij}HB_{ij}H`$ be defined as above. Define homomorphisms $`\tau _{ij}:HB_{ij}`$ by
$$\tau _{ji}(h):=(id\epsilon )\varphi _{ij}(1h).$$
(8)
(There is another possible choice, $`\tau _{ij}(h):=(id\epsilon )\varphi _{ij}(1h)`$, which correspond to another form of the cocycle condition.) One shows that this is equivalent to
$$\varphi _{ij}(ah)=a\tau _{ji}(h_1)h_2:$$
(9)
Using (4), (5) and $`(\epsilon id)\mathrm{\Delta }=id`$ it follows from (8) that
$`{\displaystyle a\tau _{ji}(h_1)h_2}`$ $`=`$ $`{\displaystyle (a1)((id\epsilon )\varphi _{ij}(1h_1)h_2)},`$
$`=`$ $`(a1)(id\epsilon id)(\varphi _{ij}id)(id\mathrm{\Delta })(1h),`$
$`=`$ $`(a1)(id\epsilon id)(id\mathrm{\Delta })\varphi _{ij}(1h),`$
$`=`$ $`(a1)\varphi _{ij}(1h),`$
$`=`$ $`\varphi _{ij}(ah).`$
Conversely, if (9) is satisfyed, the choice $`a=1`$ gives (8). $`\tau _{ii}(h)=\epsilon (h)1`$ follows from $`\varphi _{ii}=id`$. Every homomorphism $`\tau _{ij}:HB`$ is convolution invertible with convolution inverse $`\tau _{ij}^1=\tau _{ij}S`$. On the other hand from $`\varphi _{ij}\varphi _{ji}=id`$ easily follows $`\tau _{ij}^1=\tau _{ji}`$:
$$\varphi _{ij}\varphi _{ji}(1h)=\varphi _{ij}(\tau _{ij}(h_1)h_2)=\tau _{ij}(h_1)\tau _{ji}(h_2)h_3=1h.$$
Therefore $`\tau _{ij}(h_1)\tau _{ji}(h_2)=\epsilon (h)1`$, i.e. $`\tau _{ji}=\tau _{ij}S`$. $`\tau _{ij}`$ has values in the center of $`B_{ij}`$:
$`a\tau _{ij}(h)\tau _{ij}(h)a`$ $`=`$ $`a(id\epsilon )\varphi _{ji}(1h)(id\epsilon )\varphi _{ji}(1h)a`$
$`=`$ $`(id\epsilon )((a1)\varphi _{ji}(1h)\varphi _{ji}(1h)(a1))`$
$`=`$ $`(id\epsilon )(\varphi _{ji}((ah)(ah))=0.`$
To prove the last relation of the proposition, define isomorphisms $`\varphi _{ij}^k:B_{ijk}HB_{ijk}H`$ by
$$\varphi _{ij}^k((ah)+\pi _{ij}(J_k)H):=\varphi _{ij}(ah)+\pi _{ij}(J_k)H$$
(using $`B_{ijk}B_{ij}/\pi _{ij}(J_k)`$). $`\varphi _{ij}^k`$ are well defined because of $`\varphi _{ij}(a1)=a1`$. Now, a lengthy but simple computation leads to
$$\varphi _{ij}^k=\varphi _{ik}^j\varphi _{kj}^i.$$
The idea of this computation is to consider the isomorphism $`\chi _i^{ijk}:𝒫/(ker\chi _i+ker\chi _j+ker\chi _k)B_{ijk}H`$ induced by $`\chi _i`$ and to prove $`\varphi _{ij}^k=\chi _i^{ijk}\chi _{j}^{ijk}{}_{}{}^{1}`$.
Combining the definition of $`\varphi _{ij}^k`$ with (9), one obtains
$$\varphi _{ij}^k(ah)=a\pi _k^{ij}\tau _{ji}(h_1)h_2.$$
Therefore,
$$\pi _k^{ij}\tau _{ji}(h)=(id\epsilon )\varphi _{ij}^k(1h).$$
Inserting here $`\varphi _{ij}^k(1h)=\varphi _{ik}^j\varphi _{kj}^i(1h)`$ one obtains
$$\pi _k^{ij}\tau _{ji}(h)=\pi _i^{jk}\tau _{jk}(h_1)\pi _j^{ik}\tau _{ki}(h_2).$$
This ends the proof of one direction of the proposition.
We will not give the details of reconstruction of the bundle from the transition functions. We only remark, that, for a given family of transition functions $`\tau _{ij}`$, we define the isomorphisms $`\varphi _{ij}`$ by formula (9), which gives rise to the gluing
$$\stackrel{˘}{𝒫}=\{(f_i)_{iI}\underset{iI}{}(B_iH)|(\pi _j^iid)(f_i)=\varphi _{ij}(\pi _i^jid)(f_j)\}.$$
(10)
One verifies that the formulas
$`\mathrm{\Delta }_{\stackrel{˘}{𝒫}}((f_i)_{iI})`$ $`=`$ $`(id\mathrm{\Delta }(f_i))_{iI},(f_i)_{iI}\stackrel{~}{𝒫},`$ (11)
$`\stackrel{˘}{\chi }_k((f_i)_{iI})`$ $`=`$ $`f_k,(f_i)_{iI}\stackrel{~}{𝒫},`$ (12)
$`\stackrel{˘}{\iota }(a)`$ $`=`$ $`(\pi _i(a)1)_{iI},aB`$ (13)
define a locally trivial QPFB $`(\stackrel{˘}{𝒫},\mathrm{\Delta }_{\stackrel{˘}{𝒫}},H,B,\stackrel{˘}{\iota },(\stackrel{˘}{\chi }_i,J_i)_{iI})`$. If the $`\tau _{ij}`$ stem from a given locally trivial QPFB $`𝒫`$, applying the isomorphism $`\chi ^1`$ defined as above (proof of proposition 2) leads to $`𝒫_c\stackrel{˘}{𝒫}`$. $`\mathrm{}`$
## 3 Adapted covariant differential structures on locally trivial QPFB
In the sequel we will use the skew tensor product of differential calculi. Let $`\mathrm{\Gamma }(A)`$ and $`\mathrm{\Gamma }(B)`$ be two differential calculi. We define the differential calculus $`\mathrm{\Gamma }(A)\widehat{}\mathrm{\Gamma }(B)`$ as the vector space $`\mathrm{\Gamma }(A)\mathrm{\Gamma }(B)`$ equipped with the product
$$(\gamma \widehat{}\rho )(\omega \widehat{}\tau )=(1)^{mn}(\gamma \omega \widehat{}\rho \tau ),\omega \mathrm{\Gamma }^n(A),\rho \mathrm{\Gamma }^m(B),\gamma \mathrm{\Gamma }(A),\tau \mathrm{\Gamma }(B)$$
(14)
and the differential
$$d(\gamma \widehat{}\rho )=(d\gamma \widehat{}\rho )+(1)^n(\gamma \widehat{}d\rho ),\gamma \mathrm{\Gamma }^n(A),\rho \mathrm{\Gamma }(B).$$
(15)
###### Proposition 4
Let $`\mathrm{\Gamma }(A)`$ and $`\mathrm{\Gamma }(B)`$ be two differential calculi and let $`J(A)\mathrm{\Omega }(A)`$ and $`J(B)\mathrm{\Omega }(B)`$ be the corresponding differential ideals respectively. Let $`id1:AAB`$ and $`1id:BAB`$ be the embedding homomorphisms. The differential ideal $`J(AB)\mathrm{\Omega }(AB)`$ corresponding to $`\mathrm{\Gamma }(A)\widehat{}\mathrm{\Gamma }(B)`$ if it is generated by the sets
$$(id1)_\mathrm{\Omega }(J(A));(1id)_\mathrm{\Omega }(J(B))$$
$$\{(a1)d(1b)(d(1b))(a1)|aA,bB\}.$$
(16)
Proof: First we define a homomorphism $`\stackrel{~}{\psi }:\mathrm{\Omega }(AB)\mathrm{\Gamma }(A)\widehat{}\mathrm{\Gamma }(B)`$ by
$`\stackrel{~}{\psi }({\displaystyle \underset{k}{}}(a_k^01)d(a_k^11))`$ $`=`$ $`{\displaystyle \underset{k}{}}a_k^0da_k^1\widehat{}1,`$
$`\stackrel{~}{\psi }({\displaystyle \underset{k}{}}(1b_k^0)d(1b_k^1))`$ $`=`$ $`{\displaystyle \underset{k}{}}1\widehat{}b_k^0db_k^1.`$
It is easy to verify that the differential ideal $`\stackrel{~}{J}(AB)`$ generated by the sets (16) satisfies $`\stackrel{~}{J}(AB)ker\stackrel{~}{\psi }`$. Let $`\stackrel{~}{\mathrm{\Gamma }}(AB)=\mathrm{\Omega }(AB)/\stackrel{~}{J}(AB)`$. Note that there are the following relations in $`\stackrel{~}{\mathrm{\Gamma }}(AB)`$.
$`(d(a1))(1b)`$ $`=`$ $`(1b)d(a1),`$
$`d(a1)d(1b)`$ $`=`$ $`d(1b)d(a1)`$
Therefore, an element $`\gamma \stackrel{~}{\mathrm{\Gamma }}(AB)`$ has the general form
$$\gamma =\underset{k}{}\underset{l}{}(a_k^01)d(a_k^11)\mathrm{}d(a_k^n1)(1b_l^0)d(1b_l^1)\mathrm{}d(1b_l^m).$$
There exist homomorphisms $`\mathrm{{\rm Y}}_A:\mathrm{\Gamma }(A)\stackrel{~}{\mathrm{\Gamma }}(AB)`$ and $`\mathrm{{\rm Y}}_B:\mathrm{\Gamma }(B)\stackrel{~}{\mathrm{\Gamma }}(AB)`$ defined by
$`\mathrm{{\rm Y}}_A(a_0da_1\mathrm{}da_n)`$ $`:=`$ $`(a_01)d(a_11)\mathrm{}d(a_n1),`$
$`\mathrm{{\rm Y}}(b_0db_1\mathrm{}db_n)`$ $`:=`$ $`(1b_0)d(1b_1)\mathrm{}d(1b_n).`$
Because of $`\stackrel{~}{J}(AB)ker\stackrel{~}{\psi }`$ the homomorphism $`\psi :\stackrel{~}{\mathrm{\Gamma }}(AB)\mathrm{\Gamma }(A)\widehat{}\mathrm{\Gamma }(B)`$ defined by
$`\psi (d(a1))`$ $`=`$ $`da\widehat{}1,`$
$`\psi (d(1b))`$ $`=`$ $`1\widehat{}db,`$
exists. Since $`\mathrm{\Gamma }(A)\widehat{}\mathrm{\Gamma }(B)`$ is isomorphic to $`\mathrm{\Gamma }(A)\mathrm{\Gamma }(B)`$ as vector space, we can define a linear map $`\psi ^1:\mathrm{\Gamma }(A)\widehat{}\mathrm{\Gamma }(B)\stackrel{~}{\mathrm{\Gamma }}(AB)`$,
$$\psi ^1(\alpha \widehat{}\beta ):=\mathrm{{\rm Y}}_A(\alpha )\mathrm{{\rm Y}}_B(\beta ).$$
Since this linear map is a homomorphism and fulfills
$`\psi ^1d`$ $`=`$ $`d\psi ^1,`$
$`\psi ^1\psi `$ $`=`$ $`id,`$
$`\psi \psi ^1`$ $`=`$ $`id,`$
$`\psi `$ is an isomorphism, and $`\stackrel{~}{J}(AB)=ker\stackrel{~}{\psi }`$. Therefore, $`\mathrm{\Gamma }(A)\widehat{}\mathrm{\Gamma }(B)\mathrm{\Omega }(AB)/\stackrel{~}{J}(AB)\mathrm{\Omega }(AB)/J(AB)`$ and the proposition follows from uniqueness of the differential ideal corresponding to a differential calculus. $`\mathrm{}`$
Remark: If we are in the converse situation, i.e. if a differential calculus $`\mathrm{\Gamma }(AB)`$ with corresponding differential ideal $`J(AB)`$ is given, there exist differential ideals $`J(A):=J(AB)\mathrm{\Omega }(A1)`$ and $`J(B):=J(AB)\mathrm{\Omega }(1B)`$. By Proposition 4, the differential calculus is isomorphic to an algebra of the form $`\mathrm{\Gamma }(A)\widehat{}\mathrm{\Gamma }(B)`$ if and only if $`J(AB)`$ is generated by the sets (16).
In the sequel we always identify $`𝒫/ker\chi _i`$ with $`B_iH`$, by means of the isomorphisms $`\stackrel{~}{\chi }_i`$ (see (2)).
Our goal is now to define differential structures on $`𝒫`$. By Proposition (19), a family of differential calculi $`\mathrm{\Gamma }(B_i)`$ and right covariant differential calculus $`\mathrm{\Gamma }(H)`$ determine unique differential calculi $`\mathrm{\Gamma }(B)`$ and $`\mathrm{\Gamma }(𝒫)`$ such that $`(\mathrm{\Gamma }(B),(\mathrm{\Gamma }(B_i))_{iI})`$ and $`(\mathrm{\Gamma }(𝒫),(\mathrm{\Gamma }(B_i)\widehat{}\mathrm{\Gamma }(H))_{iI}`$ are adapted to $`(B,(J_i)_{iI})`$ and $`(𝒫,(ker\chi _i)_{iI})`$ respectively. $`\mathrm{\Gamma }(𝒫)`$ and $`\mathrm{\Gamma }(B)`$ are given in the following way: One has the extensions $`\chi _{i_{\mathrm{\Omega }\mathrm{\Gamma }}}:\mathrm{\Omega }(𝒫)\mathrm{\Gamma }(B_i)\widehat{}\mathrm{\Gamma }(H)`$ and $`\pi _{i_{\mathrm{\Omega }\mathrm{\Gamma }}}:\mathrm{\Omega }(B)\mathrm{\Gamma }(B_i)`$ of the $`\chi _i`$ and $`\pi _i`$ respectively. These extensions form differential ideals $`ker\chi _{i_{\mathrm{\Omega }\mathrm{\Gamma }}}\mathrm{\Omega }(𝒫)`$ and $`ker\pi _{i_{\mathrm{\Omega }\mathrm{\Gamma }}}\mathrm{\Omega }(B)`$, thus $`J(𝒫):=_{iI}ker\chi _{i_{\mathrm{\Omega }\mathrm{\Gamma }}}`$ and $`J(B):=_{iI}ker\pi _{i_{\mathrm{\Omega }\mathrm{\Gamma }}}`$ are differential ideals. By construction, $`\mathrm{\Gamma }(𝒫):=\mathrm{\Omega }(𝒫)/J(𝒫)`$ and $`\mathrm{\Gamma }(B):=\mathrm{\Omega }(B)/J(B)`$ are adapted, i.e. the extensions $`\chi _{i_\mathrm{\Gamma }}:\mathrm{\Gamma }(𝒫)\mathrm{\Gamma }(B_i)\widehat{}\mathrm{\Gamma }(H)`$ and $`\pi _{i_\mathrm{\Gamma }}:\mathrm{\Gamma }(B)\mathrm{\Gamma }(B_i)`$ of the $`\chi _i`$ and $`\pi _i`$ exist and fulfill $`_{iI}ker\chi _{i_\mathrm{\Gamma }}=0`$ and $`_{iI}ker\pi _{i_\mathrm{\Gamma }}=0`$ respectively.
###### Definition 3
A differential structure on a locally trivial QPFB is a differential calculus $`\mathrm{\Gamma }(𝒫)`$ defined by a family of differential calculi $`\mathrm{\Gamma }(B_i)`$ and a right covariant differential calculus $`\mathrm{\Gamma }(H)`$, as described above.
###### Proposition 5
Let $`\mathrm{\Gamma }(𝒫)`$ be a differential structure on $`𝒫`$, and let $`\mathrm{\Gamma }(B)`$ be determined by the corresponding $`\mathrm{\Gamma }(B_i)`$ as above. $`\mathrm{\Gamma }(𝒫)`$ is covariant. The $`\chi _{i_\mathrm{\Gamma }}`$ satisfy
$$\mathrm{\Delta }_𝒫^\mathrm{\Gamma }(ker\chi _{i_\mathrm{\Gamma }})ker\chi _{i_\mathrm{\Gamma }}HiI.$$
(17)
The extension $`\iota _\mathrm{\Gamma }:\mathrm{\Gamma }(B)\mathrm{\Gamma }(𝒫)`$ of $`\iota `$ exists, fulfills
$$\chi _{i_\mathrm{\Gamma }}\iota _\mathrm{\Gamma }(\gamma )=\pi _{i_\mathrm{\Gamma }}(\gamma )\widehat{}1,\gamma \mathrm{\Gamma }(B),$$
and is injective.
Proof: As explained before definition 3, the differential ideal corresponding to $`\mathrm{\Gamma }(𝒫)`$ is $`J(𝒫)=_iker\chi _{i_{\mathrm{\Omega }\mathrm{\Gamma }}}\mathrm{\Omega }(𝒫)`$. Using the right covariance of $`\mathrm{\Gamma }(H)`$ and Definition 14 one finds that the extensions $`\chi _{i_{\mathrm{\Omega }\mathrm{\Gamma }}}`$ fulfill
$$(\chi _{i_{\mathrm{\Omega }\mathrm{\Gamma }}}id)\mathrm{\Delta }_𝒫^\mathrm{\Omega }=(id\mathrm{\Delta }^\mathrm{\Gamma })\chi _{i_{\mathrm{\Omega }\mathrm{\Gamma }}},$$
where $`\mathrm{\Delta }^\mathrm{\Gamma }`$ is the right coaction of $`\mathrm{\Gamma }(H)`$. Due to this formula the differential ideals $`ker\chi _{i_{\mathrm{\Omega }\mathrm{\Gamma }}}`$ are covariant under the coaction of $`H`$, i.e. $`\mathrm{\Delta }_𝒫^\mathrm{\Omega }(ker\chi _{i_{\mathrm{\Omega }\mathrm{\Gamma }}})ker\chi _{i_{\mathrm{\Omega }\mathrm{\Gamma }}}H`$, thus, the differential ideal $`J(𝒫):=_{iI}ker\chi _{i_{\mathrm{\Omega }\mathrm{\Gamma }}}`$ corresponding to $`\mathrm{\Gamma }(𝒫)`$ is covariant and it follows that $`\mathrm{\Gamma }(𝒫)`$ is covariant. This also gives (17).
The differential ideal corresponding to $`\mathrm{\Gamma }(B)`$ is $`J(B)=_iker\pi _{i_{\mathrm{\Omega }\mathrm{\Gamma }}}\mathrm{\Omega }(B)`$. It is easy to see that $`\iota _\mathrm{\Omega }(J(B))J(𝒫)`$, thus the extension $`\iota _\mathrm{\Gamma }`$ of $`\iota `$ with respect to $`\mathrm{\Gamma }(B)`$ and $`\mathrm{\Gamma }(𝒫)`$ exists. Clearly $`\iota _\mathrm{\Gamma }`$ satisfies
$$\chi _{i_\mathrm{\Gamma }}\iota _\mathrm{\Gamma }(\gamma )=\pi _{i_\mathrm{\Gamma }}(\gamma )\widehat{}1,\gamma \mathrm{\Gamma }(B).$$
Because of this formula and $`_iker\pi _{i_\mathrm{\Gamma }}=0`$, $`\iota _\mathrm{\Gamma }`$ is injective. $`\mathrm{}`$
The differential structure on a locally trivial QPFB determines the covering completion $`\mathrm{\Gamma }_c(𝒫)`$ of $`\mathrm{\Gamma }(𝒫)`$ with respect to the covering $`(ker\chi _{i_\mathrm{\Gamma }})_{iI}`$ (see appendix). $`\mathrm{\Gamma }_c(𝒫)`$ is an LC differential algebra (see appendix) with local differential calculi $`\mathrm{\Gamma }(B_i)\widehat{}\mathrm{\Gamma }(H)`$. It will be shown that $`\mathrm{\Gamma }_c(𝒫)`$ is a right $`H`$-comodule algebra and that the covering completion $`\mathrm{\Gamma }_c(B)`$ of $`\mathrm{\Gamma }(B)`$ is embedded in $`\mathrm{\Gamma }_c(𝒫)`$. But first we need some facts about differential calculi over $`B_{ij}H`$ appearing in our context. For the moment we can even assume that we have a general differential calculus $`\mathrm{\Gamma }(B_iH)`$. Over the algebras $`B_{ij}H`$ there exist two isomorphic differential calculi $`\mathrm{\Gamma }^i(B_{ij}H)=\mathrm{\Gamma }(B_iH)/\chi _{i_\mathrm{\Gamma }}(ker\chi _{j_\mathrm{\Gamma }})`$ and $`\mathrm{\Gamma }^j(B_{ij}H)=\mathrm{\Gamma }(B_jH)/\chi _{j_\mathrm{\Gamma }}(ker\chi _{i_\mathrm{\Gamma }})`$, and two corresponding differential ideals $`J^i(B_{ij}H)\mathrm{\Omega }(B_{ij}H)`$ and $`J^j(B_{ij}H)\mathrm{\Omega }(B_{ij}H)`$.
###### Proposition 6
The differential ideals $`J^i(B_{ij}H)`$ and $`J^j(B_{ij}H)`$ have the following form:
$$J^i(B_{ij}H)=(\pi _j^iid)_\mathrm{\Omega }(J(B_iH))+\varphi _{ij_\mathrm{\Omega }}(\pi _i^jid)_\mathrm{\Omega }(J(B_jH))$$
(18)
$$J^j(B_{ij}H)=(\pi _i^jid)_\mathrm{\Omega }(J(B_jH))+\varphi _{ji_\mathrm{\Omega }}(\pi _j^iid)_\mathrm{\Omega }(J(B_iH)),$$
(19)
where $`\varphi _{ij_\mathrm{\Omega }}`$ are the the extensions of the isomorphisms $`\varphi _{ij}`$ corresponding to the transition functions $`\tau _{ji}`$.
For the proof we need
###### Lemma 2
$$(\pi _j^iid)\chi _i=\varphi _{ij}(\pi _i^jid)\chi _j$$
Proof of the lemma: Using the identities $`\varphi _{ij}=\chi _{ij}^i\chi _{ij}^{j}{}_{}{}^{1}`$ and $`\chi _{ij}^i\pi _{ij_𝒫}=(\pi _j^iid)\chi _i`$, one has
$`(\pi _j^iid)\chi _i`$ $`=`$ $`\chi _{ij}^i\pi _{ij_𝒫}`$
$`=`$ $`\varphi _{ij}\chi _{ij}^j\pi _{ij_𝒫}`$
$`=`$ $`\varphi _{ij}(\pi _i^jid)\chi _j.`$
$`\mathrm{}`$
Proof of the proposition: The differential calculus $`\mathrm{\Gamma }^i(B_{ij}H)=\mathrm{\Gamma }(B_iH)/\chi _{i_\mathrm{\Gamma }}(ker\chi _{j_\mathrm{\Gamma }})`$ is isomorphic to $`\mathrm{\Gamma }(𝒫)/(ker\chi _{i_\mathrm{\Gamma }}+ker\chi _{j_\mathrm{\Gamma }})`$, which in turn is isomorphic to $`\mathrm{\Omega }(𝒫)/(ker\chi _{i_{\mathrm{\Omega }\mathrm{\Gamma }}}+ker\chi _{j_{\mathrm{\Omega }\mathrm{\Gamma }}})`$. Thus the differential calculi $`\mathrm{\Gamma }^i(B_{ij}H)`$ and $`\mathrm{\Gamma }^j(B_{ij}H)`$ can be identified with $`\mathrm{\Omega }(B_iH)/\chi _{i_\mathrm{\Omega }}(ker\chi _{i_{\mathrm{\Omega }\mathrm{\Gamma }}}+ker\chi _{j_{\mathrm{\Omega }\mathrm{\Gamma }}})`$ and $`\mathrm{\Omega }(B_jH)/\chi _{j_\mathrm{\Omega }}(ker\chi _{i_{\mathrm{\Omega }\mathrm{\Gamma }}}+ker\chi _{j_{\mathrm{\Omega }\mathrm{\Gamma }}})`$ respectively, and one obtains the differential ideals
$`J^i(B_{ij}H)`$ $`=`$ $`(\pi _j^iid)_\mathrm{\Omega }\chi _{i_\mathrm{\Omega }}(ker\chi _{i_{\mathrm{\Omega }\mathrm{\Gamma }}}+ker\chi _{j_{\mathrm{\Omega }\mathrm{\Gamma }}}),`$
$`J^j(B_{ij}H)`$ $`=`$ $`(\pi _i^jid)_\mathrm{\Omega }\chi _{j_\mathrm{\Omega }}(ker\chi _{i_{\mathrm{\Omega }\mathrm{\Gamma }}}+ker\chi _{j_{\mathrm{\Omega }\mathrm{\Gamma }}}).`$
Now, $`\chi _{i_\mathrm{\Omega }}(ker\chi _{i_{\mathrm{\Omega }\mathrm{\Gamma }}})=J(B_iH)`$ and $`\chi _{j_\mathrm{\Omega }}(ker\chi _{j_{\mathrm{\Omega }\mathrm{\Gamma }}})=J(B_jH)`$ yields
$$J^i(B_{ij}H)=(\pi _j^iid)_\mathrm{\Omega }(J(B_iH)+(\pi _j^iid)_\mathrm{\Omega }\chi _{i_\mathrm{\Omega }}(\chi _{j_\mathrm{\Omega }}^1(J(B_jH))).$$
(20)
Due to Lemma 2 the two homomorphisms $`(\pi _j^iid)_\mathrm{\Omega }\chi _{i_\mathrm{\Omega }}:\mathrm{\Omega }(𝒫)\mathrm{\Omega }(B_{ij}H)`$ and $`(\pi _i^jid)_\mathrm{\Omega }\chi _{j_\mathrm{\Omega }}:\mathrm{\Omega }(𝒫)\mathrm{\Omega }(B_{ij}H)`$ are connected by
$$(\pi _j^iid)_\mathrm{\Omega }\chi _{i_\mathrm{\Omega }}=\varphi _{ij_\mathrm{\Omega }}(\pi _i^jid)_\mathrm{\Omega }\chi _{j_\mathrm{\Omega }},$$
thus
$$(\pi _j^iid)_\mathrm{\Omega }\chi _{i_\mathrm{\Omega }}(\chi _{j_\mathrm{\Omega }}^1(J(B_jH)))=\varphi _{ij_\mathrm{\Omega }}(\pi _i^jid)_\mathrm{\Omega }(J(B_jH)).$$
Inserting this formula in (20) gives (18). (19) results by exchanging $`i,j`$. $`\mathrm{}`$
Due to $`J^i(B_{ij}H)=\varphi _{ij_\mathrm{\Omega }}(J^j(B_{ij}H))`$ (immediate from Proposition 6) the isomorphism $`\varphi _{ij}`$ is differentiable with respect to $`\mathrm{\Gamma }^j(B_{ij}H)`$ and $`\mathrm{\Gamma }^i(B_{ij}H)`$.
From now on we consider the case $`\mathrm{\Gamma }(B_iH)=\mathrm{\Gamma }(B_i)\widehat{}\mathrm{\Gamma }(H)`$.
Denoting by $`(\pi _j^iid)_{\mathrm{\Gamma }^i}:\mathrm{\Gamma }(B_i)\widehat{}\mathrm{\Gamma }(H)\mathrm{\Gamma }^i(B_{ij}H)`$ the natural projection, $`\mathrm{\Gamma }_c(𝒫)`$ has the following explicit form:
$$\mathrm{\Gamma }_c(𝒫)=\{(\gamma _i)_{iI}\underset{iI}{}\mathrm{\Gamma }(B_i)\widehat{}\mathrm{\Gamma }(H)|(\pi _j^iid)_{\mathrm{\Gamma }^i}(\gamma _i)=\varphi _{ij_\mathrm{\Gamma }}(\pi _i^jid)_{\mathrm{\Gamma }^j}(\gamma _j)\}.$$
(21)
Remark: Later we will need
$$ker(\pi _j^iid)_\mathrm{\Gamma }=\chi _{i_\mathrm{\Gamma }}(ker\chi _{j_\mathrm{\Gamma }})=\chi _{i_{\mathrm{\Gamma }_c}}(ker\chi _{j_{\mathrm{\Gamma }_c}})$$
(22)
###### Proposition 7
Let $`\mathrm{\Gamma }(𝒫)`$ be a differential structure on $`𝒫`$, let $`\mathrm{\Gamma }_c(𝒫)`$ be the covering completion of $`\mathrm{\Gamma }(𝒫)`$ and let $`\mathrm{\Gamma }_c(B)`$ be the covering completion of $`\mathrm{\Gamma }(B)`$. Let $`\chi _{i_{\mathrm{\Gamma }_c}}`$ and $`\pi _{i_{\mathrm{\Gamma }_c}}`$ be the restrictions of the respective i-th projections.
Then there exist a unique right coaction $`\mathrm{\Delta }_𝒫^{\mathrm{\Gamma }_c}:\mathrm{\Gamma }_c(𝒫)\mathrm{\Gamma }_c(𝒫)H`$ and a unique injective homomorphism $`\iota _{\mathrm{\Gamma }_c}:\mathrm{\Gamma }_c(B)\mathrm{\Gamma }_c(𝒫)`$ such that
$`(\chi _{i_{\mathrm{\Gamma }_c}}id)\mathrm{\Delta }_𝒫^{\mathrm{\Gamma }_c}`$ $`=`$ $`(id_i\mathrm{\Delta }^\mathrm{\Gamma })\chi _{i_{\mathrm{\Gamma }_c}}`$ (23)
$`\chi _{i_{\mathrm{\Gamma }_c}}\iota _{\mathrm{\Gamma }_c}(\gamma )`$ $`=`$ $`\pi _{i_{\mathrm{\Gamma }_c}}(\gamma )1,\gamma \mathrm{\Gamma }_c(B).`$ (24)
Remark: Indeed, the $`\chi _{i_{\mathrm{\Gamma }_c}}:\mathrm{\Gamma }_c(𝒫)\mathrm{\Gamma }(B_i)\widehat{}\mathrm{\Gamma }(H)`$ and $`\pi _{i_{\mathrm{\Gamma }_c}}:\mathrm{\Gamma }_c(B)\mathrm{\Gamma }(B_i)`$ coincide with the differential extensions of $`\chi _i`$ and $`\pi _i`$.
Proof: The covariance of the ideals $`\chi _{i_\mathrm{\Gamma }}(ker\chi _{j_\mathrm{\Gamma }})`$ under the $`H`$-coaction $`(id_i\mathrm{\Delta }^\mathrm{\Gamma })`$ follows from the covariance of the ideals $`ker\chi _{i_\mathrm{\Gamma }}`$ under the $`H`$-coaction $`\mathrm{\Delta }_𝒫^\mathrm{\Gamma }`$. Therefore there exist $`H`$-coactions $`(id\mathrm{\Delta })^{\mathrm{\Gamma }^i}`$ on $`\mathrm{\Gamma }^i(B_{ij}H)`$ satisfying
$`(id\mathrm{\Delta })^{\mathrm{\Gamma }^i}(\pi _j^iid)_{\mathrm{\Gamma }^i}`$ $`=`$ $`((\pi _j^iid)_{\mathrm{\Gamma }^i})id)(id\mathrm{\Delta }^\mathrm{\Gamma }),`$ (25)
$`(id\mathrm{\Delta })^{\mathrm{\Gamma }^i}`$ $`=`$ $`(\varphi _{ij_\mathrm{\Gamma }}id)(id\mathrm{\Delta })^{\mathrm{\Gamma }^j}.`$ (26)
Thus there exists a $`H`$-coaction $`\mathrm{\Delta }_𝒫^{\mathrm{\Gamma }_c}`$ on $`\mathrm{\Gamma }_c(𝒫)`$ defined by
$$\mathrm{\Delta }_𝒫^{\mathrm{\Gamma }_c}((\gamma _i)_{iI})=((id_i\mathrm{\Delta }^\mathrm{\Gamma })(\gamma _i))_{iI},(\gamma _i)_{iI}\mathrm{\Gamma }_c(𝒫).$$
(27)
Further one defines an injective homomorphism $`\iota _{\mathrm{\Gamma }_c}:\mathrm{\Gamma }_c(B)\mathrm{\Gamma }_c(𝒫)`$ by
$$\iota _{\mathrm{\Gamma }_c}((\rho _i)_{iI})=(\rho _i\widehat{}1)_{iI},(\rho _i)_{iI}\mathrm{\Gamma }_c(B).$$
(28)
Both homomorphisms are uniquely determined by the assumptions of the proposition. $`\mathrm{}`$
In general the differential calculi $`\mathrm{\Gamma }^i(B_{ij}H)`$ and $`\mathrm{\Gamma }^j(B_{ij}H)`$ seem not to be isomorphic to differential calculi of the form $`\mathrm{\Gamma }(B_{ij})\widehat{}\mathrm{\Gamma }(H)`$. This is suggested by a look at the generators of the differential ideal $`J^i(B_{ij}H)`$:
Let $`\iota _{i_\mathrm{\Omega }}:\mathrm{\Omega }(B_i)\mathrm{\Omega }(B_iH)`$ be the extension of $`\iota _i:=id1`$ and let $`\varphi _{i_\mathrm{\Omega }}:\mathrm{\Omega }(H)\mathrm{\Omega }(B_iH)`$ be the extension of $`\varphi _i:=1id`$. By Proposition 4 the differential ideal $`J(B_iH)`$ corresponding to $`\mathrm{\Gamma }(B_i)\widehat{}\mathrm{\Gamma }(H)`$ is generated by the sets
$$\iota _{i_\mathrm{\Omega }}(J(B_i)),\varphi _{i_\mathrm{\Omega }}(J(H))$$
$$\{(a1)d(1h)(d(1h))(a1),aB_i,hH\},$$
(29)
where the differential ideals $`J(B_i)`$ and $`J(H)`$ correspond to the differential calculi $`\mathrm{\Gamma }(B_i)`$ and $`\mathrm{\Gamma }(H)`$. Assume that the differential ideal $`J(H)`$ is determined by a right ideal $`Rker\epsilon H`$ in the sense that $`J(H)`$ is generated by the set $`\{S^1(r_2)dr_1|rR\}`$ (see also the appendix). Using (9),(29) and (18) one obtains the following generators of $`J^i(B_{ij}H)`$:
$`(\pi _j^iid)_\mathrm{\Omega }\iota _{i_\mathrm{\Omega }}(J(B_i)),(\pi _i^jid)_\mathrm{\Omega }\iota _{j_\mathrm{\Omega }}(J(B_j)),`$ (30)
$`\{{\displaystyle }(1S^1(r_2)d(1r_1)|rR\},`$ (31)
$`\{{\displaystyle }(\tau _{ij}(r_4)S^1(r_3)d(\tau _{ji}(r_1)r_2)|rR\},`$ (32)
$`\{(a1)d(1h)(d(1h))(a1)|aB_{ij},hH\},`$ (33)
$$\{(a1)d(\tau _{ji}(h_1)h_2)(d(\tau _{ji}(h_1)h_2))(a1)|aB_{ij},hH\}.$$
(34)
Observe that
$`{\displaystyle }(\tau _{ij}(r_4)S^1(r_3)d(\tau _{ji}(r_1)r_2){\displaystyle }(\tau _{ij}(r_2)1)d(\tau _{ji}(r_1)1)`$
$`{\displaystyle (\tau _{ji}(S(r_4)r_1)S^1(r_3))d(1r_2)}J^i(B_{ij}H);rR,`$
thus one can replace the generators (32) by
$$(\tau _{ij}(r_2)1)d(\tau _{ji}(r_1)1)+(\tau _{ji}(S(r_4)r_1)S^1(r_3))d(1r_2)J^i(B_{ij}H);rR.$$
(35)
Using the Leibniz rule, the fact that the image of $`\tau _{ji}`$ lies in the center of $`B_{ij}`$, and the generators (33), one can replace (34) by the set of generators
$$\{(a1)d(\tau _{ji}(h)1)d(\tau _{ji}(h)1)(a1)|aB_{ij},hH\}.$$
###### Proposition 8
Let the differential calculus $`\mathrm{\Gamma }(H)`$ be determined by a right ideal $`Rker\epsilon H`$ and let $`\tau _{ji}`$ be the transition function corresponding to the isomorphism $`\varphi _{ij}`$. Assume that the right ideal has the property
$$\tau _{ij}(S(r_1)r_3)r_2B_{ij}RrR;i,jI.$$
(36)
Then there exist differential ideals $`J_m(B_{ij})\mathrm{\Omega }(B_{ij})`$ such that
$$\mathrm{\Gamma }^i(B_{ij}H)=\mathrm{\Gamma }^j(B_{ij}H)(\mathrm{\Omega }(B_{ij})/J_m(B_{ij}))\widehat{}\mathrm{\Gamma }(H).$$
Proof: Because of (36) the second term of (35) lies already in the part of $`J^i(B_{ij}H)`$ generated by the set (31), thus $`J^i(B_{ij}H)`$ is generated by the sets
$`(\pi _j^iid)_\mathrm{\Omega }\iota _{i_\mathrm{\Omega }}(J(B_i)),(\pi _i^jid)_\mathrm{\Omega }\iota _{j_\mathrm{\Omega }}(J(B_j)),`$
$`\{{\displaystyle }(1S^1(r_2)d(1r_1)|rR\},`$
$`\{{\displaystyle (\tau _{ij}(r_2)1)d(\tau _{ji}(r_1)1)}|rR\},`$
$`\{(a1)d(1h)(d(1h))(a1)|aB_{ij},hH\},`$
$`\{(a1)d(\tau _{ji}(h)1)(d(\tau _{ji}(h)1))(a1)|aB_{ij},hH\}.`$
One can see that the differential ideal $`J^i(B_{ij}H)`$ is of the form (16), where the differential ideal $`J_m(B_{ij})`$ corresponding to $`\mathrm{\Omega }(B_{ij})/J_m(B_{ij})`$ is generated by the following sets:
$`\pi _{j_\mathrm{\Omega }}^i(J(B_i)),\pi _{i_\mathrm{\Omega }}^j(J(B_j)),`$ (37)
$`\{{\displaystyle \tau _{ji}(r_1)d\tau _{ij}(r_2)}|rR\},`$ (38)
$`\{(d\tau _{ji}(h))aad\tau _{ji}(h)|aB_{ij};hH\}.`$ (39)
Replacing $`\tau _{ji}`$ with $`\tau _{ij}`$ we get the same differential ideal $`J_m(B_{ij})`$. This is clear because of the relation $`\tau _{ji}(S(h))=\tau _{ij}(h)`$ and the following calculation. From the identity
$$\tau _{ij}(r_1)\tau _{ji}(r_2)d(\tau _{ij}(r_3)\tau _{ji}(r_4))=0;rR$$
one obtains
$$\tau _{ji}(S(r_1)r_4)\tau _{ji}(r_2)d\tau _{ij}(r_3)+\tau _{ij}(r_1)d(\tau _{ji}(r_2)J_m(B_{ij}).$$
Due to (36) the first term lies already in $`J_m(B_{ij})`$, thus $`\{\tau _{ij}(r_1)d(\tau _{ji}(r_2)|rR\}J_m(B_{ij})`$. $`\mathrm{}`$
Remark: All right ideals $`R`$ determining a bicovariant differential calculus $`\mathrm{\Gamma }(H)`$ have the property (36), because such right ideals are Ad-invariant, i.e. $`S(r_1)r_3r_2HR;rR`$.
Observe that in the case described in the previous proposition the differential ideal $`J_m(B_{ij})`$ is in general larger than the differental ideal $`J(B_{ij})`$ (see (16)), thus the differential calculi $`\mathrm{\Gamma }_m(B_{ij}):=\mathrm{\Omega }(B_{ij})/J_m(B_{ij})`$ and $`\mathrm{\Gamma }(B)/(ker\pi _{i_\mathrm{\Gamma }}+ker\pi _{j_\mathrm{\Gamma }})`$ are in general not isomorphic.
This gives rise to the differential algebra
$$\mathrm{\Gamma }_m(B):=\{(\gamma _i)_{iI}\underset{i}{}\mathrm{\Gamma }(B_i)|\pi _{j_{\mathrm{\Gamma }_m}}^i(\gamma _i\widehat{}1)=\pi _{i_{\mathrm{\Gamma }_m}}^j(\gamma _j\widehat{}1)\},$$
where the homomorphism $`\pi _{i_{\mathrm{\Gamma }_m}}^j:\mathrm{\Gamma }(B_i)\mathrm{\Gamma }_m(B_{ij})`$ are the composition of the map $`\mathrm{\Gamma }(B_{ij})\mathrm{\Gamma }_m(B_{ij})`$ induced by the embedding $`J(B_{ij})J_m(B_{ij})`$ and $`\pi _{j_\mathrm{\Gamma }}^i`$. Because of $`J(B_{ij})J_m(B_{ij})`$ the LC differential algebra $`\mathrm{\Gamma }_c(B)`$ is a subalgebra of $`\mathrm{\Gamma }_m(B)`$. Further, $`\mathrm{\Gamma }_m(B)`$ is an LC-differential algebra naturally embedded in $`\mathrm{\Gamma }_c(𝒫)`$ by $`(\gamma _i)_{iI}(\gamma _i1)_{iI}`$. If (36) is fulfilled one has the identity
$$(\pi _j^iid)_{\mathrm{\Gamma }^i}=\pi _{j_{\mathrm{\Gamma }_m}}^iid.$$
If the right ideal $`R`$ determining $`\mathrm{\Gamma }(H)`$ does not fulfill (36), one can nevertheless construct such a LC-differential algebra $`\mathrm{\Gamma }_m(B)`$ with $`\mathrm{\Gamma }_c(B)`$ as subalgebra, and this LC-differential algebra on $`B`$ will play the role of a differential structure on $`B`$ uniquely induced from the differential structure on $`𝒫`$. For an equivalent definition of this LC-differential algebra, we need the following remark about the differential calculus induced on a subalgebra:
Let $`C`$ be an algebra and let $`AC`$ be a subalgebra. From a differential calculus $`\mathrm{\Gamma }(C)`$ one obtains a differential calculus $`\mathrm{\Gamma }(A)`$ by
$$\mathrm{\Gamma }^n(A):=\{\underset{k}{}a_0^kda_1^k\mathrm{}da_n^k\mathrm{\Gamma }(C)|a_i^kA\}.$$
Let $`J(C)\mathrm{\Omega }(C)`$ be the differential ideal corresponding to the differential calculus $`\mathrm{\Gamma }(C)`$. It is easy to verify that the differential ideal $`J(A)\mathrm{\Omega }(A)`$ corresponding to $`\mathrm{\Gamma }(A)`$ is $`J(C)\mathrm{\Omega }(A)`$.
Now recall that there are differential calculi $`\mathrm{\Gamma }^i(B_{ij}H)`$ and $`\mathrm{\Gamma }^j(B_{ij}H)`$. Since $`B_{ij}1`$ is a subalgebra of $`B_{ij}H`$ we obtain differential calculi $`\mathrm{\Gamma }^i(B_{ij})`$ and $`\mathrm{\Gamma }^j(B_{ij})`$, with corresponding differential ideals $`J^i(B_{ij})`$ and $`J^j(B_{ij})`$ defined by
$`J^i(B_{ij})`$ $`=`$ $`J^i(B_{ij}H)\mathrm{\Omega }(B_{ij}1),`$
$`J^j(B_{ij})`$ $`=`$ $`J^j(B_{ij}H)\mathrm{\Omega }(B_{ij}1).`$
Since $`\varphi _{ij_\mathrm{\Omega }}(J^j(B_{ij}H))=J^i(B_{ij}H)`$ one concludes the identity $`\varphi _{ij_\mathrm{\Omega }}(J^j(B_{ij}))=J^i(B_{ij})`$, and because of $`\varphi _{ij}(a1)=a1`$ it follows that $`J^i(B_{ij})=J^j(B_{ij})`$, i.e. $`\mathrm{\Gamma }^i(B_{ij})=\mathrm{\Gamma }^j(B_{ij})=\mathrm{\Gamma }_m(B_{ij})`$. There are injective homomorphisms $`\iota _{ij_{\mathrm{\Gamma }_m}}^i:\mathrm{\Gamma }_m(B_{ij})\mathrm{\Gamma }^i(B_{ij}H)`$ given by
$$\iota _{ij_{\mathrm{\Gamma }_m}}^i(a_0da_1da_2\mathrm{}da_n)=(a_01)d(a_11)d(a_21)\mathrm{}d(a_n1).$$
(40)
One has the idenitity
$$\iota _{ij_{\mathrm{\Gamma }_m}}^i=\varphi _{ij_\mathrm{\Gamma }}\iota _{ij_{\mathrm{\Gamma }_m}}^j.$$
(41)
Let us define the projections $`\pi _{j_{\mathrm{\Gamma }_m}}^i:\mathrm{\Gamma }(B_i)\mathrm{\Gamma }_m(B_{ij})`$ and $`\pi _{i_{\mathrm{\Gamma }_m}}^j:\mathrm{\Gamma }(B_j)\mathrm{\Gamma }_m(B_{ij})`$ by
$`\iota _{ij_{\mathrm{\Gamma }_m}}^i\pi _{j_{\mathrm{\Gamma }_m}}^i(\gamma _i)`$ $`=`$ $`(\pi _j^iid)_{\mathrm{\Gamma }^i}(\gamma _i\widehat{}1),\gamma _i\mathrm{\Gamma }(B_i),`$ (42)
$`\iota _{ij_{\mathrm{\Gamma }_m}}^j\pi _{i_{\mathrm{\Gamma }_m}}^j(\gamma _j)`$ $`=`$ $`(\pi _i^jid)_{\mathrm{\Gamma }^j}(\gamma _j\widehat{}1),\gamma _j\mathrm{\Gamma }(B_j).`$ (43)
Obviously, these projections are extensions of $`\pi _j^i`$ and $`\pi _i^j`$ respectively. In terms of these projections the LC-differential algebra $`\mathrm{\Gamma }_m(B)`$ is defined as
$$\mathrm{\Gamma }_m(B):=\{(\gamma _i)_{iI}\underset{iI}{}\mathrm{\Gamma }(B_i)|\pi _{j_{\mathrm{\Gamma }_m}}^i(\gamma _i)=\pi _{i_{\mathrm{\Gamma }_m}}^j(\gamma _j)\}.$$
(44)
$`\mathrm{\Gamma }_c(B)`$ is a subalgebra of $`\mathrm{\Gamma }_m(B)`$, and there exists an injective homomorphism $`\iota _{\mathrm{\Gamma }_m}:\mathrm{\Gamma }_m(B)\mathrm{\Gamma }_c(𝒫)`$ defined by
$$\iota _{\mathrm{\Gamma }_m}((\gamma _i)_{iI})=(\gamma _i\widehat{}1)_{iI}.$$
Example:
In this example we consider a $`U(1)`$ bundle over the sphere $`S^2`$. Assume that the algebra of differentiable functions $`C^{\mathrm{}}(U(1))`$ over $`U(1)`$ is the closure in some Fréchet topology of the algebra generated by the elements $`\alpha `$ and $`\alpha ^{}`$ satisfying
$$\alpha \alpha ^{}=\alpha ^{}\alpha =1.$$
With $`\mathrm{\Delta }(\alpha )=\alpha \alpha `$, $`\epsilon (\alpha )=1`$ and $`S(\alpha )=\alpha ^{}`$, this is a Hopf algebra. Let $`U_N`$ and $`U_S`$ be the (closed) northern and the southern hemisphere respectively, $`\{U_N,U_S\}`$ is a covering of $`S^2`$. We have a complete covering $`(I_N,I_S)`$ of $`C^{\mathrm{}}(S^2)`$, $`I_NC^{\mathrm{}}(S^2)`$ and $`I_SC^{\mathrm{}}(S^2)`$ being the functions vanishing on the subsets $`U_N`$ and $`U_S`$ respectively. Elements of $`C^{\mathrm{}}(U_N)=C^{\mathrm{}}(S^2)/I_N`$ and $`C^{\mathrm{}}(U_S)=C^{\mathrm{}}(S^2)/I_S`$ can be identified with restrictions of elements of $`C^{\mathrm{}}(S^2)`$ to the subsets $`U_N`$ and $`U_S`$ respectively. Since $`U_NU_S=S^1`$, a transition function $`\tau _{NS}:C^{\mathrm{}}(U(1))C^{\mathrm{}}(S^1)`$ defines a locally trivial QPFB $`𝒫`$. We choose
$`\tau _{NS}(\alpha )(e^{i\varphi })`$ $`=`$ $`e^{i\varphi },`$
$`\tau _{NS}(\alpha ^{})(e^{i\varphi })`$ $`=`$ $`e^{i\varphi }`$
(Hopf bundle).
Now we construct a differential structure on this bundle by giving the differential calculi
$`\mathrm{\Gamma }(C^{\mathrm{}}(U_N))`$, $`\mathrm{\Gamma }(C^{\mathrm{}}(U_S))`$ and $`\mathrm{\Gamma }(C^{\mathrm{}}(U(1)))`$. $`\mathrm{\Gamma }(C^{\mathrm{}}(U_N))`$ and $`\mathrm{\Gamma }(C^{\mathrm{}}(U_S))`$ are taken to be the usual exterior differential calculi where the corresponding differential ideals are generated by all elements of the form $`adbdba`$. For the right covariant differential calculus $`\mathrm{\Gamma }(C^{\mathrm{}}(U(1)))`$ we assume a noncommutative form. We choose as the right ideal $`R`$ determining $`\mathrm{\Gamma }(C^{\mathrm{}}(U(1)))`$ the right ideal generated by the element
$$\alpha +\nu \alpha ^{}(1+\nu )1$$
where $`0<\nu 1`$. (One obtains the usual exterior differential calculus for $`\nu =1`$.)
Now we are interested in the LC-differential algebra $`\mathrm{\Gamma }_m(C^{\mathrm{}}(S^2))`$ coming from this differential structure on $`𝒫`$ for $`q<1`$.
It is easy to verify that the right ideal $`R`$ has the property (36), thus the differential ideal $`J_m(C^{\mathrm{}}(S^1))`$ is generated by the sets (37)-(39). The sets of generators (37) and (39) give the usual exterior differential calculus on $`S^1`$, but the set of generators (38) leads to $`d\varphi =qd\varphi `$, i.e. $`d\varphi =0`$ for $`q<1`$. One obtains for the LC-differential algebra $`\mathrm{\Gamma }_m(C^{\mathrm{}}(S^2))`$
$`\mathrm{\Gamma }_m^0(C^{\mathrm{}}(S^2))`$ $`=`$ $`C^{\mathrm{}}(S^2),`$
$`\mathrm{\Gamma }_m^n(C^{\mathrm{}}(S^2))`$ $`=`$ $`\mathrm{\Gamma }^n(C^{\mathrm{}}(U_N)){\displaystyle \mathrm{\Gamma }^n(C^{\mathrm{}}(U_S))};n>0.`$
The foregoing considerations suggest the following definition.
###### Definition 4
Let $`\mathrm{\Gamma }(𝒫)`$ be a differential structure on the the locally trivial QPFB $`𝒫`$. An LC-differential algebra $`\mathrm{\Gamma }_g(B)`$ over $`B`$ is called embeddable into $`\mathrm{\Gamma }_c(𝒫)`$ if the local differential calculi of $`\mathrm{\Gamma }_g(B)`$ are $`\mathrm{\Gamma }(B_i)`$ and if there exists the extension $`\iota _{\mathrm{\Gamma }_g}:\mathrm{\Gamma }_g(B)\mathrm{\Gamma }_c(𝒫)`$ of $`\iota `$ such that
$$\chi _{i_{\mathrm{\Gamma }_c}}\iota _{\mathrm{\Gamma }_g}(\gamma )=\pi _{i_{\mathrm{\Gamma }_g}}(\gamma )\widehat{}1;\gamma \mathrm{\Gamma }_g(B)$$
(45)
($`\pi _{i_{\mathrm{\Gamma }_g}}:\mathrm{\Gamma }_g(B)\mathrm{\Gamma }(B_i)`$ is the extension of $`\pi _i`$).
Remark: From $`_{iI}ker\pi _{i_{\mathrm{\Gamma }_g}}=\{0\}`$ follows immediately that $`\iota _{\mathrm{\Gamma }_g}`$ is injective.
###### Proposition 9
The LC-differential algebra $`\mathrm{\Gamma }_m(B)`$ defined above is the maximal embeddable LC-differential algebra, i.e every embeddable LC-differential algebra $`\mathrm{\Gamma }_g(B)`$ is embedded in $`\mathrm{\Gamma }_m(B)`$ as a subalgebra of the direct sum of the $`\mathrm{\Gamma }(B_i)`$ by $`\gamma (\pi _{i_{\mathrm{\Gamma }_g}}(\gamma ))_{iI}`$.
Proof: Let $`\mathrm{\Gamma }_g(B)`$ be an embeddable LC-differential algebra. It is clear from $`_{iI}ker\pi _{i_{\mathrm{\Gamma }_g}}=\{0\}`$ that the mapping is injective. To show that its image is in $`\mathrm{\Gamma }_m(B)`$ one has to prove that for $`\gamma \mathrm{\Gamma }_g(B)`$
$$\pi _{j_{\mathrm{\Gamma }_m}}^i\pi _{i_{\mathrm{\Gamma }_g}}(\gamma )=\pi _{i_{\mathrm{\Gamma }_m}}^j\pi _{j_{\mathrm{\Gamma }_g}}(\gamma )$$
(46)
(see (44)). By (45) $`\iota _{\mathrm{\Gamma }_g}`$ has the form
$$\iota _{\mathrm{\Gamma }_g}(\gamma )=(\pi _{i_{\mathrm{\Gamma }_g}}(\gamma )\widehat{}1)_{iI};\gamma \mathrm{\Gamma }_g(B).$$
By definition, the image of $`\iota _{\mathrm{\Gamma }_g}`$ lies in $`\mathrm{\Gamma }_c(𝒫)`$, i.e.
$$(\pi _j^iid)_\mathrm{\Gamma }(\pi _{i_{\mathrm{\Gamma }_g}}(\gamma )\widehat{}1)=\varphi _{ij_\mathrm{\Gamma }}(\pi _i^jid)_\mathrm{\Gamma }(\pi _{j_{\mathrm{\Gamma }_g}}(\gamma )\widehat{}1).$$
Using (41), (42) and (43) one obtains (46). $`\mathrm{}`$
## 4 Covariant derivatives and connections on locally trivial QPFB
First we define covariant derivatives, which are more general objects then connections. This is done on the covering completion of the differential structure on $`𝒫`$, which is necessary to obtain a one to one correspondence between covariant derivatives on $`𝒫`$ and certain families of covariant derivatives on the trivializations of $`𝒫`$.
###### Definition 5
Let $`\mathrm{\Gamma }(𝒫)`$ be the differential structure on $`𝒫`$ and let $`\mathrm{\Gamma }_c(𝒫)`$ be the covering completion of $`\mathrm{\Gamma }(𝒫)`$. Let $`hor\mathrm{\Gamma }_c(𝒫)\mathrm{\Gamma }_c(𝒫)`$ be the subalgebra defined by
$$hor\mathrm{\Gamma }_c(𝒫):=\{\gamma \mathrm{\Gamma }_c(𝒫)|\chi _{i_{\mathrm{\Gamma }_c}}(\gamma )\mathrm{\Gamma }(B_i)\widehat{}HiI\}.$$
(47)
A linear map $`D_{l,r}:hor\mathrm{\Gamma }_c(𝒫)hor\mathrm{\Gamma }_c(𝒫)`$ ist called left (right) covariant derivative if it satisfies
$`D_{l,r}(hor\mathrm{\Gamma }_c^n(𝒫))hor\mathrm{\Gamma }_c^{n+1}(𝒫),`$ (48)
$`D_{l,r}(1)=0,`$ (49)
$`D_l(\iota _{\mathrm{\Gamma }_c}(\gamma )\alpha )=(d(\iota _{\mathrm{\Gamma }_c}\gamma ))\alpha +(1)^n\gamma D_l(\alpha );\gamma \mathrm{\Gamma }_c^n(B);\alpha hor\mathrm{\Gamma }_c(𝒫),`$ (50)
$`D_r(\alpha \iota _{\mathrm{\Gamma }_c}(\gamma ))=D_r(\alpha )\iota _{\mathrm{\Gamma }_c}(\gamma )+(1)^n\alpha (d\iota _{\mathrm{\Gamma }_c}(\gamma ));\gamma \mathrm{\Gamma }_c(B);\alpha hor\mathrm{\Gamma }_c^n(𝒫),`$ (51)
$`(D_{l,r}id)\mathrm{\Delta }_{𝒫^{\mathrm{\Gamma }_c}}=\mathrm{\Delta }_{𝒫^{\mathrm{\Gamma }_c}}D_{l,r},`$ (52)
$`D_{l,r}(ker\chi _{i_{\mathrm{\Gamma }_c}}|_{hor\mathrm{\Gamma }_c(𝒫)})ker\chi _{i_{\mathrm{\Gamma }_c}}|_{hor\mathrm{\Gamma }_c(𝒫)};iI.`$ (53)
In this definition the lower indices $`l`$ or $`r`$ indicate the left or the right case. The appearence $`l,r`$ means that the corresponding condition is fulfilled for both the left and the right case. This convention will be used in the sequel permanently.
Remark: In the case of trivial bundles $`BH`$ with differential structure $`\mathrm{\Gamma }(B)\widehat{}\mathrm{\Gamma }(H)`$, where $`hor(\mathrm{\Gamma }(B)\widehat{}\mathrm{\Gamma }(H))=\mathrm{\Gamma }(B)\widehat{}H`$, condition (53) is trivial. Condition (50) (respectively (51) has the form :
$`D_l(\gamma \widehat{}h)`$ $`=`$ $`d\gamma \widehat{}h+(1)^n(\gamma \widehat{}1)D_l(1h);\gamma \mathrm{\Gamma }^n(B),`$
$`D_r(\gamma \widehat{}h)`$ $`=`$ $`D_r(1h)(\gamma \widehat{}1)+d\gamma \widehat{}h.`$
###### Proposition 10
Left (right) covariant derivatives are in bijective correspondence to families of linear maps $`A_{l,r_i}:H\mathrm{\Gamma }^1(B_i)`$ with the properties
$`A_{l,r_i}(1)`$ $`=`$ $`0,`$ (54)
$`\pi _{j_{\mathrm{\Gamma }_m}}^i(A_{l,r_i}(h))`$ $`=`$ $`{\displaystyle \tau _{ij}(h_1)\pi _{i_{\mathrm{\Gamma }_m}}^j(A_{l,r_j}(h_2))\tau _{ji}(h_3)}+{\displaystyle \tau _{ij}(h_1)d\tau _{ji}(h_2)}.`$ (55)
Remark: Note that (55) is a condition in $`\mathrm{\Gamma }_m(B_{ij})`$ (See the considerations at the end of the forgoing section.).
Proof: Because of (53) a given left covariant derivative on $`hor\mathrm{\Gamma }_c(𝒫)`$ determines a family of left covariant derivatives $`D_{l_i}:\mathrm{\Gamma }(B_i)\widehat{}H\mathrm{\Gamma }(B_i)\widehat{}H`$ by
$$D_{l_i}\chi _{i_{\mathrm{\Gamma }_c}}=\chi _{i_{\mathrm{\Gamma }_c}}D_l.$$
(56)
It follows the identity $`D_l((\gamma _i)_{iI})=(D_{l_i}(\gamma _i))_{iI}`$. Since $`(D_{l_i}(\gamma _i))_{iI}\mathrm{\Gamma }_c(𝒫)`$, the $`D_{l_i}`$ satisfy
$$(\pi _j^iid)_\mathrm{\Gamma }D_{l_i}(\gamma _i)=\varphi _{ij_\mathrm{\Gamma }}(\pi _i^jid)_\mathrm{\Gamma }D_{l_j}(\gamma _j),(\gamma _i)_{iI}hor\mathrm{\Gamma }_c(𝒫).$$
(57)
One obtains a family of linear maps $`A_{l_i}:H\mathrm{\Gamma }^1(B_i)`$ by
$$A_{l_i}(h):=(id\epsilon )D_{l_i}(1h).$$
Now we need:
###### Lemma 3
$$((id\epsilon )_\mathrm{\Gamma }id)\mathrm{\Delta }_𝒫^\mathrm{\Gamma }_{\mathrm{\Gamma }(B)\widehat{}H}=id$$
Proof of the lemma: An element $`\gamma \mathrm{\Gamma }(B)\widehat{}H`$ has the general form
$$\gamma =\underset{k}{}a_0^kda_1^k\widehat{}h^k.$$
We obtain
$`((id\epsilon )_\mathrm{\Gamma }id)\mathrm{\Delta }_𝒫^\mathrm{\Gamma }(\gamma )`$ $`=`$ $`((id\epsilon )id)({\displaystyle \underset{k}{}}{\displaystyle a_0^kda_1^k\widehat{}h_1^kh_2^k})`$
$`=`$ $`{\displaystyle \underset{k}{}}a_0^kda_1^k\widehat{}h^k.`$
$`\mathrm{}`$
By the foregoing lemma, (50) and (52) one computes the identity
$$D_{l_i}(\gamma \widehat{}h)=d\gamma \widehat{}h+(1)^{n+1}\gamma A_{l_i}(h_1)\widehat{}h_2;\gamma \mathrm{\Gamma }^n(B_i);hH.$$
(58)
Because of (49) the $`A_{l_i}`$ fulfill (54). To prove the property (55) we need:
###### Lemma 4
Let $`B`$ be an algebra, $`H`$ be a Hopf algebra, $`\mathrm{\Gamma }(B)`$ be a differential calculus over $`B`$ and $`\mathrm{\Gamma }(H)`$ be a right covariant differential calculus over $`H`$. Let $`D_l:\mathrm{\Gamma }(B)\widehat{}H\mathrm{\Gamma }(B)\widehat{}H`$ be a left covariant derivative on the trivial bundle $`BH`$. Let $`J\mathrm{\Gamma }(B)\widehat{}\mathrm{\Gamma }(H)`$ be a differential ideal with the property $`(id\mathrm{\Delta }^\mathrm{\Gamma })(J)JH.`$ Then one has
$$D_l(J(\mathrm{\Gamma }(B)\widehat{}H))J(\mathrm{\Gamma }(B)\widehat{}H).$$
(59)
Proof of the Lemma: By Lemma 1 there is an ideal $`\stackrel{~}{J}\mathrm{\Gamma }(B)`$ such that
$$J(\mathrm{\Gamma }(B)\widehat{}H)=\stackrel{~}{J}\widehat{}H.$$
$`\stackrel{~}{J}`$ is an differential ideal: Let $`_k\gamma _k\widehat{}h_k\stackrel{~}{J}\widehat{}HJ`$. Since $`J`$ is a differential ideal one obtains
$$\underset{k}{}d\gamma _k\widehat{}h_k+(1)^n\underset{k}{}\gamma _k\widehat{}dh_kJ;\gamma _k\mathrm{\Gamma }^n(B).$$
The second summand lies in $`\stackrel{~}{J}\widehat{}\mathrm{\Gamma }^1(H)J`$. It follows that $`_kd\gamma _k\widehat{}h_kd\stackrel{~}{J}\widehat{}HJ(\mathrm{\Gamma }(B)\widehat{}H)`$ and one obtains $`d\stackrel{~}{J}\stackrel{~}{J}`$, thus $`\stackrel{~}{J}`$ is a differential ideal.
Applying $`D_l`$ to $`_k\gamma _k\widehat{}h_k\stackrel{~}{J}\widehat{}HJ`$ leads to
$$D_l(\underset{k}{}\gamma _k\widehat{}h_k)=\underset{k}{}d\gamma _k\widehat{}h_k+(1)^{n+1}\underset{k}{}\gamma _kD_l(1h_k),\gamma _k\mathrm{\Gamma }^n(B).$$
Since the image of $`D_l`$ lies in $`\mathrm{\Gamma }(B)\widehat{}H`$, the right hand side of this formula is an element of $`\stackrel{~}{J}\widehat{}H`$. $`\mathrm{}`$
Since the $`ker(\pi _j^iid)_\mathrm{\Gamma }\mathrm{\Gamma }(B_i)\widehat{}\mathrm{\Gamma }(H)`$ are coinvariant differential ideals (see (25)), by the foregoing lemma follows $`D_{l_i}(ker(\pi _j^iid)_\mathrm{\Gamma }(\mathrm{\Gamma }(B_i)\widehat{}H))ker(\pi _j^iid)_\mathrm{\Gamma }(\mathrm{\Gamma }(B_i)\widehat{}H)`$. This allows to define linear maps $`D_{l_i}^{ij}`$ by
$$D_{l_i}^{ij}(\pi _j^iid)_\mathrm{\Gamma }=(\pi _j^iid)_\mathrm{\Gamma }D_{l_i}.$$
Applying $`(\pi _j^iid)_\mathrm{\Gamma }`$ to (58) one obtains
$$D_{l_i}^{ij}(ah)=(d(a1))(1h)(a1)(\pi _j^iid)_\mathrm{\Gamma }(A_{l_i}(h_1)1)(1h_2).$$
(60)
Let $`(\gamma _i)_{iI}hor\mathrm{\Gamma }_c(𝒫)`$, in particular
$$(\pi _j^iid)_\mathrm{\Gamma }(\gamma _i)=\varphi _{ij_\mathrm{\Gamma }}(\pi _i^jid)_\mathrm{\Gamma }(\gamma _j).$$
(61)
Since $`D_l(\gamma _i)_{iI}=(D_{l_i}(\gamma _i))_{iI}hor\mathrm{\Gamma }_c(𝒫)`$ it follows that
$$D_{l_i}^{ij}(\pi _j^iid)_\mathrm{\Gamma }(\gamma _i)=\varphi _{ij_\mathrm{\Gamma }}D_{l_j}^{ij}(\pi _i^jid)_\mathrm{\Gamma }(\gamma _j).$$
(62)
Combining (61) and (62), one obtains
$$D_{l_i}^{ij}=\varphi _{ij_\mathrm{\Gamma }}D_{l_j}^{ij}\varphi _{ji_\mathrm{\Gamma }}.$$
(63)
Taking advantage of (60), (63), (5), (41) and (9) one computes (see also (40))
$`D_{l_i}^{ij}(1h)`$ $`=`$ $`{\displaystyle \iota _{ij_{\mathrm{\Gamma }_m}}^i(\pi _{j_{\mathrm{\Gamma }_m}}^i(A_{l_i}(h_1)))(1h_2)}`$
$`=`$ $`\varphi _{ij_\mathrm{\Gamma }}D_{l_j}^{ij}\varphi _{ji_\mathrm{\Gamma }}(1h)`$
$`=`$ $`\varphi _{ij_\mathrm{\Gamma }}D_{l_j}^{ij}({\displaystyle \tau _{ij}(h_1)h_2})`$
$`=`$ $`\varphi _{ij_\mathrm{\Gamma }}({\displaystyle }\iota _{ij_{\mathrm{\Gamma }_m}}^j(d\tau _{ij}(h_1))(1h_2){\displaystyle }\iota _{ij_{\mathrm{\Gamma }_m}}^j(\tau _{ij}(h_1)\pi _{i_{\mathrm{\Gamma }_m}}^j(A_{l_j}(h_2))(1h_3))`$
$`=`$ $`{\displaystyle \iota _{ij_{\mathrm{\Gamma }_m}}^i((d\tau _{ij}(h_1))\tau _{ji}(h_2))(1h_3)}`$
$`{\displaystyle }\iota _{ij_{\mathrm{\Gamma }_m}}^i(\tau _{ij}(h_1)(\pi _{i_{\mathrm{\Gamma }_m}}^j(A_{l_j}(h_2))\tau _{ji}(h_3))(1h_4).`$
Applying the Leibniz rule to the first term of the last row and using $`\tau _{ij}(h_1)\tau _{ji}(h_2)=\epsilon (h)1`$ one obtains the identity
$`{\displaystyle \iota _{ij_{\mathrm{\Gamma }_m}}^i(\pi _{j_{\mathrm{\Gamma }_m}}^i(A_{l_i}(h_1)))(1h_2)}`$ $`=`$ $`{\displaystyle \iota _{ij_{\mathrm{\Gamma }_m}}^i(\tau _{ij}(h_1)\pi _{i_{\mathrm{\Gamma }_m}}^j(A_{l_j}(h_2))\tau _{ji}(h_3))(1h_4)}`$
$`+{\displaystyle \iota _{ij_{\mathrm{\Gamma }_m}}^j(\tau _{ij}(h_1)d(\tau _{ji}(h_2)))(1h_3)}.`$
In order to arrive at (55) we need to kill the $`1h`$-factor. This is achieved by using a projection $`P_{inv}:\mathrm{\Gamma }^i(B_{ij}H)\{\gamma \mathrm{\Gamma }^i(B_{ij}H)|(id\mathrm{\Delta })^{\mathrm{\Gamma }^i}(\gamma )=\gamma 1\}`$ on the elements of $`\mathrm{\Gamma }^i(B_{ij}H)`$ being coinvariant under the right H coaction $`(id\mathrm{\Delta })^{\mathrm{\Gamma }^i}:\mathrm{\Gamma }^i(B_{ij}H)\mathrm{\Gamma }^i(B_{ij}H)H`$ (see also (25) and (26)). $`P_{inv}`$ is defined by
$$P_{inv}(\rho )=\rho _0S(\rho _1),\rho \mathrm{\Gamma }^i(B_{ij}H).$$
(66)
Applying $`P_{inv}`$ to the identity (4) leads to
$`\iota _{ij_{\mathrm{\Gamma }_m}}^i(\pi _{j_{\mathrm{\Gamma }_m}}^i(A_{l_i}(h)))`$ $`=`$ $`{\displaystyle \iota _{ij_{\mathrm{\Gamma }_m}}^i(\tau _{ij}(h_1)\pi _{i_{\mathrm{\Gamma }_m}}^j(A_{l_j}(h_2))\tau _{ji}(h_3))}`$
$`+{\displaystyle \iota _{ij_{\mathrm{\Gamma }_m}}^i(\tau _{ij}(h_1)d\tau _{ji}(h_2))}.`$
Due to the injectivity of $`\iota _{ij_{\mathrm{\Gamma }_m}}^i`$, this is identical to
$$\pi _{j_{\mathrm{\Gamma }_m}}^i(A_{l_i}(h))=\tau _{ij}(h_1)\pi _{i_{\mathrm{\Gamma }_m}}^j(A_{l_j})(h_2))\tau _{ij}(h_3)+\tau _{ij}(h_1)d\tau _{ji}(h_2)$$
(67)
in $`\mathrm{\Gamma }_m(B_{ij})`$.
Now we prove the converse assertion. Assume there is given a family of linear maps $`A_{l_i}:H\mathrm{\Gamma }(B_i)`$ which fulfill (54) and (55). Every $`A_{l_i}`$ defines by
$$D_{l_i}(\gamma \widehat{}h)=d\gamma \widehat{}h+(1)^{n+1}\gamma A_{l_i}(h_1)\widehat{}h_2;\gamma \mathrm{\Gamma }^n(B_i);hH$$
a left covariant derivative $`D_{l_i}`$ on $`\mathrm{\Gamma }(B_i)\widehat{}H`$. The properties (48)-(50) and (52) of $`D_{l_i}`$, are easily derived from the above formula. One has to show that $`D_l((\gamma _i)_{iI}):=(D_{l_i}(\gamma _i))_{iI},(\gamma _i)_{iI}hor\mathrm{\Gamma }_c(𝒫)`$ is a covariant derivative on $`hor\mathrm{\Gamma }_c(𝒫)`$. Because of (54), $`D_l`$ fulfills (49). The conditions (50) and (52) follows from the corresponding properties of $`D_{l_i}`$. It remains to prove, that the image of $`D_l`$ lies in $`\mathrm{\Gamma }_c(𝒫)`$, because then it also lies in $`hor\mathrm{\Gamma }_c(𝒫)`$. (This is due to the fact that all the images of the $`D_{l_i}`$ obviously are in $`\mathrm{\Gamma }(B_i)\widehat{}H`$.) Then it is also obvious from the fact that the $`\chi _{i_{\mathrm{\Gamma }_c}}`$ are the projections to the $`i`$-th components that condition (53) is fulfilled. The image of $`D_l`$ lies in $`hor\mathrm{\Gamma }_c(𝒫)`$ if the family of the $`D_{l_i}`$ fulfills
$$(\pi _j^iid)_\mathrm{\Gamma }D_{l_i}(\gamma _i)=\varphi _{ij_\mathrm{\Gamma }}(\pi _i^jid)_\mathrm{\Gamma }D_{l_j}(\gamma _j),(\gamma _i)_{iI}hor\mathrm{\Gamma }_c(𝒫).$$
(68)
By Lemma 4, the covariant derivatives $`D_{l_i}`$ give rise to maps $`D_{l_i}^{ij}`$ defined by
$$D_{l_i}^{ij}(\pi _j^iid)_\mathrm{\Gamma }=(\pi _j^iid)_\mathrm{\Gamma }D_{l_i}.$$
One has
$$D_{l_i}^{ij}(1h)=(\pi _j^iid)_\mathrm{\Gamma }(A_{l_i}(h_1)1)(1h_2),$$
(69)
and we will show that (55) yields the identity
$$D_{l_i}^{ij}=\varphi _{ij_\mathrm{\Gamma }}D_{l_j}^{ij}\varphi _{ji_\mathrm{\Gamma }}:$$
One computes for $`\gamma \mathrm{\Gamma }_m^n(B_{ij})`$
$`D_{l_i}^{ij}(\iota _{ij_{\mathrm{\Gamma }_m}}^i(\gamma )(1h))`$ $`=`$ $`(d\gamma )(1h)+(1)^{n+1}\iota _{ij_{\mathrm{\Gamma }_m}}^i(\gamma )D_{l_i}^{ij}(1h)`$
$`=`$ $`\iota _{ij_{\mathrm{\Gamma }_m}}^i(d\gamma )(1h)+(1)^{n+1}{\displaystyle \iota _{ij_{\mathrm{\Gamma }_m}}^i(\gamma \pi _{j_{\mathrm{\Gamma }_m}}^i(A_{l_i}(h_1)))(1h_2)}`$
$`=`$ $`\iota _{ij_{\mathrm{\Gamma }_m}}^i(d\gamma )(1h)`$
$`+(1)^{n+1}{\displaystyle \iota _{ij_{\mathrm{\Gamma }_m}}^i(\gamma \tau _{ij}(h_1)\pi _{i_{\mathrm{\Gamma }_m}}^j(A_{l_j}(h_2))\tau _{ji}(h_3))(1h_4)}`$
$`+(1)^{n+1}{\displaystyle \iota _{ij_{\mathrm{\Gamma }_m}}^i(\gamma (d\tau _{ij}(h_1))\tau _{ji}(h_2))(1h_3)}`$
$`=`$ $`\varphi _{ij_\mathrm{\Gamma }}D_{l_j}^{ij}\varphi _{ji_\mathrm{\Gamma }}(\gamma (1h)).`$
Thus, one obtains for $`(\gamma _i)_{iI}hor\mathrm{\Gamma }_c(𝒫)`$
$`D_{l_i}^{ij}(\pi _j^iid)_\mathrm{\Gamma }(\gamma _i)`$ $`=`$ $`D_{l_i}^{ij}\varphi _{ij_\mathrm{\Gamma }}(\pi _i^jid)_\mathrm{\Gamma }(\gamma _j)`$
$`=`$ $`\varphi _{ij_\mathrm{\Gamma }}D_{l_j}^{ij}(\pi _i^jid)_\mathrm{\Gamma }(\gamma _j),`$
and (68) follows.
It is immediate from the construction (using Lemma 3) that the correspondence is bijective.
The proof for right covariant derivatives is analogous. In this case one uses
$$D_{r_i}(\gamma \widehat{}h)=d\gamma \widehat{}h+(1)^{n+1}A_{r_i}(h_1)\gamma \widehat{}h_2$$
(70)
for $`\gamma \mathrm{\Gamma }^n(B_i)`$ $`\mathrm{}`$
Remark: Obviously, a family of linear maps $`A_i:H\mathrm{\Gamma }^1(B_i)`$ fulfilling (54) and (55) determines at the same time a left and a right covariant derivative. Consequently, there is also a bijective correspondence between left and right covariant derivatives.
###### Proposition 11
Let $`D_{l,r}:hor\mathrm{\Gamma }_c(𝒫)hor\mathrm{\Gamma }_c(𝒫)`$ be a left (right) covariant derivative and let $`\mathrm{\Gamma }_g(B)`$ be embeddable into $`\mathrm{\Gamma }_c(𝒫)`$. $`D_{l,r}`$ fulfills
$`D_l(\iota _{\mathrm{\Gamma }_g}(\gamma )\alpha )=(d(\iota _{\mathrm{\Gamma }_g}(\gamma ))\alpha +(1)^n\iota _{\mathrm{\Gamma }_g}(\gamma )D_l(\alpha );\gamma \mathrm{\Gamma }_g^n(B);\alpha hor\mathrm{\Gamma }_c(𝒫),`$ (71)
$`D_r(\alpha \iota _{\mathrm{\Gamma }_g}(\gamma ))=D_r(\alpha )\iota _{\mathrm{\Gamma }_g}(\gamma )+(1)^n\alpha (d\iota _{\mathrm{\Gamma }_g}(\gamma ));\gamma \mathrm{\Gamma }_g(B);\alpha hor\mathrm{\Gamma }_c^n(𝒫).`$ (72)
Proof: Let $`(\gamma _i)_{iI}hor\mathrm{\Gamma }_c(𝒫)`$ and $`\rho \mathrm{\Gamma }_g^n(B)`$. One has $`\iota _{\mathrm{\Gamma }_g}(\rho )=(\pi _{i_{\mathrm{\Gamma }_g}}(\rho )\widehat{}1)_{iI}`$ and $`D_l((\gamma _i)_{iI})=(D_{l_i}(\gamma _i))_{iI}`$. One calculates
$`D_l(\iota _{\mathrm{\Gamma }_g}(\rho )(\gamma _i)_{iI})`$ $`=`$ $`D_l(((\pi _{i_{\mathrm{\Gamma }_g}}(\rho )\widehat{}1)\gamma _i)_{iI})`$
$`=`$ $`(D_{l_i}(\pi _{i_{\mathrm{\Gamma }_g}}(\rho )\widehat{}1)\gamma _i))_{iI}`$
$`=`$ $`((d(\pi _{i_{\mathrm{\Gamma }_g}}(\rho )\widehat{}1))\gamma _i))_{iI}+(1)^n((\pi _{i_{\mathrm{\Gamma }_g}}(\rho )\widehat{}1)D_{l_i}(\gamma _i))_{iI}`$
$`=`$ $`(d(\iota _{\mathrm{\Gamma }_g}(\rho )))(\gamma _i)_{iI}+(1)^n\iota _{\mathrm{\Gamma }_g}(\rho )D_l((\gamma _i)_{iI}).`$
The proof for right covariant derivatives is anlog. $`\mathrm{}`$
Now we are going to define connections on locally trivial QPFB. It turns out that connections are special cases of covariant derivatives. We start with a definition dualizing the classical case in a certain sense.
###### Definition 6
Let $`\mathrm{\Gamma }(𝒫)`$ be a differential structure on $`𝒫`$ and let $`\mathrm{\Gamma }_c(𝒫)`$ be the covering completion of $`\mathrm{\Gamma }(𝒫)`$. A left (right) connection is a surjective left (right) $`𝒫`$-module homomorphism $`hor_{l,r}:\mathrm{\Gamma }_c^1𝒫hor\mathrm{\Gamma }_c^1(𝒫)`$ such that:
$$hor_{l,r}^{}{}_{}{}^{2}=hor_{l,r},$$
(73)
$$(hor_{l,r}id)\mathrm{\Delta }_𝒫^{\mathrm{\Gamma }_c}=\mathrm{\Delta }_𝒫^{\mathrm{\Gamma }_c}hor_{l,r}$$
(74)
and
$$hor_{l,r}(ker\chi _{i_{\mathrm{\Gamma }_c}})ker\chi _{i_{\mathrm{\Gamma }_c}},iI.$$
(75)
Remark: Conditions (75) in this definition are needed to have the one-to-one correspondence between connections on $`𝒫`$ and certain families of connections on the trivial bundles $`B_iH`$. On a trivial bundle $`BH`$ condition (75) is obsolete.
Remark: For a given left (right) connection there is a vertical left (right) $`𝒫`$-submodule $`ver_{l,r}\mathrm{\Gamma }_c^1(𝒫)`$ such that
$$\mathrm{\Gamma }_c^1(𝒫)=ver_{l,r}\mathrm{\Gamma }_c^1(𝒫)hor\mathrm{\Gamma }_c^1(𝒫),$$
where the projection $`ver_{l,r}:\mathrm{\Gamma }_c^1(𝒫)ver_{l,r}\mathrm{\Gamma }_c^1(𝒫)`$ is defined by $`ver_{l,r}:=idhor_{l,r}`$.
On a trivial bundle $`BH`$ with differential structure $`\mathrm{\Gamma }(B)\widehat{}\mathrm{\Gamma }(H)`$ exists always the canonical connection $`hor_c`$, which is at the same time left and right. The existence of $`hor_c`$ comes from the decomposition
$$(\mathrm{\Gamma }(B)\widehat{}\mathrm{\Gamma }(H))^1=(\mathrm{\Gamma }^1(B)\widehat{}H)(B\widehat{}\mathrm{\Gamma }^1(H))$$
(direct sum of $`(BH)`$-bimodules), which allows to define
$`hor_c(\gamma \widehat{}h)`$ $`=`$ $`\gamma \widehat{}h;\gamma \mathrm{\Gamma }^1(B),hH,`$
$`hor_c(a\widehat{}\theta )`$ $`=`$ $`0;aB,\theta \mathrm{\Gamma }^1(H).`$
###### Lemma 5
For a given connection $`hor_{l,r}`$ on $`𝒫`$ there exists a family of connections $`hor_{l,r_i}`$ on the trivilizations $`B_iH`$ such that
$$\chi _{i_{\mathrm{\Gamma }_c}}hor_{l,r}=hor_{l,r_i}\chi _{i_{\mathrm{\Gamma }_c}}.$$
(76)
Proof: The existence of linear map $`hor_{l_i}`$ satisfying (76) follows from (75). The $`hor_{l,r_i}`$ are connections on $`B_iH`$: Because of the surjectivity of the $`\chi _{i_{\mathrm{\Gamma }_c}}`$ the $`hor_{l_i}`$ map onto $`\mathrm{\Gamma }^1(B_i)\widehat{}H`$. To prove condition (73) one computes
$`hor_{l,r_i}^{}{}_{}{}^{2}\chi _{i_{\mathrm{\Gamma }_c}}`$ $`=`$ $`hor_{l,r_i}\chi _{i_{\mathrm{\Gamma }_c}}hor_{l,r}`$
$`=`$ $`\chi _{i_{\mathrm{\Gamma }_c}}hor_{l,r}^{}{}_{}{}^{2}`$
$`=`$ $`\chi _{i_{\mathrm{\Gamma }_c}}hor_{l,r}`$
$`=`$ $`hor_{l,r_i}\chi _{i_{\mathrm{\Gamma }_c}}`$
The condition (74) is fulfilled because of (23). $`\mathrm{}`$
By Definition 6 and the foregoing lemma a connection $`hor_{l,r}`$ has the following form
$$hor_{l,r}((\gamma _i)_{iI})=(hor_{l,r_i}(\gamma _i))_{iI},(\gamma _i)_{iI}hor\mathrm{\Gamma }_c(𝒫),$$
(77)
which also means that the family of linear maps $`hor_{i_{l_r}}`$ satisfies
$$(\pi _j^iid)_\mathrm{\Gamma }hor_{l,r_i}(\gamma _i)=\varphi _{ij_\mathrm{\Gamma }}(\pi _i^jid)_\mathrm{\Gamma }hor_{l,r_j}(\gamma _j)$$
(78)
for $`(\gamma _i)_{iI}\mathrm{\Gamma }_c^1(𝒫)`$.
###### Proposition 12
Let $`RH`$ be the right ideal corresponding to the right covariant differential calculus $`\mathrm{\Gamma }(H)`$. Left (right) connections on a locally trivial QPFB $`𝒫`$ are in one-to-one correspondence to left (right) covariant derivatives with the following property: The corresponding linear maps $`A_{l,r_i}`$ fulfill
$`RkerA_{l_i},iI`$ (79)
$`S^1(R)kerA_{r_i},iI.`$ (80)
Remark: Thus left (right)connections are in one to one correspondence to linear maps $`A_{l,r_i}`$ fulfilling (54), (55) and (79) (respectively (80)).
Proof: We prove the assertion only for left connections. The proof is fully analogous for right connections.
A left connection $`hor_l`$ determines a family of linear maps $`A_{l_i}:H\mathrm{\Gamma }^1(B_i)`$ by
$$A_{l_i}(h):=(id\epsilon )hor_{l_i}(1\widehat{}dh).$$
From (74) and Lemma 3 one has te identity
$$hor_{l_i}(1dh)=A_{l_i}(h_1)\widehat{}h_2.$$
(81)
Therefore $`A_{l_i}`$ have the property $`RkerA_{l_i}`$:
$$0=hor_{l_i}(1\widehat{}S^1(r_2)dr_1)=A_{l_i}(r)\widehat{}1,rR.$$
It remains to show that this family of linear maps fulfills (54) and (55).
(54) is fulfilled by definition ( $`A_{l_i}(1):=(id\epsilon )hor_{l_i}(1d1)=0`$).
Because of (75), (76) and (22) one has $`hor_{l_i}(ker(\pi _j^iid)_\mathrm{\Gamma })ker(\pi _j^iid)_\mathrm{\Gamma }`$, and the linear maps $`hor_{l_i}^{ij}`$ defined by
$$hor_{l_i}^{ij}(\pi _j^iid)_\mathrm{\Gamma }=(\pi _j^iid)_\mathrm{\Gamma }hor_{l_i}$$
exist. It follows that
$$hor_{l_i}^{ij}(d(1h))=\pi _{j_{\mathrm{\Gamma }_m}}^i(A_{l_i}(h_1))(1h_2).$$
On the other hand, by an analogoue of the computation leading to (63) (using (78)), one obtains
$$hor_{l_i}^{ij}=\varphi _{ij_\mathrm{\Gamma }}hor_{l_j}^{ij}\varphi _{ji_\mathrm{\Gamma }}.$$
Now using the last two formulas, on can repeat the arguments written after formula (63) to obtain formula (55).
Now assume that there is given a left covariant derivative $`D_l`$, whose corresponding linear maps $`A_{l_i}`$ satisfy $`RkerA_{l_i}`$. There exist left connections $`hor_{l_i}:(\mathrm{\Gamma }^1(B_i)\widehat{}H)(B_i\widehat{}\mathrm{\Gamma }^1(H))\mathrm{\Gamma }^1(B_i)\widehat{}H`$ defined by
$`hor_{l_i}(\gamma \widehat{}h)`$ $`:=`$ $`\gamma \widehat{}h,`$
$`hor_{l_i}(a\widehat{}hdk)`$ $`:=`$ $`{\displaystyle aA_{l_i}(k_1)\widehat{}hk_2}.`$ (82)
To verify this assertion we define linear maps $`hor_{l_i}^\mathrm{\Omega }:(\mathrm{\Gamma }(B_i)\widehat{}\mathrm{\Omega }(H))^1\mathrm{\Gamma }^1(B_i)\widehat{}H`$ by
$`hor_{l_i}^\mathrm{\Omega }(a_0da_1\widehat{}h)`$ $`=`$ $`a_0da_1\widehat{}h,`$
$`hor_{l_i}^\mathrm{\Omega }(a\widehat{}h^0dk)`$ $`=`$ $`{\displaystyle aA_l(k_1)\widehat{}hk_2}.`$
The $`(B_iH)`$ subbimodules $`B_i\widehat{}J^1(H)`$ are generated by the sets $`\{1\widehat{}S^1(r_2)dr_1|rR\}`$. One has
$$B_i\widehat{}\mathrm{\Gamma }^1(H)=(B_i\widehat{}\mathrm{\Omega }^1(H))/(B_i\widehat{}J^1(H))=B_i\widehat{}\mathrm{\Omega }^1(H)/J^1(H).$$
Using $`RkerA_{l_i}`$ it is easy to verify that the linear maps $`hor_{l_i}^\mathrm{\Omega }`$ sends $`B_i\widehat{}J^1(H)`$ to zero, i.e. there exist corresponding the linear maps $`hor_{l_i}`$ on $`(\mathrm{\Gamma }(B_i)\widehat{}\mathrm{\Gamma }(H))^1`$. As a consequence of there definition these linear maps are connections. One easily verifies the identity
$$hor_{l_i}d=D_{l_i}|_{B_iH},$$
(83)
where the $`D_{l_i}`$ the local left covariant derivatives defined by (56).
Now we define a linear map $`hor_l:\mathrm{\Gamma }_c^1(𝒫)_{iI}\mathrm{\Gamma }(B_i)\widehat{}H`$ by
$$hor_l((\gamma _i)_{iI}):=(hor_{l_i}(\gamma _i))_{iI},(\gamma _i)_{iI}\mathrm{\Gamma }_c^1(𝒫.$$
It remains to prove that the image of $`hor_l`$ lies in $`\mathrm{\Gamma }_c^1(𝒫)`$. Then it follows immediately from the properties of the local connections $`hor_{l_i}`$ that $`hor_l`$ is a connection.
To prove $`hor_l(\mathrm{\Gamma }_c^1(𝒫)\mathrm{\Gamma }_c^1(𝒫)`$ we need a lemma.
###### Lemma 6
$`hor_{l_i}((\chi _{i_\mathrm{\Gamma }}(ker\chi _{j_\mathrm{\Gamma }}))^1)(\chi _{i_\mathrm{\Gamma }}(ker\chi _{j_\mathrm{\Gamma }}))^1`$
Proof of the Lemma: Using the form of the generators of $`J^i(B_{ij}H)`$ (30) -(34) one finds easily that the differential calculus $`\mathrm{\Gamma }^i(B_{ij}H)`$ has the form $`\mathrm{\Gamma }^i(B_{ij}H)=(\mathrm{\Gamma }(B_{ij})\widehat{}\mathrm{\Gamma }(H))/J`$ where the differential ideal $`J`$ is generated by
$`\{{\displaystyle \tau _{ij}(r_2)d\tau _{ji}(r_1)\widehat{}1}+{\displaystyle \tau _{ij}(r_4)\tau _{ji}(r_1)\widehat{}S^1(r_3)dr_2}|rR\},`$ (84)
$`\{{\displaystyle (ad\tau _{ji}(h)(d\tau _{ji}(h))a)\widehat{}1}|hH,aB_i\}.`$ (85)
The identity $`J=(\pi _{j_\mathrm{\Gamma }}^iid)(\chi _{i_\mathrm{\Gamma }}(ker\chi _{j_\mathrm{\Gamma }}))`$ is evident.
The factorization map $`id_{ij_\mathrm{\Gamma }}^i:\mathrm{\Gamma }(B_{ij})\widehat{}\mathrm{\Gamma }(H))\mathrm{\Gamma }^i(B_{ij}H)`$ fulfills
$$id_{ij_\mathrm{\Gamma }}^i(\pi _{j_\mathrm{\Gamma }}^iid)=(\pi _j^iid)_\mathrm{\Gamma }.$$
(86)
Since $`hor_{l_i}`$ is a left modul homomorphism and $`ker(\pi _{j_\mathrm{\Gamma }}^iid)=ker\pi _{j_\mathrm{\Gamma }}^i\widehat{}\mathrm{\Gamma }(H)`$ one has
$$hor_{l_i}((ker(\pi _{j_\mathrm{\Gamma }}^iid))^1)(ker(\pi _{j_\mathrm{\Gamma }}^iid))^1,$$
(87)
thus $`hor_{l_i}`$ defines a connection $`hor_{l_i}^{ij}:(\mathrm{\Gamma }(B_{ij})\widehat{}\mathrm{\Gamma }(H))^1\mathrm{\Gamma }^1(B_{ij})\widehat{}H`$ by
$$hor_{l_i}^{ij}(\pi _{j_\mathrm{\Gamma }}^iid)=(\pi _{j_\mathrm{\Gamma }}^iid)hor_{l_i}.$$
(88)
Because of (22), (87) and
$$(\pi _j^iid)_\mathrm{\Gamma }hor_{l_i}(\chi _{i_\mathrm{\Gamma }}(ker\chi _{j_\mathrm{\Gamma }}))=id_{ij_\mathrm{\Gamma }}^ihor_{l_i}^{ij}(J)$$
(which is immediate from (86) and (88)), to prove the assertion of the lemma we have to show that $`id_{ij_\mathrm{\Gamma }}^ihor_{l_i}^{ij}(J)=0.`$ Note that the part of $`J`$ generated by (85) lies in the horizontal submodule $`\mathrm{\Gamma }^1(B_{ij})\widehat{}H`$ and is therefore invariant under $`hor_{l_i}^{ij}`$. Now let us consider the part of $`J`$ generated by (84). Since $`hor_{l_i}^{ij}`$ is a left module homomorphism, it is sufficient to consider the the product of the generators (84) with a general element $`(ah)B_{ij}H`$ on the right. Using $`Rker\epsilon `$, such an element can be written
$`{\displaystyle }\tau _{ij}(r_2)d\tau _{ji}(r_1)a\widehat{}h+{\displaystyle }\tau _{ij}(r_4)\tau _{ji}(r_1)\widehat{}S^1(r_3)dr_2)(ah)`$
$`=`$ $`{\displaystyle \tau _{ij}(r_2)d\tau _{ji}(r_1)a\widehat{}h}+{\displaystyle \tau _{ij}(r_4)\tau _{ji}(r_1)a\widehat{}S^1(r_3)d(r_2h)},rR,hH,aB_{ij}.`$
Using (82), $`Rker\epsilon `$, (55), $`RkerA_{l_j}`$ and (40) one calculates
$`id_{ij_\mathrm{\Gamma }}^ihor_{l_i}^{ij}({\displaystyle }\tau _{ij}(r_2)(d\tau _{ji}(r_1))a\widehat{}h`$
$`+{\displaystyle }\tau _{ij}(r_4)\tau _{ji}(r_1)a\widehat{}S^1(r_3)d(r_2h))`$
$`=`$ $`{\displaystyle \iota _{ij_{\mathrm{\Gamma }_m}}^i(a\tau _{ij}(r_2)d\tau _{ji}(r_1))(1h)}`$
$`{\displaystyle \iota _{ij_{\mathrm{\Gamma }_m}}^i(a\tau _{ij}(r_3)\tau _{ji}(r_1)\pi _{j_{\mathrm{\Gamma }_m}}^i(A_{l_i}(r_2h_1)))(1h_2)}`$
$`=`$ $`{\displaystyle \iota _{ij_{\mathrm{\Gamma }_m}}^i(a\tau _{ij}(r_2)d\tau _{ji}(r_1))(1h)}`$
$`{\displaystyle \iota _{ij_{\mathrm{\Gamma }_m}}^i(a\tau _{ij}(r_5)\tau _{ji}(r_1)\tau _{ij}(r_2h_1)\tau _{ji}(r_4h_3)\pi _{i_{\mathrm{\Gamma }_m}}^j(A_{l_j}(r_3h_2)))(1h_4)}`$
$`{\displaystyle \iota _{ij_{\mathrm{\Gamma }_m}}^i(a\tau _{ij}(r_4)\tau _{ji}(r_1)\tau _{ij}(r_2h_1)d\tau _{ji}(r_3h_2))(1h_3)}`$
$`=`$ $`{\displaystyle \iota _{ij_{\mathrm{\Gamma }_m}}^i(a\tau _{ij}(r_2)d\tau _{ji}(r_1))(1h)}`$
$`{\displaystyle \iota _{ij_{\mathrm{\Gamma }_m}}^i(a\tau _{ij}(h_1)\tau _{ji}(h_3)\pi _{i_{\mathrm{\Gamma }_m}}^j(A_{l_j}(rh_2)))(1h_4)}`$
$`{\displaystyle \iota _{ij_{\mathrm{\Gamma }_m}}^i(a\tau _{ij}(r_2)d\tau _{ji}(r_1))(1h)}=0.`$
The last identity comes from the fact that $`R`$ is a right ideal. $`\mathrm{}`$
Let $`(\gamma _i)_{iI}\mathrm{\Gamma }_c^1(𝒫)`$. We have to prove that
$$(\pi _j^iid)_\mathrm{\Gamma }hor_{l_i}(\gamma _i)=\varphi _{ij_\mathrm{\Gamma }}(\pi _i^jid)_\mathrm{\Gamma }hor_{l_j}(\gamma _j).$$
$`\gamma _i`$ has the general form
$$\gamma _i=\underset{k}{}\chi _i(f_k^0)d\chi _i(f_k^1);f_k^0,f_k^1𝒫.$$
Using the compability condition of (21) and (22) one verifies that $`\gamma _j`$ has the form
$$\gamma _j=\underset{k}{}\chi _j(f_k^0)d\chi _j(f_k^1)+\rho ,\rho \chi _{j_\mathrm{\Gamma }}(ker\chi _{i_\mathrm{\Gamma }}).$$
Now one obtains from Lemma 6, (57) and (83)
$`(\pi _j^iid)_\mathrm{\Gamma }hor_{l_i}(\gamma _i)`$ $`=`$ $`(\pi _j^iid)_\mathrm{\Gamma }hor_{l_i}({\displaystyle \underset{k}{}}\chi _i(f_k^0)d\chi _i(f_k^1)))`$
$`=`$ $`{\displaystyle \underset{k}{}}(\pi _j^iid)_\mathrm{\Gamma }(\chi _i(f_k^0)D_{l_i}(\chi _i(f_k^1)))`$
$`=`$ $`\varphi _{ij_\mathrm{\Gamma }}(\pi _i^jid)_\mathrm{\Gamma }({\displaystyle \underset{k}{}}\chi _j(f_k^0)D_{l_j}(\chi _j(f_k^1))`$
$`=`$ $`{\displaystyle \underset{k}{}}\varphi _{ij_\mathrm{\Gamma }}(\pi _i^jid)_\mathrm{\Gamma }hor_{l_j}(\chi _j(f_k^0)d\chi _j(f_k^1)+\rho )`$
$`=`$ $`{\displaystyle \underset{k}{}}\varphi _{ij_\mathrm{\Gamma }}(\pi _i^jid)_\mathrm{\Gamma }hor_{l_j}(\gamma _j),`$
and the assertion is proved. $`\mathrm{}`$
###### Proposition 13
There exists a bijection between left and right connections.
Proof: A left connection corresponds to a family of linear maps $`(A_{l_i})_{iI}`$ satisfying (54), (55) and (79). The linear maps $`A_{r_i}:=A_{l_i}S`$ fulfill (54) and (80), thus the $`A_{r_i}`$ define right connections on the trivilizations. One has to prove that the family $`(A_{r_i})_{iI}`$ satisfies (55). Using $`\tau _{ij}S=\tau _{ji}`$ and $`d(\tau _{ij}(h_1)\tau _{ji}(h_2))=0`$, one calculates
$`\pi _{j_{\mathrm{\Gamma }_m}}^i(A_{r_i}(h)`$ $`=`$ $`\pi _{j_{\mathrm{\Gamma }_m}}^i(A_{l_i}(S(h)))`$
$`=`$ $`{\displaystyle \tau _{ij}(S(h_3))\pi _{i_{\mathrm{\Gamma }_m}}^j(A_{l_j}(S(h_2)))\tau _{ji}(S(h_1))}{\displaystyle \tau _{ij}(S(h_2))d\tau _{ji}(S(h_1))}`$
$`=`$ $`{\displaystyle \tau _{ji}(h_3)\pi _{i_{\mathrm{\Gamma }_m}}^j(A_{l_j}(S(h_2)))\tau _{ij}(h_1)}{\displaystyle \tau _{ji}(h_2)d\tau _{ij}(h_1)}`$
$`=`$ $`{\displaystyle \tau _{ij}(h_1)\pi _{i_{\mathrm{\Gamma }_m}}^j(A_{l_j}(S(h_2)))\tau _{ji}(h_3)}+{\displaystyle \tau _{ij}(h_1)d\tau _{ji}(h_2)}`$
$`=`$ $`{\displaystyle \tau _{ij}(h_1)\pi _{i_{\mathrm{\Gamma }_m}}^j(A_{r_j}(h_2))\tau _{ji}(h_3)}+{\displaystyle \tau _{ij}(h_1)d\tau _{ji}(h_2)}.`$
$`\mathrm{}`$
Remark: A left (right) connection $`hor_{l,r}`$ and the corresponding left right covariant derivatives $`D_{l,r}`$ are connected by $`hor_{l,r}d=D_{l,r}|_𝒫`$. Note that $`hor_{l,r}`$ can be extended to the submodule
$$\mathrm{\Pi }(𝒫):=\{\gamma \mathrm{\Gamma }_c(𝒫)|\chi _{i_{\mathrm{\Gamma }_c}}(\gamma )(\mathrm{\Gamma }(B_i)\widehat{}H)(\mathrm{\Gamma }(B_i)\widehat{}\mathrm{\Gamma }^1(H))\}.$$
This means that the equation is valid on
$$hor_{l,r}d=D_{l,r}$$
(equation on $`hor\mathrm{\Gamma }_c(𝒫)`$).
To discuss curvatures of covariant derivatives and connections we introduce the notion of left (right) pre-connection forms.
###### Definition 7
A left (right) pre-connection form $`\omega _{l,r}`$ ist a linear map $`\omega _{l,r}:H\mathrm{\Gamma }_c^1(𝒫)`$ satisfying
$`\omega _{l,r}(1)`$ $`=`$ $`0,`$ (89)
$`\mathrm{\Delta }_𝒫^\mathrm{\Gamma }(\omega _l(h))`$ $`=`$ $`{\displaystyle \omega _l(h_2)S(h_1)h_3},`$ (90)
$`\mathrm{\Delta }_𝒫^\mathrm{\Gamma }(\omega _r(h))`$ $`=`$ $`{\displaystyle \omega _r(h_2)h_3S^1(h_1)},`$ (91)
$`(1hor_c)\chi _{i_{\mathrm{\Gamma }_c}}(\omega _l(h))`$ $`=`$ $`{\displaystyle 1\widehat{}S(h_1)dh_2},iI,`$ (92)
$`(1hor_c)\chi _{i_{\mathrm{\Gamma }_c}}(\omega _r(h))`$ $`=`$ $`{\displaystyle 1\widehat{}(dh_2)S^1(h_1)i}I.`$ (93)
###### Proposition 14
Left (right) covariant derivatives are in bijective correspondence to left (right) pre-connection forms.
Proof: Let $`\omega _l`$ be a left pre-connection form. $`\omega _l`$ determines a family of linear maps $`A_{l_i}`$ by
$$A_{l_i}(h):=(id\epsilon )hor_c\chi _{i_{\mathrm{\Gamma }_c}}(\omega _l(h)).$$
(94)
Because of (89) the $`A_{l_i}`$ fulfill (54).
Using
$$(1hor_c)\chi _{i_{\mathrm{\Gamma }_c}}(\omega _l(h))+hor_c\chi _{i_{\mathrm{\Gamma }_c}}(\omega _l(h))=\chi _{i_{\mathrm{\Gamma }_c}}(\omega _l(h)),$$
(92), (90), Lemma 3 and (94) one verifies easily
$$\chi _{i_{\mathrm{\Gamma }_c}}(\omega _l(h))=1\widehat{}S(h_1)dh_2A_{l_i}(h_2)\widehat{}S(h_1)h_3.$$
(95)
Since
$$(\pi _j^iid)\chi _{i_{\mathrm{\Gamma }_c}}(\omega _l(h))=\varphi _{ij_\mathrm{\Gamma }}(\pi _i^jid)\chi _{j_{\mathrm{\Gamma }_c}}(\omega _l(h)),$$
an easy calculation (using (9) and the projection $`P_{inv}`$ (66)) leads to (55).
We want to prove that the left covariant derivative $`D_l`$ determined by the $`A_{l_i}`$ is
$$D_l(\gamma )=d\gamma +(1)^n\gamma _0\omega _l(\gamma _1),\gamma hor\mathrm{\Gamma }_c^n(𝒫).$$
(96)
It is sufficiant to prove that for $`\gamma hor\mathrm{\Gamma }_c^n(𝒫)`$
$$\chi _{i_{\mathrm{\Gamma }_c}}(d\gamma +(1)^n\gamma _0\omega _l(\gamma _1))=d\gamma \widehat{}h+(1)^{n+1}\gamma A_{l_i}(h_1)\widehat{}h_2=\chi _{i_{\mathrm{\Gamma }_c}}D_l(\gamma ).$$
$`\chi _{i_{\mathrm{\Gamma }_c}}(\gamma )`$ has the general form
$$\chi _{i_{\mathrm{\Gamma }_c}}(\gamma )=\underset{k}{}\gamma _i^k\widehat{}h_i^k;\gamma _i^k\mathrm{\Gamma }^n(B_i),h_i^kH.$$
Using (95) one obtains
$`\chi _{i_{\mathrm{\Gamma }_c}}(d\gamma +(1)^n{\displaystyle \gamma _0\omega _l(\gamma _1)})`$
$`=`$ $`({\displaystyle \underset{k}{}}d\gamma _i^k\widehat{}h_i^k+(1)^n{\displaystyle \underset{k}{}}\gamma _i^k\widehat{}dh_i^k+(1)^n{\displaystyle \underset{k}{}}{\displaystyle }(\gamma _i^kh_{i1}^k)\chi _{i_{\mathrm{\Gamma }_c}}(\omega _l(h_{i2}^k))`$
$`=`$ $`{\displaystyle \underset{k}{}}d\gamma _i^k\widehat{}h_i^k+(1)^n{\displaystyle \underset{k}{}}\gamma _i^k\widehat{}dh_i^k(1)^n{\displaystyle \underset{k}{}}{\displaystyle (\gamma _i^kh_{i1}^k)(1\widehat{}S(h_{i2}^k)dh_{i3}^k)}`$
$`(1)^n{\displaystyle \underset{k}{}}{\displaystyle (\gamma _i^kh_{i1}^k)(A_{l_i}(h_{i3}^k)\widehat{}S(h_{i2}^k)dh_{i4}^k)}`$
$`=`$ $`d\gamma \widehat{}h+(1)^{n+1}\gamma A_{l_i}(h_1)\widehat{}h_2.`$
Note the identity
$$D_{l_i}(\gamma \widehat{}h)=d\gamma \widehat{}h+(1)^{n+1}\gamma A_{l_i}(h_1)\widehat{}h_2=d(\gamma \widehat{}h)+(1)^n(\gamma \widehat{}h_1)\chi _{i_{\mathrm{\Gamma }_c}}(\omega _l(h_2)).$$
(97)
Assume now there is given a left covariant derivative $`D_l`$. In terms of the corresponding linear maps $`A_{l_i}`$ one obtains a family of left pre-connection forms $`\omega _{l_i}:H(\mathrm{\Gamma }(B_i)\widehat{}\mathrm{\Gamma }(H))^1`$ by
$$\omega _{l_i}(h)=1\widehat{}S(h_1)dh_2A_{l_i}(h_2)\widehat{}S(h_1)h_3.$$
Using (55) one obtains
$$(\pi _j^iid)_\mathrm{\Gamma }(\omega _{l_i}(h))=\varphi _{ij_\mathrm{\Gamma }}(\pi _i^jid)_\mathrm{\Gamma }(\omega _{l_j}(h),$$
thus one has by
$$\omega _l(h)=(\omega _{l_i}(h))_{iI}$$
a left pre-connection form $`\omega _l:H\mathrm{\Gamma }_c^1(𝒫)`$.
One easily verifies for $`\gamma \widehat{}h\mathrm{\Gamma }^n(B_i)\widehat{}H`$
$$D_{l_i}(\gamma \widehat{}h)=d\gamma \widehat{}h+(1)^{n+1}\gamma A_{l_i}(h_1)\widehat{}h_2=d(\gamma \widehat{}h)+(1)^n(\gamma \widehat{}h_1)\omega _{l_i}(h_2).$$
(98)
Using this formula it follows that
$$D_l(\gamma )=d\gamma +(1)^n\gamma _0\omega _l(\gamma _1)$$
(99)
for $`\gamma hor\mathrm{\Gamma }_c^n(𝒫)`$. It is immediate from the formulas (97) and (98) and Proposition 10 that the correspondence is bijective.
For right covariant derivatives the proof is analogous. $`\mathrm{}`$
Remark: Note that the foregoing proof also shows the bijectiv correspondence between left (right) pre-connection forms and families of linear maps $`A_{l,r_i}:H\mathrm{\Gamma }(B_i)`$ fulfilling (54) and (55).
###### Definition 8
A left (right) pre-connection form $`\omega _{l,r}`$ is called left (right) connection form, if
$`R`$ $``$ $`ker((id\epsilon )hor_c\chi _{i_{\mathrm{\Gamma }_c}}\omega _l),iI,`$ (100)
$`S^1(R)`$ $``$ $`ker((id\epsilon )hor_c\chi _{i_{\mathrm{\Gamma }_c}}\omega _r),iI`$ (101)
is satisfied.
###### Proposition 15
Left (right) connections are in bijective correspondence to left (right) connection forms.
Proof: The claim follows immedately from Proposition 12, Proposition 14 and (94). $`\mathrm{}`$
Remark: Note that classical connection forms are related to the connection forms considered above as follows: Let a classical principal bundle be given, let $`X`$ be a vector field on the total space $`Q`$, and let $`hC^{\mathrm{}}(G)`$ where $`G`$ is the structure group. A classical connection form is a Lie algebra valued 1-form $`\stackrel{~}{\omega }`$ of type Ad on $`Q`$. Then the formula
$$\omega _l(h)(X)=\stackrel{~}{\omega }(X)(h)$$
defines a left connection form $`\omega _l`$ in the above sense. Condition (100) with $`R=(ker\epsilon )^2`$ means that $`\stackrel{~}{\omega })`$ can be interpreted as a Lie algebra valued form. In this case (90) and (92) replace the usual conditions (type $`Ad`$, condition for fundamental vectors) for connection forms.
###### Definition 9
The left (right) curvature of a given left (right) covariant derivative is the linear map $`D_{l,r}^2:hor\mathrm{\Gamma }_c(𝒫)hor\mathrm{\Gamma }_c(𝒫)`$.
###### Definition 10
Let $`\omega _{l,r}`$ be a left(right) pre-connection form of a left (right) covariant derivative $`D_{l,r}`$. The linear maps $`\mathrm{\Omega }_{l,r}:H\mathrm{\Gamma }_c^2(𝒫)`$ defined by
$`\mathrm{\Omega }_l(h)`$ $`:=`$ $`d\omega _l(h){\displaystyle \omega _l(h_1)\omega _l(h_2)},`$ (102)
$`\mathrm{\Omega }_r(h)`$ $`:=`$ $`d\omega _r(h)+{\displaystyle \omega _r(h_2)\omega _r(h_1)}`$ (103)
are called the left (right) curvature form of a given left (right) covariant derivative.
Remark: In other words we take an analogue of the structure equation as definition of the curvature form.
###### Proposition 16
The left (right) curvature of a given left (right) covariant derivative is related to the left (right) curvature form by the identity
$`D_l^2(\gamma )`$ $`=`$ $`{\displaystyle \gamma _0\mathrm{\Omega }_l(\gamma _1)},\gamma hor\mathrm{\Gamma }_c(𝒫),`$ (104)
$`D_r^2(\gamma )`$ $`=`$ $`{\displaystyle \mathrm{\Omega }_r(\gamma _1)\gamma _0},\gamma hor\mathrm{\Gamma }_c(𝒫).`$ (105)
Proof: Because of the one-to-one correspondence between covariant derivatives on $`𝒫`$ and certain families of covariant derivatives on the trivilizations $`B_iH`$ it is sufficient to prove this assertion on a trivial bundle $`BH`$. In this case the linear map $`\omega _l`$ belonging to a left covariant derivative has the form
$$\omega _l(h)=(1S(h_1)dh_2))(A_l(h_2)\widehat{}S(h_1)h_3).$$
Therefore, one obtains for $`\mathrm{\Omega }_l`$
$`d\omega _l(h){\displaystyle \omega _l(h_1)\omega _l(h_2)}`$ $`=`$ $`1\widehat{}{\displaystyle dS(h_1)dh_2}{\displaystyle dA_l(h_2)\widehat{}S(h_1)h_3}`$
$`+{\displaystyle A_l(h_2)\widehat{}(dS(h_1))h_3}+{\displaystyle A_l(h_2)\widehat{}S(h_1)dh_3}`$
$`1\widehat{}{\displaystyle S(h_1)(dh_2)S(h_3)dh_4}{\displaystyle A_l(h_2)\widehat{}S(h_1)h_3S(h_4)dh_5}`$
$`+{\displaystyle A_l(h_4)\widehat{}S(h_1)(dh_2)S(h_3)h_5}{\displaystyle A_l(h_2)A_l(h_5)\widehat{}S(h_1)h_3S(h_4)h_6}`$
$`=`$ $`{\displaystyle dA_l(h_2)\widehat{}S(h_1)h_3}{\displaystyle A_l(h_2)A_l(h_3)\widehat{}S(h_1)h_4},`$
which leads for $`\gamma \mathrm{\Gamma }^n(B)`$ to
$$(\gamma \widehat{}h_1)\mathrm{\Omega }_l(h_2)=\gamma dA_l(h_1)\widehat{}h_2\gamma A_l(h_1)A_l(h_2)\widehat{}h_3.$$
On the other hand the left hand side of (104) is
$`(D_l)^2(\gamma h)`$ $`=`$ $`D_l(d\gamma \widehat{}h{\displaystyle (1)^n\gamma A_l(h_1)\widehat{}h_2})`$
$`=`$ $`(1)^{n+1}{\displaystyle (d\gamma )A_l(h_1)\widehat{}h_2}(1)^n{\displaystyle (d\gamma )A_l(h_1)\widehat{}h_2}`$
$`{\displaystyle \gamma dA_l(h_1)\widehat{}h_2}{\displaystyle \gamma A_l(h_1)A_l(h_2)\widehat{}h_3}`$
$`=`$ $`({\displaystyle \gamma dA_l(h_1)\widehat{}h_2}+{\displaystyle \gamma A_l(h_1)A_l(h_2)\widehat{}h_3})`$
$`=`$ $`{\displaystyle (\gamma \widehat{}h_1)\mathrm{\Omega }_l(h_2)}.`$
For right covariant derivatives the proof is analogous. $`\mathrm{}`$
Remark: The proof shows that there is a linear map $`F_{l,r}:H\mathrm{\Gamma }^2(B)`$ defined by
$`F_l(h)`$ $`:=`$ $`dA_l(h)+{\displaystyle A_l(h_1)A_l(h_2)},`$ (106)
$`F_r(h)`$ $`:=`$ $`dA_r(h){\displaystyle A_r(h_2)A_r(h_1)}`$ (107)
such that the left (right) curvature form of a given left (right) covariant derivative on a trivial QPFB has the form
$`\mathrm{\Omega }_l(h)`$ $`=`$ $`{\displaystyle F_l(h_2)\widehat{}S(h_1)h_3},`$ (108)
$`\mathrm{\Omega }_r(h)`$ $`=`$ $`{\displaystyle F_r(h_2)\widehat{}h_3S^1(h_1)}.`$ (109)
Using formula (55) and the Leibniz rule (taking into account $`\tau _{ij}(h_1)\tau _{ij}(h_2)=\epsilon (h)1`$) it is easy to verify that the family of linear maps $`F_{l,r_i}:H\mathrm{\Gamma }^2(B_i)`$ corresponding to a left (right) curvature form on a locally trivial QPFB satisfies
$$\pi _{j_{\mathrm{\Gamma }_m}}^i(F_{l,r_i}(h))=\tau _{ij}(h_1)\pi _{i_{\mathrm{\Gamma }_m}}^j(F_{l,r_j}(h_2))\tau _{ji}(h_3).$$
(110)
In general, an analogue of the Bianchi identity does not exist.
Now we make some remarks about the general form of the linear maps $`A_{l_i}:H\mathrm{\Gamma }^1(B_i)`$ corresponding to connections on a locally trivial QPFB. For this we use the functionals $`𝒳_i`$ corresponding to the right ideal $`R`$, which determines the right covariant differential calculus $`\mathrm{\Gamma }(H)`$ (see , and the appendix). Let the $`h^i+Rker\epsilon /R`$ be a linear basis in $`ker\epsilon /R`$. Then every element $`h\epsilon (h)1+Rker\epsilon /R`$ has the form $`𝒳_i(h)h^i+R`$. Since $`1kerA_{l_i}`$ and $`RkerA_{l_i}`$ it follows that $`A_{l_i}`$ is determined by its values on the $`h^k`$,
$$A_{l_i}(h)=\underset{k}{}𝒳_k(h)A_{l_i}(h^k).$$
In other words, to get a connection on the trivial pieces $`B_iH`$, one chooses $`A_i^k\mathrm{\Gamma }^1(B_i)`$ and defines the linear map $`A_{l_i}`$ by
$$A_{l_i}(h)=\underset{k}{}𝒳_k(h)A_i^k.$$
The connections so defined on the trivial pieces $`B_iH`$ do in general not give a connection on the locally trivial QPFB $`𝒫`$, because they do in general not fulfill the condition (55). If the right ideal $`R`$ fulfills (36), one can rewrite the condition (55) as a condition for the one forms $`A_i^k\mathrm{\Gamma }^1(B_i)`$. Recall that in this case $`\tau _{ij}(r_1)d\tau _{ji}(r_2)=0,rR`$ (cf. (38)), thus
$$\tau _{ij}(h_1)d\tau _{ij}(h_2)=\underset{k}{}𝒳_k(h)\tau _{ij}(h_1^k)d\tau _{ji}(h_2^k).$$
Furthermore, the condition (36) leads to the identity
$$\tau _{ij}(h_1)𝒳_l(h_2)\tau _{ji}(h_3)=\tau _{ji}(S(h_1)h_3)𝒳_l(h_2)=\underset{k}{}𝒳_k(h)\tau _{ji}(S(h_1^k)h_3^k)𝒳_l(h_2^k).$$
Putting now $`A_{l_i}(h)=_k𝒳_k(h)A_i^k`$ in (55) leads to the following condition for the forms $`A_i^k`$:
$$\pi _{j_{\mathrm{\Gamma }_m}}^i(A_i^k)=\underset{l}{}\tau _{ji}(S(h_1^k)h_3^k)𝒳_l(h_2^k)\pi _{i_{\mathrm{\Gamma }_m}}^j(A_i^l)+\tau _{ij}(h_1^k)d\tau _{ji}(h_2^k).$$
Note that, in the case $`I=(1,2)`$, it follows from the last formula that there exist connections. One can choose, e.g., one forms $`A_2^l`$ on the right, and solve the remaining equation for $`A_1^k`$ due to the surjectivity of $`\pi _{2_{\mathrm{\Gamma }_g}}^1`$. One can regard the set of all left (right) covariant derivatives $`𝒟_{l,r}`$ as a set with affine structure, where the corresponding vector space is characterized by
###### Proposition 17
A linear map $`C_{l,r}:hor\mathrm{\Gamma }_c(𝒫)hor\mathrm{\Gamma }_c(𝒫)`$ is a difference of two left (right) covariant derivatives if and only if:
$`C_{l,r}(1)=0,`$ (111)
$`C_{l,r}(hor\mathrm{\Gamma }_c^n(𝒫))hor\mathrm{\Gamma }_c^{n+1}(𝒫),`$ (112)
$`C_l(\gamma \alpha )=(1)^n\gamma C_l(\alpha );\gamma \mathrm{\Gamma }_c^n(B);\alpha hor\mathrm{\Gamma }_c(𝒫),`$ (113)
$`C_r(\alpha \gamma )=(1)^nC_r(\alpha )\gamma ;\gamma \mathrm{\Gamma }_c(B);\alpha hor\mathrm{\Gamma }_c^n(𝒫),`$ (114)
$`(C_{l,r}id)\mathrm{\Delta }_{𝒫_{\mathrm{\Gamma }_c}}=\mathrm{\Delta }_{𝒫_{\mathrm{\Gamma }_c}}C_{l,r},`$ (115)
$`C_{l,r}(ker\chi _{i_{\mathrm{\Gamma }_c}}|_{hor\mathrm{\Gamma }_c(𝒫)})ker\chi _{i_{\mathrm{\Gamma }_c}}|_{hor\mathrm{\Gamma }_c(𝒫)};iI.`$ (116)
This is immediate from Definition 5.
Because of (116) such a map $`C_{l,r}`$ defines a family of local maps $`C_{l,r_i}`$ by
$$C_{l,r_i}\chi _{i_{\mathrm{\Gamma }_c}}=\chi _{i_{\mathrm{\Gamma }_c}}C_{l,r}.$$
It is immediately that the set of left (right) connections is an affine subspace of $`𝒟_{l,r}`$. The elements of the corresponding vector space have the following additional property:
$`(id\epsilon )C_{l_i}(1r)`$ $`=`$ $`0,iI,rR,`$
$`(id\epsilon )C_{r_i}(1r)`$ $`=`$ $`0,iI,rS^1(R).`$
## 5 Example
Here we present an example of a $`U(1)`$-bundle over the quantum space $`S_{pq\varphi }^2`$. The quantum space $`S_{pq\varphi }^2`$ is treated in detail in and we restrict ourselves here to a brief summary.
The algebra $`P(S_{pq\varphi }^2)`$ of all polynomials over the quantum space $`S_{pq\varphi }^2`$ is constructed by gluing together two copies of a quantum disc along its classical subspace.
###### Definition 11
The algebra $`P(D_p)`$ of all polynomials over the quantum disc $`D_p`$ is defined as the algebra generated by the elements $`x`$ and $`x^{}`$ fulfilling the relation
$$x^{}xpxx^{}=(1p)1,$$
(117)
where $`0<p<1`$.
Let $`P(S^1)`$ be the algebra generated by the elements $`\alpha ,\alpha ^{}`$ fulfilling the relation
$$\alpha \alpha ^{}=\alpha ^{}\alpha =1.$$
$`P(S^1)`$ can be considered as the algebra of all trigonometrical polynomials over the circle $`S^1`$.
There exists a surjective homomorphism $`\varphi _p:P(D_p)P(S^1)`$ defined by
$`\varphi _p(x)`$ $`=`$ $`\alpha ,`$
$`\varphi _p(x^{})`$ $`=`$ $`\alpha ^{},`$
and one can consider this homomorphism as the “pull back” of the embedding of the circle into the quantum disc.
The algebra $`P(S_{pq\varphi }^2)`$ of all polynomials over the quantum space $`S_{pq\varphi }^2`$ is defined as
$$P(S_{pq\varphi }^2):=\{(f,g)P(D_p)P(D_q)|\varphi _p(f)=\varphi _q(g)\}.$$
(118)
It was shown in that one can also regard this algebra as the algebra generated by the elements $`f_1`$, $`f_1`$ and $`f_0`$ fulfilling the relations
$`f_1f_1qf_1f_1`$ $`=`$ $`(pq)f_0+(1p)1,`$ (119)
$`f_0f_1pf_1f_0`$ $`=`$ $`(1p)f_1,`$ (120)
$`f_1f_0pf_0f_1`$ $`=`$ $`(1p)f_1,`$ (121)
$`(1f_0)(f_1f_1f_0)`$ $`=`$ $`0,`$ (122)
where the isomorphism is given by $`f_1(x,y)`$, $`f_1(x^{},y^{})`$ and $`f_0(xx^{},1)`$. (Here, the generators of $`P(D_q)`$ are denoted by $`y`$ and $`y^{}`$.) It is proved in that the $`C^{}`$-closure $`C(S_{pq\varphi }^2)`$ of $`P(S_{pq\varphi }^2)`$ is isomorphic to the $`C^{}`$-algebra $`C(S_{\mu c}^2)`$ over the Podles sphere $`S_{\mu c}^2`$ for $`c>0`$.
Now, let us construct a class of QPFB’s with structure group $`U(1)`$ and base space $`S_{pq\varphi }^2`$. The algebra of polynomials $`P(U(1))`$ over $`U(1)`$ is the same algebra as $`P(S^1)`$. With $`\mathrm{\Delta }(a)=\alpha \alpha `$, $`\epsilon (\alpha )=1`$ and $`S(\alpha )=\alpha ^{}`$, $`P(U(1))`$ is a Hopf algebra. According to Proposition 3 we need just one transition function $`\tau _{12}:P(U(1))P(S^1)`$ to obtain a locally trivial QPFB. We define a class of transition functions $`\tau _{12}^{(n)}`$ as follows:
$`\tau _{12}^{(n)}(\alpha )`$ $`:=`$ $`\alpha ^n,`$
$`\tau _{12}^{(n)}(\alpha ^{})`$ $`:=`$ $`\alpha _{}^{}{}_{}{}^{n}.`$
It follows that
$`\tau _{21}^{(n)}(\alpha )`$ $`:=`$ $`\alpha _{}^{}{}_{}{}^{n},`$
$`\tau _{21}^{(n)}(\alpha ^{})`$ $`:=`$ $`\alpha ^n.`$
We obtain a class of locally trivial QPFB’s $`(𝒫^{(n)},\mathrm{\Delta }_{𝒫^{(n)}},P(U(1)),P(S_{pq\varphi }^2),\iota ,((\chi _p,ker\pi _p),(\chi _q,ker\pi _q)))`$ corresponding to these transition functions (see formulas (9) and (10)), where $`\iota `$ is the canonical embedding $`P(S_{pq\varphi }^2)𝒫^{(n)}`$ and $`\pi _{p,q}:P(S_{pq\varphi }^2)P(D_{p,q})`$ and $`\chi _{p,q}:𝒫^{(n)}P(D_{p,q})P(U(1))`$ are the restrictions of the canonical projections on $`P(S_{pq\varphi }^2)`$ and $`𝒫^{(n)}`$ respectively.
In the following, we restrict ourselves to the case $`n=1`$.
###### Proposition 18
Let $`JP(D_p)P(D_q)`$ be the ideal generated by the element
$$(xx^{}1)(yy^{}1).$$
Then $`𝒫^{(1)}`$ is algebra isomorphic to $`(P(D_p)P(D_q))/J`$.
Proof: $`(P(D_p)P(D_q))/J`$ by
$`a`$ $`=`$ $`1y+J,`$
$`a^{}`$ $`=`$ $`1y^{}+J,`$
$`b`$ $`=`$ $`x1+J`$
$`b^{}`$ $`=`$ $`x^{}1+J.`$
It is easy to see that $`(P(D_p)P(D_q))/J`$ can be considered as the algebra $`<a,a^{},b,b^{}>/\stackrel{~}{J}`$, where the ideal $`\stackrel{~}{J}`$ is generated by the relations
$`a^{}aqaa^{}`$ $`=`$ $`(1q)1,`$
$`b^{}bpbb^{}`$ $`=`$ $`(1p)1,`$
$`ba=ab,ba^{}=a^{}b^{},b^{}a=ab^{},b^{}a^{}`$ $`=`$ $`a^{}b^{},`$ (123)
$`(1aa^{})(1bb^{})`$ $`=`$ $`0.`$
Further consider the following elements in $`𝒫^{(1)}`$:
$`\stackrel{~}{a}`$ $`=`$ $`(1\alpha ,y\alpha ),`$
$`\stackrel{~}{a}^{}`$ $`=`$ $`(1\alpha ^{},y^{}\alpha ^{}),`$
$`\stackrel{~}{b}`$ $`=`$ $`(x\alpha ^{},1\alpha ^{}),`$
$`\stackrel{~}{b}^{}`$ $`=`$ $`(x^{}\alpha ,1\alpha ).`$
A short calculation shows that these elements fulfill the same relations (123) as the $`a`$, $`a^{}`$, $`b`$ and $`b^{}`$. Thus, there exists a homomorphism $`F:(P(D_p)P(D_q))/I𝒫^{(1)}`$ defined by
$$F(a):=\stackrel{~}{a},F(b):=\stackrel{~}{b},F(a^{}):=\stackrel{~}{a}^{},F(b^{}):=\stackrel{~}{b}^{}.$$
We will show that $`F`$ is an isomorphism. For surjectivity it is sufficient to show that the elements $`\stackrel{~}{a}`$, $`\stackrel{~}{a}^{}`$, $`\stackrel{~}{b}`$ and $`\stackrel{~}{b}^{}`$ generate the algebra $`𝒫^{(1)}`$. It is shown in Lemma 2 that the elements $`x^kx_{}^{}{}_{}{}^{l},k,l0`$ form a vector space basis of $`P(D_p)`$. Analogous the elements $`\alpha ^i,i`$ ($`a^1=a^{}`$), form a vector space basis in $`P(U(1))`$. Thus a general element $`fP(D_p)P(U(1))P(D_q)P(U(1))`$ has the form
$$f=(\underset{k,l0,i}{}c_{k,l,i}^px^kx_{}^{}{}_{}{}^{l}\alpha ^i,\underset{m,n0,j}{}c_{m,n,j}^qy^my_{}^{}{}_{}{}^{n}\alpha ^j).$$
$`f𝒫^{(1)}`$ means that there is the restriction
$$\underset{k,l0,i}{}c_{k,l,i}^p\alpha ^{kl}\alpha ^i=\underset{m,n0,j}{}c_{m,n,j}^q\alpha ^{mnj}\alpha ^j,$$
which leads to the following condition for the coefficients $`c_{k,l,i}^p`$ and $`c_{m,n,j}^q`$.
$$\underset{l0,s+l0}{}c_{s+l,l,t}^p=\underset{n0,n+s+t0}{}c_{s+t+n,n,t}^q,s,t.$$
(124)
$`f𝒫^{(1)}`$ has the form $`f=_{s,t}f_{s,t}`$, where
$$f_{s,t}=(\underset{l0,l+s0}{}c_{s+l,l,t}^px^{l+s}x_{}^{}{}_{}{}^{l}\alpha ^t,\underset{n0,n+s+t0}{}c_{s+t+n,n,t}^qy^{n+s+t}y_{}^{}{}_{}{}^{n}\alpha ^t)𝒫^{(1)}$$
due to (124). Because of (124) one can write $`f_{s,t}`$ as
$`f_{s,t}`$ $`=`$ $`{\displaystyle \underset{l0,l+s0}{}}c_{s+l,l,t}^p(x^{l+s}x_{}^{}{}_{}{}^{l}\alpha ^t,y^{m+s+t}y_{}^{}{}_{}{}^{m}\alpha ^t)`$
$`+{\displaystyle \underset{n0,n+s+t0}{}}c_{s+t+n,n,t}^q(x^{k+s}x_{}^{}{}_{}{}^{k}\alpha ^t,y^{n+l+t}y_{}^{}{}_{}{}^{n}\alpha ^t)`$
$`{\displaystyle \underset{l0,l+s0}{}}c_{s+l,l,t}^p(x^{k+s}x_{}^{}{}_{}{}^{k}\alpha ^t,y^{m+s+t}y_{}^{}{}_{}{}^{m}\alpha ^t).`$
The identity
$$(x^sx^lx_{}^{}{}_{}{}^{l}\alpha ^t,y^{s+t}y^ny_{}^{}{}_{}{}^{n}\alpha ^t)=\stackrel{~}{a}^{s+t+n}\stackrel{~}{a}^{}{}_{}{}^{n}\stackrel{~}{b}_{}^{s+l}\stackrel{~}{b}^{}{}_{}{}^{l},$$
which is a direct consequence of the definition of $`\stackrel{~}{a}`$, $`\stackrel{~}{a}^{}`$, $`\stackrel{~}{b}`$ and $`\stackrel{~}{b}^{}`$, shows that $`F`$ is surjective.
To show the injectivity of $`F`$ we define the homomorphisms $`F_{p,q}:(P(D_p)P(D_q))/IP(D_{p,q})P(U(1))`$ by $`F_{p,q}:=\chi _{p,q}F`$. Because of $`ker\chi _pker\chi _q=\{0\}`$, $`kerF=\{0\}`$ if and only if $`kerF_pkerF_q=\{0\}`$. First let us describe the ideals $`kerF_{p,q}`$. Let $`I_p`$ and $`I_q`$ be the ideals generated by $`1aa^{}`$ and $`1bb^{}`$ respectively. From (123) it is immediate that the algebras $`((P(D_p)P(D_q))/I)/I_{p,q}`$ are isomorphic to $`P(D_{p,q})P(U(1))`$, where the isomorphism $`((P(D_p)P(D_q))/I)/I_pP(D_p)P(U(1))`$ is defined by $`a1\alpha `$, $`bx1`$, and the isomorphism $`((P(D_p)P(D_q))/I)/I_qP(D_q)P(U(1))`$ is defined by $`ay1`$, $`b1\alpha `$.
Moreover, there are automorphisms $`\stackrel{~}{F}_{p,q}:P(D_{p,q})P(U(1))P(D_{p,q})P(U(1))`$ defined by
$`\stackrel{~}{F}_p(1\alpha )`$ $`:=`$ $`1\alpha ,\stackrel{~}{F}_q(1\alpha ):=1\alpha ,`$
$`\stackrel{~}{F}_p(1\alpha ^{})`$ $`:=`$ $`1\alpha ^{},\stackrel{~}{F}_q(1\alpha ^{}):=1\alpha ^{},`$
$`\stackrel{~}{F}_p(x1)`$ $`:=`$ $`x\alpha ^{},\stackrel{~}{F}_q(y1):=y\alpha ,`$
$`\stackrel{~}{F}_p(x^{}1)`$ $`:=`$ $`x^{}\alpha ,\stackrel{~}{F}_q(y^{}1):=y^{}\alpha ^{}.`$
Let $`\eta _{p,q}`$ be the quotient maps with respect to the ideals $`I_{p,q}`$. A short calculation shows that
$$F_{p,q}=\stackrel{~}{F}_{p,q}\eta _{p,q},$$
thus we have found $`kerF_{p,q}=I_{p,q}`$. It remains to show $`I_pI_q=\{0\}`$. There are the following identities in $`P(D_q)P(D_p)/I`$:
$`(1aa^{})a`$ $`=`$ $`qa(1aa^{}),(1aa^{})a^{}=q^1a^{}(1aa^{}),`$
$`(1bb^{})b`$ $`=`$ $`pb(1bb^{}),(1bb^{})b^{}=p^1b^{}(1bb^{}).`$
From these relations and the definition of $`I_p=kerF_p`$ follows that for $`fkerF_p`$ there exists an element $`\stackrel{~}{f}`$ such that $`f=(1aa^{})\stackrel{~}{f}`$. $`kerF_q`$ has an anlogous property with $`1bb^{}`$, instead of $`1aa^{}`$. Using that $`1xx^{}`$ is not a zero divisor in $`P(D_p)`$, see Lemma 3, it is now easy to see that $`fkerF_pkerF_q`$ is of the form $`f=(1aa^{})(1bb^{})\stackrel{~}{f}`$. Thus $`f=0`$, i.e. $`kerF_pkerF_q=0`$. $`\mathrm{}`$.
The proof has shown that $`𝒫^{(1)}<a,a^{},b,b^{}>/J`$, where $`J`$ is the ideal generated by the relations (123). Under this identification, the mappings belonging to the bundle can be given explicitly.
$`\mathrm{\Delta }_{𝒫^{(1)}}(a)=a\alpha ,\mathrm{\Delta }_{𝒫^{(1)}}(a^{})=a^{}\alpha ^{},`$
$`\mathrm{\Delta }_{𝒫^{(1)}}(b)=b\alpha ^{},\mathrm{\Delta }_{𝒫^{(1)}}(b^{})=b^{}\alpha ,`$
$`\chi _p(a)=1\alpha ,\chi _q(a)=y\alpha ,`$
$`\chi _p(a^{})=1\alpha ^{},\chi _q(a^{})=y^{}\alpha ^{},`$
$`\chi _p(b)=x\alpha ^{},\chi _q(b)=1\alpha ^{},`$
$`\chi _p(b^{})=x^{}\alpha ,\chi _q(b^{})=1\alpha ,`$
$`\iota (f_1)=ba,\iota (f_1)=a^{}b^{},\iota (f_0)=bb^{}.`$
In the classical limit $`p,q1`$ the algebra becomes commutative and only the relation $`(1aa^{})(1bb^{})=0`$ remains. It is easy to see that this relation, together with the natural requirement $`|a|1`$, $`|b|1`$, describes a subspace of $`^4`$ homeomorphic to $`S^3`$. The right $`U(1)`$-action is a simultaneous rotation in $`a`$ and $`b`$, and the orbit through $`b=0`$ is the fibre over the top $`(0,0,0)`$ of the base space (see the discussion in ).
To build a connection on this locally trivial QPFB, first we have to construct an adapted covariant differential structure on $`𝒫^{(1)}`$. By Definition 3, the adapted covariant differential structure is defined by giving differential calculi $`\mathrm{\Gamma }(P(D_p))`$ and $`\mathrm{\Gamma }(P(D_q))`$ and a right covariant differential calculus $`\mathrm{\Gamma }(P(U(1))`$ on the Hopf algebra $`P(U(1))`$.
As the differential calculi $`\mathrm{\Gamma }((P(D_{p,q}))`$ on the quantum discs $`D_{p,q}`$ we choose the calculi already used in and described in detail in . The differential ideal $`J(P(D_p)\mathrm{\Omega }(P(D_p))`$ determining $`\mathrm{\Gamma }(P(D_p))`$ is generated by the elements
$`x(dx)`$ $``$ $`p^1(dx)x,x^{}(dx^{})p(dx^{})x^{},`$
$`x(dx^{})`$ $``$ $`p^1(dx^{})x,x^{}(dx)p(dx)x^{}.`$
Exchanging $`x`$ with $`y`$ and $`p`$ with $`q`$ one obtains the differential ideal $`J(P(D_q))\mathrm{\Omega }(P(D_q))`$ determining $`\mathrm{\Gamma }(P(D_q))`$. The corresponding calculus $`\mathrm{\Gamma }(P(S_{pq\varphi }^2))`$ on the basis was explicitely described in .
Furthermore we use the right covariant differential calculus $`\mathrm{\Gamma }(P(U(1)))`$ determined by the right ideal $`R`$ generated by the element
$$\alpha +\nu \alpha ^{}(1+\nu )1,$$
where $`0<\nu 1`$. One easily verifies that $`R`$ fulfills (36). Thus the differential ideal $`J(P(S^1))`$ is generated by the sets (37), (38) and (39). Using these generators in the present case one obtains the following relations in $`\mathrm{\Gamma }_m(P(S^1))`$:
$`(d\alpha ^{})\alpha `$ $`=`$ $`\alpha (d\alpha ^{}),`$
$`(d\alpha ^{})\alpha `$ $`=`$ $`\nu \alpha (d\alpha ^{}),`$
$`(d\alpha ^{})\alpha `$ $`=`$ $`p\alpha (d\alpha ^{}),`$
$`(d\alpha ^{})\alpha `$ $`=`$ $`q\alpha (d\alpha ^{}).`$
Therefore $`d\alpha ^{}=d\alpha =0`$ for $`\nu ,p,q1`$, and the LC differential algebra $`\mathrm{\Gamma }_m(P(S_{pq\varphi }^2))`$ has the following form:
$`\mathrm{\Gamma }_m^0(P(S_{pq\varphi }^2))`$ $`=`$ $`P(S_{pq\varphi }^2),`$
$`\mathrm{\Gamma }_m^n(P(S_{pq\varphi }^2))`$ $`=`$ $`\mathrm{\Gamma }^n(P(D_p)){\displaystyle \mathrm{\Gamma }^n(P(D_q))};n>0.`$
$`\mathrm{\Gamma }(P(S_{pq\varphi }^2))`$ coincides with $`\mathrm{\Gamma }_m(P(S_{pq\varphi }^2))`$ for $`pq`$, and is embedded as a subspace defined by the gluing for $`p=q`$ (cf. ).
Now we want to construct a connection on the bundle $`𝒫^{(1)}`$ which can be regarded as the connection corresponding to the quantum magnetic monopole with strength $`g=\frac{1}{2}`$.
The functionals $`𝒳`$ and $`f`$ on $`P(U(1))`$ corresponding to the basis element $`(\alpha 1)+Rker\epsilon /R`$ are given by
$`𝒳(\alpha )`$ $`=`$ $`1;𝒳(\alpha ^{})=\nu ^1,`$
$`f(\alpha )`$ $`=`$ $`\nu ;f(\alpha ^{})=\nu ^1,`$
$`f(hk)`$ $`=`$ $`f(h)f(k);h,k𝒜(U(1)),`$
$`𝒳(hk)`$ $`=`$ $`𝒳(h)f(k)+\epsilon (h)𝒳(k).`$
$`𝒳`$ is a linear basis in the space of functionals annihilating $`1`$ and the right ideal $`R`$ (see also the appendix and ), i.e. $`𝒳`$ is a basis of the $`\nu `$-deformed Lie algebra corresponding to the differential calculus on $`U(1)`$. We define the linear maps $`A_{l_1}:P(U(1))\mathrm{\Gamma }(P(D_p))`$ and $`A_{l_2}:P(U(1))\mathrm{\Gamma }(P(D_q))`$ corresponding to a left connection on $`𝒫^{(1)}`$ by
$`A_{l_1}(h)`$ $`=`$ $`𝒳(h){\displaystyle \frac{1}{4}}(xdx^{}x^{}dx)hP(U(1)),`$ (125)
$`A_{l_2}(h)`$ $`=`$ $`𝒳(h){\displaystyle \frac{1}{4}}(y^{}dyydy^{})hP(U(1)).`$ (126)
Because of $`𝒳(R)=0`$ and $`𝒳(1)=0`$, $`A_{l_1}`$ and $`A_{l_2}`$ fulfill the conditions (54) and (79). Since there is no gluing $`\mathrm{\Gamma }_m^1(B)`$ the condition (55) is also fulfilled.
Moreover any choice of one forms to the right of $`𝒳`$ gives a connection.
A short calculation shows (see formula (106)) that the linear maps $`F_1:P(U(1))\mathrm{\Gamma }^2(P(D_p))`$ and $`F_2:P(U(1))\mathrm{\Gamma }^2(P(D_q))`$ corresponding to the curvature have the following form:
$`F_1(h)`$ $`=`$ $`𝒳(h){\displaystyle \frac{1}{4}}(1+p)dxdx^{}+{\displaystyle 𝒳(h_1)𝒳(h_2)\frac{1}{16}(xx^{}px^{}x)dxdx^{}},`$
$`F_2(h)`$ $`=`$ $`𝒳(h){\displaystyle \frac{1}{4}}(1+q)dydy^{}+{\displaystyle 𝒳(h_1)𝒳(h_2)\frac{1}{16}(yy^{}qy^{}y)dydy^{}}`$
In the classical case, the local connection forms $`A_{l}^{}{}_{1}{}^{}`$ and $`A_{l}^{}{}_{2}{}^{}`$ can be transformed, using suitable local coordinates, from the classical unit discs to the upper and lower hemispheres of the classical $`S^2`$. The resulting local connection forms on $`S^2`$ just coincide with the well-known magnetic potentials of the Dirac monopole of charge $`1/2`$.
To explain this we will briefly describe the classical Dirac monopole (see ).
The classical Dirac monopole is defined on $`^3\{0\}`$, which is of the same homotopy type as $`S^2`$. The corresponding gauge theory is a $`U(1)`$ theory, and the Dirac monopole is described as a connection on a $`U(1)`$ principal fibre bundle over $`S^2`$.
Let $`\{U_N,U_S\}`$ be a covering of $`S^2`$, where $`U_N`$ respectively $`U_S`$ is the closed northern respectively southern hemisphere, $`U_NU_S=S^1`$. One can write $`U_N`$ and $`U_S`$ in polar coordinates (up to the poles):
$`U_N`$ $`=`$ $`\{(\theta ,\varphi ),0<\theta \pi /2,0\varphi <2\pi \}\{N\},`$
$`U_S`$ $`=`$ $`\{(\theta ,\varphi ),\pi /2\theta <\pi ,0\varphi <2\pi \}\{S\}.`$
By
$$g_{12}^{(n)}(\varphi )=\mathrm{exp}(in\varphi ),0\varphi <2\pi ,n$$
a family of transition functions $`g_{12}^{(n)}:S^1U(1),n`$ is given. A standard procedure defines a corresponding family of $`U(1)`$ principal fibre bundles $`Q^{(n)}`$.
Let $`\xi _i:S^1U_i,i=N,S`$ be the embedding defined by $`\xi _i(\varphi )=(\pi /2,\varphi )`$. A connection on $`Q^{(n)}`$ is defined by two Lie algebra valued one forms $`\stackrel{~}{A}_N`$ und $`\stackrel{~}{A}_S`$ fulfilling
$$\xi _N^{}(\stackrel{~}{A}_N)=\xi _S^{}(\stackrel{~}{A}_S)+ind\varphi .$$
The Wu-Yang forms defined by
$`\stackrel{~}{A}_N^{(n)}`$ $`=`$ $`i{\displaystyle \frac{n}{2}}(1\mathrm{cos}\theta )d\varphi ,`$
$`\stackrel{~}{A}_S^{(n)}`$ $`=`$ $`i{\displaystyle \frac{n}{2}}(1+\mathrm{cos}\theta )d\varphi `$
fulfill these condition. $`\stackrel{~}{A}_N^{(n)}`$ and $`\stackrel{~}{A}_S^{(n)}`$ are vector potentials generating the magnetic field $`B=\frac{n}{2}\frac{\stackrel{}{r}}{|\stackrel{}{r}|^3}`$. The strength of the Dirac monopole is $`n/2`$.
The classical analogue $`\stackrel{~}{P}^{(n)}`$ to the above constructed locally trivial QPFB $`𝒫^{(n)}`$ is $`U(1)`$ principal fibre bundles over a space constructed by gluing together two discs over their boundaries. A disc $`D`$ can be regarded as a subspace of $``$:
$$D:=\{x,xx^{}1\}.$$
The space resulting from the gluing together two copies of $`D`$ over $`S^1=\{xD,xx^{}=1\}`$ is topologically isomorphic to the sphere $`S^2`$. Every $`xS^1`$ has the form $`x=\mathrm{exp}(i\varphi ),0\varphi <2\pi `$. The classical $`U(1)`$ bundles $`\stackrel{~}{P}^{(n)}`$ are given by transition functions $`\stackrel{~}{g}_{12}^{(n)}:S^1U(1)`$, which are obtained by $`\tau _{12}^{(n)}=(\stackrel{~}{g}_{21}^{(n)})^{}`$ ($``$ means here the pull-back) from the above transition functions of QPFB. The exchange of the indices comes from formula (9. One has $`\stackrel{~}{g}_{12}^{(n)}(\mathrm{exp}(i\varphi ))=\mathrm{exp}(in\varphi ),n`$. Obviously, the $`\stackrel{~}{P}^{(n)}`$ are topologically isomorphic to $`Q^{(n)}`$.
The classical analogue to the above defined connection on $`𝒫^{(1)}`$ is given by the following one forms on $`D`$ (see (125) and (126)):
$`\stackrel{~}{A}_1`$ $`=`$ $`{\displaystyle \frac{1}{4}}(xdx^{}x^{}dx),`$
$`\stackrel{~}{A}_2`$ $`=`$ $`{\displaystyle \frac{1}{4}}(x^{}dxxdx^{}).`$
Let $`\xi :S^1D`$ the embedding. A short calculation shows that $`\stackrel{~}{A}_1`$ and $`\stackrel{~}{A}_2`$ fulfill
$$\xi ^{}(\stackrel{~}{A}_1)=\xi ^{}(\stackrel{~}{A}_2)id\varphi .$$
Now one defines the following maps $`\eta _N:U_N\{N\}D`$ and $`\eta _S:U_S\{S\}D`$ by
$`\eta _N(\theta ,\varphi )`$ $`:=`$ $`\sqrt{1\mathrm{cos}\theta }\mathrm{exp}(i\varphi ),`$
$`\eta _S(\theta ,\varphi )`$ $`:=`$ $`\sqrt{1+\mathrm{cos}\theta }\mathrm{exp}(i\varphi ),`$
and one easily verifies
$`\stackrel{~}{A}_N^1`$ $`=`$ $`\eta _N^{}(\stackrel{~}{A}_1),`$
$`\stackrel{~}{A}_S^1`$ $`=`$ $`\eta _S^{}(\stackrel{~}{A}_2).`$
## 6 Final remarks
We have developed the general scheme of a theory of connections on locally trivial QPFB, including a reconstruction theorem for bundles and a nice characterization of connections in terms of local connection forms. Here we make some remarks about questions and problems arising in our context, and about possible future developments.
1. It is very important to look for more examples. Our example of a $`U(1)`$ bundle over a glued quantum sphere is essentially the same as the example of of an $`SU_q(2)`$ bundle over an anlogous glued quantum sphere. (Indeed, in another quantum disc is used, which, however, is isomorphic to the disc used in our paper – both are isomorphic to the shift algebra. The bundles are equivalent in the sense of the main Theorem of , which says that a QPFB with structure group $`H`$ is determined by a bundle with the classical subgroup of $`H`$ as structure group.) For other examples, one has to look for algebras with a covering (or being a gluing) such that the $`B_{ij}`$ are “big enough” to allow for nontrivial transition functions $`\tau _{ij}:HB_{ij}`$: $`B_{ij}`$ must contain in its center subalgebras being the homomorphic image of the algebra $`H`$. This seems to be possible only if $`H`$ has nontrivial classical subgroups and $`B_{ij}`$ contains suitable classical subspaces, as in our example. The following (almost trivial) example of a gluing along two noncommutative parts indicates that one may fall back to a gluing along classical subspaces in many cases: Let $`A_1=C^{}(𝔖)_\sigma C^{}(𝔖)=A_2`$ be two copies of a quantum sphere being glued together from shift algebras via the symbol map $`\sigma `$, as described in . Then the gluing $`A_1_{pr_{1,2}}A_2:=\{((a_1,a_2),(a_1^{},a_2^{}))A_1A_2|a_2=a_1^{}\}`$ (gluing of two quantum spheres along hemispheres) is obviously isomorphic to $`\{(a_1,a_2,a_3)C^{}(𝔖)C^{}(𝔖)C^{}(𝔖)|\sigma (a_1)=\sigma (a_2)=\sigma (a_3)\}=C^{}(𝔖)_\sigma C^{}(𝔖)_\sigma C^{}(𝔖)`$. This is a glued quantum sphere with a (quantum disc) membrane inside, glued along the classical subspaces. (This corresponds perfectly to the classical picture of gluing two spheres along hemispheres.)
2. The permanent need to work with covering completions is an unpleasant feature of the theory. It would therefore be very important to find some analogue of algebras of smooth functions in the noncommutative situation which have a suitable class of ideals forming a distributive lattice with respect to $`+`$ and $``$ (cf. \[5, Proposition 2\]). It is not clear if such a class exists even in classical algebras of differentiable functions.
3. Principal bundles are in the classical case of utmost importance in topology and geometry. In the above approach, one could e. g. ask for characteristic classes (trying to generalize the Chern-Weil construction), and for a notion of parallel transport defined by a connection (a naive idea would be to call a horizontal form parallel, if its covariant derivative vanishes).
4. For locally trivial QPFB, a suitable notion of locally trivial associated quantum vector bundle (QVB) exists (). Its definition (via cotensor products) is designed to have the usual correspondence between vector valued horizontal forms (of a certain “type”) and sections of the associated bundle. To a connection on a QPFB one can also associate connections on the corresponding QVB (assuming a certain differential structure there).
5. The notion of gauge transformation in our context is considered in . Gauge transformations are defined as isomorphisms of the left (right) $`B`$-module $`𝒫`$, with natural compatibility conditions. It turns out that the set of covariant derivatives is invariant under gauge transformations, whereas connections are not always transformed into connections.
## 7 Appendix
The purpose of this appendix is to collect some results about covariant differential calculi on quantum groups (, , ) and about coverings and gluings of algebras and differential algebras .
### 7.1 Covariant calculi on Hopf algebras
We freely use standard facts about Hopf algebras, including the Sweedler notation (e. g. $`\mathrm{\Delta }(h)=h_1h_2`$). We assume that the antipode is invertible.
A differential algebra over an algebra $`B`$ is a $``$-graded algebra $`\mathrm{\Gamma }(B)=_i\mathrm{\Gamma }^i(B)`$, $`\mathrm{\Gamma }^0((B)=B`$, equipped with a differential $`d`$, i. e. a graded derivative of degree $`1`$ with $`d^2=0`$. It is called differential calculus if it is generated as an algebra by the $`db`$, $`bB`$. A differential ideal of a differential algebra is a $`d`$-invariant graded ideal. There is always the universal differential calculus $`\mathrm{\Omega }(B)`$ determined by the property that every differential calculus $`\mathrm{\Gamma }(B)`$ is of the form $`\mathrm{\Gamma }(B)\mathrm{\Omega }(B)/J(B)`$ for some differential ideal $`J(B)`$.
If two algebras $`A`$, $`B`$ and differential algebras $`\mathrm{\Gamma }(A)`$, $`\mathrm{\Gamma }(B)`$ are given, an algebra homomorphism $`\psi :AB`$ is said to be differentiable with respect to $`\mathrm{\Gamma }(A)`$, $`\mathrm{\Gamma }(B)`$, if there exists a homomorphism $`\psi _\mathrm{\Gamma }:\mathrm{\Gamma }(A)\mathrm{\Gamma }(B)`$ of differential algebras extending $`\psi `$. For $`\mathrm{\Gamma }(A)=\mathrm{\Omega }(A)`$ this extension, denoted in this case by $`\psi _{\mathrm{\Omega }\mathrm{\Gamma }}`$, always exists. If, in addition, $`\mathrm{\Gamma }(B)=\mathrm{\Omega }(B)`$, the notation $`\psi _\mathrm{\Omega }`$ is used. $`J(B)=kerid_{\mathrm{\Omega }\mathrm{\Gamma }}`$ is a differential ideal $`J(B)\mathrm{\Omega }(B)`$ such that $`\mathrm{\Gamma }(B)=Omega(B)/J(B)`$. $`J(B)`$ is called the differential ideal corresponding to $`\mathrm{\Gamma }(B)`$.
Now we list some facts about covariant differential calculi.
###### Definition 12
A differential calculus $`\mathrm{\Gamma }(H)`$ over a Hopf algebra $`H`$ is called right covariant, if $`\mathrm{\Gamma }(H)`$ is a right $`H`$ comodule algebra with right coaction $`\mathrm{\Delta }^\mathrm{\Gamma }`$ such that
$$\mathrm{\Delta }^\mathrm{\Gamma }(h_0dh_1\mathrm{}dh_n)=\mathrm{\Delta }(h_0)(did)\mathrm{\Delta }(h_1)\mathrm{}(did)\mathrm{\Delta }(h_n).$$
(127)
$`\mathrm{\Gamma }(H)`$ is called left covariant, if $`\mathrm{\Gamma }(H)`$ is a left $`H`$-comodule algebra with left coaction $`{}_{}{}^{\mathrm{\Gamma }}\mathrm{\Delta }`$ such that
$${}_{}{}^{\mathrm{\Gamma }}\mathrm{\Delta }(h_0dh_1\mathrm{}dh_n)=\mathrm{\Delta }(h_0)(idd)\mathrm{\Delta }(h_1)\mathrm{}(idd)\mathrm{\Delta }(h_n).$$
(128)
$`\mathrm{\Gamma }(H)`$ is called bicovariant if it is left and right covariant.
Because of the universality property the universal differential calculus over any Hopf algebra is bicovariant. In the sequel we list some properties of right covariant differential calculi. The construction of left covariant differential algebras is analogous.
Let $`\mathrm{\Delta }^\mathrm{\Omega }`$ be the right coaction of the universal differential calculus $`\mathrm{\Omega }(H)`$ and let $`\mathrm{\Gamma }(H)`$ be a differential algebra over the Hopf algebra $`H`$. $`\mathrm{\Gamma }(H)`$ is right covariant if and only if the corresponding differential ideal $`J(H)\mathrm{\Omega }(H)`$ has the property
$$\mathrm{\Delta }^\mathrm{\Omega }(J(H))J(H)H.$$
Let us consider a right-covariant differential calculus $`\mathrm{\Gamma }(H)`$. Let $`\mathrm{\Gamma }_{inv}^1(H):=\{\gamma \mathrm{\Gamma }(H)|\mathrm{\Delta }^\mathrm{\Gamma }(\gamma )=\gamma 1\}`$. There exists a projection $`P:\mathrm{\Gamma }^1(H)\mathrm{\Gamma }_{inv}^1(H)`$ defined by
$$P(h^0dh^1)=S^1(h_2^0h_2^1)h_1^0dh_1^1.$$
Now one can define a linear map $`\eta _\mathrm{\Gamma }:H\mathrm{\Gamma }^1(H)`$ by
$$\eta _\mathrm{\Gamma }(h):=P(dh)=S^1(h_2)dh_1.$$
By an easy calculation one obtains the identity $`dh=h_2\eta _\mathrm{\Gamma }(h_1)`$. The linear map $`\eta _\mathrm{\Gamma }`$ has the following properties:
$`\mathrm{\Delta }^\mathrm{\Gamma }(\eta _\mathrm{\Gamma }(h))`$ $`=`$ $`\eta _\mathrm{\Gamma }(h)1,`$
$`\eta _\mathrm{\Gamma }(h)k`$ $`=`$ $`{\displaystyle k_2(\eta _\mathrm{\Gamma }(hk_1)\epsilon (h)\eta _\mathrm{\Gamma }(k_1))},`$
$`d\eta _\mathrm{\Gamma }(h)`$ $`=`$ $`{\displaystyle \eta _\mathrm{\Gamma }(h_2)\eta _\mathrm{\Gamma }(h_1)}`$
In the case $`\mathrm{\Gamma }(H)=\mathrm{\Omega }(H)`$ we use the symbol $`\eta _\mathrm{\Omega }`$.
The first degrees of right-covariant differential algebras are in one-to-one correspondence to right ideals $`Rker\epsilon H`$ in the following sense: First, if a differential calculus is given, $`R:=ker\eta _\mathrm{\Gamma }ker\epsilon `$ is a right ideal with the property $`Rker\epsilon `$, and one can prove that the subbimodule $`J^1(H)`$ corresponding to $`\mathrm{\Gamma }^1(H)\mathrm{\Omega }^1(H)/J^1(H)`$ is generated by the space $`\eta _\mathrm{\Omega }(R)=\{S^1(r_2)dr_1|rR\}`$. On the other hand, every right ideal $`Rker\epsilon `$ defines a right covariant differential algebra $`\mathrm{\Gamma }(H)=\mathrm{\Omega }(H)/J(H)`$, where the differential ideal $`J(H)\mathrm{\Omega }(H)`$ is generated by the set $`\eta _\mathrm{\Omega }(R)`$. Analogously, right ideals $`Rker\epsilon `$ also correspond to left covariant differential calculi. In this case, the differential ideal $`J(H)`$ corresponding to $`R`$ is generated by $`\{S(r_1)dr_2|rR\}`$. Bicovariant differential calculi are given by right ideals $`R`$ with the property $`S(r_1)r_3r_2HR;rR`$ (Ad-invariance).
Now one can choose a linear basis $`h_i+R`$ in $`ker\epsilon /R`$. This leads to a set of functionals $`𝒳_i`$ on $`H`$ annihilating $`1`$ and $`R`$ such that $`\eta _\mathrm{\Gamma }(h)=𝒳_i(h)\eta _\mathrm{\Gamma }(h_i)`$, $`hH`$. The set of elements $`\eta _\mathrm{\Gamma }(h_i)`$ is a left and right $`H`$ module basis in $`\mathrm{\Gamma }^1(H)`$, and the set of the $`𝒳_i`$ is a linear basis in the space of all functionals annihilating $`1`$ and $`R`$. It is obvious that $`dh=h_2𝒳_i(h_1)\eta _\mathrm{\Gamma }(h_i)`$. Besides the functionals $`𝒳_i`$ the linear basis in $`ker\epsilon /R`$ determines also functionals $`f_{ij}`$ on $`H`$ satisfying
$`f_{ij}(1)`$ $`=`$ $`\delta _{ij},`$
$`f_{ij}(hk)`$ $`=`$ $`{\displaystyle \underset{l}{}}f_{il}(h)f_{lj}(k)`$
$$𝒳_i(hk)=\underset{l}{}𝒳_l(h)f_{li}(k)+\epsilon (h)𝒳_i(k).$$
###### Definition 13
Let $`A`$ be a vector space and let $`H`$ be a Hopf algebra such that there exists linear map $`\mathrm{\Delta }_A:AAH`$. $`\mathrm{\Delta }_A`$ is called right $`H`$-coaction and $`A`$ is called right $`H`$ comodule if
$`(\mathrm{\Delta }_Aid)\mathrm{\Delta }_A`$ $`=`$ $`(id\mathrm{\Delta })\mathrm{\Delta }_A,`$ (129)
$`(id\epsilon )\mathrm{\Delta }_A`$ $`=`$ $`id.`$ (130)
If $`A`$ is an algebra and $`\mathrm{\Delta }_A`$ is an homomorphism of algebras then $`A`$ is called a right $`H`$ comodule algebra. The left coaction is defined analogously.
The definition of covariant differential calculi over Hopf algebras is easily generalized to $`H`$ comodule algebras:
###### Definition 14
A differential calculus $`\mathrm{\Gamma }(A)`$ over a right $`H`$ comodule algebra $`A`$ is called right covariant if the right coaction $`\mathrm{\Delta }_A^\mathrm{\Gamma }:\mathrm{\Gamma }(A)\mathrm{\Gamma }(A)H`$ defined by
$$\mathrm{\Delta }_A^\mathrm{\Gamma }(a_0da_1\mathrm{}da_n)=\mathrm{\Delta }_A(a_0)(did)\mathrm{\Delta }_A(a_1)\mathrm{}(did)\mathrm{\Delta }_A(a_n)$$
(131)
exists.
### 7.2 Covering and gluing
Let finite families $`(B_i)_{iI}`$, $`(B_{ij})_{(i,j)I\times ID}`$, $`D`$ the diagonal in $`I\times I`$, $`B_{ij}=B_{ji}`$, and homomorphisms $`\pi _j^i:B_iB_{ij}`$ be given. Then the algebra
$$B=\{(b_i)_{iI}_iB_i|\pi _j^i(b_i)=\pi _i^j(b_j)ij\}=:_{\pi _j^i}B_i$$
is called gluing of the $`B_i`$ along the $`B_{ij}`$ by means of the $`\pi _j^i`$. Special cases of gluings arise from coverings: A finite covering of an algebra $`B`$ is a finite family $`(J_i)_{iI}`$ of ideals in $`B`$ with $`_iJ_i=0`$. Taking now $`B_i=B/J_i`$, $`B_{ij}=B/(J_i+J_j)`$, $`\pi _j^i:B_iB_{ij}`$ the canonical projections $`b+J_ib+J_i+J_j`$, one can form the gluing
$$B_c=_{\pi _j^i}B_i,$$
which is called the covering completion of $`B`$ with respect to the covering $`(J_i)_{iI}`$. $`B`$ is always embedded in $`B_c`$ via the map $`K:b(b+J_i)_{iI}`$. The covering $`(J_i)_{iI}`$ is called complete if $`K`$ is also surjective, i.e. $`B`$ is isomorphic to $`B_c`$. Every two-element covering is complete, as well as every covering of a C\*-algebra. On the other hand, if $`B=_{\pi _j^i}B_i`$ is a general gluing, and $`p_i:BB_i`$ are the restrictions of the canonical projections, then $`(kerp_i)_{iI}`$ is a complete covering of $`B`$.
If $`\mathrm{\Gamma }(B)`$ is a differential algebra, a covering $`(J_i)_{iI}`$ of $`\mathrm{\Gamma }(B)`$ is said to be differentiable if the $`J_i`$ are differential ideals. A differential algebra $`\mathrm{\Gamma }(B)`$ with differentiable covering $`(J_i)_{iI}`$ is called LC differential algebra (LC = locally calculus), if the factor differential algebras $`\mathrm{\Gamma }(B)/J_i`$ are differential calculi over $`B/J_i^0`$ ($`J_i^0`$ the degree zero component of $`J_i`$) and $`J_i^00,i`$.
###### Definition 15
Let $`(B,(J_i)_{iI})`$ be an algebra with complete covering, let $`B_i=B/J_i`$, let $`\pi _i:BB_i`$ be the natural surjections, and let $`\mathrm{\Gamma }(B)`$ and $`\mathrm{\Gamma }(B_i)`$ be differential calculi such that $`\pi _i`$ are differentiable and $`(ker\pi _{i}^{}{}_{\mathrm{\Gamma }}{}^{})_{iI}`$ is a covering of $`\mathrm{\Gamma }(B)`$. Then $`(\mathrm{\Gamma }(B),(\mathrm{\Gamma }(B_i))_{iI})`$ is called adapted to $`(B,(J_i)_{iI})`$.
The following proposition is essential for Definition 3:
###### Proposition 19
Let $`(B,(J_i)_{iI})`$ be an algebra with complete covering, and let $`\mathrm{\Gamma }(B_i)`$ be differential calculi over the algebras $`B_i`$. Up to isomorphy there exists a unique differential calculus $`\mathrm{\Gamma }(B)`$ such that $`(\mathrm{\Gamma }(B),(\mathrm{\Gamma }(B_i))_{iI}))`$ is adapted to $`(B,(J_i)_{iI})`$.
As shown in , the differential ideal corresponding to $`\mathrm{\Gamma }(B)=\mathrm{\Omega }(B)/J(B)`$ is just $`J(B)=_{iI}ker\pi _{i}^{}{}_{\mathrm{\Omega }\mathrm{\Gamma }}{}^{}`$.
Finally, there is a proposition concerning the covering completion of adapted differential calculi:
###### Proposition 20
Let $`(\mathrm{\Gamma }(B),(\mathrm{\Gamma }(B_i))_{iI})`$ be adapted to $`(B,(J_i)_{iI})`$. Then the covering completion of $`(\mathrm{\Gamma }(B),(ker\pi _{i_\mathrm{\Gamma }})_{iI})`$ is an LC differential algebra over $`B_c`$. |
warning/0002/math0002086.html | ar5iv | text | # Theorem 1
Global Stability for Holomorphic Foliations in Kaehler Manifolds
Jorge Vitório Pereira<sup>1</sup><sup>1</sup>1Supported by IMPA-CNPq
Abstract. We prove the following theorem for Holomorphic Foliations in compact complex kaehler manifolds: if there is a compact leaf with finite holonomy, then every leaf is compact with finite holonomy. As corollary we reobtain stability theorems for compact foliations in Kaehler manifolds of Edwards-Millett-Sullivan and Hollman.
1. Introduction
The question of global stability is recurrent in the theory of foliations. The work of Ehresmann and Reeb establishes the so called global stability theorem, which says that if $``$ is a transversely orientable codimension one foliation in a compact connected manifold $`M`$ that has a compact leaf $`L`$ with finite fundamental group, then every leaf of $`L`$ is compact with finite holonomy group\[Ca-LN\]. Counterexamples for codimension greater than one are known since the birth of the theorem. Here we want to abolish the hypothesis on the codimension for a special kind of foliation, namely holomorphic foliations in complex Kaehler manifolds. In other words we are going to prove the following :
###### Theorem 1
Let $``$ be a holomorphic foliation of codimension q in a compact complex Kaehler manifold. If $``$ has a compact leaf with finite holonomy group then every leaf of $``$ is compact with finite holonomy group.
Another kind of stability problem was posed by Reeb and Haefliger. The question was the stability of compact foliations, that is, if a foliation has all leaves compact is the leaf space Hausdorff? Positive answers to this problem arose in the work of Epstein\[Ep\], Edwards-Millet-Sullivan\[EMS\], Holmann\[Ho\], etc. There are plenty situations where the leaf space is not Hausdorff. Sullivan found a example in the $`C^{\mathrm{}}`$ case\[Su\], Thurston in the analytic case\[Su\] and Müller in the holomorphic case\[Ho\]. The examples of Sullivan and Thurston live in compact manifolds, and Müller’s in a non-compact non-Kaehler manifold. As corollary of the theorem we reobtain Holmann’s result and a special case of \[EMS\]’s outstanding Theorem.
###### Corollary
\[EMS,Ho\] Suppose M is a complex Kaehler manifold. If $``$ is a compact foliation, i.e., every leaf is compact, then every leaf has finite holonomy group. Consequently, there is an upper bound on the volume of the leaves, and the leaf space is Hausdorff.
The author would like to thanks B. Scárdua for valuable conversations.
2. The Leaf Volume Function
Let $``$ be a holomorphic foliation of a complex Kaehler manifold $`(M,\omega )`$. As in \[Br\] we define
$$\mathrm{\Omega }=\{pM|\text{ the leaf }L_p\text{ through p is compact with finite holonomy}}$$
By the local stability theorem of Reeb\[Ca-LN\] $`\mathrm{\Omega }`$ is an open set of $`M`$. Set, for every $`p\mathrm{\Omega },`$ $`n(p)`$ to be the cardinality of the holonomy group of $`L_p`$. If $`d`$ is the dimension of the leaves then we define volume function of $``$:
$$T:\mathrm{\Omega }^+,\text{ }T\left(p\right)=n\left(p\right)_{L_p}\omega ^d$$
###### Lemma 1
$`T`$ is a continuous locally constant function in $`\mathrm{\Omega }`$.
###### Demonstration proof
The continuity is obvious. We have to prove that $`T`$ is locally constant. To do this we have just to observe that it is constant in the residual subset of $`\mathrm{\Omega }`$, formed by the union of leaves without holonomy\[G,p. 96\]. By the Reeb local stability theorem there is a saturated neighborhood for each leaf in this set where all leaves are homologous. Then using the closedness of $`\omega ^d`$ and Stokes Theorem we prove the lemma.
Remark \- In fact, the proof of this lemma is essentially already contained in \[Ho\].
3. A Lemma about Diff($`^n,0`$)
In 1905 Burnside\[Bu\] proved that if $`G`$ is a subgroup of $`GL(n,F)`$, where $`F`$ is a field of characteristic zero, with exponent $`e`$, then $`G`$ is finite with $`cardinality(G)e^{n^3}`$. Recalling that a group has exponent $`e`$ if every element $`g`$ belonging to the group is such that $`g^e=1`$. From the generalization of this result by Herzog-Praeger\[HP\] we obtain :
###### Lemma 2
If $`G`$ is a subgroup of $`Diff(^n,0)`$ with exponent $`e`$ then $`G`$ is finite with $`cardinality\left(G\right)e^n.`$
###### Demonstration proof
If for each element of $`G`$ we consider its derivative we obtain a subgroup of $`GL(n,)`$ with exponent $`e`$. Thus we only have to prove that the normal subgroup $`G_0`$ of $`G`$, formed by its elements tangent to the identity is the trivial group.
Let $`gG_0`$, then $`g^e=Id`$. Defining $`H\left(x\right)=_{i=1}^eDg\left(0\right)^ig^i\left(x\right)`$, we see that :
$$Hg\left(x\right)=Dg\left(0\right)Dg\left(0\right)^1\underset{i=1}{\overset{e}{}}Dg\left(0\right)^ig^{i+1}\left(x\right)=Dg\left(0\right)H\left(x\right)$$
Hence g is conjugated to its linear part, and therefore g must be the identity.
4. Proof of the Theorem 1
Let $``$ be as in the theorem. Consider the connected component of $`\mathrm{\Omega }`$ containing the leaf $`L`$ that is compact and with finite holonomy, and call it $`\mathrm{\Omega }_L`$. The volume function $`T`$ is constant in $`\mathrm{\Omega }_L`$ by Lemma 1, so if $`p\mathrm{\Omega }_L`$ we have that the leaf through $`p`$ is aproximated by leaves with uniformly bounded volume, so it has bounded volume and is compact(here we use the fact that the manifold is compact to achieve the compactness of the leaf). The holonomy group of $`\mathrm{\Omega }_L`$ has finite exponent, because for any transversal $`\mathrm{\Sigma }`$ of $`L_p`$, $`\mathrm{\Sigma }\mathrm{\Omega }_L`$ will be an open set such that every leaf of $`\mathrm{\Omega }_L`$ cuts it in at most $`m`$ points. Thus for every holomy germ $`h`$ of $`L_p`$, $`(h^{m!})_{|\mathrm{\Sigma }\mathrm{\Omega }_L}=Id`$. Analytic continuation implies that $`h^{m!}=Id`$. Using Lemma 2, we see that $`\mathrm{\Omega }_L=\mathrm{}`$, and prove the theorem.
The Corollary follows observing that the set of leaves without holonomy is residual and that we don’t need the compactness of the manifold to assure that a limit leaf is compact. Then the holonomy group of each leaf is finite and by the results of Epstein\[Ep\] we get the consequences.
Remark - The same proof works in a more general context. We have just to suppose that our foliation is transversely quasi-analytic and that there is a closed form which is positive on the (n-q)-planes of the distribution associated to the foliation.
References
\[Br\] Brunella, M. :A global stability theorem for transversely holomorphic foliations, Annals of Global Analysis and Geometry 15(1997), 179-186
\[Bu\] Burnside, W. : Proc. London Math. Soc. (2) 3 (1905), 435-440
\[Ca-LN\] Camacho, C. and Lins Neto, A. : Geometric theory of foliations, Birkhauser, 1985
\[EMS\] Edwards, R., Millett K. and Sullivan D. : Foliations with all leaves compact, Topology 16(1977), 13-32
\[Ep\] Epstein, D.B.A. : Foliations with all leaves compact, Ann. Inst. Fourier 26, 1(1976), 265-282
\[G\] Godbillon, C. : Feuilletages, études géométriques, Birkhäuser, Basel, 1991
\[HP\] Herzog M. and Praeger C. : On the order of linear groups with fixed finite exponent, Jr. of Algebra 43(1976), 216-220
\[Ho\] Holmann, H. : On the stability of holomorphic foliations, LNM 798(1980), 192-202
\[Su\] Sullivan, Dennis : A counterexample to the periodic orbit conjecture, Inst. Hautes Études Sci. Publ. Math. 46(1976),5-14
Jorge Vitório Bacellar dos Santos Pereira
Email : jvpimpa.br
Instituto de Matem tica Pura e Aplicada, IMPA
Estrada Dona Castorina, 110 - Jardim Bot nico
22460-320 - Rio de Janeiro, RJ, Brasil |
warning/0002/math0002025.html | ar5iv | text | # Algebraic duality for partially ordered sets
### AMS classification:
06A06, 06A15
## Introduction
The results presented in this paper can be considered as the algebraic counterpart of the duality in the theory of linear spaces. The outline of the construction looks as follows.
Several categories occur in the theory of partially ordered sets. The most general is the category $`𝒫𝒪𝒮𝒯`$ whose objects are partially ordered sets and the morphisms are the monotone mappings. Another category which will be used is $`𝒞`$ whose objects are (bounded) complete lattices and the morhisms are the lattice homomorphisms preserving universal bounds. Evidently $`𝒞`$ is the subcategory of $`𝒫𝒪𝒮𝒯`$.
To introduce the algebraic duality (I use the term ‘algebraic’ to avoid confusion with the traditional duality based on order reversal) the two element partially ordered set $`\mathrm{𝟐}`$ is used:
$$\mathrm{𝟐}=\{0,1\},0<1$$
Let $`P`$ be an object of $`𝒫𝒪𝒮𝒯`$. Consider its dual $`P^{}`$:
$$P^{}=\mathrm{𝖬𝗈𝗋}_{𝒫𝒪𝒮𝒯}(P,\mathrm{𝟐})$$
(1)
The set $`P^{}`$ has the pointwise partial order. Moreover, it is always the complete lattice with respect to this partial order (section 1). Furthermore, starting from $`P^{}𝒞`$ (bounded compete lattices) consider the set $`P^{}`$ of all morphisms in the appropriate category:
$$P^{}=\mathrm{𝖬𝗈𝗋}_𝒞(P^{},\mathrm{𝟐})$$
(2)
And again, the set of mappings $`P^{}`$ is pointwise partaially ordered. Finally, it is proved in section 2 that $`P^{}`$ is isomorphic to the initial partially ordered set $`P`$ (the isomorphism lemma 4):
$$P^{}P$$
The account of the results is organized as follows. First it is proved that $`P^{}`$ (1) is complete lattice. Then the embeddings $`p\lambda _p`$ and $`p\upsilon _p`$ of the poset $`P`$ into $`P^{}`$ are built (6). Then it is shown that the principal ideals $`[0,\lambda _p]`$ in $`P^{}`$ are prime for all $`pP`$ (lemma 3). Moreover, it is shown that there is no more principal prime ideals in $`P^{}`$. Finally, it is observed that the principle prime ideals on $`P^{}`$ are in 1-1 correspondence with the elements of $`P^{}`$.
## 1 The structure of the dual space
First define the pointwise partial order on the elements of $`P^{}`$ (1). For any $`x,yP^{}`$
$$xypPx(p)y(p)$$
(3)
Evidently the following three statements are equivalent for $`x,yP^{}`$:
$$\begin{array}{ccc}& xy& \\ pPx(p)=1\hfill & & y(p)=1\hfill \\ pPy(p)=0\hfill & & x(p)=0\hfill \end{array}$$
(4)
To prove that $`P^{}`$ is complete lattice, consider its arbitrary subset $`KP^{}`$ and define the following mappings $`u,v:P\mathrm{𝟐}`$:
$$u(p)=\{\begin{array}{ccc}\hfill 1,& kK& k(p)=1\hfill \\ \hfill 0,& kK& k(p)=0\hfill \end{array}v(p)=\{\begin{array}{ccc}\hfill 0,& kK& k(p)=0\hfill \\ \hfill 1,& kK& k(p)=1\hfill \end{array}$$
(5)
The direct calculations show that both $`u`$ and $`v`$ are monotone mappings: $`u,vP^{}`$ and
$$u=\underset{P^{}}{sup}K,v=\underset{P^{}}{inf}K$$
which proves that $`P^{}`$ is the complete lattice. Denote by 0,1 the universal bounds of the lattice $`P^{}`$:
$$pP\mathrm{𝟎}(p)=0,\mathrm{𝟏}(p)=1$$
Let $`p`$ be an element of $`P`$. Define the elements $`\lambda _p,\upsilon _pP^{}`$ associated with $`p`$: for all $`qP`$
$$\lambda _p(q)=\{\begin{array}{ccc}\hfill 0& ,& qp\hfill \\ \hfill 1& ,& \text{otherwise}\hfill \end{array}\upsilon _p(q)=\{\begin{array}{ccc}\hfill 1& ,& qp\hfill \\ \hfill 0& ,& \text{otherwise}\hfill \end{array}$$
(6)
###### Lemma 1
For any $`xP^{},pP`$
$$\begin{array}{ccc}\hfill x(p)=0& & x\lambda _p\text{ }\text{in }P^{}\hfill \\ \hfill x(p)=1& & x\upsilon _p\text{ }\text{in }P^{}\hfill \end{array}$$
(7)
### Proof.
Rewrite the left side of the first equivalency as $`qqpx(q)=0`$, hence $`q\lambda _p(q)=0x(q)=0`$, therefore $`x\lambda _p`$ by virtue of (4). The second equivalency is proved likewise. $`\mathrm{}`$
We shall focus on the ’inner’ characterization of the elements $`\lambda _p,\upsilon _p`$ in mere terms of the lattice $`P^{}`$ itself. To do it, recall the necessary definitions.
Let $`L`$ be a complete lattice. An element $`aL`$ is called join-irreducible (meet-irreducible) if it can not be represented as the join (resp., meet) of a collection of elements of $`L`$ different from $`a`$. To make this definition more verifiable introduce for every $`aL`$ the following elements of $`L`$:
$$\begin{array}{ccc}\hfill \stackrel{ˇ}{a}& =& inf_L\{xLx>a\}\hfill \\ \hfill \widehat{a}& =& sup_L\{yLy<a\}\hfill \end{array}$$
(8)
which do exist since $`L`$ is complete. Clearly, $`\stackrel{ˇ}{a}a\widehat{a}`$ and the equivalencies
$$\begin{array}{ccc}\hfill a\stackrel{ˇ}{a}& & a\text{ is meet-irreducible}\hfill \\ \hfill a\widehat{a}& & a\text{ is join-irreducible}\hfill \end{array}$$
(9)
follow directly from the above definitions.
###### Lemma 2
An element $`wP^{}`$ is meet irreducible if and only if it is equal to $`\lambda _p`$ for some $`pP`$. Dually, $`vP^{}`$ is join irreducible iff $`v=\upsilon _p`$ for some $`pP`$.
### Proof.
First prove that every $`\lambda _p`$ is meet irreducible. To do it we shall use the criterion (9). Let $`pP`$. Define $`uP^{}`$ as:
$$u(q)=\{\begin{array}{ccc}\hfill 0& ,& q<p\hfill \\ \hfill 1& ,& \text{otherwise}\hfill \end{array}$$
then the following equivalency holds:
$$xy(qx(q)=0q<p)$$
(10)
Now let $`x>\lambda _p`$, then $`x(p)=1`$ (otherwise (7) would enable $`x\lambda _p`$). Then $`x>\lambda _p`$ implies $`x\lambda _p`$, hence $`qx(q)=0qp`$, although $`q=p`$ is excluded, hence we get exactly the right side of (10). That means that
$$u=\underset{P^{}}{inf}\{xx>\lambda _p\}=\stackrel{ˇ}{\lambda }_p$$
differs from $`\lambda _p`$, hence $`\lambda _p`$ is meet irreducible by virtue of (8). The second dual statement is proved quite analogously.
Conversely, suppose we have a meet irreducible $`wP^{}`$, hence, according to (8), there exists $`pP^{}`$ such that $`\stackrel{ˇ}{w}(p)0`$ while $`w(p)=0`$. The latter means $`w\lambda _p`$ for this $`p`$. To disprove $`w<\lambda _p`$ rewrite $`\stackrel{ˇ}{w}(p)0`$ as $`\neg (inf\{xw<x\}=\lambda _p)`$ which is equivalent to
$$y(xw<xyx)\&\neg (y\lambda _p)$$
In particular, it must hold for $`x=\lambda _p`$, thus the assumption $`w<\lambda _p`$ implies $`yy\lambda _p\&\neg (y\lambda _p)`$, and the only remaining possibility for $`w`$ is to be equal to $`\lambda _p`$. $`\mathrm{}`$
### Dual statement.
The join irreducibles of $`P^{}`$ are the elements $`\upsilon _p,pP`$ and only they.
## 2 Second dual and the isopmorphism lemma
Introduce the necessary definitions. Let $`L`$ be a lattice. An ideal in $`L`$ is a subset $`KL`$ such that
* $`kK,xkxK`$
* $`a,bKabK`$
Replacing $``$ by $``$ and $``$ by $``$ the notion of filter is introduced. An ideal (filter) $`KL`$ is called prime if its set complement $`LK`$ is a filter (resp., ideal) in $`L`$. Now return to the lattice $`P^{}`$.
###### Lemma 3
For any $`pP`$ both the principal ideal $`[0,\lambda _p]`$ and the principal filter $`[\upsilon _p,1]`$ are prime in $`P^{}`$. Moreover,
$$[\upsilon _p,1]=P^{}[0,\lambda _p]$$
### Proof.
Fix up $`pP`$, then for any $`xP^{}`$ the value $`x(p)`$ is either 0 (hence $`x\lambda _p`$) or 1 (and then $`x\upsilon _p`$) according to (7). Since $`\lambda _p`$ never equals $`\upsilon _p`$ (because their values at $`p`$ are different), the sets $`[\upsilon _p,1]`$ and $`[0,\lambda _p]`$ are disjoint, which completes the proof. $`\mathrm{}`$
The converse statement is formulated in the following lemma.
###### Lemma 4
For any pair $`u,vP^{}`$ such that
$$[0,u]=P^{}[v,1]$$
(11)
there exists an element $`pP`$ such that $`u=\lambda _p`$ and $`v=\upsilon _p`$.
### Proof.
It follows from (11) that $`u`$ and $`v`$ are not comparable, therefore $`uv<v`$. Thus there exists $`pP`$ such that $`(uv)(p)=0`$ while $`v(p)=1`$. Then (7) implies $`uv\lambda _p`$ and $`v\upsilon _p`$. Suppose $`v\upsilon _p`$, then (11) implies $`\upsilon _pu`$, which together with $`\upsilon _pv`$ implies $`\upsilon _puv\lambda _p`$ which never holds since $`\upsilon _p`$ and $`\lambda _p`$ are not comparable. So, we have to conclude that $`v=\upsilon _p`$, thus $`u=\lambda _p`$. $`\mathrm{}`$
Now introduce the second dual $`P^{}`$ as the set of all homomorphisms of complete lattices $`P^{}\mathrm{𝟐}`$ preserving universal bounds, that is, for any $`𝐩P^{},KP^{}`$
$$\begin{array}{c}𝐩(sup_K)=sup_{kK}𝐩(k)\hfill \\ 𝐩(inf_K)=inf_{kK}𝐩(k)\hfill \\ 𝐩(0)=0\text{ ; }𝐩(1)=1\hfill \end{array}$$
with the pointwise partial order as in (3).
Now we are ready to prove the following isomorphism lemma.
###### Lemma 5
The partially ordered sets $`P`$ and $`P^{}`$ are isomorphic.
### Proof.
Define the mapping $`F:PP^{}`$ by putting
$$F(p)=𝐩:𝐩(x)=x(p)xP^{}$$
Evidently $`F`$ is the order preserving injection. To build the inverse mapping $`G:P^{}P`$, for any $`𝐩P^{}`$ consider the ideal $`𝐩^1(0)`$ and the filter $`𝐩^1(1)`$ in $`P^{}`$ both being prime (see , II.4). Let $`u=sup𝐩^1(0)`$ and $`v=inf𝐩^1(1)`$. Since $`𝐩`$ is the homomorphism of complete lattices, $`u𝐩^1(0)`$ and $`v𝐩^1(1)`$, hence $`𝐩^1(0)=[0,u]`$ and $`𝐩^1(1)=[v,1]`$. Applying lemma 4 we see that there exists $`pP`$ such that $`u=\lambda _p`$ and $`v=\upsilon _p`$. Put $`G(𝐩)=p`$. The mapping $`G`$ is order preserving and injective (since the different principal ideals have different suprema). It remains to prove that $`F,G`$ are mutually inverse.
Let $`pP`$, consider $`G(F(p))`$. Denote $`𝐩=F(p)`$, then $`𝐩^1(1)=\{xP^{}x(p)=0\}=\{xx\lambda _p\}`$. Thus $`sup𝐩^1(0)=\lambda _p`$, then $`GF=\mathrm{id}_P`$ which completes the proof. $`\mathrm{}`$
## Concluding remarks
The results presented in this paper show that besides the well known duality in partially ordered sets based on order reversal, we can establish quite another kind of duality à la linear algebra. As in the theory of linear topological spaces, we see that the ‘reflexivity’ expressed as $`P=P^{}`$ can be achieved by appropriate definition of dual space.
We see that a general partially ordered set have the dual space being a complete lattice. We also see that not every complete lattice can play the rôle of dual for a poset. These complete lattices can be characterized in terms of spaces with two closure operations . For the category of orthoposets this construction was introduced in . Another approach to dual spaces when they are treated as sets of two-valued measures (in terms of this paper, as sub-posets of $`P^{}`$) is in . The main feature of the techniques suggested in the present paper is that all the constructions are formulated in mere terms of partially ordered sets and lattices.
The work was supported by the RFFI research grant (97-14.3-62). The author acknowledges the financial support from the Soros foundation (grant A97-996) and the research grant ”Universities of Russia”. |
warning/0002/math0002209.html | ar5iv | text | # Divergence operators and odd Poisson brackets
## Introduction
Graded algebras with an odd Poisson bracket – also called Gerstenhaber algebras – play an important role in the theory of deformations of algebraic structures as well as in several areas of field theory, as has been shown by Batalin and Vilkovisky , Witten , Lian and Zuckerman , Getzler , Hata and Zwiebach , among others. Generators of odd Poisson brackets, in the sense of Equation (1) below, are differential operators of order $`2`$ of the underlying graded algebra, sometimes called “odd laplacians”, and usually denoted by the letter $`\mathrm{\Delta }`$. Batalin-Vilkovisky algebras – BV-algebras, for short – are a special class of these algebras, those for whose bracket there exists a generator assumed to be of square $`0`$. The geometrical approach to odd Poisson algebras and BV-algebras in terms of supermanifolds was first developed by Leites , Khudaverdian (see also ) and Schwarz .
The purpose of this article is to study various constructions of generators of odd Poisson brackets. Our constructions will rely on the general notion of divergence operator on a graded algebra, which generalizes the concept of the divergence of a vector field in elementary analysis. Given an odd Poisson bracket, to each element in the algebra we associate the divergence of the hamiltonian derivation that it defines, and we show that such a map from the algebra to itself, multiplied by the factor $`\frac{1}{2}`$ and by an appropriate sign, is a generator of that bracket (Theorem 1.2). We then adopt the language of supermanifolds to treat two constructions which determine divergence operators on the structural sheaf of the supermanifold.
The first construction uses berezinian volumes, and is modeled after the construction of divergence operators on smooth, purely even manifolds which uses volume elements. Once a divergence operator is defined, we apply Theorem 1.2 to obtain generators of an odd Poisson bracket on the supermanifold. One can deform any generator, obtained from a berezinian volume, by a change of berezinian volume, i.e., the multiplication by an invertible, even function. The deformed generator then differs from the original one by the addition of a hamiltonian derivation. If the original generator is of square $`0`$, a sufficient condition for a deformed generator to remain of square $`0`$ is given by a Maurer-Cartan equation. (Under the name “BV quantum master equation”, this condition plays a fundamental role in the BV quantization of gauge fields.) We study two special cases in detail: (i) the cotangent bundle of a smooth manifold viewed as a supermanifold whose structural sheaf is the sheaf of multivectors on the manifold, in which case the odd Poisson bracket under consideration is the Schouten bracket, and (ii) the tangent bundle of a smooth manifold as a supermanifold whose structural sheaf is the sheaf of differential forms on the manifold, when the underlying smooth, even manifold is a Poisson manifold.
The second construction utilizes graded linear connections on supermanifolds, and generalizes to the graded case an approach to the construction of divergence operators on purely even manifolds using linear connections. Given a linear connection on a smooth, even manifold, the divergence of a vector field $`X`$ is defined as the trace of the difference of the covariant derivation in the direction of $`X`$ and of the map $`[X,.]`$, where $`[,]`$ is the Lie bracket of vector fields. In the case of a supermanifold, given a graded linear connection, the divergence of a graded vector field is defined in a similar manner, replacing the trace by the supertrace, and the Lie bracket by the graded commutator. We then apply Theorem 1.2 to obtain a generator of an odd Poisson bracket on the supermanifold. As an example, we again study the generators of the Schouten bracket of multivectors on a manifold, equipped with a linear connection. On the one hand, there exists a unique generator of the Schouten bracket whose restriction to the vector fields is the divergence defined by the linear connection. On the other hand, on the cotangent bundle considered as a supermanifold, the linear connection on the manifold defines a graded metric in a simple way. This graded metric in turn determines a graded torsionless connection on this supermanifold – the associated Levi-Civita connection –, from which we obtain a generator of the Schouten bracket, following our general procedure. We show that these two generators of the Schouten bracket coincide and that, when the connection is flat, this generator is of square $`0`$.
One can ask: what happens if we deal with an even Poisson bracket instead of an odd one? The answer is that the phenomena in the odd and in the even cases are very different, although there is a formal similarity of the constructions. In the even case, the results extend those of the purely even case, e.g., the usual case of Poisson algebras of smooth manifolds. Taking the divergence of a hamiltonian vector field with respect to a volume form if the manifold is orientable, or, more generally, to a density, yields a derivation of the algebra, i.e., a vector field. This is the modular vector field that has been studied in Poisson geometry (see , ) and in more general contexts (see , , ). One can prove that this vector field is closed in the Poisson-Lichnerowicz cohomology. A change in the volume element modifies the vector field by a hamiltonian vector field, therefore the cohomology class of the closed vector field does not change. One thus obtains a cohomology class, called “the modular class” . So, while in the odd case we get a second-order differential operator which is a generator of the bracket, in the even case we get a first-order differential operator, in fact a derivation of the structural sheaf of associative algebras.
The paper is organized as follows. In Section 1, we first recall the definition of an odd Poisson bracket on a $`_2`$-graded algebra, and, in Section 1.2, we define algebraically the notion of a divergence operator on a graded algebra and that of its curvature. We then prove Theorem 1.2, which will serve as the main tool in our constructions, and we study the deformation of divergence operators and of the associated generators of odd Poisson brackets.
In Section 2, we study the divergence operators defined by berezinian volumes (Proposition 2.2), the associated generators of odd Poisson brackets, and their deformation under a change of berezinian volume. We show that the “Batalin-Vilkovisky quantum master equation” appears as a sufficient condition for the modified generator to remain of square 0 (Proposition 2.5). In Section 2.3, we consider the example of the cotangent bundle of a manifold $`M`$, considered as a supermanifold. To a volume element element $`\mu `$ on $`M`$, there corresponds a berezinian volume, which behaves like the “square of $`\mu `$”. In Theorem 2.8, we show that the generator of the Schouten bracket furnished by the general construction outlined above coincides with the generator obtained from the de Rham differential by the isomorphism defined by $`\mu `$ relating forms to multivectors. In Section 2.4, we treat the case of the tangent bundle of $`M`$, which is an odd Poisson supermanifold when $`M`$ has a Poisson structure. We prove that the generator defined by the divergence of hamiltonian vector fields with respect to the canonical berezinian volume coincides with the Poisson homology operator, and that its square therefore vanishes (Theorem 2.11). In Section 2.5, we express the properties of the supermanifolds studied in Sections 2.3 and 2.4 in the language of QS, SP and QSP manifolds of and (Theorems 2.16 and 2.17).
In Section 3, we study the divergence operators defined by graded linear connections. The definitions of a graded linear connection, its curvature and torsion, and of the divergence operator that it defines (Proposition 3.3) are given in Section 3.1. We then study the generator associated to a torsionless graded linear connection, of an odd Poisson bracket and the effect of a change of connection on the generator (Section 3.2). The Levi-Civita connection of a graded metric on a supermanifold is introduced in Section 3.3. In the remaining part of Section 3, we study the cotangent bundle of a manifold as an odd Poisson supermanifold. More specifically we study two constructions of generators of the Schouten bracket associated to a torsionless linear connection on the base manifolds, and we show that the two constructions yield the same generator (Theorem 3.16). We conclude the paper with remarks concerning the relationships between divergence operators, right and left module structures in the theory of Lie-Rinehart algebras and right and left $`𝒟`$-modules, and their analogues in the graded case , and we formulate a conjecture regarding the existence of a unique prolongation of a divergence operator on a supermanifold into a generator of an odd bracket on the algebra of graded multivectors.
We shall usually denote a supermanifold by a pair $`(M,𝒜)`$ where $`M`$ is an ordinary smooth manifold, called the base manifold, and $`𝒜`$ is a sheaf over $`M`$ of $`_2`$-graded commutative, associative algebras. The sections of $`𝒜`$ will be denoted by $`f,g,\mathrm{}`$, but this notation will be modified in some instances. When $`a`$ is an element of a $`_2`$-graded vector space, $`|a|`$ denotes the $`_2`$-degree of $`a`$ and, whenever it appears in a formula, it is understood that $`a`$ is homogeneous. The word “graded” will often be omitted. The bracket $`[,]`$ denotes the graded commutator. Manifolds and maps are assumed to be smooth. We recall some general properties of supermanifolds and the definition of the berezinian volumes in the Appendix.
## 1. Odd Poisson brackets and divergence operators
In this section, we review the main definitions concerning odd Poisson brackets on graded algebras and we study how to construct generators of such brackets.
### 1.1. Gerstenhaber and BV-algebras
Let $`𝐀`$ be a $`_2`$-graded commutative, associative algebra over a field of characteristic $`0`$. The multiplication map of algebra $`𝐀`$, $`(f,g)𝐀\times 𝐀fg𝐀`$, will be denoted by $`m`$. By definition, an odd Poisson bracket or a $`_2`$-Gerstenhaber bracket on $`𝐀`$ is an odd bilinear map, $`\pi :(f,g)𝐀\times 𝐀[[f,g]]𝐀`$, satisfying, for any $`f,g,h𝐀`$,
* $`[[f,g]]=(1)^{(|f|1)(|g|1)}[[g,f]]`$ (skew-symmetry),
* $`[[f,[[g,h]]]]=[[[[f,g]],h]]+(1)^{(|f|1)(|g|1)}[[g,[[f,h]]]]`$ (graded Jacobi identity),
* $`[[f,gh]]=[[f,g]]h+(1)^{(|f|1)|g|}g[[f,h]]`$ (Leibniz rule) .
(“Map $`\pi `$ is odd” means that $`|[[f,g]]|=|f|+|g|1`$ modulo $`2`$.) The pair $`(𝐀,\pi )`$ is then called an odd Poisson algebra or a $`_2`$-Gerstenhaber algebra.
A linear map of odd degree, $`\mathrm{\Delta }:𝐀𝐀`$, such that, for all $`f,g𝐀`$,
(1)
$$[[f,g]]=(1)^{|f|}\left(\mathrm{\Delta }(fg)(\mathrm{\Delta }f)g(1)^{|f|}f(\mathrm{\Delta }g)\right),$$
is called a generator or a generating operator of $`\pi `$ (or of bracket $`[[,]]`$). If there exists a generator $`\mathrm{\Delta }`$ of the bracket which is of square $`0`$, then $`(𝐀,\pi ,\mathrm{\Delta })`$ is called a $`_2`$-Batalin-Vilkovisky algebra, or BV-algebra for short.
Remark. Since a Gerstenhaber bracket in the usual sense, defined on a $``$-graded algebra, is of $``$-degree $`1`$, it is clear that it can also be considered to be a $`_2`$-Gerstenhaber bracket. Similarly, Batalin-Vilkovisky algebras in the usual, $``$-graded sense, are particular cases of $`_2`$-Batalin-Vilkovisky algebras.
A generator of an odd Poisson bracket is clearly not a derivation of the graded associative algebra $`(𝐀,m)`$, unless the bracket is identically $`0`$. We shall now see under what condition a generator $`\mathrm{\Delta }`$ of an odd Poisson bracket $`\pi `$ is a derivation of the graded Lie algebra $`(𝐀,\pi )`$. A straightforward computation, using the defining relation (1) for a generator, yields the identity,
(2)
$$\mathrm{\Delta }^2(fg)(\mathrm{\Delta }^2f)gf(\mathrm{\Delta }^2g)=(1)^{|f|}\left(\mathrm{\Delta }([[f,g]])[[\mathrm{\Delta }f,g]](1)^{|f|1}[[f,\mathrm{\Delta }g]]\right),$$
for all $`f,g𝐀`$. From (2), we obtain
###### Lemma 1.1.
A generator $`\mathrm{\Delta }`$ of an odd Poisson bracket $`\pi `$ is an odd derivation of the odd Poisson algebra, $`(𝐀,\pi )`$, if and only if the map $`\mathrm{\Delta }^2`$ is an even derivation of the graded associative algebra $`(𝐀,m)`$. In particular, if $`\mathrm{\Delta }^2=0`$, then $`\mathrm{\Delta }`$ is a derivation of $`(𝐀,\pi )`$.
Remark. In the case of a $``$-graded algebra and a the linear map $`\mathrm{\Delta }`$ of degree $`1`$, the map $`\mathrm{\Delta }^2`$, which is of degree $`2`$, is a derivation of $`(𝐀,m)`$ if and only if it vanishes. Therefore, in this case, a generator $`\mathrm{\Delta }`$ of $`\pi `$ is a derivation of $`(𝐀,\pi )`$ if and only if $`\mathrm{\Delta }^2=0`$.
The following identity, of which we shall make use in Section 2.5, is the result of another computation. If $`D`$ is an odd derivation of $`(𝐀,m)`$, then
(3)
$$[D,\mathrm{\Delta }](fg)([D,\mathrm{\Delta }]f)gf([D,\mathrm{\Delta }]g)=(1)^{|f|}\left(D[[f,g]][[Df,g]](1)^{|f|1}[[f,Dg]]\right).$$
Let $`\delta _m`$ be the graded Hochschild differential of algebra $`(𝐀,m)`$. Equation (1) expresses the equality
$$\pi (f,g)=(1)^{|f|1}(\delta _m\mathrm{\Delta })(f,g).$$
Let us also introduce the graded Chevalley-Eilenberg differential of $`(𝐀,\pi )`$, denoted by $`\delta _\pi `$. Equation (2) expresses the equality
$$(\delta _\pi \mathrm{\Delta })(f,g)=(1)^{|f|1}(\delta _m(\mathrm{\Delta }^2))(f,g),$$
and Lemma 1.1 can be reformulated as follows : $`\mathrm{\Delta }`$ is a $`\delta _\pi `$-cocycle if and only if $`\mathrm{\Delta }^2`$ is a $`\delta _m`$-cocycle.
### 1.2. Divergence operators
Let $`(𝐀,m)`$ be a $`_2`$-graded commutative, associative algebra, and let $`\mathrm{Der}𝐀`$ be the graded vector space of graded derivations of $`(𝐀,m)`$. By definition, a divergence operator on $`𝐀`$ is an even linear map, $`\mathrm{div}:\mathrm{Der}𝐀𝐀`$, such that
(4)
$$\mathrm{div}(fD)=f\mathrm{div}(D)+(1)^{|f||D|}D(f),$$
for any $`D\mathrm{Der}𝐀`$ and any $`f𝐀`$.
This definition obviously generalizes the usual notion of divergence of vector fields in elementary analysis. More generally, if $`𝐀`$ is the purely even algebra of smooth real- or complex-valued functions on a smooth manifold, divergence operators on $`𝐀`$ can be defined by means of either volume forms or linear connections, the two approaches being related in a simple way. See . These two approaches will be generalized below in Sections 2 and 3, respectively.
We define the curvature of a divergence operator, as the bilinear map, $`^{\mathrm{div}}:\mathrm{Der}𝐀\times \mathrm{Der}𝐀𝐀`$ by
(5)
$$^{\mathrm{div}}(D_1,D_2)=\mathrm{div}[D_1,D_2]D_1(\mathrm{div}D_2)+(1)^{|D_1||D_2|}D_2(\mathrm{div}D_1).$$
for any derivations $`D_1,D_2\mathrm{Der}𝐀`$. A short computation shows that $`^{\mathrm{div}}`$ is $`𝐀`$-bilinear. If $`\delta _{[,]}`$ denotes the graded Chevalley-Eilenberg differential of the graded Lie algebra $`(\mathrm{Der}𝐀,[,])`$ acting on cochains on $`\mathrm{Der}𝐀`$ with values in the $`\mathrm{Der}𝐀`$-module $`𝐀`$, then the definition of $`^{\mathrm{div}}`$ in (5) can be written
$$^{\mathrm{div}}=\delta _{[,]}(\mathrm{div}).$$
### 1.3. Divergence operators and generators
On the odd Poisson algebra $`(𝐀,\pi )`$, we define the hamiltonian mapping by $`X^\pi :𝐀\mathrm{Der}𝐀`$, defined by
$$fX_f^\pi =[[f,.]].$$
The graded Jacobi identity for the bracket $`\pi `$ is equivalent to the relation
(6)
$$X_{[[f,g]]}^\pi =[X_f^\pi ,X_g^\pi ],$$
for any $`f,g𝐀`$, which expresses the fact that $`X^\pi `$ is a morphism of graded Lie algebras from $`(𝐀,\pi )`$ to $`(\mathrm{Der}𝐀,[,])`$.
We now introduce the main object of interest in this paper, the odd linear map $`\mathrm{\Delta }:𝐀𝐀`$, depending on both $`\pi `$ and the choice of a divergence operator, defined by
(7)
$$\mathrm{\Delta }f=(1)^{|f|}\frac{1}{2}\mathrm{div}(X_f^\pi ),$$
for $`f𝐀`$.
###### Theorem 1.2.
The operator $`\mathrm{\Delta }`$ on $`𝐀`$, defined by (7), is a generator of bracket $`\pi `$.
Proof. To show that Equation (1) is satisfied, we compute $`\mathrm{\Delta }(fg)`$ using the Leibniz rule for the odd Poisson bracket, and the fundamental property (4) of the divergence operators,
$`(1)`$ $`{}_{}{}^{|f|+|g|}\mathrm{\Delta }(fg)={\displaystyle \frac{1}{2}}\mathrm{div}X_{fg}^\pi ={\displaystyle \frac{1}{2}}\mathrm{div}(fX_g^\pi +(1)^{|f||g|}gX_f^\pi )`$
$`=(1)^{|g|}f\mathrm{\Delta }g+(1)^{|f|(|g|+1)}g\mathrm{\Delta }f+{\displaystyle \frac{1}{2}}(1)^{|f|(|g|+1)}[[g,f]]+{\displaystyle \frac{1}{2}}(1)^{|g|}[[f,g]]`$
$`=(1)^{|g|}f\mathrm{\Delta }g+(1)^{|f|+|g|}(\mathrm{\Delta }f)g+(1)^{|g|}[[f,g]],`$
and the result follows.∎
This fundamental result must be contrasted with a parallel, but strikingly different result valid for even Poisson brackets, and consequently also, in the usual case, for ungraded Poisson algebras. If, in (7), we replace the hamiltonian operator defined by an odd Poisson bracket by the one defined by an even Poisson bracket, a computation similar to the proof of Theorem 1.2 shows that the operator thus defined is a derivation of the associative multiplication. When the Poisson algebra is the algebra of functions of a smooth orientable Poisson manifold, in which case the divergence operator is the one associated with a volume element, this derivation is a vector field. It is easy to see that, up to the factor $`1/2`$, it coincides with the modular vector field of the Poisson manifold (also , and references cited therein).
We now establish a relation between the operator $`\mathrm{\Delta }`$ defined by $`\pi `$ and $`\mathrm{div}`$, and the curvature of the divergence operator evaluated on hamiltonian derivations.
###### Proposition 1.3.
For any $`f,g𝐀`$,
(8)
$$\mathrm{\Delta }[[f,g]][[\mathrm{\Delta }f,g]](1)^{|f|1}[[f,\mathrm{\Delta }g]]=(1)^{|f|+|g|1}\frac{1}{2}^{\mathrm{div}}(X_f^\pi ,X_g^\pi ).$$
Proof. By (6) and the definition of $`\mathrm{\Delta }`$,
$`\mathrm{\Delta }[[f,g]]`$ $`=(1)^{|f|+|g|1}{\displaystyle \frac{1}{2}}\mathrm{div}X_{[[f,g]]}^\pi =(1)^{|f|+|g|1}{\displaystyle \frac{1}{2}}\mathrm{div}[X_f^\pi ,X_g^\pi ],`$
$`[[f,\mathrm{\Delta }g]]`$ $`=(1)^{|g|}{\displaystyle \frac{1}{2}}X_f^\pi (\mathrm{div}X_g^\pi ),[[\mathrm{\Delta }f,g]]=(1)^{|f||g|1}{\displaystyle \frac{1}{2}}X_g^\pi (\mathrm{div}X_f^\pi ).`$
In view of the definition of the curvature, the proposition follows. ∎
###### Corollary 1.4.
The generator $`\mathrm{\Delta }`$ is a derivation of $`(𝐀,\pi )`$ if and only if $`^{\mathrm{div}}`$ vanishes on the hamiltonian derivations.
In terms of the differentials $`\delta _\pi `$ and $`\delta _{[,]}`$, Equation (8) can be written
$$(\delta _\pi \mathrm{\Delta })(f,g)=(1)^{|f|+|g|1}\frac{1}{2}(\delta _{[,]}\mathrm{div})(X_f^\pi ,X_g^\pi ).$$
### 1.4. Deformations of divergence operators and of generators
Since the difference of two divergence operators is an $`𝐀`$-linear map from $`\mathrm{Der}𝐀`$ to $`𝐀`$, the space of divergence operators on $`𝐀`$ is an affine space over $`\mathrm{Hom}_𝐀(\mathrm{Der}𝐀,𝐀)`$. We shall be interested in the case where the difference of two divergence operators is an evaluation map, $`D\mathrm{Der}𝐀D(2w)𝐀`$, where $`w`$ is a fixed, even element in $`𝐀`$. (The factor $`2`$ is conventional.)
Since the difference of two generators of an odd Poisson bracket $`\pi `$ on $`𝐀`$ is a derivation of $`(𝐀,m)`$, the space of generators of $`\pi `$ is an affine space over $`\mathrm{Der}𝐀`$.
###### Proposition 1.5.
Let $`\mathrm{div}`$ and $`\mathrm{div}^{}`$ be divergence operators on $`𝐀`$ such that there exists an even $`w𝐀`$ satisfying
(9)
$$\mathrm{div}^{}D=\mathrm{div}D+D(2w),$$
for all $`D\mathrm{Der}𝐀`$. For $`\pi `$ a fixed odd Poisson bracket, let $`\mathrm{\Delta }`$ and $`\mathrm{\Delta }^{}`$ be the generators of the bracket defined by (7) for $`\mathrm{div}`$ and $`\mathrm{div}^{}`$, respectively. Then
(10)
$$\mathrm{\Delta }^{}=\mathrm{\Delta }+X_w^\pi .$$
Proof. Equation (10) follows from (9) and the skew-symmetry of the bracket. ∎
Remark. In the case of an even Poisson algebra, and, in particular, in the usual case of the Poisson algebra of a smooth manifold, a similar argument is valid. As a consequence, one proves that the class of the modular vector field in the Poisson cohomology is well-defined, independently of the choice of a volume element .
We shall now consider under what condition a generator of $`\pi `$ with vanishing square remains of square $`0`$ when modified by the addition of an interior derivation $`X_w^\pi `$.
###### Proposition 1.6.
If $`\mathrm{\Delta }`$ is a generator of square $`0`$ of bracket $`\pi `$, and $`w`$ is an even element of $`𝐀`$, then
(11)
$$(\mathrm{\Delta }+X_w^\pi )^2=X_{\mathrm{\Delta }w+\frac{1}{2}[[w,w]]}^\pi .$$
Proof. From Lemma 1.1, we know that $`\mathrm{\Delta }^2=0`$ implies that $`\mathrm{\Delta }`$ is a derivation of $`\pi `$, whence $`[\mathrm{\Delta },X_w^\pi ]=X_{\mathrm{\Delta }w}^\pi `$. Using this relation and (6), we obtain
$$(\mathrm{\Delta }+X_w^\pi )^2=[\mathrm{\Delta },X_w^\pi ]+\frac{1}{2}[X_w^\pi ,X_w^\pi ]=X_{\mathrm{\Delta }w}^\pi +\frac{1}{2}X_{[[w,w]]}^\pi ,$$
whence the result. ∎
The equation
(12)
$$\mathrm{\Delta }w+\frac{1}{2}[[w,w]]=0,$$
is the Maurer-Cartan equation, familiar from deformation theory. (See .) Thus we can state
###### Corollary 1.7.
Let $`\mathrm{\Delta }`$ be a generator of square $`0`$ of bracket $`\pi `$, and let $`w`$ be an even element of $`𝐀`$. The generator $`\mathrm{\Delta }^{}=\mathrm{\Delta }+X_w^\pi `$ is of square $`0`$ if and only if the hamiltonian operator $`X_{\mathrm{\Delta }w+\frac{1}{2}[[w,w]]}^\pi `$ vanishes. If, in particular, $`w`$ satisfies the Maurer-Cartan equation (12), then the generator $`\mathrm{\Delta }^{}=\mathrm{\Delta }+X_w^\pi `$ is of square $`0`$.
## 2. Berezinian volumes and generators of odd Poisson brackets
An odd Poisson supermanifold (resp., a BV-supermanifold) is a supermanifold $`(M,𝒜)`$ whose sheaf of functions, $`𝒜`$, is a sheaf of odd Poisson algebras (resp., of BV-algebras). In the context of supermanifold theory, an odd Poisson bracket is often referred to as an antibracket or a Buttin bracket . The notions of derivations and divergence operators that we have introduced in Section 1 have obvious analogues in the case of sheaves of algebras, and we shall use the same symbols. On a supermanifold, a derivation of the sheaf of functions is called a graded vector field. We use the term “operator from $`𝒜_1`$ to $`𝒜_2`$” for a morphism of sheaves of vector spaces from $`𝒜_1`$ to $`𝒜_2`$.
In this section, we show how generators of an odd Poisson bracket on a supermanifold can be obtained from berezinian volumes.
### 2.1. Divergence operators defined by berezinian volumes
We first recall the main properties of the Lie derivatives of berezinian sections. (See the Appendix for the definition of the berezinian sheaf. See, e.g., for a proof of the following proposition.)
###### Proposition 2.1.
Let $`\xi `$ be a berezinian section on $`(M,𝒜)`$. For any graded vector field $`D`$ and any section $`f`$ of $`𝒜`$,
(13)
$$_D(\xi .f)=_D(\xi ).f+(1)^{|D||\xi |}\xi .D(f),$$
and
(14)
$$_{f.D}(\xi )=(1)^{|f|(|D|+|\xi |)}_D(\xi .f).$$
We shall now define the divergence operator associated with a berezinian volume..
###### Proposition 2.2.
Let $`\xi `$ be a berezinian volume. For any graded vector field $`D`$, there exists a unique section, $`\mathrm{div}_\xi (D)`$, of $`𝒜`$ such that
(15)
$$_D\xi =(1)^{|D||\xi |}\xi .\mathrm{div}_\xi (D).$$
The map $`D\mathrm{div}_\xi (D)`$ from $`\mathrm{Der}𝒜`$ to $`𝒜`$ defined by (15) is a divergence operator.
Proof. The map $`\mathrm{div}_\xi `$ is even, since $`|\xi |+|\mathrm{div}_\xi (D)|=|\xi .\mathrm{div}_\xi (D)|=|_D\xi |=|D|+|\xi |.`$ We must prove that $`\mathrm{div}_\xi `$ satisfies (4). Using Proposition 2.1, we obtain
$`\xi .\mathrm{div}_\xi (fD)`$ $`=(1)^{(|f|+|D|)|\xi |}_{fD}\xi =(1)^{|D|(|f|+|\xi |)}_D(\xi .f)`$
$`=(1)^{|D|(|f|+|\xi |)}(_D(\xi ).f+(1)^{|D||\xi |}\xi .D(f))`$
$`=(1)^{|f||D|}(\xi .\mathrm{div}_\xi (D)f+\xi .D(f)),`$
whence the result. ∎
Example. If $`(M,𝒜)=^{m|n}`$ with graded coordinates $`(x^1,\mathrm{},x^m,s^1,\mathrm{},s^n)`$, the section
$$\xi =[\mathrm{d}^Gx^1\mathrm{}\mathrm{d}^Gx^m\frac{}{s^1}\mathrm{}\frac{}{s^n}]$$
is a berezinian volume and, if $`D=_{i=1}^mg^i\frac{}{x^i}+_{\rho =1}^nh^\rho \frac{}{s^\rho }`$, then
$$\mathrm{div}_\xi (D)=\underset{i=1}{\overset{m}{}}\frac{g^i}{x^i}+\underset{\rho =1}{\overset{n}{}}(1)^{|h^\rho |}\frac{h^\rho }{s^\rho }.$$
A short computation shows that, for graded vector fields $`D_1`$ and $`D_2`$,
(16)
$$\mathrm{div}_\xi [D_1,D_2]=D_1(\mathrm{div}_\xi D_2)(1)^{|D_1||D_2|}D_2(\mathrm{div}_\xi D_1).$$
In view of the definition of the curvature of a divergence operator (5), we have proved
###### Proposition 2.3.
For any berezinian volume $`\xi `$, the curvature $`^{\mathrm{div}_\xi }`$ of the divergence operator $`\mathrm{div}_\xi `$ vanishes.
We now consider the effect on the divergence operator of a change of berezinian volume. When $`v`$ is an invertible, even section of $`𝒜`$, then $`\xi .v`$ is also a generator of the berezinian sheaf. We remark that any invertible, even section, $`v`$ of $`𝒜`$, can be written as $`\pm e^{2w}`$ for an even section $`w`$ of $`𝒜`$. In fact, $`v`$ can be written as the product of a nowhere vanishing function on the base manifold and a function $`1+u`$, where $`u`$ is nilpotent. Since $`u`$ is nilpotent, say of order $`k`$, $`1+u`$ is equal to $`\mathrm{exp}(\mathrm{ln}(1+u))`$, where $`\mathrm{ln}(1+u)=_{j=1}^k(1)^{j1}\frac{u^j}{j}`$.
###### Proposition 2.4.
Let $`\xi `$ be a berezinian volume. For any invertible, even section $`v`$ of $`A`$, the berezinian section $`\xi .v`$ is a berezinian volume, and, for any graded vector field $`D`$,
(17)
$$\mathrm{div}_{\xi .v}(D)=\mathrm{div}_\xi (D)+v^1D(v).$$
If $`v=e^{2w}`$, where $`w`$ is an even section of $`𝒜`$, then
(18)
$$\mathrm{div}_{\xi .e^{2w}}(D)=\mathrm{div}_\xi (D)+D(2w).$$
Proof. Using (13), we compute
$`(\xi .v).\mathrm{div}_{\xi .v}(D)`$ $`=(1)^{|D||\xi |}_D(\xi .v)`$
$`=(1)^{|D||\xi |}_D(\xi ).v+\xi .D(v)`$
$`=\xi .(\mathrm{div}_\xi (D)v)+\xi .D(v).`$
Therefore
$$v\mathrm{div}_{\xi .v}(D)=\mathrm{div}_\xi (D)v+D(v),$$
and, multiplying both sides by $`v^1`$, we obtain formula (17), since $`v`$ is even. ∎
### 2.2. Properties of generators defined by berezinian volumes
We shall now assume that there is an odd Poisson structure, $`\pi `$, on $`(M,𝒜)`$, with odd Poisson bracket $`[[,]]`$. Let $`\xi `$ be a berezinian volume on $`(M,𝒜)`$. Following the general pattern of Section 1.3, we define the operator $`\mathrm{\Delta }^{\pi ,\xi }:𝒜𝒜`$ by
(19)
$$\mathrm{\Delta }^{\pi ,\xi }f=(1)^{|f|}\frac{1}{2}\mathrm{div}_\xi X_f^\pi ,$$
for any section $`f`$ of $`𝒜`$. It follows from Proposition 2.2 and Theorem 1.2 that the odd operator $`\mathrm{\Delta }^{\pi ,\xi }`$ is a generator of bracket $`\pi `$. Thus, given an odd Poisson bracket, to any berezinian volume there corresponds a generator of this bracket. We shall now study the effect on the generator of a change of berezinian volume, and determine under which conditions the generator corresponding to a berezinian volume is of square 0.
It follows from (18) and Proposition 1.5 that, when $`\xi `$ is a berezinian volume and $`v=e^{2w}`$ an invertible, even section of $`𝒜`$,
(20)
$$\mathrm{\Delta }^{\pi ,\xi .v}=\mathrm{\Delta }^{\pi ,\xi }+X_w^\pi .$$
It also follows from Proposition 1.6 that, if $`\xi `$ is a berezinian volume such that $`(\mathrm{\Delta }^{\pi ,\xi })^2=0`$, and $`v=e^{2w}`$ is an invertible, even section of $`𝒜`$, then,
(21)
$$(\mathrm{\Delta }^{\pi ,\xi .v})^2=X_{\mathrm{\Delta }^{\pi ,\xi }w+\frac{1}{2}[[w,w]]}^\pi .$$
Moreover,
$$e^w\mathrm{\Delta }^{\pi ,\xi }(e^w)=\frac{1}{2}e^w\mathrm{div}_\xi (e^wX_w^\pi )=\mathrm{\Delta }^{\pi ,\xi }w+\frac{1}{2}e^w[[w,e^w]],$$
and therefore
(22)
$$\mathrm{\Delta }^{\pi ,\xi }w+\frac{1}{2}[[w,w]]=e^w\mathrm{\Delta }^{\pi ,\xi }e^w.$$
In the context of supermanifolds (usually infinite-dimensional), the Maurer-Cartan equation,
(23)
$$\mathrm{\Delta }^{\pi ,\xi }w+\frac{1}{2}[[w,w]]=0,$$
is referred to as the Batalin-Vilkovisky quantum master equation. In the case of odd symplectic supermanifolds, the results stated below can be found in articles dealing with the BV-quantization of gauge theories, starting with , , followed by, among others, , , , , , or of string theories . See also and . In our treatment, the more general case of possibly degenerate odd Poisson structures is included.
###### Proposition 2.5.
Let $`\xi `$ be a berezinian volume on $`(M,𝒜)`$ such that $`(\mathrm{\Delta }^{\pi ,\xi })^2=0`$, and let $`v=e^{2w}`$ be an invertible, even section of $`𝒜`$.
(i) The following conditions are equivalent
* $`\mathrm{\Delta }^{\pi ,\xi }(e^w)=0`$,
* $`w`$ is a solution of (23).
(ii) If this condition is satisfied, then $`(\mathrm{\Delta }^{\pi ,\xi .v})^2=0`$.
Proof. These implications follow from Equations (22) and (20). ∎
Conversely $`(\mathrm{\Delta }^{\pi ,\xi .v})^2=0`$ implies that there exists an odd Casimir section $`C`$ of square 0 such that $`\mathrm{\Delta }^{\pi ,\xi }(e^w)C=0`$. (A Casimir section is a section of $`𝒜`$ such that $`X_C^\pi =0`$.) In fact, if $`(\mathrm{\Delta }^{\pi ,\xi .v})^2=0`$, then there exists a Casimir section $`C`$ such that $`e^w\mathrm{\Delta }^{\pi ,\xi }(e^w)=C`$, or $`\mathrm{\Delta }^{\pi ,\xi }(e^w)=e^wC`$. Together with $`(\mathrm{\Delta }^{\pi ,\xi })^2=0`$, this condition implies that $`\mathrm{\Delta }^{\pi ,\xi }(e^w)C=0`$, whence also $`C^2=0`$.
If the odd Poisson bracket $`\pi `$ is nondegenerate, any Casimir section is a constant, therefore even, and necessarily $`C=0`$. In this case, (i) and (ii) in the proposition are equivalent.
In quantum field theory, $`e^{\frac{i}{\mathrm{}}S}`$ is the action, and the condition $`\mathrm{\Delta }^{\pi ,\xi }(e^{\frac{i}{\mathrm{}}S})`$ states that the action is closed with respect to the differential $`\mathrm{\Delta }^{\pi ,\xi }`$.
###### Proposition 2.6.
If $`v=e^{2w}`$ is an invertible, even section of $`𝒜`$ such that $`w`$ is a solution of the Equation (23), then, for any section $`f`$ of $`𝒜`$,
$$\mathrm{\Delta }^{\pi ,\xi .v}f=e^w\mathrm{\Delta }^{\pi ,\xi }(e^wf).$$
Proof. Using (1), we find that
$$e^w\mathrm{\Delta }^{\pi ,\xi }(e^wf)=e^w\left([[e^w,f]]+(\mathrm{\Delta }^{\pi ,\xi }e^w)f+e^w(\mathrm{\Delta }^{\pi ,\xi }f)\right).$$
By (22) and (20),
$$e^w\mathrm{\Delta }^{\pi ,\xi }(e^wf)=\mathrm{\Delta }^{\pi ,\xi .v}(f)+\left(\mathrm{\Delta }^{\pi ,\xi }w+\frac{1}{2}[[w,w]]\right)f.$$
The result is proved in view of (20). ∎
### 2.3. The supermanifold $`\mathrm{\Pi }T^{}M`$
For any manifold $`M`$ of dimension $`n`$, we consider the supermanifold $`\mathrm{\Pi }T^{}M`$ of dimension $`n|n`$, whose structural sheaf is the sheaf of multivectors on $`M`$. The supermanifold $`\mathrm{\Pi }T^{}M`$ has an odd Poisson bracket, the Schouten bracket of multivectors on $`M`$. It is in fact nondegenerate, i.e., the odd Poisson structure on $`\mathrm{\Pi }T^{}M`$ is symplectic. Here we revert to the usual notations for vector fields, differential forms and functions on the ordinary manifold $`M`$. See the Appendix for the definition to be used below of the map $`\alpha \stackrel{~}{\alpha }`$ from forms on a supermanifold $`(M,𝒜)`$ to forms on $`M`$.
###### Lemma 2.7.
Given a volume form $`\mu `$ on $`M`$, there is a unique berezinian volume $`\xi _\mu `$ on $`\mathrm{\Pi }T^{}M`$ such that
$$\stackrel{~}{\xi _\mu (X)}=(i_{X_{(n)}}\mu )\mu ,$$
for any field of multivectors $`X`$, where $`i_{X_{(n)}}\mu `$ is the result of the duality-pairing of the homogeneous component $`X_{(n)}`$ of degree $`n`$ of the multivector $`X`$ and the $`n`$-form $`\mu `$.
Proof. To any graded vector field, $`D`$, on $`\mathrm{\Pi }T^{}M`$ is associated a vector field, $`\stackrel{~}{D}`$, on $`M`$, defined by
$$\stackrel{~}{D}(f)=\stackrel{~}{D(f)}$$
for any function $`f`$ on $`M`$. Given a differential $`n`$-form, $`\mu `$, on $`M`$, we can define a graded $`n`$-form, $`\mu ^G`$, on $`\mathrm{\Pi }T^{}M`$ by
$$<D_1,\mathrm{},D_n,\mu ^G>=\mu (\stackrel{~}{D_1},\mathrm{},\stackrel{~}{D_n}),$$
for graded vector fields $`D_1,\mathrm{},D_n`$. Then $`\stackrel{~}{\mu ^G}=\mu `$. Since the map $`Xi_{X_{(n)}}\mu `$ is a differential operator of order $`n`$ on the structural sheaf of $`\mathrm{\Pi }T^{}M`$, the map $`\xi _\mu :X(i_{X_{(n)}}\mu )\mu ^G`$ defines a section of the berezinian sheaf, which is a berezinian volume if and only if $`\mu `$ is a volume form. ∎
We remark that, for any positive function $`v`$ on $`M`$,
$$\xi _{v\mu }=(\xi _\mu ).v^2.$$
We assume that $`M`$ is an orientable manifold, and we let $`\mu `$ be a volume form on $`M`$. In the non orientable case, densities must be used instead of volume forms. On $`\mathrm{\Pi }T^{}M`$, there is an operator $`\mathrm{\Delta }^{Schouten,\xi _\mu }`$ associated to the odd Poisson bracket and to the berezinian volume $`\xi _\mu `$ by means of (19). Let $`\mathrm{d}`$ be the de Rham differential on $`M`$, and let $`_\mu `$ be the isomorphism from multivectors to forms defined by the volume form $`\mu `$. Then we know (see, e.g., ) that $`_\mu =_\mu ^1\mathrm{d}_\mu `$ is a generator of the Schouten bracket.
###### Theorem 2.8.
For any volume form $`\mu `$ on $`M`$, the generator $`\mathrm{\Delta }^{Schouten,\xi _\mu }`$ of the Schouten bracket coincides with $`_\mu =_\mu ^1\mathrm{d}_\mu ,`$ and $`(\mathrm{\Delta }^{Schouten,\xi _\mu })^2=0`$.
Proof. It is enough to show that these operators coincide on vector fields. We prove this fact using local coordinates $`(x^1,\mathrm{},x^n,\xi _1,\mathrm{},\xi _n)`$ on $`\mathrm{\Pi }T^{}M`$. For a vector field $`X=_{i=1}^nX^i\frac{}{x^i}`$ considered as a function $`_{i=1}^nX^i\xi _i`$ on $`\mathrm{\Pi }T^{}M`$,
$$[[X,.]]=\underset{i=1}{\overset{n}{}}X^i\frac{}{x^i}\underset{i=1}{\overset{n}{}}\underset{j=1}{\overset{n}{}}\frac{X^i}{x^j}\xi _i\frac{}{\xi _j}.$$
Assume that $`\mu _0=dx^1\mathrm{}dx^n`$, so that
$$\frac{1}{2}\mathrm{div}_{\xi _{\mu _0}}[[X,.]]=\frac{1}{2}\underset{i=1}{\overset{n}{}}(\frac{X^i}{x^i}+\frac{X^i}{x^i})=\mathrm{div}_{\mu _0}X=_{\mu _0}X.$$
More generally, if $`\mu =e^w\mu _0`$, then $`\xi _\mu =\xi _{\mu _0}.e^{2w}`$, and for any vector field $`X`$,
$$\frac{1}{2}\mathrm{div}_{\xi _\mu }[[X,.]]=\frac{1}{2}\mathrm{div}_{\mu _0}XX.w=_{\mu _0}XX.w=_{e^w\mu _0}X=_\mu X.$$
The fact that $`(\mathrm{\Delta }^{Schouten,\xi _\mu })^2=0`$ is now an immediate consequence of the fact that $`\mathrm{d}^2=0`$. ∎
Remark. The equality $`\mathrm{\Delta }=_\mu `$ means that, for any differential form $`\alpha `$,
$$_\mu ^1d\alpha =\mathrm{\Delta }_\mu ^1\alpha .$$
The map $`_\mu ^1`$ coincides with the “Fourier transform with respect to the odd variables” introduced in and , p. 255. Therefore, in the case of a nondegenerate Poisson structure, our result reduces to that of and , formula (20), a fact already observed by Witten in , formula (13). In the terminology of Voronov and Schwarz, the operator $`_\mu `$ on the functions on $`\mathrm{\Pi }T^{}M`$ is the “Fourier transform” of the de Rham differential acting on functions on $`\mathrm{\Pi }TM`$. Schwarz proves that a supermanifold $`(M,𝒜)`$ of dimension $`n|n`$ with an odd symplectic structure is equivalent in a suitable sense to $`\mathrm{\Pi }T^{}M`$ with its canonical, odd symplectic structure.
Now, let $`\mu `$ be a volume form, and $`\omega `$ a differential form on $`M`$ such that $`\omega _{(n)}`$ is a volume form. If $`Q=_\mu ^1\omega `$, then $`\xi ^{}=\xi _\mu .Q^2`$ is a berezinian volume on $`\mathrm{\Pi }T^{}M`$. Setting $`\mathrm{\Delta }^{}=\mathrm{\Delta }^{Schouten,\xi ^{}}`$, we see from Proposition 2.5 that $`(\mathrm{\Delta }^{})^2=0`$ if $`\mathrm{\Delta }Q=0`$. By Theorem 2.8, this condition is equivalent to $`\mathrm{d}\omega =0`$. So $`(\mathrm{\Delta }^{})^2=0`$ if $`\omega `$ is a closed form. This result constitutes part of Theorem 5 of . Moreover, it is proved there that two closed forms in the same de Rham cohomology class yield equivalent structures.
### 2.4. The supermanifold $`\mathrm{\Pi }TM`$
We shall now consider another supermanifold attached to a smooth manifold $`M`$ of dimension $`n`$. Let $`𝒜=\mathrm{\Omega }(M)`$ be the sheaf of differential forms on $`M`$. The pair $`(M,𝒜)`$, usually denoted by $`\mathrm{\Pi }TM`$, is a supermanifold of dimension $`n|n`$. The sections of the sheaf $`\mathrm{\Omega }(M)`$ will be denoted by $`\alpha ,\beta \mathrm{}`$, but differential forms of $``$-degree $`0`$, i.e., smooth functions, will be denoted by, $`f,g,\mathrm{}`$. If $`\alpha `$ is a section of $`\mathrm{\Omega }(M)`$, then we denote the homogeneous component of $`\alpha `$ of degree $`n`$ by $`\alpha _{(n)}`$.
#### 2.4.1. Canonical berezinian volume on the supermanifold $`\mathrm{\Pi }TM`$
###### Lemma 2.9.
There is a unique berezinian volume $`\xi `$ on $`\mathrm{\Pi }TM`$, such that, for any section $`\alpha `$ of $`\mathrm{\Omega }(M)`$,
(24)
$$\stackrel{~}{\xi (\alpha )}=\alpha _{(n)}.$$
Proof. If $`\xi `$ and $`\xi ^{}`$ are berezinian volumes satisfying (24), then $`\stackrel{~}{(\xi \xi ^{})(\alpha )}=0`$ for any section $`\alpha `$ of $`\mathrm{\Omega }(M)`$, and this means, by the definition of the berezinian sheaf, that $`\xi =\xi ^{}`$, which proves the uniqueness of a berezinian volume satisfying (24). To prove its existence, we use local coordinates and we show the invariance under a change of coordinates. Let $`(x^1,\mathrm{},x^n)`$ be local coordinates on an open set $`U`$ of the manifold $`M`$. Then $`(x^1,\mathrm{},x^n,s^1=\mathrm{d}x^1,\mathrm{},s^n=\mathrm{d}x^n)`$ are graded local coordinates in $`\mathrm{\Pi }TM`$, and a local basis of derivations is
$$(_{\frac{}{x^1}},\mathrm{},_{\frac{}{x^n}},\frac{}{s^1}=i_{\frac{}{x^1}},\mathrm{},\frac{}{s^n}=i_{\frac{}{x^n}}).$$
We now consider the local section of the berezinian sheaf,
$$\xi _U=[\mathrm{d}^Gx^1\mathrm{}\mathrm{d}^Gx^ni_{\frac{}{x^1}\mathrm{}\frac{}{x^n}}].$$
A change of coordinates from $`(x^1,\mathrm{},x^n)`$ to $`(y^1,\mathrm{},y^n)`$ induces a change of graded coordinates to $`(y^1,\mathrm{},y^n,t^1=\mathrm{d}y^1,\mathrm{},t^n=\mathrm{d}y^n)`$, with matrix
$$\left(\begin{array}{cc}\left(\frac{y^i}{x^k}\right)& 0\\ 0& \left(\frac{y^i}{x^k}\right)\end{array}\right),$$
whose berezinian is equal to $`1`$. Therefore, we can define a berezinian section $`\xi `$ on $`\mathrm{\Pi }TM`$ by piecing the locally defined $`\xi _U`$’s together.
This berezinian section is also a berezinian volume. In fact, this is true locally and, if $`U`$ and $`V`$ are open sets such that $`UV\mathrm{}`$, and if local forms $`\alpha `$ and $`\beta `$ satisfy $`\xi _U.\alpha =\xi _V.\beta `$, then $`\alpha =\beta `$, on $`UV`$.
Finally, if $`\alpha _{(n)}=f\mathrm{d}x^1\mathrm{}\mathrm{d}x^n`$, where $`fC^{\mathrm{}}(U)`$, then $`\xi _U(\alpha )=f\mathrm{d}^Gx^1\mathrm{}\mathrm{d}^Gx^n`$, and therefore
$$\stackrel{~}{\xi _U(\alpha )}=\stackrel{~}{f\mathrm{d}^Gx^1\mathrm{}\mathrm{d}^Gx^n}=f\mathrm{d}x^1\mathrm{}\mathrm{d}x^n=\alpha _{(n)},$$
and $`\xi `$ satisfies condition (24). ∎
#### 2.4.2. The canonical divergence operator on $`\mathrm{\Pi }TM`$
Let us denote by $`\mathrm{div}_{can}`$ the divergence operator associated to the canonical berezinian volume on $`\mathrm{\Pi }TM`$. The graded vector fields on $`\mathrm{\Pi }TM`$ are the derivations of the sheaf of differential forms. Let us denote the sheaf of vector-valued differential $`k`$-forms by $`\mathrm{\Omega }^k(M;TM)`$, for $`k0`$. By the Frölicher-Nijenhuis theorem , we know that a derivation $`D`$ of degree $`k`$ of $`\mathrm{\Omega }(M)`$ can be uniquely written as
$$D=_K+i_L,$$
where $`K`$ is a section of $`\mathrm{\Omega }^k(M;TM)`$ and $`L`$ is a section of $`\mathrm{\Omega }^{k+1}(M;TM)`$.
We introduce the notation $`𝒞`$ for the $`(1,1)`$-contraction map from $`\mathrm{\Omega }^k(M;TM)`$ to $`\mathrm{\Omega }^{k1}(M)`$, defined on a decomposable element $`K=\omega X`$, where $`X`$ is a vector field and $`\omega `$ is a $`k`$-form, by $`𝒞K=i_X\omega `$, for $`k1`$, and by $`𝒞=0`$ on $`\mathrm{\Omega }^0(M;TM)`$.
###### Lemma 2.10.
For $`K`$ a section of $`\mathrm{\Omega }^k(M;TM)`$ and $`L`$ a section of $`\mathrm{\Omega }^{k+1}(M;TM)`$,
$$\mathrm{div}_{can}(i_L)=(1)^k𝒞L,\mathrm{div}_{can}(\mathrm{d})=0,\mathrm{div}_{can}(_K)=\mathrm{d}(𝒞K),$$
where $`\mathrm{d}`$ denotes the de Rham differential.
Proof. We shall first compute $`\mathrm{div}_{can}(i_L)`$ for a decomposable $`L=\omega X`$ where $`\omega `$ is a section of $`\mathrm{\Omega }^{k+1}(M)`$. For any differential form $`\alpha `$,
$`\stackrel{~}{(_{i_L}\xi )\alpha }`$ $`=\stackrel{~}{(\xi i_L)\alpha }=\stackrel{~}{\xi (i_L\alpha )}`$
$`=\stackrel{~}{\xi (\omega i_X\alpha )}=(\omega i_X\alpha )_{(n)}=\omega i_X(\alpha _{(nk)})`$
$`=(1)^ki_X(\omega )\alpha _{(nk)}=(1)^k(i_X(\omega )\alpha )_{(n)}`$
$`=\stackrel{~}{\xi ((1)^ki_X(\omega )\alpha )},`$
where we have used the relation
$$(i_X\omega )\alpha _{(nk)}+(1)^{k+1}\omega i_X\alpha _{(nk)}=i_X(\omega \alpha _{(nk)})=0.$$
Therefore $`_{i_L}\xi =\xi .((1)^ki_X(\omega ))=\xi .((1)^k𝒞L)`$, and $`\mathrm{div}_{can}(i_L)=(1)^k𝒞L`$.
Similarly, for the derivation $`\mathrm{d}`$,
$$\stackrel{~}{(_\mathrm{d}\xi )\alpha }=\stackrel{~}{(\xi \mathrm{d})\alpha }=\stackrel{~}{\xi (\mathrm{d}\alpha )}=\mathrm{d}(\alpha _{(n1)}),$$
which is always an exact form. Thus $`_\mathrm{d}\xi =0`$. Finally, using (16), we obtain
$$\mathrm{div}_\xi (_K)=\mathrm{div}_\xi ([i_K,\mathrm{d}])=i_K(\mathrm{div}_\xi (\mathrm{d}))(1)^{k1}\mathrm{d}(\mathrm{div}(i_K))=\mathrm{d}(𝒞K).\mathit{}$$
In particular, if $`X`$ is a vector field on $`M`$, then
(25)
$$\mathrm{div}_{can}(i_X)=0,\mathrm{div}_{can}(_X)=0.$$
#### 2.4.3. Generators of the Koszul-Schouten bracket
We shall now assume that the base manifold $`M`$ is equipped with a Poisson structure. Given a Poisson manifold $`(M,P)`$, there is an odd Poisson bracket, $`[[,]]_P`$, on the supermanifold $`\mathrm{\Pi }TM`$, called the Koszul-Schouten bracket, that is characterized by the conditions,
$$[[f,g]]_P=0,[[f,\mathrm{d}g]]_P=\{f,g\},[[\mathrm{d}f,\mathrm{d}g]]_P=\mathrm{d}\{f,g\},$$
for all $`f,gC^{\mathrm{}}(M)`$, where $`\{,\}`$ denotes the Poisson bracket on $`C^{\mathrm{}}(M)`$ defined by $`P`$, together with the graded Leibniz rule. It was shown by Koszul that a generator for this bracket is the Poisson homology operator, $`_P=[\mathrm{d},i_P]`$, sometimes called the Koszul-Brylinski operator. See , and also . On the other hand, we know that, given a berezinian volume $`\xi `$ on $`\mathrm{\Pi }TM`$, the operator $`\mathrm{\Delta }^{(P),\xi }`$ defined by
$$\mathrm{\Delta }^{(P),\xi }(\alpha )=(1)^{|\alpha |}\frac{1}{2}\mathrm{div}_\xi ([[\alpha ,.]]_P),$$
for any differential form $`\alpha `$, is also a generator of $`[[,]]_P`$.
###### Theorem 2.11.
The generator $`\mathrm{\Delta }^{(P),can}`$ of the Koszul-Schouten bracket associated to the canonical berezinian volume coincides with $`_P`$, and $`(\mathrm{\Delta }^{(P),can})^2=0`$.
Proof. It suffices to prove that $`\mathrm{\Delta }^{(P),can}`$ and $`_P`$ agree on $`1`$-forms, and it is enough to show that both vanish on exact $`1`$-forms, $`\alpha =\mathrm{d}f`$, where $`fC^{\mathrm{}}(M)`$. In fact, $`[[\mathrm{d}f,.]]_P=_{\mathrm{\#}_P\mathrm{d}f}`$, where $`\mathrm{\#}_P\mathrm{d}f=\{f,.\}.`$ From (25), it follows that $`\mathrm{\Delta }^{(P),can}(\mathrm{d}f)=\frac{1}{2}\mathrm{div}_{can}(_{\mathrm{\#}_P\mathrm{d}f})=0.`$ And clearly $`_P(\mathrm{d}f)=0.`$ Moreover, $`(_P)^2=0`$, since $`[\mathrm{d},_P]=0`$ and $`[i_P,_P]=0`$, and therefore $`(\mathrm{\Delta }^{(P),can})^2=0`$. ∎
Remark. Any nondegenerate metric $`g`$ on the manifold $`M`$ defines an isomorphism from multivectors to differential forms. Hence, from the Schouten bracket of multivectors, we obtain a $``$-graded bracket $`[[,]]_g`$ on the sheaf of differential forms on $`M`$. Then, the codifferential $`\delta _g`$ associated to $`g`$ is a generator of this bracket. (See or .) One can also consider the operator associated to the canonical berezinian on $`\mathrm{\Pi }TM`$, defined by
$$\mathrm{\Delta }^{(g),can}(\alpha )=(1)^{|\alpha |}\frac{1}{2}\mathrm{div}_{can}([[\alpha ,.]]_g),$$
for any differential form $`\alpha `$, and one can show that these two generators of the bracket $`[[,]]_g`$ coincide, $`\mathrm{\Delta }^{(g),can}=\delta _g`$.
### 2.5. QS, SP and QSP-manifolds
The following definitions, adapted from and , will be useful in order to reformulate some of our results.
###### Definition 2.12.
Let $`(M,𝒜)`$ be a supermanifold, $`D`$ an odd vector field and $`\xi `$ a berezinian volume. We say that $`((M,𝒜),D,\xi )`$ is a QS-manifold if $`D^2=0`$ and $`\mathrm{div}_\xi (D)=0`$.
###### Definition 2.13.
Let $`(M,𝒜)`$ be a supermanifold, $`\pi `$ an odd Poisson bracket and $`\xi `$ a berezinian volume. We say that $`((M,𝒜),\pi ,\xi )`$ is a weak SP-manifold if $`(\mathrm{\Delta }^{\pi ,\xi })^2=0`$, where $`\mathrm{\Delta }^{\pi ,\xi }`$ is defined by (19). If the Poisson bracket is nondegenerate, then the supermanifold is an SP-manifold.
###### Definition 2.14.
Let $`(M,𝒜)`$ be a supermanifold, $`\pi `$ an odd Poisson bracket, $`D`$ an odd vector field and $`\xi `$ a berezinian volume. We say that $`((M,𝒜),\pi ,D,\xi )`$ is a weak QSP-manifold if
* $`((M,𝒜),\pi ,\xi )`$ is a weak SP-manifold,
* $`((M,𝒜),D,\xi )`$ is a QS-manifold, and
* $`D`$ is a derivation of the odd Poisson bracket, $`\pi `$.
If the Poisson bracket is nondegenerate, and if $`D`$ is the hamiltonian vector field defined by an even section of $`𝒜`$, then the supermanifold is a QSP-manifold.
If the graded vector field $`D`$ is the hamiltonian vector field $`[[h,.]]`$, where $`h`$ is an even section of $`𝒜`$, then $`D^2=0`$ if and only if $`[[h,h]]=0.`$ In field theory, this condition appears under the name classical master equation.
It follows from identity (3) that, if $`\mathrm{\Delta }`$ is a generator of the odd Poisson bracket $`[[,]]`$ and if $`D`$ is a graded vector field on $`(M,𝒜)`$, a necessary and sufficient condition for $`D`$ to be a derivation of the odd Poisson bracket is that the graded commutator, $`[D,\mathrm{\Delta }]`$, be a derivation of the associative multiplication of $`𝒜`$. This implies
###### Proposition 2.15.
If $`\pi `$ and $`\xi `$ define a weak SP-structure on $`(M,𝒜)`$, if $`D`$ and $`\xi `$ define a QS-structure on $`(M,𝒜)`$, and if $`[D,\mathrm{\Delta }]=0`$, then $`\pi `$, D and $`\xi `$ define a weak QSP-structure on $`(M,𝒜)`$.
The following theorems follow in part from the results of Section 2.3 and 2.4. See also .
###### Theorem 2.16.
(i) For any manifold $`M`$ with a volume element, the supermanifold $`\mathrm{\Pi }T^{}M`$, with the Schouten bracket and the berezinian volume $`\xi _\mu `$, is an SP-manifold.
(ii) Let $`(M,P)`$ be a Poisson manifold, and let $`d_P=[[P,.]]`$ be the Lichnerowicz-Poisson differential. Then $`\mathrm{\Pi }T^{}M`$, with the Schouten bracket, the odd vector field $`d_P`$, and the canonical berezinian volume, is a QSP-manifold.
Proof. In fact, the odd vector field $`d_P`$ is of square $`0`$, because $`[[P,P]]=0`$, and it is a derivation of the Schouten bracket by the graded Jacobi identity. ∎
###### Theorem 2.17.
(i) For any manifold $`M`$, the supermanifold $`\mathrm{\Pi }TM`$, with the de Rham differential $`\mathrm{d}`$ and the canonical berezinian volume, is a QS-manifold.
(ii) Let $`(M,P)`$ be a Poisson manifold, and let $`[[,]]_P`$ be the Koszul-Schouten bracket. Then $`\mathrm{\Pi }TM`$ with the odd Poisson bracket $`[[,]]_P`$, the de Rham differential $`\mathrm{d}`$ and the canonical berezinian volume, is a weak QSP-manifold. If $`P`$ is a nondegenerate Poisson structure, then $`\mathrm{\Pi }TM`$ is a QSP-manifold.
Proof. To complete the proof, we only have to recall that the de Rham differential $`\mathrm{d}`$ is a derivation of $`[[,]]_P,`$ and that, when $`P`$ is nondegenerate with inverse the symplectic form, $`\omega `$, then $`\mathrm{d}=[[\omega ,.]]_P,`$ so that $`\mathrm{d}`$ is the hamiltonian vector field associated to $`\omega `$. (See or for a proof of this fact.)∎
## 3. Linear connections and generators of odd Poisson brackets
Divergence operators on smooth manifolds can be defined not only by means of volume forms, but also by means of connections. (See .) While in Section 2, we generalized the first approach to supermanifolds, replacing volume forms by their graded analogue, the berezinian volumes, in this section we generalize the second method, defining divergence operators by means of graded connections.
### 3.1. Divergence operators defined by graded connections
We first recall the notion of graded connection. See and for the definitions of left and right graded connections. Here we consider only left graded connections, which we simply call connections. Let $`(M,𝒜)`$ be a supermanifold and let $`\mathrm{Der}𝒜`$ be the sheaf of derivations of $`𝒜`$.
###### Definition 3.1.
Let $`𝒮`$ be a sheaf of $`𝒜`$-modules on $`M`$. A left graded connection, or simply a connection, on $`𝒮`$ is a morphism of sheaves of graded vector spaces from $`\mathrm{Der}𝒜𝒮`$ to $`𝒮`$, denoted $`D\alpha _D\alpha ,`$ which satisfies the identity
$$_{fD}\alpha =f_D\alpha ,$$
and the Leibniz rule,
$$_D(f\alpha )=D(f)\alpha +(1)^{|D||f|}f_D\alpha ,$$
for any section $`f`$ of $`𝒜`$, any derivation $`D`$ of $`𝒜`$, and any section $`\alpha `$ of $`𝒮`$.
A connection on the sheaf $`\mathrm{Der}𝒜`$ of derivations of $`𝒜`$ is a graded linear connection or simply a linear connection on $`(M,𝒜)`$.
###### Definition 3.2.
The curvature, $`R^{}`$, of a connection, $``$, on $`𝒮`$ is defined by
$$R^{}(D_1,D_2)=[_{D_1},_{D_2}]_{[D_1,D_2]},$$
for any derivations $`D_1,D_2`$ of $`𝒜`$.
The torsion, $`T^{}`$, of a linear connection, $``$, on $`(M,𝒜)`$ is defined by
$$T^{}(D_1,D_2)=_{D_1}D_2(1)^{|D_1||D_2|}_{D_2}D_1[D_1,D_2].$$
We shall now define the divergence operator associated with a linear connection on $`(M,𝒜)`$. Let $`sTr`$ denote the supertrace of an endomorphism of sheaves of $`𝒜`$-modules (see, e.g., or ), and let $`\mathrm{ad}_D`$ denote the endomorphism of $`\mathrm{Der}𝒜`$, $`E[D,E]`$. For any graded vector field $`D`$, we set
(26)
$$\mathrm{div}_{}(D)=sTr\left(_D\mathrm{ad}_D\right).$$
###### Proposition 3.3.
For any linear connection, $``$, on $`(M,𝒜)`$, the map, $`\mathrm{div}_{}:\mathrm{Der}𝒜𝒜,`$ defined by (26) is a divergence operator.
Proof. The map $`\mathrm{div}_{}`$ is even. It follows from $`[fD,E]=f[D,E](1)^{(|f|+|D|)|E|}E(f)D`$, that
$`\mathrm{div}_{}(fD)`$ $`=sTr(_{fD}\mathrm{ad}_{fD})`$
$`=f\mathrm{div}_{}(D)+sTr\left(E(1)^{(|f|+|D|)|E|}E(f)D\right).`$
Let the graded dimension of the supermanifold be $`m|n`$, and let us choose a system of local graded coordinates $`(x^1,\mathrm{},x^m,s^1,\mathrm{},s^n)`$. We find that
$`sTr(E`$ $`(1)^{(|f|+|D|)|E|}E(f)D)={\displaystyle }_{i=1}^m{\displaystyle \frac{f}{x^i}}D(x^i)(1)^{|f|+|D|}{\displaystyle }_{\rho =1}^n{\displaystyle \frac{f}{s^\rho }}D(s^\rho )`$
$`=(1)^{|f||D|}{\displaystyle \underset{i=1}{\overset{m}{}}}D(x^i){\displaystyle \frac{f}{x^i}}+(1)^{|f||D|}{\displaystyle \underset{\rho =1}{\overset{n}{}}}D(s^\rho ){\displaystyle \frac{f}{s^\rho }}=(1)^{|f||D|}D(f),`$
where we have used the local expression of the derivation $`D`$ in the basis of local graded vector fields, $`(\frac{}{x^1},\mathrm{},\frac{}{x^m},\frac{}{s^1},\mathrm{},\frac{}{s^n})`$. ∎
###### Proposition 3.4.
Let $``$ be a torsionless linear connection on $`𝒜`$ and let $`D_1`$ and $`D_2`$ be graded vector fields. Then
(27)
$$^{\mathrm{div}_{}}(D_1,D_2)=sTr(R^{}(D_1,D_2)).$$
Proof. We have to prove that
$`D_1(sTr(_{D_2}\mathrm{ad}_{D_2}))`$ $`(1)^{|D_1||D_2|}D_2(sTr(_{D_1}\mathrm{ad}_{D_1}))sTr(_{[D_1,D_2]}\mathrm{ad}_{[D_1,D_2]})`$
$`=sTr([_{D_1},_{D_2}]_{[D_1,D_2]}).`$
This result follows from a computation of these two expressions in local coordinates, for pairs of commuting graded vector fields, $`D_1`$ and $`D_2`$, in a local basis. ∎
### 3.2. Generators defined by graded connections
We now assume that $`(M,𝒜)`$ has an odd Poisson structure, $`\pi `$, whose bracket we denote by $`[[,]]`$. Let $``$ be a linear connection on $`(M,𝒜)`$. Following the general pattern of Section 1.3, we define the operator $`\mathrm{\Delta }^{\pi ,}:𝒜𝒜`$ by
(28)
$$\mathrm{\Delta }^{\pi ,}(f)=(1)^{|f|}\frac{1}{2}\mathrm{div}_{}([[f,.]]),$$
for any section $`f`$ of $`𝒜`$. It follows from Proposition 3.3 and Theorem 1.2 that the odd operator $`\mathrm{\Delta }^{\pi ,}`$ is a generator of bracket $`\pi `$. Therefore, to any linear connection on an odd Poisson manifold, there corresponds a generator of the odd Poisson bracket.
###### Proposition 3.5.
Let $``$ be a torsionless linear connection on $`𝒜`$. The following properties are equivalent
* $`\mathrm{\Delta }^{\pi ,}`$ is a derivation of the odd Poisson bracket,$`\pi `$,
* $`(\mathrm{\Delta }^{\pi ,})^2`$ is a derivation of the sheaf of associative algebras, $`𝒜`$,
* $`sTr(R^{})`$ vanishes on the sheaf of hamiltonian derivations.
Proof. These equivalences follow from Lemma 1.1, and from Corollary 1.4 together with Proposition 3.4. ∎
We now compare the generators associated to torsionless linear connections, $``$ and $`^{}`$, on $`𝒜`$. The difference $`_D^{}_D`$ is then a morphism of sheaves of $`𝒜`$-modules from $`\mathrm{Der}𝒜`$ to itself, which we denote by $`u(D)`$.
###### Proposition 3.6.
Let $``$ and $`^{}`$ be torsionless linear connections on $`𝒜`$. Then
$$\mathrm{\Delta }^{\pi ,^{}}(f)=\mathrm{\Delta }^{\pi ,}(f)+(1)^{|f|}\frac{1}{2}sTr(u(X_f^\pi )),$$
for any section $`f`$ of $`𝒜`$ .
Proof. This relation follows from the fact that, for any derivation $`D`$ of $`𝒜`$, $`\mathrm{div}_{^{}}(D)=\mathrm{div}_{}(D)+sTr(u(D)).`$
Remark. In the case of an ordinary manifold, the trace of the curvature of a linear connection is the curvature of the connection induced on the bundle of top-degree forms. It would be interesting to interpret the supertrace of the curvature of a graded linear connection as the curvature of a connection on the sections of the berezinian sheaf.
### 3.3. Metrics and metric connections on supermanifolds
We recall the definitions of metrics and metric linear connections on supermanifolds.
###### Definition 3.7.
A graded metric, or simply a metric, on $`(M,𝒜)`$ is a morphism of sheaves of $`𝒜`$-modules, $`,:\mathrm{Der}𝒜\mathrm{Der}𝒜𝒜`$, such that
* $`D_1,D_2=(1)^{|D_1||D_2|}D_2,D_1,`$ for derivations $`D_1`$ and $`D_2`$ (graded symmetry),
* the map $`DD,`$ is an isomorphism of sheaves of $`𝒜`$-modules from $`\mathrm{Der}𝒜`$ to $`\mathrm{Hom}_𝒜(\mathrm{Der}𝒜,𝒜)`$ (nondegeneracy).
###### Definition 3.8.
A linear connection $``$ on $`𝒜`$ is metric with respect to a metric $`,`$ if, for any derivations $`D,D_1`$, and $`D_2`$ of $`𝒜`$,
$$DD_1,D_2=_DD_1,D_2+(1)^{|D_1||D|}D_1,_D^0D_2+(1)^{|D_1|(|D|+1)}D_1,_D^1D_2,$$
where $`=^0+^1`$ is the decomposition of the linear connection into its even and odd components.
The proof of the following theorem, can be found in , p. 134, and in .
###### Theorem 3.9.
There exists a unique torsionless linear connection which is metric with respect to a given metric. It is determined by
$$\begin{array}{ccc}\hfill 2_{D_1}D_2,D_3& =& D_1D_2,D_3+[D_1,D_2],D_3\hfill \\ & +& (1)^{|D_1|(|D_2|+|D_3|)}\left(D_2D_3,D_1[D_2,D_3],D_1\right)\hfill \\ & & (1)^{|D_3|(|D_1|+|D_2|)}\left(D_3D_1,D_2[D_3,D_1],D_2\right).\hfill \end{array}$$
The linear connection defined in Theorem 3.9 is called the graded Levi-Civita connection or simply, the Levi-Civita connection of the metric $`,`$. The Levi-Civita connection of a homogeneous metric is even. (See also , where an expression in local coordinates of the Levi-Civita connection is given in the case of a homogeneous, even metric.)
### 3.4. Linear connections and Schouten bracket
We shall again consider the supermanifold $`\mathrm{\Pi }T^{}M`$, whose sheaf of functions is the sheaf of multivectors on $`M`$. We shall use the notations $`f,g\mathrm{}`$ for functions on $`M`$ or on an open set of $`M`$, and the notations $`X,Y,\mathrm{}`$ for vector fields and $`\alpha ,\beta \mathrm{}`$ for differential $`1`$-forms.
#### 3.4.1. Graded vector fields on $`\mathrm{\Pi }T^{}M`$
Whereas the derivations of the algebra of forms on a manifold can be classified by the Frölicher-Nijenhuis theorem (and see Section 2.4.2), the classification of the derivations of the algebra of multivectors on $`M`$ requires the use of an auxiliary linear connection, $``$. Let $`U`$ be an open set of $`M`$. If $`K=QX`$ is a vector-valued multivector on $`U`$, where $`Q`$ is a multivector and $`X`$ is a vector, we define $`_KV=Q_XV`$, for any multivector $`V`$ on $`U`$. If $`L=W\alpha `$ is a $`1`$-form-valued multivector on $`U`$, where $`W`$ is a multivector and $`\alpha `$ is a differential $`1`$-form, we define $`i_LV=Wi_\alpha V`$.
###### Proposition 3.10.
Let $`D`$ be a graded vector field of degree $`r`$ on $`\mathrm{\Pi }T^{}M`$, i.e., a derivation of degree $`r`$ of the sheaf of multivectors on $`M`$. On any open set $`U`$ of $`M`$, there exist a vector-valued $`r`$-vector, $`K`$, and a $`1`$-form-valued $`(r+1)`$-vector, $`L`$, each uniquely defined, such that $`D|_U=_K+i_L.`$
As a consequence, we see that, if $`(e_1,\mathrm{},e_n)`$ is a local basis of vector fields on $`U`$ and $`(ϵ^1,\mathrm{},ϵ^n)`$ is the dual basis, then $`(_{e_1},\mathrm{},_{e_n},i_{ϵ^1},\mathrm{},i_{ϵ^n})`$ generate the derivations of the algebra of multivectors over $`U`$, as a module over the algebra of multivectors over $`U`$.
#### 3.4.2. The graded connection on $`\mathrm{\Pi }T^{}M`$ associated to a linear connection on $`M`$
We shall show how to associate a metric on $`\mathrm{\Pi }T^{}M`$ to a linear connection on $`M`$, and we shall study the Levi-Civita connection of this metric.
###### Definition 3.11.
Let $``$ be a linear connection on $`M`$. We define a metric $`,_{}`$ on $`\mathrm{\Pi }T^{}M`$ by its value on derivations of type $`_X`$, where $`X`$ is a vector field, and of type $`i_\alpha `$, where $`\alpha `$ is a $`1`$-form,
$$_X,_Y_{}=0,_X,i_\alpha _{}=\alpha (X),i_\alpha ,i_\beta _{}=0.$$
To verify the nondegeneracy of the metric thus defined, we observe that, in the local basis of derivations $`(_{e_1},\mathrm{},_{e_n},i_{ϵ^1},\mathrm{},i_{ϵ^n})`$, the matrix of this metric is $`\left(\begin{array}{cc}0& Id\\ Id& 0\end{array}\right)`$. This metric is odd.
###### Proposition 3.12.
Let $``$ be a torsionless linear connection on $`M`$. The Levi-Civita connection, $``$, of the metric $`,_{}`$ on $`\mathrm{\Pi }T^{}M`$ is given by
$$__X_Y=_{_XY}+i_{R(,Y)X},__Xi_\alpha =i_{_X\alpha },_{i_\alpha }=0,$$
where $`R`$ denotes the curvature tensor of $``$.
Proof. We shall make use of the commutation relations
$$[_X,_Y]=_{[X,Y]}+i_{R(X,Y)},[_X,i_\alpha ]=i_{_X\alpha },\text{ and }[i_\alpha ,i_\beta ]=0,$$
Using Theorem 3.9, Definition 3.11 and the fact that the connection $``$ is torsionless, we obtain
$$\begin{array}{ccc}\hfill __X_Y,_Z_{}& =& \frac{1}{2}(R(X,Y)ZR(Y,Z)X+R(Z,X)Y)\hfill \\ & =& R(Z,Y)X=i_{R(.,Y)X},_Z_{},\hfill \\ & & \\ \hfill __X_Y,i_\alpha _{}& =& \alpha (_XY)=_{_XY},i_\alpha _{},\hfill \end{array}$$
and
$$__Xi_\alpha ,_Y_{}=(_X\alpha )Y=i_{_X\alpha },_Y_{},__Xi_\alpha ,i_\beta _{}=0.$$
From these relations and the nondegeneracy of the graded metric we obtain the first two formulæ, while the third follows from the fact that $``$ is torsionless.∎
###### Proposition 3.13.
The curvature $`R^{}`$ of the Levi-Civita connection, $``$, of the metric $`,_{}`$ on $`\mathrm{\Pi }T^{}M`$, satisfies
$`R^{}(_X,_Y)_Z`$ $`=_{R(X,Y)Z}+i_{(_XR)(,Z)Y}+i_{(_YR)(,Z)X},`$
$`R^{}(_X,_Y)i_\alpha `$ $`=i_{R(X,Y)^{}\alpha },`$
where $`R`$ denotes the curvature tensor of $``$, and $`R(X,Y)^{}`$ denotes the transpose of $`R(X,Y)`$. Moreover $`R^{}(_X,i_\alpha )=R^{}(i_\alpha ,i_\beta )=0`$.
Proof. The proof is a straightforward computation using Proposition 3.12.∎
###### Corollary 3.14.
Let $``$ be a torsionless linear connection on $`M`$. Then the Levi-Civita connection of the metric $`,_{}`$ on $`\mathrm{\Pi }T^{}M`$ is flat if and only if $``$ is flat.
#### 3.4.3. Generators of the Schouten bracket
We have just seen that, to a torsionless linear connection $``$ on $`M`$ we can associate the Levi-Civita connection, $``$, of the odd metric $`,_{}`$ on $`\mathrm{\Pi }T^{}M`$, and therefore, by (28), a generator of the Schouten bracket, which we shall denote by $`\mathrm{\Delta }^{Schouten,}`$. There exists another construction, due to Koszul , which associates a generator of the Schouten bracket to a torsionless linear connection $``$ on $`M`$. To $``$, he first associates the corresponding divergence operator, defined on vector fields $`X`$ by
(29)
$$\mathrm{div}_{}(X)=Tr(_X\mathrm{ad}_X).$$
This is the definition that is used in fact in (although it appears by mistake with the opposite sign in its first occurrence, page 262, before Lemma (2.1)). This map is a divergence operator, i. e., satisfies (4), on the purely even algebra $`C^{\mathrm{}}(M)`$. For a flat connection on flat space, it reduces to the elementary divergence. He then shows, using a local basis of vector fields, that there is a unique operator on the multivectors, $`\mathrm{\Delta }^{}`$, of degree $`1`$, that extends the operator $`\mathrm{div}_{}`$ and generates the Schouten bracket. We shall now show that the generators of the Schouten bracket obtained by these two constructions coincide.
###### Lemma 3.15.
For any vector field $`X`$, and for any $`1`$-form $`\alpha `$,
$$\mathrm{div}_{}(_X)=\mathrm{div}_{}(X),\mathrm{div}_{}(i_\alpha )=0.$$
Proof. If $`(x^1,\mathrm{},x^n)`$ is a system of local coordinates on $`M`$, then a local basis of graded derivations on $`\mathrm{\Pi }T^{}M`$ is $`(_{\frac{}{x^1}},\mathrm{},_{\frac{}{x^n}},i_{\mathrm{d}x^1},\mathrm{},i_{\mathrm{d}x^n}).`$ We use the relations $`<_{\frac{}{x^j}},\mathrm{d}^Gx^k>=\delta _j^k,`$ and $`<i_{\mathrm{d}x^j},\mathrm{d}^Gx^k>=0.`$ In order to compute $`\mathrm{div}_{}(_X)=sTr(__X\mathrm{ad}__X)`$, we first observe that, because $``$ is torsionless,
$$__Xi_{\mathrm{d}x^j}[_X,i_{\mathrm{d}x^j}]=_{i_{\mathrm{d}x^j}}_X=0.$$
Therefore
$`sTr(__X\mathrm{ad}__X)={\displaystyle \underset{j=1}{\overset{n}{}}}<__X_{\frac{}{x^j}}[_X,_{\frac{}{x^j}}],\mathrm{d}^Gx^j>`$
$`={\displaystyle \underset{j=1}{\overset{n}{}}}<_{_{\frac{}{x^j}}}_X,\mathrm{d}^Gx^j>={\displaystyle \underset{j=1}{\overset{n}{}}}<_{\frac{}{x^j}}X,\mathrm{d}x^j>=\mathrm{div}_{}(X).`$
###### Theorem 3.16.
For any torsionless linear connection $``$ on $`M`$, the generators $`\mathrm{\Delta }^{Schouten,}`$ and $`\mathrm{\Delta }^{}`$ of the Schouten bracket coincide. If $``$ is flat, this generator is of square $`0`$.
Proof. Since we know that both operators are generators of the Schouten bracket, we need only show that they coincide on functions and on vector fields. On functions, $`\mathrm{\Delta }^{Schouten,}`$ vanishes since $`[[f,]]=i_{df}`$ and $`\mathrm{div}_{}(i_{df})=0`$, as does $`\mathrm{\Delta }^{}`$ because it is of degree $`1`$. Now, for any vector field $`X`$ and any torsionless linear connection $``$ on an open set $`U`$ of $`M`$, $`[[X,.]]=_Xi_X`$, since both derivations of the sheaf of multivectors coincide on functions and on vectors. If $`(e_1,\mathrm{},e_n)`$ is a local basis of vector fields on $`U`$ and $`(ϵ^1,\mathrm{},ϵ^n)`$ is the dual basis, then the $`1`$-form-valued vector $`X`$ in $`U`$ can be written as $`_{j=1}^n_{e_j}Xϵ^j.`$ Thus $`i_X=_{e_j}Xi_{ϵ^j}.`$ Therefore, by Proposition 3.3 and Lemma 3.15,
$$\begin{array}{ccc}\hfill \mathrm{div}_{}(i_X)& =& \mathrm{div}_{}(_{j=1}^n_{e_j}Xi_{ϵ^j})\hfill \\ & =& _{j=1}^n_{e_j}X\mathrm{div}_{}(i_{ϵ^j})_{j=1}^nϵ^j(_{e_j}X)=\mathrm{div}_{}(X).\hfill \end{array}$$
It follows that
$$\mathrm{\Delta }^{Schouten,}(X)=\frac{1}{2}\mathrm{div}_{}([[X,]])=\mathrm{div}_{}(X).$$
It follows from Corollaries 1.4 and 3.14 together with Proposition 3.4 that, if $``$ is flat, the operator $`(\mathrm{\Delta }^{Schouten,})^2`$ is a derivation of the sheaf of multivectors with respect to the exterior product. Since, moreover, $`(\mathrm{\Delta }^{Schouten,})^2`$ is of $``$-degree $`2`$, it vanishes. ∎
Remark. Koszul proves that, conversely, any generator of the Schouten bracket of multivectors on a manifold $`M`$ is of the form $`\mathrm{\Delta }^{}`$ for some torsionless linear connection on $`M`$, and that two connections give rise to the same generator of the Schouten bracket if and only if they induce the same linear connection on $`^mTM`$, $`m`$ being the dimension of the manifold. We have not found any straightforward extension of this result to the case of odd Poisson brackets on supermanifolds in general.
#### 3.4.4. Conclusion
Given a smooth manifold $`M`$, we set $`A=C^{\mathrm{}}(M)`$, and we let $`\mathrm{Der}A`$ denote the module of vector fields on $`M`$. The definition of a divergence operator on a graded algebra reduces, in the purely even case of $`A`$, to the requirement that the linear operator, $`\mathrm{div}:\mathrm{Der}AA`$, satisfy the identity $`\mathrm{div}(fX)=f\mathrm{div}X+X(f)`$, for any $`fA`$ and $`X\mathrm{Der}A`$. The operators $`\mathrm{div}_{}`$, considered in Section 3.4.3, where $``$ are torsionless linear connections on $`M`$, are examples of divergence operators. Other examples are furnished by the operators $`\mathrm{div}_\mu `$ associated to volume forms, $`\mu `$, on an orientable manifold $`M`$. The Schouten bracket is an odd Poisson bracket on the graded commutative, associative algebra, $`𝚲=_A(\mathrm{Der}A)=_{k=0}^m_A^k(\mathrm{Der}A)`$, where $`m`$ is the dimension of the manifold $`M`$. It is the opposite of a divergence operator that can be extended into a generator of the Schouten bracket. In fact, for any divergence operator on $`A`$, the operator $`\mathrm{div}`$ can be uniquely extended to a generator of $``$-degree $`1`$, denoted $`\mathrm{\Delta }^{(\mathrm{div})}`$, of the Schouten bracket. One can characterize the generator $`\mathrm{\Delta }^{(\mathrm{div})}`$ recursively since, for any $`fA`$, it commutes with the interior product $`i_{\mathrm{d}f}`$. More generally, for any form $`\alpha `$,
$$[i_\alpha ,\mathrm{\Delta }^{(\mathrm{div})}]=i{}_{\mathrm{d}\alpha }{}^{}.$$
It is easy to see that, in the purely even case, a divergence operator $`\mathrm{div}:\mathrm{Der}AA`$ is nothing but a right $`(A,\mathrm{Der}A)`$-connection on $`A`$, in the sense of Huebschmann . In fact, if $`\mathrm{div}`$ is a divergence operator, then
$$(f,X)A\times \mathrm{Der}AfX=\mathrm{div}(fX)A$$
is a right $`(A,\mathrm{Der}A)`$-connection on $`A`$, and, conversely, if $`(f,X)A\times \mathrm{Der}AfXA`$ is a right $`(A,\mathrm{Der}A)`$-connection on $`A`$, then $`X\mathrm{Der}A1X`$, where $`1`$ is the unit of $`A`$, is a divergence operator. Moreover, the right $`(A,\mathrm{Der}A)`$-connection is a right $`(A,\mathrm{Der}A)`$-module structure if and only if the curvature of the divergence operator, defined by (5), vanishes. In fact,
$$(fX_1)X_2(fX_2)X_1f[X_1,X_2]=f^{\mathrm{div}}(X_1,X_2).$$
In the papers cited above, the notion of a divergence operator does not appear explictly, but the preceding remarks show that the 1-to-1 correspondence (, Theorem 1) between right $`(A,\mathrm{Der}A)`$-connections on $`A`$ and generating operators of the Schouten bracket of $`𝚲`$ yields a 1-to-1 correspondence between divergence operators and generators, which restricts to a 1-to-1 correspondence between divergence operators whose curvature vanishes and generators whose square vanishes. Also, the 1-to-1 correspondence (, Theorem 3) between right $`(A,\mathrm{Der}A)`$-connections on $`A`$ and left $`(A,\mathrm{Der}A)`$-connections on the top exterior power, $`_A^m\mathrm{Der}A=_A^{\mathrm{top}}\mathrm{Der}A`$, translates into a 1-to-1 correspondence between divergence operators and left $`(A,\mathrm{Der}A)`$-connections on the top exterior power. The canonical bundle $`_A^{\mathrm{top}}\mathrm{Der}A`$ is, in a natural way, a right module; equipping it with a left module structure, which can be done by choosing a volume element, is equivalent to equipping $`A`$ itself with a right module structure and therefore to selecting a divergence operator whose curvature vanishes (cf Propositon 2.3). To summarize, divergence operators, right connections on $`A`$, left connections on the top exterior power of $`\mathrm{Der}A`$, and generators of the Schouten bracket are in 1-to-1 correspondence. The definition of divergence operators and the preceding constructions extend to the framework of Lie algebroids and to that of Lie-Rinehart algebras. See , and also and . While there is a functor from Lie-Rinehart algebras to Gerstenhaber algebras, there is also a functor from Lie-Rinehart algebras with a divergence operator (resp., divergence operator with vanishing curvature) to Gerstenhaber algebras with a generator (resp., to Batalin-Vilkovisky algebras). In the case of a complex analytic manifold $`M`$ and its algebra of analytic functions, the left $`(A,\mathrm{Der}A)`$-module structures on the canonical bundle (top exterior power of holomorphic vector fields) are called Calabi-Yau structures . In this case, left (resp., right) $`(A,\mathrm{Der}A)`$-module structures coincide with left (resp., right) $`𝒟_M`$-module structures.
The extension of the above 1-to-1 correspondences valid in the purely even case to the case where $`A`$ itself is a $``$\- or $`_2`$-graded algebra, $`𝐀`$, remains to be done. The appropriate framework is that of the graded Lie-Rinehart algebras, whose theory has already been developped by Huebschmann (1990, unpublished), and left and right $`(𝐀,\mathrm{Der}𝐀)`$-connections and module structures in an appropriate sense. Sheaves of graded Lie-Rinehart algebras should then be considered, the fundamental example being $`(𝒜,\mathrm{Der}𝒜)`$, for any supermanifold $`(M,𝒜)`$. A divergence operator with vanishing curvature should define a right $`(𝒜,\mathrm{Der}𝒜)`$-module structure on $`𝒜`$, and there should be 1-to-1 correspondences between divergence operators, right structures on the structural sheaf and left structures on the berezinian sheaf. Another approach is by means of the theory of $`𝒟`$-modules. Left and right $`𝒟`$-modules on complex supermanifolds have been studied by Penkov , who showed that the berezinian sheaf of a complex analytic supermanifold is a right $`𝒟`$-module in a canonical way. Defining a left $`𝒟`$-module structure on the berezinian sheaf, which can be done by choosing a berezinian volume, is equivalent to defining a right $`𝒟`$-module structure on the structural sheaf, and should be equivalent to the choice of a divergence operator.
One can define graded analogues of the modules of multivectors on a manifold as modules of skew-symmetric multiderivations of $`𝐀`$, and one can generalize these notions to the case of sheaves of graded algebras over a manifold. Multigraded generalizations of the Schouten bracket on the space of skew-symmetric multiderivations of a graded algebra were defined by Krasil’shchik in , following his earlier paper . An analogue of the 1-to-1 correspondence between divergence operators and generators should be also valid in the graded case.
Conjecture. A divergence operator on the graded algebra $`𝐀`$, up to sign factors, can be uniquely extended to an operator on the skew-symmetric multiderivations of $`𝐀`$ that generates, in a suitable sense, the bigraded Krasil’shchik-Schouten bracket.
In particular, this construction would associate to a divergence operator on a supermanifold $`(M,𝒜)`$ a generator of the bigraded bracket on multivectors on the supermanifold. We hope to return to this question and also to study the relationship between the generators of a graded bracket and those of its derived brackets, in the sense of , in a future publication.
## Appendix. The berezinian sheaf
We shall recall the definition of the berezinian integral and some fundamental results, following , , , and mostly and .
Let $`(M,𝒜)`$ be a supermanifold of dimension $`m|n`$, in the sense of . Thus, $`M`$ is a smooth manifold and $`𝒜`$ is a sheaf of $`_2`$-graded commutative, associative $``$-algebras over $`M`$. There is an exact sequence
$$0𝒩𝒜𝒜/𝒩0,$$
where $`𝒩`$ is the sheaf of nilpotent sections of $`𝒜`$, and $`𝒜/𝒩`$ is the sheaf $`C^{\mathrm{}}(M)`$, regarded as trivially graded. The projection $`𝒜𝒜/𝒩=C^{\mathrm{}}(M)`$ is denoted by the symbol $`\stackrel{~}{},`$ and there is a unique prolongation to the module of differential forms of this projection, that commutes with the de Rham differentials. Thus, if we denote by $`\mathrm{d}`$ and $`\mathrm{d}^G`$ the de Rham differentials in $`M`$ and $`(M,𝒜)`$, then
$$\stackrel{~}{\mathrm{d}^G\alpha }=\mathrm{d}\stackrel{~}{\alpha },$$
for any differential form $`\alpha `$ on the supermanifold $`(M,𝒜)`$.
The berezinian sheaf can be described as follows. Let $`𝒫^k(𝒜)`$ be the vector space of the differential operators of order $`k`$ on $`𝒜`$. There is both a right and a left $`𝒜`$-module structure on $`𝒫^k(𝒜)`$ given by $`(f.P)(g)=f.P(g)`$ and $`(P.f)(g)=P(f.g)`$, respectively, for sections $`f,g`$ of $`𝒜`$ and $`P𝒫^k(𝒜)`$. If $`(x^1,\mathrm{},x^m,s^1,\mathrm{},s^n)`$ are graded coordinates on an open set $`U`$ in $`M`$, then $`𝒫^k(𝒜)|_U`$ is free for both structures of $`𝒜`$-module, with basis
$$(\frac{}{x^1})^{k_1}\mathrm{}(\frac{}{x^m})^{k_m}\frac{}{s^{\rho _1}}\mathrm{}\frac{}{s^{\rho _j}},$$
where $`k_1,\mathrm{},k_m`$, $`1\rho _1<\rho _2<\mathrm{}<\rho _jn`$ and $`k_1+\mathrm{}+k_m+j=k`$.
Let us now define $`𝒫^k(𝒜,\mathrm{\Omega }_𝒜^m)=\mathrm{\Omega }_𝒜^m𝒫^k(𝒜)`$, where $`\mathrm{\Omega }_𝒜^m`$ is the sheaf of differential $`m`$-forms on $`(M,𝒜)`$. Let $`𝒦^n`$ be the subsheaf of elements $`𝐏`$ in $`𝒫^n(𝒜,\mathrm{\Omega }_𝒜^m)`$ such that, for any section $`f`$ of $`𝒜`$ over an open set $`U`$ of $`M`$ with compact support, there exists an $`(m1)`$-differential form $`\omega `$ with compact support in $`U`$ such that $`\stackrel{~}{𝐏(f)}=\mathrm{d}\omega `$. Is is easy to show that $`𝒦^n`$ is a subsheaf of right $`𝒜`$-modules of $`𝒫^n(𝒜,\mathrm{\Omega }_𝒜^m)`$. The berezinian sheaf is the quotient sheaf $`𝒫^n(𝒜,\mathrm{\Omega }_𝒜^m)/𝒦^n.`$ The sections of this sheaf can be locally expressed as
$$[\mathrm{d}^Gx^1\mathrm{}\mathrm{d}^Gx^m\frac{}{s^1}\mathrm{}\frac{}{s^n}].f,$$
where $`f`$ is a section of $`𝒜`$. If $`VM`$ is an open set with graded coordinates $`(y^1,\mathrm{},y^m,t^1,\mathrm{},t^n)`$, then, on $`UV`$,
$$[\mathrm{d}^Gy^1\mathrm{}\mathrm{d}^Gy^m\frac{}{t^1}\mathrm{}\frac{}{t^n}]=[\mathrm{d}^Gx^1\mathrm{}\mathrm{d}^Gx^m\frac{}{s^1}\mathrm{}\frac{}{s^n}].Ber\left(\begin{array}{cc}A& B\\ C& D\end{array}\right)$$
where
$$\left(\begin{array}{cc}A& B\\ C& D\end{array}\right)=\left(\begin{array}{cc}\left(\frac{y^i}{x^j}\right)& \left(\frac{t^\rho }{x^j}\right)\\ \left(\frac{y^i}{s^\sigma }\right)& \left(\frac{t^\rho }{s^\sigma }\right)\end{array}\right)$$
and where $`Ber`$ denotes the berezinian. (The berezinian, or superdeterminant, of an invertible even matrix of the form $`\left(\begin{array}{cc}A& B\\ C& D\end{array}\right)`$ is $`det(ABD^1C)det(D^1)`$.)
If $`M`$ is an orientable smooth manifold, the Berezin integral$`_{(M,𝒜)}`$ , maps the sections with compact support of the berezinian sheaf to $``$, and is defined by $`_{(M,𝒜)}[𝐏]=_M\stackrel{~}{𝐏(1)}.`$ As an example, if $`(M,𝒜)=^{m|n}`$, then
$$_{^{m|n}}[\mathrm{d}^Gx^1\mathrm{}\mathrm{d}^Gx^m\frac{}{s^1}\mathrm{}\frac{}{s^n}].f=(1)^{\frac{n(n1)}{2}}_^nf_{(1,2,\mathrm{},n)}(x)\mathrm{d}x^1\mathrm{}\mathrm{d}x^m,$$
where $`f_{(1,2,\mathrm{},n)}`$ is the coefficient of $`s^1s^2\mathrm{}s^n`$ in the expansion of $`f`$ as a sum of products of the $`s^\rho `$’s.
A section, $`\xi `$, of the berezinian sheaf is called a berezinian volume if it is a generator of the berezinian sheaf, i.e., if any other section can be uniquely written as $`\xi .f`$ for some section $`f`$ of $`𝒜`$. A berezinian volume is a homogeneous section of the berezinian sheaf, whose degree depends on the parity of the dimension $`n`$. If $`\xi `$ is a berezinian volume and $`v`$ is a section of $`𝒜`$, then $`\xi .v`$ is also a berezinian volume if and only $`v`$ is invertible and even.
In order to define the Lie derivatives of berezinian volumes with respect to graded vector fields, we first observe that, in a similar way, we can define the right submodule $`𝒦^{n+k}`$ of $`𝒫^{n+k}(𝒜,\mathrm{\Omega }_𝒜^m)`$, for each $`k1`$, and that the canonical inclusion $`𝒫^n(𝒜,\mathrm{\Omega }_𝒜^m)𝒫^{n+k}(𝒜,\mathrm{\Omega }_𝒜^m)`$ induces an isomorphism of sheaves of right $`𝒜`$-modules from $`𝒫^n(𝒜,\mathrm{\Omega }_𝒜^m)/𝒦^n`$ to $`𝒫^{n+k}(𝒜,\mathrm{\Omega }_𝒜^m)/𝒦^{n+k}.`$
Let $`D`$ be a graded vector fied on $`(M,𝒜)`$. The Lie derivative of the berezinian volume $`[\omega P]`$ with respect to $`D`$ is
$$_D[\omega P]=(1)^{|D||\omega P|+1}[\omega PD].$$
The main properties of the Lie derivatives of berezinian volumes are stated in Section 2, and are used there in order to derive the properties of the divergence operators.
Acknowledgments. Y. K.-S. would like to thank P. Cartier, J. Huebschmann, D. Leites, H. Khudaverdian, A. Schwarz, J. Stasheff and T. Voronov for interesting exchanges on the topic of this work. Thanks are also due to the referee who identified various weaknesses in the first version of this paper.
J. M. is partially supported by Pla Valencià de Ciència i Tecnologia, grant $`\mathrm{\#}POST99`$-$`01`$-$`30`$, and by DGICYT, grant $`\mathrm{\#}PB97`$-$`1386`$. |
warning/0002/cond-mat0002192.html | ar5iv | text | # Heat conduction in one dimensional nonintegrable systems
## Abstract
Two classes of 1D nonintegrable systems represented by the Fermi-Pasta-Ulam (FPU) model and the discrete $`\varphi ^4`$ model are studied to seek a generic mechanism of energy transport in microscopic level sustaining macroscopic behaviors. The results enable us to understand why the class represented by the $`\varphi ^4`$ model has a normal thermal conductivity and the class represented by the FPU model does not even though the temperature gradient can be established.
preprint: HKBU-CNS-9904
Heat conduction in one-dimensional (1D) nonintegrable Hamiltonian systems is a vivid example for studying microscopic origin of the macroscopic irreversibility in terms of deterministic chaos. It is one of the oldest but a rather fundamental problem in nonequilibrium statistical mechanics. Intended to understand the underlying mechanism of the Fourier heat conduction law, the study of heat conduction has attracted increasing attention in recent years .
Based on previous studies, we can classify the 1D lattices into three categories. The first one consists of integrable systems such as the harmonic chain. It was rigorously shown that, in this category no temperature gradient can be formed, and the thermal conductivity is divergent. The second category includes a number of nonintegrable systems such as the Lorentz gas model, the ding-a-ling and alike models, and the Frenkel-Kontorova (FK) model and etc.. In this category, the heat current is proportional to $`N^1`$ and the temperature gradient $`dT/dxN^1`$, thus the thermal conductivity $`\kappa `$ is a constant independent of system size $`N`$. The Fourier heat conduction law ($`J=\kappa dT/dx`$) is justified. The third category also includes some nonintegrable systems such as the FPU chain, the diatomic Toda chain, the (mass) disorder chain, and the Heisenberg spin chain and so on. In this category, although the temperature gradient can be set up with $`dT/dxN^1`$, the heat current is proportional to $`N^{\alpha 1}`$ with $`\alpha 0.43`$, the thermal conductivity $`\kappa N^\alpha `$ which is divergent as one goes to the thermodynamic limit $`N\mathrm{}`$.
These facts suggest that the nonintegrability is necessary to have a temperature gradient, but it is not sufficient to guarantee the normal thermal conductivity in a 1D lattice. This picture brings us to ask two questions of fundamental importance: (i) Why do some nonintegrable systems have normal thermal conductivity, while the others fail? (ii) How can the temperature gradient be established in those nonintegrable systems having divergent thermal conductivity?
The reason for the divergent thermal conductivity in integrable system is that the energy transports freely along the chain without any loss so that no temperature gradient can be established. The set up of temperature gradient in nonintegrable systems implies the existence of scattering. However, the different heat conduction behavior in the two categories of nonintegrable systems indicates that the underlying mechanism must be different. To get the point, let’s write the Hamiltonian of a generic 1D lattice as
$$H=\underset{i}{}H_i,H_i=\frac{p_i^2}{2}+V(x_{i1},x_i)+U(x_i),$$
(1)
where $`V(x_{i1},x_i)`$ stands for the interaction potential of the nearest-neighbor particle and $`U(x_i)`$ is an external (or on-site) potential. The origin of external potential in real physical systems varies from model to model. For instance, in the FK model the external potential is the interaction of the adsorbed atoms with the crystal surface. It is $`U(x)`$ that distinguishs from the two categories of nonintegrable lattices. $`U(x)`$ vanishes in all 1D lattices having divergent thermal conductivity. We are thus convinced to conclude that the external potential plays a determinant role for normal thermal conduction.
In this paper we would like to study the scattering mechanism and the role of the external potential in heat conduction in the two categories of nonintegrable systems. For this purpose, we choose two representatives from these two categories, i.e. the discrete $`\varphi ^4`$ model (see, e.g. Ref. ) and the FPU model. Both models are the simplest anharmonic approximation of a monoatomic solid. In the $`\varphi ^4`$ model, $`V`$ takes the harmonic form, and the external potential $`U(x)=mx^2/2+\beta x^4/4`$, with $`m`$ fixed to be zero in this paper. In the FPU model, $`U`$ vanishes and $`V`$ takes the anharmonic form of $`(x_ix_{i1})^2/2+\beta (x_ix_{i1})^4/4`$, and $`\beta =1`$ throughout this paper. In the case of $`\beta =0`$, the FPU model reduces to the harmonic chain.
In our numerical simulations the Nosé-Hoover thermostates are put on the first and the last particles, keeping them at temperature $`T_+`$ and $`T_{}`$, respectively. The motion of these two particles are governed by
$`\ddot{x}_1=\zeta _+\dot{x}_1+f_1f_2,\dot{\zeta }_+=\dot{x}_1^2/T_+1`$ (2)
$`\ddot{x}_N=\zeta _{}\dot{x}_N+f_Nf_{N+1},\dot{\zeta }_{}=\dot{x}_1^2/T_{}1.`$ (3)
where $`f_i=(V^{}+U^{})`$ is the force acting on the $`i`$’th particle. The equation of motion of other particle is $`\ddot{x}_i=f_if_{i+1}`$. The eighth-order Runge-Kutta algorithm was used. All computations are carried out in double precision. Usually the stationary state set in after $`10^7`$ time units. We should point out that we have performed computations by using other types of thermostate, and no qualitative difference has been found.
Fig. 1(a) shows temperature profiles. In all nonintegrable systems, the temperature scales as $`T=T(i/N)`$. However, in the FPU case there is a singular behavior near the two ends, which is a typical character of 1D nonlinear lattices having divergent thermal conductivity. In the same figure we also show the temperature profiles for two integrable lattices: the harmonic and the monoatomic Toda models. In these two cases no temperature gradient could be set up and the stationary state corresponds to $`T=(T_++T_{})/2`$, which is consistent with the rigorous result .
In Fig. 1(b), we plot the quantity $`J\times N`$ versus $`N`$ for the FPU model and the $`\varphi ^4`$ model. The inset shows the same quantity for the harmonic chain and the monoatomic Toda chain. The local heat flux is defined by $`J_i=\dot{x}_i\frac{V}{x_{i+1}}`$. We found that when the system reaches a stationary state, the time average $`J_i`$ is site independent, it is denoted as $`J`$. For the harmonic chain and the monoatomic Toda chain $`J\times N`$ is expected to be proportional to $`N`$ since $`J`$ is $`N`$-independent. This is indeed the case as illustrated in the inset. In both the FPU and the $`\varphi ^4`$ models $`dT/dx`$ is proportional to $`1/N`$, the thermal conductivity $`\kappa =J/(dT/dx)J\times N`$. Fig. 1(b) tells us that, in contrast to the integrable systems and the FPU model, heat conduction in the $`\varphi ^4`$ model obeys the Fourier law.
The heat current $`J`$ in all nonintegrable systems decreases as the system size $`N`$ is increased ($`JN^{\alpha 1},0<\alpha <1`$). To clarify the underlying mechanism we decompose the interaction of the thermostat into a series of kicks and study the transport of a single kick along the chain. A free boundary condition is used in our calculation, but we should stress that the results do not depend on the types of boundary condition. In Fig. 2 we plot $`p_i`$ versus $`i`$ after a long time ($`t=800)`$ for four lattices: the harmonic chain (a); the monoatomic Toda chain (b); the FPU model (c); and the $`\varphi ^4`$ model (d). The amplitude of the wave profile in the harmonic chain decreases continuously with time, but the global profile keeps unchanged. In both the FPU and the Toda models, we observe a solitary wave separates from the long tail. Initially this wave front is connected with other low amplitude excitations. After a certain time, this wave front, moves faster, separates from the tail and goes forward and keeps its amplitude. In the mean time the tails behind it evolve in the same way as that in the harmonic chain. In the $`\varphi ^4`$ model, the head part of the profile becomes weaker and weaker. The reason is that in the first three cases (a-c) both the total energy and the total momentum are conserved, whereas in the $`\varphi ^4`$ model the momentum conservation breaks down due to the external potential. The inset in Fig. 2(d) shows that the total momentum in the $`\varphi ^4`$ model decreases at least exponentially with time. The decay of the momentum with time indicates a loss of correlation. It is thus reasonable to envisage the energy transport along the $`\varphi ^4`$ chain as a random walk-like scattering.
The solitary waves in the FPU chain exchange energy and momentum when colliding with each other. It causes the energy loss, and the heat current decreases when the system size is increased. To show this, we start two excitations at the two ends of the chain with different momentum, one moves to the right and another left. Let $`p_1=6`$, $`p_N=3`$ and $`p_i=0,i1`$ and $`N`$ as our initial excitations. We calculate the momentums of the solitary waves (by simply summing up momentums of several lattices around the peaks) and investigate its change before and after the interaction. We find that the bigger one generally transfers part of its momentum and energy to the smaller one, as is shown in Fig. 3(a). The collision takes place at $`t=850`$, where a peak is shown up.
Moreover, the interaction between solitary waves is found to depend closely on a “phase” difference. Here the “phase” difference is defined as a time lag between the excitations of two solitary waves. For instance, if we excite a solitary wave from the left end at time $`t`$, and another one from the right end at time $`t+\delta `$, then $`\delta `$ is the “phase” difference. These two solitary waves, traveling through the chain in opposite directions, will collide with each other after a certain time. Although the physical meaning of the “phase” is not obvious, it is an important and good quantity to describe the interaction. We show $`p_L^a`$ versus $`\delta `$ for two different kinds of collision in Fig 3(b), where $`p_L^a`$ is the momentum of the solitary wave from the left after collision. In the first case both left and right solitary waves have the same initial momentum $`p_L=p_R=6.27`$, which is excited by an initial condition of $`p_1=p_N=6`$ and $`p_i=0`$ for other $`i`$. In the second case, the left one has $`p_L=6.27`$ and the right one $`p_R=3.28`$, excited by an initial condition of $`p_1=6`$, $`p_N=3`$ and $`p_i=0`$ for $`i`$ $`1`$ and $`N`$. The figure shows that $`P_L^a`$ depends on the “phase” sinusoidally.
Other interesting features of the collision of solitary waves are shown in Fig. 3(c), where we plot the maximum momentum gain $`\mathrm{\Delta }p_{max}`$ versus the initial momentum $`p_0`$ for the FPU model. $`\mathrm{\Delta }p_{max}`$ is measured by subtracting the initial momentum $`p_0`$ from the maximum $`p_L`$ in Fig. 3b. First of all, this picture tells us that the exchange of momentum and energy depends on the initial momentum and energy. Secondly, there exists a critical momentum below which no energy exchange can take place. The critical $`p_0^c1.8`$ is clearly seen in the figure. For $`p_0<p_0^c`$, $`\mathrm{\Delta }p_{max}`$ is zero. This result is very significant, it indicates that there exists a threshold for the solitary wave interaction, below this threshold the interaction ceases, i.e. no momentum and energy exchanges between the solitary waves. A direct consequence of this fact is the existence of a threshold temperature below which the FPU chain should behave like a harmonic chain, namely, the excited waves travel freely along the chain without any energy loss, no temperature gradient can be set up, and the heat current remains a constant even though the size of the chain is changed. To testify this argument, we show the quantity $`J(400)/J(800)`$ versus $`T=(T_++T_{})/2`$ in Fig. 3(d), where $`J(N)`$ is the heat current flux for the system of size $`N`$. In the case of a size-independent $`J(N)`$ one should get $`J(400)/J(800)=1`$, otherwise one would get $`J(400)/J(800)>1`$. Fig. 3(d) captures this transition nicely for the FPU chain. The corresponding temperature threshold is about $`T_c0.01`$. In the region of $`T0.001`$ the numerical calculations do show that no temperature gradient is formed.
The different scattering mechanism in the FPU chain and those chains having normal thermal conductivity leads to a different temperature dependence of thermal conductivity $`\kappa (T)`$. In Fig. 4 we plot $`\kappa (T)`$ for a FPU chain with an external potential of form $`U(x)=\gamma \mathrm{cos}(2\pi x)`$ for four different values of $`\gamma =0,0.01,0.05`$ and $`0.1`$. The chain size is fixed at $`N=100`$. As pointed out above, in small $`\gamma `$ regime such as $`\gamma =0`$ and $`0.01`$, the energy transport is assisted by the solitary waves, the system has a large $`\kappa `$ which decreases as the temperature is increased. However, in the opposite regime ($`\gamma =0.05`$ and $`0.1`$), the energy transport is diffusive and obeys the Forier law, $`\kappa `$ increases with temperature, because more phonons are excited.
In summary, by studying two classes of 1D nonintegrable lattices, we have answered the two questions raised at the beginning: (i) The multiple scattering of the excited modes by the external potential leads to a decay (at least exponentially) of correlation, so that a diffusive transport process can be reached, and the heat conduction obeys the Fourier law; (ii) Although the interaction of solitary waves makes it possible to set up temperature gradient in the FPU and alike nonintegrable models, the momentum conservation prohibits the diffusive transport and consequently leading to the divergent thermal conductivity. In addition, we have uncovered an important fact in the FPU model, namely, the existence of a threshold temperature, below which the FPU mode behaves like a harmonic chain.
BL would like to thank G. Casati for useful discussions. This work was supported in part by Hong Kong Research Grant Council and the Hong Kong Baptist University Faculty Research Grant.
Note added in proof. – After submission of this paper we get to know the following results. Prosen and Campbell proved in a more rigorous way that for a 1D classical many-body lattice total momentum conservation implies anomalous conductivity. The normal thermal conductivity in the $`\varphi ^4`$ lattice has also been observed by Aoki and Kusnezov. The role of the external potential has been further studied by Tsironis et al. |
warning/0002/cond-mat0002275.html | ar5iv | text | # Silence of magnetic layers to magnetoresistive process and electronic separation at low temperatures in (La, Sm)Mn2Ge2
## Abstract
A closer look at the temperature (T) dependence of magnetoresistance (MR) of two polycrystalline magnetic compounds, LaMn<sub>2</sub>Ge<sub>2</sub> and SmMn<sub>2</sub>Ge<sub>2</sub>, previously reported by us, is made. A common feature for both these compounds is that the low temperature MR is positive (say, below, 30 K) in spite of the fact that both are ferromagnetic at such low temperatures; in addition, MR as a function of magnetic field (H) does not track magnetization (M) in the sense that M saturates at low fields, while MR varies linearly with H. These observations suggest that the magnetic layers interestingly do not dominate low temperature magnetotransport process. Interestingly enough, as the T is increased, say around 100 K, these magnetic layers dominate MR process as evidenced by the tracking of M and MR in SmMn<sub>2</sub>Ge<sub>2</sub>. These results tempts us to propose that there is an unusual "electronic separation" for MR process as the T is lowered in this class of compounds.
PACS. 72.20.My - Magnetotransport effects.
PACS. 72.15.-V - Electronic conduction in metals and alloys.
PACS. 73.61.-r - Electrical properties of layered structures.
E-mail: sampath@tifr.res.in
Keywords: LaMn<sub>2</sub>Ge<sub>2</sub>, SmMn<sub>2</sub>Ge<sub>2</sub>, Magnetoresistance, Electronic separation
The understanding of electrical resistance (R) behaviour in solids continues to be a major direction of research in condensed matter physics. In general, in metals containing magnetic impurities or a large concentration of magnetic-moment-carrying ions, one observes a dominant signature of these magnetic ions in the low temperature (T) electrical resistance, which has a characteristic response with the application of an external magnetic field (H) depending upon the magnetic state. For instance, in ferromagnetic materials R decreases with increasing H; Y (non-magnetic), if contains Ce in traces, exhibits features due to the Kondo effect resulting in a decrease of R with H . Here we argue that, in sharp contrast to this expectation, in a class of layered ternary intermetallic compounds containing a magnetically ordered layer, viz., RMn<sub>2</sub>Ge<sub>2</sub> (R= La, Sm), the non-magnetic atomic layers dominate the magnetoresistance (MR) at low temperatures with the magnetic layers apparently remaining "silent"; however, with increasing temperature, the contribution from the magnetic atomic layers is visible. This is taken as an evidence to suggest that there is a temperature dependence associated for the involvement of different layers; in other words apparently there is "an electronic separation" with decreasing temperature for the magnetoresistive process. This finding adds a new dimension to the understanding of magnetotransport phenomena in solids.
The intermetallic compounds under discussion crystallize in the well-known ThCr<sub>2</sub>Si<sub>2</sub>-type tetragonal structure , containing layers of atoms stacked in the sequence Th-Si-Cr-Si-Th along the c-axis. It is interesting to note that, among few hundred compounds known to form in this structure, the Mn is the only transition metal ion known to possess magnetic moment and magnetic ordering at very high temperatures. The nature of the magnetic ordering apparently is sensitive to Mn-Mn distances and this appears to explain the presence of multiple magnetic transitions for Sm compound with the Mn-Mn distance for this compound falling near the critical limit . Thus, while the former compound has been known to order ferromagnetically at about 300 K with the magnetism arising from Mn sublattice, the latter exhibits multiple magnetic transitions: below about 345 K, para- to ferro-magnetism; at about 140 K, ferro to anti-ferromagnetism; and antiferro to ferromagnetism below about 105 K (see, for instance, Refs. 3-13). It is surprising to note that, inspite of extensive magnetic investigations on this class of ternary compounds, virtually there has been very little MR studies on these Mn alloys. Though the compound SmMn<sub>2</sub>Ge<sub>2</sub> was investigated in the range 80 - 150 K to show giant magnetoresistance effects earlier , surprisingly these authors did not report the MR at lower temperatures. Considering a recent upsurge in MR studies in condensed matter physics, we considered it important to carry out such studies on these compounds and thus we reported the MR alongwith detailed magnetic measurements down to 4.2 K ; the major point of emphasis was to bring out the relevance of these compounds to the physics of artificial multilayers. In this letter, for the La and Sm compounds, we take into account available data in the literature and also compare the T dependence of MR to draw the present conclusion. As far as our data is concerned, we reproduce here only those data which are required to emphasize this point.
In figure 1a, we show R as a function of T in zero field as well as at 50 kOe (with the direction of H being parallel to that of the excitation current) for a specimen of LaMn<sub>2</sub>Ge<sub>2</sub> below 60 K. The magnetoresistance \[MR={R(H)-R(0)}/R(0)\] behaviour, obtained by subtracting the zero-field R data from that at 50 kOe is shown in figure 1b. The value of MR increases with decreasing temperature below 60 K eventually attaining large values at 4.2 K; at temperatures higher than 60 K the value of MR is negligibly small. What is intriguing is that the sign of MR is positive at low temperatures, inspite of the fact that the compound contains ferromagnetic Mn layers in which case one should have observed MR with a negative sign expected for polycrystalline ferromagnetic metals . (Therefore, it is the ferromagnetic nature of the magnetic ordering that enables us to draw the present conclusions on firm grounds, as positive MR for antiferromagnetics is usually expected.) We have obtained MR on several specimens of this compound prepared under different conditions of heat-treatment and the positive sign of MR is always reproducable .The
positive sign of MR is taken as an evidence for the dominant role of non-magnetic layers in controlling magnetoresistive process. Needless to mention that the positive sign of MR is characteristic of non-magnetic metals. Even if one attributes positive sign of MR to possible deviations from collinear ferromagnetism of the Mn sublattice or to any other mechanism arising from ferromagnetism, the following finding goes against these possibilities. That is, the magnetization does not track MR (e.g., at 4.2 K, see Fig. 1c) in the sense that M saturates for small applications of H while MR is a linearly varying function of H. This is a key finding in favour of our proposal that in this compound the Mn magnetic layer is not apparently involved in the magnetotransport phenomena at low temperatures. In metallic materials, at low temperatures, the disorder/impurity scattering contribution dominates as the phonon contribution tends to vanish. Therefore, one may advance an arguement that the positive sign of MR somehow arises from such crystallographic imperfections. However, a careful look at how MR varies as a function of the residual resistivity ratio, RRR \[= R(4.2K)/R(300K)\], in specimens with different degree of crystallographic imperfections clearly reveals that the increasing disorder/imperfections in fact tend to diminish the net change in R by the application of H at low T. This naturally implies that the observed positive sign with a large magnitude is intrinsic to the crystallographically well-ordered material. We therefore conclude that the positive MR arises from La and/or Ge layer only, without involving magnetic-moment-carrying Mn ions.
The present conclusion is novel as it has been generally believed that the signature of moment-carrying ions should be prominent in the low temperature transport behaviour; this contribution can be so prominent that even in the dilute limit of 3d magnetic impurities in non-magnetic matrices one usually encounters the phenomenon of the Kondo effect with corresponding magnetoresistive response. \[It may be added that the temperature dependence of mean free path (mfp) also seems to correlate with MR . Inspite of the fact that there are uncertainties in the absolute values of mfp due to approximations in its determination, it appears that the critical value for the observation of significant positive MR is of the order of interlayer spacings. We believe that this information may be useful for future theoretical work.\]
Next, one might be tempted to ask the following questions. Which of the two layers, La or Ge, dominates low temperature scattering phenomenon? Do magnetic layers participate in magnetotransport phenomena as the temperature is increased? The transport behaviour of SmMn<sub>2</sub>Ge<sub>2</sub> offers straightforward answers for these questions. The data are shown in Fig. 2a in zero and 50 kOe field. MR is distinctly positive at low temperatures even in this compound, without tracking isothermal M (see Fig. 2c), the origin of which must be the same as that in the La compound. However, as the temperature is increased, as known earlier , the jumps in R in zero field at about 110 K and 140 K attributable to magnetic transitions arising from Mn are clearly seen. The application of H wipes out these jumps due to the well-known sharp metamagnetic transition occuring at about 5 kOe, thus resulting in large MR anomalies at these temperatures (Fig. 2b). Thus, it is evident that the Mn (magnetic) layer gets involved in the transport process at these temperatures. This conclusion is further endorsed by the observation that MR tracks magnetization as a function of H at 104 K, unlike the situation at 4.2 K. In fact, the expected negative sign of MR starts appearing as the temperature is increased beyond 30 K, thereby indicating that the (ferro)magnetic layer dominates above this temperature in this compound. It is also important to note that Sm layer also orders ferromagnetically below 100 K and if this layer dominates low temperature ($`<`$ 30 K) MR, one should have seen negative MR tracking isothermal M, in sharp contrast to the experimental observations. This establishes that it is not even the rare-earth layer, but the Ge layer, that dominates low temperature magnetotransport phenomena.
Summarising, a comparison of MR behaviour of the compounds LaMn<sub>2</sub>Ge<sub>2</sub> and SmMn<sub>2</sub>Ge<sub>2</sub> suggests that the dominance of magnetic layers to the magnetotransport process diminishes with decreasing temperature, meaning thereby that there is an unusual temperature dependence associated with the relative involvement of layers of atoms for this process. Thus there appears to be an interesting "electronic separation" at the unit-cell level as the temperature is lowered as far as the magnetotransport is concerned. The above suggestion is made under the assumption that there are no unusual band structure effect on MR (which, if present, is by itself again interesting); the presently available band structure data do not reveal any unusual structure at the Fermi level within the energy scale of the magnitude of applied magnetic field and hence we advance the present line of thoughts. At this point, it is worth mentioning that, at the time of finalising this article, the idea of an electronic phase separation at a length scale much larger than the unit-cell dimensions has been proposed to explain the colossal magnetoresistance in mixed-valent manganites . It is fascinating that some kind of electronic separation (as far as magnetoresistance is concerned) may be possible with the variation of temperature even at the unit-cell level as indicated by the present data. Realisation of this possibility will go a long way in understanding the transport process in modern condensed matter physics. |
warning/0002/hep-th0002215.html | ar5iv | text | # Cosmological perturbation spectra from SL(4,R)-invariant effective actions
## I Introduction
The pre-big-bang scenario proposed by Gasperini and Veneziano is an alternative model for the very early evolution of our Universe which assumes that its initial state was a low-energy, weakly coupled, vacuum state. Such a regime is well described by the low-energy string effective action which admits two separate branches, labelled $`(+)`$ and $`()`$, for vacuum solutions of the scale factor in four-dimensional Friedmann-Robertson-Walker (FRW) cosmologies. The $`(+)`$ branch corresponds to a weakly coupled dilaton in a cold, flat universe in the $`t\mathrm{}`$ limit. For $`t0`$ we have pole driven super-inflation propelled by the dilaton kinetic energy term, with a positive Hubble parameter, $`H`$ and a singularity in the future. The $`()`$ branch corresponds to a large spatially flat universe with positive $`H`$, but decelerating and can be smoothly joined to a conventional radiation dominated universe at late times. The $`(+)`$ and $`()`$ branches are related to each other by a string symmetry called scale factor duality , but there is still no compelling dynamical model of the “graceful exit” from the $`(+)`$ to the $`()`$ branch .
In the absence of a complete theoretical understanding one may still hope to find observational evidence, such as the spectrum of primordial fluctuations that could be generated during the dilaton-driven pre-big-bang phase , in order to test the scenario against more conventional inflation models. Metric perturbations are produced on super-horizon scales during the pre big bang have a steep “blue” spectrum, strongly tilted towards small scales . This offers the interesting possibility that there might be a detectable background of relic gravitons on Laser Interferometric Gravitational wave Observatory (LIGO) scales and a related population of primordial black holes . However these metric perturbations are far from the almost scale-invariant (Harrison-Zel’dovich) spectrum of adiabatic density perturbations naturally produced by conventional slow-roll inflation models and leave effectively no metric perturbations on astrophysical scales during the pre-big-bang era.
Instead it has been proposed that a cosmic background of massless axion fluctuations could generate the observed anisotropies in the cosmic microwave background temperature at large angular scales and provide a seed for large-scale structure formation . A pre-big-bang era can produce almost scale invariant spectra of fluctuations in axion fields $`\delta \sigma ^2k^{\mathrm{\Delta }n}`$ with $`\mathrm{\Delta }n0`$, where $`k`$ is the comoving wavenumber. These are isocurvature perturbations to first-order, but assuming the axion field remains effectively massless in the subsequent post big bang era, these fluctuations give rise to a spectrum of density perturbations at horizon crossing
$$\left(\frac{\delta \rho }{\rho }\right)_{k=aH}e^\varphi \left(\frac{k}{k_s}\right)^{\mathrm{\Delta }n},$$
(1)
where $`k_s`$ is the comoving scale leaving the horizon at the end of the pre-big-bang phase. A slightly “blue” spectrum, $`\mathrm{\Delta }n+0.1`$, may be consistent with $`\delta \rho /\rho 10^5`$ on astrophysical scales ($`k10^{30}k_s`$) for a present-day string coupling $`g_s^2=e^\varphi 10^2`$.
The low energy string effective action compactified down to four-dimensions includes a dilaton and axion field, related by an SL(2,R) symmetry. In the absence of all the other moduli fields are fixed during the pre-big-bang phase then it is found that for an SL(2,R) action, with one axion field, the spectral index is fixed to be $`\mathrm{\Delta }n=2\sqrt{2}+3=0.46`$. The addition of a single moduli field gives a spectral index for the axion in the range $`0.46\mathrm{\Delta }n3`$ which allows scale-invariant or blue spectra .
The many moduli fields present in any low-energy effective action will have specific symmetry properties inherited from the higher dimensional theory and the details of compactification. For instance, the inclusion of Ramond-Ramond (RR) fields presents in the type II string theories increases the number of degrees of freedom in the four-dimensional effective theory and in Ref. it was shown that the RR 1-form and 3-form field strengths, with a single modulus field determining the size of the 6-torus, combine with the Neveu-Schwarz-Neveu-Schwarz(NS-NS) dilaton and axion to parameterise an SL(3,R)-invariant non-linear sigma model. The symmetries of this action can place constraints on the allowed spectral indices. For an SL(3,R) action , with two moduli but three axion fields, the range for each spectral index was the same as for the single axion case ($`0.46\mathrm{\Delta }n_i3,i=1,2,3`$), but there was no point at which all the spectra had $`\mathrm{\Delta }n_i>0`$. This poses a threat to the pre-big-bang scenario as all the perturbation spectra have the same normalisation at the string scale ($`\delta \rho /\rho e^\varphi `$) and if one axion field always has a red spectrum ($`\mathrm{\Delta }n_i<0`$) then there would be unacceptably large density fluctuations on large scales.
In order to study whether this remains a problem in larger symmetry groups with more moduli and axion degrees of freedom we will study the spectrum of cosmological perturbations generated in fields which parameterise an SL(4,R) non-linear sigma model in the low-energy effective action. The presence of a global SL($`n`$,R) symmetry is a completely general consequence of dimensional reduction from $`D+n`$ to $`D`$ dimensions . We will investigate the perturbation spectra generated in axion fields at late times or large scales from vacuum fluctuations at early times or small scales in a pre-big-bang era, and discuss whether these might be compatible with an almost scale invariant spectrum of small primordial density perturbations. We will show that perturbation spectra with $`\mathrm{\Delta }n>0`$ for all fields are indeed possible in models whose scalar fields parameterise an SL($`n`$,R) group where $`n4`$.
## II SL($`n`$,R) invariant actions
We begin with a discussion of the representation of the most familiar SL(2,R)-invariant effective action is string theory, and how we can extend this to provide a representation of larger SL($`n`$,R)-invariant effective actions.
The NS-NS sector of string theory contains the dilaton, $`\varphi `$, graviton, $`\widehat{G}_{AB}`$, and 2-form potential, $`B_{AB}`$, and is common to both heterotic and type II string theories. The low energy effective action is
$$S=\frac{1}{16\pi \alpha ^4}d^{10}x\sqrt{|\widehat{G}|}e^\varphi \left[\widehat{R}_{10}+(\widehat{}\varphi )^2\frac{1}{12}\widehat{H}^2\right],$$
(2)
where $`\alpha ^{}`$ is the inverse string tension and $`\widehat{H}dB`$ is a 3-form field strength. Considering the simplest Kaluza-Klein compactification on a static six-torus, the ten-dimensional line element has the form $`ds_{10}^2=e^\varphi g_{\mu \nu }dx^\mu dx^\nu +\delta _{ab}dy^ady^b`$, where $`g_{\mu \nu }`$ is the four-dimensional metric in the Einstein frame and $`\delta _{ab}`$ is the Kronecker delta-function. In four-dimensions the 3-form field strength is dual to a one-form written as the gradient of a pseudo-scalar axion field, $`H=e^\varphi \sigma `$, and the four-dimensional effective action is then
$$S=\frac{1}{2\kappa ^2}d^4x\sqrt{g}\left[R\frac{1}{2}\left(\varphi \right)^2\frac{1}{2}e^{2\varphi }\left(\sigma \right)^2\right],$$
(3)
where $`\kappa ^28\pi \mathrm{G}`$. Solutions to the equations of motion from this action respect the invariance of this action under an arbitrary global SL(2,R) transformation
$$\lambda \frac{\alpha \lambda +\beta }{\gamma \lambda +\delta },$$
(4)
where $`\lambda =\sigma +ie^\varphi `$ and the real parameters $`\alpha `$, $`\beta `$, $`\gamma `$ and $`\delta `$ obey the constraint $`\alpha \delta \beta \gamma =1`$.
In order to extend the effective action to larger symmetry groups it is convenient to re-write the action Eq. (3) in terms of the symmetric matrix
$$M=\left(\begin{array}{cc}e^\varphi \hfill & \sigma e^\varphi \hfill \\ \sigma e^\varphi \hfill & e^\varphi +\sigma ^2e^\varphi \hfill \end{array}\right)$$
(5)
which parameterises the SL(2,R)/U(1) maximal coset of SL(2,R). The SL(2,R) transformation in Eq. (4) is given by $`M\mathrm{\Theta }M\mathrm{\Theta }^T`$ where
$$\mathrm{\Theta }=\left(\begin{array}{cc}\delta & \gamma \\ \beta & \alpha \end{array}\right).$$
Any member of SL(2,R) obeys the relation $`M^TJM=J`$ where
$$J=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)$$
and hence, using the expression $`\mathrm{\Theta }^TJ\mathrm{\Theta }=J`$ we can show that the line element
$$dS^2=\frac{1}{2}\mathrm{Tr}(JdMJdM)=d\varphi ^2+e^{2\varphi }d\sigma ^2$$
(6)
is invariant under the global SL(2,R) transformation $`M\mathrm{\Theta }M\mathrm{\Theta }^T`$. Thus the effective action Eq. (3) written in the form
$$S=\frac{1}{2\kappa ^2}\mathrm{d}^4x\sqrt{g}\left[R+\frac{1}{4}\mathrm{Tr}\left(MM^1\right)\right]$$
(7)
is manifestly invariant under global SL(2,R) transformations. More generally the action for scalar fields parameterising an SL($`n`$,R)/SO($`n`$) non-linear sigma model can be written as
$$S=\frac{1}{2\kappa ^2}\mathrm{d}^4x\sqrt{g}\left[R+\frac{1}{4}\mathrm{Tr}\left(U_nU_n^1\right)\right],$$
(8)
where $`U_n`$ is a symmetric SL($`n`$,R) matrix. This action is invariant under the global transformation $`U_n\mathrm{\Theta }U_n\mathrm{\Theta }^T`$ where $`\mathrm{\Theta }`$ is a member of SL($`n`$,R).
We can build a symmetric SL(3,R) matrix $`U`$, from the SL(2,R) matrix $`M`$ in the following way :
$$U=\left(\begin{array}{cc}e^\nu M\hfill & e^\nu M\sigma \hfill \\ e^\nu \sigma ^TM\hfill & e^{2\nu }+e^\nu \sigma ^TM\sigma \hfill \end{array}\right),$$
(9)
where $`\nu `$ is a modulus field and the two additional degrees of freedom are the components of the $`2\times 1`$ vector
$$\sigma =\left(\begin{array}{c}\sigma _2\\ \sigma _3\end{array}\right).$$
(10)
The SL(3,R)-invariant trace of the $`3\times 3`$ matrix $`UU^1`$ which appears in the effective action is
$`\mathrm{Tr}\left(UU^1\right)`$ $`=`$ $`\mathrm{Tr}\left(MM^1\right)`$ (12)
$`6\left(\nu \right)^22\mathrm{T}\mathrm{r}\left(e^{3\nu }\sigma ^TM\sigma \right).`$
More generally, the same method can be used to construct an SL($`n+1`$,R) matrix $`U_{n+1}`$ from an SL($`n`$,R) matrix $`U_n`$, where
$$U_{n+1}=\left(\begin{array}{cc}e^{v_{n+1}}U_n\hfill & e^{v_{n+1}}U_n\sigma _n\hfill \\ e^{v_{n+1}}\sigma _{n}^{}{}_{}{}^{T}U_n\hfill & e^{nv_{n+1}}+e^{v_{n+1}}\sigma _{n}^{}{}_{}{}^{T}U_n\sigma _n\hfill \end{array}\right).$$
(13)
In addition to the fields in the SL($`n`$,R) matrix $`U_n`$, the SL($`n+1`$,R) matrix $`U_{n+1}`$ includes an additional modulus field $`v_{n+1}`$ and $`n`$ additional axion fields contained in the $`n\times 1`$ vector $`\sigma _n`$. This is sufficient to define our representation of the SL($`n`$,R)/SO($`n`$) coset starting from $`U_1=1`$. Thus the SL($`n`$,R) matrix $`U_n`$ contains $`n1`$ moduli and $`n(n1)/2`$ axions in total.
The SL(n+1,R)-invariant trace of the $`(n+1)\times (n+1)`$ matrix $`U_{n+1}U_{n+1}^{}{}_{}{}^{1}`$ which appears in the effective action \[see Eq. (8)\] can be calculated iteratively as
$`\mathrm{Tr}\left(U_{n+1}U_{n+1}^{}{}_{}{}^{1}\right)`$ $`=`$ $`\mathrm{Tr}\left(U_nU_{n}^{}{}_{}{}^{1}\right)`$ (16)
$`+(n^2n)\left(v_{n+1}\right)^2`$
$`2\mathrm{T}\mathrm{r}\left(e^{(n+1)v_{n+1}}\sigma _n^TU_n\sigma _n\right).`$
## III SL(4,R) Dilaton-Moduli-Vacuum Cosmologies
We first investigate the homogeneous dilaton-moduli vacuum solutions where we set $`\sigma _i=`$constant. The form of the dilaton-moduli solutions are invariant under a constant shift of the axion fields, so without loss of generality we set $`\sigma _i=0`$. This amounts to setting to zero the off-diagonal terms in the SL($`n`$,R) matrix. Thus we have the vacuum SL(2,R) matrix (putting $`\varphi =v_2`$)
$$U{}_{2}{}^{}{}_{}{}^{(0)}=\left(\begin{array}{cc}e^{v_2}& 0\\ 0& e^{v_2}\end{array}\right)$$
(17)
and the SL(3,R) matrix given in Eq. (9) becomes
$$U{}_{3}{}^{}{}_{}{}^{(0)}=\left(\begin{array}{ccc}e^{v_3+v_2}& 0& 0\\ 0& e^{v_3v_2}& 0\\ 0& 0& e^{2v_3}\end{array}\right)$$
(18)
while, from Eq. (13), an SL(n+1,R) matrix is given in terms of an SL($`n`$,R) matrix as
$$U_{n+1}^{(0)}=\left(\begin{array}{cc}e^{v_{n+1}}U_n^{(0)}\hfill & 0\hfill \\ 0\hfill & e^{nv_{n+1}}\hfill \end{array}\right).$$
(19)
The trace of $`U_n^{(0)}U_{n}^{(0)}{}_{}{}^{1}`$ for a vacuum SL($`n`$,R) matrix, using Eq. (16), is
$$\mathrm{Tr}\left(U_n^{(0)}U_{n}^{(0)}{}_{}{}^{1}\right)=\underset{i}{\overset{n}{}}i(i1)\left(v_i\right)^2.$$
(20)
Substituting this into the SL($`n`$,R)-invariant action, Eq. (8), we can now write down the effective action for vacuum SL(4,R)
$`S`$ $`=`$ $`{\displaystyle \frac{1}{2\kappa ^2}}{\displaystyle \mathrm{d}^4x\sqrt{g}}`$ (22)
$`\times \left[R3(v_4)^2{\displaystyle \frac{3}{2}}(v_3)^2{\displaystyle \frac{1}{2}}(v_2)^2\right].`$
We will assume that the external four dimensional spacetime is a spatially-flat FRW metric, with the line element
$$ds^2=a^2(\eta )\left(\mathrm{d}\eta ^2+\mathrm{d}x^2+\mathrm{d}y^2+\mathrm{d}z^2\right)$$
(23)
and scale factor, $`a(\eta )`$, where $`\eta `$ is the conformal time. The effective action in Eq. (22) can then be written for homogeneous dilaton-moduli fields as
$`S`$ $`=`$ $`{\displaystyle \frac{1}{2\kappa ^2}}{\displaystyle \mathrm{d}^3xd\eta }`$ (25)
$`\left[6a_{}^{}{}_{}{}^{2}3a^2v_{4}^{}{}_{}{}^{}{}_{}{}^{2}{\displaystyle \frac{3}{2}}a^2v_{3}^{}{}_{}{}^{}{}_{}{}^{2}{\displaystyle \frac{1}{2}}a^2v_{2}^{}{}_{}{}^{}{}_{}{}^{2}\right],`$
where a prime denotes differentiation with respect to $`\eta `$. We use the Euler-Lagrange equations derived from this action to calculate the equations of motion related to this action. The resulting evolution equations for the dilaton-moduli fields in a FRW metric are
$`12{\displaystyle \frac{a^{\prime \prime }}{a}}`$ $`=`$ $`6v_{4}^{}{}_{}{}^{}{}_{}{}^{2}3v_{3}^{}{}_{}{}^{}{}_{}{}^{2}v_{2}^{}{}_{}{}^{}{}_{}{}^{2},`$ (26)
$`v_{4}^{}{}_{}{}^{\prime \prime }+2{\displaystyle \frac{a^{}}{a}}v_{4}^{}{}_{}{}^{}`$ $`=`$ $`0,`$ (27)
$`v_{3}^{}{}_{}{}^{\prime \prime }+2{\displaystyle \frac{a^{}}{a}}v_{3}^{}{}_{}{}^{}`$ $`=`$ $`0,`$ (28)
$`v_{2}^{}{}_{}{}^{\prime \prime }+{\displaystyle \frac{a^{}}{a}}v_{2}^{}{}_{}{}^{}`$ $`=`$ $`0,`$ (29)
subject to the constraint equation (from general relativistic invariance under time-reparameterisation)
$$\left(\frac{a^{}}{a}\right)^2=\frac{1}{12}\left(6v_{4}^{}{}_{}{}^{}{}_{}{}^{2}+3v_{3}^{}{}_{}{}^{}{}_{}{}^{2}+v_{2}^{}{}_{}{}^{}{}_{}{}^{2}\right).$$
(30)
Substituting Eq. (30) into Eq. (26) yields a second-order equation for $`a(\eta )`$ which can be integrated twice to obtain
$$a=a_{}|\eta |^{\frac{1}{2}},$$
(31)
where $`a_{}`$ is one integration constant, and we have used our freedom to choose $`a=0`$ as the origin for the time coordinate $`\eta `$ in order to eliminate the other constant of integration. Equation (31) is the standard solution in the Einstein frame for a spatially flat FRW cosmology with free scalar fields. In terms of the proper time $`ta𝑑\eta `$ we have $`a|t|^{1/3}`$, which describes a non-accelerating expanding universe for $`t>0`$, or an accelerating contraction for $`t<0`$. It is this phase of accelerated contraction that is the basis of the pre-big-bang scenario .
Equations (27) to (29) can now be integrated using the solution for $`a(\eta )`$ in Eq. (31) to give
$$v_i^{}=\frac{C_i}{a^2},$$
(32)
where $`C_i`$ are constants of integration, $`i=1\mathrm{}3`$. This allows the constraint equation (30) to be rewritten in the form
$$\frac{C_1^2}{3}+C_2^2+2C_3^2=1.$$
(33)
The constants of integration $`C_1/\sqrt{3}`$, $`C_2`$ and $`\sqrt{2}C_3`$ can be interpreted as points on the surface of a sphere. It is therefore convenient to move to spherical coordinates
$`C_1/\sqrt{3}`$ $`=`$ $`\mathrm{cos}\xi _1,`$ (34)
$`C_2`$ $`=`$ $`\mathrm{sin}\xi _1\mathrm{cos}\xi _2,`$ (35)
$`\sqrt{2}C_3`$ $`=`$ $`\mathrm{sin}\xi _1\mathrm{sin}\xi _2,`$ (36)
where the constraint is automatically satisfied and the new constants of integration are $`0\xi _1\pi `$ and $`0\xi _2<2\pi `$.
Using Eq. (31) in Eqs. (32) gives monotonic power law solutions for $`v_4`$, $`v_3`$ and $`v_2`$:
$`e^{v_4}`$ $`=`$ $`e^{v_4}|\eta |^{\mathrm{sin}\xi _1\mathrm{sin}\xi _2/\sqrt{2}},`$ (37)
$`e^{v_3}`$ $`=`$ $`e^{v_3}|\eta |^{\mathrm{sin}\xi _1\mathrm{cos}\xi _2},`$ (38)
$`e^{v_2}`$ $`=`$ $`e^{v_2}|\eta |^{\sqrt{3}\mathrm{cos}\xi _1},`$ (39)
where $`v_4`$, $`v_3`$ and $`v_2`$ are constants of integration. Dilaton-moduli-vacuum solutions related by SL(3,R) transformations are recovered by setting $`\xi _2=0`$, while for $`\xi _2=0`$ and $`\xi _1=0`$ or $`\pi `$ we recover dilaton-vacuum cosmological solutions related by SL(2,R) transformations .
## IV Axion Perturbations
It is known that inhomogeneous linear perturbations in the dilaton and other moduli fields about the homogeneous FRW solutions given by Eq. (31) have the general solution $`\delta v_i=Z_0(|k\eta |)`$, where $`Z_0`$ is any linear combination of Bessel functions of order zero, independent of the various integration constants that appear in the solutions in Eq. (37). In the pre-big-bang scenario this solution inevitably leads to a cosmological spectrum of vacuum fluctuations steeply tilted towards small scales, with essentially no perturbations on super-horizon scales . If the pre big bang scenario is to produce any observable perturbations on large scales it must be through vacuum fluctuations in the axion fields which are non-minimally coupled to the dilaton-moduli fields in the four-dimensional Einstein frame and hence sensitive to the integration constants that parameterise their evolution.
Hence we will investigate inhomogeneous axion field perturbations about the dilaton-moduli vacuum solutions. We will calculate the field equations for these linear perturbations including axion fields by constructing the effective action to second-order in the perturbations<sup>*</sup><sup>*</sup>*Because $`\sigma _i^{}=0`$ in the background solution there is no metric back-reaction to lowest-order and the perturbations are automatically gauge-invariant .. This is constructed iteratively using Eq. (20) as
$`\mathrm{Tr}\left(U_{n+1}U_{n+1}^{}{}_{}{}^{1}\right)`$ $`=`$ $`\mathrm{Tr}\left(U_nU_{n}^{}{}_{}{}^{1}\right)`$ (42)
$`+(n^2n)\left(v_{n+1}\right)^2`$
$`2\mathrm{T}\mathrm{r}(e^{(n+1)v_{n+1}}\sigma _n^TU{}_{n}{}^{}{}_{}{}^{(0)}\sigma _n),`$
where we can use the vacuum solution, $`U_n^{(0)}`$, in the last term of the trace equation, in order to calculate the action to second-order in the axion fields.
Expressions can be constructed from the initial SL(2,R) matrix Eq. (5) and we can therefore write the trace of $`U_3U_{3}^{}{}_{}{}^{1}`$ for SL(3,R) to second order in the axion perturbations as
$`\mathrm{Tr}(U_3U{}_{3}{}^{}{}_{}{}^{1})`$ $`=`$ $`\mathrm{Tr}(U_2U{}_{2}{}^{}{}_{}{}^{1})6(v_2)^2`$ (44)
$`2\mathrm{T}\mathrm{r}(e^{3v_2}\sigma ^TU{}_{2}{}^{}{}_{}{}^{(0)}\sigma ).`$
Rewriting Eq. (5) as
$$U_2=\left(\begin{array}{cc}e^{v_2}\hfill & \sigma _1e^{v_2}\hfill \\ \sigma _1e^{v_2}\hfill & e^{v_2}+\sigma _1^2e^{v_2}\hfill \end{array}\right)$$
(45)
and then adding on the second order axion term with the vacuum case given in Eq. (20) in the above expression we obtain
$`\mathrm{Tr}\left(U_3U_{3}^{}{}_{}{}^{1}\right)`$ $`=`$ $`6(v_3)^22(v_2)^22e^{2v_2}(\sigma _1)^2`$ (47)
$`2e^{3v_3+v_2}(\sigma _2)^22e^{3v_3v_2}(\sigma _3)^2,`$
while for SL(4,R) we have
$`\mathrm{Tr}\left(U_4U_{4}^{}{}_{}{}^{1}\right)`$ $`=`$ (51)
$`12(v_4)^26(v_3)^22(v_2)^22e^{2v_2}(\sigma _1)^2`$
$`2e^{3v_3+v_2}(\sigma _2)^22e^{3v_3v_2}(\sigma _3)^22e^{4v_4+v_3+v_2}(\sigma _4)^2`$
$`2e^{4v_4+v_3v_2}(\sigma _5)^22e^{4v_42v_3}(\sigma _6)^2.`$
In the Einstein frame the action becomes
$`S`$ $`=`$ $`{\displaystyle \frac{1}{2\kappa ^2}}{\displaystyle }\mathrm{d}^3x{\displaystyle }\mathrm{d}\eta [6a_{}^{}{}_{}{}^{2}3a^2v_{4}^{}{}_{}{}^{}{}_{}{}^{2}{\displaystyle \frac{3}{2}}a^2v_{3}^{}{}_{}{}^{}{}_{}{}^{2}`$ (55)
$`{\displaystyle \frac{1}{2}}a^2v_{2}^{}{}_{}{}^{}{}_{}{}^{2}{\displaystyle \frac{1}{2}}a^2e^{2v_2}\sigma _{1}^{}{}_{}{}^{}{}_{}{}^{2}{\displaystyle \frac{1}{2}}a^2e^{3v_3+v_2}\sigma _{2}^{}{}_{}{}^{}{}_{}{}^{2}`$
$`{\displaystyle \frac{1}{2}}a^2e^{3v_3v_2}\sigma _{3}^{}{}_{}{}^{}{}_{}{}^{2}{\displaystyle \frac{1}{2}}a^2e^{4v_4+v_3+v_2}\sigma _{4}^{}{}_{}{}^{}{}_{}{}^{2}`$
$`{\displaystyle \frac{1}{2}}a^2e^{4v_4+v_3v_2}\sigma _{5}^{}{}_{}{}^{}{}_{}{}^{2}{\displaystyle \frac{1}{2}}a^2e^{4v_42v_3}\sigma _{6}^{}{}_{}{}^{}{}_{}{}^{2}].`$
We can now use this to derive the field equations for inhomogeneous axion perturbations evolving in the homogeneous vacuum background solutions
$`\delta \sigma _{1}^{}{}_{}{}^{\prime \prime }+\left(2{\displaystyle \frac{a^{}}{a}}+2v_{2}^{}{}_{}{}^{}\right)\delta \sigma _{1}^{}{}_{}{}^{}+k^2\delta \sigma _1`$ $`=`$ $`0,`$ (56)
$`\delta \sigma _{2}^{}{}_{}{}^{\prime \prime }+\left(2{\displaystyle \frac{a^{}}{a}}+3v_{3}^{}{}_{}{}^{}+v_{2}^{}{}_{}{}^{}\right)\delta \sigma _{2}^{}{}_{}{}^{}+k^2\delta \sigma _2`$ $`=`$ $`0,`$ (57)
$`\delta \sigma _{3}^{}{}_{}{}^{\prime \prime }+\left(2{\displaystyle \frac{a^{}}{a}}+3v_{3}^{}{}_{}{}^{}v_{2}^{}{}_{}{}^{}\right)\delta \sigma _{3}^{}{}_{}{}^{}+k^2\delta \sigma _3`$ $`=`$ $`0,`$ (58)
$`\delta \sigma _{4}^{}{}_{}{}^{\prime \prime }+\left(2{\displaystyle \frac{a^{}}{a}}+4v_{4}^{}{}_{}{}^{}+v_{3}^{}{}_{}{}^{}+v_{2}^{}{}_{}{}^{}\right)\delta \sigma _{4}^{}{}_{}{}^{}+k^2\delta \sigma _4`$ $`=`$ $`0,`$ (59)
$`\delta \sigma _{5}^{}{}_{}{}^{\prime \prime }+\left(2{\displaystyle \frac{a^{}}{a}}+4v_{4}^{}{}_{}{}^{}+v_{3}^{}{}_{}{}^{}v_{2}^{}{}_{}{}^{}\right)\delta \sigma _{5}^{}{}_{}{}^{}+k^2\delta \sigma _5`$ $`=`$ $`0,`$ (60)
$`\delta \sigma _{6}^{}{}_{}{}^{\prime \prime }+\left(2{\displaystyle \frac{a^{}}{a}}+4v_{4}^{}{}_{}{}^{}2v_{3}^{}{}_{}{}^{}\right)\delta \sigma _{6}^{}{}_{}{}^{}+k^2\delta \sigma _6`$ $`=`$ $`0,`$ (61)
where $`k`$ is the comoving wavenumber. Note that because the axion field is zero in the vacuum background solution, their back-reaction upon the moduli fields and the four-dimensional spacetime metric vanishes to first-order and the perturbations are independent of the spacetime gauge .
The field equations (56) to (56) can be written in the standard form for a free scalar field evolving in a FRW metric
$$\delta \sigma _{i}^{}{}_{}{}^{\prime \prime }+2\frac{\overline{a}_i^{}}{\overline{a}_i}\delta \sigma _{i}^{}{}_{}{}^{}+k^2\delta \sigma _i=0$$
(62)
with $`i=1\mathrm{}6`$, where we introduce the conformally rescaled scale factor, $`\overline{a}_i`$, in the corresponding “axion frame”
$`\overline{a}_1`$ $`=`$ $`e^{v_2}a,`$ (63)
$`\overline{a}_2`$ $`=`$ $`e^{(3v_3+v_2)/2}a,`$ (64)
$`\overline{a}_3`$ $`=`$ $`e^{(3v_3v_2)/2}a,`$ (65)
$`\overline{a}_4`$ $`=`$ $`e^{2v_4+(v_3+v_2)/2}a,`$ (66)
$`\overline{a}_5`$ $`=`$ $`e^{2v_4+(v_3v_2)/2}a,`$ (67)
$`\overline{a}_6`$ $`=`$ $`e^{2v_4v_3}a.`$ (68)
Substituting the solutions from Eqs. (31) and (38) into Eqs. (63)–(68) we obtain
$$\overline{a}_i|\eta |^{p_i},$$
(69)
where
$`p_1`$ $`=`$ $`\sqrt{3}\mathrm{cos}\xi _1,`$ (70)
$`p_2`$ $`=`$ $`\sqrt{3}\left({\displaystyle \frac{\sqrt{3}}{2}}\mathrm{sin}\xi _1\mathrm{cos}\xi _2+{\displaystyle \frac{1}{2}}\mathrm{cos}\xi _1\right),`$ (71)
$`p_3`$ $`=`$ $`\sqrt{3}\left({\displaystyle \frac{\sqrt{3}}{2}}\mathrm{sin}\xi _1\mathrm{cos}\xi _2{\displaystyle \frac{1}{2}}\mathrm{cos}\xi _1\right),`$ (72)
$`p_4`$ $`=`$ $`\sqrt{3}\left[{\displaystyle \frac{\sqrt{3}}{2}}\mathrm{sin}\xi _1\left({\displaystyle \frac{2\sqrt{2}}{3}}\mathrm{sin}\xi _2+{\displaystyle \frac{1}{3}}\mathrm{cos}\xi _2\right)+{\displaystyle \frac{1}{2}}\mathrm{cos}\xi _1\right],`$ (73)
$`p_5`$ $`=`$ $`\sqrt{3}\left[{\displaystyle \frac{\sqrt{3}}{2}}\mathrm{sin}\xi _1\left({\displaystyle \frac{2\sqrt{2}}{3}}\mathrm{sin}\xi _2+{\displaystyle \frac{1}{3}}\mathrm{cos}\xi _2\right){\displaystyle \frac{1}{2}}\mathrm{cos}\xi _1\right],`$ (74)
$`p_6`$ $`=`$ $`\sqrt{3}\mathrm{sin}\xi _1\left({\displaystyle \frac{\sqrt{2}}{\sqrt{3}}}\mathrm{sin}\xi _2{\displaystyle \frac{1}{\sqrt{3}}}\mathrm{cos}\xi _2\right).`$ (75)
Equation (62) has the standard form for perturbations of a free scalar field evolving in an FRW cosmology with scale factor $`\overline{a}_i`$. Thus we can define the canonically normalised variables
$$u_i=\frac{1}{\sqrt{2}\kappa }\overline{a}_i\delta \sigma _i,$$
(76)
which enables us to rearrange Eq. (62) in the form of a simple harmonic oscillator with time-dependent mass
$$u_{i}^{}{}_{}{}^{\prime \prime }+\left(k^2+\frac{(\overline{a}_{i}^{}{}_{}{}^{\prime \prime })}{\overline{a}_i}\right)u_i=0.$$
(77)
For a power-law expansion given by Eq. (69) this corresponds to Bessel’s equation
$$u_{i}^{}{}_{}{}^{\prime \prime }+\left(k^2+\frac{(1/4p{}_{i}{}^{}{}_{}{}^{2})}{|\eta |^2}\right)u_i=0.$$
(78)
Equation (78) has a standard solution
$$u_i=|k\eta |^{1/2}Z_{p_i}\left(|k\eta |\right),$$
(79)
where $`Z_{p_i}`$ is any linear combination of Bessel (or Hankel) functions and for each axion field the order of the Bessel function is $`p_i`$ given in Eq. (70).
## V Pre-Big-Bang Perturbation Spectra
We will write the solutions of the Bessel equation in terms of Hankel functions so that a general solution of Eq. (78) is given by
$$u_i=|k\eta |^{1/2}\left[u_+H_{|p_i|}^{}{}_{}{}^{(1)}\left(|k\eta |\right)+u_{}H_{|p_i|}^{}{}_{}{}^{(2)}\left(|k\eta |\right)\right].$$
(80)
For wavelengths much smaller than the horizon scale ($`|k\eta |1`$) the equation of motion Eq. (77) reduces to that for a free-scalar field $`u_i`$ in flat Minkowski spacetime with a well-defined vacuum state. Allowing only positive frequency modes in a flat-spacetime vacuum state requires that
$$u_i\frac{e^{ik\eta }}{\sqrt{2k}}.$$
(81)
The classic horizon problem of the standard hot big bang is that all modes start outside the horizon at the big bang ($`k\eta 0`$) and so there is no reason to expect modes to start in this vacuum state.
However in the pre-big-bang scenario all modes start within the horizon in the infinite past as $`\eta \mathrm{}`$. As $`\eta 0_{}`$ on the $`(+)`$ branch modes leave the horizon, giving a well-defined spectrum of super-horizon perturbations in all fields, even though there is no inflation in the conventional sense ($`\ddot{a}>0`$) in the Einstein frame.
Allowing only positive frequency modes in a flat-spacetime vacuum state at early times (as $`k\eta \mathrm{}`$) i.e. large $`\eta `$ yields
$$u_+=\frac{\sqrt{\pi }}{2\sqrt{k}}e^{i(2|p_i|+1)\pi /4},u_{}=0.$$
(82)
Therefore using Eq. (76) we can write
$$\delta \sigma _i=\kappa \sqrt{\frac{\pi }{2k}}e^{i(2|p_i|+1)\pi /4}\frac{\sqrt{k\eta }}{\overline{a}_i}H_{|p_i|}^{}{}_{}{}^{(1)}\left(k\eta \right).$$
(83)
At late times on super-horizon scales ($`|k\eta |1`$) we have
$$\delta \sigma _i=\pm i\kappa \sqrt{\frac{\pi }{k}}e^{i(2|p_i|+1)\pi /4}\left(\frac{\mathrm{\Gamma }(|p_i|)}{\pi \overline{a}_i}\right)\left(\frac{2}{k\eta }\right)^{|p_i|(1/2)}.$$
(84)
The power spectrum for these axion perturbations is denoted by
$$𝒫_{\delta \sigma _i}\frac{k^3}{2\pi ^2}|\delta \sigma _i|^2$$
(85)
which represents the dispersion $`\delta \sigma _i^2`$ due to fluctuations on comoving scales $`k^1`$ . It can be calculated from Eq. (84) to be
$$𝒫_{\delta \sigma _i}=2\kappa ^2\left(\frac{C(|p_i|)}{2\pi }\right)^2\frac{k^2}{\overline{a}_{i}^{}{}_{}{}^{2}}(k\eta )^{12|p_i|},$$
(86)
where the coefficient term
$$C(|p_i|)=\frac{2^{|p_i|}\mathrm{\Gamma }(|p_i|)}{2^{3/2}\mathrm{\Gamma }(3/2)},$$
(87)
approaches unity for $`|p_i|=3/2`$.
The spectral tilt of the perturbation spectra is defined by
$$\mathrm{\Delta }n_i\frac{\mathrm{d}\mathrm{ln}𝒫_{\delta \sigma _i}}{\mathrm{d}\mathrm{ln}k}.$$
(88)
It follows from Eq. (86) that the spectral tilt for each of the axion fields is constant and takes the values
$$\mathrm{\Delta }n_i=32|p_i|,$$
(89)
where the $`p_i`$ are given in Eq. (70) in terms of the two integration constants $`\xi _1`$ and $`\xi _2`$ which parameterise the dilaton-moduli vacuum background solutions.
From Eqs. (89) and (70) it is then possible to find the spectral tilts as functions of the integration constants $`\xi _1`$ and $`\xi _2`$. (See Figs. 1 and 2.) The maximum absolute value for any $`p_i`$ is $`\sqrt{3}`$ and thus the minimum value of the spectral tilt for any axion field is $`\mathrm{\Delta }n=2\sqrt{3}+3=0.46`$ as found previously with SL(2,R) or SL(3,R) symmetry groups. For any axion field the allowed range for the spectral tilt is
$$2\sqrt{3}+3\mathrm{\Delta }n_i3.$$
(90)
## VI Discussion
The amplitude of the axion power spectra at the end of the pre-big bang phase when the comoving horizon scale is given by $`\eta _s=1/k_s`$ can be given, from Eqs. (86) and (89), as
$$𝒫_{\delta \sigma _i}|_s=2\kappa ^2\left(\frac{C(|p_i|)}{p_i+1/2}\right)^2\left(\frac{\overline{H_i}}{2\pi }\right)_s^2\left(\frac{k}{k_s}\right)^{\mathrm{\Delta }n_i},$$
(91)
where the expression for $`C(|p_i|)`$ is given in Eq. (87) and the Hubble rate in the axion frame is given by
$$\overline{H_i}^2=\frac{(p_i+1/2)^2}{\overline{a}_i^2\eta ^2}.$$
(92)
The Hubble rate in the axion frame can also be related to $`H`$, the expansion rate in the Einstein frame, via
$$\overline{H_i}=\frac{2(p_i+1/2)}{\mathrm{\Omega }_i}H.$$
(93)
Assuming that modes remain frozen-in on large scales $`(|k\eta |1)`$ during the uncertain transition from pre to post big bang phase , then the energy density contributed by massless axion perturbations in the Einstein frame can be estimated as
$`\rho _i`$ $``$ $`\mathrm{\Omega }{}_{i}{}^{}{}_{}{}^{2}{\displaystyle \frac{k^2}{a^2}}{\displaystyle \frac{𝒫_{\delta \sigma _i}|_s}{2\kappa ^2}}`$ (94)
$``$ $`{\displaystyle \frac{k^2}{a^2}}C^2(|p_i|)\left({\displaystyle \frac{H_s}{2\pi }}\right)^2\left({\displaystyle \frac{k}{k_s}}\right)^{\mathrm{\Delta }n_i}.`$ (95)
Although massless axion fields never come to dominate the total energy density in the universe, they do contribute a perturbation to the density when a given scale re-enters the cosmological horizon
$$\frac{\delta \rho }{\rho }=\frac{\rho _i}{\rho _{\mathrm{total}}}=\frac{C^2(|p_i|)}{3}\left(\frac{\kappa H_s}{2\pi }\right)_s^2\left(\frac{k}{k_s}\right)^{\mathrm{\Delta }n_i},$$
(96)
where the total energy density during the radiation or matter dominated era is given by $`\rho _{\mathrm{total}}=3H^2/\kappa ^2`$. An analysis of the effect of this perturbation to the overall density during the radiation and matter dominated eras, and hence to the spectrum of anisotropies in the cosmic microwave background, is calculated in Refs. . Note that all the different axion fields have the same amplitude for modes crossing the Hubble scale at the start of the post big bang era, $`kk_s`$, set by the ratio of the maximum Hubble rate at the end of the pre big bang to the Planck scale $`\kappa ^2H_s^2`$. In the simplest scenario of a sudden transition to the post big bang phase where the dilaton and moduli are fixed, this gives $`\kappa ^2H_s^2e^\varphi 10^2`$ yielding the expression given in Eq. (1). On larger scales the amplitudes of the different axion fields is then fixed entirely by the different spectral tilts, with the field with the minimum tilt yielding the largest density perturbation. If $`\mathrm{\Delta }n_i<0`$ for any of the axion fields then the density perturbation given by Eq. (96) diverges on large scales.
A key question then is whether it is possible to maintain blue spectral indices ($`\mathrm{\Delta }n_i>0`$) for all the axion fields for any values of $`\xi _1`$ and $`\xi _2`$. Figure 1 shows the spectral indices for solutions restricted to the SL(3,R) dilaton-moduli-vacuum solutions where $`\xi _2=0`$, in which case it has previously been shown that there are no values of $`\xi _1`$ for which $`\mathrm{\Delta }n_i>0`$ for all $`i`$.
We define $`\mathrm{\Delta }n_{}\mathrm{min}(\mathrm{\Delta }n_1,\mathrm{\Delta }n_2,\mathrm{},\mathrm{\Delta }n_6)`$ to be the smallest of all the spectral indices for a given value of $`(\xi _1,\xi _2)`$. Figure 3 shows that there is a finite region of parameter space for which $`\mathrm{\Delta }n_{}>0`$. We find the maximum value of $`\mathrm{\Delta }n_{}`$ occurs for $`\xi _1=\pi /4`$ or $`3\pi /4`$ and $`\mathrm{tan}\xi _2=\sqrt{2}`$ (shown in Fig. 2 where we obtain $`\mathrm{\Delta }n_{}=3\sqrt{6}=0.55`$). Table 1 shows the previously calculated maximum values for $`\mathrm{\Delta }n_{}`$ for other SL($`n`$,R) groups. The maximum allowed spectral tilt gets progressively larger (bluer) as the symmetry group gets larger. Since the larger groups always contain the previous groups we see that it is indeed possible for all the axion fields to have blue spectral indices in SL($`n`$,R) non-linear sigma models where $`n4`$.
On the other hand requiring all the spectral indices to be positive represents a restriction on the allowed initial conditions.
We have focused our discussion on the cosmological perturbations induced by an effectively massless axion field proposed in Ref. , in which case the primordial axion perturbation spectra can be directly related to seed density perturbations. If the axions become non-relativistic during the radiation-dominated era then the overall amplitude of perturbations is increased by a factor $`(m/H_{\mathrm{eq}})^{1/2}`$ , where the Hubble rate at matter-radiation equality is $`H_{\mathrm{eq}}10^{27}`$eV. If any of the primordial axion spectra have a negative spectral tilt, $`\mathrm{\Delta }n_i<0`$ (as occurs in a large regime of parameter space shown in Fig. 3) then the only way to suppress the large-scale density perturbation appears to be for the axion to acquire a periodic potential where large field fluctuations may have a small effect on the energy density . In such a scenario, the massive axion could be a novel form of dark matter and lead to a very different model for large-scale structure formation.
Finally, we note that the presence of non-trivial background axion fields in the general solutions of the full SL(2,R) and SL(3,R) symmetric axion-dilaton-moduli cosmologies restricts the allowed asymptotic vacuum state . The effect of the background axion fields upon the perturbation spectra has only been calculated for a single axion field where the perturbation spectra remain invariant under SL(2,R) transformations of the background solutions . The cosmological perturbation spectra in axion-dilaton-moduli solutions with more degrees of freedom remain to be investigated.
## Acknowledgements
The authors are grateful to Carlo Ungarelli for helpful comments. H.B. is supported by the EPSRC and D.W. is supported by the Royal Society. |
warning/0002/cond-mat0002356.html | ar5iv | text | # Commensurate-Incommensurate Phase Transitions for Multi-Chain Quantum Spin Models: Exact Results
## 1 Introduction
There has recently been considerable interest on low-dimensional quantum correlated spin and electron systems. These systems, especially one-dimensional (1D), manifest the specific features of, e.g., magnetic behavior at low temperatures, which are absent for the standard, conventional 3D magnetic systems. Spin systems usually manifest 1D behavior for the temperatures higher than the temperature of the 3D magnetic ordering, but lower than the maximum characteristic energy of the interaction between spins, i.e. in our case the intra-chain spin-spin coupling. The origin of such specific features is the enhancement of the quantum fluctuations of the 1D systems due to the peculiarities of the 1D density of states together with the quantum nature of spins.
Moreover, during the last decade a large number of new quasi 1D spin compounds were created and studied experimentally. These compounds manifest at low temperatures the properties of a single or several quantum spin chains weakly coupled to each other . It is strongly believed that this class of compounds will provide the new information on the transition from 1D to 2D in quantum many-body physics. It is very important, because the 2D quantum many-body physics has been a challenge for both theorists and experimentalists since the beginning of the study of low dimensional quantum systems. On the other hand, the advantage of the 1D theoretical studies is the possibility of obtaining exact solutions by using non-perturbative methods, which are difficult to apply for the higher-dimensional quantum many-body models. The results of the exact calculations of the 1D models can serve as testing grounds for the use of perturbative and numerical methods in more realistic situations.
Recently several exactly solvable models have been introduced, in which the zigzag-like interaction between two quantum spin chains was studied exactly using the Bethe ansatz technique . This method is widely known by now, see e.g., the recent monography and references therein. The Bethe ansatz method permits exact calculation of the static characteristics of quantum many-body systems, such as the groundstate behavior, the influence of an external magnetic field, and the thermodynamic features of e.g., the temperature dependencies of the specific heat, magnetic susceptibility, etc. These results should apply to more realistic systems, but it is not obvious how the interactions between the chains modify the answers. The mean-field like approximations for the inter-chain couplings are not sufficient, because the mean field approach in any version already implies the existence of the (sometimes hidden) order parameter. It is, unfortunately, also unclear whether the numerical calculations, which can be directly applied for the quantum many-body systems of very small sizes by now (say, at most several tens of sites) describe well the properties of the real systems, in which, even in quasi-1D ones, the number of sites is at least of order of 10<sup>8</sup> or higher. On the other hand, it must be admitted that some features of the exactly solvable 1D models are far from what is observed experimentally, but these unrealistic features of the 1D models are known and simple to recognize.
The behavior of the multi-chain spin systems in an external magnetic field is especially interesting, see e.g., because of (i) the possibility of experimental observations due to recent progress in the high magnetic field measurements, and (ii) very interesting theoreticallly predictable effects which are possible to recognize in experiments, such as phase transitions in the external magnetic field. However several important issues are far from being solved in the quantum two-chain spin models. For example, there are three questions that need to be answered: (1) Are the properties of those exactly solvable two-chain spin models unique or it is possible to say something about the more general class of two-chain quantum spin models? (2) How are the multi-chain quantum models connected to the 2D many-body systems, i.e. what is the scenario of the transition from 1D to 2D when one increases the number of coupled chains while keeping the conditions of integrability? and (3) What will happen with the behavior of the non-integrable multi-chain spin models if one goes beyond the framwork of integrability i.e. adding some perturbations to the exactly solvable model? (For example, Ref. implies that namely the spin chirality, which separatly breaks the time-reversal and parity symmetries in the two-chain integrable model , is the reason for the emergence of the additional phase transitions in an external magnetic field for the two-chain spin $`\frac{1}{2}`$ model as compared to the single-chain system).
The goal of this paper is to answer these questions. First, we re-visit the exactly integrable two-chain spin $`\frac{1}{2}`$ model and show that the inclucion of the magnetic anisotropy of the “easy-plane” type, with which the system stays in the quantum critical region, will not drastically change the behavior in an external magnetic field but will shift the critical values of the magnetic fields and intra-chain couplings at which the phase transitions occur and will affect the critical exponents. We will show that these two-chain spin models share the most important features of the behavior in an external field with the wide class of the (1+1) quantum field theories. Next, we will introduce the higher-spin realizations of the two-chain spin models, e.g., investigating the important class of 1D two-chain quantum ferrimagnets with different spin values at the sites of each chain. We will also investigate the behavior of the multi-chain exactly solvable spin models in an external magnetic field and show how the additional phase transitions arising due to the increasing number of chains vanish in the quasi-2D limit. Finally, we will show how the relevant deviations from the integrability, e.g., the absence of the terms in the Hamiltonian which separately break the parity and time-reversal symmetries give rise to gaps in the spectra of low-lying excitations of the multi-chain quantum spin systems and we will calculate the scaling exponents for the gaps.
The paper is organized as follows. In Section 2 we re-visit the exactly solvable two-chain uniaxial spin model to remind the reader of the main steps of the Bethe ansatz. The investigations of Refs. of isotropic spin $`\frac{1}{2}`$ two-chain models are generalized in this section for the case of uniaxial magnetic anisotropy. The calculations in this section are rather simple, but we write them in detail because they provide the basis for the more nontrivial generalizations of this class of models, and will be used in the following Sections. In Section 3 we point out the similarities between the behavior of the uniaxial two-chain quantum spin models and a class of quantum field theories (QFT) in an external magnetic field, predicting new phases for the QFT. In Section 4 we introduce the $`SU(2)`$ generalization of the integrable two-chain model for higher values of the site spins (possibly different) in each chain, i.e., quantum ferrimagnet. We point out the similarities of the quantum ferrimagnet with the QFT with a nonzero Wess-Zumino term and predict new phases for the latter in an external magnetic field. We derive the integral eguations for the critical exponents. In Section 5 we consider the multi-chain quantum spin model and discuss how the external field behavior of the integrable multi-chain models is changed when the number of chains is increased while preserving the exact solvability. In Section 6 we briefly sketch how the deviations from integrability change the magnetic and low-temperature properties of this class of multi-chain quantum spin systems. The paper is closed with a discussion of the main results and some conclusions.
## 2 Two-chain uniaxial quantum spin model
A common property of some of the Bethe ansatz solutions is the presence of shifts $`\theta _j`$ of the spectral parameter $`\lambda `$ for the associated transfer matrix of an algebraic version of the Bethe ansatz (the Quantum Inverse Scattering Method, \[QISM\] ). Those shifts also appear in the Bethe ansatz equations (BAE) for the quantum numbers called rapidities, which parametrize the eigenfunctions and eigenvalues of the Hamiltonians. Hence, the distributions of the rapidities are also affected by the shifts. An interesting property is connected with those shifts: depending on their values and the external magnetic field, even for (quasi)particles of the same type, additional minima may appear in distributions of the rapidities. These additional minima also result in the nonmonotonic behavior of the dispersion laws of the low-lying excitations. Also, they provide additional Dirac seas for low-lying excitations, changing the structures of the physical groundstates of the models. These additional minima determine the special behavior of the models in the external magnetic field . In particular, the appearence of the new phases and new phase transitions is namely due to the emergence of these new minima in the distributions of the quantum numbers.
To set the stage, let us first remind the reader about the main steps of the QISM. The common feature of the Bethe ansatz solvable models is the factorization of a monodromy matrix (the ordered product of all two-particle scattering matrices, which depend on some spectral parameter) . Exact (Bethe ansatz) integrability requires exclusively elastic scattering between (quasi)particles. For such a theories two-particle scattering matrices and $`L`$-operators satisfy the Yang-Baxter relation . In turn, the factorization of the monodromy matrices garantees that they satisfy the Yang-Baxter equations, too. The transfer matrices of the associated statistical problem are traces over some additional, auxiliary subspace, of monodromy matrices . The most important feature of transfer matrices with different spectral parameters is their commutativity. The necessary and sufficient condition for this is the validity of the Yang-Baxter equations for two-particle scattering matrices and hence for monodromy matrices. The commutativity of transfer matrices implies that one can construct an infinite number of integrals of motion, which commute with one another and with the transfer matrix. Therefore the exact integrability is proved. Usually the structure of these integrals of motion is determined by their locality. For instance, the best known of series of integrals of motion is the series of derivatives with respect to the spectral parameter of the logarithm of a transfer matrix taken at some special value of former . Locality means that for the first derivative of the logarithm of the transfer matrix (usually called the Hamiltonian of the lattice system) only short-range particle-paricle interactions contribute.
In this paper we will see that namely the aforementioned shifts of the spectral parameters yield new phases in the groundstate behavior in an external magnetic field of a wide class of exactly solvable models, quantum spin multi-chain models and QFT. We will show that in the conformal limit these phases of the lattice models correspond to one Wess-Zumino-Witten (WZW) model or to several of them with dressed charges (proportional to the compactification radii) of scalar or matrix types for each of the phases, respectively.
Let us start with the form of the Bethe ansatz equations (BAE) for the set of rapidities $`\{u_\alpha \}_{\alpha =1}^M`$. In this paper we will concentrate only on the critical, “easy-plane” type of the magnetic anisotropy for the antiferromagnetic spin multi-chain models, $`0\gamma \pi /2`$ ($`\gamma =\pi /q`$, $`q`$ integer, parametrizing the magnetic anisotropy), and the repulsive interactions in QFT. This correspond to hyperbolic or rational solutions of the Yang-Baxter equations for two-particle scattering matrices, or to $`U(1)`$ and $`SU(2)`$ symmetries of the scattering processes, respectively. For the simplest case of one shift $`\theta `$, which is connected to the two-chain quantum spin models and most of QFT, the BAE have the form (here we use more general hyperbolic parametrization first; for the rational limit see below) :
$$\underset{\pm }{}e_1^{N_\pm }(u_\alpha \pm \theta )=e^{i\pi M}\underset{\beta =1,\beta \alpha }{\overset{M}{}}e_2(u_\alpha u_\beta ),$$
(1)
where $`N_\pm `$ are the numbers of sites in each of spin chains, $`e_n(x)=\mathrm{sinh}(x+i\gamma \frac{n}{2})/\mathrm{sinh}(xi\gamma \frac{n}{2})`$ and $`M`$ is the number of down spins. The shift $`\theta `$ determines the inter-chain coupling constant for two-chain quantum spin $`\frac{1}{2}`$ models . Please note that the Bethe ansatz equations are just the quantization conditions for the rapidities, which parametrize the eigenwaves and eigenvalues of the many-body quantum model. The Hamiltonian is the first derivative of the logarithm of the transfer matrix (note that the transfer matrix of the two coupled spin chains in this integrable model is the product of two “standard” transfer matrices of each chain with the spectral parameters $`\lambda \pm \theta `$ ):
$`\widehat{H}_{1/2}={\displaystyle \frac{1}{\mathrm{sinh}^2\theta +\mathrm{sin}^2\gamma }}{\displaystyle \underset{n}{}}(\mathrm{sinh}^2\theta \widehat{E}(\stackrel{}{S}_{n,1}\stackrel{}{S}_{n+1,1}+\stackrel{}{S}_{n,2}\stackrel{}{S}_{n+1,2})+`$
$`2\mathrm{sin}^2\gamma \widehat{I}\stackrel{}{S}_{n,1}(\stackrel{}{S}_{n,2}+\stackrel{}{S}_{n+1,2})+2\mathrm{sin}\gamma \mathrm{sinh}\theta (\widehat{J}\stackrel{}{S}_{n+1,2}\widehat{J}\stackrel{}{S}_{n,1})[\stackrel{}{S}_{n+1,1}\times \stackrel{}{S}_{n,2}]),`$ (2)
where $`diag(a,b,c)`$ is $`3\times 3`$ diagonal matrix,
$`\widehat{E}=diag(1,1,\mathrm{cos}\gamma ),`$
$`\widehat{I}=diag(\mathrm{cosh}\theta ,\mathrm{cosh}\theta ,\mathrm{cos}\gamma ),`$
$`\widehat{J}=diag(\mathrm{cos}\gamma ,\mathrm{cos}\gamma ,\mathrm{cosh}\theta ),`$
and $`[.\times .]`$ denotes the vector product. Please note that the sum runs over $`n`$ to $`N_+`$ for the chain with spins $`S_{n,1}`$ and to $`N_{}`$ for the chain with spins $`S_{n,2}`$. The parameter $`\theta `$ determines the intra-chain coupling in our two-chain spin model. For $`\theta =0`$ the Hamiltonian and BAE coincide with the ones for the single “easy-plane” antiferromagnetic spin $`\frac{1}{2}`$ chain of length $`N_++N_{}`$ with the only nearest neighbour interactions in it. The eigenvalue of the Hamiltonian (energy) is parametrized as the function of the rapidities as follows:
$$E=\mathrm{sin}\gamma \underset{\pm }{}\underset{\alpha =1}{\overset{M}{}}N_\pm (e_1(u_\alpha \pm \theta )+e_1^1(u_\alpha \pm \theta ))+E_0,$$
(3)
where $`E_0`$ is the energy of the vacuum (ferromagnetic) state (with $`M=0`$). The isotropic $`SU(2)`$-symmetric antiferromagnetic quantum spin two-chain model can be obtained from the uniaxial ($`U(1)`$-symmetric) one of Eqs. (1)-(3) by the simple change of variables in the limit: $`u_\alpha \gamma u_\alpha `$, $`\lambda \gamma \lambda `$, $`\theta \gamma \theta `$, $`\gamma 0`$. (The last limit corresponds to the rational, $`SU(2)`$-symmetric solution of the Yang-Baxter equations for two-particle scattering matrices). The two-chain isotropic ($`SU(2)`$-symmetric) spin $`\frac{1}{2}`$ Hamiltonian obtained in this limit from Eq. (2) takes the form :
$`\widehat{H}_{is}={\displaystyle \frac{1}{1+\theta ^2}}{\displaystyle \underset{n}{}}(\theta ^2(\stackrel{}{S}_{n,1}\stackrel{}{S}_{n+1,1}+\stackrel{}{S}_{n,2}\stackrel{}{S}_{n+1,2})+2\stackrel{}{S}_{n,1}(\stackrel{}{S}_{n,2}+\stackrel{}{S}_{n+1,2})+`$
$`2\theta (\stackrel{}{S}_{n+1,2}\stackrel{}{S}_{n,1})[\stackrel{}{S}_{n+1,1}\times \stackrel{}{S}_{n,2}]).`$ (4)
The summations over $`n`$ runs to $`N_\pm `$ for each kind of spins, respectively. Note that for $`\theta \mathrm{}`$ Eqs. (4) and BAE recover the Hamiltonian and BAE of two decoupled spin $`\frac{1}{2}`$ chains of lengths $`N_\pm `$ with the only nearest neighbour interactions in each of the chains.
The solution to the BAE Eqs. (1) is usually obtained in the thermodynamic limit ($`N_\pm ,M\mathrm{}`$, with the ratio $`M/(N_++N_{})`$ fixed). Here instead of the discret set of rapidities one introduces the distribution of a continuous density of rapidities. The groundstate corresponds to the solutions of the BAE with negative energies, i.e., it is connected with the filling up the Dirac sea(s) for the model. For the “easy-plane” antiferromagnetic two-chain spin $`\frac{1}{2}`$ model the groundstate corresponds to the filling of the Dirac sea for the real rapidities, i.e., no spin boundstates have negative energies. In the thermodynamic limit the real roots of Eqs. (1) are distributed continuously over some intervals, which determine the Dirac seas of the model. The set of integral equations for the dressed densities of rapidities $`u_\alpha `$ ($`\rho (u)`$) and dressed energies of low-lying quasiparticles ($`\epsilon (u)`$) are (see, e.g., Ref. for the standard procedure of deriving these integral equations from the BAE and Refs. for the isotropic two-chain spin $`\frac{1}{2}`$ model):
$$\rho (u)+_{(Q)}𝑑vK(uv)\rho (v)=\underset{\pm }{}\frac{N_\pm }{N}\rho _\pm ^0$$
(5)
and
$$\epsilon (u)+_{(Q)}𝑑vK(uv)\epsilon (v)=h\underset{\pm }{}\frac{N_\pm }{N}\epsilon _\pm ^0,$$
(6)
where the kernels of integral equations are
$$K(u)=\frac{\mathrm{ln}e_2(u)}{u}=\frac{\mathrm{sin}(2\gamma )}{2\pi [\mathrm{cosh}(u)\mathrm{cos}(2\gamma )]}.$$
(7)
and $`h`$ is an external magnetic field. The values
$$\rho _\pm ^0(u)=\frac{\mathrm{ln}e_1(u\pm \theta )}{u}\frac{p_\pm ^0(u)}{u}=\frac{\mathrm{sin}\gamma }{2\pi [\mathrm{cosh}(u\pm \theta )\mathrm{cos}(\gamma )]}$$
(8)
are bare densities of the rapidities, and
$$\epsilon _\pm ^0(u)=h\frac{\mathrm{sin}^2\gamma }{\mathrm{cosh}(u\pm \theta )\mathrm{cos}(\gamma )}$$
(9)
are bare energies (here “bare” corresponds to noninteracting particles, and the interaction “dresses” them as usual ). The integrations are performed over the domain $`(Q)`$, determined in such a way that the dressed energies inside these intervals are negative. The limits of integrations are determined by the zeros of the dressed energies, and are the Fermi points for each sea. The analysis of the integral equations Eqs. (5),(6) in an external magnetic field shows that in general, for some values of $`\theta `$ and $`h`$, there can be one Dirac sea (it corresponds to one minimum of the bare density of rapidities and, hence to one minimum of the bare energy). On the other hand, for higher values of $`\theta `$ and for some domain of $`h`$ two Dirac seas of the same type of (gapless, see below) excitations are possible (for two minima of the bare energies of the rapidities and thus two minima of the bare density). Note that for $`\theta \mathrm{}`$ at fixed $`N_\pm `$ all the roots of the integral BAE separate into two sets of “right-” and “left-moving” seas, centered at $`\pm \theta `$, respectively.
Here we briefly re-visit the analysis of Refs. , but for the case of the uniaxial two-chain model. Analytic solutions to Eqs. (5)-(6) can be easily obtained in closed form in the limit of zero field and equal lengths of the chains $`N_+=N_{}`$. The simplest non-trivial exited quasiparticle (spinon) is a hole in the Dirac sea for real rapidities with the quasimomentum
$$p(u_0)=2\mathrm{arctan}\left(\frac{\mathrm{sinh}(\pi u_0/\gamma )}{\mathrm{cosh}(\pi \theta /\gamma )}\right),$$
(10)
where $`u_0`$ is the spinon’s rapidity. Note that due to topological reasons such particles have to exist in pairs for the $`SU(2)`$-symmetric case, etc. . The energy of this spinon is given by
$$ϵ(u_0)=\mathrm{sin}\gamma \frac{p(u_0)}{u_0}.$$
(11)
It can be rewritten as function of the quasimomentum, i.e., in the form of the commonly used dispersion law
$$ϵ(p)=\frac{\pi }{\gamma }\mathrm{sin}\gamma \mathrm{tanh}\frac{\pi \theta }{\gamma }\mathrm{sin}\frac{p}{2}\left[\mathrm{cos}^2\frac{p}{2}+\mathrm{sinh}^2\frac{\pi \theta }{\gamma }\right]^{1/2}.$$
(12)
A spinon corresponds in the usual Bethe ansatz classification of BAE solutions to a string of length 1 . Naturally Eqs. (1) have string solutions of higher lengths too. Other spin excitations can be obtained as combinations of spinon quasiparticles and higher-length strings with different rapidities. However, spinons here are picked out because only their dressed energies may be negative, i.e., only spinons may form Dirac seas of the groundstate of the model.
One can see that the dispersion law Eq. (12) of the low-lying excitation of the “easy-plane” two-chain spin $`\frac{1}{2}`$ antiferromagnetic model is factorized into two parts: the gapless part at $`p=0,\pi `$ and the gapful one at $`p=\pi /2`$, cf. . The former corresponds to the oscillations of the magnetization, while the latter is connected with the oscillations of the staggered magnetization . The analysis, similar to the one of the solutions of Eqs. (5),(6) for nonzero magnetic field $`h0`$ (here we point out that according to the very accurate analysis the solution of the integral BAE in the first order approximation reproduces correctly both low- and high-coupling asymptotic behavior) shows that: (i) the dressed energy of a spinon as a function of the dressed quasimomentum has only one extremum, a maximum at $`p=\pi /2`$ for $`\theta <\theta _c`$ and (ii) for $`\theta >\theta _c`$ there are two maxima and one minimum (situated at $`p=\pi /2`$). At the (tri)critical point $`\theta _c`$, the minimum disappears and two maxima joint into one more flat (at $`p=\pi /2`$). In the limit $`\theta \mathrm{}`$ the mimimum is transformed into a cusp. It reveals that the gap of the staggered magnetization vanishes in this limit of two independent spin chains. This simple picture helps us to understand what happens if one switches on an external magnetic field $`h`$. Besides the usual phase transition to the ferromagnetic (spin-polarized) phase at
$$h_s=\underset{\pm }{}\frac{N_\pm }{N}\epsilon _\pm ^0(0)$$
(13)
there is an additional transition between two phases. One of these corresponds to one Dirac sea of spinons (at small $`\theta `$), while the other one is connected with two Dirac seas for the same kind of spinons (at large $`\theta `$). It can also be seen from the r.h.s. of Eqs. (5), (6) for the densities and dressed energies that the bare density and bare energy (corresponding to terms which do not depend on $`\rho (u)`$ and $`\epsilon (u)`$) have either one or two minima, respectively. Hence, they reproduce the same property in the dressed characteristics: The interaction simply “dresses” the (quasi)particles, as usual, but the “dressing” does not affect the picture qualitatively. The new critical field value can be approximated by $`h_c\frac{\pi }{\gamma }\mathrm{sin}\gamma \mathrm{cosh}^1\frac{\pi \theta }{\gamma }`$ in the first order approximation . In this approximation the tricritical point is the root of the equation $`1\mathrm{sinh}\frac{\pi \theta _c}{\gamma }`$. At this point two second order phase transition lines $`h_s`$ and $`h_c`$ join. Hence, the “easy-plane” magnetic anisotropy in the antiferromagnetic two-chain model does not change qualitatively the groundstate behavior in the external magnetic field, cf. . However it changes the critical values of the magnetic field and the intra-chain coupling. The difference between the two (gapless) phases is obvious: the first phase corresponds to the Néel-like antiferromagnetic groundstate for spins in both chains (along the zigzag line), while the second phase is connected with the Néel-like antiferromagnetic groundstates in each of chains, i.e. to effectively two magnetic sublattices in the two-chain model.
That is why our simple model explains in which domains of parameters the two-chain spin system behaves like one-sublattice quantum “easy-plane” antiferromagnet, and where it behaves as the two-sublattice one. Note also that the phase transitions we study here are the manifestations of the commensurate-incommensurate phase transitions for spin systems. One can obviously see this, because the intra-chain coupling for two spin chains can be interpreted as the next-nearest neighbor spin interactions for a single spin chain of higher length $`N_++N_{}`$. Here the magnetic couplings are spin-frustrated, thus the emergence of the incommensurate magnetic states is understandable.
As a consequence of the conformal invariance of (1+1)-dimensional quantum systems, the classification of universality classes is simple in terms of the central charge (conformal anomaly $`C`$) of the underlying Virasoro algebra . The critical exponents in a conformally invariant theory are scaling dimensions of the operators within the quantum model. They can be calculated considering the finite-size (mesoscopic) corrections for the energies and quasimomenta of the groundstate and low-lying excited states. Conformal invariance formally requires all gapless excitations to have the same velocity (Lorentz invariance). The complete critical theory for systems with several gapless excitations with different Fermi velocities is usually given as a semidirect product of these independent Virasoro algebras. Here we briefly sketch the procedure and write the results for the finite-size corrections to the energy, following the standard procedure, see, e.g. Refs. . One can see that for $`\theta <\theta _c`$ and for $`\theta >\theta _c`$, $`h<h_c`$, the conformal limit of our uniaxial two chain spin $`\frac{1}{2}`$ model corresponds to one level-1 Kac-Moody algebra (one WZW model of level 1 with the conformal anomaly $`C=1`$). The finite-size correction to the energy is rather standard (cf. )
$$E_{fs}(N_++N_{})=\frac{\pi }{6}v_F+2\pi v_F(\mathrm{\Delta }_l+\mathrm{\Delta }_r),$$
(14)
where $`v_F`$ is the Fermi velocity of the spinon and the conformal dimensions of primary operators are (please, pay attention: the lower indices denote the conformal dimensions for right- and left-moving quasiparticles, at the right and left Fermi point, respectively):
$$2\mathrm{\Delta }_{l,r}=\left(\frac{\mathrm{\Delta }M}{2z}\pm z\mathrm{\Delta }D\right)^2+2n_{l,r},$$
(15)
where $`\mathrm{\Delta }M`$ is an integer denoting the change of the number of particles induced by the primary operator, $`\mathrm{\Delta }D`$ is an integer (half-integer) denoting the number of transfered particles from the right to the left Fermi point (backward scattering processes), $`n_{l,r}`$ are the numbers of the particle-hole excitations of right- and left-movers. The values for the quantum numbers are restricted by $`\mathrm{\Delta }D=\mathrm{\Delta }M/2`$ (mod 1). The dressed charge $`z=\xi (Q)`$ is the solution of the (standard) integral equation
$$\xi (u)+_{(Q)}𝑑vK(uv)\xi (v)=1,$$
(16)
taken at the limits of integration (these are the Fermi points, symmetric with respect to zero). In this phase there is only one region of integration over $`v`$. The dressed charge is a scalar. The behavior of our class of models in this phase in the conformal limit is rather standard . The correlation functions decay asymptotically $`(xv_Ft)^{\mathrm{\Delta }_l}(x+v_Ft)^{\mathrm{\Delta }_r}`$. The choise of the appropriate quantum numbers of excitations $`\mathrm{\Delta }M`$, $`\mathrm{\Delta }D`$ and $`n_{l,r}`$ is determined for the leading asymptotics of correlators by taking the possible numbers with smallest exponents.
But for $`\theta >\theta _c`$, $`h>h_c`$, the conformal limit of the “easy-plane” two-chain spin $`\frac{1}{2}`$ model corresponds to the semidirect product of two level-1 Kac-Moody algebras, both with conformal anomalies $`C=1`$, i.e., to two WZW models both of level 1 . The Dirac seas (i.e. the possible spinons with negative energies) are in the intervals $`[Q^+,Q^{}]`$ and $`[Q^{},Q^+]`$ (minima in the distributions of rapidities at $`\theta `$). This can be interpreted as the symmetrically distributed (around zero) Dirac seas of “particles” for $`[Q^+,Q^+]`$ and the Dirac sea of “holes” for $`[Q^{},Q^{}]`$. In fact the valley in the density distribution for “particles” and the maximum for “holes” are in the one-to-one correspondence with the maxima and minimum of the dispersion law for spinons. The second critical field $`h_c`$ in this language corresponds to the van Hove singularity of the empty band of “holes”. Naturally, the Fermy velocities of “particles” are positive, $`v_F^+=(2\pi \rho (Q^+))^1\epsilon ^{}(u)|_{u=Q^+}`$, while the Fermy velocities of “holes” are negative $`v_F^{}=(2\pi \rho (Q^{}))^1\epsilon ^{}(u)|_{u=Q^{}}`$. The finite-size corrections to the energy for this case are
$$E_{fs}(N_++N_{})=\frac{\pi }{6}(v_F^++v_F^{})+2\pi \left(v_F^+(\mathrm{\Delta }_l^++\mathrm{\Delta }_r^+)+v_F^{}(\mathrm{\Delta }_l^{}+\mathrm{\Delta }_r^{})\right),$$
(17)
where the dispersion laws of “particles” and “holes” are linearized about the Fermi points for each Dirac sea. The conformal dimensions of the primary operators are (the upper indices denote Dirac seas; the lower indices denote right and left Fermi points of each of these two Dirac seas, cf. for the isotropic spin $`\frac{1}{2}`$ two-chain model):
$$2\mathrm{\Delta }_{l,r}^{}=\left[\frac{(x_\pm \mathrm{\Delta }M^+x_{+\pm }\mathrm{\Delta }M^{})}{2det\widehat{x}}\frac{(z_\pm \mathrm{\Delta }D^+z_{+\pm }\mathrm{\Delta }D^{})}{2det\widehat{z}}\right]^2+2n_{l,r}^{},$$
(18)
where the “minus” sign between the terms in square brackets corresponds to the right-, and “plus” sign to the left-movers. Here $`\mathrm{\Delta }M^\pm `$ denote the differences between the numbers of particles excited in the Dirac seas of “particles” and “holes”, labeled by the upper indices. $`\mathrm{\Delta }D^\pm `$ denote the numbers of backward scattering excitations, and $`n_{l,r}^\pm `$ are the numbers of the particle-hole excitations for right- and left-movers of each of Dirac seas (for “particles” and “holes”). Please pay attention that $`\mathrm{\Delta }M^\pm `$ and $`\mathrm{\Delta }D^\pm `$ are not independent. Their values are restricted by the following connections: $`\mathrm{\Delta }M^+\mathrm{\Delta }M^{}=\mathrm{\Delta }M`$, and $`\mathrm{\Delta }D^+\mathrm{\Delta }D^{}=\mathrm{\Delta }D`$, where $`\mathrm{\Delta }M`$ and $`\mathrm{\Delta }D`$ determine in a standard way the changes of the total magnetization and the total momentum of the system, respectively, due to excitations. Please note that in Refs. these restrictions were missing; this resulted in, e.g., the invalid statement that four independent backscattering low-lying excitations are possible. However one can see that only two of them are really independent. The same is true for the excitations that change the total magnetization of the system: there are only two independent of four possible such excitations. This is a direct consequence of the fact that only one magnetic field determines the filling of the Dirac seas for “particles” and “holes”, or, in other words, two Dirac seas for spinons at $`\pm \theta `$.
The dressed charges $`x_{ik}(Q^k)`$ and $`z_{ik}(Q^k)`$ ($`i,k=+,`$) are matrices in this phase. They can be expressed by using the solution of the integral equation
$$f(u|Q^\pm )=\left(_{Q^+}^{Q^+}_Q^{}^Q^{}\right)K(uv)f(v|Q^\pm )=K(uQ^\pm ),$$
(19)
with
$`z_{ik}(Q^k)=\delta _{i,k}+()^k{\displaystyle \frac{1}{2}}\left({\displaystyle _{Q^i}^{\mathrm{}}}{\displaystyle _{\mathrm{}}^{Q^i}}\right)dvf(v|Q^k)`$
$`x_{ik}(Q^k)=\delta _{i,k}()^k{\displaystyle _{Q^i}^{Q^i}}𝑑vf(v|Q^k).`$ (20)
Notice, please, that the dressed charges depend on the value of the magnetic anisotropy $`\gamma `$ via the kernels, while they depend indirectly on the value of the intra-chain coupling constant $`\theta `$, only via the limits of integrations. In the first order approximation one can write the solutions as $`x_{ik}(Q^k)\delta _{i,k}()^k_{Q^i}^{Q^i}𝑑vK(vQ^k)+\mathrm{}`$ and $`z_{ik}(Q^k)\delta _{i,k}+()^k(1/2)(_{Q^i}^{\mathrm{}}_{\mathrm{}}^{Q^i})dvK(uQ^k)+\mathrm{}`$. The Dirac sea for “holes” disappears, naturally for $`hh_c`$, $`\theta \theta _c`$. The slopes of the dressed energies of “particles” and “holes” at Fermi points of the Dirac seas (Fermi velocities) differ in general from each other. Therefore we have a semidirect product of two algebras. Hence, in this region the dressed charges are $`2\times 2`$ matrices. This means that the conformal limit of the “easy-plane” two-chain spin $`\frac{1}{2}`$ model corresponds to one or two WZW theories, depending on the values of the intra-chain coupling, magnetic anisotropy and magnetic field. At the critical line $`h_c`$ the Dirac sea of “holes” disappears as well as the components of the dressed charge matrix $`\widehat{x}`$ (with square root singularities of the critical exponents for the correlation functions). Note that the dressed charge $`z`$ becomes $`z=(2x)^1`$ at the phase transition line $`h_c`$. This corresponds to the disappearence of one of the WZW CFTs. Unfortunately it is impossible to obtain an analytic solution to Eqs. (19) in closed form for a finite inter-chain coupling $`\theta `$. Naturally in the limiting cases of two independent chains of lengths $`N_\pm `$, $`\theta \mathrm{}`$, and a single chain of length $`N_++N_{}`$, $`\theta =0`$, the solutions of Eqs. (16),(19),(20) coincide with well-known ones, see Refs. . The correlation functions of the uniaxial two-chain spin $`\frac{1}{2}`$ model decay algebraically in this phase $`(xv_F^+t)^{\mathrm{\Delta }_l^+}(x+v_F^{}t)^{\mathrm{\Delta }_l^{}}(xv_F^+t)^{\mathrm{\Delta }_r^+}(x+v_F^{}t)^{\mathrm{\Delta }_r^{}}`$ with the minimal exponents of possible quantum numbers of excitations $`\mathrm{\Delta }M^\pm `$, $`\mathrm{\Delta }D^\pm `$ and $`n_{l,r}^\pm `$. We point out once more that the same magnetic field plays the role of a chemical potential for the “particles” and “holes”, or spinons of both Dirac seas in the second phase, and hence this choise of “minimal quantum numbers” is constrained.
We must point out here that there is a crucial difference between our situation and the case of dressed charge matrices appearing for systems with the internal structure of bare particles . There the two Dirac seas of the groundstates are connected with different kinds of excitations, e.g., holons and spinons for the repulsive Hubbard model, or Cooper-like singlet pairs and spinons for the supersymmetric $`tJ`$ model. They correspond to two different kinds of Lagrange multipliers, the chemical potentials and magnetic fields. Thus the low-lying excitations of the conformal theories in the spin and charge sectors of these correlated electron models are practically independent of each other (spin-charge separation). Note that the spin and charge sectors are connected via the off-diagonal elements of the dressed charge matrix though. This is the consequence of the fact that say, holons or unbound electrons carry both charge and spin. On the other hand, two Dirac seas appear for the same kinds of particles for the models studied in this paper, which are also connected with the same magnetic field governing the filling of both Dirac seas. The latters appear due to two minima in the bare energy distribution and correspond to nonzero shift $`\theta `$ in the Bethe ansatz equations. In other words, two Dirac seas are determined by the inter-chain coupling and appear if the values of coupling and external magnetic field are higher than the threshold values $`\theta _c`$ and $`h_c`$, respectively. We believe that such a threshold behavior does not depend on the integrability of the model and is the generic feature for any multi-chain quantum spin models.
The low temperature Sommerfeld approximation shows that as usually the low temperature specific heat is proportional to $`T`$ out of critical lines. At the critical lines the van Hove singularities produce $`\sqrt{T}`$ low temperature behavior of the specific heat, while at the tricritical point we have $`T^{1/4}`$ behavior.
What are the changes due to the different lengths of the chains $`N_+N_{}`$? One can see obviously that the values of the momentum, energy and velocity of a spinon (which was $`v=(\pi /\gamma )\mathrm{sin}\gamma \mathrm{tanh}(\pi \theta /\gamma )`$) become functions of $`N_+N_{}`$. For example, the velocity renormalizes as $`vv[1+(N_+N_{})^2\mathrm{tanh}^2(\pi \theta /2\gamma )/N^2]^1`$. This introduces dependences of the critical values $`\theta _c`$ and $`h_c`$ (as well as the saturation field $`h_s`$) on the difference $`N_+N_{}`$. Also, the Fermi velocities and Fermi points for the finite-size corrections become dependent on this difference. One can in principle consider different coupling constants $`J_\pm `$ for each of the chains (overall multipliers ). This produces the renormalizations similar to the action of $`N_+N_{}`$, i.e., the velocity, e.g., renormalizes as $`vJ_+v[1+(J_{}/J_+)^2\mathrm{tanh}^2(\pi \theta /2\gamma )]^1`$.
## 3 Connections to the Quantum Field Theories
The studies presented in the previos section, being rather standard (note, though, some important new features, which were absent in the previous studies , such as the dependencies of the critical values of the inter-chain coupling and external magnetic field on the parameter of the magnetic anisotropy and on the difference in the lengths of the chains; also the important restrictions on the quantum numbers of low-lying conformal excitations). However we will use the results of that section for novel studies for a wider classes of exactly solvable models in Sections 3-5. For instance, in this section we point out the important similarities in the behaviors of the two-chain quantum spin model considered in the previous section and several models of QFT.
Really, when examinating Eqs. (1), one can see that these Bethe ansatz equations coincide with the ones, which describe the behavior of the spin (color) sector of some QFT. $`N_\pm `$ corresponds to the numbers of (bare) particles with the positive and negative chiralities. For example, for the chiral-invariant Gross-Neveu model we have to put $`\gamma 0`$ (i.e. $`SU(2)`$-symmetric case, equivalent to the $`SU(2)`$-symmetric Thirring model), and $`\theta =(1g^2)/2g`$, where $`g`$ is the coupling constant of the chiral invariant Gross-Neveu QFT . As for the Lagrange multiplier $`h`$, it can play the roles of either an external magnetic field, or the chemical potential, or an external topological field, dual to the topological Noether current in QFT. Here we point out that in fact in QFT theorists are interested in physical particles, which have the finite mass (gap). In the chiral-invariant Gross-Neveu model the gap of the staggered oscillations of the two-chain quantum spin model plays the role of the physical mass of the particle (spinor) . As for the (gapless) oscillations of the magnetization of the two-chain spin model, we point out that they are the consequences of the lattice, and play the role of the massless fermion doublers of the lattice QFT . The results of the previous section mean that the behavior of the chiral-invariant Gross-Neveu model (or $`SU(2)`$-symmetric Thirring model) in an external magnetic field depends strongly on the coupling constant $`\theta `$ (or equivalently on $`g`$). For $`\theta <\theta _c`$ the conformal limit of the QFT corresponds to one level-1 WZW model with the conformal dimension $`C=1`$. However for $`\theta >\theta _c`$ ($`\theta _c\sqrt{\theta _c^2+4}<2g<\theta _c+\sqrt{\theta _c^2+4}`$) the conformal limit of this QFT in an external magnetic field corresponds to the semidirect product of two level-1 WZW model with the conformal dimensions $`C=1`$. Two kinds of conformal points for this QFT have been mentioned already in a slightly different context. They were connected with one WZW theory or two WZW theories, coupled via a current-current interaction. This is related to right-left symmetry of the chiral invariant Gross-Neveu QFT (see, also, Refs. for the case of the QFT for the principal chiral field).
Note, that the condition $`h>h_c`$ in the QFT means that the magnetic field is larger than the mass of the physical particle (color spinor). In this sense, in the region of magnetic field values $`h<h_c`$ the results of the QFT (see, e.g., ) predict zero magnetization, however a different lattice regularization, similar to the lattice scheme used in the previous section predicts a nonzero magnetization of the chiral-invariant Gross-Neveu model in this region. This is the indirect effect of the fermion doublers. In other words, it is connected with the well-known mapping of the lattice (e.g., Thirring) model under regularization on two continuum QFTs either both bosonic (free bosonic QFT and sine-Gordon one, ), or both fermionic ones (a free one and the continuum massive Thirring model). There are necessarily two such theories because of the Nielsen-Ninomiya fermion doublers: remember that we have started from a lattice .
For other models of QFT the procedure of the lattice regularization has been used. Here $`\theta `$ plays the role of the cut-off to preserve the mass of the physical particle to be finite. For example, for the $`U(1)`$-symmetric Thirring QFT one can use the results of the previuos section with the limit $`\theta \mathrm{}`$ taken after the thermodynamic limit ($`L,N_\pm ,M\mathrm{}`$ with their ratios fixed, $`L`$ is the size of the box). In this case one can obviously obtain the conformal limit of the theory with nonzero physical masses of the particles. Naturally in the limit $`\theta \mathrm{}`$ we ever exist in an external magnetic field in the phase with two Dirac seas. Here the latters correspond to the right- and left-moving particles (with positive and negative chiralities). Actually here our point of view coincides with the one of the field theorists. Recently it was shown in Ref. that for (1+1)-dimensional sine-Gordon model the lattice regularization scheme in the “light-cone” approach gives similar to ours results for the conformal limit of the model. It was shown there that at the UV fixed point the conformal dimensions of the sine-Gordon model are determined by two $`U(1)`$ charges of excitations (the usual one and the chiral charge). The chiral charge corresponds to the number of excitations transfered from one Dirac sea to the other, similar to our results (note that the above-mentioned lattice-regularized sine-Gordon case corresponds in our notations to $`\theta \mathrm{}`$, where the integral equations for the particles with the positive and negative chiralities are totally decoupled). We point out here, that such a behavior is not unexpected, because the sine-Gordon QFT belongs to the same class of models, which is studied in our paper, i.e., its Bethe ansatz description features a shift of rapidities in the Bethe ansatz equations in the lattice-regularized theory .
## 4 Higher spin (chirality) generalizations
For the higher spin generalizations of the Bethe ansatz theory presented in Section 2 we can write down BAE in the form
$$\underset{\pm }{}e_{n_\pm }^{N_\pm }(u_\alpha \pm \theta )=e^{i\pi M}\underset{\beta =1,\beta \alpha }{\overset{M}{}}e_2(u_\alpha u_\beta ),$$
(21)
where $`n_\pm =2S_\pm `$ are the values of spins in each chain or the colors of bare particles in QFT. The eigenvalue of the transfer matrix can be writen as
$`\mathrm{\Lambda }(\lambda )={\displaystyle \underset{\alpha =1}{\overset{M}{}}}{\displaystyle \frac{\mathrm{sinh}(\lambda u_\alpha +i\gamma \frac{1}{2})}{\mathrm{sinh}(u_\alpha \lambda +i\gamma \frac{1}{2})}}+e^{i\pi M}{\displaystyle \underset{\pm }{}}\left({\displaystyle \frac{\mathrm{sinh}(\lambda \pm \theta )}{\mathrm{sinh}(i\gamma \frac{n_\pm }{2}\lambda \theta )}}\right)^{N_\pm }\times `$
$`{\displaystyle \underset{\alpha =1}{\overset{M}{}}}{\displaystyle \frac{\mathrm{sinh}(u_\alpha \lambda +i\gamma \frac{3}{2})}{\mathrm{sinh}(\lambda u_\alpha i\gamma \frac{1}{2})}}.`$ (22)
Similar new phases with one or two kinds of Dirac seas for similar kinds of low-lying excitations exist also for a number of models in which $`n_\pm 1`$, e.g. for the higher-spin ($`S>\frac{1}{2}`$) two-chain models with equal spins in each chain, $`SU(n+1)`$ CIGN QFT : there $`n_+=n_{}=n1`$; for the principal chiral field models (nonlinear $`\sigma `$-model) for $`CP`$-symmetric (there $`n_+=n_{}\mathrm{}`$) and $`CP`$-asymmetric cases (there $`n_+n_{}`$, $`(n_++n_{})\mathrm{}`$, $`(n_+n_{})`$ fixed, i.e., the symmetry $`SU(2)\times SU(2)O(4)`$); and for the $`O(3)`$-symmetric nonlinear $`\sigma `$-model as well as for spin-$`(S_+2n_+)`$ \- spin-$`(S_{}2n_{})`$ two-chain models (quantum two-chain ferrimagnet). Note that for spins $`S\frac{1}{2}`$ the procedure of the construction of the Hamiltonian is more complicated, because it corresponds to the two-chain uniaxial generalization of the Takhtajan-Babujian model, see e.g., Refs. . For the simplest case of the isotropic exchange interaction between the spins and between the chains the Hamiltonian has the form:
$`={\displaystyle \underset{n}{}}(\theta ^2(_{S_+,S_+,n_1,n_1+1}+_{S_{},S_{},n_2,n_2+1})+2(_{S_+,S_{},n_1,n_2}+_{S_+,S_{},n_1,n_2+1})+`$
$`\{(_{S_+,S_+,n_1,n_1+1}+_{S_{},S_{},n_2,n_2+1}),(_{S_+,S_{},n_1,n_2}+_{S_+,S_{},n_1,n_2+1})\}+`$
$`2i\theta [(_{S_+,S_+,n_1,n_1+1}+_{S_{},S_{},n_2,n_2+1}),(_{S_+,S_{},n_1,n_2}+_{S_+,S_{},n_1,n_2+1})]),`$ (23)
where $`[.,.]`$ ($`\{.,.\}`$) denote (anti)commutator,
$$_{S_1,S_2,n,n+1}=\underset{j=|S_1S_2|+1}{\overset{S_1+S_2}{}}\underset{k=|S_1S_2|+1}{\overset{j}{}}\frac{k}{k^2+\theta ^2}\times \underset{l=|S_1S_2|}{\overset{S_1+S_2}{}}\frac{xx_l}{x_jx_l},$$
(24)
$`x=\stackrel{}{S}_{1,n}\stackrel{}{S}_{2,n+1}`$, and $`2x_j=j(j+1)S_1(S_1+1)S_2(S_2+1)`$. The summation over $`n`$ runs to $`N_\pm `$ in each chain. One can obviously see that for $`S_\pm =\frac{1}{2}`$ the Hamiltonian Eq. (23) recovers the isotropic antiferromagnetic spin $`\frac{1}{2}`$ Hamiltonian Eq. (4) investigated in Section 2. For a single spin chain, $`\theta =0`$, $`N_+=N_{}`$ the Hamiltonian coincides with the known one of alternating spin chains . The Bethe ansatz studies of the model for $`n_\pm `$ can be performed in the complete analogy with the above mentioned case $`n_\pm =1`$, keeping in mind, of course, the main difference: for the $`SU(2)`$-symmetric or uniaxial higher spin models the groundstate corresponds to the filling up the Dirac seas for spin strings of lengths $`n_\pm `$. The well-known fusion scheme can be used for the case of a flavor-degenerate situation of the chiral invariant Gross-Neveu QFT, in the absence of flavor fields . Note that, except of the $`O(3)`$-symmetric case, $`\gamma =0`$ everywhere in the above-mentioned models of QFT. This corresponds to rational solutions of the Yang-Baxter equation for the two-particle scattering matrices. For the two spin chains the two-chain quantum ferrimagnet model corresponds to two Takhtajian-Babujian chains of different values of site spins coupled due to nonzero $`\theta `$. The total quasimomentum and the energy of the system in the framework of the lattice (local) regularization scheme for some QFT can be written as
$`2a_tE={\displaystyle \underset{\pm }{}}{\displaystyle \underset{\alpha =1}{\overset{M}{}}}{\displaystyle \frac{}{u_\alpha }}N_\pm \mathrm{ln}e_{n_\pm }(u_\alpha \pm \theta )`$
$`iaP={\displaystyle \underset{\pm }{}}{\displaystyle \underset{\alpha =1}{\overset{M}{}}}N_\pm \mathrm{ln}e_{n_\pm }(u_\alpha \pm \theta ),`$ (25)
where $`a`$ and $`a_t`$ denote the space and time lattice constants, respectively, and their ratio fixes the velocity of light (“light-cone” approach). The $`CP`$-symmetric (chiral invariant) case corresponds to the situation, in which $`n_+=n_{}=n`$. The Dirac seas are related to the dressed (quasi)particles with negative energies (strings of length $`n_\pm `$). The behavior of the dispersion law for excited particle in the $`CP`$-symmetric case ($`n_+=n_{}=n`$ and $`N_+=N_{}`$) is similar to Eq. (12): for instance, for the chiral-invariant Gross-Neveu QFT and principal chiral field model the r.h.s. of Eq. (12) must be simply multiplied by $`\mathrm{sin}(\pi r/n+1)/\mathrm{sin}(\pi /n+1)`$ and the parameter $`\theta `$ in Eq. (12) has to be replaced by $`(n+1)\theta /2`$. $`r=1,\mathrm{},n`$ is the rank of a fundamental representation of the $`SU(n+1)`$ algebra. All the previously mentioned characteristic features from the case $`n_\pm =1`$ persist. The differences are in the levels of Kac-Moody algebras in the conformal limit: The conformal anomalies are $`C=\frac{3n}{n+2}`$. Now the conformal field theory is a semidirect product of a Gaussian ($`C=1`$) and $`Z(n)`$ parafermion models : the operators identified from the scaling behavior of states consisting of Dirac sea strings only (found from finite-size corrections) are found to be composite operators formed by the product of a Gaussian-type operator and the operator in the parafermionic sector. To find a nonzero contributions from parafermions (constant shifts) one can consider the states with strings of other lengths then the Dirac sea present . For the scaling dimensions these shifts are $`\frac{2rr^2}{2n+4}`$, $`r=1,2,\mathrm{}`$.
From now on we concentrate on the $`n_+n_{}`$ situation. For the two-chain spin system the situation corresponds to the quantum ferrimagnet. Here we point out that due to the zigzag-like interactions in the system and spin frustration the ferrimagnets of this class are in the singlet groundstate (compensated phase) for $`h=0`$, unlike the standard classical ferrimagnets in uncompensated phases. The integral equations that determine the physical vacuum of the systems are similar to Eqs. (5)-(6). They reveal one or several minima of the corresponding distributions of dressed energies and densities with possible negative energy states, i.e., one or several Dirac seas:
$`\epsilon _\tau (u)+{\displaystyle 𝑑vK_{\tau \tau ^{}}(uv)\epsilon _\tau ^{}(v)}=h\frac{N_\tau }{N}n_\tau +{\displaystyle \underset{\pm }{}}\frac{N_\pm }{N}\epsilon _{\tau ,\pm }^0`$
$`\rho _\tau (u)+{\displaystyle 𝑑vK_{\tau \tau ^{}}(uv)\rho _\tau ^{}(v)}={\displaystyle \underset{\pm }{}}\frac{N_\pm }{N}\rho _{\tau ,\pm }^0.`$ (26)
The index $`\tau `$ enumerates two possible Dirac seas and appears due to $`n_+n_{}`$, and the $`\pm `$ enumerate two possible minima due to the nonzero shift $`\theta `$. The index $`\tau `$ was naturally absent for the $`CP`$-symmetric case $`n_+=n_{}`$. Note that for quantum two-chain ferrimagnets the investigated gapless phases in the groundstate in an external magnetic field are similar to the spin-compensated and uncompensated phases. Thus the phase transition between those phases is similar in nature to the well-known spin-flop phase transition in the classical theory of magnetism. Note, though, that the spin-flop transition is of the first order (“easy-axis” magnetic anisotropy), while the transition under study is the second-order one (“easy-plane” anisotropy). The Fourier transform of the kernel is given by
$`2\mathrm{coth}(\omega /2)[diag(e^{n_+|\omega /2|}\mathrm{cosh}(n_+\omega /2),e^{n_{}|\omega /2|}\mathrm{cosh}(n_{}\omega /2))+`$
$`\widehat{\sigma }_x(e^{(n_+n_{})|\omega /2|}e^{(n_++n_{})|\omega /2|})].`$ (27)
$`diag(a,b)`$ is $`2\times 2`$ diagonal matrix and $`\widehat{\sigma }_x`$ is the usual Pauli matrix. Note, that after taking the limit $`(n_++n_{})\mathrm{}`$, which is the case of the $`CP`$-asymmetric case of the QFT for the principal chiral field, i.e. with the Wess-Zumino term , the inverse kernel coincides formally (up to a constant multiplier) with the one for the case $`n_+=n_{}=1`$. This indicates the global $`O(4)`$ ($`O(3)`$) symmetry of the principal chiral field . There may be also two different behaviors, corresponding to one or several Dirac seas for $`n_+1`$ or $`n_{}1`$. Naturally in the conformal limit the associated WZW CFTs have different conformal anomalies determined by $`n_\pm `$: $`C_\pm =\frac{3n_\pm }{n_\pm +2}`$. For the determination of the Gaussian parts of the conformal dimensions of primary operators Eqs. (18) can be used. One has to add the input from the parafermionic sectors, too . The elements of the dressed charge matrices are the solutions of the following system of integral equations:
$$\xi _{\tau ,\tau ^{}}(u)+\underset{\pm }{}𝑑vK_\tau ^{}(uv)\xi _{\tau ,\pm }(v)=\delta _{\tau ,\tau ^{}},$$
(28)
in which the summation over $`\pm `$ is due to the two possible Dirac seas (two minima in the distribution of rapidities) at $`\pm \theta `$. For different values of the spins, $`n_+n_{}`$, a transition between two different phases is induced by increasing an external magnetic field to some critical value, even in the absence of the shift $`\theta `$ . It differs from the $`CP`$-symmetric case $`n_+=n_{}`$, where the phase transition is only connected with the nonzero value of the intra-chain coupling parameter $`\theta `$. For the $`CP`$-symmetric case one or two Dirac seas of the same type of excitations exist due to nonzero $`\theta `$. But in the $`CP`$-asymmetric case the existence of two Dirac seas can be related to two kinds of different low-lying excitations (particles). They are strings of lengths $`n_+`$ and $`n_{}`$, respectively. In this situation the dispersion laws may be independent (not only factorized as for the previous $`CP`$-symmetric cases). The (new) phase transition at $`h_c`$ reveals the van Hove singularity of the empty Dirac sea for the longer strings. The spin saturation field $`h_s`$ is connected with the empty Dirac sea of strings of the smaller length.
## 5 Multi-chain quantum spin models
It is worthwile to mention that phase transitions in an external magnetic field, similar to the ones studied in this paper for uniaxial spin chains and QFT, have been already studied in the 1D quantum alternating single spin chains , spin $`\frac{1}{2}`$ isotropic two-chain models , and correlated electron models with the finite concentration of magnetic impurities . The Bethe ansatz equations for those models are similar to the ones studied in the present paper, Eqs. (1),(21). Note that the energies for spin models are defined (as usual for the lattice models) as first logarithmic derivatives of the transfer matrices. The factorization of the dispersion law for the lowest excitations (spinon) reveals essentially two kinds of magnetic oscillations: excitations of the magnetization and oscillations of the staggered magnetization, i.e., the manifestation of essentially two magnetic sublattices. Naturally, the existence of the latters persists in the continuum limit of such systems too (cf., for instance, with the standard theory of antiferromagnetism). Two non-ferromagnetic phases also reveal themselves in finite-size corrections to the energies of these quantum spin models. There instead of a scalar dressed charge for the phase with one Dirac sea for spinons, $`2\times 2`$ dressed charge matrices appear in the second phase with two Dirac seas for for spin strings of different lengths in alternating spin chain , or for spinons of the same kind in zigzag-like coupled spin chains (see for the isotropic two-chain spin-$`\frac{1}{2}`$ model).
The symmetry-breaking terms \[the difference $`(n_+n_{})=2(S_1S_2)`$, or nonzero $`\theta `$\] in BAE are actually the reason for the emergence of several gapless phases (or two Dirac seas) in the groundstate in an external magnetic field. It is also interesting to note that a homogenuous shift of rapidities can be removed for one Dirac sea for the periodic boundary conditions by an appropriate unitary (gauge) transformation (shift of variables), e.g., $`u_\alpha u_\alpha \pm \theta `$. But in the case of open boundary conditions, BAE take the form (for simlicity reasons we write the free boundary situation only, without the external boundary potential):
$$\underset{\pm }{}e_{n_\pm }^{2N}(u_\alpha \pm \theta )=\underset{\pm }{}\underset{\beta }{}e_2(u_\alpha \pm u_\beta ).$$
(29)
It is clear that for the open chain one cannot remove the shift $`\theta `$ of rapidities $`u_\alpha `$ from one Dirac sea by a special choice of the gauge. From this point of view the latter case is close to the $`CP`$-asymmetric situation in QFT.
One can see from the structure of the Hamiltonians that for the two-chain spin models the parameter $`\theta `$ characterizes the intra-chain coupling for each chain (or the next-nearest-neighbor interaction in a single spin chain picture). It is obvious to introduce the series of $`\{\theta _j\}_{j=1}^J`$ (for each chain) and to construct the Hamiltonian of the exactly integrable multi-chain ($`J`$ is the number of chains) spin model. For the simplest case of all $`S=\frac{1}{2}`$ isotropic antiferromagnetic chains the Hamiltonian reads :
$`\widehat{H}_J=A{\displaystyle \underset{n}{}}(\left({\displaystyle \underset{i,k}{}}(\theta _i\theta _k)\right)\widehat{P}_{S_{n,r}S_{n+1,r}}+{\displaystyle \underset{p<q}{}}{\displaystyle \frac{\underset{i,k}{}(\theta _i\theta _k)}{(\theta _p\theta _q)}}[\widehat{P}_{S_{n,q}S_{n+1,p}},\widehat{P}_{S_{n,q}S_{n+1,q}}+`$
$`\widehat{P}_{S_{n,p}S_{n+1,p}}]+\mathrm{}+({\displaystyle \underset{j=1}{\overset{J}{}}}\widehat{P}_{S_{n,j}S_{n,j+1}}\widehat{P}_{S_{n,J}S_{n,J+1}}+\widehat{P}_{S_{n,J}S_{n+1,1}})),`$ (30)
where $`A`$ is the normalization constant (which depends on $`\theta _j`$), $`\widehat{P}_{S_aS_b}(1/2)\widehat{I}\widehat{I}+2\stackrel{}{S}_a\stackrel{}{S}_b`$ is the permutation operator and $`[.,.]`$ denotes a commutator. Note, that in the case of $`J2`$ the integrable model corresponds to the pair couplings not only between the nearest-neighbor spins but also to the next-nearest three spins, etc., couplings. All those terms are only essential in quantum mechanics, because in classical physics they are total time derivatives and do not change equations of motion. The Bethe ansatz equations have the form:
$$\underset{j=1}{\overset{J}{}}e_1^{N_j}(u_m+\theta _j\theta _1)=e^{i\pi M}\underset{k}{\overset{M}{}}e_2(u_mu_k),$$
(31)
where $`M`$ is the total number of down spins and $`N_j`$ is the number of sites in the $`j`$-th chain. The previously studied situation $`J=2`$ corresponds to the shift of the variables $`u_mu_m+\theta `$ with $`\theta _2\theta _1=2\theta `$. Now $`\theta _j\theta _1`$ determines the values of the intra-chain couplings in chain $`j`$.
The analysis of the low-temperature thermodynamics of the multi-chain spin system is analogous to the situation of $`J=2`$ studied in Sections 2-4. From the structure of the Bethe ansatz equations in the thermodynamic limit $`N_j,M\mathrm{}`$, their ratios fixed, one can see that for $`J`$-chain model (for different $`\theta _j`$) there can exist, generally speaking, $`J`$ phase transitions of the second order in the groundstate in an external magnetic field. These are nothing else than the commensurate-incommensurate phase transitions for the quantum multi-chain spin model with different couplings between the chains. The values of the critical fields $`h_{c_1},\mathrm{},h_{c_{J1}}`$ and the value of the magnetic field of the transition to the ferromagnetic state $`h_s`$ depend on the set of $`\theta _j`$, i.e., on the intra-chain couplings (and also on the values of the magnetic anisotropy constants, which can be taken different for each chain; this does not destroy the integrability). The ferromagnetic state is gapped, while all other phases are gapless in the integrable multi-chain spin quantum model. There are also $`J1`$ tricritical points at which the lines of the phase transitions $`h_{c_j}`$ join the line of the spin-saturation phase transition. Naturally, the phase that corresponds to the lowest value of the magnetic field, say $`h<h_{c_1}`$ for special values of $`\theta _j`$ (the condition is similar to $`\theta <\theta _c`$ for $`J=2`$), has in the conformal limit one scalar dressed charge. Hence, in the conformal limit our multi-chain spin model behaves as the level-1 WZW CFT. In the next phase the multi-chain quantum spin model behaves as the semidirect product of two WZW CFTs, hence their dressed charges are $`2\times 2`$ matrices, and so on, until the last gapless phase, which corresponds to the semidirect product of $`J`$ WZW CFTs with $`J\times J`$ dressed charge matrices. Note that $`J`$ in this approach also denotes the number of possible Dirac seas (each of them is connected with the same magnetic field, so the excitations in each of them are not independent), and, thus, with one-half of the number of Fermi points. In the limit $`J\mathrm{}`$ (i.e. quasi-2D spin system) one obtains the (2D) Fermi surface instead of the set of 1D Fermi points (the latters become disributed more closely to each other with the grouth of $`J`$). In this limit the differences between $`\theta _j`$ tend to zero, and that is why the differences between $`\theta _{c_j}`$, $`h_{c_j}`$ and also between $`h_{c_j}`$ and $`h_s`$ disappear, too. Therefore in this limit the only $`h_s`$ survives. It means that for the quasi-2D limit of such an integrable model of $`J`$ coupled quantum spin chains for $`J\mathrm{}`$ we expect only two phases in the groundstate in an external magnetic field: the ferromagnetic gapped one and the gapless phase, which in the conformal limit corresponds to one WZW CFT (with single scalar dressed charge). The phase transition between these two phases in the groundstate in an external magnetic field is of the second order.
## 6 Behavior of the non-integrable multi-chain spin systems
So far we have studied only integrable multi-chain quantum spin models. We have shown that the commensurate-incommensurate phase transitions of the second order have to reveal themselves in an external magnetic field due to the intra-chain interactions (or the next-nearest interactions in a single quantum spin chain picture). We have shown that the emergence of these phase transitions does not depend on the value of the site spins, they emerge in the presence of the “easy-plane” magnetic anisotropy, which keeps the system in the critical (gapless) region. It is not clear, however, which features of the behavior of the integrable models with the “fine-tuned” parameters have to exist for more realistic multi-chain models, and what are the qualitative differences, we expect to exist between the integrable multi-chain models and real multi-chain spin systems.
We have to add one more thing to clarify the situation: We study (quasi) 1D spin quantum models, for which one can use the Lieb-Schultz-Mattis theorem (and its generalizations) . However, it is obvious that due to the frustration of the interactions between neighboring spins, and the presence of additional terms in the Hamiltonians, which violate the time-reversal and parity symmetries in the systems (spin chiralities or spin currents), for all spin models studied in the paper one cannot satisfy the conditions of the theorem. Hence it cannot be applied (at least directly). That is why for all the models we study there are no spin gaps (except for the trivial one for the spin-polarized groundstate). (Here we are not talking about the gaps connected with the magnetic anisotropy, but rather about the Haldane-like spin gaps which appear even for the isotropic spin-spin interaction, and about fractional magnetization plateaux ). As we argued before , namely the presence of the chiral spin terms (or the operators of the nonzero spin currents) in the Hamiltonian (which are the total time derivatives and do not change the classical equations of motion but rather affect the topological properties, like the choise of the $`\theta `$-vacuum in Haldane’s approach) is the reason why the low-lying spin excitations (and particles for lattice QFT) for our class of models are gapless and our low energy theories are conformal. It has to be mentioned that recent results of the perturbative RG analysis of the zigzag spin $`\frac{1}{2}`$ chain without three-spin terms shows the tendency the RG currents flow to the state with the parity and time-reversal violation . By the way, one can obviously see that the XY limit of the two-chain spin model does not correspond to the free fermion point of the exactly solvable model, and this coincides with the results of Ref. . Note, though, that in the latter it was erroneously concluded that the time-reversal and parity symmetries were violated by the two-chain zigzag spin Hamiltonian with only two-spin couplings (i.e. the nearest and next-nearest-neighbor interactions in the single chain picture), without spin current terms in the Hamiltonian. Hence the symmetry of the considered state was lower than the symmetry of the Hamiltonian there.
Naturally, the relevant perturbations to our integrable models will immedeately produce spin gaps. As usual, the algebraic (power-law) decay of the correlation functions in the groundstate of the models considered in this paper determines the quantum criticality. This means that, starting from the (conformal) exact solutions obtained in this paper one can argue that the response of the more realistic spin systems to perturbations can be evaluated by using perturbative methods, e.g., in a renormalization group framework. For example, let us study the effect of relevant perturbations to the Hamiltonians considered, $`\widehat{H}_r=\widehat{H}+\delta \widehat{H}_1`$, where one can choose as $`\widehat{H}_r`$, e.g., the standard Heisenberg or uniaxial Hamiltonians for several coupled quantum spin chains, and as $`\widehat{H}`$ the Hamiltonians of spin chains considered exactly in this paper for some values of $`\theta `$, where the three-spin terms are relevant. The correction to the ground state energy and the excitation gap (mass of the particle in QFT) for the quantum critical system are: $`\mathrm{\Delta }E\delta ^{(d+z)/y}`$, and $`m\delta ^{1/y}`$, respectively, where $`d`$ is the dimension of the system, and $`z`$ is the dynamical critical exponent. For the conformally invariant systems studied here one has $`d=z=1`$. The application of the standard scaling relations yields $`y+x=2(=z+d)`$, where $`x`$ is the scaling dimension, i.e. $`x=2\mathrm{\Delta }_l+2\mathrm{\Delta }_r`$, found in the previous sections (for the phases with the dressed charge matrices the summation over upper indices is meant). Hence the gap for the low-lying excitations (the mass of the physical particles in QFT) for the perturbed systems will be $`m\delta ^{1/2(1\mathrm{\Delta }_l\mathrm{\Delta }_r)}`$. Note that because of scaling, the behavior of the critical exponents (which are related to the exponents we introduced for the integrable multi-chain spin models) in the vicinities of the lines of the phase transitions has to be universal, and this can be checked experimentally. We expect that the spin gap has to exist for the values of the isotropic zigzag inter-chain coupling higher or of order of 0.5 for the two-chain spin $`\frac{1}{2}`$ system , where the three-spin couplings are relevant and the emergence of the spin gap is known exactly .
Very recently, the density matrix renormalization group numerical studies of the two-chain zigzag spin $`\frac{1}{2}`$ model (without chiral three-spin terms in the Hamiltonian) were performed . These numerical studies strongly support the picture proposed here (see also Ref. ): the magnetization as function of the magnetic field in the groundstate reveals (i) one second order phase transition (to the spin-saturation phase) for the weak intra-chain coupling; (ii) one more second order phase transition between the magnetic (gapless) phases in the intermediate region of the intra-chain coupling and (iii) in addition to those second order phase transitions, one to the gapful phase with zero magnetization (plateau) for the intra-chain coupling value of 0.5.
We should also mention that it is not the chiral spin terms (as implied in Ref. ) but the intra-chain coupling that is responsible for the commensurate-incommensurate phase transitions between the gapless phases in this class of models. As for the aforementioned spin currents, their “fine-tuned” values produce the cancellation of the spin gap for zero magnetic field . We should also note that to our mind some features of the phase diagram obtained in Ref. are artifacts of the small number of sites involved into the numerical calculations. In Fig. 5 of Ref. in the regions of $`0.52<\kappa <0.6`$ (corresponding to intra-chain couplings, normalized to the value of the inter-chain interaction, in the domain \[0.54–0.75\]) we can obviously see that when increasing the value of the magnetic field one goes from the gapped phase with zero magnetization into the gapless one with two Dirac seas of the low-lying excitations, then reaches the gapless phase with one Dirac sea, then returnes to the gapless phase with two Dirac seas, and finally reaches the spin-saturated phase. To our mind this return to the already passed phase is non-physical. One can clearly see that the region for the values of the intra-chain couplings, where these strange returns happen, is reduced when going from 16 sites in numerical calculations to 20 sites. This confirms that presently achieved sizes of the quantum systems for numerical calculations can produce even qualitatively invalid results, and analytic calculations are necessary, too.
We point out that despite the fact that the relevant perturbations in general produce a gap for the low-lying excitations, one can apply the results of this paper to the real gapless multy-chain spin systems, too. For example, it was recently observed that even for the two-leg ladder system SrCa<sub>12</sub>Cu<sub>24</sub>O<sub>41</sub> the spin gap collapses under pressure.
## 7 Concluding remarks
In this paper, motivated by recent progress in the experimental measurements for multi-chain spin systems, we have theoretically studied the behavior in an external field of a wide class of the multi-chain quantum spin models and quantum field theories. First, we have investigated the external field behavior of the exactly integrable two-chain spin $`\frac{1}{2}`$ model and have shown that the inclucion of the magnetic anisotropy of the “easy-plane” type, with which the system stays in the quantum critical region, does not qualitatively change the behavior in an external magnetic field. However, we have shown that the magnetic anisotropy changes the critical values of the magnetic fields and intra-chain couplings, at which the phase transitions occur, and affects the critical exponents. We have shown that the external-field-induced phase transitions we discussed are the commensurate-incommensurate phase transitions due to the next-nearest-neighbor two-spin interactions, which are present in these multi-chain models with zigzag-like couplings. We have pointed out that the low-lying excitations of the conformal limit of our class of multi-chain spin models are not independent in the incommensurate phase, because they are governed by the same magnetic field. We have shown that these two-chain quantum spin models share the most important features of the behavior in an external field with the wide class of the (1+1) quantum field theories.
We have introduced higher-spin versions of the two-chain exactly solvable spin models, e.g., we have investigated the important class of 1D two-chain quantum ferrimagnets with different spin values in the sites of each chain. Here we have shown that the phase transitions in an external magnetic field in this exactly solvable two-chain quantum ferrimagnet are similar in nature to the phase transitions between the spin-compensated and uncompensated phases in ordinary classical 3D ferrimagnets.
We have also studied the behavior of the multi-chain exactly solvable spin models in an external magnetic field, and shown how the additional phase transitions arising due to the increasing number of chains vanish in the quasi-2D limit. Hence, to the best of our knowledge, we have proposed the first exact scenario of the transition from 1D to 2D quantum spin models in the presence of an external magnetic field. We have argued that the commensurate-incommensurate phase transitions in the multi-chain quantum spin models have to disappear in the limit of an infinite number of chains.
Finally, we have shown how the relevant deviations from the integrability, e.g., the absence of the three-spin (spin chiral) terms in the Hamiltonians, which separately break the parity and time-reversal symmetries, give rise to gaps in spectra of the low-lying excitations of the multi-chain quantum spin systems and we have calculated the critical scaling exponents for these gaps. We pointed out the qualitative agreement of our exact analytic calculations with recent numerical simulations for zigzag spin models.
I am grateful to A. G. Izergin, S. V. Ketov , A. Klümper, V. E. Korepin, G. I. Japaridze, A. Luther and A. A. Nersesyan for helpful discussions. I thank J. Gruneberg for his kind help. The financial support of the Deutsche Forschungsgemeinschaft and Swedish Institute is acknowledged. |
warning/0002/hep-th0002010.html | ar5iv | text | # References
February 2000 UMDEPP 00-049
Superspace Geometrical Realization of the
$`N`$-Extended Super Virasoro Algebra and its Dual<sup>1</sup><sup>1</sup>1Supported in part by National Science Foundation Grant PHY-98-02551.
C. Curto <sup>2</sup><sup>2</sup>2ccurto@fas.harvard.edu
Department of Physics, Harvard University
Cambridge, MA 02138 USA
S. James Gates, Jr.<sup>3</sup><sup>3</sup>3gatess@wam.umd.edu
Department of Physics, University of Maryland
College Park, MD 20742-4111 USA
and
V.G.J. Rodgers<sup>4</sup><sup>4</sup>4vincent-rodgers@uiowa.edu
Department of Physics and Astronomy, University of Iowa
Iowa City, Iowa 52242–1479 USA
ABSTRACT
> We derive properties of $`N`$-extended $`𝒢`$ super Virasoro algebras. These include adding central extensions, identification of all primary fields and the action of the adjoint representation on its dual. The final result suggest identification with the spectrum of fields in supergravity theories and superstring/M-theory constructed from NSR $`N`$-extended supersymmetric $`𝒢`$ Virasoro algebras.
>
> PACS: 04.65.+e, 02.20.Sv, 11.30.Pb, 11.25.-w, 03.65.Fd
>
> Keywords: Super Virasoro Algebra, Coadjoint representation, Super-
> symmetry, Supergravity
(I.) Introduction
Recently super-derivations were introduced that extend previous work of 1D, $`\mathrm{}_0`$-extended superspace. This set of super-derivations is closed under graded commutation and contains a super Virasoro-like sub-algebra for all values of $`N`$-extended supersymmetry. The smallest set of the derivations that forms a closed algebra under the action of the graded commutator contains the following:
| $`G_𝒜^\mathrm{I}`$ | $`i\tau ^{𝒜+{\scriptscriptstyle \frac{1}{2}}}\left[^\mathrm{I}i\mathrm{\hspace{0.17em}2}\zeta ^\mathrm{I}_\tau \right]+2(𝒜+\frac{1}{2})\tau ^{𝒜{\scriptscriptstyle \frac{1}{2}}}\zeta ^\mathrm{I}\zeta ^\mathrm{K}_\mathrm{K},`$ |
| --- | --- |
| $`L_𝒜`$ | $`\left[\tau ^{𝒜+1}_\tau +\frac{1}{2}(𝒜+\mathrm{\hspace{0.17em}1})\tau ^𝒜\zeta ^\mathrm{I}_\mathrm{I}\right],`$ |
| $`T_𝒜^{\mathrm{I}\mathrm{J}}`$ | $`\tau ^𝒜\left[\zeta ^\mathrm{I}^\mathrm{J}\zeta ^\mathrm{J}^\mathrm{I}\right],`$ |
| $`U_𝒜^{\mathrm{I}_1\mathrm{}\mathrm{I}_q}`$ | $`i(i)^{[{\scriptscriptstyle \frac{q}{2}}]}\tau ^{(𝒜{\scriptscriptstyle \frac{\left(q2\right)}{2}})}\zeta ^{\mathrm{I}_1}\mathrm{}\zeta ^{\mathrm{I}_{q1}}^{\mathrm{I}_q},q=\mathrm{\hspace{0.17em}3},\mathrm{},N+\mathrm{\hspace{0.17em}1},`$ |
| $`R_𝒜^{\mathrm{I}_1\mathrm{}\mathrm{I}_p}`$ | $`(i)^{[{\scriptscriptstyle \frac{p}{2}}]}\tau ^{(𝒜{\scriptscriptstyle \frac{\left(p2\right)}{2}})}\zeta ^{\mathrm{I}_1}\mathrm{}\zeta ^{\mathrm{I}_p}_\tau ,p=\mathrm{\hspace{0.17em}2},\mathrm{},N,`$ |
$`(1)`$
for any number $`N`$ of supersymmetries. (Our notational conventions can be found in .) These derivations do not depend on a specific value of $`N`$ and can therefore be used for the entire 1D, $`\mathrm{}_0`$ superspace. For low values of $`N`$, not all of the generators appear. For example, $`T_𝒜^{\mathrm{I}\mathrm{J}}`$ and $`R_𝒜^{\mathrm{I}_1\mathrm{}\mathrm{I}_p}`$ only appear for superspaces with $`N`$ $``$ 2. Generically, $`U_𝒜^{\mathrm{I}_1\mathrm{}\mathrm{I}_q}`$ only appears for superspaces with $`N`$ $``$ 3. The indices denoted by $`𝒜`$, $``$, etc. denote the level or mode number of the operators. These types of indices take their values in $`Z+\frac{1}{2}Z`$.
One of the tasks of this paper is to centrally extend the algebra generated by the above generators. We impose the Jacobi identity on all possible combinations of the generators and find that the centrally extended algebra is given by<sup>5</sup><sup>5</sup>5Some of the results here contain minor corrections to those given in .
| $`[L_𝒜,L_{}\}`$ | $`=(𝒜)L_{𝒜+}+\frac{1}{8}c(𝒜^3𝒜)\delta _{𝒜+,0},`$ |
| --- | --- |
| $`[L_𝒜,U_{}^{\mathrm{I}_1\mathrm{}\mathrm{I}_m}\}`$ | $`=[+\frac{1}{2}(m2)𝒜]U_{𝒜+}^{\mathrm{I}_1\mathrm{}\mathrm{I}_m},`$ |
| $`[G_𝒜{}_{}{}^{\mathrm{I}},G_{}{}_{}{}^{\mathrm{J}}\}`$ | $`=i\mathrm{\hspace{0.17em}4}\delta ^{\mathrm{I}\mathrm{J}}L_{𝒜+}i2(𝒜)[T_{𝒜+}^{\mathrm{I}\mathrm{J}}+2(𝒜+)U_{𝒜+}^{\mathrm{I}\mathrm{J}\mathrm{K}}{}_{\mathrm{K}}{}^{}]`$ |
| | $`ic(𝒜^2\frac{1}{4})\delta _{𝒜+,0}\delta ^{\mathrm{I}\mathrm{J}},`$ |
| $`[L_𝒜,G_{}{}_{}{}^{\mathrm{I}}\}`$ | $`=(\frac{1}{2}𝒜)G_{𝒜+}{}_{}{}^{\mathrm{I}},`$ |
| $`[L_𝒜,R_{}^{\mathrm{I}_1\mathrm{}\mathrm{I}_m}\}`$ | $`=[+\frac{1}{2}(m2)𝒜]R_{𝒜+}^{\mathrm{I}_1\mathrm{}\mathrm{I}_m}`$ |
| | $`[\frac{1}{2}𝒜(𝒜+1)]U_{𝒜+}^{\mathrm{I}_1\mathrm{}\mathrm{I}_mJ}{}_{J}{}^{},`$ |
| $`[L_𝒜,T_{}^{\mathrm{I}\mathrm{J}}\}`$ | $`=T_{𝒜+}^{\mathrm{I}\mathrm{J}},`$ |
| $`[R_𝒜^{\mathrm{I}_1\mathrm{}\mathrm{I}_m},R_{}^{\mathrm{J}_1\mathrm{}\mathrm{J}_n}\}`$ | $`=(i)^{\sigma (mn)}[𝒜\frac{1}{2}(mn)]R_{𝒜+}^{\mathrm{I}_1\mathrm{}\mathrm{I}_m\mathrm{J}_1n\mathrm{}\mathrm{J}_n},`$ |
| --- | --- |
| $`[T_𝒜^{\mathrm{I}\mathrm{J}},T_{}^{\mathrm{K}\mathrm{L}}\}`$ | $`=T_{𝒜+}^{\mathrm{I}\mathrm{K}}\delta ^{\mathrm{JL}}+T_{𝒜+}^{\mathrm{J}\mathrm{L}}\delta ^{\mathrm{IK}}T_{𝒜+}^{\mathrm{I}\mathrm{L}}\delta ^{\mathrm{JK}}T_{𝒜+}^{\mathrm{J}\mathrm{K}}\delta ^{\mathrm{IL}}`$ |
| | $`+\stackrel{~}{c}(𝒜)(\delta ^{\mathrm{IK}}\delta ^{\mathrm{JL}}\delta ^{\mathrm{IL}}\delta ^{\mathrm{JK}})\delta _{𝒜+,0},`$ |
| $`[G_𝒜{}_{}{}^{\mathrm{I}},R_{}^{\mathrm{J}_1\mathrm{}\mathrm{J}_m}\}`$ | $`=2(i)^{\sigma (m)}[+(m1)𝒜+\frac{1}{2}]R_{𝒜+}^{\mathrm{I}\mathrm{J}_1\mathrm{}\mathrm{J}_m}`$ |
| | $`(i)^{\sigma (m)}{\displaystyle \underset{r=1}{\overset{m}{}}}(1)^{r1}\delta ^{I\mathrm{J}_r}R_{𝒜+}^{\mathrm{J}_1\mathrm{}\mathrm{J}_{r1}\mathrm{J}_{r+1}\mathrm{}\mathrm{J}_m}`$ |
| | $`(i)^{\sigma (m)}[𝒜+\frac{1}{2}]U_{𝒜+}^{\mathrm{J}_1\mathrm{}\mathrm{J}_m\mathrm{I}}`$ |
| | $`+\mathrm{\hspace{0.17em}2}(i)^{\sigma (m)}[𝒜^2\frac{1}{4}]U_{𝒜+}^{I\mathrm{J}_1\mathrm{}\mathrm{J}_mK}{}_{K}{}^{},`$ |
| $`[G_𝒜{}_{}{}^{\mathrm{I}},U_{}^{\mathrm{J}_1\mathrm{}\mathrm{J}_m}\}`$ | $`=2(i)^{\sigma (m)}[+(m2)𝒜]U_{𝒜+}^{\mathrm{I}\mathrm{J}_1\mathrm{}\mathrm{J}_m}`$ |
| | $`\mathrm{\hspace{0.17em}2}(i)^{\sigma (m)}[𝒜+\frac{1}{2}]\delta ^{\mathrm{I}\mathrm{J}_m}U_{𝒜+}^{\mathrm{J}_1\mathrm{}\mathrm{J}_{m1}K}_K`$ |
| | $`(i)^{\sigma (m)}{\displaystyle \underset{r=1}{\overset{m1}{}}}(1)^{r1}\delta ^{\mathrm{I}\mathrm{J}_r}U_{𝒜+}^{\mathrm{J}_1\mathrm{}\mathrm{J}_{r1}\mathrm{J}_{r+1}\mathrm{}\mathrm{J}_m}`$ |
| | $`+\mathrm{\hspace{0.17em}2}(i)^{\sigma (m)}\delta ^{\mathrm{I}\mathrm{J}_m}R_{𝒜+}^{\mathrm{J}_1\mathrm{}\mathrm{J}_{m1}},`$ |
| $`[R_𝒜^{\mathrm{I}_1\mathrm{}\mathrm{I}_m},U_{}^{\mathrm{J}_1\mathrm{}\mathrm{J}_n}\}`$ | $`=(i)^{\sigma (mn)}{\displaystyle \underset{r=1}{\overset{m}{}}}(1)^{r1}\delta ^{\mathrm{I}_r\mathrm{J}_n}R_{𝒜+}^{\mathrm{J}_1\mathrm{}\mathrm{J}_{n1}\mathrm{I}_1\mathrm{}\mathrm{I}_{r1}\mathrm{I}_{r+1}\mathrm{}\mathrm{I}_m}`$ |
| | $`+i(i)^{\sigma (mn)}[\frac{1}{2}(n2)]U_{𝒜+}^{\mathrm{I}_1\mathrm{}\mathrm{I}_m\mathrm{J}_1\mathrm{}\mathrm{J}_n},`$ |
| $`[U_𝒜^{\mathrm{I}_1\mathrm{}\mathrm{I}_m},U_{}^{\mathrm{J}_1\mathrm{}\mathrm{J}_n}\}`$ | $`=(i)^{\sigma (mn)}\{{\displaystyle \underset{r=1}{\overset{m}{}}}(1)^{r1}\delta ^{\mathrm{I}_m\mathrm{J}_r}U_{𝒜+}^{\mathrm{I}_1\mathrm{}\mathrm{I}_{m1}\mathrm{J}_1\mathrm{}\mathrm{J}_{r1}\mathrm{J}_{r+1}\mathrm{}\mathrm{J}_{n1}\mathrm{J}_n}`$ |
| | $`(1)^{mn}{\displaystyle \underset{r=1}{\overset{m}{}}}(1)^{r1}\delta ^{\mathrm{I}_r\mathrm{J}_n}U_{𝒜+}^{\mathrm{J}_1\mathrm{}\mathrm{J}_{n1}\mathrm{I}_1\mathrm{}\mathrm{I}_{r1}\mathrm{I}_{r+1}\mathrm{}\mathrm{I}_{m1}\mathrm{I}_m}\},`$ |
| $`[T_𝒜^{\mathrm{I}\mathrm{J}},G_{}{}_{}{}^{\mathrm{K}}\}`$ | $`=2(\delta ^{\mathrm{JK}}G_{𝒜+}{}_{}{}^{\mathrm{I}}\delta ^{\mathrm{IK}}G_{𝒜+}{}_{}{}^{\mathrm{J}})`$ |
| | $`+2𝒜(\delta ^{\mathrm{JK}}U_{𝒜+}^{\mathrm{I},\mathrm{L}}{}_{\mathrm{L}}{}^{}\delta ^{\mathrm{IK}}U_{𝒜+}^{\mathrm{J}\mathrm{L}}{}_{\mathrm{L}}{}^{}+U_{𝒜+}^{\mathrm{J}\mathrm{K}\mathrm{I}}U_{𝒜+}^{\mathrm{I}\mathrm{K}\mathrm{J}}),`$ |
| $`[T_𝒜^{\mathrm{I}\mathrm{J}},R_{}^{\mathrm{I}_1\mathrm{}\mathrm{I}_p}\}`$ | $`={\displaystyle \underset{r=1}{\overset{p}{}}}(1)^{r+1}(\delta ^{\mathrm{J}\mathrm{I}_r}R_{𝒜+}^{\mathrm{I}\mathrm{I}_1\mathrm{}\mathrm{I}_{r1}\mathrm{I}_{r+1}\mathrm{}\mathrm{I}_p}\delta ^{\mathrm{I}\mathrm{I}_r}R_{𝒜+}^{\mathrm{J}\mathrm{I}_1\mathrm{}\mathrm{I}_{r1}\mathrm{I}_{r+1}\mathrm{}\mathrm{I}_p})`$ |
| | $`+i(1)^p𝒜(U_{𝒜+1}^{\mathrm{I}\mathrm{I}_1\mathrm{}\mathrm{I}_\mathrm{p}\mathrm{J}}U_{𝒜+1}^{\mathrm{J}\mathrm{I}_1\mathrm{}\mathrm{I}_\mathrm{p}\mathrm{I}}),`$ |
| $`[T_𝒜^{\mathrm{I}\mathrm{J}},U_{}^{\mathrm{I}_1\mathrm{}\mathrm{I}_p}\}`$ | $`={\displaystyle \underset{r=1}{\overset{p}{}}}(1)^{r+1}(\delta ^{\mathrm{J}\mathrm{I}_r}U_{𝒜+}^{\mathrm{I}\mathrm{I}_1\mathrm{}\mathrm{I}_{r1}\mathrm{I}_{r+1}\mathrm{}\mathrm{I}_p}\delta ^{\mathrm{I}\mathrm{I}_r}U_{𝒜+}^{\mathrm{J}\mathrm{I}_1\mathrm{}\mathrm{I}_{r1}\mathrm{I}_{r+1}\mathrm{}\mathrm{I}_p})`$ |
| | $`+(\delta ^{\mathrm{I}_\mathrm{p}\mathrm{I}}U_{𝒜+1}^{\mathrm{I}_1\mathrm{}\mathrm{I}_{\mathrm{p}1}\mathrm{J}}\delta ^{\mathrm{I}_\mathrm{p}\mathrm{J}}U_{𝒜+1}^{\mathrm{I}_1\mathrm{}\mathrm{I}_{\mathrm{p}1}\mathrm{I}}),`$ |
$`(2)`$
where the function $`\sigma (m)=0`$ if $`m`$ is even and $`1`$ if $`m`$ is odd. Here the central extensions $`c`$ and $`\stackrel{~}{c}`$ are unrelated since we have only imposed the Jacobi identity. New constraints will arise when we restrict to unitary representations. The algebra exhibits interesting properties such as a generalization of the SO(N) generators due to the presence of the $`U`$ and $`R`$ type fields. The nature of these fields will be discussed throughout as we derive transformation laws.
One way to understand the operators that arise in this new algebra is through methods used to study other infinite dimensional algebras. In particular, we will borrow techniques from coadjoint representation to help interpret these new generators. The coadjoint representation for infinite dimensional algebras has appeared in the string theory literature for some time. Its uses include the study of chiral anomalies , geometric quantization of the Virasoro group , the study of orthogonal field theories and recently in relation to $`\mathrm{AdS}_3`$ quantum gravity .
In this paper we will examine the coadjoint representation of the superspace geometrical representation (“$`𝒢`$”) of the extended super Virasoro algebras as well as some other properties. Although the algebra is quite complex, the coadjoint representation of this particular algebra can generalize many of the above mentioned aspects as well as shed light on the meaning of the new generators and the spectrum of states that may appear in a supergravity or superstring/M-theory based on this algebra.
(II.) Primary Fields
Before going into the coadjoint representation, we would like to identify the primary fields associated with this algebra and their associated conformal weights. Since $``$ is the generator of diffeomorphisms we can use its action on the other generators to determine the tensor properties of the fields. Let
$$^{}=(L_\xi ,G_{\chi ^\mathrm{I}}{}_{}{}^{\mathrm{I}},T_{t^{\mathrm{J}\mathrm{K}}}^{\mathrm{J}\mathrm{K}},_{\{\mathrm{I}_\mathrm{q}\}}^{}U_{\mu ^{\{\mathrm{I}_\mathrm{q}\}}}^{\{\mathrm{I}_\mathrm{q}\}},_{\{\mathrm{J}_\mathrm{q}\}}^{}R_{r^{\{\mathrm{J}_\mathrm{q}\}}}^{\{\mathrm{J}_\mathrm{q}\}};\alpha ),$$
$`(3)`$
represent the generators with generic functions and $`_{\{\mathrm{I}_\mathrm{q}\}}`$ represents the direct sum over all distinct generators. Then from the algebra we see that
| $`[(L_\xi ,\alpha ),(L_\zeta ,\beta )\}`$ | $`=(L_{\xi ^{}\zeta \xi \zeta ^{}},\frac{c}{i2\pi }{\displaystyle (\xi ^{\prime \prime }\zeta ^{}\zeta ^{\prime \prime }\xi ^{})𝑑x}),`$ |
| --- | --- |
| $`[L_\xi ,G_{\chi ^\mathrm{I}}^\mathrm{I}\}`$ | $`=G_{(\xi (\chi ^\mathrm{I})^{}+{\scriptscriptstyle \frac{1}{2}}\xi ^{}\chi ^\mathrm{I})}^\mathrm{I},`$ |
| $`[L_\xi ,T_{t^{\mathrm{RS}}}^{\mathrm{RS}}\}`$ | $`=T_{(\xi (t^{\mathrm{RS}})^{})}^{\mathrm{RS}},`$ |
| $`[L_\xi ,U_{w^{\{\mathrm{V}_\mathrm{r}\}}}^{\{\mathrm{V}_\mathrm{r}\}}\}`$ | $`=U_{(\xi (w^{\{\mathrm{V}_\mathrm{r}\}})^{}{\scriptscriptstyle \frac{1}{2}}(r2)\xi ^{}w^{\{\mathrm{V}_\mathrm{r}\}})}^{\{\mathrm{V}_\mathrm{r}\}},`$ |
| --- | --- |
| $`[L_\xi ,R_{\rho ^{\{\mathrm{T}_\mathrm{r}\}}}^{\{\mathrm{T}_\mathrm{r}\}}\}`$ | $`=R_{((\rho ^{\{\mathrm{T}_\mathrm{r}\}})^{}\xi {\scriptscriptstyle \frac{1}{2}}(r2)\xi ^{}(\rho ^{\{\mathrm{T}_\mathrm{r}\}}))}^{\{\mathrm{T}_\mathrm{r}\}}\frac{i}{2}U_{(\xi ^{\prime \prime }\rho ^{\{\mathrm{T}_\mathrm{r}\}})}^{\{\mathrm{T}_\mathrm{r}\}},`$ |
$`(4)`$
determines the transformation laws of the functions. In the above, we have suppressed the Grassman indices. For example $`w^{\overline{\mathrm{V}}_1\mathrm{}\overline{\mathrm{V}}_\mathrm{n}}`$ the function associated with the $`U`$ generators may be written as $`w^{\{\mathrm{V}_\mathrm{m}\}}`$ or simply as $`w`$. From the coefficient of the $`\xi ^{}`$ summand in the transformation laws we can write down the conformal weight which is also the rank of the one dimensional tensors. The quantity $`\xi `$ is a rank one contravariant tensor field, $`\chi ^\mathrm{I}`$ is a spin half field and $`t^{\mathrm{RS}}`$ is a scalar field. This is to be expected from these fields. However, notice that $`w^{\{\mathrm{V}_\mathrm{r}\}}`$ transforms with conformal weight $`\frac{1}{2}(r2)`$ where $`r`$ takes values from $`3`$ to $`N`$ when there are $`N`$ supersymmetries which corresponds to a tower of $`N2`$ fields. The quantity $`\rho ^{\{\mathrm{T}_\mathrm{r}\}}`$ to appears to transform as a rank $`\frac{1}{2}(r2)`$ tensor modulo the inhomogeneous term. However, it is this inhomogeneous term that keeps these fields from transforming like tensors. Since the transformations of $`w^{\{\mathrm{V}_\mathrm{r}\}}`$ and $`\rho ^{\{\mathrm{T}_\mathrm{r}\}}`$ are entangled, a natural question to ask is what linear combination of generators produces tensors or in the language of conformal field theory which generators are primary.
To answer this let us consider the generators
$$𝒬_𝒜^{\mathrm{I}_1\mathrm{}\mathrm{I}_\mathrm{p}}=\tau ^𝒜\zeta ^{\mathrm{I}_1}\mathrm{}\zeta ^{\mathrm{I}_\mathrm{p}}_\tau ,𝒫_𝒜^{\mathrm{I}_1\mathrm{}\mathrm{I}_{\mathrm{p}+1}}=\tau ^𝒜\zeta ^{\mathrm{I}_1}\mathrm{}\zeta ^{\mathrm{I}_\mathrm{p}}^{\mathrm{I}_{\mathrm{p}+1}}.$$
$`(5)`$
These form a closed algebra among themselves and are used to facilitate the computations below. We can write the previous generators as
| $`L_𝒜=𝒬_{𝒜+1}\frac{1}{2}(𝒜+1)𝒫_𝒜^\mathrm{I}{}_{\mathrm{I}}{}^{},`$ |
| --- |
| $`G_𝒜^\mathrm{I}=i𝒫_{𝒜+{\scriptscriptstyle \frac{1}{2}}}^\mathrm{I}+2𝒬_{𝒜+{\scriptscriptstyle \frac{1}{2}}}^\mathrm{I}+2(𝒜+\frac{1}{2})𝒫_{𝒜{\scriptscriptstyle \frac{1}{2}}}^{\mathrm{I}\mathrm{K}}{}_{\mathrm{K}}{}^{},`$ |
| $`T_𝒜^{\mathrm{I}\mathrm{J}}=𝒫_𝒜^{IJ}𝒫_𝒜^{JI},`$ |
| $`R_𝒜^{\mathrm{I}_1\mathrm{}\mathrm{I}_\mathrm{p}}=i^{[{\scriptscriptstyle \frac{\mathrm{p}}{2}}]}𝒬_{𝒜{\scriptscriptstyle \frac{\mathrm{p}}{2}}+1}^{\mathrm{I}_1\mathrm{}\mathrm{I}_\mathrm{p}},(\mathrm{p}=2,\mathrm{},N),`$ |
| $`U_𝒜^{\mathrm{I}_1\mathrm{}\mathrm{I}_\mathrm{q}}=i(i)^{[{\scriptscriptstyle \frac{\mathrm{q}}{2}}]}𝒫_{𝒜{\scriptscriptstyle \frac{\mathrm{p}}{2}}+1}^{\mathrm{I}_1\mathrm{}\mathrm{I}_\mathrm{q}},(\mathrm{q}=3,\mathrm{},N+1).`$ |
$`(6)`$
Let $`_𝒜^{\mathrm{I}_1\mathrm{}\mathrm{I}_\mathrm{p}}`$ be a primary generator. Then by definition for some particular mode dependent $`\lambda `$ this generator satisfies
$$[L_𝒜,_{}^{\mathrm{I}_1\mathrm{}\mathrm{I}_\mathrm{p}}]=\lambda (𝒜,,\mathrm{p})_{𝒜+}^{\mathrm{I}_1\mathrm{}\mathrm{I}_\mathrm{p}},$$
$`(7)`$
for fixed number of indices p. For each value of p (assuming that p is greater than 2), $`_𝒜^{\mathrm{I}_1\mathrm{}\mathrm{I}_\mathrm{p}}`$ can be generically written as
$$_𝒜^{\mathrm{I}_1\mathrm{}\mathrm{I}_\mathrm{p}}=c_0(𝒜+1)𝒬_{\mathrm{A}+1}^{\mathrm{I}_1\mathrm{}\mathrm{I}_\mathrm{p}}+c_1(𝒜+1)𝒫_{𝒜+1}^{\mathrm{I}_1\mathrm{}\mathrm{I}_\mathrm{p}}+c_2(𝒜)𝒫_𝒜^{\mathrm{I}_1\mathrm{}\mathrm{I}_\mathrm{p}\mathrm{L}}{}_{\mathrm{L}}{}^{},$$
$`(8)`$
which gives us three possible mode dependent coefficients to compute, viz $`c_0,c_1,`$ and $`c_2`$. From the commutation relations of $`𝒬_𝒜^{\mathrm{I}_1\mathrm{}\mathrm{I}_\mathrm{p}}`$ and $`𝒫_𝒜^{\mathrm{I}_1\mathrm{}\mathrm{I}_\mathrm{p}}`$ the conditions for a primary generator are
| $`\lambda c_0(𝒜++1)`$ | $`=c_0(+1)[(𝒜)+\frac{1}{2}\mathrm{p}(𝒜+1)],`$ |
| --- | --- |
| $`\lambda c_1(𝒜++1)`$ | $`=c_1(+1)[(𝒜)+\frac{1}{2}\mathrm{p}(𝒜+1)],`$ |
| $`\lambda c_2(𝒜+)`$ | $`=c_0(+1)[\frac{\mathrm{p}}{2}(𝒜+1)].`$ |
$`(9)`$
There are three classes of solutions.
1. Class 1:
Setting $`c_0=1,c_1=0`$ and $`c_2(𝒜)=a_1𝒜+a_2`$ we find that
$$_{}^{1}{}_{}{}^{\mathrm{I}_1\mathrm{}\mathrm{I}_\mathrm{p}}=𝒬_{+1}^{\mathrm{I}_1\mathrm{}\mathrm{I}_\mathrm{p}}+(\frac{+1}{2\mathrm{p}})𝒫_{}^{\mathrm{I}_1\mathrm{}\mathrm{I}_\mathrm{p}\mathrm{L}}{}_{\mathrm{L}}{}^{},$$
$`(10)`$
is a primary field with $`\lambda =(+\frac{\mathrm{p}}{2})𝒜(1\frac{\mathrm{p}}{2}).`$ This can be rewritten in terms of the original generators as
$$R_{+{\scriptscriptstyle \frac{\mathrm{p}}{2}}}^{\mathrm{I}_1\mathrm{}\mathrm{I}_\mathrm{p}}i(\frac{+1}{2\mathrm{p}})U_{+{\scriptscriptstyle \frac{p}{2}}1}^{\mathrm{I}_1\mathrm{}\mathrm{I}_\mathrm{p}\mathrm{L}}{}_{\mathrm{L}}{}^{},(\mathrm{p}2)$$
$`(11)`$
is a primary field.
2. Class 2:
Setting $`c_0=0,c_1=1`$ and $`c_2(𝒜)=a_1𝒜+a_2`$ we find that
$$_{}^{2}{}_{}{}^{\mathrm{I}\mathrm{J}}=𝒫_{+1}^{\mathrm{IJ}}+(+1)𝒫_{}^{\mathrm{I}\mathrm{J}\mathrm{K}}{}_{\mathrm{K}}{}^{}.$$
$`(12)`$
In this case $`\mathrm{p}=2`$ was forced as a condition, thus the above $`\lambda `$ simplifies to $`\lambda =+1`$.
3. Class 3:
Setting $`c_2=0,c_1=1`$ and $`c_0(𝒜)=a_1𝒜+a_2`$ we find that
$$_{}^{3}{}_{}{}^{\mathrm{I}_1\mathrm{}\mathrm{I}_\mathrm{p}}=𝒫_{+1}^{\mathrm{I}_1\mathrm{}\mathrm{I}_\mathrm{p}},$$
$`(13)`$
is a primary field. In this case $`c_0=0`$ was forced as a condition.
From these three solutions we can find all the primary fields of the original algebra.
| $`L_{}`$ | $`=_{}^{1}{}_{}{}^{},(\mathrm{p}=0)`$ |
| --- | --- |
| $`G_{}^\mathrm{I}`$ | $`=2_{}^{1}{}_{{\scriptscriptstyle \frac{1}{2}}}{}^{\mathrm{I}}+i_{}^{3}{}_{{\scriptscriptstyle \frac{1}{2}}}{}^{\mathrm{I}},(\mathrm{p}=1)`$ |
| $`T_{}^{\mathrm{I}\mathrm{J}}`$ | $`=_{}^{3}{}_{1}{}^{\mathrm{I}\mathrm{J}}_{}^{3}{}_{1}{}^{\mathrm{J}\mathrm{I}},(\mathrm{p}=2)`$ |
| $`U_{+{\scriptscriptstyle \frac{\mathrm{p}}{2}}}^{\mathrm{I}_1\mathrm{}\mathrm{I}_\mathrm{p}}`$ | $`=ii^{[{\scriptscriptstyle \frac{\mathrm{p}}{2}}]}_{}^{3}{}_{}{}^{\mathrm{I}_1\mathrm{}\mathrm{I}_\mathrm{p}},(\mathrm{p}3)`$ |
| --- | --- |
| $`R_{+{\scriptscriptstyle \frac{\mathrm{p}}{2}}}^{\mathrm{I}_1\mathrm{}\mathrm{I}_\mathrm{p}}`$ | $`i(\frac{+1}{2\mathrm{p}})U_{+{\scriptscriptstyle \frac{p}{2}}1}^{\mathrm{I}_1\mathrm{}\mathrm{I}_\mathrm{p}\mathrm{L}}{}_{\mathrm{L}}{}^{}=(i)^{[{\scriptscriptstyle \frac{\mathrm{p}}{2}}]}_{}^{1}{}_{}{}^{\mathrm{I}_1\mathrm{}\mathrm{I}_\mathrm{p}},(\mathrm{p}3).`$ |
$`(14)`$
We note that $`_𝒜^{IJ}`$ for no value of $`𝒜`$ admits a primary field. Stated in a slightly different way, the set of generators given by
| $`G_𝒜^\mathrm{I}`$ | $`i\tau ^{𝒜+{\scriptscriptstyle \frac{1}{2}}}\left[^\mathrm{I}i\mathrm{\hspace{0.17em}2}\zeta ^\mathrm{I}_\tau \right]+2(𝒜+\frac{1}{2})\tau ^{𝒜{\scriptscriptstyle \frac{1}{2}}}\zeta ^\mathrm{I}\zeta ^\mathrm{K}_\mathrm{K},`$ |
| --- | --- |
| $`L_𝒜`$ | $`\left[\tau ^{𝒜+1}_\tau +\frac{1}{2}(𝒜+\mathrm{\hspace{0.17em}1})\tau ^𝒜\zeta ^\mathrm{I}_\mathrm{I}\right],`$ |
| $`T_𝒜^{\mathrm{I}\mathrm{J}}`$ | $`\tau ^𝒜\left[\zeta ^\mathrm{I}^\mathrm{J}\zeta ^\mathrm{J}^\mathrm{I}\right],`$ |
| $`U_𝒜^{\mathrm{I}_1\mathrm{}\mathrm{I}_q}`$ | $`i(i)^{[{\scriptscriptstyle \frac{q}{2}}]}\tau ^{(𝒜{\scriptscriptstyle \frac{\left(q2\right)}{2}})}\zeta ^{\mathrm{I}_1}\mathrm{}\zeta ^{\mathrm{I}_{q1}}^{\mathrm{I}_q},q=\mathrm{\hspace{0.17em}3},\mathrm{},N+\mathrm{\hspace{0.17em}1},`$ |
| $`_𝒜^{\mathrm{I}_1\mathrm{}\mathrm{I}_p}`$ | $`(i)^{[{\scriptscriptstyle \frac{p}{2}}]}\tau ^{(𝒜{\scriptscriptstyle \frac{p}{2}})}\zeta ^{\mathrm{I}_1}\mathrm{}\zeta ^{\mathrm{I}_p}[\tau _\tau +({\displaystyle \frac{𝒜+1}{p2}})\zeta ^L_L],p=\mathrm{\hspace{0.17em}3},\mathrm{},N,`$ |
| $`R_𝒜^{\mathrm{IJ}}`$ | $`i\tau ^𝒜\zeta ^\mathrm{I}\zeta ^\mathrm{J}_\tau ,`$ |
$`(15)`$
possesses only one non-primary generator, namely $`R_𝒜^{\mathrm{IJ}}`$. We will refer to this basis as the “almost primary basis” for the $`𝒢`$ super-Virasoro algebra.
(III.) The Coadjoint Representation
In this paper we will examine the coadjoint representation of the superspace geometrical representation of the extended super Virasoro algebras as well as some other properties. Although the algebra is quite complex the coadjoint representation of this particular algebra can generalize many of the above mentioned aspects as well as shed light on the meaning of the new generators in the algebra.
(III.a)An Example:
To begin we will use the semi-direct product of the Virasoro algebra and an affine Lie algebra on the circle to fix the notation and familiarity of the coadjoint representation. In this case we have an affine Lie algebra associated with the loop group G together with the Virasoro algebra given by
| $`[J_N^\alpha ,J_M^\beta ]`$ | $`=if^{\alpha \beta \gamma }J_{N+M}^\gamma +Nk\delta _{M+N,0}\delta ^{\alpha \beta },`$ |
| --- | --- |
| $`[L_N,J_M^\alpha ]`$ | $`=MJ_{M+N}^\alpha ,`$ |
| $`[L_N,L_M]`$ | $`=(NM)L_{N+M}+{\displaystyle \frac{\widehat{c}}{12}}(N^3N)\delta _{N+M,0},`$ |
$`(16)`$
where $`\widehat{c}=\frac{2k\mathrm{Dim}(G)}{2k+c_v}`$, $`\mathrm{Dim}(G)`$ is the dimension of the group and $`c_v`$ is the value of the quadratic Casimir in the adjoint representation. Let $`(L_A,J_B^\beta ,\rho )`$ denote a centrally extended adjoint vector. Then from the commutation relations above one may write the adjoint action on the adjoint vectors as
| $`(L_A,J_B^\beta ,\rho )(L_N^{},J_M^{}^\alpha ^{},\mu )=`$ | |
| --- | --- |
| $`((AN^{})L_{A+N^{}},M^{}J_{A+M^{}}^\alpha ^{}+BJ_{B+N^{}}^\beta +if^{\beta \alpha ^{}\lambda }J_{B+M^{}}^\lambda `$ | , |
| $`{\displaystyle \frac{\widehat{c}}{12}}(A^3A)\delta _{A+N^{},0}`$ | $`+Bk\delta ^{\alpha ^{}\beta }\delta _{B+M^{},0}).`$ |
$`(17)`$
Now let $`(\stackrel{~}{L}_N,\stackrel{~}{J}_M^\alpha ,\stackrel{~}{\mu })`$ denote an element of the dual space of the algebra and let
$$(\stackrel{~}{L}_N,\stackrel{~}{J}_M^\alpha ,\stackrel{~}{\mu })|(L_N^{},J_M^{}^\alpha ^{},\mu ^{})=\delta ^{N,N^{}}+\delta ^{\alpha ,\alpha ^{}}\delta _{M,M^{}}+\mu \stackrel{~}{\mu },$$
$`(18)`$
define a suitable pairing. By requiring that this pairing be an invariant under the action of any of the adjoint elements, say $`(L_A,J_B^\beta ,\rho )`$, the coadjoint representation can be defined. The adjoint action acts as a derivation so that by Leibnitz rule one has
| | $`(\stackrel{~}{L}_N,\stackrel{~}{J}_M^\alpha ,\stackrel{~}{\mu })|(L_A,J_B^\beta ,\rho )(L_N^{},J_M^{}^\alpha ^{},\mu )=`$ |
| --- | --- |
| $``$ | $`(L_A,J_B^\beta ,\rho )(\stackrel{~}{L}_N,\stackrel{~}{J}_M^\alpha ,\stackrel{~}{\mu })|(L_N^{},J_M^{}^\alpha ^{},\mu ).`$ |
$`(19)`$
Thus the transformation properties of the coadjoint vectors are defined through,
| $`(L_A,J_B^\beta ,\rho )(\stackrel{~}{L}_N,\stackrel{~}{J}_M^\alpha ,\stackrel{~}{\mu })=`$ |
| --- |
| $`((2AN)\stackrel{~}{L}_{NA}B\delta ^{\alpha \beta }\stackrel{~}{L}_{MB}{\displaystyle \frac{\stackrel{~}{\mu }\widehat{c}}{12}}(A^3A)\stackrel{~}{L}_A,`$ |
| $`(MA)\stackrel{~}{J}_{MA}^\alpha if^{\beta \nu \alpha }\stackrel{~}{J}_{MB}^\nu \stackrel{~}{\mu }Bk\stackrel{~}{J}_B^\beta ,0).`$ |
$`(20)`$
Instead of using components, let us write $`F=(f(\theta ),\widehat{h}(\theta ),a)`$ as an arbitrary adjoint vector and $`B=(b(\theta ),h(\theta ),\mu )`$ as an arbitrary coadjoint vector, where $`f,\widehat{h},b,`$ and $`h`$ are functions. For the algebra we choose the realization,
| $`L_N`$ | $`=i\mathrm{exp}(iN\theta )_\theta ,`$ |
| --- | --- |
| $`J_N^\alpha `$ | $`=\tau ^\alpha \mathrm{exp}(iN\theta ),`$ |
$`(21)`$
and normalize the generators so that $`\mathrm{Tr}(\tau ^\alpha \tau ^\beta )=\delta ^{\alpha \beta }`$. Then Eq.(20) may be written as
| | $`\delta _FB(f(\theta ),\widehat{h}(\theta ),a)(b(\theta ),h(\theta ),\mu )=`$ |
| --- | --- |
| $``$ | $`(2f^{}b+b^{}f+i{\displaystyle \frac{\widehat{c}\mu }{12}}f^{\prime \prime \prime }+\mathrm{Tr}[h\widehat{h}^{}],h^{}f+hf^{}+[\widehat{h}hh\widehat{h}]+ik\mu \widehat{h}^{},0),`$ |
$`(22)`$
where denotes $`_\theta `$. The above equation provides an interpretation of the adjoint elements and coadjoint elements in terms of physical fields in one dimension . We already know that the Virasoro sector transforms functions as one dimensional
coordinate transformations (up to central extension). For example $`b`$ transforms as a rank two tensor field in one dimension where the infinitesimal coordinate transformation is given by $`f`$. From the second element of the triplet in Eq.(22), on sees that the function $`h`$ transforms as a one dimensional gauge field with gauge parameter $`\widehat{h}`$. The $`h^{}f+hf^{}`$ contribution to the transformation of $`h`$ simply shows that the field $`h`$ transforms as a rank one covariant tensor. The peculiar transformation is the $`Tr[h\widehat{h}^{}]`$ that appears in the transformation of $`b`$. This suggests that this rank two tensor can be shifted by fields built purely from the gauge sector. In such terms are interpreted as coming from an interaction Lagrangian. In any case the relationship between different members of the algebra juxtaposed to the dual space becomes manifest through the coadjoint representation.
Those adjoint vectors, $`F`$, that leave $`B`$ invariant will generate the isotropy group for $`B`$. Setting Eq.(1a) to zero determines the isotropy equation for $`B`$. Equation (1a) then determines the tangent space on the orbit of $`B`$. Thus for coadjoint elements $`B_1`$ and $`B_2`$, we may construct the symplectic two form by writing
$$\mathrm{\Omega }_B(B_1,B_2)=B[F_1,F_2],$$
$`(23)`$
where for example $`\delta _{F_1}B=B_1`$. In the equation of isotropy is related to constraint equations that come from a two dimensional field theory.
(III.b) $`N`$-Extended $`𝒢`$ Super Virasoro Algebra Dual Space
To proceed we will let $`\overline{}`$ denote a generic coadjoint vector and let $``$ and $`^{}`$ denote adjoint vectors. The pairing $`<\overline{}|>`$ is an invariant so we require $`^{}<\overline{}|>=0`$. By Leibnitz rule this means that $`<^{}\overline{}|>+<\overline{}|^{}>=0`$. It is from here that we can extract how $`^{}\overline{}`$ acts. The quantity $`^{}\overline{}`$ will carry the dual space back into itself. In our notation we will use as a basis for the N-extended Super Virasoro Algebra
| $`^{}`$ | $`=(L_a,G_b{}_{}{}^{\mathrm{I}},T_c^{\mathrm{J}\mathrm{K}},_{\{\mathrm{I}_\mathrm{q}\}}U_{\{d_q\}}^{\{\mathrm{I}_\mathrm{q}\}}_{\{I_p\}}R_{\{e_q\}}^{\{\mathrm{J}_q\}};\alpha ),`$ |
| --- | --- |
| $``$ | $`=(L_z,G_\gamma {}_{}{}^{\mathrm{Q}},T_X^{\mathrm{R}\mathrm{S}},_{\{\mathrm{V}_l\}}U_{\{w_l\}}^{\{\mathrm{V}_l\}}_{\{T_m\}}R_{\{h_m\}}^{\{\mathrm{T}_m\}};\beta ),`$ |
| $`\overline{}`$ | $`=\left(\overline{L}{}_{\overline{z}}{}^{},\overline{G}{}_{\overline{\gamma }}{}^{}{}_{}{}^{\overline{\mathrm{Q}}},\overline{T}{}_{\overline{X}}{}^{\overline{\mathrm{R}}\overline{\mathrm{S}}},_{\{\overline{\mathrm{V}}_\overline{\mathrm{}}\}}\overline{U}{}_{\{\overline{W}{}_{\overline{\mathrm{}}}{}^{}\}}{}^{\{\overline{\mathrm{V}}{}_{\overline{\mathrm{}}}{}^{}\}},_{\{\overline{T}_m\}}\overline{R}{}_{\overline{h}_{\overline{m}}}{}^{\{\overline{\mathrm{T}}{}_{m}{}^{}\}};\overline{\beta }\right),`$ |
$`(24)`$
We will write a generic functions $`f=f_a\tau ^a`$ and $`z=z_b\tau ^{b+1/2}`$ and realize the Virasoro generators as
| $`L_\xi \xi ^{A+1}L_A`$ | $`=\left[(\xi ^{A+1}\tau ^{A+1})_\tau +\frac{1}{2}\xi ^{A+1}\tau ^A\zeta ^\mathrm{I}_\mathrm{I}\right]`$ |
| --- | --- |
| | $`=\left[\xi _\tau +\frac{1}{2}\xi \tau ^A\zeta ^\mathrm{I}_\mathrm{I}\right].`$ |
$`(25)`$
We will use the subscript of the generators and dual for a generic function of $`\tau `$. Each generator and dual element will have a specific function and how these functions transform under the action of specific generators is the aim of this paper.
Since the action will come from $`^{}\overline{}`$ we will denote a generic function as
| $`^{}`$ | $`=(L_\xi ,G_{\chi ^\mathrm{I}}{}_{}{}^{\mathrm{I}},T_{t^{\mathrm{J}\mathrm{K}}}^{\mathrm{J}\mathrm{K}},U_{\mu ^{\mathrm{I}_1\mathrm{}\mathrm{I}_q}}^{\mathrm{I}_1\mathrm{}\mathrm{I}_q},R_{r^{\mathrm{J}_1\mathrm{}\mathrm{J}_q}}^{\mathrm{J}_1\mathrm{}\mathrm{J}_q};\alpha ),`$ |
| --- | --- |
| $`\overline{}`$ | $`=(\overline{L}{}_{D}{}^{},\overline{G}_{\psi ^{\overline{\mathrm{Q}}}}^{\overline{\mathrm{Q}}},\overline{T}{}_{\tau ^{\mathrm{R}\mathrm{S}}}{}^{\overline{\mathrm{R}}\overline{\mathrm{S}}},\overline{U}{}_{\omega ^{\overline{\mathrm{V}}_1\mathrm{}\overline{\mathrm{V}}_q}}{}^{\overline{\mathrm{V}}_1\mathrm{}\overline{\mathrm{V}}_q},\overline{R}{}_{\rho ^{\overline{\mathrm{J}}_1\mathrm{}\overline{\mathrm{J}}_p}}{}^{\overline{\mathrm{J}}_1\mathrm{}\overline{\mathrm{J}}_p};\overline{\beta }).`$ |
$`(26)`$
The coadjoint action is quite tedious but we can organize the computation by examining the outcome of each of the commutation relations in the adjoint representation. Below is a table that symbolically will summarize our results. In the notation below $`LL`$ is just the commutator of two arbitrary Virasoro generators while $`L\overline{L}\overline{L}`$ is the coadjoint action from an application of the Virasoro generator on its dual $`\overline{L}`$ that maps back into the duals of the Virasoro generators. Multiple entries in the second column correspond to the different coadjoint actions that can be extracted from the commutator in the first column.
| Table 1 | |
| --- | --- |
| $`\mathrm{Commutator}`$ | $`\mathrm{Co}\mathrm{adjoint}\mathrm{Action}(\mathrm{s})`$ |
| $`LL`$ | $`L\overline{L}\overline{L}`$ |
| $`LG`$ | $`L\overline{G}\overline{G}`$ |
| $`LT`$ | $`L\overline{T}\overline{T}`$ |
| $`LU`$ | $`L\overline{U}\overline{U}`$ |
| $`LR`$ | $`L\overline{R}\overline{R},L\overline{U}\overline{R}`$ |
| $`GL`$ | $`G\overline{G}\overline{L}`$ |
| $`GG`$ | $`G\overline{L}\overline{G},G\overline{T}\overline{G},G\overline{U}\overline{G}`$ |
| $`GT`$ | $`G\overline{G}\overline{T},G\overline{U}\overline{T},`$ |
| $`GU`$ | $`G\overline{U}\overline{U},G\overline{R}\overline{U}`$ |
| $`GR`$ | $`G\overline{R}\overline{R},G\overline{U}\overline{R}`$ |
| $`TL`$ | $`T\overline{T}\overline{L}`$ |
| $`TG`$ | $`T\overline{G}\overline{G}`$ |
| $`TT`$ | $`T\overline{T}\overline{T}`$ |
| $`TU`$ | $`T\overline{U}\overline{U}`$ |
| $`TR`$ | $`T\overline{R}\overline{R},T\overline{U}\overline{R}`$ |
| $`UL`$ | $`U\overline{U}\overline{L}`$ |
| $`UG`$ | $`U\overline{U}\overline{G},U\overline{R}\overline{G}`$ |
| $`UT`$ | $`U\overline{U}\overline{T}`$ |
| $`UU`$ | $`U\overline{U}\overline{U}`$ |
| Table 2 | |
| --- | --- |
| $`\mathrm{Commutator}`$ | $`\mathrm{Co}\mathrm{adjoint}\mathrm{Action}(\mathrm{s})`$ |
| $`UR`$ | $`U\overline{R}\overline{R},U\overline{U}\overline{R}`$ |
| $`RL`$ | $`R\overline{R}\overline{L},R\overline{U}\overline{L}`$ |
| $`RG`$ | $`R\overline{R}\overline{G},R\overline{U}\overline{G}`$ |
| $`RT`$ | $`R\overline{R}\overline{T},R\overline{U}\overline{T}`$ |
| $`RU`$ | $`R\overline{R}\overline{U},R\overline{U}\overline{U}`$ |
| $`RR`$ | $`R\overline{R}\overline{R}`$ |
From these we can see that $`^{}\overline{}`$ will lead to changes in the coadjoint vectors as:
| $`\delta \overline{L}`$ | $`=L\overline{L}+G\overline{G}+T\overline{T}+U\overline{U}+R\overline{R}+R\overline{U},`$ |
| --- | --- |
| $`\delta \overline{G}`$ | $`=L\overline{G}+G\overline{L}+G\overline{T}+G\overline{U}+T\overline{G}+U\overline{U}`$ |
| | $`+U\overline{R}+R\overline{R}+R\overline{U},`$ |
| $`\delta \overline{T}`$ | $`=L\overline{T}+G\overline{G}+G\overline{U}+T\overline{T}+U\overline{U}+R\overline{R}`$ |
| | $`+R\overline{U},`$ |
| $`\delta \overline{U}`$ | $`=L\overline{U}+G\overline{U}+G\overline{R}+T\overline{U}+U\overline{U}+R\overline{R}`$ |
| | $`+R\overline{U},`$ |
| $`\delta \overline{R}`$ | $`=L\overline{R}+L\overline{U}+G\overline{R}+G\overline{U}+T\overline{R}+T\overline{U}`$ |
| | $`+U\overline{R}+U\overline{U}+R\overline{R},`$ |
| $`\delta \overline{\beta }`$ | $`=0.`$ |
$`(27)`$
(III.c) Explicit Variations
$`LL`$ Commutator:
Starting with the invariant pairing we have:
$$(L_a,\alpha )<(\overline{L}{}_{\overline{z}}{}^{},\overline{\beta })|(L{}_{z}{}^{},\beta )>=0.$$
$`(28)`$
Then by Leibnitz rule,
$$<(L_a,\alpha )(\overline{L}{}_{\overline{z}}{}^{},\overline{\beta })|(L{}_{z}{}^{},\beta )>+<(\overline{L}{}_{\overline{z}}{}^{},\overline{\beta })|(L_a,\alpha )(L{}_{z}{}^{},\beta )>=0.$$
$`(29)`$
Since we know the adjoint action we may write
$$<(L_a,\alpha )(\overline{L}{}_{\overline{z}}{}^{},\overline{\beta })|(L{}_{z}{}^{},\beta )>=<(\overline{L}{}_{\overline{z}}{}^{},\overline{\beta })|((az)L_{a+z},+\frac{1}{8}c(a^3a)\delta _{a+z,0})>$$
$`(30)`$
which implies that,
$$<(L_a,\alpha )(\overline{L}{}_{\overline{z}}{}^{},\overline{\beta })|(L{}_{z}{}^{},\beta )>=\{(az)\delta _{\overline{z},a+z}+\frac{1}{8}c(a^3a)\overline{\beta }\delta _{a+z,0}\},$$
$$(L_a,\alpha )(\overline{L}{}_{\overline{z}}{}^{},\overline{\beta })=\left((\mathrm{\hspace{0.17em}2}a\overline{z})\overline{L}{}_{\overline{z}a}{}^{}+\frac{1}{8}c\overline{\beta }(a^3a)\overline{L}{}_{a}{}^{},\mathrm{\hspace{0.17em}0}\right),$$
$`(31)`$
where $`z=\overline{z}a`$, so that
$$(L_a,\alpha )(\overline{L}{}_{\overline{z}}{}^{},\overline{\beta })=\left((\mathrm{\hspace{0.17em}2}a\overline{z})\overline{L}{}_{\overline{z}a}{}^{}+\frac{1}{8}c\overline{\beta }(a^3a)\overline{L}{}_{a}{}^{},\mathrm{\hspace{0.17em}0}\right).$$
$`(32)`$
Rewriting in terms of functions instead of modes we have for functions $`\xi `$ and $`D`$,
$$(L_\xi ,\alpha )(\overline{L}{}_{D}{}^{},\overline{\beta })=(\overline{L}_{\stackrel{~}{D}},0),$$
$`(33)`$
where $`\stackrel{~}{D}=(\mathrm{\hspace{0.17em}2}\xi ^{}D+\xi D^{}+\frac{1}{8}c\overline{\beta }\xi ^{\prime \prime \prime })`$. This shows the usual transformation of a quadratic differential, $`D`$, with respect to the vector field $`\xi `$. Up to the inhomogeneous term $`D`$ transforms as a rank two tensor. It is the inhomogeneous term that violates tensorality. From the adjoint action one sees that $`\xi `$ transforms as a rank one contravariant tensor in one dimension making it easy to identify with $`\xi ^\alpha `$ from a Lie derivative. In the same way $`D`$ can be thought of as a two index object $`D_{\alpha \beta }`$. This suggests that a spin two type object is present in the spectrum. Throughout we will treat the action of $`\xi `$ as a one dimensional Lie derivative (up to extensions) in order to understand the type of fields that are present in the dual.
$`GL`$ Commutator:
In the same way we examine the action of the $`G_b^\mathrm{I}`$ on the pairing. Since the pairing is invariant we have
$$G_b{}_{}{}^{\mathrm{I}}<\overline{G}{}_{\overline{y}}{}^{}{}_{}{}^{\overline{\mathrm{Q}}}|L{}_{z}{}^{}>=0,$$
$`(34)`$
by Leibnitz
$$<G_b{}_{}{}^{\mathrm{I}}\overline{G}{}_{\overline{y}}{}^{}{}_{}{}^{\overline{\mathrm{Q}}}|L{}_{z}{}^{}>+<\overline{G}{}_{\overline{y}}{}^{}{}_{}{}^{\overline{\mathrm{Q}}}|G_b{}_{}{}^{\mathrm{I}}L{}_{z}{}^{}>=0,$$
$`(35)`$
which implies that
| $`<G_b{}_{}{}^{\mathrm{I}}\overline{G}{}_{\overline{y}}{}^{}{}_{}{}^{\overline{\mathrm{Q}}}|L{}_{z}{}^{}>`$ | $`=<\overline{G}{}_{\overline{y}}{}^{}{}_{}{}^{\overline{\mathrm{Q}}}|(\frac{1}{2}zb)G_b{}_{}{}^{\mathrm{I}}>=(\frac{1}{2}zb)\delta ^{\mathrm{I}\overline{\mathrm{Q}}}\delta _{\overline{y},b+z},`$ |
| --- | --- |
$`(36)`$
where $`z=\overline{y}b`$. This yields
$$G_b{}_{}{}^{\mathrm{I}}\overline{G}{}_{\overline{y}}{}^{}{}_{}{}^{\overline{\mathrm{Q}}}=(\frac{1}{2}\overline{y}\frac{3}{2}b)\overline{L}_{\overline{y}b}\delta ^{\overline{Q}I}.$$
$`(37)`$
Rewriting in terms of functions we have we have that
$$G_{\chi ^I}^I\overline{G}_{\mathrm{\Psi }^{\overline{Q}}}^{\overline{Q}}=\overline{L}_f,\mathrm{where}f=(\frac{1}{2}(\mathrm{\Psi }^{\overline{Q}})^{}\chi ^I\frac{3}{2}(\chi ^I)^{}\mathrm{\Psi }^{\overline{Q}})\delta ^{\overline{Q}I}.$$
$`(38)`$
All $``$ Commutators:
| $`L_\xi (\overline{L}_D,\overline{\beta })`$ | $`=\overline{L}_{\stackrel{~}{D}},\stackrel{~}{D}=2\xi ^{}D\xi D^{}\frac{c\overline{\beta }}{8}\xi ^{\prime \prime \prime },`$ |
| --- | --- |
| $`L_\xi \overline{G}_{\mathrm{\Psi }^{\overline{Q}}}^{\overline{Q}}`$ | $`=\overline{G}_{\stackrel{~}{\mathrm{\Psi }^{\overline{Q}}}}^{\overline{Q}},\stackrel{~}{\mathrm{\Psi }^{\overline{Q}}}=(\frac{3}{2}\xi ^{}\psi ^{\overline{Q}}+\xi (\psi ^{\overline{Q}})^{}),`$ |
| $`L_\xi \overline{T}_{\tau ^{\overline{\mathrm{R}}\overline{\mathrm{S}}}}^{\overline{\mathrm{R}}\overline{\mathrm{S}}}`$ | $`=\overline{T}_{\stackrel{~}{\tau }^{\overline{\mathrm{R}}\overline{\mathrm{S}}}}^{\overline{\mathrm{R}}\overline{\mathrm{S}}},\stackrel{~}{\tau }^{\overline{\mathrm{R}}\overline{\mathrm{S}}}=\xi ^{}\tau ^{\overline{\mathrm{R}}\overline{\mathrm{S}}}\xi (\tau ^{\overline{\mathrm{R}}\overline{\mathrm{S}}})^{},`$ |
| $`L_\xi \overline{U}_{\omega ^{\overline{\mathrm{V}}_1\mathrm{}\overline{\mathrm{V}}_\mathrm{n}}}^{\overline{\mathrm{V}}_1\mathrm{}\overline{\mathrm{V}}_\mathrm{n}}`$ | $`=\overline{U}_{\stackrel{~}{\omega }^{\overline{\mathrm{V}}_1\mathrm{}\overline{\mathrm{V}}_\mathrm{n}}}^{\overline{\mathrm{V}}_1\mathrm{}\overline{\mathrm{V}}_\mathrm{n}}+\frac{i}{2}(i)^{[{\scriptscriptstyle \frac{n2}{2}}][{\scriptscriptstyle \frac{n}{2}}]}\overline{R}_{\xi ^{\prime \prime }\omega ^{\overline{\mathrm{V}}_1\mathrm{}\overline{\mathrm{V}}_\mathrm{n}}}^{[\overline{\mathrm{V}}_1\mathrm{}\overline{\mathrm{V}}_{\mathrm{n}2}}\delta ^{\mathrm{V}_{\mathrm{n}1}]\mathrm{V}_\mathrm{n}},`$ |
| | $`\stackrel{~}{\omega }^{\overline{\mathrm{V}}_1\mathrm{}\overline{\mathrm{V}}_\mathrm{n}}=(\frac{n}{2}2)\xi ^{}\omega ^{\overline{\mathrm{V}}_1\mathrm{}\overline{\mathrm{V}}_\mathrm{n}}\xi (\omega ^{\overline{\mathrm{V}}_1\mathrm{}\overline{\mathrm{V}}_\mathrm{n}})^{},`$ |
| $`L_\xi \overline{R}_{\rho ^{\overline{\mathrm{T}}_1\mathrm{}\overline{\mathrm{T}}_\mathrm{m}}}^{\overline{\mathrm{T}}_1\mathrm{}\overline{\mathrm{T}}_\mathrm{m}}`$ | $`=\overline{R}_{\stackrel{~}{\rho }^{\overline{\mathrm{T}}_1\mathrm{}\overline{\mathrm{T}}_\mathrm{m}}}^{\overline{\mathrm{T}}_1\mathrm{}\overline{\mathrm{T}}_\mathrm{m}},`$ |
| | $`\stackrel{~}{\rho }^{\overline{\mathrm{T}}_1\mathrm{}\overline{\mathrm{T}}_\mathrm{m}}=(\frac{\mathrm{m}}{2}2)\xi ^{}\rho ^{\overline{\mathrm{T}}_1\mathrm{}\overline{\mathrm{T}}_\mathrm{m}}\xi (\rho ^{\overline{\mathrm{T}}_1\mathrm{}\overline{\mathrm{T}}_\mathrm{m}})^{},`$ |
| $`G_{\chi ^\mathrm{I}}^\mathrm{I}\overline{G}_{\psi ^{\overline{\mathrm{Q}}}}^{\overline{\mathrm{Q}}}`$ | $`=\delta ^{\mathrm{I}\overline{\mathrm{Q}}}\overline{L}_{\stackrel{~}{\xi }}+4\overline{T}_{(\chi ^\mathrm{I}\psi ^{\overline{\mathrm{Q}}})}^{\mathrm{I}\overline{\mathrm{Q}}},\stackrel{~}{\xi }=\frac{1}{2}(\psi ^{\overline{Q}})^{}\chi ^I\frac{3}{2}(\chi ^I)^{}\psi ^{\overline{Q}},`$ |
| $`G_{\chi ^\mathrm{I}}^\mathrm{I}(\overline{L}_D,\overline{\beta })`$ | $`=4i\overline{G}_{(\chi ^\mathrm{I}D\overline{\beta }c(\chi ^\mathrm{I})^{\prime \prime })}^\mathrm{I},`$ |
| $`G_{\chi ^\mathrm{I}}^\mathrm{I}\overline{T}_{\tau ^{\overline{\mathrm{R}}\mathrm{S}}}^{\overline{\mathrm{R}}\mathrm{S}}`$ | $`=\frac{i}{2}(\overline{G}{}_{\chi ^\mathrm{S}}{}^{\mathrm{S}}\delta _{}^{\overline{\mathrm{R}}\mathrm{I}}\overline{G}{}_{\chi ^{\overline{\mathrm{R}}}}{}^{\overline{\mathrm{R}}}\delta _{}^{\mathrm{I}\overline{\mathrm{S}}}),\chi ^{\overline{\mathrm{R}}}=\chi ^\mathrm{S}=2(\chi ^\mathrm{I})^{}\tau ^{\overline{\mathrm{R}}\mathrm{S}}+\chi ^\mathrm{I}(\tau ^{\overline{\mathrm{R}}\mathrm{S}})^{},`$ |
| $`G_{\chi ^\mathrm{I}}^\mathrm{I}\overline{R}_{\rho ^{\overline{\mathrm{T}}_1\mathrm{}\overline{\mathrm{T}}_\mathrm{m}}}^{\overline{\mathrm{T}}_1\mathrm{}\overline{\mathrm{T}}_\mathrm{m}}`$ | $`=2i(i)^{m+1}(i)^{[{\scriptscriptstyle \frac{m+2}{2}}][{\scriptscriptstyle \frac{m}{2}}]}\overline{U}_{(\chi ^\mathrm{I}\rho ^{\overline{\mathrm{T}}_1\mathrm{}\overline{\mathrm{T}}_\mathrm{m}})}^{[\overline{\mathrm{T}}_1\mathrm{}\overline{\mathrm{T}}_\mathrm{m}]}`$ |
| | $`2i^{[\frac{m1}{2}][{\scriptscriptstyle \frac{m2}{2}}]}\delta ^{I[\overline{T}_1}\overline{R}_{((\chi ^\mathrm{I})^{}\rho ^{\overline{\mathrm{T}}_1\mathrm{}\overline{\mathrm{T}}_\mathrm{m}}(\chi ^\mathrm{I})(\rho ^{\overline{\mathrm{T}}_1\mathrm{}\overline{\mathrm{T}}_\mathrm{m}})^{})}^{\overline{\mathrm{T}}_2\mathrm{}\overline{\mathrm{T}}_\mathrm{m}]}`$ |
| | $`(i)(i)^{[{\scriptscriptstyle \frac{m+1}{2}}][{\scriptscriptstyle \frac{m}{2}}]}{\displaystyle \underset{r=1}{\overset{m+1}{}}}(1)^{r1}\overline{R}_{(\chi ^\mathrm{I}\rho ^{\overline{\mathrm{T}}_1\mathrm{}\overline{\mathrm{T}}_\mathrm{m}})}^{\overline{\mathrm{T}}_1\mathrm{}\overline{\mathrm{T}}_{\mathrm{r}1}\mathrm{I}\overline{\mathrm{T}}_{\mathrm{r}+1}\mathrm{}\overline{\mathrm{T}}_\mathrm{m}},`$ |
| $`T_{t^{\mathrm{J}\mathrm{K}}}^{\mathrm{J}\mathrm{K}}\overline{G}_{\psi ^{\overline{\mathrm{Q}}}}^{\overline{\mathrm{Q}}}`$ | $`=2(\overline{G}_{(t^{\mathrm{J}\mathrm{K}}\psi ^{\overline{\mathrm{Q}}})}^\mathrm{K}\delta ^{\overline{\mathrm{Q}}\mathrm{J}}\overline{G}_{(t^{\mathrm{J}\mathrm{K}}\psi ^{\overline{\mathrm{Q}}})}^\mathrm{J}\delta ^{\overline{\mathrm{Q}}\mathrm{K}}),`$ |
| $`G_{\chi ^\mathrm{I}}^\mathrm{I}\overline{U}_{\omega ^{\overline{\mathrm{V}}_1\mathrm{}\overline{\mathrm{V}}_\mathrm{n}}}^{\overline{\mathrm{V}}_1\mathrm{}\overline{\mathrm{V}}_\mathrm{n}}`$ | $`=2i^{[\frac{n1}{2}][{\scriptscriptstyle \frac{n}{2}}]}\delta ^{I[\overline{V}_1}\overline{U}_{((n4)(\chi ^\mathrm{I})^{}\omega ^{\overline{\mathrm{V}}_1\mathrm{}\overline{\mathrm{V}}_\mathrm{n}}(\chi ^\mathrm{I})(\omega ^{\overline{\mathrm{V}}_1\mathrm{}\overline{\mathrm{V}}_\mathrm{n}})^{})}^{\overline{\mathrm{V}}_2\mathrm{}\overline{\mathrm{V}}_\mathrm{n}]}`$ |
| | $`+2(1)^{n1}(i)^{[{\scriptscriptstyle \frac{n1}{2}}][{\scriptscriptstyle \frac{n}{2}}]}\delta ^{\mathrm{I}[\overline{\mathrm{V}}_\mathrm{n}}\overline{U}_{((\chi ^\mathrm{I})^{}\omega ^{\overline{\mathrm{V}}_1\mathrm{}\overline{\mathrm{V}}_\mathrm{n}})}^{\overline{V}_1\mathrm{}\overline{V}_n]\mathrm{K}}_\mathrm{K}`$ |
| | $`+(i)(i)^{[{\scriptscriptstyle \frac{n1}{2}}][{\scriptscriptstyle \frac{n}{2}}]}{\displaystyle \underset{r=1}{\overset{n}{}}}\overline{U}_{(\chi ^\mathrm{I}\omega ^{\overline{\mathrm{V}}_1\mathrm{}\overline{\mathrm{V}}_\mathrm{n}})}^{[\overline{\mathrm{V}}_1\mathrm{}\overline{\mathrm{V}}_{\mathrm{r}1}\overline{\mathrm{V}}_{\mathrm{r}+1}\mathrm{}\overline{\mathrm{V}}_\mathrm{n}}\delta ^{\overline{\mathrm{V}}_\mathrm{r}]\mathrm{I}}`$ |
| | $`+\overline{G}_{(4i(\chi ^\mathrm{I})^{}(\omega ^{\overline{\mathrm{V}}_1\mathrm{}\overline{\mathrm{V}}_\mathrm{n}})^{}2i(\chi ^\mathrm{I})(\omega ^{\overline{\mathrm{V}}_1\mathrm{}\overline{\mathrm{V}}_\mathrm{n}})^{\prime \prime })}^{[\overline{\mathrm{V}}_2}\delta ^{\overline{\mathrm{V}_3}\overline{\mathrm{V}_4}}\delta ^{\overline{\mathrm{V}}_1]\mathrm{I}}\delta ^{n4}`$ |
| | $`2i(1)^n(i)^{[{\scriptscriptstyle \frac{n}{2}}][{\scriptscriptstyle \frac{n1}{2}}]}\delta ^{\mathrm{I}[\overline{\mathrm{V}}_\mathrm{n}}\overline{R}_{(\chi ^\mathrm{I}\omega ^{\overline{\mathrm{V}}_1\mathrm{}\overline{\mathrm{V}}_\mathrm{n}})}^{\overline{\mathrm{V}}_1\mathrm{}\overline{\mathrm{V}}_{\mathrm{n}1}]},`$ |
| $`T_{t^{\mathrm{J}\mathrm{K}}}^{\mathrm{J}\mathrm{K}}\overline{U}_{\omega ^{\overline{\mathrm{V}}_1\mathrm{}\overline{\mathrm{V}}_\mathrm{n}}}^{\overline{\mathrm{V}}_1\mathrm{}\overline{\mathrm{V}}_\mathrm{n}}`$ | $`={\displaystyle \underset{r=1}{\overset{n1}{}}}(1)^{n+1}(\delta ^{\mathrm{J}[\overline{\mathrm{V}}_1}\overline{U}_{(t^{\mathrm{J}\mathrm{K}}\omega ^{\overline{\mathrm{V}}_1\mathrm{}\overline{\mathrm{V}}_\mathrm{n}})}^{\overline{\mathrm{V}}_2\mathrm{}\overline{\mathrm{V}}_{\mathrm{r}1}|\mathrm{K}|\overline{\mathrm{V}}_{\mathrm{r}+1}\mathrm{}]\overline{\mathrm{V}}_\mathrm{n}}`$ |
| | $`\delta ^{\mathrm{K}[\overline{\mathrm{V}}_1}\overline{U}_{(t^{\mathrm{J}\mathrm{K}}\omega ^{\overline{\mathrm{V}}_1\mathrm{}\overline{\mathrm{V}}_\mathrm{n}})}^{\overline{\mathrm{V}}_2\mathrm{}\overline{\mathrm{V}}_{\mathrm{r}1}|\mathrm{J}|\overline{\mathrm{V}}_{\mathrm{r}+1}\mathrm{}]\overline{\mathrm{V}}_\mathrm{n}})`$ |
| | $`+\overline{U}_{(t^{\mathrm{J}\mathrm{K}}\omega ^{\overline{\mathrm{V}}_1\mathrm{}\overline{\mathrm{V}}_\mathrm{n}})}^{[\overline{\mathrm{V}}_1\mathrm{}\overline{\mathrm{V}}_{\mathrm{n}1}]\mathrm{J}}\delta ^{\overline{\mathrm{V}}_\mathrm{n}\mathrm{K}}\overline{U}_{(t^{\mathrm{J}\mathrm{K}}\omega ^{\overline{\mathrm{V}}_1\mathrm{}\overline{\mathrm{V}}_\mathrm{n}})}^{[\overline{\mathrm{V}}_1\mathrm{}\overline{\mathrm{V}}_{\mathrm{n}1}]\mathrm{K}}\delta ^{\overline{\mathrm{V}}_\mathrm{n}\mathrm{J}}`$ |
| | $`i(1)^{n2}(\delta ^{\mathrm{K}\overline{\mathrm{V}}_\mathrm{n}}\delta ^{J[\overline{V}_1}\delta ^{\mathrm{J}\overline{\mathrm{V}}_\mathrm{n}}\delta ^{\mathrm{K}[\overline{\mathrm{V}}_1})\overline{R}_{((t^{\mathrm{J}\mathrm{K}})^{}\omega ^{\overline{\mathrm{V}}_1\mathrm{}\overline{\mathrm{V}}_\mathrm{n}})}^{\overline{\mathrm{V}}_2\mathrm{}\overline{\mathrm{V}}_{\mathrm{n}1}]},`$ |
| $`R_{r^{\{J_p\}}}^{\mathrm{J}_1\mathrm{}\mathrm{J}_\mathrm{p}}\overline{U}_{\omega ^{\{V_m\}}}^{\overline{\mathrm{V}}_1\mathrm{}\overline{\mathrm{V}}_\mathrm{m}}`$ | $`=\frac{1}{2}i(i)^{\{[{\scriptscriptstyle \frac{p}{2}}][{\scriptscriptstyle \frac{p+2}{2}}]\}}\delta _{[\overline{\mathrm{J}}_1\mathrm{}\overline{\mathrm{J}}_\mathrm{p}]}^{[\overline{\mathrm{V}}_1\mathrm{}\overline{\mathrm{V}}_{\mathrm{m}2}}\delta ^{\overline{\mathrm{V}}_{\mathrm{m}1}],\overline{\mathrm{V}}_\mathrm{m}}\delta ^{\mathrm{m},\mathrm{p}+2}\overline{L}_{(r\omega )^{\prime \prime }}`$ |
| --- | --- |
| | $`+i^{\{[{\scriptscriptstyle \frac{p}{2}}][{\scriptscriptstyle \frac{p+1}{2}}]\}}\delta ^{\mathrm{p}+1,\mathrm{m}}\overline{G}_{(r\omega )^{}}^{\overline{\mathrm{V}}_\mathrm{m}}\delta _{[\mathrm{J}_1\mathrm{}\mathrm{J}_\mathrm{p}]}^{[\overline{\mathrm{V}}_1\mathrm{}\overline{\mathrm{V}}_{\mathrm{m}1}]}`$ |
| | $`+2i(1)^{p+1}\delta ^{\mathrm{p}+3,\mathrm{m}}(i)^{[{\scriptscriptstyle \frac{p}{2}}][{\scriptscriptstyle \frac{p+2}{2}}]}\overline{G}_{(r\omega )^{\prime \prime }}^{[\overline{V}_1}\delta _{[\mathrm{J}_1\mathrm{}\mathrm{J}_\mathrm{p}]}^{\overline{\mathrm{V}}_2\mathrm{}\overline{\mathrm{V}}_{\mathrm{m}2}}\delta ^{\overline{\mathrm{V}}_{\mathrm{m}1},]\overline{\mathrm{V}}_\mathrm{m}}`$ |
| | $`+i(1)^p\overline{T}_{(r\omega )^{}}^{\overline{\mathrm{V}}_1[\overline{\mathrm{V}}_\mathrm{m}}\delta _{[\mathrm{J}_1\mathrm{}\mathrm{J}_\mathrm{p}]}^{\overline{\mathrm{V}}_2\mathrm{}\overline{\mathrm{V}}_{\mathrm{m}1}]}\delta ^{p+2,m}`$ |
| | $`+(1)^{\mathrm{pm}}(i)^{\{[{\scriptscriptstyle \frac{m}{2}}]+[{\scriptscriptstyle \frac{p}{2}}][{\scriptscriptstyle \frac{m+p2}{2}}]\}}\overline{U}_{(r\omega )^{}}^{\overline{\mathrm{V}}_\mathrm{p}+1\mathrm{}\overline{\mathrm{V}}_{\mathrm{m}\mathrm{p}}}\delta _{[\mathrm{J}_1\mathrm{}\mathrm{J}_\mathrm{p}]}^{\overline{\mathrm{V}}_1\mathrm{}\overline{\mathrm{V}}_\mathrm{p}},`$ |
| $`U_{\mu ^{\{\mathrm{I}_\mathrm{q}\}}}^{\mathrm{I}_1\mathrm{}\mathrm{I}_\mathrm{q}}\overline{U}_{\omega ^{\{\mathrm{V}_\mathrm{m}\}}}^{\overline{\mathrm{V}}_1\mathrm{}\overline{\mathrm{V}}_\mathrm{m}}`$ | $`=\delta ^{\mathrm{m}\mathrm{q}}\delta _{[\overline{\mathrm{V}}_1\mathrm{}\overline{\mathrm{V}}_\mathrm{q}]}^{[\mathrm{I}_1\mathrm{}\mathrm{I}_\mathrm{q}]}\overline{L}_{(({\scriptscriptstyle \frac{4q}{2}})\mu \omega ({\scriptscriptstyle \frac{q2}{2}})\mu \omega ^{})}2\overline{T}_{(\omega \mu )}^{\overline{\mathrm{V}}_\mathrm{m}\mathrm{J}_\mathrm{q}}\delta _{[\overline{\mathrm{V}}_1\mathrm{}\overline{\mathrm{V}}_\mathrm{q}]}^{[\mathrm{I}_1\mathrm{}\mathrm{I}_\mathrm{q}]}\delta _\mathrm{m}^\mathrm{q}`$ |
| | $`2i^{([{\scriptscriptstyle \frac{\overline{V}_1}{2}}][{\scriptscriptstyle \frac{\overline{V}_1+1}{2}}])}\delta ^{\mathrm{m},(\mathrm{q}+1)}\overline{G}_{((q2)\omega ^{}\mu +(3q)\omega \mu ^{})}^{[\overline{\mathrm{V}}_1}\delta _{[\mathrm{I}_1\mathrm{}\mathrm{I}_\mathrm{q}]}^{\overline{\mathrm{V}}_2\mathrm{}\overline{\mathrm{V}}_\mathrm{m}]}`$ |
| | $`+2(1)^q(i)^{[{\scriptscriptstyle \frac{q}{2}}][{\scriptscriptstyle \frac{q+1}{2}}]}\overline{G}_{(\omega \mu )^{}}^{\mathrm{I}_\mathrm{q}}\delta ^{\mathrm{m},\mathrm{q}+1}\delta _{[\mathrm{I}_1\mathrm{}\mathrm{I}_{\mathrm{q}1}]}^{[\overline{\mathrm{V}}_1\mathrm{}\overline{\mathrm{V}}_{\mathrm{m}2}}\delta ^{\overline{\mathrm{V}}_{\mathrm{m}1}],\overline{\mathrm{V}}_\mathrm{m}}`$ |
| | $`i(i)^{[{\scriptscriptstyle \frac{q}{2}}][{\scriptscriptstyle \frac{q1}{2}}]}{\displaystyle \underset{r=1}{\overset{q1}{}}}(1)^{r1}\delta _{\mathrm{q}1}^\mathrm{m}\overline{G}_{(\omega \mu )}^{[\mathrm{I}_\mathrm{r}}\delta _{\overline{[}\mathrm{V}_1}^{\mathrm{I}_1}\mathrm{}\delta _{\overline{\mathrm{V}}_{\mathrm{r}1}}^{\mathrm{I}_{\mathrm{r}1}}\delta _{\overline{\mathrm{V}}_\mathrm{r}}^{\mathrm{I}_{\mathrm{r}+1}}\mathrm{}\delta _{\overline{\mathrm{V}}_\mathrm{m}]}^{\mathrm{I}_\mathrm{q}]}`$ |
| | $`+{\displaystyle \underset{r=1}{\overset{q1}{}}}2(1)^{r+1}(\overline{T}_{(\omega \mu )}^{\mathrm{J}_\mathrm{r}[\overline{\mathrm{V}}_1}\delta _{[\mathrm{I}_1}^{\overline{V}_2}\mathrm{}\delta _{\mathrm{I}_{\mathrm{r}1}}^{\overline{V}_r}\delta _{\mathrm{I}_{\mathrm{r}+1}}^{\overline{V}_{r+1}}\mathrm{}\delta _{]\mathrm{I}_\mathrm{q}}^{]\overline{V}_m}\delta ^{\mathrm{q},\mathrm{m}})`$ |
| | $`+i(i)^{\{[{\scriptscriptstyle \frac{q}{2}}]+[{\scriptscriptstyle \frac{mq}{2}}+2][{\scriptscriptstyle \frac{m+2}{2}}]\}}`$ |
| | $`\times \{{\displaystyle \underset{r=1}{\overset{q}{}}}(1)^{r1}\delta _{[\mathrm{I}_1\mathrm{}\mathrm{I}_{\mathrm{q}1}]}^{[\overline{\mathrm{V}}_1\mathrm{}\overline{\mathrm{V}}_{\mathrm{q}1}}\overline{U}_{\omega \mu }^{\overline{\mathrm{V}}_\mathrm{q}\mathrm{}\overline{\mathrm{V}}_{\mathrm{q}+\mathrm{r}1}\mathrm{I}_\mathrm{q}\overline{\mathrm{V}}_{\mathrm{q}+\mathrm{r}}\mathrm{}\overline{]}\mathrm{V}_\mathrm{m}}`$ |
| | $`(1)^{q(mq+2)}{\displaystyle \underset{r=1}{\overset{q}{}}}(1)^{r1}\overline{U}_{\omega \mu }^{\overline{\mathrm{V}}_1\mathrm{}\overline{\mathrm{V}}_{\mathrm{m}\mathrm{q}+1}[\mathrm{I}_\mathrm{r}}\delta _{\overline{\mathrm{V}}_{\mathrm{m}\mathrm{q}+2}\mathrm{}\overline{\mathrm{V}}_{\mathrm{m}\mathrm{q}+2+\mathrm{r}}\mathrm{}\overline{\mathrm{V}}_\mathrm{m}}^{\mathrm{I}_1\mathrm{}\mathrm{I}_{\mathrm{r}1}\mathrm{I}_{\mathrm{r}+1}\mathrm{}\mathrm{I}_\mathrm{q}]}\}`$ |
| | $`(i)^{\{[{\scriptscriptstyle \frac{q}{2}}]+[{\scriptscriptstyle \frac{mq}{2}}][{\scriptscriptstyle \frac{q+m4}{2}}]\}}\overline{R}_{(\omega \mu )^{}}^{[\overline{\mathrm{V}}_1\mathrm{}\overline{\mathrm{V}}_{\mathrm{m}\mathrm{q}}}\delta _{[\mathrm{I}_1\mathrm{}]\mathrm{I}_\mathrm{q}}^{\overline{\mathrm{V}}_{\mathrm{m}\mathrm{q}+1}\mathrm{}]\overline{\mathrm{V}}_\mathrm{m}},`$ |
| $`R_{r^{\{J_p\}}}^{\mathrm{J}_1\mathrm{}\mathrm{J}_\mathrm{p}}\overline{R}_{\rho ^{\{T_m\}}}^{\overline{\mathrm{T}}_1\mathrm{}\overline{\mathrm{T}}_\mathrm{m}}`$ | $`=\delta _{[\mathrm{J}_1\mathrm{}\mathrm{J}_\mathrm{p}]}^{[\overline{\mathrm{T}}_1\mathrm{}\overline{\mathrm{T}}_\mathrm{m}]}\delta ^{\mathrm{p}\mathrm{m}}\overline{L}_{(({\scriptscriptstyle \frac{p}{2}}2)r^{}\rho ({\scriptscriptstyle \frac{p}{2}}1)r\rho ^{})}`$ |
| | $`+(1)^p\{2(i)\overline{G}_{((2p)r^{}\rho (p1)r\rho ^{})}^{[\overline{\mathrm{T}}_1}\delta _{[\mathrm{J}_1\mathrm{}\mathrm{J}_\mathrm{p}]}^{\overline{\mathrm{T}}_2\mathrm{}\overline{\mathrm{T}}_\mathrm{m}]}\delta _{p+1}^m`$ |
| | $`+(i)(i)^{[{\scriptscriptstyle \frac{p}{2}}][{\scriptscriptstyle \frac{p1}{2}}]}{\displaystyle \underset{r=1}{\overset{p}{}}}(1)^r\overline{G}_{(r\rho )}^{\mathrm{J}_\mathrm{r}}\delta _{\mathrm{p}1}^\mathrm{m}\delta _{[\overline{\mathrm{T}}_1\mathrm{}\overline{\mathrm{T}}_{\mathrm{r}1}\overline{\mathrm{T}}_\mathrm{r}\mathrm{}\overline{\mathrm{T}}_\mathrm{m}]}^{[\mathrm{J}_1\mathrm{}\mathrm{J}_{\mathrm{r}1}\mathrm{J}_{\mathrm{r}+1}\mathrm{}\mathrm{J}_\mathrm{p}]}\}`$ |
| | $`+{\displaystyle \underset{r=1}{\overset{p}{}}}(1)^{r+1}2\overline{T}_{(r\rho )}^{[\overline{\mathrm{T}}_1|\mathrm{J}_\mathrm{r}|}\delta _{[\mathrm{J}_1\mathrm{}\mathrm{J}_{\mathrm{r}1}\mathrm{J}_{\mathrm{r}+1}\mathrm{}\mathrm{J}_\mathrm{p}]}^{\overline{\mathrm{T}}_2\mathrm{}\overline{\mathrm{T}}_\mathrm{m}]}\delta _m^p`$ |
| | $`{\displaystyle \underset{r=1}{\overset{p}{}}}(1)^{r1}\overline{U}_{(r\rho )}^{[\overline{\mathrm{T}}_1\mathrm{}\overline{\mathrm{T}}_{\mathrm{m}\mathrm{p}+1}|\mathrm{J}_\mathrm{r}|}\delta _{\mathrm{J}_1\mathrm{}\mathrm{J}_{\mathrm{r}1}\mathrm{J}_{\mathrm{r}+1}\mathrm{}\mathrm{J}_\mathrm{p}}^{\overline{\mathrm{T}}_{\mathrm{m}\mathrm{p}+2}\mathrm{}\overline{\mathrm{T}}_\mathrm{m}]}`$ |
| | $`+i^{\{[{\scriptscriptstyle \frac{p}{2}}]+[{\scriptscriptstyle \frac{mp}{2}}][{\scriptscriptstyle \frac{m}{2}}]\}}\overline{R}_{(2r^{}\rho +r\rho ^{})}^{[\overline{\mathrm{T}}_{\mathrm{p}+1}\mathrm{}\overline{\mathrm{T}}_\mathrm{m}}\delta _{[\mathrm{J}_1\mathrm{}\mathrm{J}_\mathrm{p}]}^{\overline{\mathrm{T}}_1\mathrm{}\overline{\mathrm{T}}_\mathrm{p}]},`$ |
| $`T_{t^{\mathrm{J}\mathrm{K}}}^{\mathrm{J}\mathrm{K}}\overline{R}_{\rho ^{\overline{\mathrm{T}}_1\mathrm{}\overline{\mathrm{T}}_\mathrm{m}}}^{\overline{\mathrm{T}}_1\mathrm{}\overline{\mathrm{T}}_\mathrm{m}}`$ | $`={\displaystyle \underset{r=1}{\overset{m}{}}}(1)^{r+1}(\delta ^{[\overline{\mathrm{T}}_1|\mathrm{J}|}\overline{R}_{(t^{\mathrm{J}\mathrm{K}}\rho ^{\overline{\mathrm{V}}_1\mathrm{}\overline{\mathrm{V}}_\mathrm{n}})}^{\overline{\mathrm{T}}_2\mathrm{}\overline{\mathrm{T}}_{\mathrm{r}1}|\mathrm{K}|\overline{\mathrm{T}}_{\mathrm{r}+1}\mathrm{}\overline{\mathrm{T}}_\mathrm{m}]}`$ |
| --- | --- |
| | $`\delta ^{[\overline{\mathrm{T}}_1|\mathrm{K}|}\overline{R}_{(t^{\mathrm{J}\mathrm{K}}\rho ^{\overline{\mathrm{V}}_1\mathrm{}\overline{\mathrm{V}}_\mathrm{n}})}^{\overline{\mathrm{T}}_2\mathrm{}\overline{\mathrm{T}}_{\mathrm{r}1}|\mathrm{J}|\overline{\mathrm{T}}_{\mathrm{r}+1}\mathrm{}\overline{\mathrm{T}}_\mathrm{m}]}),`$ |
| $`U_{\mu ^{\{\mathrm{I}_\mathrm{q}\}}}^{\mathrm{I}_1\mathrm{}\mathrm{I}_\mathrm{q}}\overline{R}_{\rho ^{\{\mathrm{T}_\mathrm{m}\}}}^{\overline{\mathrm{T}}_1\mathrm{}\overline{\mathrm{T}}_\mathrm{m}}`$ | $`=i(1)^{q(mq+2)}(i)^{\{[{\scriptscriptstyle \frac{mq}{2}}+2]+[{\scriptscriptstyle \frac{q}{2}}][{\scriptscriptstyle \frac{m}{2}}]\}}`$ |
| | $`\times {\displaystyle \underset{r=1}{\overset{mq+2}{}}}(1)^{r1}\delta _{[\overline{\mathrm{T}}_1\mathrm{}\overline{\mathrm{T}}_{\mathrm{q}1}]}^{[\mathrm{I}_1\mathrm{}\mathrm{I}_{\mathrm{q}1}]}\overline{R}_{\rho \mu }^{\overline{\mathrm{T}}_1\mathrm{}\overline{\mathrm{T}}_{\mathrm{q}+\mathrm{r}1}\mathrm{I}_\mathrm{q}\overline{\mathrm{T}}_{\mathrm{q}+\mathrm{r}+1}\mathrm{}\overline{\mathrm{T}}_\mathrm{m}},`$ |
| $`T_{t^{\mathrm{J}\mathrm{K}}}^{\mathrm{J}\mathrm{K}}(\overline{T}_{\tau ^{\overline{\mathrm{R}}\overline{\mathrm{S}}}}^{\overline{\mathrm{R}}\overline{\mathrm{S}}},\overline{\beta })`$ | $`=\frac{1}{2}(\delta ^{\overline{\mathrm{R}}\mathrm{J}}\delta ^{\overline{\mathrm{S}}\mathrm{K}}\delta ^{\overline{\mathrm{R}}\mathrm{K}}\delta ^{\overline{\mathrm{S}}\mathrm{J}})\overline{L}_{((t^{\mathrm{JK}})^{}\tau ^{\overline{\mathrm{R}}\overline{\mathrm{S}}})}+\frac{1}{2}\overline{T}_{(t^{\mathrm{J}\mathrm{K}}\tau ^{\overline{\mathrm{R}}\overline{\mathrm{S}}})}^{\mathrm{A}\mathrm{B}}\delta _{\mathrm{AB}}^{\mathrm{JK}\overline{\mathrm{R}}\overline{\mathrm{S}}}+4\overline{\beta }\overline{T}_{(\tau ^{\overline{\mathrm{R}}\overline{\mathrm{S}}})^{}}^{\mathrm{JK}},`$ |
| | $`\mathrm{where}\delta _{\mathrm{AB}}^{\mathrm{JK}\overline{\mathrm{R}}\overline{\mathrm{S}}}(\delta ^{\mathrm{AK}}\delta ^{\mathrm{B}\overline{\mathrm{S}}}\delta ^{\overline{\mathrm{R}}\mathrm{J}}\delta ^{\mathrm{AK}}\delta ^{\mathrm{B}\overline{\mathrm{R}}}\delta ^{\overline{\mathrm{S}}\mathrm{J}}+\delta ^{\mathrm{A}\overline{\mathrm{S}}}\delta ^{\mathrm{BJ}}\delta ^{\overline{\mathrm{R}}\mathrm{K}}`$ |
| | $`\delta ^{\mathrm{A}\overline{\mathrm{R}}}\delta ^{\mathrm{JS}}\delta ^{\overline{\mathrm{S}}\mathrm{K}}+\delta ^{\mathrm{A}\overline{\mathrm{S}}}\delta ^{\mathrm{BK}}\delta ^{\overline{\mathrm{R}}\mathrm{J}}\delta ^{\mathrm{A}\overline{\mathrm{R}}}\delta ^{\mathrm{KB}}\delta ^{\overline{\mathrm{S}}\mathrm{J}}`$ |
| | $`+\delta ^{\mathrm{AJ}}\delta ^{\mathrm{B}\overline{\mathrm{S}}}\delta ^{\overline{\mathrm{R}}\mathrm{K}}\delta ^{\mathrm{JA}}\delta ^{\overline{\mathrm{R}}\mathrm{B}}\delta ^{\overline{\mathrm{S}}\mathrm{K}}+\delta ^{\mathrm{A}\overline{\mathrm{S}}}\delta ^{\mathrm{BK}}\delta ^{\overline{\mathrm{R}}\mathrm{J}}),`$ |
$`(39)`$
where the symmetry of the indices on the left hand side should be imposed on the indices on the right side. In the above, we have sometimes suppressed the indices associated with the functions used by the generators. For example $`\omega ^{\overline{\mathrm{V}}_1\mathrm{}\overline{\mathrm{V}}_\mathrm{n}}`$ the associate of the $`\overline{U}`$ dual element may be written as $`\omega ^{\{\mathrm{V}_\mathrm{m}\}}`$ or simply as $`\omega `$. Also the notation $`\delta _{\mathrm{J}_1\mathrm{}\mathrm{J}_\mathrm{m}}^{\mathrm{I}_1\mathrm{}\mathrm{I}_\mathrm{m}}\delta _{\mathrm{J}_1}^{\mathrm{I}_1}\mathrm{}\delta _{\mathrm{J}_\mathrm{m}}^{\mathrm{I}_\mathrm{m}}`$ was utilized.
There is quite a bit of interchange between functions in various sectors suggesting that the inhomogeneous contribution in the transformation laws for $`D`$, is the interaction of the central extension with $`D`$. The coadjoint of the Virasoro algebra, i.e. the action of L on the coadjoint vectors reveals a spectrum of states containing:
* $`D`$ corresponds to a rank 2 covariant tensor when the central extension is set to zero and is otherwise a quadratic differential.
* $`\psi ^{\overline{\mathrm{I}}}`$ corresponds to $`N`$ spin-$`\frac{3}{2}`$ fields that partner with $`D`$.
* $`\tau ^{\overline{\mathrm{R}}\overline{\mathrm{S}}}`$ corresponds to the spin-1 covariant tensors that serves as the $`N(N1)/2`$ SO($`N`$) gauge potentials associated with the supersymmetries.
* Given the $`N`$ supersymmetries there are the fields $`\omega ^{\overline{\mathrm{V}}_1\mathrm{}\overline{\mathrm{V}}_\mathrm{p}}`$. For a fixed value of $`N`$, the total number of independent components is given by
$$\mathrm{\#}(U)=N(\mathrm{\hspace{0.17em}2}^NN1).$$
* Again given $`N`$ supersymmetries, there are the fields $`\rho ^{\overline{\mathrm{T}}_1\mathrm{}\overline{\mathrm{T}}_\mathrm{p}}`$. For a fixed value of $`N`$, the total number of independent components is given by
$$\mathrm{\#}(R)=(\mathrm{\hspace{0.17em}2}^NN1).$$
The spins of the fields associated with $`U`$ and $`R`$ vary according to $`(2\frac{p}{2})`$. These likely correspond to other gauge and non-gauge physical fields, auxiliary, and Stuckelberg fields that are required to close the supersymmetry algebra.
We end our discussion with a conjecture. If M-theory possesses a 1D NSR formulation, it seems likely that the $`N`$ = 32 or 16 case of the present discussion determines the structure of its representation. We conjecture that the spectrum of the 1D, $`N`$ = 32 or 16 $`𝒢`$ super Virasoro theory provides a set of fundamental NSR variables to describe M-theory.
“Equations were drawn up in paisley form.” – Rakim (1997)
Acknowledgment
We thank Takeshi Yasuda for discussion. V.G.J.R. and C.C. thank the University of Maryland for hospitality where this work was initiated. S.J.G. likewise acknowledges the University of Iowa for hospitality during the completion of this work. This work was supported in part by NSF grant # PHY-96-43219. |
warning/0002/cond-mat0002199.html | ar5iv | text | # Critical behaviour of annihilating random walk of two species with exclusion in one dimension
## I Introduction
Non-universal dynamical behaviour seems to be a controversial issue in non-equilibrium models. An outstanding example is the debated behaviour of systems exhibiting infinitely many absorbing states . There is no analytic treatment up to now; argumentation of various authors, in most of the cases, is based on simulation results. Despite intensive study, the critical behaviour of such systems is poorly understood, non-universality remains an unresolved problem and even scaling behaviour is questioned. Roughly speaking, in these coupled processes the ’primary’ particles follow a branching and annihilating random walk while the other species just provide a slowly changing environment that effects the branching rates of the primaries. The spreading exponents of the primaries depend on the initial conditions of the environment.
A possible way which might lead to a deeper understanding of the mechanism behind non-universal spreading could be the study of simpler coupled systems. Perhaps the simplest case is the coupled annihilating random walk of two species ( $`A+A\mathrm{}`$, $`B+B\mathrm{}`$). Naively one would expect that this could be described by the exactly solved field theory of the $`A+A\mathrm{}`$ process (ARW). In one dimension, however, the situation is more subtle than in higher dimensions. Particles of different type can block the motion of each other. The difference between one and two dimensions has been found to give rise to different phase diagrams in the case of the general epidemic model . The question now arises how relevant the exclusion perturbation caused by this blocking mechanism is to a fixed point of the kind determined in .
An other motivation of this study originates from the investigations of Hinrichsen , who found, by simulations, a strange scaling behaviour in some special case of his model ,for which, however, an explanation is still lacking. In section II Hinrichsen’s model will be introduced. It is easy to see, that the kinks in this model at the symmetry point corresponding to the compact directed percolation point of the Domany-Kinzel automaton, exhibit the process described above. In sections III and IV we present our high precision time dependent simulation results from random and seed initial conditions. In section V these results are compared with those obtained by rigid (i.e. parabolic) boundaries. We further investigate this analogy on the mean-field level in section VI, while section VII is devoted to results in the explicit two-species Annihilating Random Walk model with exclusion (ARW2e). We summarise our numerical results in section VIII and give an outlook toward $`N`$-species generalisation in IX. A qualitative description of the behaviour of Hinrichsen’s model outside the symmetry point on the line of compactness is presented in section X and finally in section XI we summarise and discuss our results.
## II The Generalised Domany-Kinzel SCA
The Domany-Kinzel (DK) stochastic cellular automaton (SCA) is one of the simplest models which show a non-equilibrium phase transition into an absorbing state. This one dimensional SCA is defined on a ring with two states ’1’ and ’0’ with the following rule of update:
```
t: 0 0 0 1 1 0 1 1
t+1: 0 p p q
```
where at $`t+1`$ the probability of ’1’-s is shown .
In the plane of the parameters $`(p,q)`$, the phase diagram of the DK SCA is as follows. A line of critical points separates the active phase ( with a finite concentration of ’1’-s) from the absorbing (vacuum) phase (with zero steady state density of ’1’-s). This continuous transition belongs to the universality class of directed percolation (DP) . The endpoint of this line ($`q=1`$, $`p=1/2`$) describes a transition, however, outside the DP class; it corresponds to compact directed percolation. Here the model exhibits Ising symmetry and can be solved exactly .
In 1997 Hinrichsen introduced a generalized version of the DK model including more than one symmetric inactive states ($`I1`$, $`I2`$, …) and one active state ($`A`$). The motivation for this study was to look for a possible change in the universality class of the line separating the active and passive steady states. This generalized DK model (GDK in the following), in its simplest form with two absorbing states $`I1`$ and $`I2`$ has been defined by the rules given below:
| $`s_1,s_2`$ | $`P(A|s_1,s_2)`$ | $`P(I_1|s_1,s_2)`$ | $`P(I_2|s_1,s_2)`$ |
| --- | --- | --- | --- |
| $`AA`$ | $`q`$ | $`(1q)/2`$ | $`(1q)/2`$ |
| $`AI_1`$ | $`p`$ | $`1p`$ | $`0`$ |
| $`AI_2`$ | $`p`$ | $`0`$ | $`1p`$ |
| $`I_1A`$ | $`p`$ | $`1p`$ | $`0`$ |
| $`I_2A`$ | $`p`$ | $`0`$ | $`1p`$ |
| $`I_1I_1`$ | $`0`$ | $`1`$ | $`0`$ |
| $`I_1I_2`$ | $`1`$ | $`0`$ | $`0`$ |
| $`I_2I_1`$ | $`1`$ | $`0`$ | $`0`$ |
| $`I_2I_2`$ | $`0`$ | $`0`$ | $`1`$ |
The geometry of updating is the same as in the case of the DK SCA. It has been shown by simulation that the phase diagram which emerges is similar to that in the DK SCA : an active phase is separated from an inactive one by a line of continuous phase transitions. The inactive phase, however is symmetrically degenerated ($`I1`$ or $`I2`$) and the phase transition line now belongs to the parity conserving (PC) universality class. This class has been studied by many authors as the first exception from the robust DP class .
The phase diagram exhibited on Fig1.b) shows that the line of PC transitions ends at $`q=1`$, $`p=1/2`$, a point which corresponds the Ising symmetry point of the DK automaton. The primary aim of the present work is to investigate the scaling properties of GDK at this point, which will be called CDP2 transition point. A typical time evolution of the GDK model at this special point when starting from a random initial arrangement of $`I1`$-s, $`I2`$-s and $`A`$-s is shown on Fig. 2. Here active islands can be spatially extended, thus three kinds of compact clusters can grow. Nevertheless only the $`I1`$ and $`I2`$ phases are $`Z_2`$-symmetric while the active phase plays a special role. (The situation is different from a 3-states Potts model with Glauber kinetics ).
It is well known that the CDP process in 1d is equivalent to an annihilating random walk process of kinks separating compact domains of $`0`$-s and $`1`$-s. In the model investigated here two types of kinks can be defined, namely kink $`K1`$ between domains $`AI1`$ (and $`I1A`$) and kink $`K2`$ between neighbouring $`AI2`$-s (and $`I2A`$-s). The rules of the model inhibit occurrence of kinks between domains of absorbing phases, i.e. between $`I1`$-$`I1`$ and $`I2`$-$`I2`$.
Kinks $`K1`$ and $`K2`$ perform annihilating random walks: $`K1+K1\mathrm{}`$, $`K2+K2\mathrm{}`$, while the processes $`K1+K2\mathrm{}`$, $`K2+K1\mathrm{}`$ are, however, forbidden. In other words, upon meeting, a $`K1`$ and a $`K2`$ “block” each other (do not annihilate and do not exchange sites). To our knowledge such kind of kinetics has not been studied before. Motivated by this fact we have decided to explore the critical behaviour of the above described system, on the line of compactness ($`q=1`$), by computer simulation. In this study special attention will be paid to the $`p=1/2`$ symmetry point CDP2.
## III Simulations from random initial state
### A Kink decay simulations
We have performed time dependent simulations starting from states with uniformly distributed species $`A`$, $`I1`$, and $`I2`$, with respective densities: $`\rho _0(A)`$, $`\rho _0(I1)`$ and $`\rho _0(I2)`$. At the CDP2 point unusual scaling behaviour of the density of kinks has been observed previously : a deviation from the ordinary annihilation-diffusion process with kink-density decay $`\rho (t)\frac{1}{\sqrt{t}}`$. Instead, $`\rho (t)t^\alpha `$ with $`\alpha 0.55`$ has resulted from the first simulations.
To check whether the observed deviation from standard ARW behaviour is only a cross-over effect, or it heralds some basic feature of altered kinetics we have performed very long-time ($`t_{max}=10^6`$ MCS) simulations on systems with $`L=24000`$ (Fig. 3) . Throughout the whole paper $`t`$ is measured in units of Monte-Carlo sweeps.
Figure 4 shows the results of simulations. It is seen that the deviation from the standard ARW value of the decay exponent remains present asymptotically as well: the local slopes of the decay curves
$$\alpha (t)=\frac{\mathrm{ln}\left[\rho (t)/\rho (t/m)\right]}{\mathrm{ln}(m)}$$
(1)
(where we use $`m=8`$ usually) go to constant values. Moreover, another interesting feature has become apparent: the kink-decay exponent depends on the initial concentrations of the components $`\rho _0(I1)=\rho _0(I2)`$ and in such a way that for higher initial kink density (lower average distance between the kinks) the decay is faster.
Asymptotically, as $`\rho _00`$, the average distance of dissimilar kinks goes to zero and the decay exponent tends to the ARW value: $`\alpha 0.5`$.
In the case of asymmetric initial condition ($`\rho _0(I1)\rho _0(I2)`$) $`K1`$-s and $`K2`$-s decay with different rates. The type that has smaller initial density decays faster. Example: in case of $`\rho _0(I1)=1/9`$, $`\rho _0(I2)=1/3`$, $`K2`$ decays roughly like $`t^{0.5}`$ (unperturbed by $`K1`$-s) but the local slopes of the log-log $`\rho (K1)t`$ dependence decrease from $`0.5`$ strongly.
## IV Simulations from an active seed
The cluster simulations were started from a state with uniformly distributed $`A`$-s and $`I1`$-s except a single $`I2`$ pair in the middle and the following characteristic quantities for the $`I2`$-s were followed:
* the average number of $`I2`$-s, $`N_{I2}(t)`$,
* their survival probability $`P_{I2}(t)`$,
* and the average mean square distance of spreading of $`I2`$’s from the center $`R_{I2}^2(t)`$
The above quantities were averaged over $`N_s`$ independent runs at the CDP2 point ( in case of $`R_{I2}^2(t)`$ only for surviving samples ). At the critical point we expect these quantities to behave for $`t\mathrm{}`$, as
$$N_{I2}(t)t^\eta ,$$
(2)
$$P_{I2}(t)t^\delta ,$$
(3)
$$R_{I2}^2(t)t^z.$$
(4)
Upon varying the initial density $`\rho _0(I1)`$, for the exponents $`\delta `$ and $`\eta `$ (defined similarly to eq.(1), the local slopes of $`N_{I2}(t)`$ and $`P_{I2}(t)`$) continuously changing values have been observed (Figs. 5, 6). The deviation of these exponents from those of the single-species annihilation random walk process: $`1/2`$ and $`0`$, respectively is remarkable. The spreading exponent, $`z`$, on the other hand, seems to be constant within numerical accuracy and equals that of the single species annihilating random walk : $`z=2/Z=1`$ such that the generalised hyper-scaling law of the compact directed percolation
$$\eta +\delta =z/2$$
(5)
is satisfied. In this respect it is important that $`\eta `$ has been found to be negative.
## V Cluster simulations of Compact directed percolation confined in a parabola.
To understand the physics of our numerical results up to now we set up a parallelism with an other case where DP process is bounded by parabolic space-time boundary conditions. We perform simulations on the compact cluster version of this and compare the results with those of the GDK model in section VIII.
Kaiser and Turban have investigated the $`1+1`$ dimensional DP process limited by a special, parabolic boundary condition in space and time directions:
$$y=\pm Ct^k$$
(6)
where $`C`$ changes under uniform length rescaling (by $`b`$) to:
$$C^{^{}}=b^{Zk1}C$$
(7)
Here $`Z`$ is the dynamical critical exponent. By referring to conformal mapping of the parabola to straight lines and showing it in the mean-field approximation Kaiser and Turban claim that for $`k<1/Z`$ this surface gives relevant, for $`k>1/Z`$ irrelevant and for $`k=1/Z`$ marginal perturbation to the DP process. The marginal case results in $`C`$ dependent non-universal power-law decay, (for details see next section), while for the relevant case stretched exponential functions have been obtained. The above authors have given support to this claim by numerical simulations.
We have investigated the effect of parabolic and reflecting boundary conditions for the CDP2 process numerically. Time-dependent cluster spreading simulations have been performed in the GDK model with parabolic boundaries such that at each time step the simulation region is bounded by two $`I1`$-s at $`y_{min}`$ and $`y_{max}`$, where
$$y_{min}=L/22Ct^k$$
(8)
$$y_{max}=L/2+2+Ct^k$$
(9)
Two $`I2`$-s have been put initially at the centre $`(L/2,L/2+1)`$ and some initial space (two $`A`$-s to the left and right) between $`K1`$-s and $`K2`$-s has been added. Therefore the role of $`I1`$-s now is purely the formation of parabolic boundaries around $`I2`$-s and in fact we investigate the plain CDP process with reflective boundary conditions. A typical 1+1 dimensional run looks like as shown below (1, 2 and 0 stand for $`I1`$, $`I2`$ and $`A`$, respectively):
```
<-- y -->
| 10022001
| 10222001
1002200001
t 1022000001
1222200001
| 1022220001
| 1022220001
| 100222200001
| 100222220001
V 100222200001
100222000001
102222000001
```
When we fixed the exponent at $`k=1/2`$, to make the situation marginal we found continuously changing exponents for the exponents of the survival probability $`P_{I2}`$ and the number of $`I2`$-s,$`N_{I2}`$, by varying the shape $`C`$ (Fig. 7). One can see that the exponent-slopes of $`N_{I2}(t)`$ (Fig. 8) and those of $`P_{I2}(t)`$ (Fig. 9) change by varying $`C`$. The spreading exponent of the $`I2`$-s, $`z`$, seems to be constant: equal to unity. (Fig. 10). These results are very similar to those of the seed simulations in GDK of section IV.
The analysis based on local slopes (Figs. 8, 9, 10) shows again plateaus for high values of $`t`$, indicating true power-law behaviours. The magnitude of the exponent characterizing the decay of the density of $`I2`$-s decreases as $`C`$ is increased reminiscent of a similar situation in .
The survival exponent changes in such a way that the hyper-scaling relation valid in case of compact directed percolation:
$$z/2=\eta +\delta =1/2$$
is fulfilled. In this case it is important, again, that $`\eta `$ takes negative values.
## VI Theoretical considerations for CDP confined in a parabola
### A Anisotropic scaling
In 1+1 dimensional anisotropic systems the correlation length diverges as $`\xi _{}t^\nu _{}`$ in time and as $`\xi _{}t^\nu `$ in the space direction with a dynamical exponent $`Z=\nu _{}/\nu `$ ($`\nu `$ is also denoted as $`\nu _{}`$ in the literature). Covariance under a change of the length scales then requires two different scaling factors, $`b_{}=b^Z`$ and $`b_{}=b`$ .
We consider now a system displaying anisotropic critical behaviour and limited by a free surface in the $`(t,y)`$ plane as given in eq(6). Under rescaling, with $`t^{}=t/b^Z`$ and $`y^{}=y/b`$, $`C`$ transforms according to eq.(7),as discussed in the previous section.
In the marginal case, which we will consider now, $`Z=1/k`$, the scaling dimension $`x_m`$ of the tip order parameter becomes $`C`$–dependent $`x_m(C)`$.
The order parameter correlation function between the origin and a point at $`(t,y)`$ transforms as
$$G(\mathrm{\Delta },t,y,\frac{1}{C})=b^{2x_m}G(b^{1/\nu }\mathrm{\Delta },\frac{t}{b^Z},\frac{y}{b},\frac{b^{1Zk}}{C})$$
(10)
when $`L`$ is infinite. With $`b=t^{1/Z}`$, equation (10) leads to:
$$G(\mathrm{\Delta },t,y,\frac{1}{C})=t^{2x_m/Z}g(\frac{t}{\mathrm{\Delta }^\nu _{}},\frac{y^Z}{t},\frac{t}{l_C}).$$
(11)
Here $`l_C=C^{\frac{Z}{1Zk}}`$, $`\mathrm{\Delta }=\frac{(pp_c)}{p_c}`$ and $`x_m`$ is the scaling dimension of the order parameter. The latter is connected to $`\beta `$, the critical exponent of the order parameter via $`\beta =\nu x_m`$. We will use this scaling form in the following. $`\beta `$ is the usual order-parameter exponent, defined, for the DKCA, through $`\rho _1(pp_c)^\beta `$, for $`p>p_c`$; $`\rho _1`$ is the stationary density of 1’s. In case of a first order transition as is the case with compact directed percolation the following considerations hold.
As already mentioned,
$$P(t)t^\delta ,$$
(12)
is the survival probability of $`1^{}s`$ for spread of particles (1’s, in our notation) about the origin. Away from the critical point $`\beta ^{}`$ governs the ultimate survival probability (starting from a localized source): $`P_{\mathrm{}}lim_t\mathrm{}P(t)(pp_c)^\beta ^{}`$. It is known that $`\beta ^{}=1`$ in CDP. The order-parameter exponent, $`\beta `$, however, is zero. This is because $`p=1/2`$ marks a discontinuous transition, by symmetry. $`\rho _1=0`$ for $`p<1/2`$ and $`\rho _1=1`$ for $`p>1/2`$; strictly speaking, $`\beta `$ is not defined here but it is natural to associate the value $`\beta =0`$ with the discontinuous transition.
This problem with the ill-defined exponent $`\beta `$ can be avoided following the lines of Grassberger and de la Torre’s scaling argument for discontinuous transitions as it has been presented by Dickman and Tretyakov. Consider a model with a transition from an absorbing to an active state at $`\mathrm{\Delta }=0`$, with exponents $`\delta `$, $`\eta `$, $`z`$, and $`\beta ^{}`$ defined as above. Suppose, however, that the order parameter $`\rho `$ is discontinuous, being zero for $`\mathrm{\Delta }<0`$, and
$$\rho =\rho _0+f(\mathrm{\Delta }),$$
(13)
for $`\mathrm{\Delta }>0`$, where $`\rho _0>0`$, and $`f`$ is continuous and vanishes at $`\mathrm{\Delta }=0`$. According to the scaling hypothesis for spreading from a source there exist two scaling functions, defined via
$$\rho (y,t)t^{\eta dz/2}\stackrel{~}{G}(y^2/t^z,\mathrm{\Delta }t^{1/\nu _{||}}),$$
(14)
and
$$P(t)t^\delta \mathrm{\Phi }(\mathrm{\Delta }t^{1/\nu _{||}}).$$
(15)
(Here $`\rho (y,t)`$ is the local order-parameter density. $`\tau \mathrm{\Delta }^{\nu _{||}}`$.) Existence of the limit $`P_{\mathrm{}}`$ implies that $`\mathrm{\Phi }(x)x^\beta ^{}`$ as $`x\mathrm{}`$, with $`\beta ^{}=\delta \nu _{||}`$. In a surviving trial, the local density must approach the stationary density $`\rho `$ as $`t\mathrm{}`$, so $`\rho (y,t)\mathrm{\Delta }^\beta ^{}\rho _0`$, for $`t\mathrm{}`$ with fixed $`y`$, and $`\mathrm{\Delta }`$ small but positive. It follows that $`\stackrel{~}{G}(0,x)x^\beta ^{}`$ for large $`x`$.
An important consequence is that we can use as scaling dimension of the order parameter for CDP the value $`\beta ^{}`$ in the relation $`x_m=\frac{\beta }{\nu }`$ instead of $`\beta `$. Via scaling relations $`\beta ^{}=\delta \nu _{||}`$, the values obtained by computer simulations for $`\delta `$ will be compared with results for CDP+parabolic boundary conditions. In this context the connection, again via scaling relation, between $`\delta `$ and the decay exponent of the density of kinks when starting from a random initial state $`\alpha `$ will also be made use of .
### B Mean field analysis for CDP confined in a parabola
In this section we will follow the lines of the mean-field analysis of the $`1+1`$ dimensional DP process confined by a parabola as given in ref., but now applied to compact directed percolation.
The basic processes are:
The order parameter correlation function is the probability density $`P(t,y)`$ for a site at $`(t,y)`$ to be connected to the origin.
First we consider the case without confinement. In mean-field approximation one can set up an equation for the connectedness at $`(t+1,y)`$:
$`P(t+1,y)`$ $`=`$ $`p\{P(t,y+1)[1P(t,y1)]+`$ (16)
$`P(t,y1)`$ $`[`$ $`1P(t,y+1)]\}+P(t,y+1)P(t,y1)`$ (17)
Going to the continuum limit the following differential equation is obtained:
$$\frac{P}{t}=p\frac{^2P}{y^2}+(2p1)P+(12p)P^2$$
(18)
The homogeneous, stationary solution of eq.(18) is:
$`P_0=\{\begin{array}{ccc}1\hfill & \text{for}\hfill & p>1/2\hfill \\ 0\hfill & \text{for}\hfill & p1/2\hfill \end{array}`$ (21)
describing a first order transition for CDP at $`p_c=1/2`$, as it is the case. At the transition, $`p=p_c`$, eq.(18) reduces to
$$\frac{P}{t}=\frac{1}{2}\frac{^2P}{y^2}$$
(22)
This is the ordinary diffusion equation with RW solution
$$P(t,y)=\frac{\mathrm{exp}\left(\frac{y^2}{2t}\right)}{\sqrt{2\pi t}}$$
(23)
which is exact in the CDP case. From comparison with the scaling form in the previous subsection, the following (well-known) exponents for CDP arise:
$$\nu _{}=1\nu =1/2Z=2x_m=\frac{1}{2}$$
(24)
On a parabolic system, we use the new variables $`t`$ and $`\zeta (t,y)=y/t^k`$ for which the free surface is shifted to $`\zeta =\pm C`$ and equation (22) is changed into
$$\frac{P}{t}=\frac{1}{2t^{2k}}\frac{^2P}{\zeta ^2}+k\frac{\zeta }{t}\frac{P}{\zeta }$$
(25)
with the boundary condition $`P(t,\zeta =\pm C)=0`$. Through the change of function
$$P(t,\zeta )=\mathrm{exp}\left[\frac{k}{2}\zeta ^2t^{2k1}\right]Q(t,\zeta )$$
(26)
equation (25) leads to
$$\frac{Q}{t}=\frac{1}{2t^{2k}}\frac{^2Q}{\zeta ^2}+\frac{k}{2}\left[(k1)\zeta ^2t^{2k2}\frac{1}{t}\right]Q$$
(27)
for which the variables separate when $`k=1,1/2`$ or $`0`$. These values of $`k`$ just correspond to irrelevant, marginal and relevant perturbations.
For $`k=1`$ the critical behaviour is the same as for un–confined percolation as expected for an irrelevant perturbation.
For the true parabola which is the marginal geometry, one may use equation (25) with $`k=1/2`$ to obtain
$$t\frac{P}{t}=\frac{1}{2}\frac{^2P}{\zeta ^2}+\frac{\zeta }{2}\frac{P}{\zeta }\zeta =\frac{y}{t^{1/2}}$$
(28)
which is of the form studied in for the directed walk problem. Writing $`Q(t,\zeta )=\varphi (t)\psi (\zeta )`$ leads to the following eigenvalue problem for $`\psi (\zeta )`$
$$\frac{1}{2}\frac{d^2\psi }{d\zeta ^2}+\frac{\zeta }{2}\frac{d\psi }{d\zeta }=\lambda ^2\psi $$
(29)
with $`\varphi (t)t^{\lambda ^2}`$. The solution is obtained as the eigenvalue expansion
$$P(t,y)=\underset{n=0}{\overset{\mathrm{}}{}}B_nt_1^{\lambda _n^2}F_1[\lambda _n^2,\frac{1}{2};\frac{y^2}{2t}].$$
(30)
The behaviour at large $`t`$ is governed by the first term in this expansion which decays as $`t^{\lambda _0^2}`$, i. e. with a $`C`$–dependent exponent as expected for a marginal perturbation. The dimension of the tip–to–bulk correlation function is the sum of the tip and bulk order parameter dimensions, the first one being variable. Comparing with the form of the decay in eq.(11) gives $`\lambda _0^2=[x_m^{mf}(C)+x_m]/Z`$ and, using (24), the tip order parameter dimension is given by
$$x_m^{mf}(C)=2\lambda _0^2\frac{1}{2}.$$
(31)
Its dependence on $`C`$ is shown in figure 2. of
Analytical results can be obtained only in limiting cases which have already been discussed in . When $`C`$ is infinite, $`\lambda _0^2=1/2`$, only the first term in the expansion remains, which satisfies the initial and boundary conditions, giving back the free solution in equation (23). For large C-values the tip exponent is $`x_m^{mf}(C)=\frac{1}{2}+\sqrt{\frac{2}{\pi }}C\mathrm{exp}\left(\frac{C^2}{2}\right)[1+O(\epsilon )]`$ where $`\epsilon `$ is the correction term itself. For narrow systems, the hyper-geometric function gives a cosine to leading order in $`C^2`$. One obtains the following asymptotic behaviour in $`t`$
$$P(t,y)t^{\pi ^2/8C^2}\mathrm{cos}\left(\frac{\pi y}{2C\sqrt{t}}\right)$$
(32)
and the tip exponent diverges as $`\pi ^2/4C^2`$.
For $`0<k<1/2`$ the dependence on $`t`$ is expected to be a stretched exponential function. For details see .
## VII Annihilating random walk of two species with exclusion
To check our results concerning the scaling properties of kinks in the GDK model at the CDP2 point we have carried out an explicit simulation of the annihilating random walk of two species ($`A`$, $`B`$) with exclusion. The model we have investigated has been suggested by Hinrichsen and is as follows. A(B) will hop to a neighbouring empty site with probability $`p1A`$ ($`p1B`$) or annihilate with a neighbour A (B) with probability $`p2A`$ ($`p2B`$) while A and B do not react when getting into neighbouring positions. The initial configuration was chosen in such a way that allows pairs of the same kind to annihilate always within some finite time interval (i.e. the system evolves into an empty state),namely:
```
... A.A.B.B.A.A.A.A.B.B....
```
That means that $`AA`$ and $`BB`$ pairs have been put in a 1d ring with initial probability $`\rho (0)`$. Had we not chosen the initial state like this the system would have ended up in some finite particle configuration where $`A`$-s and $`B`$-s follow each other alternatingly, separated by arbitrary empty regions. (This initial configuration is in agreement with the arrangement of the two kinds of kinks in some random initial state of the GDK model, too. ) The probabilities $`p1A`$, $`p1B`$, $`p2A`$ and $`p2B`$ have been chosen to be unity, to achieve maximum simulation effectiveness; no qualitative difference in the results have been found upon lowering them.
Clearly this process is different from the simple annihilating random walk of two species $`A+B\mathrm{}`$ , therefore we may expect that a field theory describing this model (which, however, is still missing) would result in a different fixed point with different critical exponents as well. Furthermore one can argue that when comparing the simple random walk and the random walk+exclusion (SEP) processes one also observes different dynamical behaviours. This latter case is nothing else but the $`T=0`$ dynamics of the $`1d`$ Ising Model with Kawasaki exchange, where we have different domain growth properties than in case of a simple random walk.
An extensive numerical simulation with look-up table algorithm seems to confirm this expectation. As Figure 12 shows the slopes of the density decay started from the special pairwise random states described above depend on the initial density $`\rho (0)`$.
The local slopes tend to constant values greater than $`\alpha =0.5`$ in agreement with the GDK kink results. The level-off in case of $`\rho _0=0.5`$ happens only for $`t>1.5\times 10^6`$ MCS. The average $`AA`$ and $`BB`$ distances confining an other type of particle have also been measured during the simulations, that enables us to extract the amplitude ($`C`$) of the confinement in the function fitted ($`C\times t^\alpha `$). These values will be used to compare the results with those of the GDK (see next section).
## VIII Summary of time dependent results
Since in all of the previously shown cases we found non-universal scaling depending on the initial conditions and the generic model to account for such behavior seems to be the CDP2 with parabolic boundary condition, we have decided to measure the region of confinement in all cases and plotted the survival probability exponents $`\delta `$ and the kink decay exponents $`\alpha `$ as a function of the shape of the measured parabola.
In the present case $`\beta ^,`$ the final survival probability of a cluster plays the role of the order parameter exponent $`\beta `$ as explained at the end of Sec.VI and for the characteristic exponents we have: $`\delta =\beta ^{^{}}/\nu _{}`$. Thus we have plotted the results for $`\delta (C)`$. In a common graph the fitted values for $`\alpha (C)`$ are also shown; on the level of kinks the order parameter $`\beta `$ is connected to $`\alpha `$ in the same way as $`\beta ^{}`$ to $`\delta `$ for ’spins’ (see eg. ).
For random initial conditions in the GDK model the characteristic distance between two neighbouring kinks of a given type has also been measured. The average neighbour distance $`l_{K2K2}`$ shown on Fig. 13 have been obtained for initial densities: $`\rho _0(I1)=\rho _0(I2)=1/3`$. The power law increase for large $`t`$ (see the plateau for $`t>30.000`$) with the same scaling exponent as the decay exponent is not very surprising because $`\rho _{kink}(t)1/l_{K2K2}`$.
Since the $`K2K2`$ and the $`K1K1`$ pairs confine the motion of each other ( a $`K1K2`$ pair can not exchange to $`K2K1`$) this power-law increasing length scale imposes a ’stochastic’ boundary condition (pressure on kinks) with a mean value of a parabola similarly that was investigated by Kaiser and Turban in case of $`1+1`$ d DP processes and adapted for the case of a CDP-like first order transition in section VI. As discussed before, the scaling dimension of the order parameter changes continuously with the amplitude of the parabolically growing confining box size if it grows with the same exponent as the cluster inside.
In our case we encounter similar situation. The kink density decay exponent $`\alpha `$ seems to vary continuously in case of symmetrical initial conditions. The initial conditions effect the amplitudes of the density decays (as was shown to be valid by field theory for pure the reaction diffusion of A,B particles ) and therefore the amplitudes of the confinement region sizes ($`C`$) (see Fig. 14). To compare our results with those of the form $`A+Ct^\alpha `$ has been fitted to the $`l_{K2K2}(t)`$ distances determined from the density decay simulations (assuming $`l_{K2K2}(t)=2/\rho _{kink}(t)`$). The following table summarises the results for GDK with random initial conditions:
| $`\rho _0`$ | $`C`$ | $`\alpha `$ |
| --- | --- | --- |
| $`0.0`$ | $`\mathrm{}`$ | 0.5000(3) |
| $`0.05`$ | $`14.09`$ | $`0.517(2)`$ |
| $`0.10`$ | $`7.62`$ | $`0.528(2)`$ |
| $`0.15`$ | $`6.11`$ | $`0.534(2)`$ |
| $`0.2`$ | $`5.85`$ | $`0.537(2)`$ |
| $`0.3`$ | $`4.18`$ | $`0.540(1)`$ |
This is in agreement with Fig. 8 of , where an increasing $`C`$ causes a decreasing exponent. Note that the first line in the table corresponds to the simple ARW process, therefore there is no confinement (amplitude $`C`$ is $`\mathrm{}`$).
The main difference between the rigid parabola boundary case and the ”stochastic” confinement is that in the latter case the boundary generates an additional noise to the motion of confined particles. Therefore we don’t simple have ”free” particles confined in a parabola, but they are also perturbed by the noise in such a $`K1K2`$ symmetrical way that the outcome perturbation is marginal.
In ARW simulations again the form $`A+Ct^\alpha `$ has been fitted for the measured $`AA`$, $`BB`$ distances.
In cluster simulations we fitted the form: $`y=C\times t^\delta `$ for the region of confinements and determined the respective $`C`$-s in all cases.
As Figure 15 shows we obtained similar monotonically decreasing curves in all cases that also agrees with the results of . The GDK-uniform and the ARW2 results seem to lie on the same curve. The spreading simulation results of GDK are different from those of the CDP+parabolic boundary condition case.
This can be understood, however, since in the former case the confined particles ($`I2`$-s) have a back effect on the bulk ($`I1`$-s) particles, while this is not the case when the boundary is fixed.
We also show the $`\pi ^2/8C^2`$ curve, determined as the asymptotic solution $`C0`$ of the mean-field approximation, see section VI eq.(32). This seems to be in fair agreement with the case of CDP+fixed parabolic boundary condition.
## IX Generalisation for $`N>2`$: symmetric annihilating exclusion process of $`N`$ species
We have carried out preliminary simulations in the generalized version of the model introduced in section VII. The system was started from configurations like:
```
.... A..A.B...B..C.C..DD....E..EFF ...
```
where species of the same type can annihilate each other but different types can not exchange. Our results show that concerning the time dependence of the density the deviations from the square root decay persist for $`N>2`$ and this property seems to remain valid also for $`N\mathrm{}`$.
Fig. 16 shows that a similar level-off can be observed in the local slopes as in the $`N=2`$ case with $`\alpha >0.5`$. A tendency of increasing $`\alpha `$ with increasing $`N`$ is apparent in our simulations; the growing effect of finite size corrections, however, prevented us from going further, for higher values of $`N`$, in this study.
## X GDK on the line of compactness
On the line $`q=1`$ the role of the absorbing states ($`I1`$ and $`I2`$) is symmetric. The above reported simulation results for the GDK model refer to $`p=1/2`$, the CDP2 point. Now we will discuss the situation for $`p1/2`$. For $`p>1/2`$ the creation of new $`A`$-s happens with probability greater than 0.5, the active domain size grows exponentially and the inactive regions die out quickly (and symmetrically); the all-A- phase plays the same role as the all- 0- phase of the original DK automaton . The deviation from the DK picture is quite apparent,however, for $`p<1/2`$, as instead of the all-1-phase of DK, for all values of $`p`$ with the exception of $`p=0`$, Glauber-Ising -like kinetics governs the motion of kinks. The kinks here are extended objects (A-s), somewhat similarly to those in Grassberger’s CA models( for $`p<p_c`$, where $`p_c`$ is the critical point of a parity conserving phase transition ) where also kinks of different extension (and there even of different structure) separate absorbing-phase clusters. (At $`p=0`$ diffusion stops and a striped space-time picture of I1 and I2 domains freezes in, again like in Grassberger’s model A cited above.) The average size of $`A`$-s goes to zero between two domains of the same type quickly, the kink ’particles’ of same type perform a biased random walk toward each other. On the other hand between domains of different types remains a film of $`A`$-s of average size $`1`$, since a collision of an $`I1`$ and $`I2`$ domain always creates a new $`A`$ at the next time step. That means that kinks of different types still block the motion of each other. Therefore the role of $`A`$-s is similar to the kinks of the $`T=0`$ Glauber Ising model. On the whole line of $`1/2>p>0`$, in the long time limit, therefore one can expect the number of such kinks to decrease as ($`t^{1/2}`$). On the other hand this is a line of compactness, as all clusters growing from a seed are compact, the characteristic exponents ($`\eta `$, $`\delta `$ and $`z`$) though strongly dependent on $`p`$ and the composition of the initial state, satisfy the hyper-scaling law valid for compact cluster . This statement involves that this line is a line of first order transition points with order parameter exponent $`\beta _s=0`$. This first order transition occurs in a symmetry-breaking (’magnetic’) field coupled to the $`I1`$ and $`I2`$ spins. For the detailed description of a similar situation see . All these features have been supported by simulations. As an example we can give some results obtained at $`p=0.4`$ starting with a single $`I2`$ in the sea of $`I1`$, $`A`$-s ($`25\%`$ of $`I1`$, $`75\%`$ of $`A`$ ); $`\delta =0.45`$, $`z/2=0.47`$, $`\eta _{I2}=0.02`$. It is worth mentioning, that for hyper-scaling to hold it is again important that $`\eta _{I2}`$ is negative.
## XI Discussion
We have investigated numerically the one dimensional generalized Domany-Kinzel cellular automaton on a line in the plane of its parameters where only compact clusters grow. The two types of kinks in the simplest version of this compact GDK model (two absorbing phases) follow annihilating random walk with exclusion (no reaction) between different types. The equivalence with an explicit two-species ARW model with exclusion is shown provided the initial state is prepared in such a way that the kinks are arranged in pairs with some density. High precision simulations revealed that this system relaxes in a non-trivial way: the decay exponent of the kink density depends on the initial density of kinks. We argue that this is a kind of (internal) surface effect; similar to the ARW process confined by a rigid space-time parabola provided the power of the parabola is chosen to be marginal. This case has been explicitly investigated with the result that the spreading exponents behave qualitatively the same way as expected from the corresponding mean-field approximation. We have no proof of the marginality for the theory including fluctuations, but rely on symmetry arguments. If we assumed that particles would exert relevant perturbation ($`k<1/2`$) on each other, the corresponding parabola picture would predict stretched exponential decay ( a behaviour that is very difficult to differentiate from power-laws by simulations) and the local slopes should go to some higher value as a function of time (meaning faster that any power-law decay). However our high precision data show just the opposite, the local slopes decrease as a function of time tending to a value somewhat greater than $`1/2`$.
Nevertheless, the possibility of pure square-root decay masked by some tremendously long crossover function can not be ruled out. One could still expect a non-universal scaling of the survival probability of particles in the same way as was observed in or in an other similar situation where a diffusing ”prisoner” confined by marginally growing cage was investigated. In the latter case the the boundary condition was absorbing and an exact solution was possible giving an exponent for the survival probability which is a continuous function of the amplitude of the marginal parabola. Furthermore the survival of a diffusing prisoner (with diffusivity D) inside a cage where both walls diffuse (with diffusivity A) has been solved exactly and the decay exponent was found to be $`\pi /2\mathrm{cos}^1(D/(D+A)`$ .
Non-standard scaling in a 1d ARW model was also observed by Frachebourg et al. . They have shown that the survival probability of particles in an ARW with one free boundary depends on the location of the particles. If we count the particles from the free boundary the survival probability of odd particles decay with exponent $`0.225`$, while those of even numbered decay with exponent $`0.865`$. The explanation for this is based on the fact that even numbered particles always have left and right neighbours during the process, while odd numbered particles lack one of the neighbours and since the ARW in 1d is diffusion limited they can escape. One can notice a similarity of this mechanism to the one in ARW2 models we investigated. Namely in our case there are infinitely many internal boundaries (generated by particles of different types which can not exchange sites).
Recently Bray has shown that the relaxation towards the critical state in the 2d XY model depends on the initial state. This is very different from what is expected from field-theoretical RG predictions that can not take into account low-dimensional topological effects. Moreover, Bray has shown that the non-universal behavior of the persistence exponent in this case can be described by the random walk of a particle moving under an attractive central power-law force that creates marginal perturbation as compared to free random walk This scenario is similar to ours since we also have particles with RW exhibiting pressure of marginal strength on each other.
Right after our submission two other preprints appeared on cond-mat , dealing with models very similar to those presented here and reporting results which are in accord with ours for those quantities they also investigated.
Acknowledgements:
The authors would like to thank Z. Rácz, S. Redner, P. Arndt, U. Tauber for useful remarks and H. Hinrichsen for taking part in the early stages of this work. Support from Hungarian research fund OTKA (Nos. T-23791, T-25286 and T-23552) is acknowledged. One of us (N.M.) would like to thank R. Graham for hospitality at the Fachbereich Physik of Universität-GHS Essen, where this work has been completed. G. Ó. acknowledges support from Hungarian research fund Bólyai (No. BO/00142/99) as well. The simulations were performed partially on Aspex’s System-V parallel processing system (www.aspex.co.uk). |
warning/0002/astro-ph0002081.html | ar5iv | text | # The Contribution of Faint Blue Galaxies to the Sub-mm Counts and Background
## 1 Introduction
The SCUBA camera \[Holland et al. 1999\] on the James Clerk Maxwell Telescope has transformed our knowledge of dusty galaxies in the distant Universe as a result of the discovery of a new population of luminous, dusty, infra-red galaxies (Smail et al. 1997; Ivison et al 1998). It has been proposed that these galaxies may be similar to IRAS ULIRGs (ultra-luminous infra-red galaxies) which appear to be starbursting/AGN galaxies, containing large amounts of dust. The possibility that much star-formation is hidden by dust means that sub-mm observations can give an invaluable insight into the star-formation history of the Universe. This view aided by the redshifting of the thermal dust emission peak in starbursting galaxies into the FIR, which results in a negative k-correction in the sub-mm. By this route, we can therefore study our Universe all the way back to very early times and gain unprecedented insight into the formation and evolution of galaxies.
The first sub-mm galaxy to be detected by SCUBA was SMM J02399-0136 \[Ivison et al. 1998\], which is a massive starburst/AGN at z=2.8 and the current situation is that the complete 850$`\mu `$m sample from all the various groups consists of well over 50 sources (Blain et al. 1999; Eales et al. 1999; Hughes et al. 1998; Holland et al. 1998; Barger et al. 1998; Smail et al. 1997). Optical and near infra-red(NIR) counterparts have been identified for about a third of the sources, although the reliability of these identifications varies greatly. This problem is due to the fact that the $`15^{\prime \prime }`$ FWHM of the SCUBA beam results in $`{}_{}{}^{+}3_{}^{\prime \prime }`$ positional errors on a sub-mm source, so there is a reasonable chance that several candidates could lie within these errors. Also, there is no guarantee that the true source will be detected down to the optical flux limit as, for example, many of the sources have been shown to be very red objects (Dey et al. 1999; Smail et al. 1999: Ivison et al. 2000) and therefore have not been found in optical searches for sub-mm sources.
What has proved extremely enlightening is that radio counterparts at 1.4GHz have now been identified for many of the sub-mm sources (Smail et al. 2000: Ivison et al. 2000) providing much more accurate angular positions ($`<1^{\prime \prime }`$ in some cases) and reasonably accurate photometric redshifts. Various groups have obtained redshift distributions of sub-mm samples (Hughes et al. 1998; Barger et al. 1999a; Lilly et al. 1999; Smail et al. 2000) and they all derive results that are consistent with a mean redshift in the range $`1<z<3`$. The fact that almost all of the sources are associated with mergers or interactions seems to confirm that the population of sources contributing at the ‘bright’ ($`>2`$mJy) sub-mm fluxes (since most of the sources so far discovered are ‘bright’) are similar to local IRAS ULIRG’s, ie massive, starbursting/AGN galaxies which are extremely luminous in the far-infra-red. This hypothesis is strengthened further by the fact that the only two sub-mm sources (SMM J02399-0136 and SMM J14011+0252) with reliable redshifts have been followed up with millimeter wave observations (Frayer et al. 1998, 1999), resulting in CO emission being detected at the redshifts of both sources (z=2.8 and z=2.6), a characteristic indicator of large quantities of molecular gas present in IRAS galaxies.
The nature of the fainter ($`1`$mJy) sub-mm population is, however, the focus of this paper. It has been claimed by Peacock et al. (2000) and Adelberger et al. (2000) that the Lyman Break Galaxy(LBG) population could not only contribute significantly to the faint sub-mm number counts, but could also account for a substantial proportion of the background at $`850\mu `$m. This may indicate that ULIRG’s cannot explain all of the sub-mm population and that the UV-selected galaxy population, which are predicted to be evolved spirals by the Bruzual & Charlot models, may in fact make a substantial contribution. It is exactly this hypothesis our paper addresses.
In this paper we will first review the situation regarding the optical galaxy counts, focusing in particular on the models of Metcalfe et al. (1996). These simple models which use a $`\tau =9`$Gyr SFR for spirals and include the effects of dust give good fits to galaxy counts and colours from U to K. The idea is then to see whether this combination of exponential SFR and relatively small amounts of dust in the first instance (A$`{}_{B}{}^{}=0.3`$ mag. for the $`1/\lambda `$ law), which would re-radiate the spiral ultra-violet (UV) radiation into the FIR, could cause a significant contribution to the sub-mm galaxy number counts and background at $`850\mu `$m. Our modelling will be described in section 3 and then in section 4 our predicted contribution to the $`850\mu `$m and $`60\mu m`$ galaxy counts and the extra-galactic background in the sub-mm will be shown. Also in this section we demonstrate how to get a fit to the background in the $`100300\mu m`$ range by using warmer, optically-thicker dust in line with that typically seen in ULIRG’s. We will then discuss the implications of our predictions in section 5 and conclude in section 6.
## 2 The Optical Counts
It is well known that non-evolving galaxy count models, where number density and luminosity of galaxies remain constant with look-back time, do not fit the optical number counts e.g. \[Shanks et al. 1984\], as there is always a large excess of galaxies faintwards of $`B22^m`$. One way to account for this excess of ’faint blue galaxies’ is to investigate the way galaxy evolution will influence the optical number counts. Metcalfe et al (1996) showed that by assuming that the number density of galaxies remains constant, the Bruzual and Charlot(1993) evolutionary models of spiral galaxies with a $`\tau =9`$Gyr SFR give excellent fits to the optical counts. The galaxy number counts are normalised at $`B18^m`$ so that the non-evolving models give good fits to the $`B`$ band data and redshift distributions in the range $`18^m<B<22.^m5`$. With this high normalisation, the models of the galaxy counts represent both spiral and early-type galaxies extremely well for $`17^m<I<22^m`$ (Glazebrook et al. 1995a, Driver et al 1995) and also the less steep $`H/K`$ counts out to $`K20^m`$. The evolution model then produces a reasonable fit to the fainter counts to $`B27^m`$, $`I26^m`$, $`H28^m`$.
Metcalfe et al (1996) included a $`1/\lambda `$ internal dust absorption law with $`A_B=0.3`$ for spirals to prevent the $`\tau =9`$ Gyr SFR from over-predicting the numbers of high redshift galaxies detected in faint B$`<24`$ redshift surveys (Cowie et al 1995). This $`1/\lambda `$ dust law differs from the Calzetti(1997) dust law derived for starburst galaxies, in that for a given $`A_B`$, more radiation is absorbed in the UV. The Calzetti dust law is used by Steidel et al(1999) to model their ‘Lyman Break’ galaxies; they find an average E(B-V)=0.15 which gives $`A_B=0.87`$mag and $`A_{1500}=1.7`$mag. This compares to our $`A_{1500}=0.9`$mag with $`A_B=0.3`$mag. Both models also fail to predict as red colours as observed for the U-B colours of spirals in the Herschel Deep Field (Metcalfe et al 1996). However, if we assumed E(B-V)=0.15 for our z=0 spirals, as compared to our E(B-V)=0.05, then the rest colours of spirals as predicted by the Bruzual & Charlot model might start to look too red as compared to what is observed. Otherwise, the main difference between these two dust laws is that the Calzetti law would produce more overall absorption and hence a higher FIR flux from the faint blue galaxies. Thus in some ways our first use of the $`1/\lambda `$ law appears conservative in terms of the predicting the faint blue galaxy FIR flux. Later, we shall experiment by replacing the $`1/\lambda `$ law with the Calzetti(1997) law in our model.
So this pure luminosity evolution (PLE) model with $`1/\lambda `$ dust and $`q_0=0.05`$ then slightly under-estimates the faintest optical counts but otherwise fits the data well, whereas for $`q_0=0.5`$ the underestimate (with or without dust) is far more striking. An extra population of galaxies has to be invoked at high redshift to attempt to explain this more serious discrepancy for the high $`q_0`$ model. This new population was postulated to have a constant SFR from their formation redshift until $`z1`$ and then the Bruzual & Charlot models predict a dimming of $`5^m`$ in $`B`$ to form a red dwarf elliptical (dE) by the present day and therefore has the form of a ’disappearing dwarf’ model \[Babul & Rees 1992\]. No dust was previously assumed in the dE population but this assumption is somewhat arbitrary.
The $`\tau =9`$Gyr SFR was inconsistent with the early observations at low redshift from Gallego et al.(1996) and this is partly accounted for by the high normalisation of the optical number counts at $`B18^m`$. There is still a problem with the UV estimates from the CFRS UV data of Lilly et al at z=0.2. More recent estimates of the global SFR at low redshift based on the \[OII\] line (Gronwall et al.1998; Tresse & Maddox 1998; Hammer and Flores 1998) indicates that the decline from z=1 to the present day may not be as sharp as first thought and that the $`\tau =9`$Gyr SFR in fact provides a better fit to this low redshift data. Metcalfe et al (2000) have further found that this model also agrees well with recent estimates of the luminosity function of the z=3 Lyman break galaxies detected by Steidel et al.(1999).
The main question then that we will address in this paper is whether the small amount of internal spiral dust absorption assumed in these models which give an excellent fit to the optical galaxy counts, could cause a significant contribution to the sub-mm number counts and background at $`850\mu `$m.
## 3 Modelling
Using the optical B band parameters for spiral galaxies, we attempt to predict the contribution to the sub-mm galaxy counts and background at $`850\mu m`$ by using a 1/$`\lambda `$ absorption law for the dust and re-radiating the spiral UV radiation into the FIR. We use the Bruzual & Charlot(1993) galaxy evolution models with $`H_0=50`$km$`s^1`$Mpc<sup>-1</sup> and a $`\tau `$=9 Gyr SFR - with a galaxy age of 16 Gyr in the low $`q_0`$ case, and 12.7 Gyr in the high $`q_0`$ case to produce our 1M$``$ galactic spectral energy distribution(SED) as a function of redshift. We then use the equation
$$G_{abs}(z)=F_\lambda (z)(110^{0.4A_B(4500/\lambda )})𝑑\lambda $$
(1)
as used by Metcalfe et al(1996), which is used to calculate the radiation absorbed by the dust, $`G_{abs}`$($`ergss^1`$), for our 1M$``$ model spiral galaxy as a function of z, using our $`1/\lambda `$ absorption law with A$`{}_{B}{}^{}=0.3`$. Since Bruzual & Charlot provides us with a 1M$``$ SED at each redshift increment, we need to calculate the factor required to scale this SED (after the effect of absorption from the dust) to obtain that of a galaxy with absolute magnitude $`M_B`$ at zero redshift, and this factor will then remain constant for $`M_B`$ galaxies at all other redshifts. This then provides a zero point from which to calculate scaling factors for all the other galaxies in our luminosity functions. We find the scaling factor for an $`M_B`$ galaxy by making use of a relation from Allen(1995)
$$m_B=2.5log(B_\lambda f_\lambda 𝑑\lambda )12.97$$
(2)
where $`f_\lambda `$ is the received flux($`ergs^1\AA ^1cm^2`$) and $`B_\lambda `$ is the B band filter function. By re-arranging, setting $`m_B`$=$`M_B`$ and then multiplying by $`4\pi (10pc)^2`$ we obtain the total emitted power, $`L_B`$($`ergs^1`$) in the B band from an $`M_B`$ galaxy
$$L_B=4\pi (10pc)^2.10^{[0.4(M_B+12.97)]}$$
(3)
The intensity emitted in the B band, after absorption by the dust from our 1M$``$ galaxy, $`L_{BM_{}}`$ is then calculated by integrating the SED, assuming a flat B band filter, between $`4000\AA `$ and $`5000\AA `$.
$$L_{BM_{}}=F_\lambda (z)10^{0.4A_B(4500/\lambda )}B_\lambda 𝑑\lambda $$
(4)
The scaling factor to scale a Bruzual & Charlot $`1M_{}`$spectral energy distribution for a galaxy of absolute magnitude, M<sub>B</sub>, is then defined by the ratio $`L_B/L_{BM_{}}`$.
The way the dust will re-radiate this absorbed flux depends on its temperature, particle size and chemical composition. However the normalisation of the re-radiated flux from a galaxy with absolute magnitude M<sub>B</sub>, at redshift z, is already determined (the quantity $`G_{abs}`$$`E_B/E_{BM_{}}`$). We will adopt a simple model by assuming a mean interstellar dust temperature of 15K, \[Bianchi et al. 1999\] and also a modest warmer component of 45K, (the actual luminosity ratio we use is $`L_{45K}/L_{15K}=0.162`$), which would come from circumstellar dust \[Dominigue et al. 1999\] and is needed in order to fit counts at shorter wavelengths eg. $`60\mu m`$. The effect of varying the dust parameters is explored in section 4. We then simply scale the Planck function so that
$$C(z,M_B)_{\mathrm{}}^{\mathrm{}}\beta (\lambda ,T)𝑑\lambda =G_{abs}L_B/L_{BM_{}}$$
(5)
where C(z,M<sub>B</sub>) is the scaling factor, which is a function of z and M<sub>B</sub>, $`\beta (\lambda ,T)`$ is the Planck function (in this case a sum of two Planck functions) and $`\kappa _d(\lambda )\lambda ^\beta `$, where $`\kappa _d(\lambda )`$ is an opacity law (we use $`\beta =2.0`$ for each Planck function to model optically thin dust).
We then calculate the received $`850\mu m`$ flux, S(z,M<sub>B</sub>), from a galaxy with absolute magnitude $`M_B`$ and redshift z using the equation
$$S(z,M_B)=\frac{C(z,M_B)\lambda _e^\beta \beta (\lambda _e,T)}{4\pi (1+z)d_L^2}$$
(6)
where C(z,M<sub>B</sub>) is defined from (4) and $`\lambda _e`$ is equal to $`850\mu m`$/(1+z). We can then obtain the number count of galaxies with absolute magnitude between M<sub>B</sub> and M$`{}_{B}{}^{}+dM_B`$ and redshift between z and z+dz for which we measure the same flux density S(z,M<sub>B</sub>) at $`850\mu m`$ (see (4)).
$$dN(z,M_B)=\varphi (M_B)\frac{dV}{dz}dM_Bdz$$
(7)
where $`\varphi (M_B)`$ is the optical Schechter function and $`\frac{dV}{dz}`$ is the cosmological volume element. Then the integral source counts N$`(>S_{lim})`$ are obtained, for each value of $`S_{lim}`$, by integrating (5) over the range of values of $`M_B`$ and $`z`$ such that S(z,M$`{}_{B}{}^{})>S_{lim}`$, where S(z,M<sub>B</sub>) is defined in (4).
$$N(>S_{lim})=_{M_B}_z\varphi (M_B)\frac{dV}{dz}𝑑M_B𝑑z$$
(8)
It is straightforward to then obtain model predictions of the FIR background for a given wavelength. The intensity, dI, at $`850\mu m`$ from galaxies with absolute magnitudes between $`M_B`$ and $`M_B+dM_B`$ and redshifts between z and z+dz is given by multiplying the number of galaxies with these z’s and M<sub>B</sub>’s by the flux density which we would measure from each
$$dI_{850}=S(z,M_B)\varphi (M_B)\frac{dV}{dz}dM_Bdz$$
(9)
and then we simply integrate over all absolute magnitudes and all redshifts ($`0<z<4`$ in this case)
$$I_{850}=_{M_B}_zS(z,M_B)\varphi (M_B)\frac{dV}{dz}𝑑M_B𝑑z$$
(10)
## 4 Predictions
Fig. 1 shows our model predictions for the $`60\mu m`$ differential number counts of IRAS galaxies (Saunders et al. 1990). This was an all sky local survey carried out with the IRAS satellite down to a flux limit of 0.6Jy. It therefore provides an important test of our model since spiral galaxies contribute significatly to IRAS counts (Neugebauer et al. 1984) and so if we are going to assume PLE out to redshifts of 4 then our local galaxy count predictions at $`60\mu m`$ need to be reasonably consistent with the data. The figure shows our evolution and no evolution model(the $`q_0`$ makes no difference) and because the IRAS survey was probing redshifts out to z=0.2 we can see that there is very little difference between the two models and that they both fit the data reasonably well. The IRAS counts below 0.2Jy are slightly under-predicted using both dust laws, which could possibly be due to the fact our model doesn’t include any fast-evolving AGN/ULIRG population. We use the Calzetti dust law with three dust components of 15, 25, and 32K and this failure of the fainter IRAS counts is greater than when the 1/$`\lambda `$ law is used because of the absence of the 45K dust component, which dominates the thermal emission at $`60\mu `$m..
We then go on to show in Fig. 2 our sub-mm predictions using the Bruzual & Charlot evolution model with low and high $`q_0`$ ($`q_0=0.05,q=0.5`$) and also for the corresponding no-evolution models where we use the Bruzual & Charlot SED at $`z=0`$ for all redshifts. We have used a two-component dust temperature, as described in the previous section and a galaxy formation redshift, $`z_f=4`$. The low $`q_0`$ model reproduces the faint counts well, but fails the very bright counts. This makes sense since these very luminous sources would require ULIRG’s, having SFR’s of order $``$ 100-1000Myr<sup>-1</sup>, and/or AGN, in order to produce these huge FIR luminosities. Indeed, the $`850\mu m`$ integral log N:log S appears flat between 2-10mJy before rising again at fainter fluxes, suggesting that 2 populations may be contributing to the counts.
The high $`q_0`$ model contains a dwarf elliptical population in order to fit the optical counts, as already explained, but no dust was invoked in these galaxies in the optical models and so they do not contribute to our $`850\mu m`$ predictions. Contrary to the optical number counts, the high $`q_0`$ models predict more galaxies greater than a given flux limit than low $`q_0`$ models. The reason for this is illustrated in Fig. 3, which shows how the received flux density from a M$`{}_{B}{}^{}=22.5`$ galaxy would vary with redshift in the high and low $`q_0`$ case, with and without $`\tau `$=9Gyr. Bruzual & Charlot evolution. In the no-evolution cases the two factors involved are the cosmological dimming and the effect of the negative k-correction, since we are effectively looking up the black body curve as we look out to higher redshift. The high $`q_0`$ model( dotted line) predicts greater flux densities for a given redshift than with low $`q_0`$, explaining why the integral number counts are higher for a given flux density. When the Bruzual & Charlot evolution is invoked (solid and dashed lines), we predict more flux than in the corresponding no-evolution cases at high redshift, because a galaxy is significantly brighter than at the present day. The high $`q_0`$ model(with evolution) is virtually flat in the redshift range $`0.5<z<2`$ and the low $`q_0`$ model again predicts slightly lower flux densities for a given redshift compared to high $`q_0`$. It may be noted that the no-evolution models in this plot differ slightly from that of Hughes et al.(1998). This discrepancy is a result of the different assumed dust temperature and beta parameter. The colder temperature means that the peak of the thermal emission from the dust is probed at lower redshifts and so we lose the benefit of the negative k-correction at z$``$2-3 instead of at z$``$7-9 as in Hughes & Dunlop(1998).
Fig. 4 shows the effect of altering the interstellar dust temperature (where we have used the low $`q_0`$ evolving model). The interstellar dust temperature, $`T_{int}`$ makes a big difference to our $`850\mu m`$ number count predictions and the variation is perhaps contrary to what one may expect in that the lower $`T_{int}`$ means that we expect to see more galaxies above a given flux limit S<sub>lim</sub>. This is because, as we lower the dust temperature, although the integrated energy ie the area under the Planck curve goes down, the flux density at $`850\mu m`$ goes up slightly because we are seeing the majority of radiation at much longer wavelengths. Now recall from the previous section that the normalisation of the Planck emission curve is already defined from the amount of flux absorbed by the dust and the Planck curve is simply scaled accordingly. So because the normalisation is fixed, when we lower the dust temperature, we have to scale the Planck curve up by a much larger factor and therefore find that we obtain much larger flux densities at $`850\mu m`$, explaining why our models are very sensitive to $`T_{int}`$.
We have used a galaxy formation redshift, $`z_f=4`$ which is reasonable since sub-mm sources seem to exist out to at least that, but we do in fact find that adopting $`z_f=4`$ or $`z_f=6`$ or indeed $`z_f=10`$ does not make any difference to the number counts. Fig. 3 illustrates this, since at $`z>4`$ we are observing radiation that was emitted beyond the peak of the black-body curve, and so cosmological dimming is no longer compensated for and all the curves begin to fall away very quickly explaining why increasing $`z_f`$ beyond about z=4 makes essentially no difference to the $`850\mu m`$ number counts. Of course, a higher assumed $`T_{int}`$ would extend this redshift range to beyond z=4.
Fig. 6 shows what sort of contribution we get to the extra-galactic background, simply by integrating over the number counts in each wavelength bin. The plot shows the low and high $`q_0`$ models with and without evolution, and with our standard parameters of $`T_{int}=15K`$, $`T_{circ}=45K`$, $`\beta =2.0`$ and $`z_f=4`$. All the models predict the same intensity at short wavelengths($`\lambda =60\mu `$m), as low redshift objects would dominate making the evolution and $`q_0`$ dependence less significant. The low $`q_0`$ model is able to account for all of the background at $`850\mu m`$, the high $`q_0`$ model in fact overpredicts it by about a factor of 2 and the no evolution models, although underpredicting it, are still well within an order of magnitude. Although we can fit the background at $`850\mu m`$, we noticeably fail the data between about $`100`$ and $`300\mu m`$. We find that the only way to fit these observations using our model is to use higher values of $`A_B`$ and higher dust temperatures, as this means dust is absorbing more energy from each galaxy and so the contribution to the background in the wavelength range where warmer dust emission dominates($`100\mu m<\lambda <500\mu m`$) is much greater. The solid curve in Fig. 6 shows a prediction where we have tried the Calzetti dust model which gives more overall absorption with similar amounts of reddening; this model might also be expected to fit the B optical counts. We see that its larger amount of absorbed flux allows more flexibility in terms of using more dust components. By using three dust temperature components results we obtain a better (though still not perfect) fit to Fig. 6 in the $`100\mu m<\lambda <300\mu m`$ range, while still giving fits to the IRAS $`60\mu m`$ (Fig. 1) and faint $`850\mu m`$ number counts (Fig. 2).
## 5 Discussion
We have taken a different approach from the standard way in which sub-mm flux’s are estimated using UV luminosities \[Meurer et al. 1999\]. Instead of assuming a relationship between the UV slope $`\beta `$ and the ratio $`L_{FIR}/L_{UV}`$, we proceed directly from the spiral galaxy UV luminosity functions and simply re-radiate into the FIR by assuming a simple dust law constrained from the optical counts. A direct result of this, as has already been illustrated in the previous section, is that decreasing the interstellar dust temperature actually increases the received flux density at $`850\mu m`$, firstly because the peak in the Planck emission curve moves towards longer wavelengths and secondly because (as the absorbed flux from the dust is fixed) the normalisation scaling factor goes up. The fact then that we model the dust using a dominant interstellar component of 15K, which is significantly colder than that used in models of starburst galaxies (typically 30-50K), means that we are able to show that the evolution of normal spiral galaxies like our own Milky Way, using the Bruzual model with an exponential SFR of $`\tau =9`$Gyr, could make a very significant contribution to the sub-mm number counts in the $`S_{850}<2`$mJy range. Indeed this sort of temperature for spirals has been given recent support from observations of ISO at $`200\mu m`$ \[Alton et al. 1998a\] where, for a sample of 7 spirals, a mean temperature of 20K was found, about 10K lower than previous estimates from IRAS at shorter wavelengths. They found that 90 percent of the FIR emission came from very cold dust at temperatures of 15K. Sub-mm observations of spirals (Alton et al. 1998b; Bianchi et al. 1998) and observations of dust in our own galaxy (Sodroski et al. 1994; Reach et al. 1995; Boulanger et al. 1996; Sodroski et al. 1997) also support the claims of these sorts of dust temperatures. Of course, at z=4 our assumed interstellar dust temperature of 15K is comparable to that of the microwave background.
Our models show that normal spiral galaxies (ie those that evolve into galaxies like our own Milky Way assuming the Bruzual model) fail to provide the necessary FIR flux of the most luminous sources($`>2`$mJy) and this is not surprising since the $`\tau =9`$Gyr SFR at high redshift($`z>1`$), which is consistent with the UV data, is lower than that inferred by other models which fit the sub-mm counts by a factor of about 5 or so \[Blain et al. 1998a\]. The LBG galaxies at high redshift are predicted to be evolved spirals by the Bruzual models and the dust we invoke ($`A_B`$=0.3 implies an attenuation factor at 1500Åof 2.3) is enough to make them low luminosity sub-mm sources at flux levels of around 0.5mJy. This amount of dust, though, is not enough to account for the factor of 5 discrepancy and there are several possible reasons for this.
The first is the possible additional contribution to the sub-mm counts from AGN. Modelling of the obscured QSO population has shown that they could contribute, at most, about 30% of the background at $`850\mu m`$ but they can get much closer to the bright end of the sub-mm number counts \[Gunn & Shanks 1999\]. This is shown in Fig. 2 where we also show the $`q_0`$=0.5 model of Gunn & Shanks. Although the slope of the QSO count at the faintest limits is too flat, at brighter fluxes the QSO model fits better than the faint blue galaxy model and the combination of the two gives a better fit overall.
It is also possible that the optical and sub-mm observations are sampling a completely different population of galaxies as the obscured galaxies sampled by the sub-mm observations may well just be too red or too faint to be detected in the UV at the current flux limits (Smail et al. 1999, 2000; Dey et al. 1999). That may mean that the most luminous sub-mm sources or ULIRG’s($`>10^{13}L_{}`$) are not the LBG galaxies (which the Bruzual model predicts as evolved spirals) and so then it would not be surprising if the current sub-mm and UV derived star-formation histories at high redshift were different. However, the evidence is growing that the faint blue galaxies are significant contributors to the faint sub-mm counts. Chapman et al.(1999) carried out sub-mm observations of 16 LBG’s and found, with one exception, null detections down to their flux limit of $`0.5`$mJy. But their one detection may suggest that with enough SCUBA integration time it might be possible to detect LBG’s that are particularly luminous in the FIR and indeed, while this paper was in preparation, work from Peacock et al (1999) suggests that faint blue galaxies may be detected at $`850\mu m`$ at around the 0.2mJy level. This is below the SCUBA confusion limit of $`2`$mJy (Hughes et al. 1998; Blain et al. 1998b) and highlights the problem faced by Chapman et al.(1999) in performing targetted sub-mm observations of LBG’s. The conclusions of Peacock et al.(1999) suggest that the LBG population (the faint blue galaxies in our model) contribute at least 25 percent of the background at $`850\mu m`$ and Adelberger et al.(2000) also come to similar conclusions, namely that the UV-selected galaxy population could account for all the $`850\mu `$m background and the shape of the number counts at $`850\mu `$m. However, the conclusions of Adelberger et al.(2000) are based on the fact that the SED of SMM J14011+0252 is representative of both the LBG and sub-mm population. At present, they are only assumptions, but nevertheless the conclusions of all these authors seem to suggest that ULIRG’s may not contribute to the faint sub-mm number counts and background as much as was first thought.
The spectral slope of the UV continuum and the strength of the H$`\beta `$ emission line in Lyman Break Galaxies support the fact that interstellar dust is present (Chapman et al. 1999), but the physics of galactic dust and the way it obscures the optical radiation from a source is still very poorly understood. We started by adopting a very simplistic model for the dust, treating it as a spherical screen around our model spiral galaxy. The dust might, in reality, be concentrated in the plane of the disk for spiral galaxies and may also tend to clump around massive stars. This would make the extinction law effectively grayer as suggested by observations of local starburst galaxies (Calzetti, 1997). Indeed, we have investigated the effect of the grayer Calzetti extinction law and found that it would produce a larger sub-mm count contribution due to the higher overall absorption it would imply. Metcalfe et al (2000) have also suggested that there may be evidence for evolution of the extinction law from the U-B:B-R diagram of faint blue galaxies in the Herschel Deep Field.
We have assumed pure luminosity evolution (PLE) throughout this paper. The assumption that the number density of spiral galaxies remains constant might certainly not be the case if dynamical galaxy merging is important for galaxy formation. However, as we have seen it is relatively easy to fit the sub-mm number counts with PLE models whereas it is in fact impossible to fit the counts using pure density evolution models without hugely overpredicting the background by 50 or 100 times \[Blain et al. 1998a\]. So, if existing sub-mm observations are correct then although density evolution may also occur, luminosity evolution may be dominant. It is also striking how well the PLE models do in the optical number counts and colour-magnitude diagrams and together with the fact that we observe highly luminous objects in the sub-mm out to at least $`z=3`$ , this could indicate that the biggest galaxies could have formed relatively quickly, on timescales of about $`1`$Gyr or so. If this were true, then the PLE models may be a fair approximation to the galaxy number density and evolution in the Universe out to $`z3`$ in both the optical/near-IR and FIR.
We have not taken into account early-type galaxies as no dust was invoked in these in the optical galaxy count models. In particular, we have not included any contribution from dust in the dE population which is invoked to fit the faint optical counts in the $`q_0=0.5`$ model (Metcalfe et al. 1996). If we were to include their possible contribution this would increase our $`850\mu m`$ counts predictions at the faint end since in our models both early-type and dE star formation occurs at high redshift which is the region of greatest sensitivity for the sub-mm counts. At brighter fluxes though, where, in our models, low redshift galaxies are the only possible influence, the inclusion of early-type galaxies would be negligible.
## 6 Conclusions
The aim of this paper was to investigate whether, by re-radiating the absorbed spiral galaxy UV flux into the FIR, the dust invoked in the faint blue spirals at high z from the optical galaxy count models of Metcalfe et al.(1996) could have a significant contribution to the sub-mm galaxy counts and also the FIR background at $`850\mu `$m. We have found that, using a interstellar dust temperature of $`15K`$, a modest circumstellar component of $`45K`$, a beta parameter of 2.0 and a galaxy formation redshift of $`z_f4`$ we can account for a very significant fraction of the faint $`850\mu m`$ source counts, both in the low and high $`q_0`$ cases when we invoke Bruzual & Charlot evolution (see Fig. 2). These evolutionary models give 5-10 times more contribution to the faint sub-mm counts than the corresponding no-evolution models. At brighter fluxes, we find that the SFR and dust assumed in our normal spiral model are too low to produce the FIR fluxes of the most luminous sources. In the no-evolution cases, we underpredict the number counts, even at the faint end. Our predicted redshift distribution of sub-mm selected faint blue galaxies suggests that the main contribution to the faint counts is in the range $`0.5<z<3`$, peaking at $`z1.8`$. We have shown that our model fits the $`60\mu m`$ IRAS data well, an important local test if we want to assume PLE and extrapolate our optical spiral galaxy luminosity functions out to higher redshift. With the evolution models we can easily account for 50-100% of the FIR background at $`850\mu m`$ but fail the data by nearly an order of magnitude in the $`100300\mu m`$ range. We have shown that the only way to fit these observations using this optically based model is to use assume more dust obscuration ($`A_B=0.6`$) and much warmer dust (T=30K). Effectively gray extinction laws such as that of Calzetti et al (1997) may also provide more overall absorption and hence allow more dust temperature components to allow the flexibility to fit the FIR background from 60-850 $`\mu m`$. However, the bright sub-mm counts will still require a further contribution from QSO’s or ULIRGs to complement the contribution of the faint blue galaxies at fainter fluxes. |
warning/0002/astro-ph0002492.html | ar5iv | text | # Neutrino astronomy with the MACRO detector
## 1 Motivations for neutrino astronomy
The origin of cosmic rays is still largely an open question. The cosmic ray spectrum extends up to $`10^{20}`$ eV and the nature of the mechanisms capable of explaining such high energies is still unknown. Due to magnetic fields, charged cosmic particles are deflected from their original direction, hence the information on the position of their source is lost. On the other hand, protons of energies $`10^7`$ TeV and neutral particles, such as photons and neutrinos, point back to their sources since they are not deflected by magnetic fields. However, the universe should become opaque to protons with energies above $`510^{19}`$ eV at distances of $`30`$ Mpc due to photo-pion production when they interact with the Cosmic Microwave Radiation (CMBR) GZK . Photons are currently the main observation channel of our universe and the field of gamma-ray astronomy is now well established.
The idea of using neutrinos as probes of the deep universe was introduced in the sixties. Already in those years the first calculations on the diffuse neutrino flux from interactions of cosmic rays in the Galaxy Greisen and of high energy neutrinos from the Crab Bahcall64 were performed.
Neutrinos are weakly interacting particles and are therefore much less absorbed than gamma rays, which are not only absorbed during their propagation, but can even be absorbed by the source producing them. Neutrinos can bring information on the deep interior of sources and on the far Universe. Several examples of detection of cosmic neutrinos already exist: solar neutrinos (from 0.1 MeV up to around 10 MeV), first detected by Homestake Davis , and neutrinos from SN1987A (from $`10`$ MeV up to $`50`$ MeV) detected by Kamiokande and IMB Hirata87 ; Bionta87 . Nevertheless, neutrinos of astrophysical origin with energies larger than 100 MeV have not yet been observed. This observation would open the new field of high energy neutrino astronomy complementary to gamma ray astronomy. Moreover, an important hint on the existence of neutrino astronomy would come from the detection of photons of energies above 100 TeV. Such energies, in fact, cannot be explained by electron energy loss mechanisms (e.g. synchrotron radiation, bremsstrahlung and inverse Compton scattering) because electron acceleration is limited by the intense synchrotron radiation produced in the ambient magnetic fields. Therefore, alternate acceleration mechanisms which involve neutrinos are required.
Satellites, ground based imaging Cherenkov telescopes and extensive air shower arrays are currently investigating the universe, cosmic ray sources and acceleration mechanisms using photons as probes Jackson93 . Current Space experiments typically work in the energy range up to about 30 GeV and ground-based terrestrial experiments have typical threshold energies of about 250 GeV. Ground based experiments have larger surfaces and longer time exposures than space experiments, therefore they can observe higher energies where fluxes are low. Nevertheless, they are limited at low energies by the large background due to gamma rays produced in the electromagnetic cascades induced by cosmic ray interactions in the atmosphere.
The EGRET detector on board the Compton Gamma Ray Observatory (CGRO) satellite has thus far furnished the largest amount of information on sources up to $`30`$ GeV. The recent third EGRET catalogue EGRET , covering the observations made from 1991 to 1995, contains 271 sources observed with energies greater than 100 MeV. Between them, there are 5 pulsars, one probable radio galaxy (Cen A), 66 high-confidence identifications of blazars (BL Lac objects and radio quasars), 27 lower-confidence potential blazar identifications and a large number of identified supernova remnants (SNRs), and also 170 sources not yet identified firmly with known objects.
Satellite based detectors are providing observations on $`\gamma `$-ray bursts (GRBs) capable of solving the mystery concerning their nature. The BATSE experiment BATSE on CGRO satellite has now detected more than 2500 $`\gamma `$-ray bursts and the Italian-Dutch BeppoSAX satellite Frontera is providing breakthroughs thanks to the precise measurement (the error box radius is at the level of 4’) of the position of the bursts. Ground based experiments are looking for emissions above the TeV from GRBs: recently the Milagrito detector has found a correlation with the BATSE GRB970417a with chance probability $`1.5\times 10^3`$ Milagrito .
Cherenkov telescopes at ground level such as the Whipple observatory, HEGRA, Cangaroo and University of Durham Mark 6 telescopes, have so far detected 8 sources emitting $`\gamma `$-rays well above 300 GeV: the supernova remnants Crab Lang , Vela Pulsar (at a distance of only $`500`$ pc)Cangaroo0 , SN1006 Cangaroo , the extra-galactic BL Lac objects (highly variable active galactic nuclei) Mkn 421 ($`z=0.031`$)Punch ; Hegra , Mkn 501 ($`z<0.034`$) Quinn , PKS2155-304 ($`z=0.116`$) Chadwick , and the pulsars PSR1706-44 Kifune , PSR1259-63 Sako . The Whipple group detected the first source, the Crab supernova remnant Lang . The Crab is considered now as a standard candle for high energy gamma ray astronomy due to its gamma ray steady emission.
The number of sources detected so far by ground based experiments is much smaller than the number of sources detected by EGRET. One of the possible explanations is that high energy gamma rays are absorbed: TeV $`\gamma `$-rays suffer absorption through pair production in intergalactic space on infrared light, PeV $`\gamma `$s on the CMBR and EeV $`\gamma `$s on radio-waves. This can even explain why the only BL Lac objects observed till now are also the nearest ones.
The discovery of TeV gamma rays emitted from the 8 sources quoted above shows the possibility of production from $`\pi ^o`$ decay and the possible existence of beam dump sources (see Sec. 2) producing high energy neutrinos. Nevertheless, “a few TeV” are energies not high enough to exclude a synchrotron radiation production mechanism. It could be completely excluded only for sources of energies above 100 TeV, but up to now, no source emitting at such energies has been discovered. There were some claims in the past, particularly about Cygnus X3, but they have not been confirmed Hillas95 .
A different kind of neutrino production from astrophysical sources has been suggested by Gondolo and Silk Gondolo . If cold dark matter exists in the Galactic Center, it can be accreted by the black hole which probably is there. The cold dark matter is redistributed by the black hole into a cusp, which they call “central spike”. If dark matter is made of neutral particles that can annihilate, such as the supersymmetric neutralino, the annihilation rate in the spike is strongly increased as it depends on the square of the matter density. Neutrinos can escape and produce relevant fluxes. For the neutralino, the fluxes are very high in the case of the presence of a central spike at the level of $`10^{15}÷10^{14}`$ cm<sup>-2</sup> s<sup>-1</sup> for $`m_\chi 50`$ GeV.
During the ’70s and ’80s the first generation underground detectors of surface $`100`$ m<sup>2</sup> have been measuring neutrinos. Those detectors were aimed at detecting proton decay for which atmospheric neutrinos were considered a background. Nevertheless, neutrinos were soon considered as an interesting signal themselves and results on searches for astrophysical sources of neutrinos were made. Previous results on the search for point-like sources have been published by the Kolar Gold Field experiment Adarkar and by the water Cherenkov detectors IMB IMB and Kamiokande Kamiokande . Other experiments (Baksan Baksan and AMANDA AMANDA ) have presented preliminary results at conferences. MACRO has been detecting muon neutrinos since 1989 while it was still under construction. We present here the results of the search for astrophysical neutrino sources with MACRO during the period March 1989 - September 1999.
## 2 Neutrino astronomy
Astrophysical neutrinos can be produced in the interactions of protons accelerated by compact sources with a target around the source (gas of matter or photons). This is the most plausible model for a neutrino source, the so called “beam dump model” Berezinsky85 ; Gaisser95 . The acceleration process requires the presence of a strong magnetic field with sufficient local gas to act as a beam dump. The column density of the gas in the source is assumed to be larger than the nuclear depth ($`x_N70`$ g/cm<sup>2</sup>), but smaller than the neutrino absorption depth due mainly to $`\nu N`$ interactions ($`x_\nu 310^{12}100`$ GeV$`/E_\nu `$ g/cm<sup>2</sup> and $`x_{\overline{\nu }}610^{12}100`$ GeV$`/E_{\overline{\nu }}`$ g/cm<sup>2</sup> Berezinsky85 ). The chain of reactions is:
$`\begin{array}{cc}p+N(\gamma )\pi ^0+\pi ^\pm +\mathrm{}& \\ & \gamma +\gamma \mu ^\pm +\nu _\mu (\overline{\nu }_\mu )\\ & e^++\nu _e+\overline{\nu }_\mu (e^{}+\overline{\nu }_e+\nu _\mu ).\end{array}`$ (4)
Neutral pions produce the observed photons; from the same chain it is expected the production of charged pions and kaons which can decay producing neutrinos and muons. Moreover, muons decay too. The result are neutrinos and antineutrinos of electron and muon flavors. Neglecting the photon absorption effect, which is subject to very large uncertainties, the neutrino flux have at least the same spectral shape and intensity with respect to the gamma ray flux; hence very low neutrino event rates are expected due to the small neutrino cross section. The presence of $`100`$ TeV gamma ray sources should guarantee the existence of neutrino sources, but no reliable information could be drawn on neutrino fluxes from gamma ray ones because they are subject to non negligible absorption.
Cosmic accelerators produce a power law spectrum:
$$\frac{d\varphi }{dE}E^{(\gamma +1)}$$
(5)
where $`\gamma 1+\epsilon `$, with $`\epsilon `$ a small number. The first order Fermi acceleration mechanism in strong shock waves has the attractive feature of resulting in this kind of power spectrum Longair and it predicts a spectral index $`\gamma 1`$. The primary cosmic ray spectrum is thought to be steeper than the one resulting from a cosmic accelerator because of the energy dependence of the cosmic ray diffusion out of the Galaxy, as explained in Gaisser95 .
Primary cosmic rays interact with the nuclei in the atmosphere and produce cascades from which atmospheric neutrinos of muon and electron flavor originate from the decays of pions, kaons and muons. Up to now, only atmospheric neutrinos with energies above 100 MeV have been detected by underground detectors. If all the parent mesons of atmospheric neutrinos decay the neutrino spectrum follows the spectrum of the parent particles (($`\frac{d\varphi _\nu }{dE})_{atm}E^{2.7}`$ for $`E_\nu 10`$ GeV). For higher energies, since the path length in the atmosphere is not large enough to allow the decay of all pions and kaons, interactions of mesons begin to dominate and the atmospheric neutrino spectrum becomes steeper (($`\frac{d\varphi _\nu }{dE})_{atm}E^{3.7}`$ for $`E_\nu 100`$ GeV) due to the change of the spectral index of the meson spectra. These neutrinos originating in the Earth atmosphere are a background for the search for astrophysical neutrinos which, on the other hand, are produced by cosmic rays at their acceleration sites and hence should follow the hard cosmic ray source spectrum of the form:
$$\left(\frac{d\varphi _\nu }{dE_\nu }\right)_{source}E^{(2.0÷2.5)}.$$
(6)
Thus the signal to noise ratio becomes larger at increasing energies, and above some tens of TeV the neutrinos from sources start to dominate.
The search presented here uses only the direction information of the neutrinos. Other searches could maximize the signal to noise ratio using the energy information on the detected particle and looking to the diffuse neutrino events from the whole sky. This was done by the Frejus experiment Frejus and some results have been recently presented by the Baikal collaboration Baikal . Preliminary MACRO results were presented elsewhere Corona95 and will be the subject of a future paper.
### 2.1 Candidate sources and expected rates
High energy neutrinos are expected to be emitted from a wide class of possible celestial objects which can be divided into two wide classes: galactic sources and extragalactic sources Gaisser95 .
Galactic sources are energetic systems, such as binary systems and supernova remnants, in which cosmic rays (CRs) are accelerated and interact with matter (mainly protons). The most interesting sources are SNRs, which are the most likely sources to be observed by a detector of the MACRO size. In such systems, the target is the material of the expanding shell and the accelerating mechanism is originated by the intense magnetic field of the pulsar. There are however possibilities to have neutrino emissions originate by acceleration at the supernova blast waves and therefore neutrino emission even without pulsars. The neutrino emission should be in an active time of up to a few years. Of course the disadvantage of galactic supernovae as neutrino emitters is that their rate is low (of the order of 1/30 years). According to detailed calculations made for several historical supernova Gaisser96 , the most intense source should be the supernova remnant Vela Pulsar with a rate of upward-going muons induced by neutrinos in the rock surrounding a detector of the order 0.1 ev/yr/1000 m<sup>2</sup> for $`E_\mu >`$1 GeV. Another model for young SNRs with a pulsar having high magnetic field and short period ($`5`$ ms) is suggested in Ref. Protheroe98 : for a beaming solid angle of neutrino emission of 1 sr, about 5 events/yr are expected in 1000 m<sup>2</sup> for $`E_\nu 100`$ GeV after 0.1 yr from an explosion at a distance of 10 kpc.
A different kind of galactic source is suggested in Ref. Gondolo due to WIMP annihilations in the core of the Galactic Center. The rates would be very promising, even for detectors like MACRO, being of the order of 1-20 events/yr/1000 m<sup>2</sup>.
Possible extragalactic sources are active galactic nuclei and gamma-ray bursters. For these sources the dominant mechanism for producing neutrinos is accelerated protons interacting on ambient photons. Possible alternative mechanisms are the so called Top-Down models Sigl .
Active galactic nuclei (AGNs), being among the most luminous objects in the Universe with luminosities ranging from 10<sup>42</sup> to 10<sup>48</sup> erg/s, have been recognized for a long time as promising possible sources of neutrinos. Present models assume that they consist of a central engine (massive black hole) with an accretion disk and jets Gaisser95 . Accretion onto the central black hole provides the total power. Two possible sources of high energy neutrino fluxes within AGNs have been suggested. The first is associated with the central engine and the second with the production in jets associated with several blazars (radio-loud AGNs in which the beam intersects the observer line of sight). AGNs could emit neutrinos up to $`10^{10}`$ GeV.
Even considering the highest luminosities and the presence of jets, single AGNs are difficult to detect. Jets carry about 10$`\%`$ of the AGN luminosity and AGNs may appear brighter because of the motion of the emitting matter toward the observer (for an observer looking along the jet axis $`E_{obs}=\mathrm{\Gamma }E_{jet}`$ and $`L_{obs}=\mathrm{\Gamma }^4L_{jet}`$, where $`\mathrm{\Gamma }`$ is the Lorentz factor). Expected event rates for blazars are of the order of $`10^210^1`$ /1000 m<sup>2</sup>/yr for $`\mathrm{\Gamma }=1010^2`$ and $`E_\nu >1`$ TeV Halzen98 .
Stecker et al. Stecker suggested the possibility to integrate the neutrino flux from single generic AGNs to obtain a diffuse flux from all cosmological AGNs. Various models have been suggested Stecker ; Szabo <sup>1</sup><sup>1</sup>1Most of the Szabo and Protheroe Szabo models are excluded by the Frejus limit Frejus . and the event rates in upward-going muons vary between $`10^110`$ /1000 m<sup>2</sup>/yr for $`E_\nu >1`$ TeV.
Gamma Ray Bursters are considered as promising sources of high energy neutrinos. They yield transient events originating beyond the solar system, with typical durations of $`10^2÷10^3`$ s. The BATSE BATSE experiment has now collected more than 2500 events which appear isotropically distributed. This feature suggests that they are located at cosmological distances. The recent observations by BeppoSAX of GRB970228 have allowed the precise measurement of the position which for the first time led to the identification of a fading optical counterpart Costa . Immediately after, the direct measurement of the redshift in the optical afterglow at $`z=0.835`$ for GRB970508 Metzger and other identifications of the distances of GRBs have given support to the cosmological origin hypothesis. These observations make GRBs the most luminous objects observed in our universe with emitted energies $`10^{51}`$ erg and a spectrum peaked between 100 keV $`÷`$ 1 MeV.
One of the most plausible models is the “fireball model”, which solves the compactness problem introducing a beamed relativistic motion with $`\mathrm{\Gamma }100`$ of an expanding fireball Piran .
The question of the energy of the engine of GRBs, which is strictly connected to that of beaming, is still under discussion. Evidence for beaming are a break and a steepening of the spectrum. They have been found in spectra of some bursts, e.g. GRB980519 and GRB990123; in the case of GRB990123 at $`z=1.6`$ for isotropic emission the emitted energy would be the highest ever observed ($`210^{54}`$ erg) while if there is a beam the emitted energy would be reduced to $`10^{52}`$ erg due to the Lorentz factor.
Several authors have suggested a possible correlation between Gamma Ray Bursters and emissions of high energy neutrinos Halzen ; Bahcall ; Meszaros ; Vietri produced by accelerated protons on photons. Expected rates could be up to $`10^6`$ muon induced events in a 1000 m<sup>2</sup> detector for muon energies above $``$ 30 TeV for emissions lasting $`<1`$ s Halzen . In other scenarios, such as for fireballs, rates are of the order of $`10^3`$ upward-going muons in a 1000 m<sup>2</sup>-size detector for a burst at a distance of 100 Mpc producing $`0.410^{51}`$ erg in $`10^{14}`$ eV neutrinos WB . Considering a rate of $`10^3`$ bursts per year over 4$`\pi `$ sr, averaging over burst distances and energies, $`210^2`$ upward-going muons are expected in 1000 m<sup>2</sup> per 1 yr for 4$`\pi `$ sr Bahcall . It is important to consider that the uncertainty on the Lorentz factor $`\mathrm{\Gamma }`$ produces high variations in the expected rates: the higher the $`\mathrm{\Gamma }`$, the larger the luminosity at the observer ($`L_{obs}\mathrm{\Gamma }^4L_{jet}`$), but the smaller are the rates of events because the actual photon target density in the fireball is diluted by large Lorentz factors (the fraction of total energy going into pion production in the source and hence into neutrinos varies approximately as $`\mathrm{\Gamma }^4`$ Halzen99 ).
According to Waxman and Bahcall (WB) WB an energy independent upper bound on diffuse fluxes of neutrinos with $`E_\nu 10^{14}`$ eV produced by photo-meson or p-p interactions in sources from which protons can escape can be estimated at the level of $`E_\nu ^2\varphi _\nu <2\times 10^8`$ GeV cm<sup>-2</sup> s<sup>-1</sup> sr<sup>-1</sup>. This bound relies on the flux measurement of extremely high energy cosmic rays in extensive air showers, which are assumed to be of extragalactic origin. Their limit would exclude most of the present models of neutrino production in AGNs which are commonly normalized to the extragalactic MeV-GeV gamma-ray background. Contrary to WB, Mannheim, Protheroe and Rachen MPR find an energy dependent upper limit which agrees within a factor of 2 with WB in the limited range of $`E_\nu 10^{1618}`$ eV, while at other energies the neutrino flux is mainly limited by their contributions to extragalactic gamma-ray background which is at a level of about 2 orders of magnitude higher than the WB limit.
## 3 The MACRO detector and the data selection
In the range of energies from several GeV to several TeV the neutrinos produced by astrophysical sources can be detected in underground detectors as upward-going muons produced by neutrino charged current (CC) interactions in the rock surrounding the detector. Neutrino events can be discriminated from among the background of atmospheric muons of many orders of magnitude larger ($`510^5`$ at MACRO depth) recognizing that they travel from the bottom to the top of the apparatus after having transversed the Earth. Neutrino detection is experimentally much more difficult than the gamma ray one; because of the low neutrino interaction cross section it requires very large detectors.
The MACRO detector, shown in Fig. Neutrino astronomy with the MACRO detector and described in detail in Ahlen93 is located in the Hall B of the Gran Sasso underground laboratory at a minimum rock depth of 3150 hg/cm<sup>2</sup> and an average rock depth of 3700 hg/cm<sup>2</sup>. The detector, 76.6 m long, 12 m wide and 9.3 m high, is divided longitudinally in six similar supermodules and vertically in a lower part (4.8 m high) and an upper part (4.5 m high).
The active detectors include 14 horizontal and 12 vertical planes of 3 cm wide limited streamer tubes for particle tracking, and liquid scintillation counters for fast timing. In the lower part, the eight inner planes of limited streamer tubes are separated by passive absorbers (iron and rock $`50`$ g cm<sup>-2</sup>) in order to set a minimum threshold of $`1`$ GeV for vertical muons crossing the detector. The upper part of the detector is an open volume containing electronics and other equipment. The horizontal streamer tube planes are instrumented with external 3 cm pick-up strips at an angle of $`26.5^{}`$ with respect to the streamer tube wires, providing stereo readout of the detector hits. The transit time of particles through the detector is measured by the time of flight technique (T.o.F.) using scintillation counters. The mean time at which signals are observed at the two ends of each counter is measured and the difference in the measured mean time between counters located in different planes gives the T.o.F.. The time resolution of the scintillation counter system is about 500 ps.
In order to achieve the largest reconstruction efficiency for all directions, three algorithms for muon tracking are used in this analysis. The first kind of tracking is for events with aligned hits in at least 4 horizontal planes; the second is for events with at least 2 horizontal planes in coincidence with at least 3 vertical planes; the third is for events having at least 3 vertical planes in coincidence with two scintillation counters (this tracking is useful for almost horizontal tracks).
The angular resolution depends on the wire and strip cluster widths and on the track length. The average errors on the slopes of tracks are $`0.14^{}`$ for the wires and $`0.29^{}`$ for the strips Ahlen93 . Our pointing capabilities for point-sources has been checked with the observation of the Moon shadowing effect using atmospheric down-going muons Moon98 .
The data used for the upward-going muon search belong to three running periods with different apparatus configurations: 26 events have been detected with the lower half of the first supermodule from March 1989 until November 1991 (about 1/12 of the full acceptance, livetime of 1.38 years, efficiencies included), 55 with the full lower half of the detector (about 60$`\%`$ of the full acceptance) from December 1992 until June 1993 (0.41 years, efficiencies included). Starting from April 1994 the apparatus has been running in the final configuration. From April 1994 until Sep. 1999 we have measured 1000 events with the full detector (4.41 live years, including efficiencies). We also consider events which were measured during periods when the detector acceptance was changing with time due to construction works (19 events during 1992, 0.2 yr).
The selection of upward-going muons using the T.o.F. technique has been described in detail in Ref. Ambrosio95 ; Ambrosio98 . The velocity and direction of muons is determined from the T.o.F. between at least 2 scintillation layers combined with the path length of a track reconstructed using the streamers. Taking as a reference the upper counter which measures the time $`T_1`$, the time of flight $`\mathrm{\Delta }T=T_1T_2`$ is positive if the muon travels downward und it is negative if it travels upwards. Fig. Neutrino astronomy with the MACRO detector shows the $`1/\beta =c\mathrm{\Delta }T/L`$ ($`L`$ is the track-length, and $`c`$ the speed of light) distribution for the entire data set. In this convention, muons going down through the detector have $`1/\beta 1`$, while muons going upwards have $`1/\beta 1`$. Several cuts are imposed to remove backgrounds caused by radioactivity in coincidence with muons and multiple muons. The main cut requires that the position of a muon crossing a scintillator agrees within 70 cm (140 cm for slanted tracks with $`\mathrm{cos}\theta 0.2`$) with the position along the counter determined by the more precise streamer system. Other cuts apply only to events which cross 2 scintillator planes only. These cuts tend to remove high multiplicity events because when more than one track crosses the same scintillator box the reconstructed time of the event is wrong. Events which cross more than 2 scintillator planes (about 50$`\%`$ of the total) have a more reliable time determination thanks to the possibility to evaluate the velocity of the particle from a linear fit of times as a function of the height of the scintillator counters. In this case, the only cut then is on the quality of the fit ($`\chi ^210`$).
Events in the range $`1.25<1/\beta <0.75`$ are defined to be upward-going muon events. There are 1100 events which satisfy this definition summed over all running periods. One event is shown in Fig. Neutrino astronomy with the MACRO detector. In order to maximize the acceptance for this search, we do not require a minimum amount of material be crossed by the muon track as was done to select the sample used for the neutrino oscillation analysis Ambrosio95 ; Ambrosio98 . Without this requirement we introduce some background due to large angle pions produced by down-going muons Spurio98 . We also include events with an interaction vertex inside the lower half of the detector. All of these data can be used for the point-like astrophysical source search since the benefit of a greater exposure for setting flux limits offsets the slight increase of the background and of the systematic error in the acceptance. Moreover, one can notice that for neutrino oscillation studies upward-going muons are mostly signal and background rejection is very critical, while for neutrino astronomy upward-going muons are mostly background due to atmospheric neutrinos and background rejection is less critical.
## 4 Neutrino signal in upward-going muons
Muon neutrinos are detected as upward-going muons through CC interactions:
$$\nu _\mu (\overline{\nu }_\mu )+N\mu ^{}(\mu ^+)+X.$$
(7)
The probability that a neutrino (or antineutrino) with energy $`E_\nu `$ interacts in the rock below the detector and gives rise to a muon which crosses the apparatus with energy $`E_\mu E_\mu ^{th}`$ ($`E_\mu ^{th}`$ is the energy threshold of the apparatus) is:
$$P_\nu (E_\nu ,E_{th}^\mu )=N_A_0^{E_\nu }𝑑E_\mu ^{^{}}\frac{d\sigma _\nu }{dE_\mu ^{^{}}}(E_\mu ^{^{}},E_\nu )R_{eff}(E_\mu ^{^{}},E_{th}^\mu )$$
(8)
where $`N_A`$ is Avogadro’s number. This probability is a convolution of the $`\nu `$ cross sections and of the muon effective range $`R_{eff}(E_\mu ,E_\mu ^{th})`$ described below; the computed probability is shown in Fig. Neutrino astronomy with the MACRO detector and some values are given in Tab. 1. The trend of the probability at energies $``$ 1 TeV reflects the cross-section linear rise with neutrino energy ($`\sigma _\nu E_\nu `$) and that of the muon range ($`R_{eff}E_\mu `$), while at higher energies it reflects the damping effect of the propagator ($`\sigma _\nu E_\nu ^{0.4}`$ for $`E_\nu 10^3`$ TeV) and the logarithmic rise of the muon range with its energy.
It is relevant to notice that, thanks to the recent HERA measurements Wolf , our knowledge of the high energy deep inelastic neutrino cross section has improved significantly. There is good agreement between various sets of parton functions which provide confident predictions of the cross-sections up to $`10^6`$ GeV Gandhi . For the calculation of the probability shown in Fig. Neutrino astronomy with the MACRO detector we have used the CTEQ3-DIS CTEQ parton function, available in the PDFLIB CERN library PDFLIB , which have been considered by Gandhi et al. Gandhi and in good agreement with the more recent CTEQ4-DIS.
The technique of detecting upward-going muons generated in the rock surrounding a detector has the advantage to increase the effective detector mass, which in fact is a convolution of the detector area and of the muon range in the rock. The gain increases with energy: for example, for TeV muons, the range is of the order of 1 km. The effective muon range is given by the probability that a muon with energy $`E_\mu `$ survives with energy above threshold after propagating a distance X:
$$R_{eff}(E_\mu ,E_{th}^\mu )=_0^{\mathrm{}}𝑑XP_{surv}(E_\mu ^{^{}},E_\mu ^{th},X)$$
(9)
where the integral is evaluated from the $`\mu `$ energy losses. We have used the energy loss calculation by Lohmann et al. Lohmann85 using standard rock for muon energies up to $`10^5`$ GeV. For higher energies we use the approximate formula:
$$\frac{dE_\mu }{dX}=\alpha +\beta E_\mu ,$$
(10)
where $`\alpha 2.0`$ MeV g<sup>-1</sup> cm<sup>2</sup> takes into account the continuous ionization losses and $`\beta 3.910^6`$ g<sup>-1</sup> cm<sup>2</sup> takes into account the stochastic losses due to bremsstrahlung, pair production and nuclear interactions.
The flux of neutrino induced muons detected by an apparatus for a source of declination $`\delta `$ and for a neutrino spectrum $`\mathrm{\Phi }_\nu (E_\nu )E^\gamma `$ is:
$`\mathrm{\Phi }_\mu (E_\mu ^{th},E_\nu ,\delta )=N_A{\displaystyle _{E_\mu ^{th}}^{E_\mu ^{max}}}{\displaystyle \frac{d\sigma _\nu }{dE_\mu ^{^{}}}}(E_\mu ^{^{}},E_\nu )R_{eff}(E_\mu ^{^{}},E_{th}^\mu )Area(E_\mu ^{^{}},\delta )\mathrm{\Phi }_\nu (E_\nu )𝑑E_\mu ^{^{}}.`$ (11)
The effective area of the detector $`Area(E_\mu ^{^{}},\delta )`$, averaged over 24 hours, depends on the source declination. In the low energy region, the effective area increases with increasing muon energy because not all muons are detected depending on their track length in the detector. At higher energies (in MACRO for $`E_\mu 3`$ GeV) it reaches a plateau when all muons from all directions have enough energy to be detected. At very high energies the effective area can decrease due to electromagnetic showers. As a matter of fact, the efficiency of the analysis cuts can decrease due to high track multiplicities for high energy events. Moreover the presence of showers could lead to a bad reconstruction of the neutrino induced muon with another track of the shower. From Monte Carlo studies, the MACRO average effective area begins to decrease for $`E_\mu 1`$ TeV and it is about 20$`\%`$ (42$`\%`$) lower at 10 TeV (100 TeV) with respect to 10 GeV. The average area as a function of declination for various energies is shown in Fig. Neutrino astronomy with the MACRO detector. It has been obtained using the detector simulation based on GEANT Brun87 , but modified to properly treat the stochastic muon energy losses above 10 TeV Perrone . To obtain large Monte Carlo statistics we have used “beams” of monoenergetic muons intercepting isotropically from the lower hemisphere a volume containing MACRO and more than 2 m of the surrounding rock (to evaluate the effect of electromagnetic showers induced by high energy muons). For each beam energy, we have simulated about 10<sup>5</sup> muons.
Due to the increasing value of the $`\nu `$ cross-section, at high energies neutrinos are “absorbed” by the large amount of material they cross through the Earth. Neutrino absorption in the Earth can be taken into account introducing in the integral in eq. 11 the exponential factor:
$$e^{N_A\sigma _\nu (E)X(\mathrm{cos}\theta )}$$
(12)
which depends on $`X(\mathrm{cos}\theta )`$, the quantity of matter transversed by the incident neutrino in the Earth and hence on its zenith angle. The differential number of neutrinos as a function of the neutrino energy (response curves) for a source of differential spectral index $`\gamma =2.1`$ at two different declinations with and without absorption in the Earth is shown in Fig. Neutrino astronomy with the MACRO detector. The median neutrino energy is about 15 TeV, while for the atmospheric neutrinos it is between 50-100 GeV Ambrosio98 . It is noticeable how the absorption becomes negligible for sources seen near the horizon ($`\delta 0^{}`$). In Fig. Neutrino astronomy with the MACRO detector the same response curves are shown for 3 spectral indices. In these plots, the normalization of the neutrino fluxes is arbitrary.
It is relevant to notice that if muon neutrinos oscillate into tau neutrinos, as atmospheric neutrino experiment results suggest Ambrosio98 ; SK , $`\nu _\tau `$ are subject to considerably less absorption than muon neutrinos Bottai . Tau neutrinos are subject to a regeneration effect in the Earth: $`\nu _\tau `$ interacts and the produced tau lepton immediately decays with negligible energy loss; hence from $`\tau `$ decay another $`\nu _\tau `$ is produced. This effect, more noticeable for harder spectra, has been neglected here.
The fluxes of detectable upward-going muons for sources with $`\gamma =2.1`$ and $`\delta =60^{},0^{}`$ are shown with and without absorption in Fig. Neutrino astronomy with the MACRO detector. If one assumes that the normalization of the neutrino flux is of the order of the upper limits from $`\gamma `$-ray experiments at $`100`$ TeV ($`2\times 10^{13}`$ cm<sup>-2</sup> s<sup>-1</sup> for the Galactic Center), the expected rate of neutrino induced muons varies between $`10^2÷10^3`$ ($`10^1÷10^2`$) events/yr/1000 m<sup>2</sup> for $`\gamma =2.1`$ ($`2.5`$) depending on the declination of the source. Note that the softer the source spectrum, the higher the neutrino event rates.
An important quantity in the search for celestial point sources is the effective angular spread of the detected muons with respect to the neutrino direction. We have computed the angle between the neutrino and the detected muon using a Monte Carlo simulation. We have assumed a neutrino energy flux of the form $`dN/dE_\nu =constant\times E^\gamma `$, for several neutrino spectral indices $`\gamma `$, and considered the neutrino cross-sections, the muon energy loss in the rock and the detector angular resolution. Tab. 2 shows the fraction of the events in a $`3^{}`$ search half-cone for two different spectral indices as a function of the zenith angle. With the simulations of monoenergetic muon beams on a box larger than the detector including $`2`$ m of rock, we have calculated the effective area and even checked that our intrinsic resolution does not worsen with energy due to the effect of increasing electromagnetic showers induced by stochastic energy losses of muons. Up to 100 TeV the average angle between the generated muons and the reconstructed ones is less than $`1^{}`$.
## 5 Search for point-like sources
The MACRO data sample is shown in equatorial coordinates (right ascension in hours and declination in degrees) in Fig. Neutrino astronomy with the MACRO detector. For the point-like source search using the direction information of upward-going muons, we evaluate the background due to atmospheric neutrino induced muons randomly mixing for 100 times the local angles of upward-going events with their times. The number of mixings is chosen to have a statistical error for the background about 10 times smaller than the data fluctuations. The local angles are then smeared by $`\pm 10^{}`$ in order to avoid repetitions, particularly in the declination regions where there is small acceptance. The value of $`10^{}`$ is chosen to have variations larger than the dimensions of the search cones.
For a known candidate point-like source $`S`$ the background in the search cone $`\mathrm{\Delta }\mathrm{\Omega }=\pi \omega ^2`$, with $`\omega `$ the half width of the search cone in radians, is evaluated counting the events in a declination band around the source declination $`\delta _S`$ of $`\mathrm{\Delta }\delta =\pm 5^{}`$:
$$N_{back}=\frac{N(\mathrm{\Delta }\delta )\mathrm{\Delta }\mathrm{\Omega }}{2\pi [\mathrm{sin}(\delta _S+5^{})\mathrm{sin}(\delta _S5^{})]}.$$
(13)
We have considered the case of a possible detection of an unknown source represented by an excess of events clustered inside cones of half widths $`1.5^{}`$, $`3^{}`$ and $`5^{}`$. Hence we have looked at the number of events falling inside these cones around the direction of each of the 1100 measured events. The cumulative result of this search is shown in Fig. Neutrino astronomy with the MACRO detector for the data (full circles) and the simulation of atmospheric events (solid line). We find 60 clusters of $`4`$ muons around a given muon (including the event itself), to be compared with 56.3 expected from the background of atmospheric neutrino-induced muons. The largest cluster is made of 7 events in the $`3^{}`$ half-cone and it is located around the equatorial coordinates (right ascension, declination)=($`222.5^{},72.7^{}`$). Other 2 clusters of 6 events in $`3^{}`$ are located around = ($`188.1^{},48.1^{}`$) and ($`342.5^{},74.4^{}`$), respectively. Nevertheless, they are not statistically significant.
For our search among known point-sources, we have considered several existing catalogues: the recent EGRET catalogue EGRET , a catalogue of BL Lacertae objects Padovani of which 181 fall in the visible sky of MACRO ($`90^{}\delta 50^{}`$), the list of 8 sources in the visible sky emitting photons above TeV already mentioned in Sec. 1, the Green catalogue Green of SNRs, the BATSE BATSE catalogues, 32 BeppoSAX GRBs Pian , a compilation of 29 Novae X Masetti , which are binaries with a compact object and a companion star which transfers mass into an accretion disk. Novae X are characterized by sudden increases of luminosities in the X range ($`L10^{37}10^{38}`$ erg s<sup>-1</sup> reached after 20-90 days). Among these catalogues we have selected 42 sources we consider interesting because they have the features required by the “beam dump” model. In Fig. Neutrino astronomy with the MACRO detector the distribution of the numbers of events falling in the search cones is shown for the data and the simulation for the 42 sources. We find no statistically significant excess from any of the considered sources with respect to the atmospheric neutrino background. For the 42 selected sources we find 11 sources with $`2`$ events in a search cone of $`3^{}`$ to be compared to 12.0 sources expected from the simulation.
Upper limits on muon fluxes from sources can be calculated at a given confidence level, e.g. 90$`\%`$ c.l. as:
$$\mathrm{\Phi }(90\%c.l.)=\frac{\mathrm{Upper}\mathrm{limit}(90\%c.l.)}{\mathrm{effective}\mathrm{area}\times \mathrm{livetime}},$$
(14)
where the numerator is the upper limit calculated from the number of measured events and from the number of expected background events and the denominator is the exposure of the detector which is the area of the apparatus seen by the source during the livetime. Different methods to evaluate upper limits are described in DPB . We have calculated the upper limits (the numerator in eq. 14) using the recent and well motivated unified approach by Feldman and Cousins Feldman98 . It is possible to calculate neutrino flux upper limits from muon flux upper limits because they are related (see eq. 11). The $`90\%`$ c.l. muon and neutrino flux limits are given in Tab. 3 for the 42 selected sources. These limits are valid for muon energies $`>1`$ GeV. They include the effect of the absorption of muon neutrino in their propagation through the Earth. The limits are obtained assuming a neutrino spectrum from a source with $`\gamma =2.1`$. Moreover, the effect of the decrease in efficiency at very high energies and the reduction factors for a search half-cone of $`3^{}`$ and a spectral index $`\gamma =2.1`$ (given in Tab. 2) are included. For comparison we include the best limits from previous experiments. In order to see how the limits depend on the spectral index $`\gamma `$ we report in Tab. 4 the percentage difference of the exposure as a function of declination calculated for a source with $`\gamma =2.1`$ and a source with $`\gamma =2.3`$, 2.5, 2,7, 3.7.
To evaluate the physical implications of our limits we recall that a muon flux of the order of $`0.03\times 10^{14}`$ cm<sup>-2</sup> s<sup>-1</sup> is expected from the supernova remnant Vela Pulsar Gaisser96 , which predict a yield of neutrinos at the level of about one order of magnitude lower than present limits.
We notice that there are 6 events from GX339-4 in a $`3^{}`$ with chance probability $`P=610^3`$. Considering that we have looked at 42 sources the probability to find such an excess from at least one of these sources is $`8.6\%`$ (evaluated from Fig. Neutrino astronomy with the MACRO detector).
Between the selected 42 candidate sources, Mkn 421 and Mkn 501 are particularly interesting due to the strong emissions (in the TeV region) they present. These emissions have variable intensity during time. Mkn 421 shows a strongly variable emission with peak flares during June 1995, May 1996 and April 1998 Krennrich99 . Mkn 501 had a high state emission during about 6 months in 1997, particularly intense between April and September Protheroe97 . The strongest flare in 1998 occurred on March 5. Unfortunately, the MACRO esposure for this sources is not favoured because they are seen almost at the horizon where the acceptance is lower. No event from both sources is found inside a search cone as large as $`5^{}`$. Only two events for each of the sources are found inside $`10^{}`$. They are of marginal interest due to the large angle with respect to the source directions. They have been measured in periods in which there were not known intense flares (for Mkn 421: 10 Sep. 1996 and 27 Jun. 1998; for Mkn 501: 29 Sep. 1996 and 26 Jun. 1998).
We have also made a search for neutrino signals using a cumulative analysis: for each of several catalogues of source types, we set a limit on flux from sources from that catalogue. In some situations (for example for a uniform distribution in space of sources having the same intensity) this method could give a better sensitivity than the search for a single source. It depends on the spatial distribution and on the intensity of the sources.
We consider the average value $`N_0`$ of the distribution for the data and the average value $`N_B`$ of the distribution for the simulation in Fig. Neutrino astronomy with the MACRO detector for the 42 sources in the MACRO list, in Fig. Neutrino astronomy with the MACRO detector for the 220 SNRs, in Fig. Neutrino astronomy with the MACRO detector for the 181 blazars and in analogous plots for the other catalogues. Then we estimate the cumulative upper limits for $`N`$ sources in the catalogue as:
$$\mathrm{\Phi }_{cumulative}(90\%\mathrm{c}.\mathrm{l}.)=\frac{\mathrm{Upper}\mathrm{limit}(90\%\mathrm{c}.\mathrm{l}.)}{\mathrm{Average}\mathrm{Area}\mathrm{livetime}}$$
(15)
where the average area is $`\frac{_{i=1}^NArea(\delta _i)}{N}`$, with $`Area(\delta _i)`$ the area seen by a source with declination $`\delta _i`$. The upper limit is evaluated for $`N_0>N_B`$ as:
$$\mathrm{Upper}\mathrm{limit}(90\%\mathrm{c}.\mathrm{l}.)=N_0N_B+1.28RMS/\sqrt{N}$$
(16)
where RMS is the root mean square value of the considered expected distributions and if $`N_0<N_B`$ as:
$$\mathrm{Upper}\mathrm{limit}(90\%\mathrm{c}.\mathrm{l}.)=1.28RMS/\sqrt{N}$$
(17)
We obtain $`\mathrm{\Phi }_{lim}(90\%)=3.0610^{16}`$ cm<sup>-2</sup> s<sup>-1</sup> for the 42 sources in the MACRO list. This can be considered a limit on a diffuse flux.
In the case of the 220 SNRs in the Green catalogue we obtain the cumulative upper limit from the cumulative analysis shown in Fig. Neutrino astronomy with the MACRO detector of $`2.6310^{16}`$ cm<sup>-2</sup> s<sup>-1</sup>. This can be considered a diffuse muon flux limit for neutrino production from supernova remnants. For the 181 blazars in Padovani (see Fig. Neutrino astronomy with the MACRO detector) we find a cumulative upper limit of the muon flux $`5.4410^{16}`$ cm<sup>-2</sup> s<sup>-1</sup>. In Tab. 5 we summarize the upper limits for the various catalogues considered.
Finally, it is interesting to note, that most of the models for neutralino annihilation on the Galactic Center in Gondolo are excluded by our experimental upper limit of $`310^{15}`$ cm<sup>-2</sup> s<sup>-1</sup> when there is a “central spike” for a $`3^{}`$ cone. In Tab. 6 the muon flux limits (90$`\%`$ c.l.) for various search cones around the direction of the Galactic Center ($`3^{}`$, $`5^{}`$ and $`10^{}`$) for 5 values of the neutralino mass from 60 GeV to 1 TeV are calculated. The dependence of the effective area of the detector as a function of the neutrino energy has been calculated folding with the neutrino flux from neutralino annihilation calculated by Bottino et al. Bottino . As a first approximation, the difference of using the neutrino fluxes by Silk and Gondolo in the effective area calculation should be negligible. The limits are calculated for $`E_\mu >1`$ GeV.
## 6 Search for correlations with gamma ray bursts
We look for correlations with the gamma ray bursts given in the BATSE Catalogues 3B and 4B BATSE containing 2527 gamma ray bursts from 21 Apr. 1991 to 5 Oct. 1999. They overlap in time with 1085 upward-going muons collected by MACRO during this period. The effective area for upward-going muon detection in the direction of the bursts averaged over all the bursts in the catalogue is 121 m<sup>2</sup>. Its value is small because our detector is sensitive to neutrinos only in one hemisphere and because it was not complete in the period 1991-1994. Fig. Neutrino astronomy with the MACRO detector shows our neutrino events and BATSE GRBs as a function of the year.
We find no statistically significant correlation between neutrino events and gamma burst directions for search cones of $`10^{}`$, $`5^{}`$ and $`3^{}`$ half widths. The width of the search cones is related to the BATSE angular resolution; these cones include 96.9$`\%`$, 85.1$`\%`$ and 70.5$`\%`$ neutrinos respectively if emitted from GRBs. These numbers do not include the contributions due to the muon-neutrino angle, to the muon propagation in the rock or to the MACRO angular resolution, which are small with respect to BATSE angular resolution.
We also consider possible time correlations between MACRO and BATSE events. For the temporal coincidences we use both the position information and the time information. In order to calculate the background we add 200 temporal shifts to the time difference between the event detected by other experiments and the $`\nu `$ event in MACRO considering various time intervals (the minimum interval is \[-4000 s, 4000 s\], the maximum interval is \[-80000 s, 80000 s\]). We consider time windows of $`\pm 400`$ s every 20 s.
We find one event after 39.4 s from the 4B950922 $`\gamma `$-ray burst of 22 Sep. 1995 at an angular distance of 17.6 and another very horizontal event in coincidence with the 4B940527 $`\gamma `$-ray burst of 27 May 1994 inside 280 s at $`14.9^{}`$. The $`90\%`$ c.l. muon flux limit is calculated for a search cone of $`10^{}`$ around the gamma burst direction and in an arbitrary time window of $`\pm 200`$ s. The choice of this time window is arbitrary because one does not know a priori what the duration of the neutrino emission is. Models of GRB emitters are not yet clear in predicting when and for how long neutrinos are emitted. This is in fact the reason why we have considered even a directional analysis of GRBs using no time information (see previous section). On the other hand, in this section we are using the time information and our choice of the time window where we set the upper limit is only motivated by the fact that this window is larger than the duration of 97.5$`\%`$ of the 3B Catalogue GRBs (for which the measured durations are available). In the chosen search window we find no events to be compared to 0.04 expected background events. Fig. Neutrino astronomy with the MACRO detector shows the difference in time between the detection of an upward-going muon and a GRB as a function of the cosine of their angular separation. Two scales are shown: the upper plot is an expanded scale of the lower one.
The corresponding flux upper limit (90$`\%`$ c.l.) is $`0.79\times 10^9`$ cm<sup>-2</sup> upward-going muons per average burst. The limit is almost eight orders of magnitude lower than the flux coming from an “extreme” topological defect model reported in Halzen , while according to a model in the context of the fireball scenario Bahcall a burst at a distance of 100 Mpc producing $`0.4\times 10^{51}`$ erg in neutrinos of about $`10^{14}`$ eV would produce $`6\times 10^{11}`$ cm<sup>-2</sup> upward going muons.
The same analysis on space and time correlations has been performed for 32 BeppoSAX events; the result is compatible with the atmospheric neutrino background.
## 7 Conclusions
We have investigated the possibility that the sample of 1100 upward-going muons detected by MACRO since 1989 shows evidence of a possible neutrino astrophysics source. We do not find any significant signal with respect to the statistical fluctuations of the background due to atmospheric neutrinos from any of the event directions or from any candidate sources. We also used the time information to look for correlations with gamma-ray bursts detected by BATSE and BeppoSAX. Having found no excess of events with respect to the expected background we set muon and neutrino flux upper limits for point-like sources and for the cumulative search for catalogues of sources. These limits have been calculated taking into account the response of MACRO to various neutrino fluxes from candidate sources until energies $`100`$ TeV. These limits are for almost all of the considered sources the most stringent ones compared to other current experiments. They are about 1 order of magnitude higher than values quoted by most plausible neutrino source models except for the model in Gondolo which is seriously constrained.
Acknowledgements
We gratefully acknowledge the support of the director and of the staff of the Laboratori Nazionali del Gran Sasso and the invaluable assistance of the technical staff of the Institutions participating in the experiment. We thank the Istituto Nazionale di Fisica Nucleare (INFN), the U.S. Department of Energy and the U.S. National Science Foundation for their generous support of the MACRO experiment. We thank INFN, ICTP (Trieste) and NATO for providing fellowships and grants for non Italian citizens. |
warning/0002/hep-ph0002029.html | ar5iv | text | # Is it possible to determine the S–factor of the hep process from a laboratory experiment?
## I Introduction
The reaction
$$p+{}_{}{}^{3}\mathrm{He}{}_{}{}^{4}\mathrm{He}+e^++\nu _e,$$
(1)
the so–called hep reaction, is the source of solar neutrinos with the highest energies (up to 18.8 MeV). It is expected that the cross section of this process at solar energies ($`30`$ keV) is extremely small and cannot be measured in laboratory conditions . The most detailed calculations of the cross section of the hep process were done in Refs. .
The process (1) proceeds via a Gamov–Teller transition (a change in isospin from 1 to 0 is involved). In ref. there are two reasons for the strong suppression of the hep cross section. The first one is due to the fact that the matrix element of the process vanishes in the allowed approximation if only the main $`s`$–state components of <sup>4</sup>He and <sup>3</sup>He wave functions are taken into account . The second reason lies in strong cancellations between different matrix elements of the weak nuclear current: it was found in Ref. that the contribution to the matrix element of the hep process from the diagrams with $`\pi `$ and $`\rho `$ meson exchange currents (in which the transition of a nucleon into the $`\mathrm{\Delta }`$ isobar is also taken into account) is comparable in modulus with the contribution of the one–body current, but has opposite sign. As a result there is a cancellation of the contributions of these two terms and the calculated cross section is about 5 times smaller than the cross section predicted by the one–nucleon term only.
The result of the calculation also depends on the two-body nuclear potential: by using two different, typical NN–potentials, for the S–factor of hep process in Ref. the following values were obtained:
$`S_1(hep)`$ $`=`$ $`1.44\times 10^{20}\mathrm{keV}\mathrm{b},`$ (2)
$`S_2(hep)`$ $`=`$ $`3.14\times 10^{20}\mathrm{keV}\mathrm{b}.`$ (3)
The average between these two values,
$$S_0(hep)=2.3\times 10^{20}\mathrm{keV}\mathrm{b},$$
(4)
is used in the Standard Solar Model (SSM) . If the value of the astrophysical S–factor of the hep process is given by Eq.(4) then the total flux of hep–neutrinos ,
$$\mathrm{\Phi }(hep)=2.1\times 10^3\mathrm{cm}^2\mathrm{s}^1,$$
(5)
is more than three orders of magnitude smaller than the flux of <sup>8</sup>B neutrinos,
$`\mathrm{\Phi }({}_{}{}^{8}\mathrm{B})=5.15\times \left(1.00_{0.14}^{+0.19}\right)\times 10^6\mathrm{cm}^2\mathrm{s}^1,`$
and hep neutrinos give a negligible contribution to the event rates observed in solar neutrino experiments. Let us stress, however, that the value (4) of $`S_0(hep)`$ is the result of a very complicated and model-dependent calculations.
The recent interest in hep neutrinos was triggered by the results of the Super-Kamiokande experiment in which the spectrum of recoil electrons in the solar-neutrino-induced process
$`\nu +e^{}\nu +e^{}`$
was measured. The spectrum of neutrinos from the decay $`{}_{}{}^{8}\text{B}{}_{}{}^{8}\text{Be}+e^++\nu _e`$ is determined by weak interactions and well known. The recoil electron spectrum measured in the Super-Kamiokande experiment is in agreement with the predicted spectrum in the whole energy range starting from 5.5 MeV, with the exception of the highest energy region, in which two data points with large errors (from the bins 13.5 – 14 MeV and 14 – 20 MeV) are above the prediction .
The Super-Kamiokande recoil electron spectrum can be fitted with the assumption that there is no distortion of the spectrum ($`\chi ^2=24.3`$ at 17 d.o.f.) . In Ref. attention was payed, however, to the fact that the high energy points in the Super-Kamiokande data could be due to the contribution of hep neutrinos . If one considers $`S(hep)`$ as a free parameter, then from the fit of the data (504 days of Super-Kamiokande) in Ref. in the hypothesis of no oscillations $`S(hep)/S_0(hep)=26`$ was obtained. In a more recent fit (825 days of Super-Kamiokande) it was found $`S(hep)/S_0(hep)=16`$ ($`\chi ^2=19.5`$ at 16 d.o.f.).
Let us notice that the distortion of the recoil electron spectrum could also be due to the MSW effect or to vacuum oscillations (VO). The largest enhancement of the high energy part of the spectrum is expected for the VO solution. By fitting the data with $`S(hep)=S_0(hep)`$ in the case of VO solution it was found $`\mathrm{sin}^22\theta =0.79`$, $`\mathrm{\Delta }m^2=4.3\times 10^{10}`$ eV<sup>2</sup> ($`\chi ^2=44.1`$ at 35 d.o.f.) (see also Refs. ).
The investigation of the problem of hep neutrinos in the solar neutrino Super-Kamiokande experiment will be continued. The results of a new measurement of the spectrum of solar $`\nu `$’s in the region $`E5`$ MeV will be soon available from the SNO experiment . In this experiment the spectrum of solar $`\nu _e`$’s will be determined from the measurement of the electron spectrum in the process $`\nu _e+de^{}+p+p`$.
One of the central theoretical problems connected with hep neutrinos is the astrophysical S–factor of the hep process . “The most important unsolved problem in theoretical nuclear physics related to solar neutrinos is the range of values allowed by fundamental physics for the hep production cross section”(Bahcall ). In this letter we consider the possibility to determine $`S(hep)`$ from experimental data. We will obtain here a relation between $`S(hep)`$ and the total cross section of the process
$$e^{}+{}_{}{}^{4}\mathrm{He}{}_{}{}^{3}\mathrm{H}+n+\nu _e$$
(6)
near threshold. The relation we obtain is based on the isotopic invariance of the strong interactions (we neglect the Coulomb interaction in the region of nuclear forces). From the existing nuclear data it follows that the violation of isotopic invariance for light nuclei cannot be larger than $`1020\%`$.<sup>*</sup><sup>*</sup>*This estimate follows from nuclear mass differences, mirror nuclei spectra and so on.
From the point of view of a possible investigation of the process (6) at small energies, we would like to stress two points:
1. The cross section of the process (6) does not contain the Coulomb penetration factor, which suppresses at small energies the cross sections of processes with initial particles having charges of equal sign.
2. There exist electron accelerators (microtrons) which allow to obtain high intensity electron beams in the range of energies which are appropriate for the investigation of the process (6).
## II The relation between $`S(hep)`$ and $`\sigma (e^{}{}_{}{}^{4}\mathrm{He}{}_{}{}^{3}\mathrm{H}n\nu _e)`$
Let us start by considering the process (1) at small solar energies ($`30`$ keV). In Ref. the general arguments are given that at small energies the cross sections of the reactions with charged initial particles have the form
$$\sigma (E)=\frac{1}{E}e^{2\pi \eta }S(E).$$
(7)
Here $`E`$ is the kinetic energy of the initial particles in the C.M. system and
$$\eta =\frac{Z_1Z_2e^2}{v},$$
(8)
where $`Z_1e`$, $`Z_2e`$ are the charges of the initial particles, $`v=\sqrt{2E/\mu }`$ is their relative velocity and $`\mu `$ is the reduced mass. In the expression (7),
$$Pe^{2\pi \eta }$$
(9)
is the probability of penetration of the incident particle through the Coulomb barrier and the factor $`S(E)`$ is determined mostly by strong interactions. If there are no resonances at small energies, the function $`S(E)`$ depends very weakly on the energy $`E`$.
The relation (7) with $`S\mathrm{const}.`$ allows to describe the existing low energy data and is used for the extrapolation of laboratory data to the energy region which is relevant for solar reactions (see, for example, Ref. ). Our further considerations will be based on this relation.
The standard weak interaction Hamiltonian density is given by
$$_I=\frac{G_F}{\sqrt{2}}\overline{\nu }_e\gamma ^\alpha (1\gamma _5)ej_\alpha +\mathrm{h}.\mathrm{c}.,$$
(10)
where the hadronic $`VA`$ current $`j_\alpha =j_\alpha ^1ij_\alpha ^2j_\alpha ^{1i2}`$ is the “minus” component of the isovector $`j_\alpha ^a`$ ($`a=1,2,3`$).
For the matrix element of the process (1) we have
$`f|S|i`$ $`=`$ $`i{\displaystyle \frac{G_F}{\sqrt{2}}}{\displaystyle \frac{1}{(2\pi )^3}}{\displaystyle \frac{1}{\sqrt{4k_0k_0^{}}}}\mathrm{}^\alpha (k,k^{})`$ (11)
$`\times `$ $`{\displaystyle d^4xe^{iqx}{}_{}{}^{4}\mathrm{He}|T\left(j_\alpha (x)e^{i{\scriptscriptstyle d^4y_I^0(y)}}\right)|p{}_{}{}^{3}\mathrm{He}}.`$ (12)
Here $`\mathrm{}^\alpha (k,k^{})=\overline{u}(k^{})\gamma ^\alpha (1\gamma _5)v(k)`$ is the matrix element of the weak leptonic current, $`q=k+k^{}`$ ($`k`$ and $`k^{}`$ being the momenta of the $`e^+`$ and $`\nu _e`$, respectively) and $`_I^0=_I^h+_I^{em}`$ is the Hamiltonian density of strong ($`_I^h`$) and electromagnetic ($`_I^{em}`$) interactions.
Let us first consider only that part of the matrix element of the process (1) which is determined by the strong interactions and gives the major contribution to the $`S`$–factor. Neglecting the Coulomb interaction in the region of nuclear forces, we have for the hadronic part of the matrix element (12)
$$f|S|i=i\frac{G_F}{\sqrt{2}}\frac{1}{(2\pi )^3}\frac{1}{\sqrt{4k_0k_0^{}}}\mathrm{}^\alpha (k,k^{}){}_{}{}^{4}\mathrm{He}|J_\alpha ^{()}(0)|p{}_{}{}^{3}\mathrm{He}(2\pi )^4\delta (P^{}P),$$
(13)
where $`J_\alpha ^{()}(x)J_\alpha ^{1i2}(x)`$ is the hadronic weak $`VA`$ current in the Heisenberg representation and $`P`$ ($`P^{}`$) is the total four-momentum of the initial (final) states. It is evident that $`{}_{}{}^{4}\mathrm{He}|J_\alpha ^{()}(0)|p{}_{}{}^{3}\mathrm{He}`$ includes all possible contributions coming from strong interactions (for example, in addition to the one–body nucleonic current, also two–body exchange currents, effects of the $`\mathrm{\Delta }`$ isobar and so on).
Using the charge symmetry of strong interactions we have
$`{}_{}{}^{4}\mathrm{He}|J_\alpha ^{1i2}|p{}_{}{}^{3}\mathrm{He}={}_{}{}^{4}\mathrm{He}|𝒰^1𝒰J_\alpha ^{1i2}𝒰^1𝒰|p{}_{}{}^{3}\mathrm{He}`$ (14)
$`={}_{}{}^{4}\mathrm{He}|J_\alpha ^{1+i2}|n{}_{}{}^{3}\mathrm{H}=n{}_{}{}^{3}\mathrm{H}|J_\alpha ^{1i2}|{}_{}{}^{4}\mathrm{He}^{}.`$ (15)
Here $`𝒰=\mathrm{exp}\{i\pi T_2\}`$ is the unitary operator of rotation by an angle $`\pi `$ around the second axis in isospace. In Eq.(15) we took into account that
$`𝒰J_\alpha ^{1i2}𝒰^1`$ $`=`$ $`J_\alpha ^{1+i2},`$ (16)
$`𝒰|p{}_{}{}^{3}\mathrm{He}`$ $`=`$ $`|n{}_{}{}^{3}\mathrm{H},`$ (17)
$`𝒰|{}_{}{}^{4}\mathrm{He}`$ $`=`$ $`|{}_{}{}^{4}\mathrm{He}.`$ (18)
Thus the hadronic part of the matrix elements of the process $`e^{}+{}_{}{}^{4}\mathrm{He}{}_{}{}^{3}\mathrm{H}+n+\nu _e`$ is connected with the matrix element of the hadronic weak current for the hep process by the simple charge symmetry relation (15). With the help of Eq.(15) we will obtain a relation which connects the $`S`$–factor of the hep process (1) with the cross section of the process (6).
Let us continue considering the hep process. In the region of small energies we are interested in, only the contribution of the $`s`$–wave of the initial $`p^3`$He system is relevant. Taking into account the Coulomb interaction between the initial $`p`$ and <sup>3</sup>He, for the total cross section of the hep process in the center of mass system we have
$`\sigma (hep)`$ $`={\displaystyle \frac{(2\pi )^4}{v}}{\displaystyle \frac{G_F^2}{2}}{\displaystyle \frac{1}{4}}{\displaystyle \underset{\mathrm{spins}}{}}{\displaystyle \frac{d^3k}{2k_0}\frac{d^3k^{}}{2k_0^{}}}`$ (20)
$`\times \left|\mathrm{}^\alpha {}_{}{}^{4}\mathrm{He}|J_\alpha ^{()}|p{}_{}{}^{3}\mathrm{He}\right|^2\delta (k_0+k_0^{}\mathrm{\Delta }){\displaystyle \frac{|\psi _\stackrel{}{p}^{(+)}(0)|^2}{|\psi _\stackrel{}{p}|^2}}.`$
Here $`\mathrm{\Delta }=m_p+m_{{}_{}{}^{3}\mathrm{He}}m_{{}_{}{}^{4}\mathrm{He}}=19.284`$ MeV (see, e.g., Ref. for the values of the nuclear masses), $`v=\sqrt{2E/\mu }`$ is the relative velocity, $`E`$ being the initial energy and $`\mu `$ the reduced mass of the $`p^3`$He system, $`\psi _\stackrel{}{p}^{(+)}(0)`$ is the Coulomb wave function of the initial $`p^3`$He system at $`r=0`$ and $`\psi _\stackrel{}{p}(\stackrel{}{r})=\mathrm{exp}(i\stackrel{}{p}\stackrel{}{r})/(2\pi )^{3/2}`$. Notice that in (20) we have neglected the small recoil energy of $`{}_{}{}^{4}\mathrm{He}`$ and that the factor $`1/4`$ is due to averaging over the spin states of the initial particles.
For the hep process the non-relativistic expression
$$\frac{|\psi _\stackrel{}{p}^{(+)}(0)|^2}{|\psi _\stackrel{}{p}|^2}=\frac{2\pi \eta }{e^{2\pi \eta }1}$$
(21)
holds, where
$$\eta =\frac{2e^2}{v}8.66\frac{1}{\sqrt{E[\mathrm{keV}]}}.$$
(22)
In the region of small energies $`2\pi \eta 1`$, we have, for the Coulomb factor (21),
$$\frac{|\psi _\stackrel{}{p}^{(+)}(0)|^2}{|\psi _\stackrel{}{p}|^2}2\pi \eta e^{2\pi \eta }.$$
(23)
For the small energies we are interested in, we have $`|\stackrel{}{q}|R1`$, $`R`$ being the radius of nuclear forces. Moreover, the parity of the initial and final nuclear states is the same and in the matrix element the linear term in the expansion of $`\mathrm{exp}\{i\stackrel{}{q}\stackrel{}{x}\}`$ vanishes. Thus, in Eq.(12) we can put $`e^{i\stackrel{}{q}\stackrel{}{x}}1`$ (allowed transition). In this approximation, the matrix element of the hadronic current does not depend on $`\stackrel{}{q}`$. This independence of the matrix element upon $`\stackrel{}{q}`$ is confirmed by the detailed calculations made in Refs. , in which not only the one–nucleon term but also two–body terms due to the exchange of $`\pi `$ and $`\rho `$ mesons were taken into account. Nevertheless we think that further investigation of the $`q`$–dependence of the nuclear matrix element is an important and interesting issue. Moreover, in that region the matrix element $`{}_{}{}^{4}\mathrm{He}|J_\alpha |p{}_{}{}^{3}\mathrm{H}`$ does not depend on the relative momentum $`\stackrel{}{p}`$ of the initial particles either. Since only the axial vector current contributes to the matrix element (Gamow–Teller transition) we obtain
$`{\displaystyle \frac{1}{4k_0k_0^{}}}{\displaystyle \underset{\mathrm{spins}}{}}\left|\mathrm{}^\alpha {}_{}{}^{4}\mathrm{He}|J_\alpha ^{()}|p{}_{}{}^{3}\mathrm{He}\right|^2=`$ (25)
$`{\displaystyle \frac{1}{4k_0k_0^{}}}{\displaystyle \underset{\mathrm{spins}}{}}\mathrm{}^i\mathrm{}_{}^{k}{}_{}{}^{}\delta _{ik}{\displaystyle \frac{1}{3}}\left|{}_{}{}^{4}\mathrm{He}|\stackrel{}{J}^{()}|p{}_{}{}^{3}\mathrm{He}\right|^2=2\left(1{\displaystyle \frac{1}{3}}{\displaystyle \frac{\stackrel{}{k}\stackrel{}{k}^{}}{k_0k_0^{}}}\right){\displaystyle \underset{\mathrm{spins}}{}}\left|{}_{}{}^{4}\mathrm{He}|\stackrel{}{J}^{()}|p{}_{}{}^{3}\mathrm{He}\right|^2.`$
It is obvious that the second term in the last equality does not give a contribution to the total cross section. Taking into account the Coulomb interaction of the final $`e^+`$ and <sup>4</sup>He, from Eqs.(20), (8) and (25) we obtain the total cross section of the hep process
$$\sigma (hep)=\frac{(2\pi )^7}{E}G_F^2e^2m_e^5\mu \underset{\mathrm{spins}}{}\left|{}_{}{}^{4}\mathrm{He}|\stackrel{}{J}^{()}|p{}_{}{}^{3}\mathrm{He}\right|^2f(\epsilon _0)e^{4\pi e^2/v}.$$
(26)
Here $`f(\epsilon _0)`$ is given by
$$f(\epsilon _0)=_1^{\epsilon _0}F(2,\epsilon )(\epsilon _0\epsilon )^2\sqrt{\epsilon ^21}\epsilon 𝑑\epsilon ,$$
(27)
where $`\epsilon =k_0/m_e`$, $`\epsilon _0=\mathrm{\Delta }/m_e`$ and $`F(2,\epsilon )`$ is the Fermi function (ratio of the modulus squared of the positron wave function in the Coulomb field of the final nucleus, calculated at $`r=R`$, to the modulus squared of the plane wave). Tables of the Fermi function are given in Ref. . For the hep reaction we have $`\epsilon _037.7`$. Using the approximation $`F(2,\epsilon )1`$ valid for small $`Z`$ and large positron energies, we obtain $`f(\epsilon _0)\epsilon _0^5/302.55\times 10^6`$.
Finally, from Eqs.(7) and (26), we derive the following expression for the S–factor for the hep process:
$$S(hep)=(2\pi )^7G_F^2e^2\mu m_e^5f(\epsilon _0)\underset{\mathrm{spins}}{}\left|{}_{}{}^{4}\mathrm{He}|\stackrel{}{J}^{()}|p{}_{}{}^{3}\mathrm{He}\right|^2.$$
(28)
Let us consider now the process (6). Neglecting the electromagnetic interaction of hadrons, for the matrix element of the process we get an expression similar to Eq.(13):
$$f|S|i=i\frac{G_F}{\sqrt{2}}\frac{1}{(2\pi )^3}\frac{1}{\sqrt{4k_0k_0^{}}}\mathrm{}^\alpha (k,k^{})n{}_{}{}^{3}\mathrm{H}|J_\alpha ^{()}(0)|{}_{}{}^{4}\mathrm{He}(2\pi )^4\delta (P^{}P),$$
(29)
where $`\mathrm{}^\alpha (k,k^{})=\overline{u}(k^{})\gamma ^\alpha (1\gamma _5)u(k)`$ and $`k`$ and $`k^{}`$ are the four–momenta of $`e^{}`$ and $`\nu _e`$, respectively.
The threshold for the process (6) is given by (see, e.g., Ref. for a table of nuclear masses)
$$E_{\mathrm{th}}=\frac{(m_n+m_{{}_{}{}^{3}\mathrm{H}})^2m_{{}_{}{}^{4}\mathrm{He}}^2}{2m_{{}_{}{}^{4}\mathrm{He}}}21.167\text{MeV}.$$
(30)
We will consider the process (6) at electron energies close to the threshold energy $`E_{\mathrm{th}}`$. For the total cross section we have
$`\sigma (e^{}{}_{}{}^{4}\mathrm{He}{}_{}{}^{3}\mathrm{H}n\nu _e)=`$ (31)
$`{\displaystyle \frac{1}{v_e}}(2\pi )^4{\displaystyle \frac{G_F^2}{2}}F(2,E_e/m_e){\displaystyle \frac{1}{2E_e}}{\displaystyle \frac{d^3k^{}}{2k_0^{}}d^3pd^3p^{}\frac{1}{2}\underset{\mathrm{spins}}{}\left|n{}_{}{}^{3}\mathrm{H}|J_\alpha ^{()}|{}_{}{}^{4}\mathrm{He}\mathrm{}^\alpha \right|^2\delta (P^{}P)},`$ (32)
where $`E_ek_0`$ is the electron energy, $`v_e`$ the velocity of the electron, $`p^{}`$ and $`p`$ are the total and relative momenta, respectively, of the $`n{}_{}{}^{3}\mathrm{H}`$ system and the Fermi function $`F(2,E_e/m_e)`$ takes into account the Coulomb interaction between the initial $`e^{}`$ and $`{}_{}{}^{4}\mathrm{He}`$. Within the same approximations we have used in the derivation of the relation (28), the cross section for the process (6) turns out to be
$`\sigma (e^{}{}_{}{}^{4}\mathrm{He}{}_{}{}^{3}\mathrm{H}n\nu _e)=`$ (34)
$`{\displaystyle \frac{32}{105}}(2\pi )^6G_F^2{\displaystyle \underset{\mathrm{spins}}{}}\left|n{}_{}{}^{3}\mathrm{H}|\stackrel{}{J}^{()}|{}_{}{}^{4}\mathrm{He}\right|^2(E_eE_{\mathrm{th}})^{7/2}\mu \sqrt{2\mu }F(2,E_e/m_e),`$
where $`\mu `$ is the reduced mass of the $`n{}_{}{}^{3}\mathrm{H}`$ system (704.1 MeV) which we identify numerically with that of the $`p{}_{}{}^{3}\mathrm{He}`$ system (703.3 MeV), in agreement with our assumption of isotopic invariance of the strong interactions. Using now the isotopic relation (15), which connects the hadronic parts of the S–matrix elements of the processes (1) and (6), we obtain the following relation between the total cross section of the process (6) and the astrophysical S–factor of the hep–process:
$$\sigma (e^{}{}_{}{}^{4}\mathrm{He}{}_{}{}^{3}\mathrm{H}n\nu _e)=\frac{32}{105}\frac{1}{(2\pi )e^2}\sqrt{\frac{2\mu }{m_e}}\frac{F(2,E_e/m_e)}{f(\epsilon _0)}\left(\frac{E_eE_{\mathrm{th}}}{m_e}\right)^{7/2}\frac{S(hep)}{m_e}.$$
(35)
For the case of electron energies $`E_e>20`$ MeV, we can set the Fermi function equal to one. Then we get for the cross section of the reaction (6)
$$\sigma (e^{}{}_{}{}^{4}\mathrm{He}{}_{}{}^{3}\mathrm{H}n\nu _e)0.62\times 10^{50}\mathrm{cm}^2\times \left(\frac{E_eE_{\mathrm{th}}}{m_e}\right)^{7/2}\frac{S(hep)}{S_0(hep)}.$$
(36)
The relation (35) allows to determine the astrophysical S–factor of the hep process $`p+{}_{}{}^{3}\mathrm{He}{}_{}{}^{4}\mathrm{He}+e^++\nu _e`$ directly from the measurements of the cross section of the process $`e^{}+{}_{}{}^{4}\mathrm{He}{}_{}{}^{3}\mathrm{H}+n+\nu _e`$ near threshold.Note that the kinetic energy of the $`n{}_{}{}^{3}\mathrm{H}`$ system in its C.M. system, which corresponds to the kinetic C.M. system energy of $`p{}_{}{}^{3}\mathrm{He}`$ in the hep reaction, is given by $`E(n{}_{}{}^{3}\mathrm{H})E_eE_{\mathrm{th}}k_0^{}`$ , where $`k_0^{}`$ is the final neutrino energy. If, for example, $`S(hep)=20S_0(hep)`$, a representative value among the ones indicated by the analysis of the Super–Kamiokande data, then at $`E_eE_{\mathrm{th}}10`$ MeV, for the cross section of the process (6) we get the value
$`\sigma (e^{}{}_{}{}^{4}\mathrm{He}{}_{}{}^{3}\mathrm{H}n\nu _e)4.0\times 10^{45}`$ cm$`^2.`$
Obviously, the cross section of the process (6) is small (weak interactions and small energies). However, taking into account the importance of obtaining direct experimental information on $`S(hep)`$, it is worthwhile from our point of view to consider the possibility of performing such a measurement. One can use the advantage of high intensity beams of low energy electrons from microtrons and the possibility to detect the process (6) by radiochemical or mass-spectrometric methods.
## III Conclusions
The problem of hep neutrinos is one of the most important issues in solar neutrino physics and in solar neutrino oscillations. The direct measurement of the cross section of the process (1) at solar energies does not seem to be possible with the present techniques and the calculation of the cross section of the hep process is a very complicated problem.
In this paper we have obtained a relation between the astrophysical S–factor of the hep process and the total cross section of the process (6) near the threshold ($`E_{\mathrm{th}}21.167`$ MeV). The relation is based on the isotopic invariance of strong interactions. The measurement of the cross section of this process near threshold would allow to determine the major hadronic part of $`S(hep)`$. Though the smallness of the cross section (36) precludes going so close to the threshold such that $`E_eE_{\mathrm{th}}`$ is of the order of the temperature in the solar core, measuring the process (6) at the more realistic energies $`E_eE_{\mathrm{th}}10`$ MeV might allow to extrapolate $`S(hep)`$ to smaller energies. Therefore, we believe that it is worthwhile to consider the possibility of measuring the cross section (36) at microtron facilities.
###### Acknowledgements.
We are indebted to F. von Feilitzsch, D. Harrach, A. Molinari, T. Ohlsson and M. Rho for very useful discussions and suggestions on the problems considered in this work.
The research has been supported in part by CICYT, under the Grant AEN/99-0692, by Italian MURST and INFN research funds. |
warning/0002/astro-ph0002124.html | ar5iv | text | # Fractal structures and the large scale distribution of galaxies aafootnote aIn the proceedings of the 7th Course in astrofundamental physics, Nato Advanced Study Institute, International Euroconference Erice, 5-16 December 1999
## 1 Introduction
Nowadays there is a general agreement about the fact that galactic structures are fractal up to a distance scale of $`50h^1Mpc`$ $`^{\mathrm{?},\mathrm{?}}`$ and the increasing interest about the fractal versus homogeneous distribution of galaxy in the last year $`^{\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?}}`$ has focused, mainly on the determination of the homogeneity scale $`\lambda _0`$.<sup>b</sup><sup>b</sup>bSee the web page http://pil.phys.uniroma1.it/debate.html where all these materials have been collected The main point in this discussion is that galaxy structures are fractal no matter what is the crossover scale, and this fact has never been properly appreciated. Clearly, qualitatively different implications are related to different values of $`\lambda _0`$.
* Characterization of scaling properties.
Given a distribution of points, the first main question concerns the possibility of defining a physically meaningful average density. In fractal-like systems such a quantity depends on the size of the sample, and it does not represent a reference value, as in the case of an homogeneous distribution. Basically a system cannot be homogeneous below the scale of the maximum void present in a given sample. However the complete statistical characterization of highly irregular structures is the objective of Fractal Geometry$`^\mathrm{?}`$.
The major problem from the point of view of data analysis is to use statistical methods which are able to properly characterize scale invariant distributions, and hence which are also suitable to characterize an eventual crossover to homogeneity. Our main contribution $`^{\mathrm{?},\mathrm{?},\mathrm{?}}`$, in this respect, has been to clarify that the usual statistical methods (correlation function, power spectrum, etc.) are based on the assumption of homogeneity and hence are not appropriate to test it. Instead, we have introduced and developed various statistical tools which are able to test whether a distribution is homogeneous or fractal, and to correctly characterize the scale-invariant properties. Such a discussion is clearly relevant also for the interpretation of the properties of artificial simulations. The agreement about the methods to be used for the analysis of future surveys such as the Sloan Digital Sky Survey (SDSS) and the two degrees Fields (2dF) is clearly a fundamental issue.
Then, if and only if the average density is found to be not sample-size dependent, one may study the statistical properties of the fluctuations with respect to the average density itself. In this second case one can study basically two different length scales. The first one is the homogeneity scale ($`\lambda _0`$), which defines the scale beyond which the density fluctuations become to have a small amplitude with respect to the average density ($`\delta \rho <\rho `$). The second scale is related to the typical length scale of the structures of the density fluctuations, and, according to the terminology used in statistical mechanics $`^\mathrm{?}`$, it is called correlation length $`r_c`$. Such a scale has nothing to do with the so-called ”correlation length” used in cosmology and corresponding to the scale $`\xi (r_0)=1`$$`^\mathrm{?}`$, which is instead related to $`\lambda _0`$ if such a scale exists.
* Implication of the fractal structure up to scale $`\lambda _0`$.
The fact that galactic structures are fractal, no matter what is the homogeneity scale $`\lambda _0`$, has deep implication on the interpretation of several phenomena such as the luminosity bias, the mismatch galaxy-cluster, the determination of the average density, the separation of linear and non-linear scales, etc., and on the theoretical concepts used to study such properties. We discuss in detail some of these points. We then review some of the main consequences of the power law behavior of the galaxy number density, by relating various observational quantities (e.g. $`r_0`$, $`\sigma _8`$, $`\mathrm{\Omega }`$, etc.) to the length scale $`\lambda _0`$.
We also note that the properties of dark matter are inferred from the ones of visible matter, and hence they are closely related. If now one observes different statistical properties for galaxies and clusters, this necessarily implies a change of perspective on the properties of dark matter.
* Determination of the homogeneity scale $`\lambda _0`$.
This is, clearly, a very important point which is at the basis of the understanding of galaxy structures and more generally of the cosmological problem. We distinguish here two different approaches: direct tests and indirect tests. By direct tests, we mean the determination of the conditional average density in three dimensional surveys, while with indirect tests we refer to other possible analyses, such as the interpretation of angular surveys, the number counts as a function of magnitude or of distance or, in general, the study of non-average quantities, i.e. when the fractal dimension is estimated without making an average over different observes (or volumes). While in the first case one is able to have a clear and unambiguous answer from the data, in the second one is only able to make some weaker claims about the compatibility of the data with a fractal or a homogeneous distribution. For example the papers of Wu et al$`^\mathrm{?}`$ and Nusser & Lahav$`^\mathrm{?}`$ mainly concern with compatibility arguments, rather than with direct tests. However, also in this second case, it is possible to understand some important properties of the data, and to clarify the role and the limits of some underlying assumptions which are often used without a critical perspective. We do not enter here in the details of the discussion about real data (see e.g. Joyce et al. 1999, Wu et al. 1999), however in the last section we consider separately the case (i) $`\lambda _050h^1Mpc`$, (ii) $`\lambda _0300h^1Mpc`$ and (iii) $`\lambda _01000h^1Mpc`$, briefly discussing the main theoretical consequences.
## 2 Characterization of scaling properties
In this section we describe in detail the correlation properties of a fractal distribution of points having eventually a crossover towards homogeneity (see Gabrielli & Sylos Labini, 2000 for a more exhaustive discussion of the subject), as the distribution of galaxies is thought to be (see Sylos Labini, Montuori & Pietronero, 1998 and Wu, Lahav & Rees, 1999 for two opposite views on the matter).
Let
$$n(\stackrel{}{r})=\underset{i}{}\delta (\stackrel{}{r}\stackrel{}{r}_i)$$
(1)
be the number density of points in the system (the index $`i`$ runs over all the points) and let us suppose to have an infinite system. If the presence of an object at the point $`\stackrel{}{r}_1`$ influences the probability of finding another object at $`\stackrel{}{r}_2`$, these two points are correlated. Hence there is a correlation at the scale distance $`r`$ if
$$G(r)=n(\stackrel{}{0})n(\stackrel{}{r})n^2$$
(2)
where we average over all occupied points of the system chosen as origin and on the total solid angle, supposing statistical isotropy. On the other hand, there is no correlation if
$$G(r)=n^2.$$
(3)
### 2.1 Homogeneity scale and correlation length
The proper definition of $`\lambda _0`$, the homogeneity scale, is the length scale beyond which the average density becomes to be well-defined, i.e. there is a crossover towards homogeneity with a flattening of $`G(r)`$. The length-scale $`\lambda _0`$ is related to the typical dimension of the largest voids in the system. On the other hand, the correlation length $`r_c`$ separates correlated regimes of the fluctuations with respect to the average density from uncorrelated ones, and it can be defined only if a crossover towards homogeneity is shown by the system, i.e. if $`\lambda _0`$ exists$`^\mathrm{?}`$. In other words $`r_c`$ defines the organization in geometrical structures of the fluctuations with respect to the average density. Clearly $`r_c>\lambda _0`$: only if the average density can be defined one may study the correlation length of the fluctuations from it. In the case $`\lambda _0`$ is finite and then $`n>0`$, in order to study the correlation properties of the fluctuations around the average and then the behaviour of $`r_c`$, we can introduce the correlation function
$$\xi (r)=\frac{n(0)n(r)n^2}{n^2}.$$
(4)
In the case of a fractal distribution, the average density $`n`$ in the infinite system is zero, then $`G(r)=0`$ and $`\lambda _0=\mathrm{}`$ and consequently $`\xi (r)`$ is not defined. In this case the only well defined quantity characterizing the two point correlations is the function $`\mathrm{\Gamma }(r)`$ $`^{\mathrm{?},\mathrm{?}}`$:
$$\mathrm{\Gamma }(r)=\underset{R_s\mathrm{}}{lim}\frac{n(r)n(0)_{R_s}}{n_{R_s}}$$
(5)
where $`R_s`$ is the size of the a generic finite sample of the system, $`\mathrm{}_{R_s}`$ indicates the average over all the points of the sample as origins, hence $`n_{R_s}`$ is the average density of the sample. This function measures the average density of points at a distance $`r`$ from another occupied point, and this is the reason why it is called the conditional average density $`^\mathrm{?}`$. Obviously in the case of a distribution for which $`\lambda _0`$ is finite $`\mathrm{\Gamma }(r)`$ provides the same information of $`G(r)`$, i.e. it characterizes the correlation properties for $`r<\lambda _0`$ and the crossover to homogeneity.
A very important point is represented by the kind of information about the correlation properties of the infinite system which can be extracted from the analysis of a finite sample of it. In Pietronero (1987) it is demonstrated that even in the super-correlated case of a fractal the estimate of $`\mathrm{\Gamma }(r)`$ extracted from a finite sample, is not dependent on the sample size $`R_s`$, providing a good approximation of that of the whole system. Clearly this is true a part from statistical fluctuations due to the fact that in a finite sample the average over the all possible origins is an average over a finite number of points, while in the global infinite system the average is over an infinite number of points. In fact, $`\mathrm{\Gamma }(r)`$ extracted from a sample can be written in the following way:
$$\mathrm{\Gamma }(r)=\frac{1}{N}\underset{i=1}{\overset{N}{}}\frac{1}{4\pi r^2\mathrm{\Delta }r}_r^{r+\mathrm{\Delta }r}n(\stackrel{}{r}_i+\stackrel{}{r}^{})d^3r^{},$$
(6)
where $`N`$ is the number of points in the sample, $`n(\stackrel{}{r}_i+\stackrel{}{r}^{})`$ is the number of points in the volume element $`d^3r^{}`$ around the point $`\stackrel{}{r}_i+\stackrel{}{r}^{}`$ and $`\mathrm{\Delta }r`$ is the thickness of the shell at distance $`r`$ from the point at $`\stackrel{}{r}_i`$.
Therefore, from an operative point of view, having a finite sample of points (e.g. galaxy catalogs), the first analysis to be done concerns the determination of $`\mathrm{\Gamma }(r)`$ of the sample itself. Such a measurement is necessary to distinguish between the two cases: (1) a crossover towards homogeneity in the sample shown by a flattening of $`\mathrm{\Gamma }(r)`$, and hence an estimate of $`\lambda _0<R_s`$ and $`n`$; (2) a continuation of the fractal behavior. Obviously only in the case (1), it is physically meaningful to study the correlation function $`\xi (r)`$ (Eq.4), and extract from it the length scale $`r_0`$ ($`\xi (r_0)=1`$), which is related to the intrinsic homogeneity scale $`\lambda _0`$. The functional behavior of $`\xi (r)`$ with distance gives instead information on the correlation length of the density fluctuations.
### 2.2 The case of a fractal distribution ($`R_s\lambda _0`$)
Hereafter we study the three-dimensional case, i.e. $`d=3`$, and we suppose that the sample is a sphere of radius $`R_s`$. Obviously, this choice is not a restriction.
Let us analyze the case $`R_s\lambda _0<r_c`$. This is the so called “fractal” case, and it is compatible with both the situation of $`\lambda _0`$ finite, but $`R_s\lambda _0`$ (a sample-size which is smaller than the homogeneity scale), or the situation in which $`\lambda _0\mathrm{}`$, i.e. the case of a fractal distribution at any scale.
It is simple to show$`^{\mathrm{?},\mathrm{?},\mathrm{?}}`$ that in this case (and in a spherical sample), Eq.6 becomes
$$\mathrm{\Gamma }(r)=\frac{BD}{4\pi }r^{D3}$$
(7)
with $`B=N/R_s^D`$. Note that $`B`$ is independent on the sample size: in fact, by changing $`R_s`$, $`N`$ in average scales as $`R_s^D`$. This shows the aforementioned assertion that $`\mathrm{\Gamma }(r)`$ is practically independent on the sample-size. On the other hand, it is possible to show that $`B`$ is related approximatively to the average distance between nearest neighbors points in the system $`^\mathrm{?}`$:
$$\mathrm{}\left(\frac{1}{B}\right)^{\frac{1}{D}}\mathrm{\Gamma }_e\left(1+\frac{1}{D}\right)$$
(8)
where $`\mathrm{\Gamma }_e`$ is the Euler’s gamma-function.
### 2.3 The “standard” correlation function for a fractal distribution
As already mentioned, in the fractal case ($`R_s\lambda _0`$), the sample estimate of the homogeneity scale, through the value of $`r`$ for which the sample-dependent correlation function $`\xi (r)`$ (given by Eq.9) is equal to $`1`$, is meaningless: This estimate is the so-called “correlation length” $`r_0`$ $`^\mathrm{?}`$ in the standard approach of cosmology. As we discuss below, $`r_0`$ has nothing to share with the true correlation length $`r_c`$. Let us see why $`r_0`$ is unphysical in the case $`R_s\lambda _0`$. $`\xi (r)`$ is given operatively by
$$\xi (r)=\frac{n(r)n(0)_{R_s}}{n_{R_s}^2}1=\frac{\mathrm{\Gamma }(r)}{n_{R_s}}1.$$
(9)
The basic point in the present discussion$`^\mathrm{?}`$, is that the mean density of the sample, $`n_{R_s}`$, used in the normalization of $`\xi (r)`$, is not an intrinsic quantity of the system, but it is a function of the finite size $`R_s`$ of the sample.
In fact, from Eq.7, the expression of the $`\xi (r)`$ of the sample in the case of fractal distributions is $`^\mathrm{?}`$
$$\xi (r)=\frac{D}{3}\left(\frac{r}{R_s}\right)^{D3}1.$$
(10)
From Eq.10 it follows that $`r_0`$ is a linear function of the sample size $`R_s`$
$$r_0=\left(\frac{D}{6}\right)^{\frac{1}{3D}}R_s$$
(11)
and hence it is a spurious quantity without physical meaning but it is simply related to the sample’s finite size.
We note that the amplitude of $`\mathrm{\Gamma }(r)`$ (Eq.7) is related to the lower cut-off of the fractal $`\mathrm{}`$ by Eq.8, while the amplitude of $`\xi (r)`$ is related to the upper cut-off (sample size $`R_s`$) of the distribution. This crucial difference has never been appreciated appropriately.
Finally we stress that in the standard analysis of galaxy catalogs the fractal dimension is estimated by fitting $`\xi (r)`$ with a power law, which instead, as one can see from Eq.10, it is power law only for $`rr_0`$ (or $`\xi 1`$). For larger distances there is a clear deviation from the power law behavior due to the definition of $`\xi (r)`$. Again this deviation is due to the finite size of the observational sample and does not correspond to any real change in the correlation properties. It is easy to see that, if one estimates the exponent at distances $`r\mathrm{}<r_0`$, one systematically obtains a higher value of the correlation exponent due to the break of $`\xi (r)`$ in a log-log plot. To illustrate more clearly this we compute the log derivative of Eq.10 with respect to $`\mathrm{log}(r)`$, indicating $`D3`$ with $`\gamma `$ and its estimate with $`\gamma ^{}`$:
$$\gamma ^{}=\frac{d(\mathrm{log}(\xi (r))}{d\mathrm{log}(r)}=\frac{2r_0^\gamma r^\gamma }{2r_0^\gamma r^\gamma 1}\gamma ,$$
(12)
where $`r_0`$ is defined by Eq.11. The tangent to $`\xi (r)`$ at $`r=r_0`$ has a slope $`\gamma ^{}=2\gamma `$. This explain why it has been found in galaxy and cluster catalogs that $`\gamma 2`$ by the $`\xi (r)`$ analysis $`^{\mathrm{?},\mathrm{?},\mathrm{?}}`$ instead of $`\gamma 1`$ found with the $`\mathrm{\Gamma }(r)`$ analysis $`^\mathrm{?}`$.
### 2.4 The case of a fractal distribution with a crossover to homogeneity ($`R_s\lambda _0`$)
Let us now analyze the case of a fractal with a crossover to homogeneity. In Coleman & Pietronero (1992) a very simple approximation has been used to describe such a situation which we discuss in more detail below.
By defining $`n_{R_s}=N/V`$ with $`V=4\pi R_s^3/3`$, it is simple to see$`^\mathrm{?}`$ that the behavior of $`\mathrm{\Gamma }(r)`$ in our sample is fractal (i.e. $`\mathrm{\Gamma }(r)`$ is a power law) up to a certain distance $`\lambda _0`$, and then it flattens, becoming homogeneous at scales $`R_s>r\lambda _0`$:
$$\{\begin{array}{c}\mathrm{\Gamma }(r)=\frac{DB}{4\pi }r^{D3}\text{for}\mathrm{}r\lambda _0\hfill \\ \\ \mathrm{\Gamma }(r)n_{R_s}\text{for}\lambda _0rR_s,\hfill \end{array}$$
(13)
where $`n_{R_s}`$ is the estimation of the average density in the sample of size $`R_s`$. That is, $`n_{R_s}`$ does not depend on $`r`$ if $`\lambda _0r\mathrm{}<R_s`$, apart small amplitude fluctuations. In Eq.13 the detailed approach to homogeneity depends on the specific properties of the fluctuations around the average density, i.e. it is determined by $`r_c`$. Hence, the statistical properties of the density fluctuations determine how good is the estimation of the average density throught $`n_{R_s}`$.
From the definition of the function $`\xi (r)`$ we can find$`^\mathrm{?}`$
$$\xi (r)\left(\frac{r}{\lambda _0}\right)^{D3}f\left(\frac{r}{r_c}\right).$$
(14)
Note that the amplitude of $`\xi (r)`$ is determined by the homogeneity scale $`\lambda _0`$ which has been previously extracted from $`\mathrm{\Gamma }(r)`$, and that in this approximation $`r_0\lambda _0`$ . The function $`\xi (r)`$ characterizes the correlations among the fluctuations of the distribution with respect to the average density. It is important to clarify that these fluctuations must be both positive and negative, in fact the integral over the whole sample of Eq.9 must be $`0`$. Hence the function $`f\left(\frac{r}{r_c}\right)`$ must be oscillating, and in the case $`\lambda _0<r_c,R_s`$, it should present an exponential cut-off at $`rr_c`$. We have that, when $`\xi (\lambda _0)1`$, the density fluctuations begin to become small with respect to the average density $`n`$, but if $`\lambda _0<r<r_c`$ they are still well correlated among them. Only for $`rr_c`$ the fluctuations are not correlated.
Let us now consider the case $`\lambda _0R_s<r_c`$. This situation is compatible with the following two situations: $`r_c`$ finite, but larger than $`R_s`$, and the case $`r_c\mathrm{}`$. In both cases, $`\xi (r)`$ of our sample should be a power law modulated by an oscillating function $`g(r)`$ which describes the positive and negative fluctuations with respect to the average density,
$$\xi (r)=\left(\frac{r}{\lambda _0}\right)^\gamma g(r).$$
(15)
In such a situation the (positive and negative) fluctuations from the average density are of all sizes and they do not have any intrinsic characteristic scale: this is a critical system (see Gaite et al., 1999 for a more detailed discussion). The only intrinsic scale of the system is then $`\lambda _0`$, the length-scale beyond which $`\mathrm{\Gamma }(r)`$ flattens and the fluctuations are small with respect to the average.
Let us suppose to be in the case in which $`\mathrm{\Gamma }(r)`$ flattens at a certain $`\lambda _0R_s`$. We can then evaluate the correlation function $`\xi (r)`$ of the sample via Eq.9. At this point we can clarify how to interpret the eventual cut-off shown by $`\xi (r)`$.
* If the cut-off scale is well below $`R_s`$, we can be sure that it is a good estimate of the intrinsic correlation length $`r_c`$;
* if the cut-off is at a scale $`rR_s`$ we can have two cases: it represents an “intrinsic” cut-off with $`r_cR_s`$, or it is only a finite-size effect due to the fact that from Eq.9 $`\left|\xi (R_s)\right|=0`$. In order to distinguish between these two possibilities, it is necessary to increase the sample size and to look at the behavior of the cut-off scale. If it increases proportionally to $`R_s`$, then it is a finite size effect. Otherwise if it does not change, it represents the estimate of the intrinsic correlation length $`r_c`$.
In Fig.1 we show two possible behaviours of the flattening of $`\mathrm{\Gamma }(r)`$, while in Fig.2 it is shown the corresponding $`\xi (r)`$ (we neglect for simplicity the oscillating term which must be present, and we have considered the situation $`R_s\mathrm{}`$).
### 2.5 About the amplitude of $`\xi (r)`$
We note that if $`\lambda _0R_s`$, $`\lambda _0`$ has nothing to share with questions like “which is the typical size of structures in the system?” or “up to which length-scale the system is clusterised?”. The answer to this question is strictly related to $`r_c`$ and not to $`\lambda _0`$. The length scale $`r_c`$ characterizes the distance over which two different points are correlated (clusterised)$`^\mathrm{?}`$. In fact, this property is not related to how large are the fluctuations with respect to the average ($`\lambda _0`$), but to the length extension of their correlations ($`r_c`$).
To be more specific, let us consider a fixed set of density fluctuations. They can be superimposed to different values of a uniform density background. The larger is this background the lower $`\lambda _0`$, but obviously the length scale of the correlations ($`r_c`$) among these fluctuations is not changed, i.e. they are clusterised independently of the background.
One can see$`^{\mathrm{?},\mathrm{?}}`$ that a linear amplification of $`\xi (r)`$ such that
$$\xi ^{}(r)=A\xi (r)$$
(16)
doesn’t change $`r_c`$ (which can be finite or infinite) but only $`\lambda _0`$, i.e. if $`A>1`$ we need larger subsamples to have a good estimation of $`n`$, but the characteristic length (correlation length) of the structures is not changed.
### 2.6 Homogeneity scale and the size of voids
Basically $`\lambda _0`$ is related to the maximum size of voids: the average density will be constant, at least, on scales larger than the maximum void in a given sample. Several authors have approached this problem by looking at voids distribution. For example El-Ad and Piran (1997) have shown that the SSRS2 and IRAS 1.2 Jy. redshift surveys are dominated by voids: they cover the $`50\%`$ of the volume. Moreover the two samples show very similar properties even if the IRAS voids are $`33\%`$ larger than SSRS2 ones because they are not bounded by narrow angular limits as the SSRS2 voids. The voids have a scale of at least $`40÷50h^1Mpc`$ and the largest void in the SSRS2 sample has a diameter of $`60h^1Mpc`$, i.e. comparable to the Bootes void. The problem is to understand whether such a scale has been fixed by the samples’ volume, or whether there is a tendency not to find larger voids: in this case one would have a (weaker evidence) for the homogeneity scale. In any case, we note that the homogeneity scale cannot be smaller than the scale of the largest void found in these samples and that one has to be very careful when comparing the size of the voids to the effective depth of catalogs. For example in the Las Campanas Redshift Survey, even if it is possible to extract sub-samples limited at $`500h^1Mpc`$, the volume of space investigated is not so large, as the survey is made by thin slices. In such a situation a definitive answer to the dimension of the of voids, and hence to the existence of the homogeneity scale, is rather difficult and uncertain.
### 2.7 Luminosity Bias
We would like to stress again that, even if the fractal behavior breaks at a certain scale $`\lambda _0`$, the use of $`\xi (r)`$ is in anyhow inconsistent at scales smaller than $`\lambda _0`$. We illustrate below an example of the confusion due to the use of $`\xi (r)`$ when $`r]ll\lambda _0`$.
From the use of the $`\xi (r)`$ analysis, it has been found that $`r_0`$ is different in different volume (hereafter VL) samples. In particular it has been found $`^\mathrm{?}`$ that deeper is the VL sample, larger is the value of $`r_0`$. As the deeper VL samples contain brighter galaxies, this fact has been interpreted as a real physical phenomenon, leading to the idea that more brighter galaxies are more strongly clustered than fainter ones, in view of their larger correlation amplitude: this is the so-called luminosity segregation phenomenon $`^{\mathrm{?},\mathrm{?}}`$. In other words, the fact that the giant galaxies are ”more clustered” than the dwarf ones, i.e. that they are located in the peaks of the density field, has given rise to the proposition that larger objects may correlate up to larger length scales and that the amplitude of $`\xi (r)`$ is larger for giants than for dwarfs one. The deeper VL samples contain galaxies which are in average brighter than those in the VL samples with smaller depths. As the brighter galaxies should have a larger correlation length the shift of $`r_0`$ in different samples can be related, at least partially, with the phenomenon of luminosity segregation.
As previously discussed there are two problems with such a model: (i) The amplitude of $`\xi (r)`$ in an homogeneous distribution, does not give any information about the clustering ”strength”. It is instead related to the local amplitude of the fluctuations with respect to the average density. (ii) The amplitude of $`\xi (r)`$ has a physical meaning only in the case $`\lambda _0`$ is found to be finite and smaller than the sample’s size. This is clearly not the case up, at least, to $`50h^1Mpc`$.
A natural explanation of the scaling of $`r_0`$ is then the fractal behavior of galaxy distribution, and more specifically the fact that $`r_0`$ is a fraction of the sample’s size in the fractal case. The fact that giant elliptical galaxies are located in the core of rich clusters, and other morphological properties of this kind, can be naturally related to the multifractal properties of matter distribution$`^{\mathrm{?},\mathrm{?}}`$. In such a case, bright galaxies are more strongly clustered than fainter ones in view of the fact that their fractal dimension is smaller$`^\mathrm{?}`$.
### 2.8 Power Spectrum of density fluctuations
The problems with the standard correlation analysis also show that the properties of fractal correlations have not been really appreciated. These problems are actually far more serious and fundamental than mentioned, for example, by Landy $`^\mathrm{?}`$ and the idea that they can be solved by simply taking the Fourier transform is once more a proof of the superficiality of the discussion. We have extensively shown $`^{\mathrm{?},\mathrm{?}}`$ that the power spectrum of the density fluctuations has the same kind of problems which $`\xi (r)`$ has, because it is normalized to the average density as well. The density contrast $`\delta (r)=\delta \rho (r)/\rho `$ is not a physical quantity unless the average density is demonstrated to exist. More specifically, like in the case of $`\xi (r)`$, the power spectrum (Fourier Transform of the correlation function) is affected by finite size effects at large scale: even for a fractal distribution the power spectrum has not a power law behavior but it shows a large scale (small $`k`$) cut-off which is due to the finiteness of the sample $`^\mathrm{?}`$. Hence the eventual detection of the turnover of the power spectrum, which is expected in CDM-like models to match the galaxy clustering to the anisotropies of the CMBR, must be considered a finite size effect, unless a clear determination of the average density in the same sample has been done.
Essentially all the currently elaborated models of galaxy formation $`^\mathrm{?}`$ assume large scale homogeneity and predict that the galaxy power spectrum (hereafter PS), which is the PS of the density contrast, decreases both toward small scales and toward large scales, with a turnaround somewhere in the middle, at a scale $`\lambda _f`$ that can be taken as separating “small” from “large” scales. Because of the homogeneity assumption, the PS amplitude should be independent on the survey scale, any residual variation being attributed to luminosity bias (or to the fact that the survey scale has not yet reached the homogeneity scale). However, the crucial clue to this picture, the firm determination of the scale $`\lambda _f`$, is still missing, although some surveys do indeed produce a turnaround scale around 100 $`h^1Mpc`$ $`^{\mathrm{?},\mathrm{?}}`$. Recently, the CfA2 survey analyzed by $`^\mathrm{?}`$ (hereafter PVGH) (and confirmed by SSRS2 $`^\mathrm{?}`$ \- hereafter DVGHP), showed a $`n=2`$ slope up to $`30h^1Mpc`$, a milder $`n1`$ slope up to 200 $`h^1Mpc`$, and some tentative indication of flattening on even larger scales. PVGH also find that deeper subsamples have higher power amplitude, i.e. that the amplitude scales with the sample depth.
In the following we argue that both features, bending and scaling, are a manifestation of the finiteness of the survey volume, and that they cannot be interpreted as the convergence to homogeneity, nor to a PS flattening. The systematic effect of the survey finite size is in fact to suppress power at large scale, mimicking a real flattening. Clearly, this effect occurs whenever galaxies have not a correlation scale much larger than the survey size, and it has often been studied in the context of standard scenarios $`^{\mathrm{?},\mathrm{?}}`$. We push this argument further, by showing that even a fractal distribution of matter, which never reaches homogeneity, shows a sharp flattening and then a turnaround. Such features are partially corrected, but not quite eliminated, when the correction proposed by $`^\mathrm{?}`$ is applied to the data. We show also how the amplitude of the PS depends on the survey size as long as the system shows long-range correlations.
The standard PS (SPS) measures directly the contributions of different scales to the galaxy density contrast $`\delta \rho /\rho `$. It is clear that the density contrast, and all the quantities based on it, is meaningful only when one can define a constant density, i.e. reliably identify the sample density with the average density of all the Universe. In other words in the SPS analysis one assumes that the survey volume is large enough to contain a homogeneous sample. When this is not true, and we argue that is indeed an incorrect assumption in all the cases investigated so far, a false interpretation of the results may occur, since both the shape and the amplitude of the PS (or correlation function) depend on the survey size.
Let us recall the basic notation of the PS analysis. Following Peebles $`^\mathrm{?}`$ we imagine that the Universe is periodic in a volume $`V_u`$, with $`V_u`$ much larger than the (presumed) maximum homogeneity scale. The survey volume $`VV_u`$ contains $`N`$ galaxies at positions $`\stackrel{}{r_i}`$, and the galaxy density contrast is
$$\delta (\stackrel{}{r})=\frac{n(\stackrel{}{r})}{\widehat{n}}1$$
(17)
where it is assumed that exists a well defined constant density $`\widehat{n}`$, obtained averaging over a sufficiently large scale. The density function $`n(\stackrel{}{r})`$ can be described by a sum of delta functions: $`n(\stackrel{}{r})=_{i=1}^N\delta ^{(3)}(\stackrel{}{r}\stackrel{}{r_i}).`$ Expanding the density contrast in its Fourier components we have
$$\delta _\stackrel{}{k}=\frac{1}{N}\underset{jϵV}{}e^{i\stackrel{}{k}\stackrel{}{r_j}}W(\stackrel{}{k}),$$
(18)
where
$$W(\stackrel{}{k})=\frac{1}{V}_V𝑑\stackrel{}{r}W(\stackrel{}{r})e^{i\stackrel{}{k}\stackrel{}{r}}$$
(19)
is the Fourier transform of the survey window $`W(\stackrel{}{r})`$, defined to be unity inside the survey region, and zero outside. If $`\xi (\stackrel{}{r})`$ is the correlation function of the galaxies ($`\xi (\stackrel{}{r})=<n(\stackrel{}{r})n(0)>/\widehat{n}^21`$), the true PS $`P(\stackrel{}{k})`$ is defined as the Fourier conjugate of the correlation function $`\xi (r)`$. Because of isotropy the PS can be simplified to
$$P(k)=4\pi \xi (r)\frac{\mathrm{sin}(kr)}{kr}r^2𝑑r.$$
(20)
The variance of $`\delta _\stackrel{}{k}`$ is $`^{\mathrm{?},\mathrm{?},\mathrm{?}}`$
$$<|\delta _\stackrel{}{k}|^2>=\frac{1}{N}+\frac{1}{V}\stackrel{~}{P}(\stackrel{}{k}).$$
(21)
The first term is the usual additional shot noise term while the second is the true PS convoluted with a window function which describe the geometry of the sample (PVGH)
$$\stackrel{~}{P}(\stackrel{}{k})=\frac{V}{(2\pi )^3}<|\delta _\stackrel{}{k^{}}|^2>|W(\stackrel{}{k}\stackrel{}{k^{}})|^2d^3\stackrel{}{k^{}}.$$
(22)
$$\stackrel{~}{P}(\stackrel{}{k})=𝑑\stackrel{}{k^{}}P(\stackrel{}{k^{}})F(\stackrel{}{k}\stackrel{}{k^{}}),$$
(23)
with
$$F(\stackrel{}{k}\stackrel{}{k^{}})=\frac{V}{(2\pi )^3}|W(\stackrel{}{k}\stackrel{}{k^{}})|^2.$$
(24)
We apply now this standard analysis to a fractal distribution. We recall the expression of $`\xi (r)`$ in this case is
$$\xi (r)=[(3\gamma )/3](r/R_s)^\gamma 1,$$
(25)
where $`\gamma =3D`$. A key point of our discussion is that that on scales larger that $`R_s`$ the $`\xi (r)`$ cannot be calculated without making assumptions on the distribution outside the sampling volume.
As we have already mentioned, in a fractal quantities like $`\xi (r)`$ are scale dependent: in particular both the amplitude and the shape of $`\xi (r)`$ depend the survey size. It is clear that the same kind of finite size effects are also present when computing the SPS, so that it is very dangerous to identify real physical features induced from the SPS analysis without first a firm determination of the homogeneity scale.
The SPS for a fractal distribution model described by Eq.25 inside a sphere of radius $`R_s`$ is
$$P(k)=_0^{R_s}4\pi \frac{\mathrm{sin}(kr)}{kr}\left[\frac{3\gamma }{3}\left(\frac{r}{R_s}\right)^\gamma 1\right]r^2𝑑r=\frac{a_k(R_s)R_s^{3D}}{k^D}\frac{b_k(R_s)}{k^3}.$$
(26)
Notice that the integral has to be evaluated inside $`R_s`$ because we want to compare $`P(k)`$ with its estimate in a finite size spherical survey of scale $`R_s`$. In the general case, we must deconvolve the window contribution from $`P(k)`$; $`R_s`$ is then a characteristic window scale. Eq.26 shows the two scale-dependent features of the PS. First, the amplitude of the PS depends on the sample depth. Secondly, the shape of the PS is characterized by two scaling regimes: the first one, at high wavenumbers, is related to the fractal dimension of the distribution in real space, while the second one arises only because of the finiteness of the sample. In the case of $`D=2`$ in Eq.26 one has:
$$a_k(R_s)=\frac{4\pi }{3}(2+\mathrm{cos}(kR_s))$$
(27)
and
$$b_k(R_s)=4\pi \mathrm{sin}(kR_s).$$
(28)
The PS is then a power-law with exponent $`2`$ at high wavenumbers, it flattens at low wavenumbers and reaches a maximum at $`k4.3/R_s`$, i.e. at a scale $`\lambda 1.45R_s`$. The scale at which the transition occurs is thus related to the sample depth. In a real survey, things are complicated by the window function, so that the flattening (and the turnaround) scale can only be determined numerically.
In practice one has several complications. First, the survey in general is not spherical. This introduces a coupling with the survey window which is not easy to model analytically. For instance, we found that windows of small angular opening shift to smaller scales the PS turnaround. This is analogous to what happens with the correlation function of a fractal: when it is calculated in small angle surveys, the correlation length $`r_0`$ decreases. Second, the observations are in redshift space, rather than in real space. The peculiar velocities generally make steeper the PS slope $`^\mathrm{?}`$ with respect to the real space. Third, in a fractal the intrinsically high level of fluctuations makes hard a precise comparison with the theory when the fractal under study is composed of a relatively small number of points.
## 3 Implications for cosmology
We now consider some implications for cosmology of the scaling properties of galaxy distribution, up to a lenght scale $`\lambda _0`$. For example, we consider more specifically $`\lambda _050h^1Mpc`$.
### 3.1 Estimation of the average luminosity and mass density
From the studies of Large Scale Structures (LSS) of galaxies and galaxy clusters, one would like to estimate the average density of visible matter and then to infer the one of the whole (visible plus dark, i.e. all the matter in clusterised objects) matter distribution. While for the first we have direct estimations, for the second we have only indirect methods, especially at large scales, based on some assumptions, which can be tested by looking at the distribution of what is observable, i.e. visible matter.
We briefly describe how to do such a measurement in galaxy redshift catalogs directly, from the knowledge of galaxy positions and luminosity (for a more detailed discussion see Sylos Labini, 2000). In this case, and by measuring the Mass-to-Luminosity ratio, one can infer, from the average luminosity density, the average mass density. The new point we address more specifically is that galaxies are fractally distributed up to a certain crossover scale $`\lambda _0`$. As there is still some controversy about the value of $`\lambda _0`$ $`^{\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?}}`$ we give the estimation as a function of $`\lambda _0`$. We stress that the way this estimation is performed, is substantially different from the usual one $`^\mathrm{?}`$, because in such a case the fractal behavior is not considered at all, and one assumes a perfect homogeneous distribution at relatively small scale ($`\lambda _05÷10h^1Mpc`$). This situation is clearly not the one corresponding to the more ”optimistic” estimation of the homogeneity scale $`\lambda _0`$ $`^{\mathrm{?},\mathrm{?},\mathrm{?}}`$. The other assumption usually made is that galaxy positions are independent on their luminosity. We have shown$`^\mathrm{?}`$ that, although such an assumption cannot describe local morphological properties of galaxy distribution$`^\mathrm{?}`$, it works rather well in the available galaxy redshift surveys.
The estimation of the the average density we are able to make depends hence on two parameters. The first one is the homogeneity scale $`\lambda _0`$ and the second is the Mass-to-Luminosity ratio. We can give an upper limit to $`\mathrm{\Omega }`$ by taking the highest $`(/)_c=300h`$ observed up to now $`^\mathrm{?}`$ (in clusters of galaxies) to be universal across all the scales, and by considering a lower limit for the homogeneity scale $`\lambda _0=50h^1Mpc`$. We compute the critical $`(/)_{crit}`$, i.e. the Mass-to-Luminosity ratio needed to have $`\mathrm{\Omega }=1`$. As the others, also this parameter depends on the homogeneity scale $`\lambda _0`$.
### 3.2 Average luminosity density from galaxy catalogs
Let
$$\nu (r,L)dLd^3r=\varphi (L)\mathrm{\Gamma }(r)dLd^3r=Ar^{D3}L^\alpha e^{\frac{L}{L_{}}}d^3rdL$$
(29)
be the average number of galaxies in the volume element $`d^3r`$ at distance $`r`$ from a observer located on a galaxy, and with luminosity in the range $`[L,L+dL]`$. In Eq.29 we have used the fact that the galaxy luminosity function has been observed to have the so-called Schechter shape with parameters $`L_{}`$ (luminosity cut-off) and $`\alpha `$ (power law index) which can be determined experimentally. The conditional average space density $`\mathrm{\Gamma }(r)`$ has a power law behavior corresponding to a fractal dimension $`D`$ (which eventually can be a function of scale, and hence can approach to $`D=3`$ at a scale $`\lambda _0`$). Both the fractal dimension $`D`$ and the overall amplitude $`A`$ can be determined in redshift surveys. Hence $`\nu (r,L)`$ is a function of four parameters: $`L_{},\alpha ,D,A`$. Moreover we note that by writing $`\nu (r,L)`$ as a product of the space density and of the luminosity function we have implicitly assumed that galaxy positions are independent on galaxy luminosity.
We would like to estimate the average luminosity density in a sphere of radius $`R`$ and volume $`V(R)`$ placed around a galaxy, and defined as
$$j(<R)=\frac{1}{V(R)}_0^R_0^{\mathrm{}}L\nu (r,L)𝑑Ld^3rj(10)\left(\frac{R}{10h^1Mpc}\right)^{D3}$$
(30)
which is $`R`$ dependent as long as the space density shows power law behavior (i.e. $`D<3`$). By considering $`M_{}=19.53`$ (i.e. $`L_{}=1.010^{10}h^2L_{}`$), $`\alpha =1.05`$ $`^\mathrm{?}`$ and by estimating the prefactor $`A`$ (Eq.29) and the fractal dimension ($`D2`$) in galaxy redshift samples we obtain$`^\mathrm{?}`$
$$j(10)210^8hL_{}/Mpc^3.$$
(31)
We now estimate the density parameter in terms of the critical density,
$$\rho _c=2.7810^{11}h^2M_{}/Mpc^3$$
(32)
where $`M_{}`$ is the solar mass. By considering the product of the mass-to-luminosity ratio (in solar and $`h`$ units) and the average luminosity density given by Eqs.30-31, we obtain
$$\mathrm{\Omega }(\lambda _0)=(6\pm 2)10^4\frac{\frac{}{}}{h}\left(\frac{\lambda _0}{10h^1}\right)^1,$$
(33)
where $`\lambda _0`$ is the scale where the crossover to homogeneity occurs (it can also be $`\lambda _0=\mathrm{}`$, and in such a case $`\mathrm{\Omega }(\mathrm{})=0`$). Note that in view of the dependence of $`/`$ on $`h`$, Eq.33 does not depend on the Hubble’s constant.
Let us now suppose that $`/10h`$ as it has been derived by Faber and Gallengher $`^\mathrm{?}`$. We obtain
$$\mathrm{\Omega }(\lambda _0)610^3\left(\frac{\lambda _0}{10h^1}\right)^1.$$
(34)
If galaxy distribution turns out to be homogeneous at $`\lambda _010h^1Mpc`$ then $`\mathrm{\Omega }610^3`$ as it is obtained in the standard treatment $`^\mathrm{?}`$. If, instead, the crossover to homogeneity lies at $`100h^1Mpc`$, we obtain $`\mathrm{\Omega }610^4`$.
From Eq.33 and if galaxy distribution turns out to be homogeneous at $`\lambda _0`$, we obtain that the critical Mass-to-Luminosity ratio (such that $`\mathrm{\Omega }(\lambda _0)=1`$) is given by
$$\left(\frac{}{}\right)_{crit}1600h\left(\frac{\lambda _0}{10h^1}\right),$$
(35)
so that if $`\lambda _0=10h^1Mpc`$ one obtains $`(/)_{crit}1600h`$ (which is again consistent with the usual adopted value $`^\mathrm{?}`$) while if $`\lambda _0=100h^1Mpc`$ $`(/)_{crit}16000h`$, which is about two orders of magnitude larger than the highest $`/`$ observed in clusterised objects. For an intermediate value of $`\lambda _0=50h^1Mpc`$ one obtains $`(/)_{crit}8000h`$.
For what concerns the analysis of galaxy clusters it is often used a value $`(/)_c300h`$ $`^{\mathrm{?},\mathrm{?}}`$ which, by using Eq.33, gives
$$\mathrm{\Omega }(\lambda _0)210^1\left(\frac{\lambda _0}{10h^1}\right)^1.$$
(36)
Such an estimation is based on the fact that the Mass-to-Luminosity ratio found in clusters is representative of all the field galaxies. This is a very strong assumption: this means that $`(/)_g=(/)_c`$ which is not supported by any observation. The usually adopted value of $`\mathrm{\Omega }=0.2`$ $`^\mathrm{?}`$ can be derived from Eq.36 by assuming $`\lambda _0=10h^1Mpc`$. In the case the crossover to homogeneity occurs at $`\lambda _0=100h^1Mpc`$ we have that $`\mathrm{\Omega }(\lambda _0)=0.02`$ if we consider $`(/)_c300h`$ to be ”universal”.
Under the assumptions:
(i) galaxies are homogeneously distributed at scales larger than $`\lambda _05h^1Mpc`$,
(ii) the $`/`$ of galaxies is the same of clusters, that is galaxies should contain a factor $`30`$ more dark matter than what it is observed with the study of the rotation curves $`^\mathrm{?}`$,
(iii) Galaxy positions are independent of galaxy luminosity: such a assumption is not strictly valid, but it has been tested$`^\mathrm{?}`$ to hold rather well in the available samples,
we get the following upper limit to $`\mathrm{\Omega }`$. If we assume that $`/=300h`$ across all the scales, and that $`\lambda _050h^1Mpc`$, which we consider to be a lower limit for the homogeneity scale, we get from Eq.36
$$\mathrm{\Omega }(50Mpc/h,300hM_{}/L_{})0.04.$$
(37)
The direct estimation with galaxies (Eqs.33-34) gives lower a value, the reason being the strong assumption of taking $`/=300h`$ as representative of field galaxies (see Fig.3).
### 3.3 Homogeneity scale and primordial nucleosynthesis constraints
Let us now discuss the previous results in relation with the nucleosynthesis constrain $`\mathrm{\Omega }_b^{BBN}=0.015h^2`$, deriving in such a way an upper limit for $`\lambda _0`$ which is consistent with such a scenario. We may see$`^\mathrm{?}`$ that if $`\lambda _0100h^2Mpc`$ then $`\mathrm{\Omega }(\lambda _0,(/)_c)\mathrm{\Omega }_b^{BBN}`$. Such a fact has two important implications.
First of all one does not need the presence of non baryonic dark matter to reconcile local observations of matter contained in galaxies and galaxy clusters ($`\mathrm{\Omega }_{local}`$) with the primordial nucleosynthesis constraint. This is an important point as the existence of non baryonic dark matter has been inferred also, but not only, by the discrepancy $`\mathrm{\Omega }_{local}/\mathrm{\Omega }_b^{BBN}`$ $`^\mathrm{?}`$, which, we show, is not observed if $`\lambda _0100h^2Mpc`$. There is still place, in this picture, for a uniform background of non-baryonic dark matter, which has completely different clustering properties from the ones of visible matter.
The second point concerns the baryon fraction in clusters. According to the usual root, in order to be consistent with $`\mathrm{\Omega }_b^{BBN}`$, one has to force $`\mathrm{\Omega }=0.1÷0.3`$ from clusters analysis$`^\mathrm{?}`$. Here we can argue as follows. By assuming a very large $`/300h`$, and that it is universal across all scales, and by assuming that all dark matter is made of baryons, if the crossover to homogeneity occurs at scales larger than $`150h^2Mpc`$ then there is not enough baryonic matter to satisfy the nucleosynthesis constraint. The situation gets clearly worst if one takes into account that not all the mass in clusters is baryonic or that $`/<300h`$, i.e. it is not considered to be universal across all scales, or that $`\lambda _0100h^2Mpc`$. There are then two different possibilities: the first is to study the case of a non homogeneous primordial nucleosynthesis which could lower the limit on $`\mathrm{\Omega }_b^{BBN}`$, given the observed abundances of light elements, and the second would be to have a more uniform background of baryonic dark matter. This seems rather unlikely because it would have very different clustering properties with respect to visible mass and it is very difficult to find a dynamical explanations for such a segregation.
### 3.4 Where does linear approximation hold ?
In the standard picture, the properties of dark matter on cosmological relevant scales $`r>5h^1Mpc`$, are inferred from the observed properties of visible matter and radiation. Now one studies change in these properties, i.e. the presence of fractal correlations, and in this respect they will have consequences on dark matter too$`^{\mathrm{?},\mathrm{?}}`$. For example, the determination of the mass density including dark matter has been performed on the basis of the linear theory $`^\mathrm{?}`$. Here the problem is: beyond which scale can linear theory be considered as a useful approximation ?
In other words, the dynamical estimates use gravitational effects, of departure from a strictly homogeneous distribution on the motion of objects such as galaxies considered as a test particle. A completely different situation occurs if $`\lambda _0`$ is larger than the scales at which linear approximation is usually adopted. For example, methods based on the Cosmic Virial Theorem $`^\mathrm{?}`$, distortions in redshift surveys $`^\mathrm{?}`$, local group dynamics $`^\mathrm{?}`$, the ”reconstructed” peculiar velocity field from the density field (e.g POTENT-like methods$`^\mathrm{?}`$), are clearly not useful at scales $`r\lambda _0`$. Up to now it is implicitly assumed $`\lambda _05÷10h^1Mpc`$ and all these methods are considered to be valid on larger scales. If, for example $`\lambda _050h^1Mpc`$ then it is not possible to interpret peculiar velocities in the range $`1÷30h^1Mpc`$ through the linear approximation (as it is usually done$`^\mathrm{?}`$), unless there is a background of dark matter which is uniform beyond a certain small scale. In such a situation however the estimates would be different $`^\mathrm{?}`$ from the usual ones $`^\mathrm{?}`$.
Let us review some simple relations between the conditional average density $`\mathrm{\Gamma }(r)`$ and the usual correlation function $`\xi (r)`$. If $`\mathrm{\Gamma }(r)`$ has a power law behavior up to a scale $`\lambda _0`$ and then it presents a crossover to homogeneity, it is simple to show that (in the case of $`D=2`$) the length scale at which $`\xi (r_0)=1`$ is of the order of
$$r_0\frac{\lambda _0}{3}$$
(38)
where the exact relation between $`r_0`$ and $`\lambda _0`$ depends on the details of the crossover $`^\mathrm{?}`$. However Eq.38 gives a reasonable order of magnitude in the general case. In such a situation it is also simple to compute $`\sigma (r,\lambda _0)=N(r)N/N`$, i.e the amplitude of the fluctuation with respect to the average at the scale $`r\lambda _0`$ $`^\mathrm{?}`$
$$\sigma (r,\lambda _0)\frac{\lambda _0}{2r}.$$
(39)
It has been found in various nearby surveys that $`\sigma (8h^1Mpc)=1`$. However if the crossover to homogeneity occurs at $`\lambda _0`$ we have that
$$\sigma (8h^1Mpc,\lambda _0)\frac{\lambda _0}{16}$$
(40)
For example if $`\lambda _015h^1Mpc`$ then $`r_05h^1Mpc`$ and $`\sigma (8h^1Mpc,15h^1Mpc)=1`$; otherwise if $`\lambda _050h^1Mpc`$ then $`r_016h^1Mpc`$ and $`\sigma (8h^1Mpc,15h^1Mpc)=3`$.
Linear approximation holds in the linear regime, when the amplitude of the density contrast is small, i.e. for $`\sigma 1`$. We have that $`\sigma =1`$ at a scale $`r_{\sigma =1}(\lambda _0)`$
$$r_{\sigma =1}(\lambda _0)\frac{\lambda _0}{2}$$
(41)
For instance, if the crossover to homogeneity occurs at $`\lambda _0=50h^1Mpc`$ one has that $`\sigma (\lambda _0)=1`$ at $`r_{\sigma =1}=25h^1Mpc`$. In such a situation all the estimations of the density parameter at smaller scales based on the linear approximation$`^\mathrm{?}`$, and hence based on the untested assumption of linearity, are not correct.
We have studied in detail$`^\mathrm{?}`$ the gravitational force distribution in a fractal structure. Its behaviour can be understood as the sum of two parts, a local or ‘nearest neighbours’ piece due to the smallest cluster (characterised by the lower cut-off $`\mathrm{\Lambda }`$ in the fractal) and a component coming from the mass in other clusters. The latter is bounded above by the scalar sum of the forces
$$|\stackrel{}{F}|\underset{L\mathrm{}}{lim}_\mathrm{\Lambda }^L\frac{G\rho _m(r)}{r^2}4\pi r^2𝑑rL^{D2}$$
(42)
so that for $`D<2`$ it is convergent, while for $`D>2`$ it may diverge. If there is a divergence, it is due to the presence of angular fluctuations at large scales, described by the three-point correlation properties of the fractal. For the difference in the force between two points the local contribution will be irrelevant well beyond the scale $`\mathrm{\Lambda }`$, while it is easy to see that the ‘far-away’ contribution will now converge as $`L^{D3}`$, and its being non-zero is a result of the absence of perfect spherical symmetry. We have then applied such a result to the case of an open universe$`^\mathrm{?}`$ in order to compute the expected deviations from a pure linear Hubble flow.
## 4 Discussion and Conclusions
We now present a short discussion about the perspective of our work: the cosmological implications of the fractal behavior of visible matter crucially depend on the crossover scale $`\lambda _0`$, but, no matter what is the actual value of such a scale, we have some important consequences from the theoretical point of view. We may identify three different scenarios.
* (1) The fractal extends only up to $`30÷100h^1Mpc`$. This is the minimal concept which begins to be absorbed in the literature$`^\mathrm{?}`$ but, sometimes, without considering its real consequences. The standard approach to galaxy distribution has identified very small ”correlation lengths”, namely $`5h^1Mpc`$ for galaxies and $`25h^1Mpc`$ for clusters. These numbers (which were supposed to know with high precision$`^\mathrm{?}`$) are anyhow inconsistent with the fractal extending a factor $`2÷4`$ more. We have shown that this inconsistency is conceptual and not due to incomplete data or week statistics. Hence, in this hypothesis one has to abandon all the concepts related to these length scales. These are:
(i) The estimate of the matter density in clusterised objects (visible + dark), which has been claimed to be $`\mathrm{\Omega }0.2÷0.3`$, decreases by one order of magnitude or more.
(ii) The normalization of N-body simulations is usually performed to some length-scale or amplitude of fluctuations, which are related to $`5h^1Mpc`$ and $`25h^1Mpc`$.
(iii) Concepts like the galaxy-cluster mismatch and the related luminosity bias, as well as the understanding of the clustering via the bias parameter $`b`$ (i.e. linear or non-linear - ”stochastic bias” - amplification of $`\xi (r)`$) loose any physical meaning.
(iv) The interpretation of the velocity field is also based on the linear approximation which cannot certainly hold at scales smaller than $`30÷50h^1Mpc`$.
(v) The reconstruction of the three dimensional properties from the angular data suffers of the same untested assumption of homogeneity.
In summary major modifications are necessary for the origin and dynamics of large scale structures and for the role of dark matter. However the structures may still formed via gravitational instability, in the sense that they are not necessarily primeval$`^\mathrm{?}`$.
* (2) The fractal extends up to $`300÷500h^1Mpc`$. In this case the standard picture of gravitationally induced structures after the electromagnetic decoupling is untenable. There is no time to create such large scale correlated structures via gravitational instability, starting from Gaussian initial conditions. More string consequences are clearly important for what concerns the amount of matter in clusterised objects.
* (3) The fractal extends up to $`\mathrm{}>1000h^1Mpc`$, and homogeneity does not exist, at least for what concerns galaxies. In this extreme case a new picture for the global metric$`^\mathrm{?}`$ is then necessary.
For some questions the fractal structure leads to a radically new perspective and this is hard to accept. But it is based on the best data and analyses available. It is neither a conjecture nor a model, it is a fact. The theoretical problem is that there is no dynamical theory to explain how such a fractal Universe could have arisen from the pretty smooth initial state we know existed in the big bang. However this is a different question. The fact that something can be hard to explain theoretically has nothing to do with whether it is true or not. Facing a hard problem is far more interesting than hiding it under the rug by an inconsistent procedure. For example some interesting attempts to understand why gravitational clustering generates scale-invariant structures have been recently proposed by de Vega et al$`^{\mathrm{?},\mathrm{?},\mathrm{?}}`$. Indeed this will be the key point to understand in the future, but first we should agree on how these new 3d data should be analyzed. In addition, the eventual crossover to homogeneity has also to be found with our approach. If for example homogeneity would really be found say at $`100h^1Mpc`$, then clearly all our criticism to the previous methods and results still holds fully. In summary the standard method cannot be used neither to disprove homogeneity, nor to prove it. One has simply to change methods.
## Acknowledgements
We warmly thank R. Durrer, A. Gabrielli, M. Joyce and M. Montuori with whom various parts of this work have been done in fruitful collaborations. FSL is grateful to Pedro Ferreira for very useful comments and discussions. We thank Prof. N. Sanchez for the organization of this very interesting School. This work has been partially supported by the EC TMR Network ”Fractal structures and self-organization” ERBFMRXCT980183 and by the Swiss NSF.
## References |
warning/0002/astro-ph0002166.html | ar5iv | text | # A spatially resolved plerionic X-ray nebula around PSR B0540-69
## 1. Introduction
The X-ray-bright 50 ms pulsar PSR B0540$``$69 was discovered in the N158A nebula of the Large Magellanic Cloud (LMC) and has long been compared to the Crab pulsar (Seward, Harnden, & Helfand 1984). Based on its timing and spectral properties, the two rotation powered pulsars are very similar with a spin period of 50 vs. 33 ms and a spin-down rate of $`4.8\times 10^{13}`$ vs. $`4.0\times 10^{13}`$ s/s for PSR B0540$``$69 and the Crab, respectively. From these quantities, assuming the standard magnetic dipole pulsar model, one can infer a characteristic age (1.7 vs. 1.2 kyr), spin-down energy ($`2.0\times 10^{38}`$ vs. $`6.5\times 10^{38}`$ ergs s<sup>-1</sup>), and surface magnetic field strength ($`5.0\times 10^{12}`$ vs. $`3.8\times 10^{12}`$ G) for PSR B0540$``$69 and the Crab, respectively. This similarity suggests that PSR B0540$``$69 could be accompanied by a “plerion”, a pulsar driven wind nebula (Weiler & Sramek 1988), reminiscent of that seen for the Crab.
Indeed, there are several lines of evidence indicating the presence of a plerion in the vicinity of PSR B0540$``$69. Chanan, Helfand, & Reynolds (1984) detected a polarized optical nebula of half-power diameter $`4^{\prime \prime }`$ around the pulsar. This apparent synchrotron nebula was also resolved (though barely) in a radio image presented by Manchester et al. (1993) and in a ROSAT high resolution imager observation by Seward & Harnden (1994). Furthermore, the overall X-ray spectrum of PSR B0540$``$69 and its remnant (SNR B0540$``$69) is well characterized by a power law, as expected if the emission is predominantly non-thermal (Clark et al. 1982; Wang & Gotthelf 1998a).
The ROSAT observation also revealed a faint X-ray emitting shell, $`15`$ pc in size surrounding the pulsar. This shell contributes less than $`20\%`$ to the total luminosity of $`1.0\times 10^{37}\mathrm{ergs}\mathrm{s}^1`$ in the ROSAT 0.1-2 keV band and is likely the SNR associated with the pulsar (Seward & Harden 1994). However, no evidence has yet been found for a similar X-ray-emitting shell or a shell-like SNR around the Crab (e.g. see discussion in Jones et al. 1998).
In this Letter, we report new results on SNR B0540$``$69 based on a recent observation acquired with the Chandra High Resolution Camera. This observation enables us for the first time to distinguish morphological details of the nebula around PSR B0540$``$69. We analyze phase dependent images and resolve the expected plerion-like nebula from the point-like pulsar emission. This allows us to identify features similar to those seen from the Crab nebula; we present morphological evidence for a torus of X-ray emission, which most likely represents shocked pulsar wind materials and a likely X-ray jet emanating from the pulsar. We discuss the implications of the results in the context of the pulsar-nebula system. Throughout the paper we adopt a distance of $`51`$ kpc for the LMC.
## 2. Observation
The Chandra observatory (Weisskopf et al. 1996) observed SNR B0540$``$69 on Aug 31 1999 as part of the initial calibration of the High Resolution Camera (HRC; Murray et al. 1997). A total of 19.4 ksec of data were collected during a portion of the orbit which avoided regions of high background contamination such as from bright Earth and radiation belt passages. The remnant was centered on the on-axis position of the HRC where the point-spread function (PSF) has a half power radius (the radius enclosing 50% of total source counts) of $`0\stackrel{}{\mathrm{.}}5`$. Time-tagged photons were acquired with 15.6 $`\mu `$s precision, and the arrival times were corrected to the solar system barycenter using a beta version of AXBARY available from the Chandra X-ray Center FTP site (A. Rots 1999, private communication). While the detector is sensitive to X-rays over a $`0.110.0`$ keV range, there is essentially no energy information available. We analyzed event data calibrated by the initial processing and dated 1999 September 12, which was made available through the Chandra public archive. In addition to the standard processing, the event data was further filtered to reduce the instrumental background and to remove “ghost image” artifacts using a beta version of HRC\_SCREEN (S. Murray 2000, private communication). We extracted $`1024\times 1024`$ pixel images centered on the pulsar rebinned by a factor of two from the native HRC pixel size of 0$`\stackrel{}{\mathrm{.}}`$13175 per side.
## 3. Results
A global view of N158A and its environment as seen by the Chandra HRC is shown in Figure $`1`$ (contours) and Figure $`2`$ (greyscale). The large-scale X-ray enhancement on scales up to $`1^{}`$ is the previously resolved shell-like emission (Seward & Harnden 1994). In fact, the X-ray and radio emission together outlines a nearly circular morphology around PSR B0540$``$69. Clearly, the shell represents the blastwave of SNR B0540$``$69. The X-ray intensity distribution within the remnant appears rather patchy. While the southwest X-ray enhancement is a good tracer of the radio and optical emission peaks, there is no general correlation between fainter radio and X-ray features.
The superb spatial resolution of the Chandra observation further allows for a close-up of the immediate vicinity of PSR B0540$``$69 (Fig. $`2b`$). The presence of a diffuse plerion-like nebula around the pulsar is apparent. To decompose the nebula emission from the pulsar contribution, we conducted phase-resolved image analysis. This enables us to estimate the local PSF based on the pulsed, point-like emission from the pulsar and to quantify the extended, unpulsed nebula radiation.
First, we must determine the pulse period at the current epoch. We constructed a periodigram around a narrow range of periods centered on the expected period $`\pm 0.1`$ ms, sampled in increments of $`0.05(P^2/T)`$, were $`T`$ is the observation duration, and $`P`$ is the test period. For each trial period, we folded photons extracted from a 1$`\stackrel{}{\mathrm{.}}`$0 aperture centered on the bootstrapped pulsar position (see below) into 20 phase bins and computed the $`\chi ^2`$ of the resultant profile. We find a highly significant signal ($`>56\sigma `$) at $`P=50.508132(6)`$ ms at Epoch 51421.630741 MJD; the uncertainty is estimated according to the method of Leahy (1987). We have assumed a period derivative of $`\dot{P}=4.789342\times 10^{13}`$ for the data epoch (Deeter et al. 1999). In Figure $`3`$ we display the resultant light-curve folded at the peak period which, as expected, is roughly sinusoidal and modulated with a $`40\%`$ pulse fraction (defined as the amplitude divided by the mean).
Next, we defined two regions in the phase space, on- and off-pulse, by selecting eight adjacent phase bins corresponding to the peak and trough of the pulse profile. The on-pulse image with the off-pulse image subtracted is presented in Figure $`4b`$. This image reproduced the expected PSF with no evidence of asymmetric deviations, as might be caused by poor aspect reconstruction, like that typically found for ROSAT images. Figure $`5`$ presents average radial intensity distributions around the centroid of the point-like source.
By subtracting the normalized pulsar image (Fig. $`4b`$), we are able to construct an image of the nebula emission (Fig. $`4c`$) alone. The subtracted image is scaled to compensated for both the relative phase coverage and for a 21% unpulsed emission contribution, estimated by minimize a pointlike contribution at the pulsar position of the nebula image. As shown in Figure $`4c`$, the extended emission is distinctly different from the pointlike image of the pulsed emission from the pulsar. The primarily feature is the NE-SW elongated feature, which is morphologically symmetric relative to the pulsar and extends about $`2\stackrel{}{\mathrm{.}}5`$ on both sides of the pulsar. However, the observed emission on the SW side appears twice as strong compared to the NE side, with an average intensity of $`7.5\times 10^2\mathrm{counts}\mathrm{s}^1\mathrm{arcsec}^2`$. Because the central core of the distribution is significantly brighter than the extended features and the subtraction of the pulsar contribution is somewhat arbitrary, the exact intensity distribution is uncertain.
There is also marginal evidence for a jet-like feature emanating from the pulsar. This emission, most apparent in the NW and extending about $`3^{\prime \prime }`$, is nearly perpendicular to the NE-SW elongated nebula and is slightly bent toward North. The integrated emission of the jet-like feature is roughly $`3.1\times 10^2\mathrm{counts}\mathrm{s}^1`$. The configuration of the jet feature relative to the nebula is remarkably similar to that of the Crab nebula as seen by ROSAT and which is now clearly resolved with Chandra (see Chandra publicity photo).
In short, the X-ray emission can be decomposed into three major morphological components: a point-like source, the surrounding nebula which shows evidence for a jet feature, and a patchy supernova remnant shell.
## 4. Discussions
A comparison between the X-ray-emitting nebula around PSR B0540$``$69 and the Crab nebula (see Chandra publicity photo<sup>1</sup><sup>1</sup>1Available at http://xrtpub.harvard.edu/photo/0052/0052\_hand.html) is very informative. The Crab nebula image shows a torus of X-ray-emitting loops, which most likely represents shocked pulsar wind materials consisting of magnetic waves and ultra-relativistic particles. Also clear are the two jets of X-ray-emitting material emanating from the Crab pulsar, in the direction perpendicular to the major axis of the torus. We speculate that the nebula around PSR B0540$``$69 has a similar structure. In fact, at the spatial resolution of Chandra at the LMC, the size, morphology, and surface intensity of the two nebulae are all remarkably similar (Fig. $`2b`$).
Assuming a power law spectrum for the X-ray emission from the SNR B0540$``$69 nebula of photon index 2.0 and N$`{}_{H}{}^{}=4\times 10^{21}\mathrm{cm}^2`$ (Finley et al. 1993, Wang & Gotthelf 1998a), the conversion between the count rate and the energy flux in the standard 1.0-10 keV band is $`1\times 10^{10}\mathrm{ergs}\mathrm{s}^1\mathrm{cm}^2/\mathrm{counts}\mathrm{s}^1`$. The corresponding total luminosity of the nebula is $`2.7\times 10^{37}\mathrm{ergs}\mathrm{s}^1`$, which is 13% of the spin-down energy of PSR B0540$``$69. The fraction is again similar to that of the Crab.
N157B (PSR J0537$``$6910) is the only other LMC SNR with a detected pulsar (16 ms) and also shows both an extended (resolved by ROSAT HRI), non-thermal nebula and a partial X-ray-emitting shell (Wang & Gotthelf 1998a; 1998b). The upcoming Chandra observations will make a detailed comparison between these two young Crab-like SNRs possible.
We gratefully acknowledge the Chandra team for making available the public data used herein. In particular we thank A. Rots and S. Murray for kindly making available their beta software. We thank U. Hwang for pointing out an instrumental artifact (“ghost image”) in the original HRC image of SNR B0540-69. We thank Dick Manchester for sending us the radio image and Fernando Camilo, David Helfand, and Jules Halpern for carefully reading the manuscript. This work was funded in part by NASA LTSA grants NAG5-7935 (E.V.G.) and NAG5-6413 (Q.D.W.). This is contribution #690 of the Columbia Astrophysics Laboratory. |
warning/0002/hep-ph0002243.html | ar5iv | text | # Angular Distribution of Charming 𝐵→𝑉𝑉 Decays and Time Evolution Effects
## I Introduction
The decays of the $`B`$ meson into two vector mesons $`BV_1+V_2`$, either with charm quarks in the final state particles, such as $`BJ/\mathrm{\Psi }\rho `$, or with particles without charm quarks, such as $`B\rho K^{}`$, have been calculated in many models . The time evolution effects in neutral $`B`$ meson decays are also discussed in . In this work, we would like to extend the general discussion on time evolving observables and to emphasize the charming decays in numerical analysis.
One major advantage of analyzing $`BVV`$ decays is that the interference of CP-even and CP-odd final states appear in the angular distributions. These interference terms provide good opportunities to observe CP or T violating effects. Since it is possible to measure all nine observables in certain decays , the physically interesting quantities such as $`\beta `$ and $`\eta `$ can be determined from experiments given sufficient statistics. In addition, relations among the nine observables provide a consistency check for the amplitude bilinears obtained experimentally.
The decay amplitude involves the hadronic matrix element of a $`B`$ meson decaying to two vector mesons through a weak current, which at present can not be calculated from first principles. Thus, our numerical evaluation of these observables at $`t=0`$ is based upon the assumptions of factorization and no final state interactions and form factor models, where the updated Wilson coefficients are used.
This paper is organized as follows: In Section II, we review the observables in the angular distributions of $`BVV`$ decays. In Section III, we derive the time-dependent formulas for the observables and list the complete results in the Appendix. Section IV considers the situation of no time evolution or at $`t=0`$. The case with no strong phases is discussed in Section V. Results of single weak amplitude decays are presented in Section VI, wherein CP asymmetries are also extensively discussed. In Section VII, we present the numerical estimation of the nine observables. We summarize this paper in Section VIII.
## II Observables and Angular Distributions in $`BVV`$ Decays
To extract the $`CP`$-odd and $`CP`$-even or $`T`$-odd and $`T`$-even components more easily, the angular distribution is often written in the transversity basis. Let us define the amplitude of $`BV_1V_2`$ in the rest frame of $`V_1`$. According to their polarization combinations, the amplitude can be decomposed into
$$A(BV_1V_2)=A_0ϵ_{V_1}^Lϵ_{V_2}^L\frac{A_{}}{\sqrt{2}}\stackrel{}{ϵ}_{V_1}^T\stackrel{}{ϵ}_{V_2}^Ti\frac{A_{}}{\sqrt{2}}\stackrel{}{ϵ}_{V_1}^{}\times \stackrel{}{ϵ}_{V_2}^{}\widehat{𝐩},$$
(1)
and similarly for $`\overline{B}\overline{V}_1\overline{V}_2`$. In Eq. (1), $`\stackrel{}{ϵ}_{V_1}`$ and $`\stackrel{}{ϵ}_{V_2}`$ are the unit polarization vectors of $`V_1`$ and $`V_2`$, respectively. $`\widehat{𝐩}`$ is the unit vector along the direction of motion of $`V_2`$ in the rest frame of $`V_1`$, $`ϵ_{V_i}^L\stackrel{}{ϵ}_{V_i}^{}\widehat{𝐩}`$ and $`\stackrel{}{ϵ}_{V_i}^T=\stackrel{}{ϵ}_{V_i}^{}ϵ_{V_i}^L\widehat{𝐩}`$. It is easy to see that $`A_{}`$ is odd under the parity transformation because of the appearance of $`\stackrel{}{ϵ}_{V_1}^{}\times \stackrel{}{ϵ}_{V_2}^{}\widehat{𝐩}`$, whereas $`A_0`$ and $`A_{}`$ are even.
The nine observables in the squared amplitude $`A^{}A`$ are
$`K_1(t)=|A_0(t)|^2,`$ $`K_4(t)=Re\left[A_0^{}(t)A_{}(t)\right],`$ $`L_4(t)=Im\left[A_0^{}(t)A_{}(t)\right],`$ (2)
$`K_2(t)=|A_{}(t)|^2,`$ $`K_5(t)=Im\left[A_0^{}(t)A_{}(t)\right],`$ $`L_5(t)=Re\left[A_0^{}(t)A_{}(t)\right],`$ (3)
$`K_3(t)=|A_{}(t)|^2,`$ $`K_6(t)=Im\left[A_{}^{}(t)A_{}(t)\right],`$ $`L_6(t)=Re\left[A_{}^{}(t)A_{}(t)\right].`$ (4)
So we have
$`A^{}(t)A(t)`$ $`=`$ $`K_1(t)X_1(\mathrm{\Omega })+K_2(t)X_2(\mathrm{\Omega })+K_3(t)X_3(\mathrm{\Omega })`$ (7)
$`+K_4(t)X_4(\mathrm{\Omega })+L_5(t)X_5(\mathrm{\Omega })+L_6(t)X_6(\mathrm{\Omega })`$
$`+L_4(t)Y_4(\mathrm{\Omega })+K_5(t)Y_5(\mathrm{\Omega })+K_6(t)Y_6(\mathrm{\Omega }),`$
where the quantities $`X_i(\mathrm{\Omega })`$ and $`Y_i(\mathrm{\Omega })`$ represent polarizations or polarization correlations of the final vector mesons and $`\mathrm{\Omega }`$ stands for the angles of the outgoing particles.
In general, the angular distribution of the decay in the transversity basis can be written as
$$\frac{d^3\mathrm{\Gamma }(t)}{d\mathrm{cos}\theta _1d\mathrm{cos}\theta _2d\varphi }=\underset{i}{}K_i(t)f_i(\theta _1,\theta _2,\varphi ),$$
(8)
where $`K_i`$’s are the amplitude bilinears that contain the dynamics and generally evolve with time, and $`f_i(\theta _1,\theta _2,\varphi )`$ are the corresponding angular distribution functions.
One can classify the decays into three types of processes according to the properties of the final product particles as follows:
Type I : For the case in which the decays of $`V_1`$ and $`V_2`$ are both into two pseudoscalar mesons, one can immediately translate the tensor correlations into angular distributions . The normalized angular distribution of the decays $`BV_1(P_1P_1^{})V_2(P_2P_2^{})`$, where $`P_1^{()}`$ and $`P_2^{()}`$ denote pseudoscalar mesons, is:
$`{\displaystyle \frac{1}{\mathrm{\Gamma }_0}}{\displaystyle \frac{d^3\mathrm{\Gamma }(t)}{d\mathrm{cos}\theta _1d\mathrm{cos}\theta _2d\varphi }}`$ $`=`$ $`{\displaystyle \frac{9}{8\pi }}\{{\displaystyle \frac{K_1(t)}{\mathrm{\Gamma }_0}}\mathrm{cos}^2\theta _1\mathrm{cos}^2\theta _2+{\displaystyle \frac{K_2(t)}{2\mathrm{\Gamma }_0}}\mathrm{sin}^2\theta _1\mathrm{sin}^2\theta _2\mathrm{cos}^2\varphi .`$ (11)
$`+{\displaystyle \frac{K_3(t)}{2\mathrm{\Gamma }_0}}\mathrm{sin}^2\theta _1\mathrm{sin}^2\theta _2\mathrm{sin}^2\varphi +{\displaystyle \frac{K_4(t)}{2\sqrt{2}\mathrm{\Gamma }_0}}\mathrm{sin}2\theta _1\mathrm{sin}2\theta _2\mathrm{cos}\varphi `$
$`.{\displaystyle \frac{K_5(t)}{2\sqrt{2}\mathrm{\Gamma }_0}}\mathrm{sin}2\theta _1\mathrm{sin}2\theta _2\mathrm{sin}\varphi {\displaystyle \frac{K_6(t)}{2\mathrm{\Gamma }_0}}\mathrm{sin}^2\theta _1\mathrm{sin}^2\theta _2\mathrm{sin}2\varphi \}.`$
Here $`\theta _1`$ ($`\theta _2`$) is the angle between the $`P_1`$ ($`P_2`$) three-momentum vector in the $`V_1(V_2)`$ rest frame and the $`V_1`$ ($`V_2`$) three-momentum vector defined in the $`B`$ rest frame, and $`\varphi `$ is the angle between the normals to the planes defined by $`P_1P_1^{}`$ and $`P_2P_2^{}`$, in the $`B`$ rest frame. Examples of such decays are $`B^+\overline{D}^{}(\overline{D}^0\pi ^0)\rho ^+(\pi ^+\pi ^0)`$, $`B_dD^{}(\overline{D}^0\pi ^{})\rho ^+(\pi ^+\pi ^0)`$, and $`B_dD^{}(\overline{D}^0\pi ^{})D^+(D^0\pi ^+)`$.
Type II : For the case of the decay $`BV_1(P_1P_1^{})V_2(l^+l^{})`$, suppose we observe that $`l^{}`$ is a right-handed particle and comes out in the direction $`\stackrel{}{k_2}=(\mathrm{sin}\theta _2\mathrm{cos}\varphi ,\mathrm{sin}\theta _2\mathrm{sin}\varphi ,\mathrm{cos}\theta _2)`$ and the momentum of $`P_1`$, $`\stackrel{}{k_1}=(\mathrm{sin}\theta _1,0,\mathrm{cos}\theta _1)`$ with angles defined in the same fashion as in the previous type, we have instead the differential angular distribution :
$`{\displaystyle \frac{1}{\mathrm{\Gamma }_0}}{\displaystyle \frac{d^3\mathrm{\Gamma }}{d\mathrm{cos}\theta _1d\mathrm{cos}\theta _2d\varphi }}`$ $`=`$ $`{\displaystyle \frac{9}{16\pi \mathrm{\Gamma }_0}}\{K_1\mathrm{cos}^2\theta _1\mathrm{sin}^2\theta _2+{\displaystyle \frac{K_2}{2}}(\mathrm{sin}^2\theta _1\mathrm{cos}^2\theta _2\mathrm{cos}^2\varphi +\mathrm{sin}^2\theta _1\mathrm{sin}^2\varphi ).`$ (15)
$`+{\displaystyle \frac{K_3}{2}}\left(\mathrm{sin}^2\theta _1\mathrm{cos}^2\theta _2\mathrm{sin}^2\varphi +\mathrm{sin}^2\theta _1\mathrm{cos}^2\varphi \right)+{\displaystyle \frac{K_4}{2\sqrt{2}}}\mathrm{sin}2\theta _1\mathrm{sin}2\theta _2\mathrm{cos}\varphi `$
$`{\displaystyle \frac{K_5}{2\sqrt{2}}}\mathrm{sin}2\theta _1\mathrm{sin}2\theta _2\mathrm{sin}\varphi {\displaystyle \frac{K_6}{2}}\mathrm{sin}^2\theta _1\mathrm{sin}^2\theta _2\mathrm{sin}2\varphi `$
$`.+{\displaystyle \frac{L_4}{\sqrt{2}}}\mathrm{sin}2\theta _1\mathrm{sin}\theta _2\mathrm{sin}\varphi {\displaystyle \frac{L_5}{\sqrt{2}}}\mathrm{sin}2\theta _1\mathrm{sin}\theta _2\mathrm{cos}\varphi +{\displaystyle \frac{L_6}{2}}\mathrm{sin}^2\theta _1\mathrm{cos}\theta _2\}`$
To obtain the result for the other possible final state with a left-handed outgoing $`l^{}`$, one only needs to flip the signs of $`L_4`$, $`L_5`$, and $`L_6`$. The muon polarization is equal to the sum of the terms $`L_4`$, $`L_5`$, $`L_6`$ divided by the sum of the other 6 terms. For the case of $`L_6`$ it is seen that the polarization does not vanish after integrating over $`\theta _1`$ and $`\varphi `$ and so the observation can be made without observing the $`V_1`$ decay. Such decay modes are $`B_u^+J/\mathrm{\Psi }(l^+l^{})K^+(\pi ^0K^+)`$, $`B_u^+J/\mathrm{\Psi }(l^+l^{})\rho ^+(\pi ^+\pi ^0)`$, $`B_dJ/\mathrm{\Psi }(l^+l^{})K^{}(\pi K)`$, $`B_dJ/\mathrm{\Psi }(l^+l^{})\rho ^0(\pi \pi )`$, $`B_dJ/\mathrm{\Psi }(l^+l^{})\omega (\pi ^+\pi ^{}\pi ^0)`$, $`B_sJ/\mathrm{\Psi }(l^+l^{})\overline{K}^{}(\pi \overline{K})`$, and $`B_dJ/\mathrm{\Psi }(l^+l^{})\varphi (K^+K^{})`$. Although $`\omega `$ decays into three pions, they are still correlated so that one can pick the normal direction to the decay plane formed by the three pions in the $`\omega `$ rest frame to define the direction $`\theta _2`$ and $`\varphi `$.
Although the $`BV(PP)V(P\gamma )`$ modes have a different decay pattern from that of $`BV(PP)V(l^+l^{})`$, they share the same differential angular distribution, with the direction of $`l^{}`$ in the latter case replaced by that of $`\gamma `$. For instance, for the decay with a right-handed circularly polarized photon in the final state, the angular distribution is the same as Eq. (15). Such examples are $`B_u^+D_s^+(D_s^+\gamma )\overline{D}^0(\overline{D}^0\pi )`$, $`B_dD_s^+(D_s^+\gamma )D^{}(\overline{D}^0\pi ^{})`$, $`B_sD_s^{}(D_s^{}\gamma )D^+(D^0\pi ^+)`$. If one does not measure the polarization of the product particles, the angular distribution would be the one by doubling Eq. (15) and eliminating the $`L_{4,5,6}`$ terms.
Type III : Next we consider the decay $`BV(P\gamma )V(P\gamma )`$. Since it is experimentally impractical to measure the polarizations of both photons in the final state, we just give here the differential angular distribution with no polarization measured:
$`{\displaystyle \frac{1}{\mathrm{\Gamma }_0}}{\displaystyle \frac{d^3\mathrm{\Gamma }}{d\mathrm{cos}\theta _1d\mathrm{cos}\theta _2d\varphi }}`$ $`=`$ $`{\displaystyle \frac{9}{8\pi \mathrm{\Gamma }_0}}\{K_1\mathrm{sin}^2\theta _1\mathrm{sin}^2\theta _2.`$ (20)
$`+{\displaystyle \frac{K_2}{2}}\left(\mathrm{cos}^2\theta _1\mathrm{cos}^2\theta _2\mathrm{cos}^2\varphi +\mathrm{cos}^2\theta _1\mathrm{sin}^2\varphi +\mathrm{cos}^2\theta _2\mathrm{sin}^2\varphi +\mathrm{cos}^2\varphi \right)`$
$`+{\displaystyle \frac{K_3}{2}}\left(\mathrm{cos}^2\theta _1\mathrm{cos}^2\theta _2\mathrm{sin}^2\varphi +\mathrm{cos}^2\theta _1\mathrm{cos}^2\varphi +\mathrm{cos}^2\theta _2\mathrm{cos}^2\varphi +\mathrm{sin}^2\varphi \right)`$
$`{\displaystyle \frac{K_4}{2\sqrt{2}}}\mathrm{sin}2\theta _1\mathrm{sin}2\theta _2\mathrm{cos}\varphi +{\displaystyle \frac{K_5}{2\sqrt{2}}}\mathrm{sin}2\theta _1\mathrm{sin}2\theta _2\mathrm{sin}\varphi `$
$`.+{\displaystyle \frac{K_6}{2}}\mathrm{sin}^2\theta _1\mathrm{sin}^2\theta _2\mathrm{sin}2\varphi \}`$
## III Time Evolution of the Amplitude Bilinears
The time evolution of an arbitrary neutral $`B`$ meson state $`a|B^0(t)+b|\overline{B}^0(t)`$ is governed by the Schrödinger equation
$$i\frac{d}{dt}\left(\begin{array}{c}a(t)\\ b(t)\end{array}\right)=\left(\begin{array}{c}a(t)\\ b(t)\end{array}\right).$$
(21)
If we write the mass eigenstates, $`|B_{L,H}`$, with eigenvalues $`m_{L,H}\frac{i}{2}\mathrm{\Gamma }_{L,H}`$ in terms of $`|B^0`$ and $`|\overline{B}^0`$ as
$$|B_{L,H}=p|B^0\pm q|\overline{B}^0,$$
(22)
then the time evolutions of $`B^0`$ and $`\overline{B}^0`$ are
$`|B^0(t)`$ $`=`$ $`g_+(t)|B^0(0)+{\displaystyle \frac{q}{p}}g_{}(t)|\overline{B}^0(t),`$ (23)
$`|\overline{B}^0(t)`$ $`=`$ $`{\displaystyle \frac{p}{q}}g_{}(t)|B^0(0)+g_+(t)|\overline{B}^0(t),`$ (24)
where
$$g_\pm (t)=\frac{1}{2}\left(e^{im_Lt}e^{\frac{1}{2}\mathrm{\Gamma }_Lt}\pm e^{im_Ht}e^{\frac{1}{2}\mathrm{\Gamma }_Ht}\right).$$
(25)
Suppose $`|f_\eta `$ is a state with definite CP property, namely, $`CP|f_\eta =\eta _i|f_\eta `$ for $`i=1,2,3`$ and $`\eta =0,,`$, respectively. The CP eigenvalues $`\eta _1=\eta _2=+1`$ and $`\eta _3=1`$. Suppose we write the decay matrix element of $`B^0`$ decaying into the final states $`f_\eta `$ at time $`t=0`$ as
$$A_\eta (0)f_\eta |B^0(0)=Y_{CKM}^Te^{i\theta _\eta }(T_\eta +P_\eta e^{i\varphi _w}e^{i\delta _\eta }).$$
(26)
$`Y_{CKM}^T`$ is the overall CKM factors appearing in the amplitudes. $`\theta _\eta `$ are the factored strong phases of $`A_\eta `$, but only the relative phases are essential. Conventionally, we take $`\theta _{}=0`$. $`T_\eta `$ and $`P_\eta `$ are the absolute values of two types of amplitudes that differ by a relative weak phase $`\varphi _w`$ and a relative strong phase $`\delta _\eta `$. We will refer to them by “tree” and “penguin” amplitudes, respectively. Similarly, for the CP conjugate mode we have
$$\overline{A}_\eta (0)\overline{f}_\eta |\overline{B}^0(0)=Y_{CKM}^{T}{}_{}{}^{}e^{i\theta _\eta }(T_\eta +P_\eta e^{i\varphi _w}e^{i\delta _\eta }).$$
(27)
Here we may assume that $`|f_\eta `$ and $`|\overline{f}_\eta `$ are the same state that both $`|B^0`$ and $`|\overline{B}^0`$ can decay into (e.g., $`B^0,\overline{B}^0J/\mathrm{\Psi }\varphi ,D^+D^{}`$). They can also be conjugate states so that only $`|B^0`$ (or $`|B^+`$) can decay into $`|f_\eta `$ and only $`|\overline{B}^0`$ (or $`|B^{}`$) to $`|\overline{f}_\eta `$. According to the time evolution, the decay amplitude at time $`t`$ would be
$$A_\eta (t)=f_\eta |B^0(t)=A_\eta (0)\left[g_+(t)+\eta _i\lambda _\eta g_{}(t)\right],$$
(28)
where
$$\lambda _\eta =\frac{q}{p}\frac{Y_{CKM}^{T}{}_{}{}^{}}{Y_{CKM}^T}\frac{T_\eta +P_\eta e^{i\varphi _w}e^{i\delta _\eta }}{T_\eta +P_\eta e^{i\varphi _w}e^{i\delta _\eta }}.$$
(29)
It is convenient to define a phase $`\varphi `$ by
$$e^{i\varphi }\frac{q}{p}\frac{Y_{CKM}^{T}{}_{}{}^{}}{Y_{CKM}^T},$$
(30)
and
$$R_\eta Re\left[\frac{T_\eta +P_\eta e^{i\varphi _w}e^{i\delta _\eta }}{T_\eta +P_\eta e^{i\varphi _w}e^{i\delta _\eta }}\right],I_\eta Im\left[\frac{T_\eta +P_\eta e^{i\varphi _w}e^{i\delta _\eta }}{T_\eta +P_\eta e^{i\varphi _w}e^{i\delta _\eta }}\right].$$
(31)
Note that $`R_\eta ^2+I_\eta ^2=1`$ if and only if $`\delta _\eta ,\varphi _w=0`$ (mod $`\pi `$). If either (i) no nontrivial relative weak phase ($`0`$ or $`\pi `$), (ii) negligible tree contributions ($`T_\eta 0`$), or (iii) negligible penguin contributions ($`P_\eta 0`$) happens, then $`R_\eta =1`$ and $`I_\eta =0`$, apart from a possible overall phase.
With the above definitions, one can get, for example, the time evolving $`|A_\eta (t)|^2`$ as follows:
$`|A_\eta (t)|^2=|A_\eta (0)|^2e^{\mathrm{\Gamma }t}\{{\displaystyle \frac{1+R_\eta ^2+I_\eta ^2}{2}}\mathrm{cosh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)+{\displaystyle \frac{1R_\eta ^2I_\eta ^2}{2}}\mathrm{cos}\left(\mathrm{\Delta }mt\right)`$ (32)
$`+\eta _i[(R_\eta \mathrm{cos}\varphi I_\eta \mathrm{sin}\varphi )\mathrm{sinh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)(R_\eta \mathrm{sin}\varphi +I_\eta \mathrm{cos}\varphi )\mathrm{sin}\left(\mathrm{\Delta }mt\right)]\},`$ (33)
where $`\mathrm{\Delta }mm_Hm_L`$ and $`\mathrm{\Delta }\mathrm{\Gamma }\mathrm{\Gamma }_H\mathrm{\Gamma }_L`$.
Similarly, one uses the time evolution for the conjugate mode to get, along with Eq. (27), for example, the corresponding time evolution formulas for $`|\overline{A}_\eta (t)|^2`$:
$`|\overline{A}_\eta (t)|^2=|A_\eta (0)|^2e^{\mathrm{\Gamma }t}\{{\displaystyle \frac{1+R_\eta ^2+I_\eta ^2}{2}}\mathrm{cosh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right){\displaystyle \frac{1R_\eta ^2I_\eta ^2}{2}}\mathrm{cos}\left(\mathrm{\Delta }mt\right).`$ (34)
$`.+\eta _i[(R_\eta \mathrm{cos}\varphi I_\eta \mathrm{sin}\varphi )\mathrm{sinh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)+(R_\eta \mathrm{sin}\varphi +I_\eta \mathrm{cos}\varphi )\mathrm{sin}\left(\mathrm{\Delta }mt\right)]\}.`$ (35)
A complete list of all the observable amplitude bilinears and their CP conjugates is given in the Appendix.
Before we proceed the discussion, let’s define the CP asymmetry parameters, $`\zeta _i(t)K_i(t)\overline{K}_i(t)`$ for $`i=1,2,3\mathrm{},6`$ and $`\xi _i(t)L_i(t)\overline{L}_i(t)`$ for $`i=4,5,6`$. These nine parameters measure the changes of the amplitude bilinears under the CP transformation. For instance, from Eqs. (32) and (34), we obtain
$$\zeta _1(t)=K_1(0)e^{\mathrm{\Gamma }t}\left[\left(1R_0^2I_0^2\right)\mathrm{cos}\left(\mathrm{\Delta }mt\right)2\left(R_0\mathrm{sin}\varphi +I_0\mathrm{cos}\varphi \right)\mathrm{sin}\left(\mathrm{\Delta }mt\right)\right].$$
(36)
This relation along with others for $`K_{2,3}(t)`$ provide information on $`\varphi `$ given $`\mathrm{\Delta }m`$ and $`\mathrm{\Gamma }`$ extracted from other experiments and theoretical estimates of $`K_{1,2,3}(0)`$, $`R_{0,,}`$ and $`I_{0,,}`$.
## IV Case I: No Time Evolution
If we take $`t=0`$ in Eqs. (76)-(106) and (134)-(152), we get the bilinear formulas for neutral $`B`$ meson decays at time $`t=0`$, or the charged $`B`$ meson decays. The relations between the conjugate amplitude bilinears and amplitude bilinears are
$`\overline{K}_i`$ $`=`$ $`\left(R_\eta ^2+I_\eta ^2\right)K_i,\mathrm{for}i=1,2,3,`$ (37)
$`\overline{K}_4`$ $`=`$ $`\left(R_{}R_0+I_{}I_0\right)K_4\left(I_{}R_0R_{}I_0\right)L_4,`$ (38)
$`\overline{K}_{5,6}`$ $`=`$ $`\left(R_{}R_{0,}+I_{}I_{0,}\right)K_{5,6}+\left(I_{}R_{0,}R_{}I_{0,}\right)L_{5,6},`$ (39)
$`\overline{L}_4`$ $`=`$ $`\left(R_{}R_0+I_{}I_0\right)L_4+\left(I_{}R_0R_{}I_0\right)K_4,`$ (40)
$`\overline{L}_{5,6}`$ $`=`$ $`\left(R_{}R_{0,}+I_{}I_{0,}\right)L_{5,6}\left(I_{}R_{0,}R_{}I_{0,}\right)K_{5,6},`$ (41)
As discussed in the paragraph after Eq. (31), if none of the relative strong and weak phases are trivial, i.e., $`0`$ or $`\pi `$, CP asymmetry exists in the above bilinears. However, if there are no strong phases (including all the factored strong phases and relative phases) but the relative weak phase is nontrivial, then one can simplify the above equations to get $`\overline{K}_{1,2,3,4}=K_{1,2,3,4}`$, $`\overline{L}_{5,6}=L_{5,6}`$, $`\overline{K}_{5,6}=K_{5,6}`$, and $`\overline{L}_4=L_4`$. This effect is purely due to that fact that there is a relative weak phase and $`Im[A_0^{}A_{}]`$, $`Im[A_0^{}A_{}]`$, and $`Im[A_{}^{}A_{}]`$ are CP odd quantities .
The observation of CP asymmetries in any of the bilinears indicates that nontrivial strong and weak phases are involved in the decay. Therefore, if the relative weak phase within the Standard Model is trivial, that is, effectively only one weak amplitude dominates, then no CP asymmetry will be observed among all the bilinears.
The formulae presented in this section can be applied to $`B_u^+D^+\overline{D}^{}`$, $`B_u^+J/\mathrm{\Psi }\rho ^+`$, $`B_u^+D_s^+\overline{D}^{}`$, and $`B_u^+J/\mathrm{\Psi }K^+`$. One can only measure $`K_{16}`$ in the first decay mode because it is a Type I decay. There is no nontrivial weak phases in the latter two decays. Therefore, one should not expect to observe CP asymmetries in the observables; but $`K_{5,6}`$ and $`L_4`$ may be nonzero, and provide evidence for strong phases due to final state interactions. The observation of CP asymmetries in such modes indicates new CP violating source from physics beyond the Standard Model.
## V Case II: No Strong Phases
If there is no strong phases involved in the decays, then $`R_\eta ^2+I_\eta ^2=1`$. One can write
$$R_\eta =\mathrm{cos}2\alpha _\eta ,I_\eta =\mathrm{sin}2\alpha _\eta ,$$
(42)
where $`\alpha _\eta `$ is the phase of $`T_\eta +P_\eta e^{i\varphi _w}`$. With Eq. (42) and the definition of the phase $`\varphi `$ in Eq. (30), one can get the nine CP asymmetry parameters
$`\zeta _i(t)`$ $`=`$ $`2\eta _iK_i(0)e^{\mathrm{\Gamma }t}\mathrm{sin}(2\alpha _\eta \varphi )\mathrm{sin}(\mathrm{\Delta }mt),\mathrm{for}i=1,2,3,`$ (43)
$`\zeta _4(t)`$ $`=`$ $`K_4(0)e^{\mathrm{\Gamma }t}\left[\mathrm{sin}(2\alpha _{}\varphi )+\mathrm{sin}(2\alpha _0\varphi )\right]\mathrm{sin}(\mathrm{\Delta }mt)`$ (45)
$`L_4(0)e^{\mathrm{\Gamma }t}\left[\mathrm{cos}(2\alpha _{}\varphi )\mathrm{cos}(2\alpha _0\varphi )\right]\mathrm{sin}(\mathrm{\Delta }mt),`$
$`\zeta _5(t)`$ $`=`$ $`K_5(0)e^{\mathrm{\Gamma }t}\{\mathrm{cosh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)+\mathrm{cos}(\mathrm{\Delta }mt)`$ (50)
$`+\left[\mathrm{cos}(2\alpha _0\varphi )\mathrm{cos}(2\alpha _{}\varphi )\right]\mathrm{sinh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)`$
$`\mathrm{cos}(2\alpha _02\alpha _{})[\mathrm{cosh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)\mathrm{cos}(\mathrm{\Delta }mt)]\}`$
$`+L_5(0)e^{\mathrm{\Gamma }t}\{[\mathrm{sin}(2\alpha _0\varphi )+\mathrm{sin}(2\alpha _{}\varphi )]\mathrm{sinh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)`$
$`\mathrm{sin}(2\alpha _02\alpha _{})[\mathrm{cosh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)\mathrm{cos}(\mathrm{\Delta }mt)]\},`$
$`\xi _4(t)`$ $`=`$ $`L_4(0)e^{\mathrm{\Gamma }t}\{\mathrm{cosh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)+\mathrm{cos}(\mathrm{\Delta }mt)`$ (55)
$`+\left[\mathrm{cos}(2\alpha _0\varphi )+\mathrm{cos}(2\alpha _{}\varphi )\right]\mathrm{sinh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)`$
$`+\mathrm{cos}(2\alpha _02\alpha _{})[\mathrm{cosh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)\mathrm{cos}(\mathrm{\Delta }mt)]\}`$
$`K_4(0)e^{\mathrm{\Gamma }t}\{[\mathrm{sin}(2\alpha _{}\varphi )\mathrm{sin}(2\alpha _0\varphi )]\mathrm{sinh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)`$
$`\mathrm{sin}(2\alpha _{}2\alpha _0)[\mathrm{cosh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)\mathrm{cos}(\mathrm{\Delta }mt)]\},`$
$`\xi _5(t)`$ $`=`$ $`L_5(0)e^{\mathrm{\Gamma }t}\left[\mathrm{sin}(2\alpha _{}\varphi )+\mathrm{sin}(2\alpha _0\varphi )\right]\mathrm{sin}(\mathrm{\Delta }mt)`$ (57)
$`+K_5(0)e^{\mathrm{\Gamma }t}\left[\mathrm{cos}(2\alpha _{}\varphi )\mathrm{cos}(2\alpha _0\varphi )\right]\mathrm{sin}(\mathrm{\Delta }mt).`$
The formulas for $`\zeta _6(t)`$ and $`\xi _6(6)`$ can be obtained by replacing “0” in Eq. (50) and (57) by “$``$”. Thus, in principle, one may extract information about the phase combinations $`2\alpha _{0,,}\varphi `$. One should notice that $`\xi _4`$ and $`\zeta _{5,6}`$ can be nonzero at $`t=0`$, whereas the others are identically zero. Although the assumption of no strong phases is unlikely to be true in charming decays, it may be applied to charmless decays such as $`B\rho \rho `$.
## VI Case III: No Relative Weak Phase
If there is no relative weak phase in each transversity amplitude, namely, $`\varphi _w=0`$, then one gets $`R_\eta =1`$ and $`I_\eta =0`$. This case is equivalent to the cases where only one type of amplitude dominates the decay process. For completeness, we list time evolutions of the nine observables in Tables I and II While we agree with in $`K_{16}`$ and $`\overline{K}_{14}`$, our $`\overline{K}_{5,6}`$ differ from theirs by an overall minus sign..
So the nine CP asymmetry parameters are
$`\zeta _i(t)`$ $`=`$ $`2\eta _iK_i(0)e^{\mathrm{\Gamma }t}\mathrm{sin}\varphi \mathrm{sin}(\mathrm{\Delta }mt),\mathrm{for}i=1,2,3,`$ (58)
$`\zeta _4(t)`$ $`=`$ $`2K_4(0)e^{\mathrm{\Gamma }t}\mathrm{sin}\varphi \mathrm{sin}(\mathrm{\Delta }mt),`$ (59)
$`\zeta _{5,6}(t)`$ $`=`$ $`2L_{5,6}(0)e^{\mathrm{\Gamma }t}\mathrm{sin}\varphi \mathrm{sinh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right),`$ (60)
$`\xi _4(t)`$ $`=`$ $`2L_4(0)e^{\mathrm{\Gamma }t}\mathrm{sin}\varphi \mathrm{sin}(\mathrm{\Delta }mt),`$ (61)
$`\xi _{5,6}(t)`$ $`=`$ $`2K_{5,6}(0)e^{\mathrm{\Gamma }t}\mathrm{sin}\varphi \mathrm{sinh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right).`$ (62)
These equations hold even if there are nontrivial strong phases.
If we fix the overall strong phases of the transversity amplitudes by the following convention: $`A_{}(0)=|A_{}(0)|`$, $`A_0(0)=|A_0(0)|e^{i\delta _0}`$, and $`A_{}(0)=|A_{}(0)|e^{i\delta _{}}`$ Here we ignore the common weak factor that will be cancelled in all amplitude bilinears., then $`K_{4,5,6}(0)`$ and $`L_{4,5,6}(0)`$ can be rewritten as
$`K_4(0)=\sqrt{K_1(0)K_2(0)}\mathrm{cos}(\delta _0\delta _{}),`$ $`L_4(0)=\sqrt{K_1(0)K_2(0)}\mathrm{sin}(\delta _0\delta _{}),`$ (63)
$`K_5(0)=\sqrt{K_1(0)K_3(0)}\mathrm{cos}(\delta _0),`$ $`L_5(0)=\sqrt{K_1(0)K_3(0)}\mathrm{sin}(\delta _0),`$ (64)
$`K_6(0)=\sqrt{K_2(0)K_3(0)}\mathrm{cos}(\delta _{}),`$ $`L_6(0)=\sqrt{K_2(0)K_3(0)}\mathrm{sin}(\delta _{}).`$ (65)
One can readily reach four relations among them:
$`K_1(0)K_2(0)=K_4(0)^2+L_4(0)^2,`$ $`K_2(0)K_3(0)=K_6(0)^2+L_6(0)^2,`$ (66)
$`K_3(0)K_1(0)=K_5(0)^2+L_5(0)^2,`$ $`{\displaystyle \frac{L_4(0)}{K_4(0)}}={\displaystyle \frac{L_5(0)K_6(0)K_5(0)L_6(0)}{K_5(0)K_6(0)+L_5(0)L_6(0)}}.`$ (67)
All experimentally measured nine amplitude bilinears should obey the above consistency relations. If the strong phases $`\delta _0`$ and $`\delta _{}`$ are nontrivial, one could possibly get sizeable $`L_{4,5,6}`$ that can be observed experimentally. We can then obtain information on the strong phases $`\delta _0`$, $`\delta _{}`$, the mass difference $`\mathrm{\Delta }m`$, the decay width difference $`\mathrm{\Delta }\mathrm{\Gamma }`$, and $`\mathrm{sin}\varphi `$ from Eqs. (58). Since some of them share the same time evolution pattern, they also provide a consistency check for the experimental results.
At $`t=0`$, there is no $`CP`$ asymmetry at all. So for charged $`B`$ decays where one weak amplitude dominates in the Standard Model, one should get the same amplitude bilinears for the particle and its conjugate modes. However, for neutral $`B`$ decays, the asymmetries develop as time goes on due to the mixing effect. In particular, $`\zeta _{14}(t)`$ and $`\xi _4(t)`$ have a sinusoidal time dependence, while $`\zeta _{5,6}`$ and $`\xi _{5,6}(t)`$ decay exponentially at $`B_L`$’s decay rate, $`\mathrm{\Gamma }_L`$, in the large $`t`$ limit.
It is, nevertheless, interesting to look at the time integrated quantities for sizeable CP or T violating effects. The particle total decay rate, after time integration, is
$`{\displaystyle _0^{\mathrm{}}}𝑑t\left[K_1(t)+K_2(t)+K_3(t)\right]`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{\Gamma }}}\{{\displaystyle \frac{4}{4y^2}}[K_1(0)+K_2(0)+K_3(0)].`$ (69)
$`.(\mathrm{cos}\varphi {\displaystyle \frac{2y}{4y^2}}\mathrm{sin}\varphi {\displaystyle \frac{x}{1+x^2}})[K_1(0)+K_2(0)K_3(0)]\}.`$
Similarly, the time integrated anti-particle total decay rate can be obtained by simply reversing the sign of $`\varphi `$ in Eq. (69). One can obtain $`\mathrm{sin}\varphi `$ from the asymmetry between the time integrated total rates of conjugate modes and $`\mathrm{cos}\varphi `$ from the untagged analysis. This then eliminates the discrete ambiguity in the angle $`\varphi `$.
By integrating Eq. (58) from $`t=0`$ to $`t=\mathrm{}`$, we find the asymmetries to be
$`{\displaystyle _0^{\mathrm{}}}𝑑t\zeta _i(t)`$ $`=`$ $`2\eta _iK_i(0){\displaystyle \frac{1}{\mathrm{\Gamma }}}{\displaystyle \frac{2x}{1+x^2}}\mathrm{sin}\varphi ,\mathrm{for}i=1,2,3,`$ (70)
$`{\displaystyle _0^{\mathrm{}}}𝑑t\zeta _4(t)`$ $`=`$ $`2K_4(0){\displaystyle \frac{1}{\mathrm{\Gamma }}}{\displaystyle \frac{2x}{1+x^2}}\mathrm{sin}\varphi ,`$ (71)
$`{\displaystyle _0^{\mathrm{}}}𝑑t\zeta _{5,6}(t)`$ $`=`$ $`2L_{5,6}(0){\displaystyle \frac{1}{\mathrm{\Gamma }}}{\displaystyle \frac{2y}{4y^2}}\mathrm{sin}\varphi ,`$ (72)
$`{\displaystyle _0^{\mathrm{}}}𝑑t\xi _4(t)`$ $`=`$ $`2L_4(0){\displaystyle \frac{1}{\mathrm{\Gamma }}}{\displaystyle \frac{2x}{1+x^2}}\mathrm{sin}\varphi ,`$ (73)
$`{\displaystyle _0^{\mathrm{}}}𝑑t\xi _{5,6}(t)`$ $`=`$ $`2K_{5,6}(0){\displaystyle \frac{1}{\mathrm{\Gamma }}}{\displaystyle \frac{2y}{4y^2}}\mathrm{sin}\varphi ,`$ (74)
In the above equations, $`x\mathrm{\Delta }m/\mathrm{\Gamma }`$ and $`y\mathrm{\Delta }\mathrm{\Gamma }/\mathrm{\Gamma }`$. For $`B_d`$, $`x=0.73`$ and $`y`$ is negligibly small; for $`B_s`$, $`x>14.0(\mathrm{CL}=95\%)`$ and $`y<0.67(\mathrm{CL}=95\%)`$ . From these relations, one can also directly extract $`\mathrm{sin}\varphi `$ given the information about the amplitude bilinears at initial time.
In principle, one can extract information about $`\varphi `$ and strong phases $`\delta _0`$, $`\delta _{}`$ either from the time-dependent CP asymmetries $`\zeta `$’s and $`\xi `$’s or from the integrated asymmetries once the bilinears are determined experimentally or from models. $`\mathrm{sin}2\beta `$ has been measured from the mixing-induced CP asymmetry of $`B_dJ/\mathrm{\Psi }K_S`$ . For $`B_dJ/\mathrm{\Psi }K^{}(\pi K_S)`$, the $`B_d\overline{B}_d`$ and $`K\overline{K}`$ mixings and the CKM factor in the weak decay amplitude also give $`\varphi =2\beta `$ . Therefore, this offers an alternative way of measuring $`\mathrm{sin}2\beta `$ through the angular distribution analysis of tagged $`B_d`$ decays. In addition, the $`\mathrm{cos}2\beta `$ dependence in $`K_{5,6}`$ and $`L_{5,6}`$ helps resolving the discrete ambiguity of the CKM angle $`\beta `$ .
For $`B_sD_s^+D_s^{}`$ and $`B_sJ/\mathrm{\Psi }\varphi `$, $`\varphi =2\lambda ^2\eta =𝒪(0.03)`$ to an extremely good approximation, where $`\lambda =0.22`$ is the Cabibbo angle and $`\eta `$ is one of the Wolfenstein parameters . In this case, the CP asymmetries $`\zeta _{1,2,3}(t)`$ can be used to provide an unambiguous determination of the sign of $`\varphi `$, and therefore the sign of $`\eta `$.
The decays that one may apply the results in this section to include: $`B_sD_s^+D_s^{}`$, $`B_dJ/\mathrm{\Psi }K^{}(\pi K_S)`$, $`B_dJ/\mathrm{\Psi }\varphi `$. The first mode is a Type III decay, whereas the latter two are Type II decays.
Notice that $`\mathrm{sin}\varphi `$ appears in all CP asymmetries, where $`\varphi `$ is the phase of mixing and the single CKM factor involved in the decay amplitude. Since the amplitude has only one CKM factor, no CP violation effects would be found in the nine observables at $`t=0`$. Yet the mixing will produce differences between the particle and anti-particle decay modes as time goes on. So any observation of the CP asymmetries in such modes indicates CP violation due to mixing and decay.
## VII Numerical Calculation
In this section, we apply the factorization hypothesis to the calculation of hadronic decay amplitudes. In general, factorization is expected to hold more strongly for color-allowed processes, such as $`B_qD_s^+\overline{D}_q^{}`$ with $`qs,d,u`$, though it is doubtful in color-suppressed modes, such as $`B_qJ/\mathrm{\Psi }V`$ with $`(q,V)(s,\varphi ),(d,K^0),(u,K^+)`$ . Throughout the calculations, we ignore the strong phases $`\theta _\eta `$ in Eq. (26) for simplicity but keep the strong phases in the Wilson coefficients . These strong phases may be extracted from experimental data as mentioned in the previous section.
In our calculations, the Wolfenstein parameters are $`(\rho ,\eta )=(0.18,0.37)`$. The decay constants used are $`F_D^{}=230MeV`$, $`F_{D_s^{}}=275MeV`$, $`F_{J/\mathrm{\Psi }}=394MeV`$, $`F_K^{}=221MeV`$, and $`F_\varphi =237MeV`$. Extracting dominant Wilson coefficients, $`a_1`$ and $`a_2`$, from experimental decay rates has been performed in . We apply their results to $`B^\pm `$ and $`B_d`$ decays with $`bs`$ quark level transitions, i.e., $`B_uD_s^+\overline{D}^0`$, $`B_dD_s^+D^{}`$, $`B_uJ/\mathrm{\Psi }K^+`$, and $`B_dJ/\mathrm{\Psi }K^0`$. We then extend the results to other decays according to the final state configurations, color-allowed ($`B_qD_s^+\overline{D}_q^{}`$ and $`B_qD^+D^{}`$) or color-suppressed ($`B_qJ/\mathrm{\Psi }V`$), charged or neutral. Assuming heavy quark symmetry, we use the $`BD^{}`$ decay form factors for the $`BD_s^{}`$ transitions.
The bilinears $`K_i`$ and $`L_i`$ in the following tables are normalized by dividing with $`\mathrm{\Gamma }_0K_1+K_2+K_3`$. The branching ratio asymmetry is defined by
$$a_{CP}\frac{𝒜\overline{𝒜}}{𝒜+\overline{𝒜}},$$
(75)
where $`𝒜`$ and $`\overline{𝒜}`$ are the branching ratios for the particle and antiparticle decays, respectively. In the following tables, we list all nine amplitude bilinears for each mode even if $`L_{4,5,6}`$ may not be able to be observed from the angular distributions of some of them (Type I decays).
In Tables III and IV, we take the modified BSW (or BSW II) model for the form factors in the evaluation of hadronic matrix elements. In Tables V and VI, the Neubert-Stech (NS) model is used. The relativistic light-front (LF) model is applied to the calculations in Tables VII and VIII.
We see that in general (1) there are no CP asymmetries for the nine normalized observables at the initial time (yet the CP asymmetries do exist for the unnormalized observables); (2) the branching ratio CP asymmetry is larger in decays involving the $`bd`$ quark level processes because of the relative phase between the CKM factors of two weak amplitudes; and (3) $`L_4`$, $`K_5`$, and $`K_6`$ that involve the imaginary parts of the amplitude bilinears are essentially zero because the tree amplitudes dominates over the penguin contributions in these decays and we ignore possible final state interaction phases.
Although the branching ratios and nine observables may vary as one uses different form factor models, the asymmetries are roughly the same. It is found that $`bd`$ type decays have larger asymmetries than $`bs`$ type ones, as one would expect. From the models we analyze, the asymmetries for $`B_u^+D^+\overline{D}^0`$ and $`B_dD^+D^{}`$ range from $`3.29\%`$ to $`3.53\%`$ and from $`4.14\%`$ to $`4.46\%`$, respectively. Unlike the charmless decays where $`A_0`$ is the dominant component in the transversity amplitudes, both $`A_0`$ and $`A_{}`$ are about the same size. The parity odd component $`A_{}`$ is still small in $`BD^{}D^{}`$ type transitions, but larger in $`BJ/\mathrm{\Psi }V`$ decays.
## VIII Summary
The decays of a $`B`$ meson into two vector mesons, which subsequently decay into two lighter particles via CP conserving currents, have a specific pattern in the differential angular distributions. The coefficient of each angular function in the distribution is a bilinear of amplitudes with certain CP properties. The knowledge of these amplitude bilinears can help us observe and understand CP violating effects. In particular, the time evolution of these bilinears in the cases of neutral $`B`$ meson decays further reveals the information such as the mass and decay width differences and CKM parameters. In situations where we can measure the polarization of the final product particles, all the nine combinations of amplitude bilinears are observable.
Under certain special circumstances, one can find simple relations among the nine observables. Therefore, experimental determination of them is valuable in testing our theoretical assumptions in the calculations, such as the factorization hypothesis and form factor models. The results are particularly simplified when there is only one weak amplitude dominating in the decay process. In such cases, one can test the Standard Model from the CP asymmetries at $`t=0`$. The time development of CP asymmetries provides a window for observing CP violations due to mixing effects.
We provide numerical estimates of the observables in 12 sets of charming $`BVV`$ decays using three different form factor models. We find that the results do not depend strongly on the models used. In particular, we find bigger branching ratio asymmetries in $`bd`$ type decays, and those for $`B_u^+D^+\overline{D}^0`$ and $`B_dD^+D^{}`$ are as large as $`3\%`$ and $`4\%`$, respectively.
ACKNOWLEDGMENT
This research work is supported by the Department of Energy under Grant No. DE-FG02-91ER40682. The author is grateful to F. Gilman, L. Wolfenstein for useful discussion and to A. Leibovich for his comments. He also would like to thank the hospitality of Academia Sinica in Taiwan where part of this work is done.
## Time-Dependent Amplitude Bilinears
The time evolutions of the amplitude bilinears are as follows:
$`|A_\eta (t)|^2=|A_\eta (0)|^2e^{\mathrm{\Gamma }t}\{{\displaystyle \frac{1+R_\eta ^2+I_\eta ^2}{2}}\mathrm{cosh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)+{\displaystyle \frac{1R_\eta ^2I_\eta ^2}{2}}\mathrm{cos}\left(\mathrm{\Delta }mt\right)`$ (76)
$`+\eta _i[(R_\eta \mathrm{cos}\varphi I_\eta \mathrm{sin}\varphi )\mathrm{sinh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)(R_\eta \mathrm{sin}\varphi +I_\eta \mathrm{cos}\varphi )\mathrm{sin}\left(\mathrm{\Delta }mt\right)]\},`$ (77)
$`\mathrm{where}i=1,2,3\mathrm{for}\eta =0,,,\mathrm{respectively},\mathrm{and}\eta _{1,2}=\eta _3=1;`$ (78)
$`Re\left[A_0^{}(t)A_{}(t)\right]=Re\left[A_0^{}(0)A_{}(0)\right]{\displaystyle \frac{e^{\mathrm{\Gamma }t}}{2}}\times `$ (79)
$`\{[\mathrm{cosh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)+\mathrm{cos}\left(\mathrm{\Delta }mt\right)].`$ (80)
$`+(R_{}\mathrm{cos}\varphi I_{}\mathrm{sin}\varphi )\mathrm{sinh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)(R_{}\mathrm{sin}\varphi +I_{}\mathrm{cos}\varphi )\mathrm{sin}\left(\mathrm{\Delta }mt\right)`$ (81)
$`+(R_0\mathrm{cos}\varphi I_0\mathrm{sin}\varphi )\mathrm{sinh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)(R_0\mathrm{sin}\varphi +I_0\mathrm{cos}\varphi )\mathrm{sin}\left(\mathrm{\Delta }mt\right)`$ (82)
$`.+(R_{}R_0+I_{}I_0)[\mathrm{cosh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)\mathrm{cos}\left(\mathrm{\Delta }mt\right)]\}`$ (83)
$`Im\left[A_0^{}(0)A_{}(0)\right]{\displaystyle \frac{e^{\mathrm{\Gamma }t}}{2}}\times `$ (84)
$`\{(R_{}\mathrm{sin}\varphi +I_{}\mathrm{cos}\varphi )\mathrm{sinh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)+(R_{}\mathrm{cos}\varphi I_{}\mathrm{sin}\varphi )\mathrm{sin}\left(\mathrm{\Delta }mt\right).`$ (85)
$`(R_0\mathrm{sin}\varphi +I_0\mathrm{cos}\varphi )\mathrm{sinh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)(R_0\mathrm{cos}\varphi I_0\mathrm{sin}\varphi )\mathrm{sin}\left(\mathrm{\Delta }mt\right)`$ (86)
$`.+(I_{}R_0R_{}I_0)[\mathrm{cosh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)\mathrm{cos}\left(\mathrm{\Delta }mt\right)]\};`$ (87)
$`Im\left[A_0^{}(t)A_{}(t)\right]=Im\left[A_0^{}(0)A_{}(0)\right]{\displaystyle \frac{e^{\mathrm{\Gamma }t}}{2}}\times `$ (88)
$`\{[\mathrm{cosh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)+\mathrm{cos}\left(\mathrm{\Delta }mt\right)].`$ (89)
$`+(R_{}\mathrm{cos}\varphi I_{}\mathrm{sin}\varphi )\mathrm{sinh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)(R_{}\mathrm{sin}\varphi +I_{}\mathrm{cos}\varphi )\mathrm{sin}\left(\mathrm{\Delta }mt\right)`$ (90)
$`+(R_0\mathrm{cos}\varphi I_0\mathrm{sin}\varphi )\mathrm{sinh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)(R_0\mathrm{sin}\varphi +I_0\mathrm{cos}\varphi )\mathrm{sin}\left(\mathrm{\Delta }mt\right)`$ (91)
$`.+(R_{}R_0+I_{}I_0)[\mathrm{cosh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)\mathrm{cos}\left(\mathrm{\Delta }mt\right)]\}`$ (92)
$`+Re\left[A_0^{}(0)A_{}(0)\right]{\displaystyle \frac{e^{\mathrm{\Gamma }t}}{2}}\times `$ (93)
$`\{(R_{}\mathrm{sin}\varphi +I_{}\mathrm{cos}\varphi )\mathrm{sinh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)+(R_{}\mathrm{cos}\varphi I_{}\mathrm{sin}\varphi )\mathrm{sin}\left(\mathrm{\Delta }mt\right).`$ (94)
$`(R_0\mathrm{sin}\varphi +I_0\mathrm{cos}\varphi )\mathrm{sinh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)(R_0\mathrm{cos}\varphi I_0\mathrm{sin}\varphi )\mathrm{sin}\left(\mathrm{\Delta }mt\right)`$ (95)
$`.+(I_{}R_0R_{}I_0)[\mathrm{cosh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)\mathrm{cos}\left(\mathrm{\Delta }mt\right)]\};`$ (96)
$`Re\left[A_0^{}(t)A_{}(t)\right]=Re\left[A_0^{}(0)A_{}(0)\right]{\displaystyle \frac{e^{\mathrm{\Gamma }t}}{2}}\times `$ (97)
$`\{[\mathrm{cosh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)+\mathrm{cos}\left(\mathrm{\Delta }mt\right)].`$ (98)
$`(R_{}\mathrm{cos}\varphi I_{}\mathrm{sin}\varphi )\mathrm{sinh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)+(R_{}\mathrm{sin}\varphi +I_{}\mathrm{cos}\varphi )\mathrm{sin}\left(\mathrm{\Delta }mt\right)`$ (99)
$`+(R_0\mathrm{cos}\varphi I_0\mathrm{sin}\varphi )\mathrm{sinh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)(R_0\mathrm{sin}\varphi +I_0\mathrm{cos}\varphi )\mathrm{sin}\left(\mathrm{\Delta }mt\right)`$ (100)
$`.(R_{}R_0+I_{}I_0)[\mathrm{cosh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)\mathrm{cos}\left(\mathrm{\Delta }mt\right)]\}`$ (101)
$`+Im\left[A_0^{}(0)A_{}(0)\right]{\displaystyle \frac{e^{\mathrm{\Gamma }t}}{2}}\times `$ (102)
$`\{(R_{}\mathrm{sin}\varphi +I_{}\mathrm{cos}\varphi )\mathrm{sinh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)+(R_{}\mathrm{cos}\varphi I_{}\mathrm{sin}\varphi )\mathrm{sin}\left(\mathrm{\Delta }mt\right).`$ (103)
$`+(R_0\mathrm{sin}\varphi +I_0\mathrm{cos}\varphi )\mathrm{sinh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)+(R_0\mathrm{cos}\varphi I_0\mathrm{sin}\varphi )\mathrm{sin}\left(\mathrm{\Delta }mt\right)`$ (104)
$`.+(I_{}R_0R_{}I_0)[\mathrm{cosh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)\mathrm{cos}\left(\mathrm{\Delta }mt\right)]\};`$ (105)
$`Im\left[A_0^{}(t)A_{}(t)\right]=Im\left[A_0^{}(0)A_{}(0)\right]{\displaystyle \frac{e^{\mathrm{\Gamma }t}}{2}}\times `$ (106)
$`\{[\mathrm{cosh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)+\mathrm{cos}\left(\mathrm{\Delta }mt\right)].`$ (107)
$`(R_{}\mathrm{cos}\varphi I_{}\mathrm{sin}\varphi )\mathrm{sinh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)+(R_{}\mathrm{sin}\varphi +I_{}\mathrm{cos}\varphi )\mathrm{sin}\left(\mathrm{\Delta }mt\right)`$ (108)
$`+(R_0\mathrm{cos}\varphi I_0\mathrm{sin}\varphi )\mathrm{sinh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)(R_0\mathrm{sin}\varphi +I_0\mathrm{cos}\varphi )\mathrm{sin}\left(\mathrm{\Delta }mt\right)`$ (109)
$`.(R_{}R_0+I_{}I_0)[\mathrm{cosh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)\mathrm{cos}\left(\mathrm{\Delta }mt\right)]\}`$ (110)
$`Re\left[A_0^{}(0)A_{}(0)\right]{\displaystyle \frac{e^{\mathrm{\Gamma }t}}{2}}\times `$ (111)
$`\{(R_{}\mathrm{sin}\varphi +I_{}\mathrm{cos}\varphi )\mathrm{sinh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)+(R_{}\mathrm{cos}\varphi I_{}\mathrm{sin}\varphi )\mathrm{sin}\left(\mathrm{\Delta }mt\right).`$ (112)
$`+(R_0\mathrm{sin}\varphi +I_0\mathrm{cos}\varphi )\mathrm{sinh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)+(R_0\mathrm{cos}\varphi I_0\mathrm{sin}\varphi )\mathrm{sin}\left(\mathrm{\Delta }mt\right)`$ (113)
$`.+(I_{}R_0R_{}I_0)[\mathrm{cosh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)\mathrm{cos}\left(\mathrm{\Delta }mt\right)]\}.`$ (114)
Similar formulas for $`Re\left[A_{}^{}(0)A_{}(0)\right]`$ and $`Im\left[A_{}^{}(0)A_{}(0)\right]`$ can be obtained from Eq. (97) and (106) by replacing “$`0`$” with “$``$”, respectively.
The time evolution formulas for the CP conjugate amplitude bilinears are:
$`|\overline{A}_\eta (t)|^2=|A_\eta (0)|^2e^{\mathrm{\Gamma }t}\{{\displaystyle \frac{1+R_\eta ^2+I_\eta ^2}{2}}\mathrm{cosh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right){\displaystyle \frac{1R_\eta ^2I_\eta ^2}{2}}\mathrm{cos}\left(\mathrm{\Delta }mt\right).`$ (134)
$`.+\eta _i[(R_\eta \mathrm{cos}\varphi I_\eta \mathrm{sin}\varphi )\mathrm{sinh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)+(R_\eta \mathrm{sin}\varphi +I_\eta \mathrm{cos}\varphi )\mathrm{sin}\left(\mathrm{\Delta }mt\right)]\};`$
$`Re\left[\overline{A}_0^{}(t)\overline{A}_{}(t)\right]=[Re\left[A_0^{}(0)A_{}(0)\right](R_{}R_0+I_{}I_0)Im\left[A_0^{}(0)A_{}(0)\right](I_{}R_0R_{}I_0)]{\displaystyle \frac{e^{\mathrm{\Gamma }t}}{2}}\times `$
$`\{[\mathrm{cosh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)+\mathrm{cos}\left(\mathrm{\Delta }mt\right)].`$
$`+{\displaystyle \frac{1}{R_{}^2+I_{}^2}}\left[(R_{}\mathrm{cos}\varphi I_{}\mathrm{sin}\varphi )\mathrm{sinh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)+(R_{}\mathrm{sin}\varphi +I_{}\mathrm{cos}\varphi )\mathrm{sin}\left(\mathrm{\Delta }mt\right)\right]`$
$`+{\displaystyle \frac{1}{R_0^2+I_0^2}}\left[(R_0\mathrm{cos}\varphi I_0\mathrm{sin}\varphi )\mathrm{sinh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)(R_0\mathrm{sin}\varphi +I_0\mathrm{cos}\varphi )\mathrm{sin}\left(\mathrm{\Delta }mt\right)\right]`$
$`.+{\displaystyle \frac{1}{\left(R_{}^2+I_{}^2\right)\left(R_0^2+I_0^2\right)}}(R_{}R_0+I_{}I_0)[\mathrm{cosh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)\mathrm{cos}\left(\mathrm{\Delta }mt\right)]\}`$
$`[Re\left[A_0^{}(0)A_{}(0)\right](I_{}R_0R_{}I_0)+Im\left[A_0^{}(0)A_{}(0)\right](R_{}R_0+I_{}I_0)]{\displaystyle \frac{e^{\mathrm{\Gamma }t}}{2}}\times `$
$`\{{\displaystyle \frac{1}{R_{}^2+I_{}^2}}[(R_{}\mathrm{sin}\varphi +I_{}\mathrm{cos}\varphi )\mathrm{sinh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)+(R_{}\mathrm{cos}\varphi I_{}\mathrm{sin}\varphi )\mathrm{sin}\left(\mathrm{\Delta }mt\right)].`$
$`{\displaystyle \frac{1}{R_0^2+I_0^2}}\left[(R_0\mathrm{sin}\varphi +I_0\mathrm{cos}\varphi )\mathrm{sinh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)+(R_0\mathrm{cos}\varphi I_0\mathrm{sin}\varphi )\mathrm{sin}\left(\mathrm{\Delta }mt\right)\right]`$
$`.{\displaystyle \frac{1}{\left(R_{}^2+I_{}^2\right)\left(R_0^2+I_0^2\right)}}(I_{}R_0R_{}I_0)[\mathrm{cosh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)\mathrm{cos}\left(\mathrm{\Delta }mt\right)]\};`$
$`Im\left[\overline{A}_0^{}(t)\overline{A}_{}(t)\right]=[Re\left[A_0^{}(0)A_{}(0)\right](I_{}R_0R_{}I_0)+Im\left[A_0^{}(0)A_{}(0)\right](R_{}R_0+I_{}I_0)]{\displaystyle \frac{e^{\mathrm{\Gamma }t}}{2}}\times `$
$`\{[\mathrm{cosh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)+\mathrm{cos}\left(\mathrm{\Delta }mt\right)].`$
$`+{\displaystyle \frac{1}{R_{}^2+I_{}^2}}\left[(R_{}\mathrm{cos}\varphi I_{}\mathrm{sin}\varphi )\mathrm{sinh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)+(R_{}\mathrm{sin}\varphi +I_{}\mathrm{cos}\varphi )\mathrm{sin}\left(\mathrm{\Delta }mt\right)\right]`$
$`+{\displaystyle \frac{1}{R_0^2+I_0^2}}\left[(R_0\mathrm{cos}\varphi I_0\mathrm{sin}\varphi )\mathrm{sinh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)+(R_0\mathrm{sin}\varphi +I_0\mathrm{cos}\varphi )\mathrm{sin}\left(\mathrm{\Delta }mt\right)\right]`$
$`.+{\displaystyle \frac{1}{\left(R_{}^2+I_{}^2\right)\left(R_0^2+I_0^2\right)}}(R_{}R_0+I_{}I_0)[\mathrm{cosh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)\mathrm{cos}\left(\mathrm{\Delta }mt\right)]\}`$
$`[Re\left[A_0^{}(0)A_{}(0)\right](R_{}R_0+I_{}I_0)Im\left[A_0^{}(0)A_{}(0)\right](I_{}R_0R_{}I_0)]{\displaystyle \frac{e^{\mathrm{\Gamma }t}}{2}}\times `$
$`\{{\displaystyle \frac{1}{R_{}^2+I_{}^2}}[(R_{}\mathrm{sin}\varphi +I_{}\mathrm{cos}\varphi )\mathrm{sinh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)(R_{}\mathrm{cos}\varphi I_{}\mathrm{sin}\varphi )\mathrm{sin}\left(\mathrm{\Delta }mt\right)].`$
$`{\displaystyle \frac{1}{R_0^2+I_0^2}}\left[(R_0\mathrm{sin}\varphi +I_0\mathrm{cos}\varphi )\mathrm{sinh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)(R_0\mathrm{cos}\varphi I_0\mathrm{sin}\varphi )\mathrm{sin}\left(\mathrm{\Delta }mt\right)\right]`$
$`.+{\displaystyle \frac{1}{\left(R_{}^2+I_{}^2\right)\left(R_0^2+I_0^2\right)}}(I_{}R_0R_{}I_0)[\mathrm{cosh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)\mathrm{cos}\left(\mathrm{\Delta }mt\right)]\};`$
$`Re\left[\overline{A}_0^{}(t)\overline{A}_{}(t)\right]=[Re\left[A_0^{}(0)A_{}(0)\right](R_{}R_0+I_{}I_0)Im\left[A_0^{}(0)A_{}(0)\right](I_{}R_0R_{}I_0)]{\displaystyle \frac{e^{\mathrm{\Gamma }t}}{2}}\times `$ (152)
$`\{[\mathrm{cosh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)+\mathrm{cos}\left(\mathrm{\Delta }mt\right)].`$
$`{\displaystyle \frac{1}{R_{}^2+I_{}^2}}\left[(R_{}\mathrm{cos}\varphi I_{}\mathrm{sin}\varphi )\mathrm{sinh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)+(R_{}\mathrm{sin}\varphi +I_{}\mathrm{cos}\varphi )\mathrm{sin}\left(\mathrm{\Delta }mt\right)\right]`$
$`+{\displaystyle \frac{1}{R_0^2+I_0^2}}\left[(R_0\mathrm{cos}\varphi I_0\mathrm{sin}\varphi )\mathrm{sinh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)+(R_0\mathrm{sin}\varphi +I_0\mathrm{cos}\varphi )\mathrm{sin}\left(\mathrm{\Delta }mt\right)\right]`$
$`.{\displaystyle \frac{1}{\left(R_{}^2+I_{}^2\right)\left(R_0^2+I_0^2\right)}}(R_{}R_0+I_{}I_0)[\mathrm{cosh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)\mathrm{cos}\left(\mathrm{\Delta }mt\right)]\}`$
$`[Re\left[A_0^{}(0)A_{}(0)\right](I_{}R_0R_{}I_0)+Im\left[A_0^{}(0)A_{}(0)\right](R_{}R_0+I_{}I_0)]{\displaystyle \frac{e^{\mathrm{\Gamma }t}}{2}}\times `$
$`\{{\displaystyle \frac{1}{R_{}^2+I_{}^2}}[(R_{}\mathrm{sin}\varphi +I_{}\mathrm{cos}\varphi )\mathrm{sinh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)(R_{}\mathrm{cos}\varphi I_{}\mathrm{sin}\varphi )\mathrm{sin}\left(\mathrm{\Delta }mt\right)].`$
$`+{\displaystyle \frac{1}{R_0^2+I_0^2}}\left[(R_0\mathrm{sin}\varphi +I_0\mathrm{cos}\varphi )\mathrm{sinh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)(R_0\mathrm{cos}\varphi I_0\mathrm{sin}\varphi )\mathrm{sin}\left(\mathrm{\Delta }mt\right)\right]`$
$`.+{\displaystyle \frac{1}{\left(R_{}^2+I_{}^2\right)\left(R_0^2+I_0^2\right)}}(I_{}R_0R_{}I_0)[\mathrm{cosh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)\mathrm{cos}\left(\mathrm{\Delta }mt\right)]\};`$
$`Im\left[\overline{A}_0^{}(t)\overline{A}_{}(t)\right]=[Re\left[A_0^{}(0)A_{}(0)\right](I_{}R_0R_{}I_0)+Im\left[A_0^{}(0)A_{}(0)\right](R_{}R_0+I_{}I_0)]{\displaystyle \frac{e^{\mathrm{\Gamma }t}}{2}}\times `$
$`\{[\mathrm{cosh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)+\mathrm{cos}\left(\mathrm{\Delta }mt\right)].`$
$`{\displaystyle \frac{1}{R_{}^2+I_{}^2}}\left[(R_{}\mathrm{cos}\varphi I_{}\mathrm{sin}\varphi )\mathrm{sinh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)+(R_{}\mathrm{sin}\varphi +I_{}\mathrm{cos}\varphi )\mathrm{sin}\left(\mathrm{\Delta }mt\right)\right]`$
$`+{\displaystyle \frac{1}{R_0^2+I_0^2}}\left[(R_0\mathrm{cos}\varphi I_0\mathrm{sin}\varphi )\mathrm{sinh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)+(R_0\mathrm{sin}\varphi +I_0\mathrm{cos}\varphi )\mathrm{sin}\left(\mathrm{\Delta }mt\right)\right]`$
$`.{\displaystyle \frac{1}{\left(R_{}^2+I_{}^2\right)\left(R_0^2+I_0^2\right)}}(R_{}R_0+I_{}I_0)[\mathrm{cosh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)\mathrm{cos}\left(\mathrm{\Delta }mt\right)]\}`$
$`[Re\left[A_0^{}(0)A_{}(0)\right](R_{}R_0+I_{}I_0)Im\left[A_0^{}(0)A_{}(0)\right](I_{}R_0R_{}I_0)]{\displaystyle \frac{e^{\mathrm{\Gamma }t}}{2}}\times `$
$`\{{\displaystyle \frac{1}{R_{}^2+I_{}^2}}[(R_{}\mathrm{sin}\varphi +I_{}\mathrm{cos}\varphi )\mathrm{sinh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)+(R_{}\mathrm{cos}\varphi I_{}\mathrm{sin}\varphi )\mathrm{sin}\left(\mathrm{\Delta }mt\right)].`$
$`+{\displaystyle \frac{1}{R_0^2+I_0^2}}\left[(R_0\mathrm{sin}\varphi +I_0\mathrm{cos}\varphi )\mathrm{sinh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)+(R_0\mathrm{cos}\varphi I_0\mathrm{sin}\varphi )\mathrm{sin}\left(\mathrm{\Delta }mt\right)\right]`$
$`.{\displaystyle \frac{1}{\left(R_{}^2+I_{}^2\right)\left(R_0^2+I_0^2\right)}}(I_{}R_0R_{}I_0)[\mathrm{cosh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }t}{2}}\right)\mathrm{cos}\left(\mathrm{\Delta }mt\right)]\}.`$
Similar formulas for $`Re\left[\overline{A}_{}^{}(0)\overline{A}_{}(0)\right]`$ and $`Im\left[\overline{A}_{}^{}(0)\overline{A}_{}(0)\right]`$ can be obtained from Eq. (152) and (152) by replacing “$`0`$” with “$``$”, respectively. |
warning/0002/nucl-th0002016.html | ar5iv | text | # Ground states of the Wick-Cutkosky model using light-front dynamics
## I Introduction
Recent experiments at Thomas Jefferson National Accelerator Facility have measured the $`A(Q^2)`$ structure function of the deuteron for momentum transfers up to 6 (GeV/c)<sup>2</sup> , and measurements for $`B(Q^2)`$ are planned. At such large momentum transfers, a relativistic description of the deuteron is required. One approach that gives such a description is light-front dynamics, which we will examine here. To separate the effects of the using light-front dynamics from the effects of the model, we choose to use the massive Wick-Cutkosky model. This is a “toy model” investigation, instead of the full nuclear theory calculation. Using this model, the light-front Hamiltonian approach is used to solve for the bound-state wavefunction. The results of our calculation can then be compared to other calculations done with the same model but different approaches. The simplest observable that can be compared is the relationship between the bound-state mass and the coupling constant.
The utility of the light-front dynamics was first discussed by Dirac . We start by expressing the four-vector $`x^\mu `$ in terms of the light-front variables $`x^\mu =(x^+,x^{},x^1,x^2)`$, where $`x^\pm =x^0\pm x^3`$. This is simply a change of variables, but an especially convenient one. Using this coordinate system and defining the commutation relations at equal light-front time ($`x^+=t_{\text{LF}}`$), we obtain a light-front Hamiltonian . The Hamiltonian is used in the light-front Schrödinger equation to solve for the ground state.
There are many desirable features of the light-front dynamics and the use of light-front coordinates. First of all, high-energy experiments are naturally described using light-front coordinates. The wave front of a beam of high-energy particles traveling in the (negative) three-direction is defined by a surface where $`x^+`$ is (approximately) constant. Such a beam can probe the wavefunction of a target described in terms of light-front variables : the Bjorken $`x`$ variable used to describe high-energy experiments is simply the ratio of the plus momentum of the struck constituent particle to the total plus momentum ($`p^+`$) of the bound state. Secondly, the vacuum for a theory with massive particles can be very simple on the light front. This is because all massive particles and anti-particles have positive plus momentum, and the total plus momentum is a conserved quantity. Thus, the naïve vacuum (with $`p^+=0`$) is empty, and diagrams that couple to this vacuum are zero. This greatly reduces the number of non-trivial light-front time-ordered diagrams. Thirdly, the generators of boosts in the one, two, and plus directions are kinematical, meaning they are independent of the interaction. Thus, even when the Hamiltonian is truncated, the wavefunctions will transform correctly under boosts. Thus, light-front dynamics is useful for describing form factors at high momentum transfers. A drawback of the light-front formalism is that the Hamiltonian is not manifestly rotationally invariant, since the generators of rotations about the one and two directions are dynamical. A study of the effects of the loss of rotational invariance of the excited states in the model being used here was made in Ref. , which shows that there is less breaking of the degeneracy in the spectrum when higher order potentials are used as opposed to lower order potentials.
There are other approaches that can be used to obtain relativistic wavefunctions and bound-state energies, including the Feynman-Schwinger representation (FSR) of the two-particle Green’s function and the Bethe-Salpeter equation (BSE) . The FSR is useful since it can be constructed so it is equivalent to the Bethe-Salpeter equation using a kernel where all two-particle-to-two-particle ladder and crossed ladder diagrams are included. However, it requires a path integral to be done numerically, so it is computationally intensive to obtain an accurate answer. The FSR result is to be considered the full solution that the Bethe-Salpeter and Hamiltonian equations approximate. The BSE can be solved much quicker than the FSR, however a truncation of the BSE kernel is required, causing the BSE results to differ from the FSR results. It is well known that any finite truncation of the kernel yields bound-state wavefunctions with problems, such as the incorrect one-body limit . For our scalar model the potential is always attractive, so truncation of the kernel gives binding energies that are too small . Another approach is explicitly covariant light-front dynamics , where manifest covariance is kept at the price of using a null-plane whose orientation is not fixed.
Here we study light-front dynamics because of its close connection to experimental observables. Using field theory, light-front potentials can be derived that give results physically equivalent to those of the Bethe-Salpeter equation. We define “physically equivalent” in section III C. Depending on the diagrams used to construct the potential, one can argue that certain potentials are physically equivalent to the Bethe-Salpeter equations with certain kernels. The best that these potentials can do is reproduce the results of the corresponding Bethe-Salpeter equations. On the other hand, light-front potentials can be constructed that attempt to incorporate non-perturbative physics. It is possible that these potentials can give results that agree with the full theory results better than the Bethe-Salpeter equation using low-order kernels. We will consider both types of potentials in this paper, and see how well they perform.
A brief discussion about the approximation used is in order. It is well known that the vacuum of the full Wick-Cutkosky model is unstable due to the cubic coupling which provides the interaction. However, when the bound-state calculation is restricted to the two-particle sector, the quenched approximation is used, and the self-energy and vertex-renormalization diagrams are neglected, the theory has a well defined ground state. In this paper, we compare the results of our light-front Hamiltonian calculation to the Bethe-Salpeter and FSR calculations, both of which use the same approximations. The use of these simplifying approximations allows us to highlight the differences between the various approaches. The inclusion of the self-energy diagrams and counterterms for the light-front Hamiltonian and for the FSR will not be discussed here.
### A Outline of the paper
The objective of this paper is to obtain the bound-state energy for the ground state in our theory. In our model, neutral scalar nucleons interact via a Yukawa interaction, which is mediated by a neutral scalar meson. The light-front Hamiltonian derived from the Lagrangian is used in the light-front Schrödinger equation for a two-nucleon bound state. The rules for light-front time-ordered perturbation theory (LFTOPT) are then derived for this Schrödinger equation, along with the Feynman rules for the effective potential. All of this is discussed in more detail in section II.
The full potential for the light-front Schrödinger equation is given by an infinite sum of diagrams. Using the LFTOPT rules in section III, we derive the one-boson-exchange (OBE) and two-boson-exchange (TBE) potentials, where the diagrams that give rise to mass and vertex renormalization are not included. The two-boson-exchange stretched-box (TBE:SB) diagrams, a subset of the TBE diagrams, are used to construct the TBE:SB potential, and the utility of this potential is commented on. After the discussion of the perturbative potentials, we define several potentials in section IV which attempt to incorporate non-perturbative physics in a OBE type of potential. Three potentials are obtained by approximating the OBE potential directly (giving the symmetrized mass, instantaneous, and retarded potentials), and one potential is obtained by approximating the Bethe-Salpeter equation and reducing it to a three-dimensional equation (giving the modified-Green’s-function potential).
In section V, the potentials are used to numerically obtain the spectra, the coupling constant versus bound-state mass curves. This is done by solving the light-front Schrödinger equation with each of the truncated potentials (OBE, TBE, and TBE:SB), and the approximate potentials (symmetrized mass, instantaneous, retarded, and modified-Green’s-function). We compare these results to those in the literature obtained with other approaches.
We summarize our findings in section VI. Only a few low-order terms from the perturbative potentials are needed to approximate the results of the physically equivalent Bethe-Salpeter equations . Since the light-front potentials are calculated without including the mass and vertex renormalization diagrams, the Bethe-Salpeter kernels used also do not include those diagrams. Since the interaction in the Wick-Cutkosky model is strictly attractive, the spectra calculated using the perturbative potentials will underestimate the the binding energy compared with the spectra for the physically equivalent Bethe-Salpeter equations. As progressively higher-order terms in the potential are kept, the spectra calculated will agree better with the BSE spectra. However, when a truncated kernel is used in the Bethe-Salpeter equation, the solutions obtained are known to be a poor approximation of the full solution , which for our theory (where no renormalization graphs are kept) is given by the Feynman-Schwinger representation of the Green’s function . Hence, our truncated potentials cannot give the full results. The non-perturbative approximations of the potential give spectra that more closely match the spectrum of the full solution. This suggests that these approximations are reproducing the physics more accurately than the perturbative potentials.
The conventions and notations employed in this paper are summarized in Appendix A. The conversion of the light-front Schrödinger equation into matrix form suitable for numerical evaluation is discussed in Appendix B. Azimuthal-angle integrations of the OBE potentials, which help simplify the evaluation of the bound-state wavefunctions, are given in Appendix C. The loop integrations and azimuthal-angle integrations needed for the TBE potentials are discussed in Appendix D. A check of the validity of the uncrossed approximation used in section IV D is done in Appendix E.
This study is related to the work of Sales et al. , who computed bound states using the light-front Hamiltonian with the OBE and TBE:SB approximations, and compared to the ladder Bethe-Salpeter equation results. Here, we consider also the effect of including the crossed graph part of the TBE potential as well as several non-perturbative approximations.
## II Our model
We consider an isospin doublet of two uncharged scalars $`\varphi =(\varphi _1,\varphi _2)`$ with mass $`M`$ (which we will refer to as nucleons), that couple to a third, uncharged scalar $`\chi `$ with mass $`\mu `$ (which we will refer to as a meson) by a $`\varphi ^2\chi `$ interaction. This is the massive extension of the Wick-Cutkosky model , which has been used on the light front to study scattering states as well as bound states . The Lagrangian is
$``$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(_\mu \varphi ^\mu \varphi M^2\varphi ^2\right)+{\displaystyle \frac{1}{2}}\left(_\mu \chi ^\mu \chi \mu ^2\chi ^2\right)+g{\displaystyle \frac{M}{2}}\varphi ^2\chi ,`$ (1)
where $`g`$ is a dimensionless coupling constant and $`\varphi ^2=\varphi _1^2+\varphi _2^2`$.
### A Light-front Hamiltonian
To obtain the light-front Hamiltonian from the Lagrangian in Eq. (1), we follow the approach used by Miller and many others (see the review ) to write the light-front Hamiltonian ($`P^{}`$) as the sum of a free, non-interacting part and a term containing the interactions. We use the conventions given in Appendix A. The operators we use can be expressed in terms of Fock space operators since for this theory in light-front dynamics, the physical vacuum is the Fock space vacuum, and thus the Hilbert space is simply the Fock space. The Hamiltonian is obtained by using the energy-momentum tensor in
$`P^\mu ={\displaystyle \frac{1}{2}}{\displaystyle 𝑑x^{}d^2x_{}T^{+\mu }(x^+=0,x^{},𝒙_{})},`$ (2)
The usual relations determine $`T^{+\mu }`$, with
$`T^{\mu \nu }=g^{\mu \nu }+{\displaystyle \underset{r}{}}{\displaystyle \frac{}{(_\mu \varphi _r)}}^\nu \varphi _r,`$ (3)
in which the degrees of freedom (the fields $`\varphi `$ and $`\chi `$) are labeled by $`\varphi _r`$.
It is worthwhile to consider the limit in which the interactions between the fields are removed. This will allow us to define the free Hamiltonian $`P_0^{}`$ and to display the necessary commutation relations. The energy-momentum tensor of the non-interacting fields is defined as $`T_0^{\mu \nu }`$. Use of Eq. (3) leads to the result
$`T_0^{\mu \nu }`$ $`=`$ $`^\mu \varphi ^\nu \varphi {\displaystyle \frac{g^{\mu \nu }}{2}}\left[_\sigma \varphi ^\sigma \varphi M^2\varphi ^2\right]+^\mu \chi ^\nu \chi {\displaystyle \frac{g^{\mu \nu }}{2}}\left[_\sigma \chi ^\sigma \chi \mu ^2\chi ^2\right],`$ (4)
with
$`T_0^+=\mathbf{}_{}\varphi \mathbf{}_{}\varphi +M^2\varphi ^2+\mathbf{}_{}\chi \mathbf{}_{}\chi +\mu ^2\chi ^2.`$ (5)
The scalar nucleon fields can be expressed in terms of creation and destruction operators:
$`\varphi _i(x)`$ $`=`$ $`{\displaystyle \frac{d^2k_{}dk^+\theta (k^+)}{(2\pi )^{3/2}\sqrt{2k^+}}\left[a_i(𝒌)e^{ikx}+a_i^{}(𝒌)e^{ikx}\right]},`$ (6)
where $`i=1,2`$ is a particle index, $`kx=\frac{1}{2}(k^{}x^++k^+x^{})𝒌_{}𝒙_{}`$ with $`k^{}=\frac{M^2+𝒌_{}^2}{k^+}`$, and $`𝒌(k^+,𝒌_{})`$. Note that $`k^{}`$ is such that the particles are on the mass shell, which is a consequence of using a Hamiltonian theory. The $`\theta `$ function restricts $`k^+`$ to positive values. Likewise, the scalar meson field is given by
$`\chi (x)`$ $`=`$ $`{\displaystyle \frac{d^2k_{}dk^+\theta (k^+)}{(2\pi )^{3/2}\sqrt{2k^+}}\left[a_\chi (𝒌)e^{ikx}+a_\chi ^{}(𝒌)e^{ikx}\right]},`$ (7)
where $`k^{}=\frac{\mu ^2+𝒌_{}^2}{k^+}`$, so that the mesons are also on the mass shell. The non-vanishing commutation relations are
$`[a_\alpha (𝒌),a_\alpha ^{}(𝒌^{})]`$ $`=`$ $`\delta (𝒌_{}𝒌_{}^{})\delta (k^+k^+),`$ (8)
where $`\alpha =1,2,\chi `$ is a particle index. The commutation relations are defined at equal light-front time, $`x^+=0`$. It is useful to define
$`\delta ^{(2,+)}(𝒌𝒌^{})`$ $``$ $`\delta (𝒌_{}𝒌_{}^{})\delta (k^+k^+),`$ (9)
which will be used throughout this paper.
We write a ket in the two-distinguishable-particle sector of the Fock space as
$`|k_1,k_2`$ $`=`$ $`a_1^{}(k_1)a_2^{}(k_2)|0.`$ (10)
This implies that the identity operator in this Fock space sector can be written as
$`I_2`$ $`=`$ $`{\displaystyle d^2k_{1,}𝑑k_1^+d^2k_{2,}𝑑k_2^+|k_1,k_2k_1,k_2|}.`$ (11)
The derivatives appearing in the quantity $`T_0^+`$ are evaluated and then one sets $`x^+`$ to 0 to obtain the result
$`P_0^{}`$ $`=`$ $`{\displaystyle _k}\left[{\displaystyle \frac{M^2+𝒌_{}^2}{k^+}}\left(a_1^{}(k)a_1(k)+a_2^{}(k)a_2(k)\right)+{\displaystyle \frac{\mu ^2+𝒌_{}^2}{k^+}}a_\chi ^{}(k)a_\chi (k)\right],`$ (12)
with $`_k=d^2k_{}𝑑k^+\theta (k^+)`$. Eq. (12) has the interpretation of an operator that counts the light-front energy $`k^{}`$ (which is $`\frac{M^2+𝒌_{}^2}{k^+}`$ for the nucleons and $`\frac{\mu ^2+𝒌_{}^2}{k^+}`$ for the mesons) of all of the particles.
We now consider the interacting part of the Lagrangian, $`_I`$. An analysis similar to that for the non-interacting parts yields the interacting part of the light-front Hamiltonian $`P_I^{}`$;
$`P_I^{}`$ $`=`$ $`{\displaystyle \underset{i=1,2}{}}{\displaystyle \frac{M}{2}}{\displaystyle _k}{\displaystyle _k^{}}{\displaystyle \frac{1}{(2\pi )^{3/2}\sqrt{2k^+k^+(k^++k^+)}}}`$ (15)
$`\times \{[2a_i^{}(k+k^{})a_\chi (k^{})a_i(k)+a_\chi ^{}(k+k^{})a_i(k^{})a_i(k)]`$
$`+\text{Hermitian conjugate}\}.`$
The interaction Hamiltonian is self-adjoint since the Hilbert space is the Fock space. The total light-front Hamiltonian is given by $`P^{}=P_0^{}+gP_I^{}`$.
### B Hamiltonian bound-state equations
We will be studying the bound states of two distinguishable nucleons. The technology of time-ordered (old-fashioned) perturbation theory is used to construct the light-front time-ordered perturbation theory (LFTOPT) for our Hamiltonian. We start with the light-front Schrödinger equation in the full Fock space,
$`\left(P_0^{}+gP_I^{}\right)|\psi _F^{\text{GS}}`$ $`=`$ $`|\psi _F^{\text{GS}}P_{\text{GS}}^{},`$ (16)
where $`P_0^{}+gP_I^{}`$ is the Hamiltonian in the full Fock-space basis, $`|\psi _F^{\text{GS}}`$ is the ground-state wavefunction in the full Fock space, and $`P_{\text{GS}}^{}`$ is the light-front energy of that state. Recall that $`P_0^{}`$, the non-interacting part of the Hamiltonian, is diagonal in the momentum basis, while $`P_I^{}`$, which contains the interaction, has only off-diagonal elements.
A serious drawback of this equation is that the wavefunction $`\psi _F^{\text{GS}}`$ has support from infinitely many sectors of the Fock space, since $`P_I^{}`$ changes the total number of particles. However, the components of the wavefunction with many particles will be small compared to the two-particle component if the coupling constant is not too large. We will construct the two-particle light-front Schrödinger equation which the two-particle component of the wavefunction satisfies. From this construction, we will obtain the rules for the LFTOPT.
We start by introducing the projection operators $`𝒫`$ and $`𝒬`$. The operator $`𝒫`$ projects out the sector of Fock space with two distinguishable nucleons and no mesons, while $`𝒬=I𝒫`$ projects out all the other sectors. We define
$`𝒫|\psi _F^{\text{GS}}`$ $``$ $`|\psi ^{\text{GS}}`$ (17)
$`𝒬|\psi _F^{\text{GS}}`$ $``$ $`|\psi _Q^{\text{GS}},`$ (18)
so that $`|\psi _F^{\text{GS}}=|\psi ^{\text{GS}}+|\psi _Q^{\text{GS}}`$. Since the free Hamiltonian does not change the number of particles, $`[𝒫,P_0^{}]=[𝒬,P_0^{}]=0`$. The interaction Hamiltonian changes the particle number, so it cannot connect the two-particle sector to itself, thus $`𝒫P_I^{}𝒫=0`$.
Using these projection operators, Eq. (16) can be broken up into two parts,
$`P_0^{}|\psi ^{\text{GS}}+g𝒫P_I^{}𝒬|\psi _Q^{\text{GS}}`$ $`=`$ $`|\psi ^{\text{GS}}P_{\text{GS}}^{}`$ (19)
$`\left(P_0^{}+g𝒬P_I^{}𝒬\right)|\psi _Q^{\text{GS}}+g𝒬P_I^{}𝒫|\psi ^{\text{GS}}`$ $`=`$ $`|\psi _Q^{\text{GS}}P_{\text{GS}}^{}.`$ (20)
Eliminating the $`|\psi _Q^{\text{GS}}`$ and using the expression of the identity given in Eq. (11) we obtain the two-particle effective light-front Schrödinger equation
$`{\displaystyle d^2p_{1,}𝑑p_1^+d^2p_{2,}𝑑p_2^+𝒌_1,𝒌_2|\left[P_0^{}+V(g,P_{\text{GS}}^{})\right]|𝒑_1,𝒑_2𝒑_1,𝒑_2|\psi ^{\text{GS}}}`$ (21)
$`=𝒌_1,𝒌_2|\psi ^{\text{GS}}P_{\text{GS}}^{},`$ (22)
where $`P_0^{}`$ and the potential $`V`$ act in the two-nucleon basis. The two-particle potential is given by
$`V(g,P^{})`$ $`=`$ $`g^2𝒫P_I^{}{\displaystyle \frac{𝒬}{P^{}P_0^{}g𝒬P_I^{}𝒬}}P_I^{}𝒫.`$ (23)
Note that Eq. (22) is similar to Eq. (16), except for two main differences. Here we have a two-nucleon wavefunction, which makes it simpler. However, the potential is light-front energy dependent, which makes it more complicated.
The denominator in the definition of the potential is non-diagonal in the full Fock space, so the matrix inversion that it represents is highly non-trivial. This problem is avoided by expanding the inversion in powers of the coupling constant $`g`$ to get
$`V(g,P^{})`$ $`=`$ $`𝒫P_I^{}\left[{\displaystyle \frac{g^2𝒬}{P^{}P_0^{}}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\left(P_I^{}{\displaystyle \frac{g𝒬}{P^{}P_0^{}}}\right)^n\right]P_I^{}𝒫.`$ (24)
This can be simplified further by noting that in the two-nucleon sector of our theory, every meson emitted must be absorbed, so there must be an even number of interactions. Thus, the full potential can be written as the sum of $`n`$ meson exchange potentials,
$`V(P^{},g)`$ $`=`$ $`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}g^{2n}V_{(2n)}(P^{}),`$ (25)
where $`V_{(2n)}`$ is the potential due to the exchange of $`n`$ mesons, given by
$`V_{(2n)}(P^{})`$ $`=`$ $`𝒫\left(P_I^{}{\displaystyle \frac{𝒬}{P^{}P_0^{}}}\right)^{2n1}P_I^{}𝒫.`$ (26)
It is easy to see how to write a sum of diagrams for the potential when one says what Eq. (26) represents in words. We start off with two particles, then the interaction occurs. There are two possibilities of what can happen; nucleon 1 or 2 can emit a meson. Each possibility has a separate diagram. After the interaction, there is propagation with the light-front Green’s function,
$`G_{\text{LF}}(P^{})`$ $`=`$ $`{\displaystyle \frac{1}{P^{}P_0^{}}},`$ (27)
until another interaction occurs, and so on. We simply sum up all of the possible orderings of the interaction to get the full potential. The $`n^{\text{th}}`$ order potential is simply the sum of all possible diagrams with $`n`$-meson exchanges.
Each intermediate state in Eq. (26) has more than two particles, so the diagrams are two-particle irreducible with respect to the two-particle Green’s function $`G_{2\text{LF}}=𝒫G_{\text{LF}}𝒫`$. We can represent $`G_{2\text{LF}}`$ by its diagonal matrix elements,
$`G_{2\text{LF}}(𝒌_1,𝒌_2;P^{})={\displaystyle \frac{1}{P^{}k_1^{}k_2^{}}}.`$ (28)
In the diagrams we draw, the nucleons will be represented by solid lines and the mesons by the dashed lines. Although the states we will be considering consist of two distinguishable nucleons, we will not label the nucleon lines. We will be using the quenched approximation (so there are no nucleon loops) and neglect the mass and vertex renormalization diagrams (so the physical masses and coupling constant are used, and each meson emitted from one nucleon must be absorbed by the other nucleon). It is not expected that these restrictions will lead to qualitatively different results than the true full solution when the states are not too deeply bound. The quenched approximation is reasonable when the masses of the nucleon fields are large compared to the binding energy. Use of the physical masses and coupling constant are reasonable as well when the momenta are not too large.
We stress again that we will compare various truncations of the light-front Hamiltonian to other calculations which do not include renormalization diagrams. This is because we want to determine the effect of truncation on the light-front Hamiltonian. The differences between our calculation and those which include the self-energy graphs , which may be large for deeply-bound states, are not considered here.
The rules for drawing the $`n`$-meson-exchange graphs that correspond to this approximation are:
1. Draw all topologically distinct time-ordered diagrams with $`n`$ mesons. Use solid lines for the nucleons and dashed lines for the mesons.
2. Delete all graphs which couple particles to the vacuum. In the massive theory we consider here, these diagrams always vanish since the vacuum has zero plus momentum, and massive particles always have positive plus momentum.
3. Our quenched approximation and use of the physical masses and coupling constant requires us to delete all graphs that have nucleon loops or have mesons that are emitted and absorbed from the same nucleon.
4. Delete all other graphs which are not allowed in the particular approximation that is being considered. For example, consider the potential from the the Hamiltonian theory that can be obtained from the ladder Bethe-Salpeter equation. That potential will not have any graphs where the meson lines cross.
Once the diagrams are drawn, we use the following rules to convert the sum of diagrams into the potential $`𝒌_1,𝒌_2|V_{(2n)}(P^{})|𝒑_1,𝒑_2`$:
1. Overall factor of $`\frac{\delta ^{(2,+)}(𝒌_1+𝒌_2𝒑_1𝒑_2)}{2(2\pi )^3\sqrt{k_1^+k_2^+p_1^+p_2^+}}`$. This delta function says that the total light-front three-momentum is conserved. We define the light-front three-momentum $`𝑷𝒌_1+𝒌_2`$.
2. To each internal line, assign a light-front three-momentum $`𝒒_i`$ where $`i=1,2,\mathrm{},N`$ and $`N`$ is the number of internal lines. The light-front energy for particle $`i`$ with mass $`m_i`$ is $`q_i=\frac{m_1^2+𝒒_{i,}^2}{q_i^+}`$. It is useful to define $`z_i=q_i^+/P^+`$.
3. A factor of $`\frac{\theta (z_i^+)}{z_i^+}`$ for each internal line.
4. An extra factor of $`\frac{M^2}{P^+P^{}}`$ for each internal meson line.
5. A factor of $`\frac{P^{}}{\left(P^{}_iq_i^{}\right)}`$ between consecutive vertices, where the sum is over only the particles that exist in the intermediate time between those vertices.
6. Use light-front three-momentum conservation to eliminate all the independent momenta.
7. Integrate with $`\frac{d^2q_{i,}dz_i}{2P^+P^{}(2\pi )^3}`$ over all remaining free internal momenta.
8. Symmetry factor of $`\frac{1}{2}`$ when two nucleons are created or destroyed at the same time.
With these rules, one can calculate the effective potential for any order.
### C Further development of the light-front Schrödinger equation
Once the potential is calculated, we can plug it into Eq. (22), which we as
$`{\displaystyle d^2p_{1,}𝑑p_1^+d^2p_{2,}𝑑p_2^+𝒌_1,𝒌_2|\left[P_0^{}+V(g(P^{}),P^{})\right]|𝒑_1,𝒑_2𝒑_1,𝒑_2|\psi ^{\text{GS}}}`$ (29)
$`=𝒌_1,𝒌_2|\psi ^{\text{GS}}P^{},`$ (30)
where $`P^{}`$ is an arbitrary light-front energy and $`g(P^{})`$ is the coupling constant which yields the bound-state wavefunction with $`P^{}`$ as the bound-state energy. We call this $`g(P^{})`$ the spectrum of the light-front Schrödinger equation for the corresponding wavefunction.
The total momentum $`𝑷=𝒌_1+𝒌_2`$ is conserved by the potential given in Eq. (25), so the wavefunction in Eq. (30) can be parameterized by the total momentum. To make the calculations easier later, we choose to be in the center-of-momentum frame, where the components of the total momentum can be written as $`𝑷_{}=0`$ and $`P^+=P^{}=E`$. The ground-state energy, $`E`$, is the same as the mass of the bound state. In terms of the binding energy $`B`$, $`E=2MB`$. In the center-of-momentum frame, the ground-state wavefunction is parameterized by $`E`$, so we can define
$`𝒌_1,𝒌_2|\psi _M^{\text{GS}}`$ $`=`$ $`\delta ^{(2,+)}(𝒌_1+𝒌_2𝑷)\psi ^{\text{GS}}(𝒌_1)`$ (31)
$`𝒌_1,𝒌_2|V_{(2n)}(P^{})|𝒑_1,𝒑_2`$ $`=`$ $`\delta ^{(2,+)}(𝒌_1+𝒌_2𝒑_1𝒑_2)V_{(2n)}(E;𝒌_1;𝒑_1).`$ (32)
With these, Eq. (30) effectively becomes a one-particle equation, where particle 2’s momentum is determined by $`𝒌_2=𝑷𝒌_1`$. The minus component (the light-front energy) of particle 2 is defined by the requirement that the particle 2 is on mass shell, so $`k_2^{}=(M^2+𝒌_{2,}^2)/k_2^+`$. We also define $`xk_1^+/P^+=x_{Bj}`$, where $`x_{Bj}`$ is the Bjorken $`x`$ variable, so that $`k_2^+/P^+=1x`$. Likewise, we write the Bjorken variables that correspond to the momenta $`𝒑_1`$ and $`𝒒_1`$ as $`yp_1^+/P^+`$ and $`zq_1^+/P^+`$.
Using Eqs. (31), (32), and the fact that the plus momentum of both nucleons is positive, we can write the light-front Schrödinger equation Eq. (30) as
$`{\displaystyle d^2p_{1,}_0^E𝑑p_1^+V(g(E),E;𝒌_1;𝒑_1)\psi ^{\text{GS}}(𝒑_1)}`$ $`=`$ $`\psi ^{\text{GS}}(𝒌_1)(Ek_1^{}k_2^{}).`$ (33)
It useful to convert from light-front coordinates $`𝒌_1=(k_1^+,𝒌_{})`$ to equal-time coordinates $`𝒌_{\text{ET}}=(𝒌_{},k^3)`$, using an implicit definition of $`k^3`$
$`k_1^+`$ $`=`$ $`{\displaystyle \frac{E}{2k^0(𝒌_{\text{ET}})}}\left[k^0(𝒌_{\text{ET}})+k^3\right]`$ (34)
$`k^0(𝒌_{\text{ET}})`$ $`=`$ $`\sqrt{M^2+𝒌_{\text{ET}}^2}.`$ (35)
Often the explicit dependence of $`k^0`$ on $`𝒌_{\text{ET}}`$ will not be shown. It is worth emphasizing that this is just a convenient change of variables; $`\psi ^{\text{GS}}(𝒌_{\text{ET}})`$ is not the usual equal-time ground-state wavefunction. With this transformation, we can express $`k_1^\pm `$ and $`k_2^\pm `$ as
$$k_1^+=k_2^{}\left(\frac{E}{2k^0}\right)^2=\left(1+\frac{k^3}{k^0}\right)\frac{E}{2},$$
(37)
$$k_2^+=k_1^{}\left(\frac{E}{2k^0}\right)^2=\left(1\frac{k^3}{k^0}\right)\frac{E}{2}.$$
(38)
Using these, Eq. (33) becomes
$`{\displaystyle d^3p_{\text{ET}}\frac{2p_1^+p_2^+}{p^0}V(g(E),E,𝒌_{\text{ET}};𝒑_{\text{ET}})\psi ^{\text{GS}}(𝒑_{\text{ET}})}`$ $`=`$ $`\psi ^{\text{GS}}(𝒌_{\text{ET}})\left[E^2(2k^0)^2\right].`$ (39)
Now consider the exchange of the particle labels 1 and 2. This causes
$`𝒌_{1,}`$ $``$ $`𝒌_{2,}=𝒌_{1,}`$ (40)
$`k_1^+`$ $``$ $`k_2^+=Ek_1^+,`$ (41)
which means that $`k^3`$ as defined in Eq. (34) transforms as $`k^3k^3`$, so $`𝒌_{\text{ET}}𝒌_{\text{ET}}`$. Consequently, exchange of particle labels 1 and 2 is the same as parity for $`𝒌_{\text{ET}}`$.
Since the two nucleons are identical except for the particle label, the effective potential commutes with parity to all orders in $`g^2`$. Furthermore, the light-front Hamiltonian is explicitly invariant under rotations about the three-axis. These considerations allow us to classify the wavefunctions as eigenfunctions of parity and the three-component of the angular momentum operator. The ground state will have even parity and be invariant under rotation about the three-axis, so we can write
$`\psi ^{\text{GS}}(𝒌_{\text{ET}})=\psi ^{\text{GS}}(k_{\text{ET}},\theta _k,\varphi _k)=\psi ^{\text{GS}}(k_{\text{ET}},\theta _k)=\psi ^{\text{GS}}(k_{\text{ET}},\pi \theta _k),`$ (42)
where $`\theta _k`$ and $`\varphi _k`$ are the polar and azimuthal angles, respectively, for the vector $`𝒌_{\text{ET}}`$.
The cylindrical symmetry and parity of the wavefunction can be used to rewrite Eq. (39) as
$`{\displaystyle _0^{\mathrm{}}}𝑑p_{\text{ET}}{\displaystyle _0^{\pi /2}}𝑑\theta _p{\displaystyle \frac{2p_1^+p_2^+p_{\text{ET}}^2\mathrm{sin}\theta _p}{p^0}}V^+(k_{\text{ET}},\theta _k;p_{\text{ET}},\theta _p)\psi ^{\text{GS}}(p_{\text{ET}},\theta _p)`$ (43)
$`=\psi ^{\text{GS}}(k_{\text{ET}},\theta _k)\left[E^24(k^0)^2\right],`$ (44)
where
$`V^+(k_{\text{ET}},\theta _k;p_{\text{ET}},\theta _p)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[V(k_{\text{ET}},\theta _k;p_{\text{ET}},\theta _p)+V(k_{\text{ET}},\theta _k;p_{\text{ET}},\pi \theta _p)\right]`$ (45)
$`V(k_{\text{ET}},\theta _k;p_{\text{ET}},\theta _p)`$ $`=`$ $`{\displaystyle _0^{2\pi }\frac{d\varphi _kd\varphi _p}{2\pi }V(𝒌_{\text{ET}};𝒑_{\text{ET}})}.`$ (46)
We call $`V(k_{\text{ET}},\theta _k;p_{\text{ET}},\theta _p)`$ the azimuthal-angle-averaged potential. The light-front vectors can also be expressed in terms of the azimuthal angle $`\varphi `$, where $`𝒌=(k^+,k_{},\varphi _k)`$, which allows the azimuthal-angle-averaged potential to be written in light-front coordinates
$`V(k_1^+,k_{1,};p_1^+,p_{1,})`$ $`=`$ $`{\displaystyle _0^{2\pi }\frac{d\varphi _kd\varphi _p}{2\pi }V(𝒌_1;𝒑_1)}.`$ (47)
All of the simplifications of Eq. (44) based on physical considerations have been addressed. However, further rearrangements need to be done before Eq. (44) is fit to be solved on the computer. Since these involve only numerical techniques, they are relegated to Appendix B.
## III Perturbative Potentials
A truncation must be made of the expansion of the potential given in Eq. (25), since it is not feasible to calculate the infinite sum of graphs for the potential in this Hamiltonian theory. In this paper we consider three truncations of the potential derived from the field theory. First the OBE potential and the TBE potentials are calculated. Then, we note that a subset of the TBE diagrams, the stretched-box diagrams, correspond to the truncated potential derived from the ladder Bethe-Salpeter equation. Thus, three truncated potentials are obtained that have a physical interpretation.
The matrix elements of these potentials are written in the two-particle momentum basis, denoting the momentum of the incoming particles by $`𝒑_1`$ and $`𝒑_2`$, and the outgoing particles $`𝒌_1`$ and $`𝒌_2`$. For simplicity, we choose to work in the center-of-momentum frame. By inspecting the rules for converting a light-front time-ordered diagram into a potential given in section II B, and looking at Eq. (32), we find that each piece of the effective-one-particle potential $`V(E;𝒌_1,𝒑_1)`$ is proportional to
$`{\displaystyle \frac{E}{2(2\pi )^3\sqrt{k_1^+k_2^+p_1^+p_2^+}}}.`$ (48)
(An extra factor of $`E`$ is included to simplify later equations.) This term will by suppressed in all of the potentials written in this paper.
### A OBE Potential
We start by drawing all the allowed and non-vanishing time-ordered diagrams with one meson exchange. These diagrams are shown in Fig. 1. The light-front time-ordered perturbation theory rules given in section II B are used to calculate the potential due to OBE potential,
$`V_{\text{OBE}}(E;𝒌_1;𝒑_1)`$ $`=`$ $`\left({\displaystyle \frac{M}{E}}\right)^2[{\displaystyle \frac{\theta (xy)/|xy|}{Ep_1^{}k_2^{}\omega ^{}(𝒌_1𝒑_1)}}`$ (50)
$`+{\displaystyle \frac{\theta (yx)/|yx|}{Ek_1^{}p_2^{}\omega ^{}(𝒑_1𝒌_1)}}].`$
We have introduced the notation that meson with light-front three-momentum $`𝒒`$ has a light-front energy given by
$`\omega ^{}(𝒒)`$ $`=`$ $`{\displaystyle \frac{\mu ^2+𝒒_{}^2}{q^+}}.`$ (51)
The azimuthal-angle average of $`V_{\text{OBE}}`$ is discussed in Appendix C.
The potential given in Eq. (50) can also be used for scattering states. In that case, $`E=k_1^{}+k_2^{}=p_1^{}+p_2^{}`$, which allows the potential to be written as
$`V_{\text{OBE}}(E_{\text{scat}};𝒌_1;𝒑_1)`$ $`=`$ $`{\displaystyle \frac{M^2/E_{\text{scat}}}{(k_1p_1)^2\mu ^2}}.`$ (52)
The scattering potential is the same as the usual equal-time OBE potential. This must be the case, since the scattering potential is also given by covariant Feynman diagrams, which have the same form independent of the form of dynamics.
Returning to the bound-state regime, we note that the OBE potential can easily be written in terms of the equal-time coordinates. A reorganization of Eq. (50) yields
$`V_{\text{OBE}}(E;𝒌_{\text{ET}};𝒑_{\text{ET}})`$ $`=`$ $`\left({\displaystyle \frac{M}{E}}\right)^2E[{\displaystyle \frac{\theta (xy)}{(k_1^+p_1^+)(Ep_1^{}k_2^{})\mu ^2(𝒑_{}𝒌_{})^2}}`$ (54)
$`+{\displaystyle \frac{\theta (yx)}{(p_1^+k_1^+)(Ek_1^{}p_2^{})\mu ^2(𝒑_{}𝒌_{})^2}}].`$
Using the relations in Eq. (II C), we find
$`(k_1^+p_1^+)(Ep_1^{}k_2^{})`$ $`=`$ $`\left({\displaystyle \frac{k^3}{k^0}}{\displaystyle \frac{p^3}{p^0}}\right)\left({\displaystyle \frac{E^24M^2}{2}}p_{\text{ET}}^2k_{\text{ET}}^2\right)`$ (56)
$`+{\displaystyle \frac{k^3}{k^0}}{\displaystyle \frac{p^3}{p^0}}\left(k^0p^0\right)^2(k^3p^3)^2.`$
Under the exchange $`kp`$, the only thing that changes in Eq. (56) is that the first term picks up a minus sign. This observation allows Eq. (54) to be rewritten as
$`V_{\text{OBE}}(E;𝒌_{\text{ET}};𝒑_{\text{ET}})`$ $`=`$ $`\left({\displaystyle \frac{M}{E}}\right)^2{\displaystyle \frac{E}{\left|\frac{k^3}{k^0}\frac{p^3}{p^0}\right|\mathrm{\Delta }+\frac{k^3}{k^0}\frac{p^3}{p^0}(q_{\text{ET}}^0)^2𝒒_{\text{ET}}^2\mu ^2}},`$ (57)
where
$`\mathrm{\Delta }`$ $`=`$ $`{\displaystyle \frac{E^24M^2}{2}}p_{\text{ET}}^2k_{\text{ET}}^2`$ (58)
$`q_{\text{ET}}^\mu `$ $`=`$ $`k_{\text{ET}}^\mu p_{\text{ET}}^\mu .`$ (59)
Note that $`q_{\text{ET}}^{}`$ is not the light-front energy of the meson, since in a Hamiltonian theory only the light-front three-momenta are conserved; the four-momenta is not conserved. Equation (57) will be useful in the context of approximations based on the physical arguments that we will discuss in section IV.
### B TBE Potential
As in the previous section, we start by drawing all the allowed, non-vanishing time-ordered diagrams with two meson exchanges shown in Fig. 2. The diagrams are classified according to the behavior of the intermediate particles. The total TBE potential is given by the sum of all the diagrams, so
$`V_{\text{TBE}}`$ $`=`$ $`V_{\text{TBE:SB}}+V_{\text{TBE:SX}}+V_{\text{TBE:TX}}+V_{\text{TBE:WX}}+V_{\text{TBE:ZX}}.`$ (60)
In the diagrams for the TBE potential in Fig. 2, the intermediate loop momenta can be parameterized by $`𝒒_1`$ or $`𝒒_2`$. The dependent variable is defined by the relation $`𝑷=𝒒_1+𝒒_2`$. The Bjorken $`x`$ variable that corresponds to $`𝒒_1`$ ($`𝒒_2`$) is labeled with $`z`$ ($`1z`$). We use the Feynman rules to calculate all of these potentials, starting with
$`V_{\text{TBE:SB}}(E;𝒌_1;𝒑_1)`$ $`=`$ $`\left({\displaystyle \frac{M}{E}}\right)^4{\displaystyle }{\displaystyle \frac{d^2q_{}}{2(2\pi )^3}}[{\displaystyle _0^1}dz{\displaystyle \frac{\theta (zy)\theta (xz)}{z(1z)(zy)(xz)}}`$ (65)
$`\times {\displaystyle \frac{1}{Eq_1^{}k_2^{}\omega ^{}(𝒌_1𝒒_1)}}`$
$`\times {\displaystyle \frac{1}{Ep_1^{}k_2^{}\omega ^{}(𝒌_1𝒒_1)\omega ^{}(𝒒_1𝒑_1)}}`$
$`\times {\displaystyle \frac{1}{Ep_1^{}q_2^{}\omega ^{}(𝒒_1𝒑_1)}}]`$
$`+\{12\}.`$
The symbol $`\{12\}`$ means that all labels 1 are replaced with 2 and vice versa, as well as replacing the Bjorken variables $`x`$, $`y`$, and $`z`$ with $`1x`$, $`1y`$, and $`1z`$. This is a way of explicitly stating the symmetry of the potential under exchange of particles 1 and 2. A detailed discussion of the evaluation of the loop integral in Eq. (65) is given in Appendix D.
It is straightforward to calculate the other parts of the TBE potential,
$`V_{\text{TBE:SX}}(E;𝒌_1;𝒑_1)`$ $`=`$ $`\left({\displaystyle \frac{M}{E}}\right)^4{\displaystyle }{\displaystyle \frac{d^2q_{}}{2(2\pi )^3}}[{\displaystyle _0^1}dz{\displaystyle \frac{\theta (xz)\theta (zy)}{z(xz)(1+zxy)(zy)}}`$ (70)
$`\times {\displaystyle \frac{1}{Eq_1^{}k_2^{}\omega ^{}(𝒌_1𝒒_1)}}`$
$`\times {\displaystyle \frac{1}{Ep_1^{}k_2^{}\omega ^{}(𝒌_1𝒒_1)\omega ^{}(𝒒_1𝒑_1)}}`$
$`\times {\displaystyle \frac{1}{Ep_1^{}\widehat{q}_2^{}\omega ^{}(𝒌_1𝒒_1)}}]`$
$`+\{12\},`$
where we have denoted the light-front energy of particle 2 by $`\widehat{q}_2^{}`$, given by
$`\widehat{q}_2^{}`$ $`=`$ $`ϵ^{}(𝑷+𝒒_1𝒑_1𝒌_1)`$ (71)
$`ϵ^{}(𝒒)`$ $`=`$ $`{\displaystyle \frac{M^2+𝒒_{}^2}{q^+}}.`$ (72)
The rest of the TBE potential is given by
$`V_{\text{TBE:TX}}(E;𝒌_1;𝒑_1)`$ $`=`$ $`\left({\displaystyle \frac{M}{E}}\right)^4{\displaystyle }{\displaystyle \frac{d^2q_{}}{2(2\pi )^3}}[{\displaystyle _0^1}dz{\displaystyle \frac{\theta (xz)\theta (1+zxy)\theta (yz)}{z(xz)(1+zxy)(yz)}}`$ (77)
$`\times {\displaystyle \frac{1}{Eq_1^{}k_2^{}\omega ^{}(𝒌_1𝒒_1)}}`$
$`\times {\displaystyle \frac{1}{Eq_1^{}\widehat{q}_2^{}\omega ^{}(𝒌_1𝒒_1)\omega ^{}(𝒑_1𝒒_1)}}`$
$`\times {\displaystyle \frac{1}{Ep_1^{}\widehat{q}_2^{}\omega ^{}(𝒌_1𝒒_1)}}]`$
$`+\{12\}`$
$`V_{\text{TBE:WX}}(E;𝒌_1;𝒑_1)`$ $`=`$ $`\left({\displaystyle \frac{M}{E}}\right)^4{\displaystyle }{\displaystyle \frac{d^2q_{}}{2(2\pi )^3}}[{\displaystyle _0^1}dz{\displaystyle \frac{\theta (xz)\theta (1+zxy)\theta (yz)}{z(xz)(1+zxy)(yz)}}`$ (82)
$`\times {\displaystyle \frac{1}{Eq_1^{}k_2^{}\omega ^{}(𝒌_1𝒒_1)}}`$
$`\times {\displaystyle \frac{1}{Eq_1^{}\widehat{q}_2^{}\omega ^{}(𝒌_1𝒒_1)\omega ^{}(𝒑_1𝒒_1)}}`$
$`\times {\displaystyle \frac{1}{Eq_1^{}p_2^{}\omega ^{}(𝒑_1𝒒_1)}}]`$
$`+\{12\}`$
$`V_{\text{TBE:ZX}}(E;𝒌_1;𝒑_1)`$ $`=`$ $`\left({\displaystyle \frac{M}{E}}\right)^4{\displaystyle }{\displaystyle \frac{d^2q_{}}{2(2\pi )^3}}[{\displaystyle _0^1}dz{\displaystyle \frac{\theta (x+y1z)}{z(xz)(x+y1z)(yz)}}`$ (87)
$`\times {\displaystyle \frac{1}{Eq_1^{}k_2^{}\omega ^{}(𝒌_1𝒒_1)}}`$
$`\times {\displaystyle \frac{1}{Eq_1^{}k_2^{}p_2^{}ϵ^{}(𝒑_1+𝒌_1𝑷𝒒_1)}}`$
$`\times {\displaystyle \frac{1}{Eq_1^{}p_2^{}\omega ^{}(𝒑_1𝒒_1)}}]`$
$`+\{12\}.`$
The loop integrals in the expressions for the TBE potentials and the azimuthal-angle averaging are discussed in Appendix D.
### C TBE:SB Potential: Connection to the ladder Bethe-Salpeter equation
It is well known that the full, untruncated Bethe-Salpeter equation can be reduced to the full, untruncated Hamiltonian (Schrödinger-type) equation by integration over the energy or light-front energy variables. If a truncated kernel is used for the Bethe-Salpeter equation, then the physically equivalent Hamiltonian equation will not include all the graphs that the full theory allows. By physically equivalent, we mean that the spectra of the potential $`V`$ should reproduce the spectrum for the states of the Bethe-Salpeter equation, excluding the so-called “abnormal” states . For an extensive discussion of this in the equal-time case see, for instance, Klein , Phillips and Wallace , Lahiff and Afnan , and for examples on the light front, Chang and Ma and Ligterink and Bakker .
In particular, consider the Bethe-Salpeter equation when the ladder kernel is used. The physically equivalent light-front potential will not include any graphs where the meson lines cross, so to order $`g^4`$, the potential is given by $`g^2V_{\text{OBE}}+g^4V_{\text{TBE:SB}}`$. Therefore, by considering the TBE:SB truncation, we can test how well the light-front Hamiltonian approach approximates the full ladder Bethe-Salpeter equation. This is idea discussed more throughly in .
## IV Non-perturbative Potentials
The potentials discussed in this section are derived from the OBE field theory potential, but additional approximations are made to simplify the expressions.
### A Symmetrized-mass approximation
Krautgärtner, Pauli and Wölz and Trittmann and Pauli studied positronium with a large coupling constant in light-front dynamics. The one-photon-exchange potential they obtain has a colinear singularity due to the sum of the instantaneous photon exchange graph and a gauge-dependent factor from the spin sum. They argue that the singularity is not physical, and therefore must be canceled by higher-order terms in the potential. The effect of those terms can be simulated by choosing the bound-state energy so that the coefficient of the singular term vanishes. They find that this condition is met when the light-front energy $`P^{}`$ in the one-photon-exchange potential is replaced with the operator $`\omega `$, expressed here in the two-particle basis,
$`P^{}\omega (𝒌_1,𝒌_2;𝒑_1,𝒑_2)`$ $``$ $`{\displaystyle \frac{1}{2}}\left(p_1^{}+p_2^{}+k_1^{}+k_2^{}\right).`$ (88)
This is called the symmetrized mass , the average of the total $`P^{}`$ in the initial and final states. It is important to note that this approximation affects not only the singular term, but also the energy denominator in the rest of the OBE potential. The modified denominators simulate the effects of the non-perturbative higher-order terms that are not included explicitly in the OBE potential. Potentials obtained with this approximation are similar to those given by the unitary transformation method , where the potentials depend explicitly on the initial- and final-state energies.
In our model, there are no singularities associated with the OBE graphs because we deal only with scalar fields. However, we may use their approximation to obtain a new light-front OBE potential that should incorporate some non-perturbative effects. Recalling that the only place where $`P^{}`$ occurred in Eq. (50) was in the denominator, the $`E`$ in the denominator of the OBE potential is replaced with $`\omega `$ to get
$`V_\omega (𝒌_1;𝒑_1)`$ $`=`$ $`\left({\displaystyle \frac{M}{E}}\right)^2[{\displaystyle \frac{\theta (xy)/|xy|}{\frac{1}{2}\left(p_1^{}+p_2^{}+k_1^{}k_2^{}\right)\omega ^{}(𝒌_1𝒑_1)}}`$ (90)
$`+{\displaystyle \frac{\theta (yx)/|yx|}{\frac{1}{2}\left(+p_1^{}p_2^{}k_1^{}+k_2^{}\right)\omega ^{}(𝒑_1𝒌_1)}}]`$
$`=`$ $`\left({\displaystyle \frac{M}{E}}\right)^2{\displaystyle \frac{E}{\frac{1}{2}(k_1^+p_1^+)\left(p_1^{}+p_2^{}+k_1^{}k_2^{}\right)\mu ^2(𝒌_{}𝒑_{})^2}}.`$ (91)
Writing the light-front variables in the denominator in terms of the equal-time variables, as prescribed in Eq. (II C), we find
$`{\displaystyle \frac{1}{2}}(k_1^+p_1^+)\left(p_1^{}+p_2^{}+k_1^{}k_2^{}\right)`$ $`=`$ $`(k^3p^3)^2+{\displaystyle \frac{k^3}{k^0}}{\displaystyle \frac{p^3}{p^0}}(k^0p^0)^2.`$ (92)
Thus, Eq. (91) can be rewritten as
$`V_\omega (𝒌_{\text{ET}};𝒑_{\text{ET}})`$ $`=`$ $`\left({\displaystyle \frac{M}{E}}\right)^2{\displaystyle \frac{E}{\frac{k^3}{k^0}\frac{p^3}{p^0}(k^0p^0)^2(𝒌_{\text{ET}}𝒑_{\text{ET}})^2\mu ^2}}.`$ (93)
This result can also by obtained more directly by considering the first term in Eq. (56). Recall that the $`E^2`$ that appears in the denominator is written as $`P^+P^{}`$ in an arbitrary frame, so in the symmetrized-mass approximation, the $`E^2`$ term is replaced with $`E\omega `$. This causes $`\mathrm{\Delta }`$ term in the denominator of Eq. (57) to vanish, so the equation reduces to Eq. (93). Also, note that by writing this new potential, we attempt to incorporate physics from higher-order graphs than just the OBE graphs.
The singularity structure of the symmetrized-mass potential is easily analyzed. When scattering states are used, in the center-of-momentum frame the total energy of the state is $`E=2k^0=2p^0`$, so the relations in Eq. (II C) become
$$k_1^\pm =k^0\pm k^3,$$
(95)
$$k_2^\pm =k^0k^3.$$
(96)
Using these relations, the symmetrized mass is $`\omega =E`$. Thus, for scattering states, this potential is same as the OBE scattering potential and the singularity structure is the same.
### B Instantaneous and Retarded approximations
For our bound states, $`k_{\text{ET}}^2M^2`$, so that the $`\frac{k^3p^3}{k^0p^0}(k^0p^0)^2`$ will be much smaller than the other terms in the denominator of Eq. (93). Therefore, we may approximate the symmetrized-mass potential $`V_\omega `$ by the instantaneous potential,
$`V_{\text{Inst}}(𝒌_{\text{ET}};𝒑_{\text{ET}})`$ $`=`$ $`\left({\displaystyle \frac{M}{E}}\right)^2{\displaystyle \frac{E}{(𝒌_{\text{ET}}𝒑_{\text{ET}})^2+\mu ^2}}.`$ (97)
Alternatively, we may also argue that since the energy difference term is small, we can also approximate $`V_\omega `$ by the retarded potential,
$`V_{\text{Ret}}(𝒌_{\text{ET}};𝒑_{\text{ET}})`$ $`=`$ $`\left({\displaystyle \frac{M}{E}}\right)^2{\displaystyle \frac{E}{(k^0p^0)^2(𝒌_{\text{ET}}𝒑_{\text{ET}})^2\mu ^2}}`$ (98)
$`=`$ $`\left({\displaystyle \frac{M}{E}}\right)^2{\displaystyle \frac{E}{(k_{\text{ET}}p_{\text{ET}})^2\mu ^2}},`$ (99)
where $`k_{\text{ET}}`$ and $`p_{\text{ET}}`$ represent four-vectors, defined by the equal-time three-vectors and the condition that $`k_{\text{ET}}^2=p_{\text{ET}}^2=M^2`$. These potentials resemble the three-dimensional Blankenbecler-Sugar or Gross quasi-potentials.
Both of these approximations are reasonable if the energy difference between the initial and final states is small, which is valid for lightly-bound states. The instantaneous potential is a better approximation of the symmetrized-mass potential, since if we expand the symmetrized-mass potential to second-order in perturbation theory about $`k^0=p^0`$, we get $`V_\omega =V_{\text{Inst}}`$. Also, note that these potentials are explicitly rotationally invariant in terms of our equal-time parameterization, which provides significant computational advantages.
### C Three-dimensional reduction of the Bethe-Salpeter equation
We now consider a non-perturbative approximation used by Wallace and Mandelzweig . The basic idea is to first make an approximation of the Bethe-Salpeter equation, then reduce that modified Bethe-Salpeter equation to the physically equivalent Hamiltonian equation. This approach was used by Phillips and Wallace for the model we use, however, they obtained an equal-time Hamiltonian, while we seek a light-front Hamiltonian. Before we do this, we first review the basic mechanics of the three-dimensional reduction, as presented in Sales et al. and specialized to our particular case. We postpone the discussion of the approximation until section IV D.
The Bethe-Salpeter equation can be written in matrix form as $`\mathrm{\Gamma }=KG_0\mathrm{\Gamma }`$, or explicitly in function form in the momentum basis as
$`\mathrm{\Gamma }(k_1;P)`$ $`=`$ $`{\displaystyle \frac{d^4p_1}{(2\pi )^4}K(k_1,p_1;P)G_0(p_1;P)\mathrm{\Gamma }(p_1;P)}.`$ (100)
In these equations, $`\mathrm{\Gamma }`$ is the four-dimensional vertex function, $`K`$ is the four-dimensional kernel, and $`G_0`$ is the two-particle four-dimensional Green’s function. The momenta are $`P`$, the total four-momentum, and $`p_1`$, the four-momentum of particle 1. Particle 2’s momentum is implicitly $`Pp_1`$. The four-dimensional Green’s function is given by
$`G_0(k_1;P)`$ $`=`$ $`id(k_1)d(Pk_1),`$ (101)
where $`d`$ is the one-particle Green’s function. On the light front, $`d`$ can be written as
$`d(p)`$ $`=`$ $`\left({\displaystyle \frac{1}{p^+}}\right){\displaystyle \frac{1}{p^{}\text{Sign}(p^+)ϵ^{}(𝒑)}},`$ (102)
where the light-front energy $`ϵ^{}`$ is given by
$`ϵ^{}(𝒑)`$ $`=`$ $`{\displaystyle \frac{M^2+𝒑_{}^2}{|p^+|}}i\eta .`$ (103)
The real part of $`ϵ^{}`$ is a positive definite quantity, and $`\eta `$ is positive infinitesimal.
The Bethe-Salpeter equation can be rewritten as
$`\mathrm{\Gamma }`$ $`=`$ $`W\widehat{G}_0\mathrm{\Gamma },`$ (104)
where $`\widehat{G}_0`$ is an auxiliary Green’s function, and $`W`$ is defined by
$`W`$ $`=`$ $`K+K(G_0\widehat{G}_0)W.`$ (105)
The advantage of this rearrangement is that we are free to choose the form of the auxiliary Green’s function, $`\widehat{G}_0`$. The choice of $`\widehat{G}_0`$ advocated in Ref. is
$`\widehat{G}_0(k_1,p_1;P)`$ $`=`$ $`G_0(k_1;P){\displaystyle \frac{\delta ^{(2,+)}(𝒌_1𝒑_1)}{g_0(𝒌_1;P)}}G_0(p_1;P),`$ (106)
where
$`g_0(𝒌_1,P)`$ $`=`$ $`{\displaystyle \frac{dk_1^{}}{2(2\pi )}G_0(k_1;P)}.`$ (107)
There is an extra factor of 2 in the denominator of Eq. (107) when compared to the equal-time formalism. This is due to the Jacobian of the light-front coordinates.
Using the definition of $`\widehat{G}_0`$ given in Eq. (106), we can integrate the modified Bethe-Salpeter equation, in Eq. (104), over the light-front energy to get
$`\gamma (𝒌_1;P)`$ $`=`$ $`{\displaystyle \frac{d^2p_{1,}dp_1^+}{(2\pi )^3}w(𝒌_1,𝒑_1;P)g_0(𝒑_1;P)\gamma (𝒑_1;P)},`$ (108)
where
$`w(𝒌_1,𝒑_1;P)`$ $``$ $`{\displaystyle \frac{1}{g_0(𝒌_1;P)}}G_0WG_0(𝒌_1,𝒑_1;P){\displaystyle \frac{1}{g_0(𝒑_1;P)}}`$ (109)
$`\gamma (𝒌_1;P)`$ $``$ $`{\displaystyle \frac{1}{g_0(𝒌_1;P)}}{\displaystyle \frac{dk_1^{}}{2(2\pi )}G_0(k_1;P)\mathrm{\Gamma }(k_1;P)}.`$ (110)
The functional $`f`$ is defined by its action on an arbitrary function $`f(k_1,p_1)`$, where $`k_1`$ and $`p_1`$ are four-vectors, as
$`f(𝒌_1,𝒑_1)`$ $`=`$ $`{\displaystyle \frac{dk_1^{}}{2(2\pi )}\frac{dp_1^{}}{2(2\pi )}f(k_1,p_1)}.`$ (111)
We proceed by calculating the specific form of $`g_0`$,
$`g_0(𝒌_1;P)`$ $`=`$ $`{\displaystyle \frac{\theta (k_1^+)\theta (k_2^+)}{2k_1^+k_2^+}}{\displaystyle \frac{1}{P^{}k_1^{}k_2^{}}},`$ (112)
where $`k_i^{}=ϵ^{}(𝒌_i)`$. As $`g_0`$ is a three-dimensional quantity, it is clear that $`k_i^{}`$ is not the independent minus component of a momentum four-vector. With this expression for $`g_0`$, we can specialize Eq. (108) to the center-of-momentum frame and obtain
$`\left(Ek_1^{}k_2^{}\right)\psi (𝒌_1;E)`$ $`=`$ $`{\displaystyle d^2p_{1,}_0^E𝑑p_1^+\frac{w(𝒌_1,𝒑_1;E)}{2(2\pi )^3\sqrt{k_1^+k_2^+p_1^+p_2^+}}\psi (𝒑_1;E)},`$ (113)
where
$`\psi (𝒌_1;E)`$ $`=`$ $`{\displaystyle \frac{g_0(𝒌_1;E)}{\sqrt{k_1^+k_2^+}}}\gamma (𝒌_1;E).`$ (114)
By comparing this equation to Eq. (33), we find that the full light-front two-nucleon effective potential that corresponds to the kernel $`K`$, after suppressing the coefficient given in Eq. (48), is
$`V(𝒌_1,𝒑_1;E)`$ $`=`$ $`{\displaystyle \frac{1}{E}}w(𝒌_1,𝒑_1;E).`$ (115)
Thus, we can calculate light-front potentials directly from the Bethe-Salpeter equation using this method.
The potential $`V`$ can be expanded in powers of the coupling constant, as done in LFTOPT. We find that when the auxiliary Green’s function given in Eq. (106) and the OBE kernel are used, the lowest order parts of the potential (as calculated in ) are the same as our OBE and TBE:SB potentials. Thus, we conclude that this method produces the physically equivalent Hamiltonian theory to the Bethe-Salpeter equation being used. We will use this in the next section to derive a Hamiltonian potential for a situation where LFTOPT cannot be used.
### D The modified-Green’s-function approach
Now that the technology for the three-dimensional reduction has been reviewed, we derive an approximate kernel for the Bethe-Salpeter equation. We will follow the approach of Phillips and Wallace and works cited therein. The idea is to start with the Bethe-Salpeter equation where the kernel is truncated to only include ladder (one-boson-exchange) and crossed (two-boson-exchange) parts,
$`\mathrm{\Gamma }`$ $`=`$ $`(K_{\text{ladder}}+K_{\text{cross}})G_0\mathrm{\Gamma }.`$ (116)
An uncrossed approximation is used where the crossed part of the kernel is approximated by $`K_{\text{cross}}K_{\text{ladder}}G_CK_{\text{ladder}}`$. Our job is to find a valid modified Green’s function, $`G_C`$. Using this uncrossed approximation,
$`\mathrm{\Gamma }`$ $``$ $`(K_{\text{ladder}}+K_{\text{ladder}}G_CK_{\text{ladder}})G_0\mathrm{\Gamma }.`$ (117)
One can attempt to rewrite Eq. (116) as an equation linear in $`K_{\text{ladder}}`$, to obtain the modified-Green’s-function Bethe-Salpeter equation,
$`\mathrm{\Gamma }_{\text{MGF}}`$ $`=`$ $`K_{\text{ladder}}\left(G_0+G_C\right)\mathrm{\Gamma }_{\text{MGF}}.`$ (118)
By iterating this integral equation for $`\mathrm{\Gamma }_C`$, we obtain
$`\mathrm{\Gamma }_{\text{MGF}}`$ $`=`$ $`\left[K_{\text{ladder}}+{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}K_{\text{ladder}}\left(G_CK_{\text{ladder}}\right)^n\right]G_0\mathrm{\Gamma }_{\text{MGF}}.`$ (119)
The part of Eq. (118) that plays the role of the kernel includes the uncrossed approximation of the original kernel $`K_{\text{ladder}}+K_{\text{ladder}}G_CK_{\text{ladder}}`$ as well as many more terms. We note that the higher-order terms approximate some of the higher-order terms that should be included in the full kernel, such as three-boson-exchange diagrams where several meson lines cross. However, this approach undercounts the higher-order terms which it approximates, and also leaves out some terms completely. Therefore, this new Bethe-Salpeter equation will give results that are closer to the full solution than Eq. (116), but will not give the exact solution. The articles by Wallace and Mandelzweig demonstrate that this approach, by effectively summing an infinite set of interactions, gives the correct one-body limit, which is something that the usual Bethe-Salpeter equation with a truncated kernel cannot do.
The modified Bethe-Salpeter equation in Eq. (118) is reduced to a Hamiltonian equation via the technique discussed in the previous section. The equal-time Hamiltonian has been derived by Phillips and Wallace . They found that this modified-Green’s-function approach gave a spectra that lies closer to the full ground-state spectra than the other approximations they considered. We will use the light-front reduction to obtain the light-front potential for the Hamiltonian equation physically equivalent to Eq. (118).
To clearly see what role $`G_C`$ plays, we compare the crossed and uncrossed Feynman graphs in Fig. 3. Using the Feynman rules,
$`K_{\text{crossed}}`$ $``$ $`{\displaystyle \frac{d^4q_1}{(2\pi )^4}\frac{1}{(k_1q_1)^2\mu ^2}d(q_1)d(\widehat{q}_2)\frac{1}{(p_1q_1)^2\mu ^2}}`$ (120)
$`K_{\text{uncrossed}}`$ $``$ $`{\displaystyle \frac{d^4q_1}{(2\pi )^4}\frac{1}{(k_1q_1)^2\mu ^2}d(q_1)d(q_2)\frac{1}{(p_1q_1)^2\mu ^2}},`$ (121)
where $`d`$ is the one-particle propagator given in Eq. (102), and $`\widehat{q}_2=q_2+p_2+k_2`$. The only difference between these two graphs is that the crossed one has $`d(\widehat{q}_2)`$ while the uncrossed one has $`d(q_2)`$.
We want an approximate one-particle Green’s function $`d_C`$ that only depends on $`q_1`$ and $`P`$, so that
$`d_C(q_1;P)`$ $``$ $`d(\widehat{q}_2).`$ (122)
Substitution of $`d_C(q_1;P)`$ for $`d(\widehat{q}_2)`$ in the crossed graph causes the graph to become uncrossed. The penalty for this simplification is that a modified Green’s function propagates in the intermediate state, namely $`id(q_1)d_C(q_1;P)`$. It is important that this approximation is invariant under relabeling particle labels, so we explicitly symmetrize by defining
$`G_C`$ $`=`$ $`{\displaystyle \frac{i}{2}}\left[d(q_1)d_C(q_1;P)+d(q_2)d_C(q_2;P)\right]`$ (123)
How should we approximate $`d_C`$? Since we are interested in obtaining a bound state, a low-energy approximation is chosen. Specializing to the center-of-momentum frame in this limit, the external momenta are half the total momentum, so $`p_2=k_2=P/2`$ and $`\widehat{q}_2=Pq_2=q_1`$. This approximation is similar to the one used by Phillips and Wallace. Thus, we define $`d_C(q_1;P)d(q_1)`$ so
$`G_C(k_1;P)`$ $`=`$ $`G_1(k_1;P)+G_2(k_1;P),`$ (124)
where we define $`G_1`$ and $`G_2`$ by
$`G_1(k_1;P)`$ $`=`$ $`{\displaystyle \frac{i}{2}}d(k_1)^2`$ (125)
$`G_2(k_1;P)`$ $`=`$ $`{\displaystyle \frac{i}{2}}d(Pk_1)^2.`$ (126)
This approximation for $`G_C`$ is valid for this model for the energy range we study, as discussed in Appendix E.
We can write the modified-Green’s-function Bethe-Salpeter equation as
$`\mathrm{\Gamma }_{\text{MGF}}=K_{\text{ladder}}\stackrel{~}{G}_0\mathrm{\Gamma }_{\text{MGF}},`$ (127)
where
$`\stackrel{~}{G}_0`$ $`=`$ $`G_0+G_1+G_2.`$ (128)
This is used as the starting point of three-dimensional reduction discussed in Section IV C, where $`\stackrel{~}{G}_0`$ is considered as the Green’s function. Before doing the reduction, note that the two poles of $`G_1`$ and $`G_2`$ lie in the same half plane for each function, so
$`\stackrel{~}{g}_0(𝒌_1,P)`$ $``$ $`{\displaystyle }{\displaystyle \frac{dk_1^{}}{2(2\pi )}}\stackrel{~}{G}_0(k_1,;P)`$ (129)
$`=`$ $`g_0(𝒌_1,P)`$ (130)
Proceeding with the three-dimensional reduction of Eq. (127) in the center-of-momentum frame, we obtain
$`\left(Ek_1^{}k_2^{}\right)\stackrel{~}{\psi }(𝒌_1;E)`$ $`=`$ $`{\displaystyle d^2p_{1,}_0^E𝑑p_1^+\frac{\stackrel{~}{w}(𝒌_1,𝒑_1;E)}{2(2\pi )^3\sqrt{k_1^+k_2^+p_1^+p_2^+}}\stackrel{~}{\psi }(𝒑_1;E)},`$ (131)
where
$`\stackrel{~}{w}(𝒌_1,𝒑_1;P)`$ $`=`$ $`{\displaystyle \frac{1}{g_0(𝒌_1;P)}}\stackrel{~}{G}_0\stackrel{~}{W}\stackrel{~}{G}_0(𝒌_1,𝒑_1;P){\displaystyle \frac{1}{g_0(𝒑_1;P)}}`$ (132)
$`\stackrel{~}{\psi }(𝒌_1;E)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{k_1^+k_2^+}}}{\displaystyle \frac{dk_1^{}}{2(2\pi )}\stackrel{~}{G}_0(k_1;P)\mathrm{\Gamma }_{\text{MGF}}(k_1;P)},`$ (133)
and the modified kernel $`\stackrel{~}{W}`$ is given by
$`\stackrel{~}{W}`$ $`=`$ $`K_{\text{ladder}}+K_{\text{ladder}}\left[\stackrel{~}{G}_0\widehat{\stackrel{~}{G}}_0\right]\stackrel{~}{W}`$ (134)
$`\widehat{\stackrel{~}{G}}_0(k_1,p_1;P)`$ $`=`$ $`\stackrel{~}{G}_0(k_1;P){\displaystyle \frac{\delta ^{(2,+)}(𝒌_1𝒑_1)}{g_0(𝒌_1;P)}}\stackrel{~}{G}_0(p_1;P).`$ (135)
It is a feature of the light front that $`\stackrel{~}{g}_0=g_0`$, so that the uncrossed approximation only affects the potential, and Eq. (131) has the same form as Eq. (108). In the equal-time calculation $`\stackrel{~}{g}_0g_0`$, so the approximation changes both the Green’s function as well as the potential.
We now expand $`\stackrel{~}{W}`$ in powers of the coupling constant, and keep only the lowest order term, $`K_{\text{ladder}}`$. According to Eq. (115), the light-front potential that corresponds to this truncation of the kernel is the modified-Green’s-function (MGF) potential $`V_{\text{MGF}}`$,
$`V_{\text{MGF}}(𝒌_1,𝒑_1;P)`$ $`=`$ $`{\displaystyle \frac{1}{E}}{\displaystyle \frac{1}{g_0(𝒌_1;P)}}\stackrel{~}{G}_0K_{\text{ladder}}\stackrel{~}{G}_0(𝒌_1,𝒑_1;P){\displaystyle \frac{1}{g_0(𝒑_1;P)}}.`$ (136)
The one-boson-exchange kernel $`K_{\text{ladder}}`$ is given by the Feynman diagram, so
$`K_{\text{ladder}}(k_1,p_1;P)`$ $`=`$ $`{\displaystyle \frac{(iM)^2}{(k_1p_1)^2\mu ^2+i\eta }}`$ (137)
$`=`$ $`\left({\displaystyle \frac{1}{k_1^+p_1^+}}\right){\displaystyle \frac{M^2}{(k_1^{}p_1^{})\text{Sign}(k_1^+p_1^+)\omega ^{}(𝒌_1𝒑_\mathrm{𝟏})}},`$ (138)
where the light-front energy of the meson is given by
$`\omega ^{}(𝒒)`$ $`=`$ $`{\displaystyle \frac{\mu ^2𝐪_{}^2}{|q^+|}}i\eta .`$ (139)
By examining the locations of all the poles in the $`k^{}`$ integrals for $`V_{\text{MGF}}`$, we find the integrals are non-vanishing only when both $`x`$ and $`y`$ are between $`0`$ and $`1`$. The sign functions in the denominator of $`K_{\text{ladder}}`$ naturally divide $`V_{\text{MGF}}`$ into two parts, one for $`x<y`$ and the other for $`x>y`$. The integrals in $`V_{\text{MGF}}`$ are straightforward, but quite lengthy and tedious. Therefore, we show only the final answer,
$`V_{\text{MGF}}(𝒌_1,𝒑_1;P)`$ $`=`$ $`\left({\displaystyle \frac{M}{E}}\right)^2[{\displaystyle \frac{\theta (xy)}{|xy|}}({\displaystyle \frac{1}{D_1}}+{\displaystyle \frac{N_{p,21}+N_{k,12}}{2D_1^2}}+{\displaystyle \frac{N_{p,21}N_{k,12}}{2D_1^3}})`$ (141)
$`+{\displaystyle \frac{\theta (yx)}{|yx|}}({\displaystyle \frac{1}{D_2}}+{\displaystyle \frac{N_{k,21}+N_{p,12}}{2D_2^2}}+{\displaystyle \frac{N_{p,12}N_{k,21}}{2D_2^3}})],`$
where
$`N_{k,12}`$ $`=`$ $`{\displaystyle \frac{k_1^+}{k_2^+}}(Ek_1^{}k_2^{})`$ (142)
$`N_{k,21}`$ $`=`$ $`{\displaystyle \frac{k_2^+}{k_1^+}}(Ek_1^{}k_2^{})`$ (143)
$`D_1`$ $`=`$ $`Ep_1^{}k_2^{}\omega ^{}(𝒌_1𝒑_2)`$ (144)
$`D_2`$ $`=`$ $`Ek_1^{}p_2^{}\omega ^{}(𝒑_1𝒌_2).`$ (145)
The expressions for $`N_{p,12}`$ and $`N_{p,21}`$ are obtained by replacing $`k`$ with $`p`$ in $`N_{k,12}`$ and $`N_{k,21}`$.
What is the physical interpretation of this modified-Green’s-function potential? The first term multiplying each $`\theta `$ function gives the OBE potential we derived before from the perturbation theory. There the $`D`$ in the denominators corresponds to one meson exchange. The second and third terms multiplying the $`\theta `$ functions, with $`D^2`$ and $`D^3`$ in the denominators appear to be effective two- and three-meson-exchange terms. Since time-ordered perturbation theory does not apply to the modified Bethe-Salpeter equation that we use, the exact nature of these terms is not easy to understand. However, it is clear that these terms increase the strength of the potential, and should mimic the higher-order diagrams that are not being included explicitly.
The only dependence on the direction of the perpendicular components of $`k`$ and $`p`$ comes from the $`D`$’s. This allows the azimuthal-angle integration of $`V_{\text{MGF}}`$ to be done easily, as shown in Appendix C.
## V Results
For our numerical work, we pick the meson mass to be $`0.15`$ times that of the nucleon, so $`\mu =0.15M`$. This is chosen so that our ground state can be considered a toy model of deuterium, and also to facilitate comparison with the results of Nieuwenhuis and Tjon and Phillips and Afnan . Nieuwenhuis and Tjon used the Feynman-Schwinger representation (FSR) of the two-particle Green’s function in the quenched approximation without the mass and vertex renormalization terms . Their result is to be considered the full solution that the Bethe-Salpeter and Hamiltonian equations approximate. For the Bethe-Salpeter equation, computation of the bound-state energies for models similar to ours have been done for the ladder and ladder plus crossed kernels over 30 years ago. More recent results are found in , where the solutions are compared to those given by the FSR approach.
Now consider how the light-front Hamiltonian approach fits in with the other approaches. As discussed in section III C, different light-front potentials can be derived from Bethe-Salpeter equations with different kernels. We have mentioned that the OBE+TBE:SB potential should approximate the ladder Bethe-Salpeter equation, and similarly the OBE+TBE potential should approximate the Bethe-Salpeter equation when the ladder plus crossed kernel is used. The best that these truncated Hamiltonians can do is approximate their respective Bethe-Salpeter equations.
With this in mind, we evaluate the coupling constant versus bound-state energy curves (which we will call the spectrum) for the Hamiltonian equation with the OBE potential, the OBE+TBE:SB potential, and the OBE+TBE potential. For the range of values we use here, we find numerical errors in the value of $`g^2`$ are less than 2%. Our results (without error bars) are plotted along with the results obtained with the ladder BSE , and ladder plus crossed BSE in Fig. 4. We note that the OBE+TBE:SB potential agrees well with the ladder BSE, and the OBE+TBE potential agrees with the ladder plus crossed BSE. This is the best that a Hamiltonian can do, so this result is interpreted as evidence that, in general, the higher-order diagrams are very small for the ground state on the light front.
If all one wanted was a way to approximate the spectra for Bethe-Salpeter equations with different kernels, one could just use the truncated potentials that the LFTOPT provide. However, the true goal is to approximate the spectra for the full ground state, which in this model is given by the FSR approach . We expect that the non-perturbative potentials should give a better approximation of the full solution than the perturbative potentials, since the non-perturbative potentials attempt to incorporate physics from higher-order diagrams, although this is not immediately clear by looking at the forms of the potentials used. We plot the results for all of the light-front potentials described in this paper, the three truncated potentials (OBE, OBE+TBE:SB, and OBE+TBE) and the four non-perturbative potentials (symmetrized mass, retarded, instantaneous, and modified Green’s function) along with the results for the full theory in Fig. 5. For deeply-bound states, there is considerable disagreement between the perturbative results and the full results, while the non-perturbative results do better, with the modified-Green’s-function (MGF) potential achieving the closest agreement. For lightly-bound states, the results for all of the potentials appear converge to each other, close to the full result.
For the modified-Green’s-function potential, only the first term of the expansion in $`g^2`$ was kept. In principle, higher-order terms could be calculated. However, since there was fairly good agreement between the OBE potential and the ladder Bethe-Salpeter equation, it is expected that the MGF potential will give results that are close to the ladder Bethe-Salpeter equation using the modified-Green’s-function.
## VI Conclusions
In this paper, we consider the ground state of a massive Wick-Cutkosky model using a Hamiltonian derived from light-front dynamics. We examine three different truncations of the effective potential derived from the perturbative field theory, and four approximations that attempt to incorporate non-perturbative physics. For each of these potentials, we calculate the coupling constant that gives the ground-state for a given bound-state energy (the spectra), and compare to the spectra from different approaches found in the literature. We find fairly good agreement between all the methods for lightly-bound systems.
For the full range of binding energies studied, the results for calculations including one- and two-boson-exchange potentials agree with the Bethe-Salpeter equation using the physically equivalent truncation of the kernel within the numerical errors. (This is a consequence of examining the ground state. For the excited states more higher-order light-front time-ordered graphs are required to get the same level of agreement .) The agreement for the case with the stretched box diagrams (OBE+TBE:SB) has been shown previously by Sales et al. ; the result with the crossed-box contribution (OBE+TBE) is new. This excellent agreement with the Bethe-Salpeter results has an undesirable consequence: The BSE results are known to be a poor approximation of the full solution for deeply bound systems, so the truncated Hamiltonian approach cannot provide a good approximation to the full solution in that regime.
The non-perturbative potentials based on physical considerations give a better approximation of the full solution than the potentials obtained from LFTOPT. For all binding energies, the modified-Green’s-function potential achieves the closest agreement with the full solution of all the potentials considered here. However, there is still considerable disagreement between approximate potentials and the full result for deeply bound states. We interpret this as an indication that the approximations, while incorporating some non-perturbative physics, do not go far enough. In the weakly-bound regime, which is of relevance for deuteron calculations, the spectra for all of the potentials are close together, indicating that light-front dynamics provides a good description of lightly-bound systems.
###### Acknowledgements.
We are grateful to D.R. Phillips for extensive discussions and for providing us with unpublished material. This work is supported in part by the U.S. Dept. of Energy under Grant No. DE-FG03-97ER4014.
## A Notation, conventions, and useful relations
This is patterned after the review by Harindranath. For a general four-vector $`a`$, we define the light-front variables
$`a^\pm `$ $`=`$ $`a^0\pm a^3,`$ (A1)
$`𝒂_{}`$ $`=`$ $`(a^1,a^2),`$ (A2)
so the 4-vector $`a^\mu `$ can be denoted
$`a`$ $`=`$ $`(a^+,a^{},𝒂_{}).`$ (A3)
Using this, we find that the scalar product is
$`ab`$ $`=`$ $`a^\mu b_\mu ={\displaystyle \frac{1}{2}}\left(a^+b^{}+a^{}b^+\right)𝒂_{}𝒃_{}.`$ (A4)
This defines $`g_{\mu \nu }`$, with $`g_+=g_+=1/2`$, $`g_{11}=g_{22}=1`$, and all other elements of $`g`$ vanish. The elements of $`g^{\mu \nu }`$ are obtained from the condition that $`g^{\mu \nu }`$ is the inverse of $`g_{\mu \nu }`$, so $`g^{\alpha \beta }g_{\beta \lambda }=\delta _\lambda ^\alpha `$. Its elements are the same as those of $`g_{\mu \nu }`$, except for $`g^+=g^+=2`$. Thus,
$`a^\pm `$ $`=`$ $`2a_{}.`$ (A5)
and the partial derivatives are similarly given by
$`^\pm `$ $`=`$ $`2_{}=2{\displaystyle \frac{}{x^{}}}.`$ (A6)
Moving to the physical consequences of this coordinate system, the commutation relations $`[p^\mu ,x^\nu ]=ig^{\mu \nu }`$ yields
$`[p^\pm ,x^{}]`$ $`=`$ $`2i`$ (A7)
$`[𝒑_{}^i,𝒙_{}^j]`$ $`=`$ $`i\delta _{i,j},`$ (A8)
with the other commutators equal to zero. Thus, $`𝒙_{}^i`$ is canonically conjugate to $`𝒑_{}^i`$, and $`x^\pm `$ is conjugate to $`p^{}`$. In light-front dynamics, $`x^+`$ plays the role of time (the light-front time), so $`p^{}`$ is the light-front energy and the light-front Hamiltonian is given by $`P^{}`$.
Particles have the light-front energy defined by the on-shell constraint $`k^2=m^2`$. This implies that the light-front energy is
$`k^{}`$ $`=`$ $`{\displaystyle \frac{m^2+𝒌_{}^2}{k^+}}.`$ (A9)
The free components of the momentum can be written as the light-front three-vector $`𝒌`$, denoted by
$`𝒌`$ $`=`$ $`(k^+,𝒌_{}).`$ (A10)
## B Conversion to matrix form
To solve for the bound-state wavefunction numerically, the light-front Schrödinger equation given in Eq. (44) must be discretized and cast in matrix form. The equation is first symmetrized to get
$`{\displaystyle _0^{\mathrm{}}}𝑑p_{\text{ET}}{\displaystyle _0^{\pi /2}}𝑑\theta _pV_S^+(k_{\text{ET}},\theta _k;p_{\text{ET}},\theta _p)\psi _S^{\text{GS}}(p_{\text{ET}},\theta _p)`$ $`=`$ $`\psi _S^{\text{GS}}(k_{\text{ET}},\theta _k),`$ (B1)
where
$`V_S^+(k_{\text{ET}},\theta _k;p_{\text{ET}},\theta _p)`$ $`=`$ $`B(k_{\text{ET}})A(k_{\text{ET}},\theta _k)V^+(k_{\text{ET}},\theta _k;p_{\text{ET}},\theta _p)A(p_{\text{ET}},\theta _p)B(p_{\text{ET}})`$ (B2)
$`\psi _S^{\text{GS}}(k_{\text{ET}},\theta _k)`$ $`=`$ $`B(k_{\text{ET}})^1A(k_{\text{ET}},\theta _k)\psi ^{\text{GS}}(k_{\text{ET}},\theta _k)`$ (B3)
$`A(k_{\text{ET}},\theta _k)`$ $`=`$ $`\sqrt{{\displaystyle \frac{2k_1^+k_2^+k_{\text{ET}}^2\mathrm{sin}\theta _k}{(k^0)}}}`$ (B4)
$`B(k_{\text{ET}})`$ $`=`$ $`\sqrt{{\displaystyle \frac{1}{E^24(k^0)^2}}}.`$ (B5)
Before discretizing the integrals, note that
$`{\displaystyle _0^{\mathrm{}}}𝑑pf(p)`$ $`=`$ $`a{\displaystyle _0^1}𝑑u\left(f(au)+{\displaystyle \frac{f(a/u)}{u^2}}\right).`$ (B6)
Using this trick, the $`p_{\text{ET}}`$ integral in Eq. B1 can be written as an integral over a finite range. Since we are concerned with a bound state, the wavefunction is exponentially damped for large momenta, and the second term of Eq. (B6) converges as $`u`$ approaches zero.
All the integrals in Eq. B1 then are over a finite range, and can be discretized using Gauss-Legendre quadrature. The specific routines for the quadrature are given by Numerical Recipes in C . This conversion gives a matrix equation that approximates the original Eq. B1,
$`V_S^+(g(E),E)\psi _S^{\text{GS}}`$ $`=`$ $`\psi _S^{\text{GS}},`$ (B7)
where the explicit dependence of $`V_S^+`$ on the symmetrized potential on the binding energy $`E`$ and the coupling constant $`g`$ is shown. This equation must be solved self-consistently for the spectrum $`g(E)`$.
The approach we use is to first solve for the spectrum for the OBE potential. The eigenvalue equation
$`V_{S,\text{OBE}}^+(E)\psi _S^{\text{GS}}`$ $`=`$ $`\alpha \psi _S^{\text{GS}},`$ (B8)
where
$`\alpha `$ $`=`$ $`{\displaystyle \frac{1}{g_{\text{OBE}}(E)^2}}.`$ (B9)
The ground-state wavefunction is the eigenvector that corresponds to the smallest eigenvalue $`\alpha `$. We calculate the wavefunction and the smallest coupling constant using EISPACK routines for a range of energies to map out the spectrum.
Using the coupling constant for the OBE potential as a starting point, we can use Eq. B7 for higher-order potentials that include $`N`$ meson exchanges. For a given energy, the coupling constant $`g(E)`$ is initially chosen as $`g_{\text{OBE}}(E)`$, then we solve
$`\left[{\displaystyle \underset{n=1}{\overset{N}{}}}g(E)^{2n}V_{S,(2n)}^+(E)\right]\psi _S^{\text{GS}}`$ $`=`$ $`\beta \psi _S^{\text{GS}},`$ (B10)
as an eigenvalue equation for $`\beta `$. The coupling constant $`g(E)`$ is varied until the the lowest eigenvalue is $`\beta =1`$, at which point $`g(E)`$ is the correct value of the spectrum corresponding to the ground-state wavefunction $`\psi _S^{\text{GS}}`$.
## C Azimuthal-angle integration of the OBE and MGF potentials
In this section, we evaluate the azimuthal-angle integration of the OBE potential in Eq. (54) and the first term in MGF potential in Eq. (141), using the prescription for azimuthal-angle integration given in Eq. (47). One of the integrals is easily done since since the potential is independent of the azimuthal angle between the two perpendicular momenta, so
$`V(k^+,k_{};p^+,p_{})`$ $``$ $`\left[\theta (xy){\displaystyle _0^{2\pi }}{\displaystyle \frac{d\varphi }{A_1+B\mathrm{cos}\varphi }}+\theta (yx){\displaystyle _0^{2\pi }}{\displaystyle \frac{d\varphi }{A_2+B\mathrm{cos}\varphi }}\right],`$ (C1)
where
$`A_1`$ $`=`$ $`(k_1^+p_1^+)(Ep_1^{}k_2^{})\mu ^2p_{}^2k_{}^2`$ (C2)
$`A_2`$ $`=`$ $`(p_1^+k_1^+)(Ek_1^{}p_2^{})\mu ^2p_{}^2k_{}^2`$ (C3)
$`B`$ $`=`$ $`2k_{}p_{}.`$ (C4)
The integrals in Eq. C1 are easily done to give, since the $`A`$’s are negative,
$`{\displaystyle _0^{2\pi }}{\displaystyle \frac{d\varphi }{A+B\mathrm{cos}\varphi }}`$ $`=`$ $`{\displaystyle \frac{2\pi }{\sqrt{A^2B^2}}}.`$ (C5)
Using this, the azimuthal-angle-averaged OBE potential is given by
$`V_{\text{OBE}}(k^+,k_{};p^+,p_{})`$ $`=`$ $`2\pi \left({\displaystyle \frac{M}{E}}\right)^2E\left[{\displaystyle \frac{\theta (xy)}{\sqrt{A_1^2B^2}}}+{\displaystyle \frac{\theta (yx)}{\sqrt{A_2^2B^2}}}\right].`$ (C6)
It is straightforward to rewrite this equation for the potential in terms of the equal-time coordinates.
When the other terms in the MGF potential are azimuthal-angle averaged, integrations similar to the one given in Eq. C5 are encountered, with the denominator squared or cubed. We note that
$`{\displaystyle _0^{2\pi }}{\displaystyle \frac{d\varphi }{(A+B\mathrm{cos}\varphi )^2}}`$ $`=`$ $`{\displaystyle \frac{2\pi }{\sqrt{A^2B^2}}}{\displaystyle \frac{A}{A^2B^2}}`$ (C7)
$`{\displaystyle _0^{2\pi }}{\displaystyle \frac{d\varphi }{(A+B\mathrm{cos}\varphi )^3}}`$ $`=`$ $`{\displaystyle \frac{2\pi }{\sqrt{A^2B^2}}}{\displaystyle \frac{2A^2+B^2}{(A^2B^2)^2}},`$ (C8)
so the azimuthal-angle-averaged MGF potential is given by
$`V_{\text{MGF}}(𝒌_1,𝒑_1;P)`$ $`=`$ $`2\pi \left({\displaystyle \frac{M}{E}}\right)^2[{\displaystyle \frac{\theta (xy)}{\sqrt{A_1^2B^2}}}(1+{\displaystyle \frac{N_{p,21}+N_{k,12}}{2D_{1,2}}}+{\displaystyle \frac{N_{p,21}N_{k,12}}{2D_{1,3}}})`$ (C10)
$`+{\displaystyle \frac{\theta (yx)}{\sqrt{A_2^2B^2}}}(1+{\displaystyle \frac{N_{k,21}+N_{p,12}}{2D_{2,2}}}+{\displaystyle \frac{N_{p,12}N_{k,21}}{2D_{2,3}}})],`$
where
$`D_{i,2}`$ $``$ $`{\displaystyle \frac{A^2B^2}{A}}`$ (C11)
$`D_{i,3}`$ $``$ $`{\displaystyle \frac{(A^2B^2)^2}{2A^2+B^2}},`$ (C12)
and $`i=1,2`$.
## D Azimuthal-Angle Integration and Loop integration of the TBE potentials
As in the previous section, we want the azimuthal-angle integrals of the TBE potentials given in Eqs. (65-87). For these potentials, there is also a loop integral that has to be done. We start by analyzing the equations schematically. Each of the terms in the TBE potentials can be written in the following form,
$`V_{\text{TBE}}(k^+,k_{};p^+,p_{})`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{q_{}dq_{}}{2(2\pi )^3}}{\displaystyle _0^1}𝑑zJ(k^+,q^+,p^+)I(k^+,k_{},q^+,q_{},p^+,p_{})`$ (D1)
$`I(k^+,k_{},q^+,q_{},p^+,p_{})`$ $`=`$ $`{\displaystyle _0^{2\pi }}𝑑\varphi _q𝑑\varphi _p{\displaystyle \frac{1}{A_1+B_1\mathrm{cos}\varphi _q}}`$ (D4)
$`\times {\displaystyle \frac{1}{A_2+B_2\mathrm{cos}\varphi _q+C_2\mathrm{cos}\varphi _p+D_2\mathrm{cos}(\varphi _p\varphi _q)}}`$
$`\times {\displaystyle \frac{1}{A_3+B_3\mathrm{cos}\varphi _q+C_3\mathrm{cos}\varphi _p+D_3\mathrm{cos}(\varphi _p\varphi _q)}},`$
where the $`A`$’s, $`B`$’s, $`C`$’s, and $`D`$’s may have dependence on $`k^+`$, $`k_{}`$, $`p^+`$, $`p_{}`$, $`q^+=zE`$, and $`q_{}`$; they are independent of the azimuthal angles. These functions can be easily determined for each potential by examining the forms of the original equations. The rotational invariance of the potential about the three-axis allows the $`\varphi _k`$ integration to be done trivially.
In the integrand of $`I`$, only the last two terms depend on $`\varphi _p`$. To emphasize this, we write
$`I(k^+,k_{},q^+,q_{},p^+,p_{})`$ $`=`$ $`{\displaystyle _0^{2\pi }}𝑑\varphi _q{\displaystyle \frac{I_2(k^+,k_{},𝒒,p^+,p_{})}{(A_1+B_1\mathrm{cos}\varphi _q)(A_2+B_2\mathrm{cos}\varphi _q)(A_3+B_3\mathrm{cos}\varphi _q)}}`$ (D5)
$`I_2(k^+,k_{},𝒒,p^+,p_{})`$ $`=`$ $`{\displaystyle _0^{2\pi }}{\displaystyle \frac{d\varphi _p}{(1+a_2\mathrm{cos}\varphi _p+b_2\mathrm{sin}\varphi _p)(1+a_3\mathrm{cos}\varphi _p+b_3\mathrm{sin}\varphi _p)}}`$ (D6)
where, for $`i=2,3`$,
$`a_i`$ $`=`$ $`{\displaystyle \frac{C_i+D_i\mathrm{cos}\varphi _q}{A_i+B_i\mathrm{cos}\varphi _q}}`$ (D7)
$`b_i`$ $`=`$ $`{\displaystyle \frac{D_i\mathrm{sin}\varphi _q}{A_i+B_i\mathrm{cos}\varphi _q}}.`$ (D8)
The integral in $`I_2`$ is evaluated to obtain
$`I_2(k^+,k_{},𝒒,p^+,p_{})`$ $`=`$ $`{\displaystyle \frac{2\pi }{(a_2a_3)^2+(b_2b_3)^2(a_2b_3a_3b_2)^2}}`$ (D10)
$`\times \left({\displaystyle \frac{a_2(a_2a_3)+b_2(b_2b_3)}{\sqrt{1a_2^2b_2^2}}}+{\displaystyle \frac{a_3(a_3a_2)+b_3(b_3b_2)}{\sqrt{1a_3^2b_3^2}}}\right)`$
The remaining three-dimensional loop integral in $`V_{\text{TBE}}`$ on $`𝒒`$ is done using numeric techniques. The trick introduced in Appendix B to convert the semi-infinite $`q_{}`$ integration into an integration on a compact range. Before doing the $`z`$ integral, the range of integration is limited by using the step functions. Gauss-Legendre quadrature, given by Numerical Recipes in C , is used to evaluate all the integrals.
Since each of the parts of the full TBE potential (TBE:SB, TBE:SX, …) should be hermitian and invariant under interchange of particle 1 and 2, these invariances can be used as a self-consistency check. Each matrix element is calculated twice, first by using the straightforward approach, then particle labels 1 and 2 are interchanged and it is calculated again. The results are compared, and if they differ by an unacceptable amount, the number of quadrature points is increased and the element is recalculated. In order to get the numerical accuracy of the potentials correct to within 1%, we start with ten points for the $`q_{}`$ integral, six points for the $`\varphi _q`$ integral, and three points for the $`z`$ integral, resulting in a three-dimensional integral using 180 points.
## E Check of the Uncrossed Approximation
In this section, we want to check that how well the approximation
$`K_{\text{cross}}`$ $``$ $`K_{\text{ladder}}G_CK_{\text{ladder}},`$ (E1)
works. Since we are using a Hamiltonian theory and are interested in the potentials, we compare the potentials defined by
$`V_{\text{TBE:X}}`$ $`=`$ $`{\displaystyle \frac{1}{E}}g_0^1G_0K_{\text{cross}}G_0g_0^1`$ (E2)
$`V_{\text{TBE:UX}}`$ $`=`$ $`{\displaystyle \frac{1}{E}}g_0^1G_0K_{\text{ladder}}G_CK_{\text{ladder}}G_0g_0^1.`$ (E3)
The notation used here is defined in sections IV C and IV D.
The TBE crossed potential (TBE:X) can be written as
$`V_{\text{TBE:X}}`$ $`=`$ $`V_{\text{TBE:SX}}+V_{\text{TBE:TX}}+V_{\text{TBE:WX}}+V_{\text{TBE:ZX}},`$ (E4)
where the potentials on the right-hand side are defined in section III B. Calculation of the TBE approximate uncrossed potential (TBE:UX) is straightforward, but tedious. We find that
$`V_{\text{TBE:UX}}(E;𝒌_1,𝒑_1)`$ $`=`$ $`{\displaystyle \frac{1}{2}}[V_{\text{TBE:UX1}}(E;𝒌_1,𝒑_1)+V_{\text{TBE:UX1}}(E;𝒑_1,𝒌_1)`$ (E6)
$`+V_{\text{TBE:UX2}}(E;𝒌_1,𝒑_1)+V_{\text{TBE:UX2}}(E;𝒑_1,𝒌_1)],`$
where
$`V_{\text{TBE:UX1}}(E;𝒌_1,𝒑_1)`$ $`=`$ $`\left({\displaystyle \frac{M}{E}}\right)^4{\displaystyle }{\displaystyle \frac{d^2q_{}}{2(2\pi )^3}}[{\displaystyle _0^1}dz{\displaystyle \frac{\theta (xz)\theta (zy)}{(xz)z^2(zy)}}`$ (E10)
$`\times {\displaystyle \frac{1}{Ep_1^{}k_2^{}\omega ^{}(𝒒_1𝒑_1)\omega ^{}(𝒌_1𝒒_1)}}`$
$`\times \left({\displaystyle \frac{1}{Eq_1^{}k_2^{}\omega ^{}(𝒌_1𝒒_1)}}\right)^2]`$
$`+\{12\}`$
$`V_{\text{TBE:UX2}}(E;𝒌_1,𝒑_1)`$ $`=`$ $`\left({\displaystyle \frac{M}{E}}\right)^4{\displaystyle }{\displaystyle \frac{d^2q_{}}{2(2\pi )^3}}[{\displaystyle _0^1}dz{\displaystyle \frac{\theta (xz)\theta (yz)}{(xz)z^2(yz)}}`$ (E14)
$`\times {\displaystyle \frac{1}{Eq_1^{}p_2^{}\omega ^{}(𝒑_1𝒒_1)}}`$
$`\times \left({\displaystyle \frac{1}{Eq_1^{}k_2^{}\omega ^{}(𝒌_1𝒒_1)}}\right)^2]`$
$`+\{12\}.`$
Now consider the azimuthal-angle and loop integrals for these potentials. The approach used is similar to that of section D. Analyzing the potentials reveals that each of the terms can be written in the following schematic form,
$`V_{\text{TBE:UX,i}}(k^+,k_{};p^+,p_{})`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{q_{}dq_{}}{2(2\pi )^3}}{\displaystyle _0^1}𝑑zJ(k^+,q^+,p^+)I_i(k^+,k_{},q^+,q_{},p^+,p_{})`$ (E15)
$`I_i(k^+,k_{},q^+,q_{},p^+,p_{})`$ $`=`$ $`{\displaystyle _0^{2\pi }}{\displaystyle \frac{d\varphi _qd\varphi }{A_{1,i}+B_{1,i}\mathrm{cos}\varphi _q+C_{1,i}\mathrm{cos}\varphi }}\left({\displaystyle \frac{1}{A_2+B_2\mathrm{cos}\varphi _q}}\right)^2,`$ (E16)
where $`\varphi =\varphi _q+\varphi _p`$, and the $`A`$’s, $`B`$’s, and $`C`$’s may have dependence on $`k^+`$, $`k_{}`$, $`p^+`$, $`p_{}`$, $`q^+=zE`$, and $`q_{}`$; they are independent of the azimuthal angles. These factors can be easily determined for each potential by examining the forms of the original equations. The rotational invariance of the potential about the three-axis allows the $`\varphi _k`$ integration to be done trivially. The $`\varphi `$ integral is easily done to obtain
$`I(k^+,k_{},q^+,q_{},p^+,p_{})`$ $`=`$ $`2\pi {\displaystyle _0^{2\pi }}𝑑\varphi _q{\displaystyle \frac{1}{\sqrt{a_i^2C_{1,i}^2}}}\left({\displaystyle \frac{1}{A_2+B_2\mathrm{cos}\varphi _q}}\right)^2,`$ (E17)
where
$`a_i`$ $`=`$ $`A_{1,i}+B_{1,i}\mathrm{cos}\varphi _q.`$ (E18)
Further simplification is possible for $`V_{\text{TBE:UX2}}`$, since for that potential $`B_{1,2}=0`$,
$`I(k^+,k_{},q^+,q_{},p^+,p_{})`$ $`=`$ $`{\displaystyle \frac{(2\pi )^2A_2}{\sqrt{A_{1,2}^2C_{1,2}^2}\sqrt{A_2^2B_2^2}(A_2^2B_2^2)}}.`$ (E19)
The techniques discussed in the previous section are used to do the remaining loop integrals.
The spectra for the OBE+TBE:SB+TBE:UX potential can be calculated and compared to the OBE+TBE:SB+TBE:X potential (which is the same as the OBE+TBE potential), the modified-Green’s-function potential, and the ladder plus crossed Bethe-Salpeter equation. The spectra are plotted in Fig. 6. The spectra for the TBE:UX, TBE:X and BSE all lie close to each other, which shows that that the uncrossed approximation is valid.
This shows that the important approximation in the modified-Green’s-function approach is not the uncrossed approximation, but the addition of the extra interaction terms in Eq. (118) which serve to mimic the higher order interactions. |
warning/0002/astro-ph0002172.html | ar5iv | text | # A wavelet analysis of QSO spectra
## 1 Introduction
Resonant absorption by neutral hydrogen in the intergalactic medium along the line of sight to a distant quasar is responsible for the many absorption lines seen in the Ly$`\alpha `$ forest, blueward of the quasar’s Ly$`\alpha `$ emission line (Bahcall & Salpeter 1965, Gunn & Peterson 1965; see Rauch 1998 for a review). The general properties of these Ly$`\alpha `$ absorption lines are remarkably well reproduced by hydrodynamic simulations of cold dark matter (CDM) dominated cosmologies (Cen et al. 1994, Zhang, Anninos & Norman 1995, Miralda-Escudé et al. 1996, Hernquist et al. 1996, Wadsley & Bond 1996, Zhang et al. 1997, Theuns et al. 1998).
On large scales where pressure is unimportant, gas traces the dark matter and the Ly$`\alpha `$ spectrum can be used to infer the underlying density perturbations in the dark matter (Croft et al. 1997, Nusser & Haehnelt 1999). On small scales however, pressure gradients oppose the infall of gas into small potential wells (Jeans smoothing), leaving the absorber more extended in space than the underlying dark matter. The width of the absorption line is then determined by residual Hubble expansion across the absorber (Hernquist et al. 1996), Jeans smoothing and thermal broadening. Theuns, Schaye & Haehnelt (2000) analysed various line broadening mechanisms and demonstrated the importance of the gas temperature in controlling the line-widths.
The strong dependence of the small-scale properties of the Ly$`\alpha `$ forest on the temperature of the gas allows one to reconstruct the thermal evolution of the IGM. The gas temperature is set by the balance between adiabatic cooling caused by expansion and photo-heating by the UV-background. This introduces a tight relation between density and temperature, $`T=T_0(\rho /\rho )^{\gamma 1}`$ (Hui & Gnedin 1997). The parameters $`T_0`$ and $`\gamma `$ of this ‘equation of state’ are very sensitive to the reionization history of the IGM (Haehnelt & Steinmetz 1998). This is because thermal time scales are long in the low density IGM probed by the Ly$`\alpha `$ forest, hence that gas retains a memory about the past history of the ionising background. Consequently, the Ly$`\alpha `$ forest provides us with a fossil record of the history of reionization, which can be explored by unravelling its thermal history as deduced from the Ly$`\alpha `$ forest.
Schaye et al. (1999, see also Ricotti, Gnedin & Shull 2000) developed and tested a method to infer $`T_0`$ and $`\gamma `$ based on the line-widths of the absorption lines. Applying this method to high resolution QSO spectra for a range of redshifts, they found (Schaye et al. 2000) that the temperature $`T_0`$ decreases with decreasing redshift as expected, however, there is a large increase in $`T_0`$ round $`z=3`$, together with a decrease in the value of $`\gamma `$. They attributed this change in the equation of state to late reionization of helium II. They also noted that the temperature at higher redshifts is still fairly high, which might be an indication that we are approaching the epoch of hydrogen reionization.
The method of Schaye et al. to characterise line-widths is based on Voigt profile fitting of absorption lines (Webb 1987, Carswell et al. 1987). The rationale behind fitting absorption lines with a Voigt profile is partly historical, and stems from earlier theoretical models in which the forest was produced by a set of Ly$`\alpha `$ ‘clouds’ . The line-width of these absorbers was assumed to be set by thermal and ‘turbulent’ broadening, which would produce a Voigt profile, and line blending was responsible for the lines with large deviations from the Voigt profile. In the new paradigm of the Ly$`\alpha `$ forest absorption in the general IGM is responsible for lines, and there is no a priori reason to expect lines to have the Voigt shape.
In this paper we discuss a different method of characterising line-widths, based on discrete wavelets (see e.g. Press et al. 1992 for an introduction and further references). Wavelets provide an orthogonal basis for a unique decomposition of a signal (the spectrum) in terms of localised functions with a finite bandwidth. Thus they are a compromise between characterising a signal in terms of its individual pixel values and in terms of Fourier modes. In the first case, the characterisation has no information on correlations between different pixels (no frequency information) but perfect positional information. A Fourier decomposition, on the other hand, has perfect frequency information but no positional information. The analysis of a spectrum in terms of wavelets has the advantage that one can study the clustering of lines (‘positional information’), as a function of their widths (‘frequency information’).
The usage of wavelets to analyse QSO spectra was pioneered by Pando & Fang (1996, 1998), who used a wavelet analysis of Ly$`\alpha `$ absorption lines to describe the clustering of those lines. The wavelet analysis detected large scale structure in the Ly$`\alpha `$ forest, which had proved difficult using more traditional methods. In contrast to Pando & Fang, we will use wavelets to analyse the absorption spectrum directly, thereby eliminating the somewhat subjective step of first decomposing the continuous spectrum in absorption lines. The advantage of this new method is that it allows us to objectively characterise the typical width of absorption features as a function of position along the spectrum<sup>1</sup><sup>1</sup>1We will usually refer to absorption features as ‘lines’, but this is just a convenient name for what the eye picks out. The wavelet decomposition itself is unique and has no prejudice as to what should be considered a line..
We will show using hydrodynamic simulations that the probability distribution of wavelet amplitudes can be used to characterise the equation of state of the absorbing medium, in terms of the temperature at the mean density, $`T_0`$, and the slope, $`\gamma `$, of the temperature-density relation. In addition we use the fact that wavelets are localised in position along the spectrum, thereby allowing us to detect spatial variations in $`T_0`$ and/or $`\gamma `$, which might be present as a result of late helium II reionization or local effects.
This paper is organised as follows. In Section 2 we first give a brief description of the generation of mock spectra from our simulations and illustrate the decomposition of the spectra in discrete wavelets. The statistics of the wavelet amplitudes for different simulations is discussed in Section 3 and the results are summarised in Section 4. Recently, Meiksin (2000) discussed indepently the application of wavelets to QSO spectra.
## 2 Wavelet analysis of mock spectra
### 2.1 Mock spectra
We use the L1 simulation described before in Theuns et al. (2000). Briefly, this is a simulation of a flat, vacuum energy dominated cold dark matter model with matter density $`\mathrm{\Omega }_m=0.3`$, baryon fraction $`\mathrm{\Omega }_bh^2=0.019`$ and Hubble constant $`H_0=65`$ km s<sup>-1</sup> Mpc<sup>-1</sup>. Density fluctuations in this model are normalised to the abundance of galaxy clusters (Eke et al. 1996) and we have used CMBFAST (Seljak & Zaldarriaga 1996) to compute the appropriate linear transfer function. The IGM in this model is photo-ionised and photo-heated by the UV-background from QSOs, as computed by Haardt & Madau (1996).
We simulated this cosmological model with a modified version of the HYDRA simulation code (Couchman et al. 1995), which combines hierarchical P3M gravity (Couchman 1991) with smoothed particle hydrodynamics (SPH, Lucy 1977, Gingold & Monaghan 1977). We simulate a periodic, cubic box of size 7.7 co-moving Mpc using 128<sup>3</sup> particles of each species, which gives us sufficient resolution to compute line-widths reliably (Theuns et al. 1998). To investigate other effects, we also make use of simulations of a model with the same numerical resolution, cosmology and thermal history, but with a smaller box size (3.8 Mpc), and a set of simulations with a smaller normalisation $`\sigma _8=0.775`$ and $`\sigma _8=0.4`$.
In the analysis stage, we impose a particular equation of state on the gas at low overdensities ($`\rho /\rho <20`$) of the form $`T=T_0(\rho /\rho )^{\gamma 1}`$, varying the values of $`T_0`$ and $`\gamma `$. We then compute mock spectra that mimick the actual observed HIRES spectrum of the $`z_{\mathrm{em}}=3.0`$ QSO 1107+485, discussed by Rauch et al. (1997), using the following procedure. We divide the observed spectrum in three redshifts bins, $`z=2.52.625`$, $`z=2.6252.875`$ and $`z=2.8753`$ and scale the mean absorption of the simulations at $`z=2.5`$, $`z=2.75`$ and $`z=3`$ to the corresponding observed value. The simulated spectra are resampled to the observed resolution, and convolved with a Gaussian to mimick instrumental broadening. We have analysed the noise statistics of the QSO 1107 spectrum as a function of flux, and add noise with these properties to the simulated spectra. By randomly combining individual sight lines through the simulation volume, we generate a single long spectrum of length 37 492 km s<sup>-1</sup>. Velocity $`v`$ is related to redshift $`z`$ via $`vc\left[log_e(1+z)log_e(1+z_1)\right]`$, where $`c`$ is the speed of light, $`z`$ is redshift and $`z_1`$ is the redshift where Ly$`\alpha `$ starts to be confused with Ly$`\beta `$ for QSO 1107. In order to perform the wavelet analysis, we resample the spectrum to $`2^{15}`$=32768 pixels, equally spaced in velocity. In what follows, we will refer to a simulation with a particular equation of state by giving $`T_0/10^4`$K and $`\gamma `$, so the model $`(1.5,5/3)`$ has the imposed equation of state $`T=1.5\times 10^4(\rho /\rho )^{2/3}`$. We will present results for four equations of state, using $`T_0=1.5`$ and $`2.2\times 10^4`$K and $`\gamma =1`$ and 5/3.
### 2.2 Wavelets
The decomposition of a mock spectrum in terms of discrete wavelets is unique, once a particular wavelet basis has been chosen. Here we will use the Daubechies 20 wavelet (Daubechies 1988; see e.g. Press et al. 1992 for a general discussion on wavelets, and an example of the Daubechies 20 wavelet). Just as fast Fourier transforms, (discrete) wavelets come in powers of two, but unlike Fourier modes, a given wavelet has finite bandwidth and hence corresponds to a range of frequencies. Nevertheless we will refer to a wavelet of a particular ‘resolution’, for example quoting its full width at half maximum. The simulated spectrum has a length of $`V=`$37492 km s<sup>-1</sup> and the wavelet resolutions correspond to $`2^{i15}\times V`$. Here we will use the exponent $`i`$ to refer to wavelets of a particular resolution, e.g. $`i=9`$ corresponds to a wavelet of width 18.3 km s<sup>-1</sup>. Analysing a signal in terms of the amplitudes of wavelets with different resolutions was pioneered in a different context by Mallat (1989).
An example of a wavelet decomposition of a simulated spectrum is shown in Figure 1. Using wavelets with only four resolutions ($`i=912`$) already gives a relatively good description of the strong absorption features in the spectrum. Note how every line in the top panel is ‘detected’ on most resolution levels, indicating that each individual absorption line is also made-up of a range of frequencies. This is of course because these lines are relatively well approximated by Voigt profiles, which also have extended bandwidth. However, some lines are only weakly detected in the $`i=9`$ narrow wavelet, while some of the narrower lines lead to large amplitudes at this high resolution. It is this feature, namely that some narrow lines are picked-up strongly by the narrow wavelets while the broader lines are not, that allows us to characterise objectively the typical line-widths of absorption lines.
For a smaller value $`T_0`$ of the IGM temperature, there will be a larger fraction of narrow lines in the absorption spectrum. For a given pixel at velocity $`v`$ in the spectrum, let
$$𝒜(v;i,W)_{vW/2}^{v+W/2}A(v;i)^2𝑑v/W$$
(1)
denote the mean rms amplitude of the wavelet at resolution $`i`$, box-car smoothed over a window of size $`W`$ (km s<sup>-1</sup>). We will usually drop the indices $`i`$ and $`v`$ in what follows, and assume $`i=9`$ unless stated otherwise. For a spectrum with a larger fraction of narrow lines, $`𝒜`$ will be larger on average, hence we can in principle use the statistics of $`𝒜`$ as a measure of $`T_0`$, once the relation between them is calibrated with simulations.
In addition to this mean trend, $`𝒜`$ will fluctuate along the spectrum, due to (random) fluctuations in the strengths of lines. Here we give an example showing that averaging $`A^2`$ over a relatively short stretch of spectrum is already enough to distinguish between models with different $`T_0`$. This suggests it might be possible to detect fluctuations in $`T_0`$ (and $`\gamma `$), which might be a relic of a recent epoch of reionization or local effects. We will present a more detailed analysis of how this can be done below and restrict ourselves here to a typical example illustrated in Figure 2. To make the shown spectrum, we have combined spectra of the $`(1.5,5/3)`$ model on scales of 6000 km s<sup>-1</sup> with spectra of the 30 per cent hotter model $`(2.2,5/3)`$, into one long spectrum of length $`V`$. (In what follows, we will refer to this model as the mixed-temperature model.) The rms amplitude $`𝒜(v;9,1000)`$ of the $`i=9`$ (18.3 km s<sup>-1</sup>) wavelet, smoothed on 1000 km s<sup>-1</sup>, is sufficiently different between these two equations of state that stretches of the colder model can readily be distinguished from the hotter one as regions with larger $`𝒜`$.
In this example, both models have been scaled independently to have the same mean optical depth, corresponding to the observed value for QSO 1107. In reality, regions of higher temperature would tend to have smaller optical depth because of the $`T^{0.7}`$ temperature dependence of the recombination coefficient. This would tend to decrease the amplitude of the wavelets in the hotter regions even more, making it easier to distinguish between hot and cold regions.
## 3 Wavelet statistics
### 3.1 measuring the equation of state
In the previous section we showed that the rms amplitude of the $`i=9`$ narrow wavelet is strongly anti-correlated with the temperature of the absorbing gas. Consequently we can characterise the temperature distribution of the IGM over the spectrum using the corresponding distribution of wavelet amplitudes. For each of 100 realisations of models with a specified equation of state, we have computed the cumulative distribution of $`𝒜`$,
$$C(<𝒜)=_0^𝒜P(𝒜)𝑑𝒜,$$
(2)
where $`P(𝒜)`$ is the probability distribution of $`𝒜`$, and we plot the mean over 100 realisations, $`\overline{C}(<𝒜)`$, in figures 3 and 4 for $`W=500`$ and 2000 km s<sup>-1</sup>, respectively.
As expected, the colder models are systematically shifted to larger values of $`𝒜`$, since they contain a large number of narrow lines and consequently have larger values of $`𝒜`$. Note, however, that the dependence on the slope $`\gamma `$ is also quite strong, but this may be partly a consequence of using the mean density as the pivot point around which we change the slope. We have also superposed the mixed-temperature model, which stays close to the hot component for small values of $`𝒜`$ before veering away to the locus of the cold component for large values of the amplitude.
Having shown that the mean cumulative distribution $`\overline{C}(𝒜)`$ depends on the equation of state, we now want to characterise how well different models can be distinguished from each other, based on a single spectrum. Hence, we want to characterise to what extent the cumulative distribution $`C_j(𝒜)`$ for a single spectrum of model $`j`$ differs from the mean, $`\overline{C}_i`$, for model $`i`$. To this end, we compute the dispersion
$$\sigma _{ij}^2_0^{\mathrm{}}(\overline{C}_i(𝒜)C_j(𝒜))^2𝑑𝒜.$$
(3)
For a single realisation of a spectrum of model $`j`$, $`\sigma _{ij}^2`$ is just a number. In order to be able to distinguish between two models $`i`$ and $`j`$ based on a single spectrum, it is necessary that the dispersion $`\sigma _{ii}^2`$ be much smaller than the mean difference $`\sigma _{ij}^2`$ between the models.
Figure 5 shows the cumulative probability distribution $`C(>\sigma _{ij}^2)`$ for $`W=500`$ for three different equations of state. The confidence level at which a single spectrum of the model with equation of state say (1.5,5/3) (model $`j`$) can be distinguished from the model with equation of state (1.5,1) (model $`i`$) can be directly read-off from this figure. For example, in $`>95`$ per cent of cases $`\sigma _{ij}^2>0.004`$ for $`ij`$, whereas in only 2 per cent of cases, a model which really has the equation of state (1.5,1) will differ from the mean of this model to such a large extent.
A more usual statistic to judge whether a single realisation of a model is drawn from a given probability distribution is the Kolmogorov-Smirnov test, based on the maximum absolute difference $`dC=\mathrm{max}|\overline{C}_i(𝒜)C_j(𝒜)|`$ between two cumulative distributions. The cumulative distribution of the KS-statistic is shown in figure 6, where we compare it for models (1.5,5/3) and (2.2,1), which resemble each other most in figure 5. For 20 (5) per cent of realisations of model (1.5,5/3), $`dC>0.1`$ ($`dC>0.12`$), and at this level of contamination, 60 (40) per cent of models (2.2,1) have $`dC`$ larger than that.
Finally, figure 7 illustrates how well the mixed-temperature model can be distinguished from either the cold or the hot model with $`\gamma =5/3`$. This model is most likely mistaken with the colder single temperature counterpart. In 70 (25) per cent of cases, the mixed model has $`\sigma _{ij}^2>0.004`$ ($`\sigma _{ij}^2>0.01)`$. This happens for the cold model in only 10 (5) per cent of realisations.
### 3.2 other effects
Absorption features are broader in models with a smaller amplitude of the dark matter fluctuations (Theuns et al. 2000), thereby resembling more clustered but hotter models. This may lead to a degeneracy between $`T_0`$ and $`\sigma _8`$ (Bryan & Machacek 1999; note that Theuns, Schaye & Haehnelt (2000) showed that their Voigt profile analysis does not suffer from such a degeneracy). For the statistic presented here, this degeneracy is not very strong, as shown in figure 8. The model with $`\sigma _8=0.775`$ does not differ much from its more clustered counterpart with $`\sigma _8=0.9`$. Only for very low levels of clustering, $`\sigma _8=0.4`$, is the effect important. All models have been scaled to a mean effective optical depth of 0.26 at a redshift $`z=3`$.
Finally we have investigated the influence of the small box size in our numerical simulations, and the result is shown in figure 9. Lack of long wavelength perturbations decreases the observed range in $`𝒜`$, as expected, but the effect of this purely numerical artifact is relatively weak.
## 4 Conclusions
Clues to the thermal history of the Universe are hidden in the small scale structure of the Ly$`\alpha `$ forest. There are two reasons for this. Firstly, the widths of absorption lines are very sensitive to the temperature of gas, and secondly, thermal time scales are long in the low-density IGM that is responsible for the Ly$`\alpha `$ forest. Since the temperature of the photo-ionised IGM is determined by the evolution of the ionising background, unravelling the thermal history will have the added benefit of putting strong limits on the sources of UV light at high redshifts.
We have presented a new way of analysing the small scale structure of the Ly$`\alpha `$ forest, based on the unique decomposition of a spectrum in discrete wavelets. We have shown that the rms amplitude $`A^2`$ of narrow wavelets (18.3 km <sup>-1</sup>) correlates strongly with the temperature of the IGM, and also depends on the slope of the equation of state. We have quantified to what extent different models can be distinguished, using statistics of $`A^2`$
Our mock spectra have been designed to mimick an observed spectrum of QSO 1107+485 as much as possible. In particular, we have imposed on our simulated spectra the same large scale optical depth fluctuations as are observed in QSO 1107, making our mock spectra quite realistic. Even so, we can still easily distinguish between models that differ in temperature by less than 30 per cent. We have quantified the dependence of these statistics on numerical artifacts (missing long wavelength perturbations due to the smallness of our simulation box) and on the amplitude of the dark matter fluctuations ($`\sigma _8`$).
Wavelets are also localised in space, making it possible to study $`T_0`$ and $`\gamma `$ as a function of position along the spectrum. We characterised the extent to which we can distinguish models with a single value of $`T_0`$ from a model with temperature fluctuations, as might result from late reionization or local effects.
## Acknowledgments
We acknowledge simulating discussion with Martin Haehnelt, Michael Rauch, Joop Schaye and Simon White, This work has been supported by the ‘Formation and Evolution of Galaxies’ network set up by the European Commission under contract ERB FMRX-CT96086 of its TMR programme. This research was conducted in cooperation with Silicon Graphics/Cray Research utilising the Origin 2000 super computer at DAMTP, Cambridge. |
warning/0002/hep-th0002131.html | ar5iv | text | # Untitled Document
hep-th/0002131 DUK-CGTP-00-04, IASSNS–HEP–99/117
The Structure of the D0-D4 Bound State
Savdeep Sethi<sup>1</sup> sethi@sns.ias.edu and Mark Stern<sup>2</sup> stern@math.duke.edu
$``$ School of Natural Sciences, Institute for Advanced Study, Princeton, NJ 08540, USA
$``$ Department of Mathematics, Duke University, Durham, NC 27706, USA
We derive a set of equations for the wavefunction describing the marginal bound state of a single D0-brane with a single D4-brane. These are equations determining the vacuum of an $`N=8`$ abelian gauge theory with a charged hypermultiplet. We then solve these equations for the most general possible zero-energy solution using a Taylor series. We find that there are an infinite number of such solutions of which only one must be normalizable. We explore the structure of a normalizable solution under the assumption of an asymptotic expansion. Even the leading terms in the asymptotic series, which should reflect the supergravity solution, are unusual. Through the $`Spin(5)`$ flavor symmetry, the modes which are massive at long distance actually influence the leading behavior. Lastly, we show that the vacuum equations can quite remarkably be reduced to a single equation involving one unknown function. The resulting equation has a surprisingly simple and suggestive form.
2/00
1. Introduction
A single D0-brane and a single D4-brane form a marginal bound state . The low-energy dynamics of a D0-brane in the presence of a D4-brane is described by a quantum mechanical Yang-Mills theory with eight supercharges. The theory has a $`U(1)`$ vector multiplet coupled to a charged hypermultiplet . With a single D4-brane, there is only a Coulomb branch. The same quantum mechanics appears in the problem of counting H-monopole ground states in the toroidally compactified heterotic string . While the structure of vacuum wavefunctions in marginally bound systems is typically very difficult to analyze, this particular theory has a number of simplifying features. The aim of this paper is to study the vacuum wavefunction of this $`0+1`$-dimensional gauge theory with eight supercharges.
Our goal is to gain insight into a number of issues. For example, how do we go about uncovering the structure of threshold wavefunctions? It is not even clear how to formulate reasonable questions about a system as complex as the quantum mechanics describing many D0-branes. Another major issue is how the full quantum mechanics resolves the singularity of the moduli space metric. The vector multiplet contains five scalars $`x^\mu `$. For large $`r=|x|`$, the effective action describing the Coulomb branch dynamics should be a reasonable description of the physics . The metric on the Coulomb branch is protected by supersymmetry and takes the form,
$$ds^2=\left(\frac{1}{g^2}+\frac{1}{r^3}\right)(dx)^2,$$
$`(1.1)`$
where $`g^2`$ is the Yang-Mills coupling constant. We can express $`g^2`$ in terms of the type IIA string scale $`M_s`$ and coupling constant $`g_s`$,
$$g^2=g_sM_s^3.$$
This is the only scale in the theory. In the following sections, we set $`g^2=1`$ for simplicity. The tube-like metric (1.1) has a singularity at $`r=0`$ which is resolved by the full quantum mechanics. Metrics with a similar structure appear in D1-D5 systems.
In the following section, we present the supercharges and describe the symmetries of the problem. The flavor symmetry is $`Spin(5)\times SU(2)_R`$, and the unique vacuum is invariant under this symmetry . In section three, we derive the general form of a gauge invariant and flavor invariant wavefunction. A general wavefunction depends on $`11`$ functions of two variables, $`r`$ and $`y`$. As above, $`r`$ is a radial coordinate for the $`5`$ scalars $`x^\mu `$ of the vector multiplet. The hypermultiplet has $`4`$ scalars $`q_i`$, which parametrize the massive directions. We take $`y=|q|`$. We then derive a set of differential equations that any zero energy wavefunction must obey.
In section four, we analyze the implications of these differential equations. We can immediately reduce the number of unknown functions from $`11`$ to $`7`$. These $`7`$ functions satisfy $`14`$ first order coupled partial differential equations in $`2`$ variables. We point out some intriguing features of these equations: in particular, there is an interesting formal method of reducing the number of functions and equations in which the harmonic oscillator plays a key role. This method for collapsing the system of equations has similarities to the technique used to find non-renormalization theorems . We, however, proceed in a different direction. We solve the $`14`$ equations exactly by Taylor expanding in the $`y`$ variable. This is an expansion in the massive directions. The first interesting point is that the differential equations alone do not determine the solution. There are actually an infinite number of zero energy solutions. The ‘gauge’ degrees of freedom are not a finite set of parameters as we might have expected, but an entire function of degrees of freedom.
Essentially, from the perspective of the Taylor series, the problem boils down to solving $`2`$ ordinary differential equations in $`3`$ unknown functions. All the other terms in the wavefunction are determined by these $`3`$ functions of $`r`$. The condition that must uniquely specify the solution is normalizability. Rather remarkably, this global condition is strong enough to fix an entire arbitrary function. This seems to hint that the kind of principle that should underly M theory involves global rather than local constraints. In a vague sense, this is reminiscent of holography.
In section five, we turn to the practical problem of determining the normalizable solution. It turns out to be difficult to implement the global constraint of normalizability in any nice way. Instead, we expand the solution in an asymptotic series. This is akin to solving the M theory equations of motion for the geometry of an M5-brane in a derivative expansion. We find that even the leading terms in the solution are unusual. These terms should match a supergravity analysis. However, the structure of these terms is strongly dictated by invariance under the $`Spin(5)`$ flavor symmetry. Invariance under $`Spin(5)`$ is not a statement about long distance, moduli space physics. It is a statement that requires knowledge of both long and short distance physics because the $`Spin(5)`$ generators act on both the massless and massive degrees of freedom. Somewhat contrary to intuition, we find that the vacuum state for the massive degrees of freedom at large $`r`$ is a sum of $`3`$ representations of $`Spin(5)`$: a spherically symmetric $`\mathrm{𝟏}`$, a $`\mathrm{𝟓}`$ and a $`\mathrm{𝟏𝟒}`$.
The leading terms in the bound state wavefunction $`\mathrm{\Psi }`$ have the form,
$$\mathrm{\Psi }\frac{1}{r^3}e^{\frac{ry^2}{2}}|b_1>+\frac{x^\mu }{r^4}e^{\frac{ry^2}{2}}|b_2>^\mu +\frac{x^\mu x^\nu }{r^5}e^{\frac{ry^2}{2}}|b_3>^{\mu \nu },$$
$`(1.2)`$
where the $`|b_i>`$ are constructed from fermions. It seems unlikely that this asymptotic form could have been determined from low-energy considerations alone. In this sense, the massive degrees of freedom, through the $`Spin(5)`$ flavor symmetry, are important even at arbitrarily long distances. The structure of $`\mathrm{\Psi }`$ in (1.2) really begs for an interpretation both in terms of the supergravity solution of the D0-D4 brane , and in terms of the DLCQ description of an M5-brane via Matrix theory .
We proceed to compute the general form of $`\mathrm{\Psi }`$ in an asymptotic expansion. This takes us well beyond supergravity. The corrections to the leading terms (1.2) take the form of a perturbation series in the coupling constant $`g^2`$. Each correction depends on an a priori unknown constant, and we give a prescription for determining this constant. This amounts to summing up the higher derivative corrections to the supergravity solution for an M5-brane. Is the asymptotic solution actually convergent as $`r0`$? We know of no non-perturbative effects in the abelian gauge theory that could be relevant at short distances. However, this does not prove that the solution is convergent. It would be interesting to sum up a sufficient number of terms in the asymptotic series to see whether the solution is well-behaved as $`r`$ becomes small. This would clarify how the singularity in the metric is resolved from the perspective of a derivative expansion.
We note that finding the bound state wavefunction in an asymptotic series is much like trying to understand M theory in a derivative expansion. In section six, we present a quite different reduction of our initial $`14`$ vacuum equations, one that perhaps an M theorist might use. The result is quite incredible. The entire problem reduces to solving a scalar equation of the form,
$$\left(\mathrm{\Delta }+\stackrel{}{B}+W\right)u=0,$$
$`(1.3)`$
where $`\mathrm{\Delta }=_r^2+_y^2`$ and $`u`$ is a particular combination of the functions that appear in the bound state wavefunction. The vector field $`\stackrel{}{B}`$ and potential $`W`$ are rational functions of $`r`$ and $`y`$. Equation (1.3) is both simple and highly suggestive. It would be very interesting to find the solution to (1.3) either analytically or numerically. Supergravity and the structure of the derivative expansion should emerge from the asymptotics of the resulting solution for the bound state wavefunction. It would also be extremely interesting to generalize this reduction to non-abelian gauge theories. This might possibly help us understand how to define M theory through Matrix theory . It does not seem too unlikely to us that the ability to ‘deprolong’ the initial vacuum equations to get (1.3), as described in section six, is tied to supersymmetry and invariance theorems .
There are many additional directions to explore. Some of the simplifying features of the D0-D4 system remain when we add more hypermultiplets. However, there will now be a Higgs branch and the zero energy wavefunctions will spread in an interesting way onto the Higgs branch. Turning on $`B`$-fields makes the gauge theory on the D4-brane non-commutative . Certain choices of $`B`$-field should change the asymptotic behavior from polynomial decay to exponential decay. It would also be interesting to actually match the asymptotic structure of the bound state wavefunction with higher derivative corrections to the supergravity solution, like those generated by the $`R^4`$ terms \[11,,12,,13\].
2. The D0-D4 Quantum Mechanics
2.1. The vector multiplet supercharge
The D0-D4 system is obtained by dimensionally reducing $`N=1`$ abelian Yang-Mills with a single charged hypermultiplet from six dimensions. The symmetry group consists of the R-symmetries<sup>1</sup> The symmetry group, including both gauge and flavor symmetries, is not globally a product. There are discrete identifications. However, for this analysis we only need the Lie algebra generators so we can ignore global identifications. $`Spin(5)\times SU(2)_RSp(2)\times Sp(1)_R`$. The Hamiltonian is invariant under the symmetry group while the eight real supercharges transform in the $`(\mathrm{𝟒},\mathrm{𝟐})`$ representation.
Let us begin with the vector multiplet which contains the five scalars $`x^\mu `$ transforming in the $`(\mathrm{𝟓},\mathrm{𝟏})`$ of the symmetry group. Let $`p^\mu `$ be the associated canonical momenta obeying,
$$[x^\mu ,p^\nu ]=i\delta ^{\mu \nu }.$$
$`(2.1)`$
Associated to these bosons are eight real fermions $`\lambda _a`$ where $`a=1,\mathrm{},8`$ transforming in the $`(\mathrm{𝟒},\mathrm{𝟐})`$ representation of the symmetry group. These fermions obey the usual quantization relation,
$$\{\lambda _a,\lambda _b\}=\delta _{ab}.$$
$`(2.2)`$
Let $`\gamma ^\mu `$ be hermitian real gamma matrices which obey,
$$\{\gamma ^\mu ,\gamma ^\nu \}=2\delta ^{\mu \nu }.$$
$`(2.3)`$
An explicit basis for these gamma matrices along with a discussion of the symmetry group action is given in Appendix A.
To write the vector multiplet supercharge, we introduce an auxiliary field $`D`$ which transforms as $`(\mathrm{𝟏},\mathrm{𝟑})`$ under the symmetry group. The $`D`$-term is independent of $`x^\mu `$. The vector multiplet supercharge is given by:
$$Q_a^v=\left(\gamma ^\mu p^\mu \lambda \right)_a+D_{ab}\lambda _b.$$
$`(2.4)`$
The real anti-symmetric matrix $`D`$ commutes with $`\gamma ^\mu `$ because the $`Sp(1)_R`$ and $`Sp(2)`$ actions commute. The $`D`$-term must also satisfy,
$$D_{ac}D_{bc}=\delta _{ab}|D|^2.$$
$`(2.5)`$
It is then not hard to check that:
$$\{Q_a^v,Q_b^v\}=\delta _{ab}\left\{p^2+|D|^2\right\}.$$
$`(2.6)`$
Under a symmetry transformation $`(U,g)Sp(2)\times Sp(1)_R`$, we note that
$$\begin{array}{cc}\hfill \gamma ^\mu p^\mu U\gamma ^\mu p^\mu U^1,\lambda Ug\lambda ,DgDg^1,& \end{array}$$
$`(2.7)`$
so that
$$Q^vUgQ^v.$$
$`(2.8)`$
2.2. The hypermultiplet supercharge
A hypermultiplet contains four real scalars which we can package into a quaternion $`q`$ with components $`q^i`$ where $`i=1,2,3,4`$. This field transforms as $`(\mathrm{𝟏},\mathrm{𝟐})`$ under the symmetry group. We again introduce canonical momenta $`p_i`$ satisfying the usual commutation relations.
The hypermultiplet is charged under the $`U(1)`$ gauge symmetry so we need to determine how $`U(1)`$ acts on $`q`$. The $`q^i`$ parametrize $`\mathrm{IR}^4`$ so the symmetry group acting on the hypermultiplet must sit inside,
$$SO(4)Sp(1)_L\times Sp(1)_R.$$
Gauge transformations and $`Sp(1)_R`$ transformations commute. Therefore, the $`U(1)`$ gauge symmetry sits inside $`Sp(1)_L`$. We choose to generate the gauge symmetry by left multiplication on $`q`$ by $`I`$. The hermitian generator of gauge transformations on the bosons is then given by,
$$G_b=W_{12}+W_{34},$$
$`(2.9)`$
where
$$W_{ij}=q_ip_jq_jp_i.$$
$`(2.10)`$
The superpartner to $`q`$ is a real fermion $`\psi _a`$ with $`a=1,\mathrm{},8`$ satisfying,
$$\{\psi _a,\psi _b\}=\delta _{ab},$$
$`(2.11)`$
and transforming in the $`(\mathrm{𝟒},\mathrm{𝟏})`$ representation. In terms of the $`s^j`$ operators given in Appendix A, the free hypermultiplet charge takes the form
$$Q_a^{h_f}=s_{ab}^j\psi _bp_j.$$
$`(2.12)`$
Note that since the $`s^j`$ implement right multiplication by a quaternion, they commute with $`\gamma ^\mu `$. This free charge obeys the algebra,
$$\{Q_a^{h_f},Q_b^{h_f}\}=\delta _{ab}p_ip_i.$$
Invariance of (2.12) under the $`U(1)`$ gauge symmetry requires that
$$G_f=\frac{i}{2}s_{ab}^2\psi _a\psi _b$$
$`(2.13)`$
generate gauge transformations on $`\psi `$. The total generator of the $`U(1)`$ gauge symmetry is then given by,
$$G=G_b+G_f=W_{12}+W_{34}\frac{i}{2}\psi s^2\psi .$$
$`(2.14)`$
The full hypermultiplet supercharge $`Q^h`$ also includes couplings to the vector multiplet,
$$Q_a^h=s_{ab}^j\psi _bp_j+(\gamma ^\mu s^js^2)_{ab}\psi _bx^\mu q_j.$$
$`(2.15)`$
The form of the interaction term in (2.15) is fixed up to an overall constant by symmetry. The $`s^2`$ appearing in the interaction term is needed to ensure that $`Q^h`$ is gauge-invariant. The charge obeys the algebra:
$$\{Q_a^h,Q_b^h\}=\delta _{ab}\left\{p_ip_i+|x|^2|q|^2\frac{i}{4}x^\mu \psi \gamma ^\mu s^2\psi \right\}+2\gamma _{ab}^\mu x^\mu G.$$
$`(2.16)`$
As we expect, the supersymmetry algebra only closes on the Hamiltonian up to gauge transformations.
2.3. The coupled system
The full supercharge $`Q`$ is the given by,
$$Q=Q^v+Q^h,$$
$`(2.17)`$
where we define the $`D`$-term in the following way: using the components $`q^i`$ and matrices $`s^i`$, we can define a quaternion and its conjugate which act by right multiplication on $`\lambda `$:
$$\begin{array}{cc}\hfill q^R& =s^1q^1+s^2q^2+s^3q^3+s^4q^4\hfill \\ \hfill \overline{q}^R& =s^1q^1s^2q^2s^3q^3s^4q^4.\hfill \end{array}$$
We can then write the $`D`$-term in the form,
$$D_{ab}=\frac{1}{2}\left(q^Rs^2\overline{q}^R\right)_{ab}.$$
$`(2.18)`$
This $`D`$-term obeys (2.5),
$$\begin{array}{cc}\hfill \left(D^2\right)_{ab}& =\delta _{ab}|D|^2\hfill \\ & =\delta _{ab}\frac{1}{4}\left(q\overline{q}\right)^2.\hfill \end{array}$$
The full charge obeys the algebra:
$$\begin{array}{cc}\hfill \{Q_a,Q_b\}& =\delta _{ab}\left\{p^\mu p^\mu +|D|^2+p_ip_i+|x|^2|q|^2+\mathrm{}\right\}+2\gamma _{ab}^\mu x^\mu G\hfill \\ & =\delta _{ab}\mathrm{\hspace{0.17em}2}H+2\gamma _{ab}^\mu x^\mu G.\hfill \end{array}$$
$`(2.19)`$
The omitted terms are bilinears in the fermions whose exact form we will not need. The bosonic potential $`V`$ appearing in (2.19) is given by,
$$V=|x|^2|q|^2+\frac{1}{4}|q|^4.$$
$`(2.20)`$
Since we have coupled a single hypermultiplet to the $`U(1)`$ vector multiplet, the only flat direction is $`q=0`$ and there is no Higgs branch.
3. Deriving Equations for the Vacuum Wavefunction
To be consistent with predictions from string duality, there should be a unique vacuum wavefunction for this quantum mechanical gauge theory . An index argument proves that there is at least one normalizable vacuum wavefunction . Coupled with a recent invariance theorem , the index result implies that the ground state is unique.
On quantization, the fermions $`\lambda `$ and $`\psi `$ act as gamma matrices on a $`256`$-dimensional spinor wavefunction. A priori, the vacuum wavefunction then consists of $`256`$ complex functions of the $`9`$ bosonic variables $`x^\mu `$ and $`q^i`$. However, we can significantly simplify the problem by using symmetries. First note that any state in the Hilbert space $`|s>`$ must be gauge-invariant,
$$G|s>=0.$$
$`(3.1)`$
Further, all states can be grouped into representations of the global $`Sp(2)\times Sp(1)_R`$ symmetry group. The $`Sp(2)`$ is generated by the operators,
$$T^{\mu \nu }=X^{\mu \nu }\frac{i}{4}\gamma _{ab}^{\mu \nu }\left(\lambda _a\lambda _b+\psi _a\psi _b\right),$$
$`(3.2)`$
where
$$X^{\mu \nu }=x^\mu p^\nu x^\nu p^\mu .$$
$`(3.3)`$
The three generators of $`Sp(1)_R`$ correspond to right multiplication by $`I,J,K`$ and in accord with prior notation, we will denote them by $`\stackrel{~}{s}^i`$:
$$\begin{array}{cc}\hfill \stackrel{~}{s}^2& =W_{12}W_{34}+\frac{i}{2}\lambda s^2\lambda \hfill \\ \hfill \stackrel{~}{s}^3& =W_{13}+W_{24}+\frac{i}{2}\lambda s^3\lambda \hfill \\ \hfill \stackrel{~}{s}^4& =W_{14}W_{23}+\frac{i}{2}\lambda s^4\lambda .\hfill \end{array}$$
$`(3.4)`$
The unique ground state $`\mathrm{\Psi }`$ must be invariant under the actions of $`Q_a,G,T^{\mu \nu }`$ and $`\stackrel{~}{s}^i`$,
$$Q_a\mathrm{\Psi }=G\mathrm{\Psi }=T^{\mu \nu }\mathrm{\Psi }=\stackrel{~}{s}^i\mathrm{\Psi }=0.$$
These constraints are quite powerful; for example, they allow us to replace a differential operator $`X^{\mu \nu }`$ by an algebraic one $`\frac{i}{4}\gamma _{ab}^{\mu \nu }\left(\lambda _a\lambda _b+\psi _a\psi _b\right).`$ There are multiple ways to derive equations for the vacuum wavefunction. We will describe two approaches which we used to derive these equations.<sup>2</sup> Using two different approaches helped enormously in the search for errors.
3.1. Radial coordinates
In the first approach, we can rewrite the supercharge (2.17) in terms of radial coordinates:
$$r^2=|x|^2,y^2=|q|^2.$$
The charge takes the form,
$$\begin{array}{cc}\hfill Q_a& =(\gamma ^\mu \lambda )_a\frac{x^\mu }{r}p_r+(s^j\psi )_a\frac{q^j}{y}p_y+(\gamma ^\mu \lambda )_a\frac{x^\nu }{r^2}X^{\nu \mu }+(s^j\psi )_a\frac{q^k}{y^2}W_{kj}\hfill \\ & +(\gamma ^\mu s^js^2\psi )_ax^\mu q_j+(D\lambda )_a\hfill \\ & =(\gamma ^\mu \lambda )_a\frac{x^\mu }{r}p_r+(s^j\psi )_a\frac{q^j}{y}p_y+M_a.\hfill \end{array}$$
$`(3.5)`$
We have lumped the angular derivatives and non-derivative terms into the operator $`M_a`$. What is particularly nice about $`M_a`$ is that we can replace all the derivative operators by bilinears in fermions. As we noted before, $`Sp(2)`$ invariance allows us to replace $`X^{\nu \mu }`$ by
$$\frac{i}{4}\gamma _{ab}^{\nu \mu }\left(\lambda _a\lambda _b+\psi _a\psi _b\right).$$
However, we can also replace $`q^kW_{kj}`$ by a bilinear in fermions using $`Sp(1)_R`$ invariance. Using (3.4), we note that:
$$\begin{array}{cc}\hfill q^kW_{k1}& =\frac{i}{2}\left\{q^2\lambda s^2\lambda +q^3\lambda s^3\lambda +q^4\lambda s^4\lambda \right\}\hfill \\ \hfill q^kW_{k2}& =\frac{i}{2}\left\{q^4\lambda s^3\lambda q^1\lambda s^2\lambda q^3\lambda s^4\lambda \right\}\hfill \\ \hfill q^kW_{k3}& =\frac{i}{2}\left\{q^2\lambda s^4\lambda q^1\lambda s^3\lambda q^4\lambda s^2\lambda \right\}\hfill \\ \hfill q^kW_{k4}& =\frac{i}{2}\left\{q^3\lambda s^2\lambda q^1\lambda s^4\lambda q^2\lambda s^3\lambda \right\}.\hfill \end{array}$$
$`(3.6)`$
Therefore, $`M_a`$ is a completely algebraic operator.
Using the $`Sp(2)`$ symmetry, we can then rotate $`x`$ to the special point where $`x^10`$ and $`x^\mu =0`$ for $`\mu >1`$. Likewise, we can rotate $`q`$ using $`Sp(1)_R`$ to the point $`q^10`$ and $`q^i=0`$ for $`i>1`$. At this point, $`r=|x^1|`$ and $`y=|q^1|`$. Since all angular derivatives in $`Q_a`$ are replaced by algebraic operators, there is no difficulty in restricting $`\mathrm{\Psi }`$ to this point. The question of determining $`\mathrm{\Psi }`$ at this point then reduces to finding coupled differential equations in two variables. The form of $`\mathrm{\Psi }`$ at an arbitrary choice of $`x`$ and $`q`$ can then be obtained by applying the rotation generators (3.2) and (3.4).
3.2. Symmetries and the fermion Hilbert space
In the second approach which we will use for the rest of the paper, we will first solve the invariance conditions explicitly for the most general possible invariant wavefunction. As we will show, the most general wavefunction depends on $`11`$ functions of $`r`$ and $`y`$. We will then derive coupled equations for these functions from the requirement that the wavefunction have zero energy.
The first step is to construct the fermion Hilbert space. We need to complexify our real fermions and build a Fock space:
$$\begin{array}{cc}\hfill \sqrt{2}du_1=\lambda _1+i\lambda _2& \sqrt{2}dv_1=\psi _1+i\psi _2,\hfill \\ \hfill \sqrt{2}du_2=\lambda _3i\lambda _4& \sqrt{2}dv_2=\psi _3i\psi _4,\hfill \\ \hfill \sqrt{2}du_3=\lambda _5+i\lambda _6& \sqrt{2}dv_3=\psi _5+i\psi _6,\hfill \\ \hfill \sqrt{2}du_4=\lambda _7i\lambda _8& \sqrt{2}dv_4=\psi _7i\psi _8.\hfill \end{array}$$
$`(3.7)`$
It is natural to think of $`du_a`$ and $`dv_a`$ as one-forms obeying the relation,
$$\{du_a,du_b^{}\}=\delta _{ab}\{dv_a,dv_b^{}\}=\delta _{ab},$$
where $`a=1,\mathrm{},4`$. Wavefunctions in the Hilbert space are then $`(p,q)`$ forms where $`p`$ and $`q`$ are the $`du_a`$ and $`dv_a`$ degrees, respectively. We choose the Fock vacuum or $`(0,0)`$ form to satisfy,
$$du_a^{}|0>=dv_a^{}|0>=0.$$
Note that the complex conjugate of a $`(p,q)`$-form is a $`(4p,4q)`$-form. If the ground state is unique then it is bosonic so the form degree must be even.
With our choice of complexification (3.7), the $`Sp(2)`$ generators $`T^{\mu \nu }`$ acting on forms preserve degree. Actually, the generators preserve $`p`$ and $`q`$ separately so a $`(p,q)`$ form is mapped to a $`(p,q)`$ form. The $`Sp(2)`$ generators naturally split into commuting generators for an $`Sp(2)_b`$ acting on bosons and an $`Sp(2)_f`$ acting on fermions. In turn, the $`Sp(2)_f`$ splits into an $`Sp(2)_{f_p}`$ acting on $`du`$ with generators,
$$\frac{i}{4}\underset{a,b}{}\gamma _{ab}^{\nu \mu }\lambda _a\lambda _b,$$
and an $`Sp(2)_{f_q}`$ acting on $`dv`$ with generators,
$$\frac{i}{4}\underset{a,b}{}\gamma _{ab}^{\nu \mu }\psi _a\psi _b.$$
We can now employ some group theory to see how the various $`128`$ bosonic forms transform under $`Sp(2)_f`$. Let us start with the $`(p,0)`$ forms which appear in the following representations:
$`p`$ $`Sp(2)_f`$ rep. 0 $`\mathrm{𝟏}`$ 1 $`\mathrm{𝟒}`$ 2 $`\mathrm{𝟓}\mathrm{𝟏}`$ 3 $`\mathrm{𝟒}`$ 4 $`\mathrm{𝟏}`$
We wedge the (odd) even $`(p,0)`$ forms with the (odd) even $`(0,q)`$ forms to get the $`128`$ bosonic forms. The following representations appear from wedging even forms with even forms,
$$(\mathrm{𝟏})^{10}(\mathrm{𝟓})^6\mathrm{𝟏𝟎}\mathrm{𝟏𝟒},$$
while from wedging odd forms with odd forms, we find:
$$(\mathrm{𝟏})^4(\mathrm{𝟓})^4(\mathrm{𝟏𝟎})^4.$$
We can immediately discard forms transforming in the $`\mathrm{𝟏𝟎}`$ representation. A tensor say $`a^{\mu \nu }`$ transforming in the $`\mathrm{𝟏𝟎}`$ is antisymmetric in $`\mu ,\nu `$ so contraction with $`x^\mu x^\nu `$ to get a singlet of the full $`Sp(2)`$ gives zero.
Let us now constrain our Hilbert space further by imposing invariance under $`Sp(1)_R`$ and the gauge symmetry. We can rewrite the generators (3.4) in terms of our complex fermions:
$$\begin{array}{cc}\hfill \stackrel{~}{s}^2& =W_{12}W_{34}+\underset{a}{}du_adu_a^{}2\hfill \\ \hfill \stackrel{~}{s}^3& =W_{13}+W_{24}i\left(du_1du_2+du_1^{}du_2^{}+du_3du_4+du_3^{}du_4^{}\right)\hfill \\ \hfill \stackrel{~}{s}^4& =W_{14}W_{23}+\left(du_1du_2du_1^{}du_2^{}+du_3du_4du_3^{}du_4^{}\right).\hfill \end{array}$$
$`(3.8)`$
Likewise for the gauge symmetry,
$$G=W_{12}+W_{34}\underset{a}{}dv_adv_a^{}+2.$$
$`(3.9)`$
Note that the operator $`W_{ij}`$ has eigenvalues $`n`$ and $`n`$ with corresponding eigenfunctions,
$$z_{ij}^n=(q_i+iq_j)^n,\overline{z}_{ij}^n=(q_iiq_j)^n.$$
What does invariance under (3.8) and (3.9) imply? By taking the sum and difference of $`G`$ and $`\stackrel{~}{s}^2`$, we see that we should restrict to $`(p,q)`$ forms
$$(z_{12})^{\frac{qp}{2}}(z_{34})^{\frac{(p+q)}{2}2}|p,q>,$$
$`(3.10)`$
which can be multiplied by a function of both $`|z_{12}|^2`$ and $`|z_{34}|^2`$. The remaining two generators $`\stackrel{~}{s}^3`$ and $`\stackrel{~}{s}^4`$ change the value of $`p`$. It is natural to study the complex combinations $`\stackrel{~}{s}^3i\stackrel{~}{s}^4`$ and $`\stackrel{~}{s}^3+i\stackrel{~}{s}^4`$, which raise and lower the value of $`p`$:
$$\begin{array}{cc}\hfill s^+=\stackrel{~}{s}^3i\stackrel{~}{s}^4& =z_{12}p_{34}\overline{z}_{34}\overline{p}_{12}2i\left(du_1du_2+du_3du_4\right)\hfill \\ \hfill s^{}=\stackrel{~}{s}^3+i\stackrel{~}{s}^4& =\overline{z}_{12}\overline{p}_{34}z_{34}p_{12}2i\left(du_1^{}du_2^{}+du_3^{}du_4^{}\right),\hfill \end{array}$$
$`(3.11)`$
where $`p_{ij}=p_iip_j`$ and $`\overline{p}_{ij}=p_i+ip_j`$.
Again the generators of $`Sp(1)_R`$ split into an $`Sp(1)_b`$ acting on bosons and an $`Sp(1)_f`$ acting on fermions. It is easy to see how the $`(p,q)`$-forms fall into representations of $`Sp(1)_f`$. The three singlets under $`Sp(2)_{f_p}`$ denoted $`|0,q>,|2,q>_{\mathrm{𝟏}_p},|4,q>`$ transform in the $`\mathrm{𝟑}`$ of $`Sp(1)_f`$. For the choice $`q=0,4`$, we can construct one singlet under the full $`Sp(1)_R`$ which can be multiplied by an arbitrary function of $`y`$. For the case $`q=2`$, we can construct two singlets under $`Sp(1)_R`$ by tensoring with either the $`\mathrm{𝟏}`$ or the $`\mathrm{𝟓}`$ of $`Sp(2)_{f_q}`$. Let us denote the $`\mathrm{𝟓}`$ of $`Sp(2)_{f_q}`$ by $`|2>_{\mathrm{𝟓}_q}^\mu `$. The explicit $`Sp(1)_R`$ singlets are then given by the forms,
$$\begin{array}{cc}& \left\{\overline{z}_{34}^2+\overline{z}_{12}\overline{z}_{34}(du_1du_2+du_3du_4)+\overline{z}_{12}^2du_1du_2du_3du_4\right\}|0>,\hfill \\ & \{z_{12}\overline{z}_{34}+\frac{1}{2}(|z_{12}|^2|z_{34}|^2)(du_1du_2+du_3du_4)\overline{z}_{12}z_{34}du_1du_2du_3du_4\}\times \hfill \\ & (dv_1dv_2+dv_3dv_4)|0>,\hfill \\ & \left\{z_{12}\overline{z}_{34}+\frac{1}{2}(|z_{12}|^2|z_{34}|^2)(du_1du_2+du_3du_4)\overline{z}_{12}z_{34}du_1du_2du_3du_4\right\}|2>_{\mathrm{𝟓}_q}^\mu ,\hfill \\ & \left\{z_{12}^2z_{12}z_{34}(du_1du_2+du_3du_4)+z_{34}^2du_1du_2du_3du_4\right\}dv_1dv_2dv_3dv_4|0>.\hfill \end{array}$$
The $`\mathrm{𝟓}`$ of $`Sp(2)_{f_p}`$ denoted $`|2>_{\mathrm{𝟓}_p}^\mu `$ decomposes into five singlets under $`Sp(1)_f`$. The form $`|2>_{\mathrm{𝟓}_p}`$ can therefore only appear with a function of $`y`$. To satisfy the constraint (3.10), we must then tensor $`|2>_{\mathrm{𝟓}_p}^\mu `$ with a $`q=2`$ form constructed from $`dv`$. The two choices are either the $`\mathrm{𝟏}`$ or the $`\mathrm{𝟓}`$ of $`Sp(2)_{f_q}`$. This gives three additional possibilities denoted,
$$|2,2>_\mathrm{𝟏},|2,2>_\mathrm{𝟓}^\mu ,|2,2>_{\mathrm{𝟏𝟒}}^{\mu \nu },$$
where the subscript denotes the representation under the full $`Sp(2)_f`$. The construction of these forms is described in Appendix B.
Lastly, we need to consider the case of odd $`p`$. The $`|1>_{\mathrm{𝟒}_p}`$ and $`|3>_{\mathrm{𝟒}_p}`$ forms combine to form a doublet under $`Sp(1)_f`$. By tensoring with either the $`|1>_{\mathrm{𝟒}_q}`$ or the $`|3>_{\mathrm{𝟒}_q}`$ forms, we can construct the following four $`Sp(1)_R`$ invariants:
$$\begin{array}{cc}& \left\{z_{12}z_{34}(du_1du_2+du_3du_4)\right\}|1,3>_\mathrm{𝟏}\hfill \\ & \left\{z_{12}z_{34}(du_1du_2+du_3du_4)\right\}|1,3>_\mathrm{𝟓}^\mu \hfill \\ & \left\{\overline{z}_{34}+\overline{z}_{12}(du_1du_2+du_3du_4)\right\}|1,1>_\mathrm{𝟏}\hfill \\ & \left\{\overline{z}_{34}+\overline{z}_{12}(du_1du_2+du_3du_4)\right\}|1,1>_\mathrm{𝟓}^\mu .\hfill \end{array}$$
Again the subscript denotes the representation under the full $`Sp(2)_f`$. After imposing all the invariance constraints, we are therefore left with $`11`$ complex functions $`f_i=f_i(r,y)`$ appearing in the following way:
$$f_1|0,0>=f_1\left\{\overline{z}_{34}^2+\overline{z}_{12}\overline{z}_{34}(du_1du_2+du_3du_4)+\overline{z}_{12}^2du_1du_2du_3du_4\right\}|0>$$
$`(3.12)`$
$$\begin{array}{cc}\hfill f_2|0,4>=& f_2\left\{z_{12}^2z_{12}z_{34}(du_1du_2+du_3du_4)+z_{34}^2du_1du_2du_3du_4\right\}\hfill \\ & \times dv_1dv_2dv_3dv_4|0>\hfill \end{array}$$
$`(3.13)`$
$$f_3|2,2>_\mathrm{𝟏},f_4x^\mu |2,2>_\mathrm{𝟓}^\mu ,f_5x^\mu x^\nu |2,2>_{\mathrm{𝟏𝟒}}^{\mu \nu },$$
$`(3.14)`$
$$\begin{array}{cc}& f_6|1,3>=f_6\left\{z_{12}z_{34}(du_1du_2+du_3du_4)\right\}|1,3>_\mathrm{𝟏}\hfill \\ & f_7|1,1>=f_7\left\{\overline{z}_{34}+\overline{z}_{12}(du_1du_2+du_3du_4)\right\}|1,1>_\mathrm{𝟏}\hfill \end{array}$$
$`(3.15)`$
$$\begin{array}{cc}& f_8x^\mu |1,3>^\mu =f_8x^\mu \left\{z_{12}z_{34}(du_1du_2+du_3du_4)\right\}|1,3>_\mathrm{𝟓}^\mu \hfill \\ & f_9x^\mu |1,1>^\mu =f_9x^\mu \left\{\overline{z}_{34}+\overline{z}_{12}(du_1du_2+du_3du_4)\right\}|1,1>_\mathrm{𝟓}^\mu \hfill \end{array}$$
$`(3.16)`$
$$\begin{array}{cc}\hfill f_{10}|0,2>=& f_{10}\{z_{12}\overline{z}_{34}+\frac{1}{2}(|z_{12}|^2|z_{34}|^2)(du_1du_2+du_3du_4)\hfill \\ & \overline{z}_{12}z_{34}du_1du_2du_3du_4\}(dv_1dv_2+dv_3dv_4)|0>\hfill \end{array}$$
$`(3.17)`$
$$\begin{array}{cc}\hfill f_{11}x^\mu |0,2>^\mu =& f_{11}x^\mu \{z_{12}\overline{z}_{34}+\frac{1}{2}(|z_{12}|^2|z_{34}|^2)(du_1du_2+du_3du_4)\hfill \\ & \overline{z}_{12}z_{34}du_1du_2du_3du_4\}|2>_{\mathrm{𝟓}_q}^\mu .\hfill \end{array}$$
$`(3.18)`$
Our choice of normalization in constructing these forms is described in Appendix C. We take the ground state $`\mathrm{\Psi }`$ to be the sum of these eleven forms.
3.3. Dynamical constraints
What remains is to determine the consequences of the eight equations,
$$Q_a\mathrm{\Psi }=0.$$
$`(3.19)`$
First note that each term in $`Q_a`$ can be assigned a parity $`(\pm ,\pm )`$ according to whether it changes the parity of the wavefunction in $`(x,q)`$ respectively. For example, the term
$$(\gamma ^\mu s^js^2\psi )_ax^\mu q_j$$
$`(3.20)`$
has parity $`(,)`$ since it is odd in $`x`$ and odd in $`q`$. Likewise, each term in $`\mathrm{\Psi }`$ has a definite parity. We can therefore isolate all terms in $`Q_a\mathrm{\Psi }`$ with a definite parity. It is also sufficient to restrict to the case $`a=1`$ because our wavefunction $`\mathrm{\Psi }`$ is $`Sp(2)`$ invariant, but an $`Sp(2)`$ transformation rotates us from one choice of charge to another.
We can then ask: what combinations give terms with parity $`(+,)`$? A quick check of $`Q_a`$ acting on the possible forms composing $`\mathrm{\Psi }`$ gives the following equation,
$$\begin{array}{cc}& (s^j\psi )_ap_j\left\{f_1\right|0,0>+f_{10}|0,2>+f_2|0,4>+f_3|2,2>_\mathrm{𝟏}+f_5x^\mu x^\nu |2,2>_{\mathrm{𝟏𝟒}}^{\mu \nu }\}+\hfill \\ & (\gamma ^\mu s^js^2\psi )_ax^\mu q^j\left\{f_4x^\rho \right|2,2>_\mathrm{𝟓}^\rho +f_{11}x^\rho |0,2>^\rho \}\hfill \\ & +(\gamma ^\mu \lambda )_ap^\mu \left\{f_8x^\rho \right|1,3>^\rho +f_9x^\rho |1,1>^\rho \}+(D\lambda )_a\left\{f_6\right|1,3>+f_7|1,1>\}=0.\hfill \end{array}$$
$`(3.21)`$
The terms giving $`(,)`$ satisfy:
$$\begin{array}{cc}& (\gamma ^\mu s^js^2\psi )_ax^\mu q^j\{f_1|0,0>+f_{10}|0,2>+f_2|0,4>+f_3|2,2>_\mathrm{𝟏}+\hfill \\ & f_5x^\mu x^\nu |2,2>_{\mathrm{𝟏𝟒}}^{\mu \nu }\}+(s^j\psi )_ap_j\{f_4x^\rho |2,2>_\mathrm{𝟓}^\rho +f_{11}x^\rho |0,2>^\rho \}\hfill \\ & +(D\lambda )_a\left\{f_8x^\rho \right|1,3>^\rho +f_9x^\rho |1,1>^\rho \}+(\gamma ^\mu \lambda )_ap^\mu \left\{f_6\right|1,3>+f_7|1,1>\}=0.\hfill \end{array}$$
$`(3.22)`$
From $`(,+)`$, we find:
$$\begin{array}{cc}& (\gamma ^\mu \lambda )_ap^\mu \left\{f_1\right|0,0>+f_{10}|0,2>+f_2|0,4>+f_3|2,2>_\mathrm{𝟏}+f_5x^\mu x^\nu |2,2>_{\mathrm{𝟏𝟒}}^{\mu \nu }\}+\hfill \\ & (D\lambda )_a\left\{f_4x^\rho \right|2,2>_\mathrm{𝟓}^\rho +f_{11}x^\rho |0,2>^\rho \}+(s^j\psi )_ap_j\left\{f_8x^\rho \right|1,3>^\rho +f_9x^\rho |1,1>^\rho \}\hfill \\ & +(\gamma ^\mu s^js^2\psi )_ax^\mu q^j\left\{f_6\right|1,3>+f_7|1,1>\}=0.\hfill \end{array}$$
$`(3.23)`$
The last equation follows from considering the $`(+,+)`$ terms,
$$\begin{array}{cc}& (D\lambda )_a\left\{f_1\right|0,0>+f_{10}|0,2>+f_2|0,4>+f_3|2,2>_\mathrm{𝟏}+f_5x^\mu x^\nu |2,2>_{\mathrm{𝟏𝟒}}^{\mu \nu }\}+\hfill \\ & (\gamma ^\mu \lambda )_ap^\mu \left\{f_4x^\rho \right|2,2>_\mathrm{𝟓}^\rho +f_{11}x^\rho |0,2>^\rho \}+\hfill \\ & (\gamma ^\mu s^js^2\psi )_ax^\mu q^j\left\{f_8x^\rho \right|1,3>^\rho +f_9x^\rho |1,1>^\rho \}+\hfill \\ & (s^j\psi )_ap_j\left\{f_6\right|1,3>+f_7|1,1>\}=0.\hfill \end{array}$$
$`(3.24)`$
In each equation, we set $`a=1`$ as a first simplification. After evaluating angular derivatives, we are free to rotate $`q`$ using $`Sp(1)_R`$ so that $`q^10`$ and $`q^i=0`$ for $`i>1`$. In a similar way, we can consider the point $`x^10`$ with $`x^\mu =0`$ for $`\mu >1`$ after evaluating the $`x`$ angular derivatives. With this choice of coordinates, $`y=|q^1|`$ and $`r=|x^1|`$.
The first set of equations relate $`f_1,f_7,f_9,f_{10}`$ and $`f_{11}`$. These follow from considering the $`(4,1)`$ forms and the $`(3,0)`$ forms in (3.21), (3.22) and (3.23), (3.24) respectively:
$$\begin{array}{cc}& r\frac{f_9}{r}y\frac{f_1}{y}\frac{y^2}{2}f_7+5f_94f_1+2f_{10}=0,\hfill \\ & \frac{f_7}{r}+ry^2\left\{f_1\frac{f_9}{2}\right\}+2rf_{11}=0,\hfill \\ & r\frac{f_9}{y}+y\frac{f_1}{r}+ryf_7=0,\hfill \\ & \frac{f_7}{y}+r^2yf_9+\frac{y^3}{2}f_1=0.\hfill \end{array}$$
$`(3.25)`$
Likewise, by considering the $`(0,3)`$ forms in (3.21), (3.22) and the $`(1,4)`$ forms in (3.23), (3.24), we find the following equations relating $`f_2,f_6,f_8,f_{10}`$ and $`f_{11}`$:
$$\begin{array}{cc}& r\frac{f_8}{r}+y\frac{f_2}{y}+\frac{y^2}{2}f_6+5f_8+4f_22f_{10}=0,\hfill \\ & \frac{f_6}{r}+ry^2\left\{f_2+\frac{f_8}{2}\right\}+2rf_{11}=0,\hfill \\ & y\frac{f_2}{r}r\frac{f_8}{y}+ryf_6=0,\hfill \\ & \frac{f_6}{y}r^2yf_8+\frac{y^3}{2}f_2=0.\hfill \end{array}$$
$`(3.26)`$
Note that equations (3.26) are the same as (3.25) under the identification:
$$f_2f_1f_6f_7f_8f_9.$$
$`(3.27)`$
It is easy to check there are no non-vanishing $`(0,1),(4,3)`$ and $`(1,0),(3,4)`$ forms in (3.21), (3.22) and (3.23), (3.24). This leaves equations coming from forms with degree $`(2,1),(2,3)`$ in (3.21) and (3.22), and $`(1,2),(3,2)`$ in (3.23) and (3.24). Let us start by considering the $`(2,1)`$ parts of (3.21) which give the following constraints,
$$\begin{array}{cc}& f_{10}+\frac{y}{2}\frac{f_{10}}{y}2f_1+4f_9+r^2f_4+\frac{1}{y}(\frac{f_3}{y}+\frac{4}{5}r^2\frac{f_5}{y})+\frac{y^2}{2}r^2f_{11}=0,\hfill \\ & f_{10}+\frac{y}{2}\frac{f_{10}}{y}2f_1r^2f_4+r\frac{f_9}{r}+f_9\frac{1}{y}(\frac{f_3}{y}+\frac{4}{5}r^2\frac{f_5}{y})\hfill \\ & +\frac{y^2}{2}f_7+\frac{y^2}{2}r^2f_{11}=0,\hfill \\ & \frac{2}{5}\frac{r^2}{y}\frac{f_5}{y}\frac{2}{y}\frac{f_3}{y}r\frac{f_9}{r}3f_9+\frac{y^2}{2}f_7=0.\hfill \end{array}$$
$`(3.28)`$
The $`(2,3)`$ terms give the equations,
$$\begin{array}{cc}& f_{10}+\frac{y}{2}\frac{f_{10}}{y}2f_24f_8+r^2f_4+\frac{1}{y}(\frac{f_3}{y}+\frac{4}{5}r^2\frac{f_5}{y})+\frac{y^2}{2}r^2f_{11}=0,\hfill \\ & f_{10}+\frac{y}{2}\frac{f_{10}}{y}2f_2r\frac{f_8}{r}f_8r^2f_4\frac{1}{y}(\frac{f_3}{y}+\frac{4}{5}r^2\frac{f_5}{y})\hfill \\ & +\frac{y^2}{2}f_6+\frac{y^2}{2}r^2f_{11}=0,\hfill \\ & \frac{2}{5}\frac{r^2}{y}\frac{f_5}{y}\frac{2}{y}\frac{f_3}{y}+r\frac{f_8}{r}+3f_8+\frac{y^2}{2}f_6=0.\hfill \end{array}$$
$`(3.29)`$
Note that equations (3.29) are the same as (3.28) under the identification (3.27). This is a nice check that the equations are correct.
We next need the $`(2,1)`$ parts of (3.22) which give the relations,
$$\begin{array}{cc}& \frac{y^2}{2}f_{10}+f_3+\frac{4}{5}r^2f_5+\frac{1}{y}\frac{f_4}{y}+\frac{y}{2}\frac{f_{11}}{y}+f_{11}=0,\hfill \\ & \frac{y^2}{2}f_{10}f_3\frac{4}{5}r^2f_5\frac{1}{y}\frac{f_4}{y}+\frac{1}{r}\frac{f_7}{r}+\frac{y^2}{2}f_9+\frac{y}{2}\frac{f_{11}}{y}+f_{11}=0,\hfill \\ & 2f_3\frac{2}{5}r^2f_5\frac{1}{r}\frac{f_7}{r}+\frac{y^2}{2}f_9=0.\hfill \end{array}$$
$`(3.30)`$
From the $`(2,3)`$ components, we find the equations:
$$\begin{array}{cc}& \frac{y^2}{2}f_{10}+f_3+\frac{4}{5}r^2f_5+\frac{1}{y}\frac{f_4}{y}+\frac{y}{2}\frac{f_{11}}{y}+f_{11}=0,\hfill \\ & \frac{y^2}{2}f_{10}f_3\frac{4}{5}r^2f_5\frac{1}{y}\frac{f_4}{y}+\frac{1}{r}\frac{f_6}{r}\frac{y^2}{2}f_8+\frac{y}{2}\frac{f_{11}}{y}+f_{11}=0,\hfill \\ & 2f_3\frac{2}{5}r^2f_5\frac{1}{r}\frac{f_6}{r}\frac{y^2}{2}f_8=0.\hfill \end{array}$$
$`(3.31)`$
These equations are again consistent with (3.27).
We now turn to the $`(1,2)`$ parts of (3.23) which imply that,
$$\begin{array}{cc}& \frac{y^2}{2r}\frac{f_{10}}{r}+\frac{1}{r}\frac{f_3}{r}+\frac{4}{5}r\frac{f_5}{r}+\frac{28}{5}f_5+\frac{1}{2}y^2f_4+2f_9+\frac{y^4}{4}f_{11}=0,\hfill \\ & \frac{y^2}{2r}\frac{f_{10}}{r}\frac{1}{r}\frac{f_3}{r}\frac{4}{5}r\frac{f_5}{r}\frac{28}{5}f_5+\frac{1}{2}y^2f_4+y^2f_6\hfill \\ & +y\frac{f_8}{y}+2f_8\frac{y^4}{4}f_{11}=0,\hfill \\ & \frac{2}{r}\frac{f_3}{r}\frac{2}{5}r\frac{f_5}{r}\frac{14}{5}f_5y^2f_6+y\frac{f_8}{y}+2f_82f_9=0.\hfill \end{array}$$
$`(3.32)`$
The $`(3,2)`$ forms give the following set of equations:
$$\begin{array}{cc}& \frac{y^2}{2r}\frac{f_{10}}{r}+\frac{1}{r}\frac{f_3}{r}+\frac{4}{5}r\frac{f_5}{r}+\frac{28}{5}f_5+\frac{1}{2}y^2f_42f_8+\frac{y^4}{4}f_{11}=0,\hfill \\ & \frac{y^2}{2r}\frac{f_{10}}{r}\frac{1}{r}\frac{f_3}{r}\frac{4}{5}r\frac{f_5}{r}\frac{28}{5}f_5+\frac{1}{2}y^2f_4+y^2f_7y\frac{f_9}{y}\hfill \\ & 2f_9\frac{y^4}{4}f_{11}=0,\hfill \\ & \frac{2}{r}\frac{f_3}{r}\frac{2}{5}r\frac{f_5}{r}\frac{14}{5}f_5y^2f_7y\frac{f_9}{y}2f_9+2f_8=0.\hfill \end{array}$$
$`(3.33)`$
Again, (3.33) and (3.32) are identical under (3.27).
The $`(1,2)`$ parts of (3.24) give the following equations,
$$\begin{array}{cc}& \frac{y^4}{4}f_{10}+\frac{y^2}{2}f_3+\frac{2y^2r^2}{5}f_5+r\frac{f_4}{r}+5f_4+2f_7+\frac{y^2}{2}f_{11}+\frac{y^2r}{2}\frac{f_{11}}{r}=0,\hfill \\ & \frac{y^4}{4}f_{10}\frac{y^2}{2}f_3\frac{2y^2r^2}{5}f_5+r\frac{f_4}{r}+5f_4+y\frac{f_6}{y}+2f_6\hfill \\ & +r^2y^2f_8\frac{y^2}{2}f_{11}\frac{y^2r}{2}\frac{f_{11}}{r}=0,\hfill \\ & y^2f_3\frac{y^2r^2}{5}f_5y\frac{f_6}{y}2f_6+2f_7+r^2y^2f_8+y^2f_{11}=0.\hfill \end{array}$$
$`(3.34)`$
Lastly, the $`(3,2)`$ parts give the equations:
$$\begin{array}{cc}& \frac{y^4}{4}f_{10}+\frac{y^2}{2}f_3+\frac{2y^2r^2}{5}f_5+r\frac{f_4}{r}+5f_4+2f_6+\frac{y^2}{2}f_{11}+\frac{y^2r}{2}\frac{f_{11}}{r}=0,\hfill \\ & \frac{y^4}{4}f_{10}\frac{y^2}{2}f_3\frac{2y^2r^2}{5}f_5+r\frac{f_4}{r}+5f_4+y\frac{f_7}{y}+2f_7\hfill \\ & r^2y^2f_9\frac{y^2}{2}f_{11}\frac{y^2r}{2}\frac{f_{11}}{r}=0,\hfill \\ & y^2f_3\frac{y^2r^2}{5}f_5y\frac{f_7}{y}+2f_62f_7r^2y^2f_9+y^2f_{11}=0.\hfill \end{array}$$
$`(3.35)`$
Again, note that (3.34) and (3.35) are identical under the exchange (3.27).
4. The Structure of the Bound State Wavefunction
4.1. Reducing the number of functions
Initially, we have eleven independent functions obeying a set of coupled differential equations. To make progress, we need to reduce the number of functions in a systematic fashion. We can begin to whittle down the number of independent functions in the following way: the difference of the third equations of (3.30) and (3.31) together with (3.25) imply that,
$$f_1=f_2.$$
$`(4.1)`$
Let us now turn to (3.25). By taking $`_y`$ of the second equation and $`_r`$ of the fourth equation, we find two equations for $`_{ry}^2f_7`$. For these equations to be compatible, we require that:
$$f_{10}f_1+f_9+\frac{1}{y}\frac{f_{11}}{y}=0.$$
$`(4.2)`$
This equation is actually not a new addition to our list of constraints. It follows from (3.30) and (3.25). The same analysis applied to (3.26) gives,
$$f_{10}f_1f_8+\frac{1}{y}\frac{f_{11}}{y}=0.$$
$`(4.3)`$
These equations together require that,
$$f_8=f_9.$$
$`(4.4)`$
The difference of the fourth equations in (3.25) and (3.26) then implies the equivalence,
$$f_6=f_7.$$
$`(4.5)`$
In this way, we are reduced to eight functions $`\{f_1,f_3,f_4,f_5,f_7,f_9,f_{10},f_{11}\}`$. Half the equations we derived are now redundant since the symmetry (3.27) is an actual identity.
We obtain an algebraic relation between the remaining functions by using the third equation in (3.30) together with (3.25):
$$f_3+\frac{y^2}{2}f_1\frac{r^2}{5}f_5+f_{11}=0.$$
$`(4.6)`$
Using the algebraic relation, we can eliminate one function. The remaining $`7`$ functions of $`2`$ variables are constrained by $`14`$ equations, which is the minimal number we could have expected.
4.2. A first reduction of the equations
Let us summarize the equations that describe the bound state. We have taken some simpler linear combinations of the previous equations:
$$\begin{array}{cc}\hfill (1)& r\frac{f_9}{r}y\frac{f_1}{y}\frac{y^2}{2}f_7+5f_94f_1+2f_{10}=0,\hfill \\ \hfill (2)& \frac{f_7}{r}+ry^2\left\{f_1\frac{f_9}{2}\right\}+2rf_{11}=0,\hfill \\ \hfill (3)& r\frac{f_9}{y}+y\frac{f_1}{r}+ryf_7=0,\hfill \\ \hfill (4)& \frac{f_7}{y}+r^2yf_9+\frac{y^3}{2}f_1=0,\hfill \\ \hfill (5)& 2f_{10}+y\frac{f_{10}}{y}4f_1+5f_9+y^2r^2f_{11}+r\frac{f_9}{r}+\frac{y^2}{2}f_7=0,\hfill \\ \hfill (6)& 2r^2yf_4\frac{y^3}{2}f_7+3yf_9ry\frac{f_9}{r}+2\frac{f_3}{y}+\frac{8}{5}r^2\frac{f_5}{y}=0,\hfill \\ \hfill (7)& \frac{2}{5}\frac{r^2}{y}\frac{f_5}{y}\frac{2}{y}\frac{f_3}{y}r\frac{f_9}{r}3f_9+\frac{y^2}{2}f_7=0,\hfill \\ \hfill (8)& y^2f_{10}+\frac{y^2}{2}f_9+2f_{11}+y\frac{f_{11}}{y}+\frac{1}{r}\frac{f_7}{r}=0,\hfill \\ \hfill (9)& 2yf_3+\frac{8}{5}r^2yf_5\frac{y^3}{2}f_9+2\frac{f_4}{y}\frac{y}{r}\frac{f_7}{r}=0,\hfill \\ \hfill (10)& ry(f_4+f_7)r\frac{f_9}{y}+y\frac{f_{10}}{r}=0,\hfill \\ \hfill (11)& \frac{56}{5}rf_5ry^2f_7+4rf_9+\frac{1}{2}ry^4f_{11}+ry\frac{f_9}{y}+2\frac{f_3}{r}+\frac{8}{5}r^2\frac{f_5}{r}=0,\hfill \\ \hfill (12)& \frac{2}{r}\frac{f_3}{r}\frac{2}{5}r\frac{f_5}{r}\frac{14}{5}f_5y^2f_7y\frac{f_9}{y}4f_9=0,\hfill \\ \hfill (13)& 10f_4+4f_7r^2y^2f_9+\frac{y^4}{2}f_{10}+y\frac{f_7}{y}+2r\frac{f_4}{r}=0,\hfill \\ \hfill (14)& yf_3+\frac{4}{5}r^2yf_5+r^2yf_9+yf_{11}\frac{f_7}{y}+ry\frac{f_{11}}{r}=0.\hfill \end{array}$$
$`(4.7)`$
In these equations, we can remove one function using (4.6). This is a complicated set of coupled equations. To uncover the structure of the wavefunction, let us begin by simplifying as much as possible.
After staring at these equations for sometime, a pair of equations $`(3)`$ and $`(4)`$ in (4.7) appear distinguished. At our special point, these equations involve only the $`(4,0)`$ and $`(3,1)`$ forms given explicitly in Appendix C. They do not involve any $`(2,2)`$ forms. These are the analogues of the ‘top forms’ which played a crucial role in proving non-renormalization theorems . We can eliminate $`f_9`$ from these two equations giving,
$$\begin{array}{cc}\hfill \left\{_y^2+\frac{1}{y}_y+r^2y^2\right\}f_7& =y^2\left\{1+\frac{y}{2}_yr_r\right\}f_1,\hfill \\ & =y^2S(f_1).\hfill \end{array}$$
$`(4.8)`$
Note that if we set the source term $`S(f_1)=0`$, then the homogeneous solutions for $`f_7`$ are $`e^{\pm ry^2/2}`$. The plus sign is not normalizable. However, both solutions suffer from a more serious problem. The ground state wavefunction must be a smooth function. It must have a convergent Taylor series about the origin. This implies that each $`f_i`$ must be function of $`r^2`$ and $`y^2`$ near the origin. The homogeneous solutions alone are therefore ruled out.
What this teaches us is that $`f_7`$ is determined in terms of $`f_1`$. It is convenient to make the following redefinition,
$$f_7=y^2\stackrel{~}{f}_7.$$
Note that $`\stackrel{~}{f}_7`$ can have a $`1/y^2`$ term near the origin. The equation (4.8) now takes the form,
$$\begin{array}{cc}\hfill \left\{_y^2\frac{3}{y}_y+r^2y^2\right\}\stackrel{~}{f_7}& =S(f_1).\hfill \end{array}$$
$`(4.9)`$
The left hand side of (4.9) is the Hamiltonian for a four-dimensional harmonic oscillator! Now there is a very pretty collapse. Equation $`(4)`$ in (4.7) determines $`f_9`$ in terms of $`f_1`$. Likewise $`(1)`$ and $`(2)`$ determine $`f_{10}`$ and $`f_{11}`$. Equation $`(10)`$ determines $`f_4`$, while $`(14)`$ determines $`f_3`$. The algebraic constraint (4.6) then fixes $`f_5`$. Clearly, there many other ways to collapse the problem. The main point is that all the remaining functions are given in terms of $`f_1`$. This leaves $`8`$ equations which must determine $`f_1`$.
4.3. Solving for $`f_7`$
We can now express $`f_7`$ in terms of $`f_1`$ in the following way:
$$f_7=f_7^0(r^2)e^{ry^2/2}+y^2\left\{_y^2\frac{3}{y}_y+r^2y^2\right\}^1S(f_1).$$
$`(4.10)`$
The $`y=0`$ component $`f_7^0`$ is determined by requiring that $`f_7`$ be smooth, as we discussed previously. Smoothness of $`f_7`$ also requires that $`f_7^0`$ be a function of $`r^2`$. In turn, we can expand $`f_1`$ as follows:
$$f_1=\underset{n=0}{\overset{\mathrm{}}{}}f_1^n(r)|n>.$$
$`(4.11)`$
The $`|n>`$ are radial harmonic oscillator eigenstates which obey,
$$\left\{_y^2\frac{3}{y}_y+r^2y^2\right\}|n>=E_n|n>,$$
where $`E_n=4(n+1)r`$. The construction and properties of these eigenstates are described in Appendix D. These eigenstates have the nice feature that acting on $`|n>`$, the operators $`_r,y_y`$ and $`y^2`$ involve only $`|n1>,|n>,|n+1>`$. Using the relations from Appendix D, we see that the source term has a beautifully simple form:
$$S(f_1)=\underset{n}{}\left(f_1^nr_rf_1^n\right)|n>.$$
$`(4.12)`$
It is now easy to solve for $`f_7`$ in terms of $`|n>`$,
$$f_7=f_7^0(r^2)|0>+y^2\underset{n=0}{\overset{\mathrm{}}{}}\frac{1}{E_n}\left(f_1^nr_rf_1^n\right)|n>.$$
$`(4.13)`$
We have left $`y^2`$ in (4.13) for later convenience. By considering the coefficient of $`y^2`$ in (4.13) and imposing smoothness, we obtain the following relation:
$$\frac{r}{2}f_7^0+\underset{n=0}{\overset{\mathrm{}}{}}\frac{1}{E_n}\left(1r_r\right)f_1^n=0.$$
$`(4.14)`$
Since $`f_7^0`$ is non-singular as $`r0`$, we obtain the sum rule
$$\underset{n=0}{\overset{\mathrm{}}{}}\frac{1}{4(n+1)}f_1^n(0)=0,$$
$`(4.15)`$
and the relation:
$$f_7^0=\frac{2}{r^2}\underset{n=0}{\overset{\mathrm{}}{}}\frac{1}{4(n+1)}\left(1r_r\right)f_1^n.$$
$`(4.16)`$
Note that these sum rules may be largely formal since we do not know whether the sums are absolutely convergent.
4.4. Equations for the physics near the flat directions
Instead of proceeding to reduce the number of functions, let us take a different tack. The most interesting physics in this problem occurs in a neighbourhood of the flat directions. So let us consider a Taylor expansion about the flat direction $`y=0`$. This approach turns out to be more useful than reducing the number of functions, which increases the complexity of the resulting equations. We expand each $`f_i`$ in the following way,
$$f_i=t_i^0(r^2)+t_i^2(r^2)y^2+t_i^4(r^2)y^4+\mathrm{}.$$
$`(4.17)`$
The algebraic constraint (4.6) together with the equations of (4.7) give the following set of relations on the $`t_i^0`$:
$$\begin{array}{cc}\hfill (1)& t_3^0\frac{r^2}{5}t_5^0+t_{11}^0=0,\hfill \\ \hfill (2)& r_rt_9^0+5t_9^04t_1^0+2t_{10}^0=0,\hfill \\ \hfill (3)& _rt_7^0+2rt_{11}^0=0,\hfill \\ \hfill (4)& \frac{56}{5}rt_5^0+4rt_9^0+2_rt_3^0+\frac{8}{5}r^2_rt_5^0=0,\hfill \\ \hfill (5)& \frac{2}{r}_rt_3^0\frac{2}{5}r_rt_5^0\frac{14}{5}t_5^04t_9^0=0,\hfill \\ \hfill (6)& 10t_4^0+4t_7^0+2r_rt_4^0=0,\hfill \\ \hfill (7)& t_4^0+2t_7^0+\frac{1}{r}_r\left\{t_1^0+t_{10}^0\right\}=0.\hfill \end{array}$$
$`(4.18)`$
This gives $`7`$ equations for $`8`$ unknown functions. As we might have expected, this is not sufficient to determine the flat direction physics without input from higher $`y`$ terms. We can similarly derive equations involving only $`t_i^2`$ and $`t_i^0`$ which are given in Appendix E.
Note the critical observation that we can solve for all $`t_i^2`$ in terms of the $`t_i^0`$ using $`8`$ of the equations of $`E.1`$. There are no new independent functions at order $`y^2`$ in the Taylor expansion. We can express the $`t_i^2`$ in terms of the $`t_i^0`$ and at most their second derivatives using the first $`7`$ equations and equation $`(10)`$ of $`E.1`$,
$$\begin{array}{cc}\hfill (1)& t_7^2=\frac{1}{2}r^2t_9^0,\hfill \\ \hfill (2)& t_9^2=\frac{1}{2r}\left\{_rt_1^0+rt_7^0\right\},\hfill \\ \hfill (3)& t_{11}^2=\frac{1}{4}\left\{3t_9^0+r_rt_9^02t_1^0\right\},\hfill \\ \hfill (4)& t_5^2=\frac{1}{2r}\left\{rt_4^0_rt_9^0\right\},\hfill \\ \hfill (5)& t_3^2=\frac{1}{20}\left\{2r^2t_4^015t_9^03r_rt_9^0\right\},\hfill \\ \hfill (6)& t_4^2=\frac{1}{20r}\left\{10rt_3^0+8r^3t_5^05_rt_7^0\right\},\hfill \\ \hfill (7)& t_{10}^2=\frac{1}{16}\left\{6r^2t_{11}^0+\frac{4}{r}_rt_1^0+r_rt_7^0+_r^2t_1^0\right\},\hfill \\ \hfill (8)& t_1^2=\frac{1}{16r}\left\{8rt_7^0+2r^3t_{11}^0+4_rt_1^0+r^2_rt_7^0+r_r^2t_1^0\right\}.\hfill \end{array}$$
$`(4.19)`$
What we need to check is whether any of the remaining $`7`$ equations of $`E.1`$ are new relations. After some algebra that we will spare the reader, it turns out that all the remaining equations of $`E.1`$ are consequences of (4.18). This is not too surprising. The interactions involve order $`y^4`$ terms so we should not be able to completely determine the physics on the flat directions without expanding to higher order in $`y`$. So let us expand to order $`y^4`$ giving equations relating the $`t_i^4,t_i^2`$ and $`t_i^0`$. These equations are again given explicitly in Appendix E.
The equations of $`E.2`$ are very similar to $`E.1`$, except for mixing with some $`t_i^0`$ terms through the $`y^3`$ and $`y^4`$ interactions in (4.7). It is easy to see that we can again solve for all the $`t_i^4`$ in terms of the $`t_i^0`$. It should be clear that all the coefficients of the higher $`y^{2n}`$ terms in the Taylor expansion are determined in terms of the $`t_i^0`$. The problem is then to determine the $`t_i^0`$, and we need one more relation in addition to those of (4.18). Again using the first $`7`$ equations and equation $`(10)`$ of $`E.2`$, we can solve for the $`t_i^4`$:
$$\begin{array}{cc}\hfill (1)& t_7^4=\frac{1}{4}\left\{r^2t_9^2+\frac{1}{2}t_1^0\right\},\hfill \\ \hfill (2)& t_9^4=\frac{1}{4r}\left\{_rt_1^2+rt_7^2\right\},\hfill \\ \hfill (3)& t_{11}^4=\frac{1}{16}\left\{8t_9^2+2r_rt_9^28t_1^2+\frac{1}{r}_rt_1^0\right\},\hfill \\ \hfill (4)& t_5^4=\frac{1}{16r^3}\left\{rt_7^04r^3t_4^2+2rt_9^2+_rt_1^0+4r^2_rt_9^2\right\},\hfill \\ \hfill (5)& t_3^4=\frac{1}{80r}\left\{rt_7^04r^3t_4^238rt_9^24_rt_1^06r^2_rt_9^2\right\},\hfill \\ \hfill (6)& t_4^4=\frac{1}{80}\left\{5t_9^020t_3^216r^2t_5^2+\frac{10}{r}_rt_7^2\right\},\hfill \\ \hfill (7)& t_{10}^4=\frac{1}{200}\left\{20r^2t_1^0+40r^2t_3^28r^4t_5^25t_7^2+\frac{20}{r}_rt_1^2+5r_rt_7^2+5_r^2t_1^2\right\},\hfill \\ \hfill (8)& t_1^4=\frac{1}{200r}\{5r^3t_1^010r^3t_3^2+2r^5t_5^2+45rt_7^2+20_rt_1^2\hfill \\ & +5r^2_rt_7^2+5r_r^2t_1^2\}.\hfill \end{array}$$
$`(4.20)`$
Note that the expression for $`t_7^4`$ agrees with the expression coming from the earlier relation (4.8) that we derived between $`f_1`$ and $`f_7`$. Once again we are left with $`7`$ additional equations. It turns out that these additional equations again give no new relations. After checking higher order Taylor coefficients, we find that there are no further relations. It appears that any choice of $`t_i^0`$ satisfying (4.18) give a zero energy solution. However, most of these solutions are not normalizable.
It is straightforward to derive the general recursion relation for $`t_i^n`$ in terms of lower Taylor coefficients:
$$\begin{array}{cc}\hfill (1)& t_7^n=\frac{1}{n}\left\{r^2t_9^{n2}+\frac{1}{2}t_1^{n4}\right\},\hfill \\ \hfill (2)& t_9^n=\frac{1}{nr}\left\{_rt_1^{n2}+rt_7^{n2}\right\},\hfill \\ \hfill (3)& t_{11}^n=\frac{1}{2}\left\{\frac{1}{2}t_9^{n2}t_1^{n2}\frac{1}{r}_rt_7^n\right\},\hfill \\ \hfill (4)& t_5^n=\frac{1}{n}\left\{\frac{1}{r}_rt_9^{n2}t_4^{n2}\right\},\hfill \\ \hfill (5)& t_3^n=\left\{\frac{r^2}{5}t_5^nt_{11}^n\frac{1}{2}t_1^{n2}\right\},\hfill \\ \hfill (6)& t_4^n=\frac{1}{2n}\left\{\frac{1}{r}_rt_7^{n2}+\frac{1}{2}t_9^{n4}2t_3^{n2}\frac{8}{5}r^2t_5^{n2}\right\},\hfill \\ \hfill (7)& t_{10}^n=\frac{1}{2n(n+6)}\left\{(n+8)t_7^{n2}+10nt_9^n+r^2(2n+8)t_{11}^{n2}+2nr_rt_9^n\right\},\hfill \\ \hfill (8)& t_1^n=\frac{1}{n}\left\{t_7^{n2}+nt_{10}^n+r^2t_{11}^{n2}\right\}.\hfill \end{array}$$
$`(4.21)`$
Let us close this discussion of the Taylor expansion by pointing out that the entire Taylor series depends only on $`3`$ of the $`t_i^0`$. To see this, let us express (4.18) in a form which will be more convenient for later manipulation,
$$\begin{array}{cc}\hfill (1)& t_5^0=\frac{5}{r^2}\left(t_3^0+t_{11}^0\right),\hfill \\ \hfill (2)& t_{11}^0=\frac{1}{2r}_rt_7^0,\hfill \\ \hfill (3)& t_7^0=\frac{1}{2}\left(5t_4^0+r_rt_4^0\right),\hfill \\ \hfill (4)& t_9^0=\frac{5}{6r}_rt_3^0,\hfill \\ \hfill (5)& t_{10}^0=2t_1^0\frac{5}{2}t_9^0\frac{r}{2}_rt_9^0,\hfill \\ \hfill (6)& 20_r(r^3t_3^0)+r\left\{72_r+33r_r^2+3r^2_r^3\right\}t_4^0=0,\hfill \\ \hfill (7)& _rt_1^0+\frac{1}{48r^4}\{64r^5+72_r16r^6_r72r_r^2\hfill \\ & +12r^2_r^3+12r^3_r^4+r^4_r^5\}t_4^0=0.\hfill \end{array}$$
$`(4.22)`$
The first $`5`$ equations of (4.22) are definitions of various $`t_i^0`$. The last $`2`$ are relations on the $`3`$ independent functions $`t_1^0,t_3^0`$ and $`t_4^0`$.
5. An Asymptotic Expansion of the Wavefunction
5.1. Matching Taylor and oscillator expansions
It is hard to see how to implement the normalizability condition in a Taylor series. So although we have found the zero energy solution in terms of the $`t_i^0`$, we need a practical procedure to construct the $`t_i^0`$. We will determine the $`t_i^0`$ under the assumption that the bound state wavefunction admit an asymptotic expansion in powers of $`1/r`$. The primary motivation for studying the asymptotic expansion is that the asymptotic form should be interpretable in terms of supergravity plus higher derivative corrections. What we do not know is whether the asymptotic expansion converges to the actual bound state wavefunction. This issue is closely related to the following two questions: if we sum the effects of all higher derivative terms in the effective action for this gauge theory, does the result give non-singular physics as $`r0`$? If we sum the effects of all higher derivative terms beyond supergravity on the spacetime solution for a $`5`$-brane in M theory, do we find a smooth convergent solution near the $`5`$-brane?
These questions are intrinsically tied to the problem of gauge-fixing the Taylor series solution in a way that results in a normalizable solution. We want to construct the $`t_i^0`$ in a useful systematic expansion. The dominant terms in an asymptotic expansion are those that decay polynomially in $`1/r`$. We point out that for large $`r`$, an approximate asymptotic bound state can be constructed in a $`1/r`$ expansion using a method described in . An analytic expansion in $`1/r`$ near infinity is essentially a perturbative expansion, although not in $`g^2`$ a priori but in $`g^{2/3}`$. A similar technique was used in \[15,,16,,17\] to further explore the long distance dynamics and the asymptotic structure of the bound state wavefunction for $`2`$ D0-branes. Unfortunately, the effective long distance Hamiltonian has only been constructed to order $`1/r^2`$, which is the required order for an index computation \[14,,18\]. With our knowledge of the Taylor series solution (4.21) for this problem, we can do significantly better than those approximate constructions.
Can there be non-perturbative terms? These are terms which are not visible in a $`1/r`$ expansion, like $`e^{r^2}`$, but which become important as $`r`$ becomes small. We do not actually know whether there are any such terms, and there are no candidates like instanton configurations that could generate these terms in the abelian gauge theory. This leads the first author to suspect that the analytic expansion in $`1/r`$ might well be exact. Nevertheless, we cannot prove that non-perturbative terms are not present. To really rule out such terms, we need a technique for finding the bound state solution which is inherently more global than an asymptotic expansion. We shall discuss a more global approach in section $`6`$.
What is hard to see in the Taylor expansion of section $`4`$ is normalizability in the $`y`$-direction. This is much easier to see in the oscillator expansion so we will match the two expansions to get control over the question of normalizability. When expanded in a harmonic oscillator basis with frequency proportional to $`r`$,
$$f_i=\underset{n=0}{\overset{\mathrm{}}{}}f_i^n(r)|n>,$$
$`(5.1)`$
the $`f_i^n`$ must decay as $`r\mathrm{}`$. Let us expand each $`f_i^n`$ in powers of $`1/r`$. This implies that the $`t_i^n`$ have an expansion in $`1/r`$. We can then reorganize the Taylor series for $`f_i`$ in the following way:
$$\begin{array}{cc}\hfill f_i& =t_i^0+t_i^2y^2+t_i^4y^4+\mathrm{},\hfill \\ & =\underset{p}{}\frac{1}{r^p}\underset{k}{}b_k^{p,i}|k>.\hfill \end{array}$$
$`(5.2)`$
The $`b_k^{p,i}`$ are just some numbers which determine the collection of harmonics contributing to a given power $`1/r^p`$. Note that the oscillator eigenstates depend only on the combination $`ry^2`$ so that $`_kb_k^{p,i}|k>`$ is a power series in $`ry^2`$.
For example, suppose that the harmonic $`|m>`$ is the only harmonic with a non-zero coefficient in the sum $`_kb_k^{p,i}|k>`$ for the $`1/r^p`$ term of $`f_i`$. We stress that in general there can be many harmonics contributing to a given term, but for simplicity, let us assume there is just one. The Taylor series for $`f_i`$ must then contain the terms:
$$\begin{array}{cc}\hfill f_i=& b_m^{p,i}\frac{1}{r^p}\left(1+a_1^{(m)}y^2+\mathrm{}+a_m^{(m)}y^{2m}\right)e^{ry^2/2}+\mathrm{},\hfill \\ \hfill =& b_m^{p,i}\frac{1}{r^p}\left(1+\left\{a_1^{(m)}\frac{1}{2}r\right\}y^2+\left\{a_2^{(m)}\frac{1}{2}ra_1^{(m)}+\frac{1}{2^22!}r^2\right\}y^4+\mathrm{}\right)+\mathrm{},\hfill \\ \hfill =& b_m^{p,i}\frac{1}{r^p}\left(1\frac{1}{2}(1+m)ry^2+\frac{1}{24}(3+4m+2m^2)r^2y^4+\mathrm{}\right)+\mathrm{}.\hfill \end{array}$$
$`(5.3)`$
There are specific relations between the Taylor coefficients in (5.3). We want to impose relations of this kind on the $`t_i`$ to satisfy normalizability in the $`y`$-direction for each choice of $`p`$ in the $`1/r`$ expansion.
5.2. The structure of the solution
Our Taylor series solution is completely determined by the three functions $`t_1^0,t_3^0`$ and $`t_4^0`$. These three functions must obey equations $`(6)`$ and $`(7)`$ of (4.22). To ensure that $`t_4^0`$ is decaying, we see from $`(6)`$ that $`t_3^0`$ must take the form,
$$t_3^0=\frac{c_1}{r^3}+\mathrm{},$$
$`(5.4)`$
where omitted terms decay more rapidly. For a given choice of $`p`$ in the $`1/r`$ expansion, we wish to extract the terms in each $`t_3^n`$ which contribute to $`_kb_k^{p,3}|k>`$. Again, this is just the statement that we can organize the Taylor series for $`f_3`$ so that,
$$f_3=\frac{1}{r^3}\underset{k}{}b_k^{3,3}|k>+\mathrm{},$$
$`(5.5)`$
for some coefficients $`b_k^{3,3}`$. This is formally true for any Taylor series. What is generally not true is that generic $`b_k^{3,3}`$ give a wavefunction normalizable in the $`y`$ direction. Heuristically, the norm of $`f_3`$ should be dominated by the leading term in the $`1/r`$ expansion. If we compute the norm of $`f_3`$ under this assumption, we see that:
$$|f_3|^2r^4𝑑r\left\{\left(\frac{c_1}{r^3}\right)^2\frac{1}{r^2}\underset{k}{}\frac{|b_k^{3,3}|^2}{(2+2k)}\right\},$$
where we have taken the normalization of $`|k>`$ given in Appendix D.
By looking at (4.19), (4.20), (4.21) and (4.22), we can see what terms in the Taylor series for $`f_3`$ have the right structure to determine the $`b_k^{3,3}`$ coefficients of (5.5). To quickly answer this question, let us list some order of magnitude relations which follow from (4.22) for the perturbative expansion. All relations are given in terms of our $`3`$ independent functions $`t_1^0,t_3^0`$ and $`t_4^0`$,
$$\begin{array}{cc}\hfill (1)& t_5^0O(t_3^0/r^2+t_4^0/r^4),\hfill \\ \hfill (2)& t_{11}^0O(t_4^0/r^2),\hfill \\ \hfill (3)& t_7^0O(t_4^0),\hfill \\ \hfill (4)& t_9^0O(t_3^0/r^2),\hfill \\ \hfill (5)& t_{10}^0O(t_1^0+t_3^0/r^2).\hfill \end{array}$$
$`(5.6)`$
There are similar order of magnitude relations for $`t_i^2`$ from (4.19),
$$\begin{array}{cc}\hfill (1)& t_7^2O(t_3^0),\hfill \\ \hfill (2)& t_9^2O(t_1^0/r^2+t_4^0),\hfill \\ \hfill (3)& t_{11}^2O(t_3^0/r^2+t_1^0),\hfill \\ \hfill (4)& t_5^2O(t_4^0+t_3^0/r^4),\hfill \\ \hfill (5)& t_3^2O(r^2t_4^0+t_3^0/r^2),\hfill \\ \hfill (6)& t_4^2O(t_3^0+t_4^0/r^2),\hfill \\ \hfill (7)& t_{10}^2O(t_4^0+t_1^0/r^2),\hfill \\ \hfill (8)& t_1^2O(t_4^0+t_1^0/r^2).\hfill \end{array}$$
$`(5.7)`$
It is easy to continue and list order of magnitude relations for $`t_i^4`$ from (4.20), and higher $`t_i^n`$ using (4.21). These relations are useful for easily determining which terms in $`t_i^n`$ are relevant for determining the oscillator coefficients $`b_k^{p,i}`$.
Returning to our specific case of $`f_3`$, we see that since $`t_3^01/r^3`$, the only terms in $`t_3^2`$ relevant for computing $`b_k^{3,3}`$ are those proportional to $`1/r^2`$. Looking at $`t_3^2`$ from (5.7), we see that the only way to have a term contributing to the $`b_k^{3,3}|k>`$ is if,
$$t_4^0=\frac{c_1^{}}{r^4}+\mathrm{}.$$
$`(5.8)`$
It could be the case that $`c_1^{}=0`$, which means that the $`b_k^{3,3}`$ sum up in such a way that the $`y^2/r^2`$ term in (5.5) vanishes. Actually if $`c_1^{}=0`$, the situation is much worse: we can see after some work that the $`b_k^{3,3}`$ have to be chosen so that the $`(ry^2)^{1+2n}`$ terms in the Taylor expansion of,
$$\underset{k}{}b_k^{3,3}|k>,$$
vanish for all $`n`$. The $`(ry^2)^{2n}`$ terms are all proportional to $`c_1`$ and give the constraint,
$$c_1\left(1+\frac{1}{8}(ry^2)^2+\frac{1}{348}(ry^2)^4+\frac{1}{46080}(ry^2)^6+\mathrm{}\right)=\underset{k}{}b_k^{3,3}|k>.$$
$`(5.9)`$
However, the left hand side of (5.9) is the expansion of the lowest oscillator $`|0>`$ with the $`(ry^2)^{1+2n}`$ terms missing. Since $`|0>=e^{ry^2/2}`$, we immediately see that the left hand side of (5.9) sums to the expression:
$$\frac{1}{2}\left\{e^{ry^2/2}+e^{ry^2/2}\right\}.$$
This is clearly not normalizable so $`c_1^{}`$ is necessarily non-zero. We will use this kind of argument repeatedly to determine the asymptotic solution.
With $`c_1^{}0`$, we get a modified constraint:
$$c_1\left(1+\frac{1}{8}(ry^2)^2+\mathrm{}\right)c_1^{}\left(\frac{1}{10}(ry^2)+\frac{1}{240}(ry^2)^3+\mathrm{}\right)=\underset{k}{}b_k^{3,3}|k>.$$
$`(5.10)`$
It is easy to see that the terms proportional to $`c_1^{}`$ sum up to give,<sup>3</sup> To see this, let us introduce new variables $`(r,u=ry^2/2)`$. In terms of these variables and simple redefinitions of the $`f_i`$, the equations of (4.7) can be put in a triangular form with respect to the grading induced by the $`1/r`$ expansion. In doing this, we treat $`u`$ as independent of $`r`$. The resulting triangular system is exactly solvable in the $`1/r`$ expansion.
$$\frac{1}{10}\left\{e^{ry^2/2}e^{ry^2/2}\right\}.$$
In this case, it is completely clear that we only get a normalizable solution if we pick,
$$c_1^{}=5c_1.$$
$`(5.11)`$
More generally, we will encounter the situation where we have a Taylor series of the form,
$$\alpha \underset{n}{}d_n(ry^2)^{2n}+\alpha ^{}\underset{n}{}d_n^{}(ry^2)^{2n+1},$$
$`(5.12)`$
where we know both $`d_n`$ and $`d_n^{}`$, and either $`\alpha `$ or $`\alpha ^{}`$. Let us assume we know $`\alpha `$ and we wish to determine the values of $`\alpha ^{}`$ for which the Taylor series is normalizable. If the Taylor series can be fitted by a finite number of oscillators for some $`\alpha ^{}`$ then that choice of $`\alpha ^{}`$ is unique. To see this, suppose there were two distinct choices of $`\alpha ^{}`$. We could take the difference between the two series to obtain a Taylor series proportional to,
$$\underset{n}{}d_n^{}(ry^2)^{2n+1}.$$
However, this series cannot be generated by any sum over a finite number of oscillators since all the $`(ry^2)^{2n}`$ terms must vanish. Therefore the choice of $`\alpha ^{}`$ is unique.
From equation $`(6)`$ of (4.22), we see that the $`O(1/r^4)`$ term in $`t_4^0`$ induces an $`O(1/r^6)`$ correction to $`t_3^0`$. This begins to suggest a perturbation expansion in $`1/r^3`$, and indeed that seems to be the case. From $`t_4^2O(1/r^6)`$, we see that we need a $`O(1/r^7)`$ correction to $`t_4^0`$ and so on. It is easy to add $`t_1^0`$ into the story: from the expression for $`t_1^2O(1/r^4)`$, we learn that $`t_1^01/r^5`$ at leading order. From $`(7)`$ of (4.22), we see that $`t_1^0`$ mixes with the first correction to $`t_4^0`$ and so on. What we have uncovered is the minimal required form for the solution given $`c_10`$. The choice of $`c_1`$ then determines the normalization of the wavefunction. The solution is given by the expansion:
$$\begin{array}{cc}\hfill t_3^0& =\frac{c_1}{r^3}+\frac{c_2}{r^6}+\frac{c_3}{r^9}+\mathrm{},\hfill \\ \hfill t_4^0& =\frac{c_1^{}}{r^4}+\frac{c_2^{}}{r^7}+\frac{c_3^{}}{r^{10}}+\mathrm{},\hfill \\ \hfill t_1^0& =\frac{c_1^{\prime \prime }}{r^5}+\frac{c_2^{\prime \prime }}{r^8}+\frac{c_3^{\prime \prime }}{r^{11}}+\mathrm{}.\hfill \end{array}$$
$`(5.13)`$
It is interesting that the $`t_i^0`$ have an expansion in powers of $`1/r^3`$. Restoring the coupling constant, we see that this is an expansion in $`g^2/r^3`$. This is the natural expansion parameter in the gauge theory. From the perspective of the effective action on the Coulomb branch, what we are doing is summing the effects of the metric and all higher derivative corrections on the vacuum state of the gauge theory.
The expansion (5.13) represents the minimal required terms in the solution. It is natural to ask whether other powers of $`1/r`$ are possible. Equation $`(6)`$ of (4.22) requires that,
$$t_4^0O(r^2t_3^0),$$
except for the special case where $`t_3^0O(1/r^3)`$. Likewise, equation $`(7)`$ of (4.22) requires that,
$$t_1^0O(r^2t_4^0+\frac{1}{r^4}t_4^0),$$
again except for the case $`t_4^0O(1/r^4)`$. Let us be concrete: suppose $`t_3^0`$ has a term of order $`O(1/r^4)`$. This requires $`t_4^0O(1/r^2)`$ and therefore $`t_1^0O(1)`$, which is not possible. Suppose $`t_3^0`$ has a term of order $`O(1/r^5)`$, which implies that $`t_4^0O(1/r^3)`$. In turn, we see that $`t_3^2O(1/r)`$ which requires that $`t_3^0=\alpha /r^2`$. This is impossible unless $`\alpha =0`$. In that case, the Taylor expansion of $`f_3`$ has terms of the schematic form:
$$f_3O(\frac{y^2}{r})+O(ry^6)+\mathrm{}.$$
As in our earlier discussion, it is not hard to check that this sum is not normalizable. This kind of argument extends to higher powers: suppose $`t_3^0`$ has a term of order $`O(1/r^7)`$, which implies that $`t_4^0O(1/r^5)`$. From $`t_3^2O(1/r^3)`$, we require a $`O(1/r^4)`$ correction to $`t_3^0`$ which brings us back to an earlier case. This kind of reasoning suggests that only those terms generated by the original $`1/r^3`$ term in $`t_3^0`$ appear in the solution. This also agrees with our gauge theory intuition. We shall therefore restrict our attention to the terms of (5.13).
5.3. Supergravity and the leading terms of the solution
So far, we have found that the leading order term of $`f_3`$ is proportional purely to the oscillator ground state $`|0>`$. We also know all the leading powers for the $`t_i^0`$. A glance at Appendix C tells us that the dominant terms for large $`r`$ in the bound state $`\mathrm{\Psi }`$ are $`f_3,f_4`$ and $`f_5`$. Note that on the flat directions where $`y=0`$, the only non-vanishing forms in $`\mathrm{\Psi }`$ are those with coefficients $`f_3,f_4`$ and $`f_5`$. It is not hard to check that both $`f_3`$ and $`f_5`$ are also proportional to $`|0>`$ at leading order. Recalling that $`r=|x^1|`$, we have learnt that these leading terms are given by,
$$\begin{array}{cc}& f_3=\frac{c_1}{r^3}|0>+\mathrm{},\hfill \\ & x^1f_4=\frac{5c_1x^1}{r^4}|0>+\mathrm{},\hfill \\ & (x^1)^2f_5=\frac{5c_1(x^1)^2}{r^5}|0>+\mathrm{}.\hfill \end{array}$$
$`(5.14)`$
Although each term transforms very differently under $`Spin(5)`$ in either the $`\mathrm{𝟏},\mathrm{𝟓}`$ or $`\mathrm{𝟏𝟒}`$, each decays at the same rate at leading order.
We can compare what we have learned about the exact Taylor series solution with what we might expect from an effective Hamiltonian construction like the one given in . To construct the approximate bound state, we note that the wavefunction is sharply localized near the flat directions for large $`r`$. The potential term (2.20) is dominated by the $`r^2y^2`$ term. We conclude that for large $`r`$, the wavefunction can be expanded in a harmonic oscillator basis, with higher oscillator modes suppressed by powers of $`1/r`$. The approximate bound state wavefunction takes the form of a product: there is a wavefunction of the light vector multiplet degrees of freedom multiplied by the ground state for the massive hypermultiplet degrees of freedom. We can now construct the effective Hamiltonian governing the long distance physics in a $`1/r`$ expansion. The hypermultiplet ground state is easily determined from (2.16) to have the form,<sup>4</sup> The natural decomposition of $`x^\mu `$ and $`q_i`$ into light and heavy variables at large $`r`$ makes the construction of the approximate ground state much simpler than the non-abelian D0-D0 theory considered in .
$$re^{ry^2/2}dv_1dv_2.$$
$`(5.15)`$
In the most primitive approximation where the effective interactions are cancelled to order $`1/r`$, the effective Hamiltonian is $`(1/2)p^\mu p^\mu `$ and acts on (5.15) multiplied by some function of $`x`$ and $`du`$. In fact, to determine the leading decay, we should go to order $`1/r^2`$. Even then there will be a degenerate set of approximate ground states.
However, this rough approximation is good enough for the purpose of comparison with the structure following from the Taylor series solution. Using the forms given in Appendix C, we see that the coefficients in (5.14) precisely conspire to give agreement with the approximate construction. The leading terms in $`\mathrm{\Psi }`$ sum to give,
$$\begin{array}{cc}\hfill \mathrm{\Psi }& =f_3|2,2>_\mathrm{𝟏}+f_4x^1|2,2>_\mathrm{𝟓}^1+f_5(x^1)^2|2,2>_{\mathrm{𝟏𝟒}}^{11}+\mathrm{},\hfill \\ & \frac{10c_1}{r^3}e^{ry^2/2}(du_1du_2du_3du_4)(dv_1dv_2)+\mathrm{},\hfill \end{array}$$
$`(5.16)`$
where we take $`x^1`$ positive. There are a number of strange features of this asymptotic solution. First note that the vacuum for the massive fermions $`dv_1dv_2`$ is not a single representation of $`Spin(5)`$. Let us contrast this with the case of $`2`$ D0-branes where the vacuum for the massive fermions transforms in a single representation, the $`\mathrm{𝟒𝟒}`$, of $`Spin(9)`$ \[15,,17\]. The appearance of three different representations, the $`\mathrm{𝟏},\mathrm{𝟓}`$ and $`\mathrm{𝟏𝟒}`$, in the asymptotic solution really cries out for an interpretation both in terms of the DLCQ M theory $`5`$-brane, and in terms of the supergravity solution for the D0-D4 bound state.
What this seems to suggest is that the massive degrees of freedom never really decouple at large $`r`$. Through the $`Spin(5)`$ flavor symmetry, the vacuum always knows about the existence of massive degrees of freedom. This issue is intimately tied to uniqueness of the bound state: the statement that the bound state is unique involves knowledge about both long and short distance physics. Uniqueness however forces invariance under the full $`Spin(5)`$ acting on both light and heavy degrees of freedom. That $`dv_1dv_2`$ is not an irreducible representation then requires particular combinations of $`Spin(5)`$ representations for the light degrees of freedom. A deeper understanding of this issue is certainly in order.
5.4. Beyond supergravity
Let us return to our general solution (5.13). We determined $`c_1^{}`$ in terms of $`c_1`$ in (5.11). How do we determine $`c_1^{\prime \prime }`$? The straightforward way to determine $`c_1^{\prime \prime }`$ is not particularly elegant. We can take the Taylor series for $`f_1`$ and impose normalizability in the $`y`$-direction,
$$\begin{array}{cc}\hfill f_1& =\left(\frac{c_1^{\prime \prime }}{r^5}\right)+\left(\frac{5c_1}{4r^4}\right)y^2+\left(\frac{c_1^{\prime \prime }10c_1}{40r^3}\right)y^4+\left(\frac{5c_1}{96r^2}\right)y^6\hfill \\ & +\left(\frac{c_1^{\prime \prime }}{4480r}\frac{c_1}{168r}\right)y^8+O(y^{10})+\mathrm{}.\hfill \end{array}$$
$`(5.17)`$
Note the general structure: half the terms in the Taylor series are known. The other half depend on the unknown constant that we wish to determine. It is not hard to see that the $`(ry^2)^{1+2n}`$ terms come from expanding the lowest oscillator $`|0>`$ state. The situation is then essentially the same as in (5.10), and we find that:
$$c_1^{\prime \prime }=\frac{5}{2}c_1.$$
$`(5.18)`$
From equation $`(6)`$ of (4.22), we find that $`c_2`$ is determined by $`c_1^{}`$:
$$c_2=c_1.$$
$`(5.19)`$
Equation $`(7)`$ of (4.22) relates $`c_2^{}`$ to $`c_1^{\prime \prime }`$,
$$c_2^{}=\frac{25}{2}c_1.$$
$`(5.20)`$
We are again back to the question of studying $`f_1`$ to determine $`c_2^{\prime \prime }`$. The relevant terms in $`f_1`$ begin at order $`1/r^8`$,
$$\begin{array}{cc}\hfill f_1=& \frac{1}{r^8}(c_2^{\prime \prime }+\frac{125c_1}{16}(ry^2)+\frac{c_2^{\prime \prime }25c_1}{40}(ry^2)^2+\frac{25c_1}{384}(ry^2)^3\hfill \\ & +\{\frac{c_2^{\prime \prime }}{4480}+\frac{5c_1}{384}\}(ry^2)^4+O(ry^2)^5)+\mathrm{}.\hfill \end{array}$$
$`(5.21)`$
Apart from more complicated coefficients, the general pattern is the same. Half the terms in the Taylor series are proportional to $`c_1`$ while the other half depend on the unknown constant $`c_2^{\prime \prime }`$. We can try to match (5.21) with an oscillator expansion involving a finite number of oscillators. Indeed, for the choice
$$c_2^{\prime \prime }=\frac{175}{8}c_1,$$
$`(5.22)`$
we can fit (5.21) by the first two oscillators:
$$f_1=\frac{225}{8}\frac{c_1}{r^8}\left(|0>\frac{2}{9}|1>\right)+\mathrm{}.$$
As we might expect, higher excited oscillator states appear in the solution as we study more rapidly decreasing terms. This also suggests that only a finite number of oscillators will appear at any given order in the $`1/r`$ expansion.
Let us iterate the argument one more time. From $`(6)`$ and $`(7)`$ of (4.22), we immediately obtain the following relations:
$$\begin{array}{cc}& c_3=\frac{35}{2}c_1,\hfill \\ & c_3^{}=\frac{195}{2}c_1.\hfill \end{array}$$
$`(5.23)`$
The relevant terms in $`f_1`$ take the form,
$$\begin{array}{cc}\hfill f_1=& \frac{1}{r^{11}}(c_3^{\prime \prime }+\frac{2825c_1}{16}(ry^2)+\{\frac{c_3^{\prime \prime }}{40}\frac{595c_1}{32}\}(ry^2)^2\hfill \\ & +\frac{425c_1}{128}(ry^2)^3+O(ry^2)^4)+\mathrm{},\hfill \end{array}$$
$`(5.24)`$
and can be fit by the first three oscillator modes,
$$f_1=\frac{525}{32}\frac{c_1}{r^{11}}\left(|2>\frac{208}{21}|1>+\frac{115}{3}|0>\right)+\mathrm{},$$
for the choice,
$$c_3^{\prime \prime }=\frac{7725}{16}c_1.$$
$`(5.25)`$
6. A Remarkable Reduction
6.1. Prolongation
Studying the bound state solution in an asymptotic series is analogous to studying the M5-brane in a derivative expansion around the supergravity solution. Ideally, we want a more powerful technique to solve the vacuum equations. The aim of this final section is to present a more global approach to solving the vacuum equations. Hopefully, this approach is closer to the method an M theorist might use to study the M5-brane.
We shall present a surprising reduction of the long list of equations (4.7) to a single scalar elliptic equation. The equation takes the form,
$$\left(\mathrm{\Delta }+\stackrel{}{B}+W\right)u=0,$$
$`(6.1)`$
where $`\mathrm{\Delta }=_r^2+_y^2`$, the vector field $`\stackrel{}{B}`$ with components $`(B_r,B_y)`$ and the potential $`W`$ are rational functions of $`r`$ and $`y`$. The function $`u=u(r,y)`$ is a particular combination of the $`f_i`$. In order to explain how we reduce our first order system to a single second order equation, let us first consider the inverse process: prolongation. To understand the procedure, let us begin with an illustrative example: we take an equation of the form,
$$F_{xx}+F_{yy}+UF=0.$$
$`(6.2)`$
We can define new functions $`p=F_x`$ and $`q=F_y`$ with which we can ‘prolong’ our scalar second order equation into a system of first order equations:
$$\begin{array}{cc}& p_x+q_y+UF=0,\hfill \\ & F_x=p,\hfill \\ & F_y=q.\hfill \end{array}$$
$`(6.3)`$
We also have a compatibility relation,
$$F_{xy}=F_{yx},$$
which implies an additional fourth equation:
$$p_yq_x=0.$$
$`(6.4)`$
We can express these relations as a differential system,
$$\begin{array}{cc}& dF=pdx+qdy,\hfill \\ & dp=adx+bdy,\hfill \\ & dq=bdx(a+UF)dy.\hfill \end{array}$$
$`(6.5)`$
We could repeat the procedure and prolong again by adding the equations,
$$\begin{array}{cc}& da=cdx+edy,\hfill \\ & db=edx+((UF)_xc)dy.\hfill \end{array}$$
$`(6.6)`$
Each time we prolong a system of equations like this, we add two new unknown functions. These functions are the unknown derivatives of the functions comprising the previous system of equations.
6.2. Deprolongation
Now we would like to ‘deprolong’ our system of equations (4.7).<sup>5</sup> We are especially grateful to Robert Bryant for suggesting and explaining this reduction to us. We shall see that the system (4.7) containing $`7`$ independent functions can be obtained by prolonging a single second order equation three times. In order to deprolong (4.7), we write the equations as an exterior system as before. We choose $`5`$ functions $`F_i`$ which are linear combinations of the initial seven $`f_i`$ so that $`dF_i`$ is expressible algebraically in terms of the original seven. We keep the $`5`$ equations defining the $`dF_i`$ and discard the remaining two equations. For example, for this first step, we can make the following choice:
$$\begin{array}{cc}\hfill F_1& =f_1+f_{10},\hfill \\ \hfill F_2& =r^2f_5y^2f_1,\hfill \end{array}$$
$$F_3=f_7,F_4=f_4,F_5=f_{11}.$$
So far the equations remain first order. We then iterate this procedure until we arrive at the first prolongation of the scalar equation. We then make the obvious substitution to transform the first prolongation into a second order scalar equation.
Let us summarize the results of the deprolongation. We obtain the following equation:
$$u_{rr}+u_{yy}+\frac{\stackrel{~}{B}_r}{rF}u_r+\frac{\stackrel{~}{B}_y}{yF}u_y+\frac{\stackrel{~}{W}}{F}u=0.$$
$`(6.7)`$
We define $`u`$ in terms of the variables $`s=y^2/2`$ and $`t=r^2/2`$ and the functions $`f_1,f_4,f_5,f_7,f_{10},f_{11}`$:
$$\begin{array}{cc}\hfill u=& \left(2st\right)f_7\left(\frac{4}{3}ts^2\frac{1}{3}t^2s\frac{1}{6}t^31\right)f_{10}+\left(\frac{4}{3}ts^2\frac{1}{3}t^2s\frac{1}{6}t^3+1\right)f_1\hfill \\ & sf_4+\left(\frac{2}{3}ts^2\frac{8}{3}s^3+\frac{1}{3}t^2s\right)f_5\left(\frac{2}{3}ts\frac{8}{3}s^2+\frac{5}{6}t^2\right)f_{11}.\hfill \end{array}$$
$`(6.8)`$
It would be interesting to find a geometric interpretation for (6.7) – perhaps in terms of some bundle over either $`\mathrm{IR}^2`$ or $`\mathrm{IR}^9`$. Let us list the rational functions which appear in (6.7). For $`\stackrel{~}{B}_r`$, we find the expression:
$$\begin{array}{cc}\hfill \stackrel{~}{B}_r=& 224t^3s^3+576t^2s^4+32t^4s^212t^5s+4t^612t^2s\hfill \\ & +144ts^224t^3448s^5t48s^3+36,\hfill \end{array}$$
$`(6.9)`$
while $`\stackrel{~}{B}_y`$ and $`F`$ are given by,
$$\begin{array}{cc}\hfill \stackrel{~}{B}_y=& 272t^2s^4192t^3s^340t^4s^2+264ts^2216t^2s+68t^5s+42t^315t^6+964s^5t,\hfill \\ \hfill F=& 9+64t^3s^38t^4s^26t^3+24t^2s112t^2s^4+t^624ts^2+64s^5t4t^5s.\hfill \end{array}$$
$`(6.10)`$
Lastly, the function $`\stackrel{~}{W}`$ which determines the potential $`W`$ takes the form:
$$\begin{array}{cc}\hfill \stackrel{~}{W}=& 105t^2+80s^3t^2+38t^5+144st96t^3s^2+24s^2144t^4s^432t^5s^3\hfill \\ & +24t^6s^2256s^6t^2+224s^4t+384s^5t^3t^876st^4.\hfill \end{array}$$
$`(6.11)`$
Equation (6.7) is remarkably simple by comparision with (4.7). That the equations reduce this way opens up the possibility of answering a host of otherwise intractable questions.
Acknowledgements
It is our pleasure to thank Robert Bryant for suggesting and explaining deprolongation which led to the results of section six. The work of S.S. is supported by the William Keck Foundation and by NSF grant PHY–9513835; that of M.S. by NSF grant DMS–9870161.
Appendix A. Quaternions and Symplectic Groups
We will summarize some useful relations between quaternions and symplectic groups. Let us label a basis for our quaternions by $`\{\mathrm{𝟏},I,J,K\}`$ where,
$$I^2=J^2=K^2=\mathrm{𝟏},IJK=\mathrm{𝟏}.$$
A quaternion $`q`$ can then be expanded in components
$$q=q^1+Iq^2+Jq^3+Kq^4.$$
The conjugate quaternion $`\overline{q}`$ has an expansion
$$q=q^1Iq^2Jq^3Kq^4.$$
The symmetry group $`Sp(1)_RSU(2)_R`$ is the group of unit quaternions. Let $`\mathrm{\Lambda }`$ be a field transforming in the $`\mathrm{𝟐}`$ of $`Sp(1)_R`$. If we view $`Sp(1)_R`$ acting on $`\mathrm{\Lambda }`$ as right multiplication by a unit quaternion $`g`$ then,
$$\mathrm{\Lambda }\mathrm{\Lambda }g.$$
In this formalism, $`\mathrm{\Lambda }`$ is valued in the quaternions. Equivalently, we can expand $`\mathrm{\Lambda }`$ in components and express the action of $`g`$ in the following way,
$$\mathrm{\Lambda }_ag_{ab}\mathrm{\Lambda }_b,$$
where $`g_{ab}`$ implements right multiplication by the unit quaternion $`g`$. For example, right multiplication by $`I`$ on $`q`$ gives
$$\begin{array}{cc}\hfill q& qI\hfill \\ & q^1Iq^2q^3K+q^4J,\hfill \end{array}$$
which can be realized by the matrix
$$I^R=\left(\begin{array}{cccc}0& 1& 0& 0\\ 1& 0& 0& 0\\ 0& 0& 0& 1\\ 0& 0& 1& 0\end{array}\right)$$
$`(\text{A.}1)`$
acting on $`q`$ in the usual way $`q_aI_{ab}^Rq_b`$. The matrices $`J^R`$ and $`K^R`$ realize right multiplication by $`J,K`$ while $`\mathrm{𝟏}^R`$ is the identity matrix:
$$J^R=\left(\begin{array}{cccc}0& 0& 1& 0\\ 0& 0& 0& 1\\ 1& 0& 0& 0\\ 0& 1& 0& 0\end{array}\right),K^R=\left(\begin{array}{cccc}0& 0& 0& 1\\ 0& 0& 1& 0\\ 0& 1& 0& 0\\ 1& 0& 0& 0\end{array}\right).$$
$`(\text{A.}2)`$
We define operators $`s^j`$ in terms of $`\{\mathrm{𝟏}^R,I^R,J^R,K^R\}`$
$$s^1=\left(\begin{array}{cc}\mathrm{𝟏}^R& 0\\ 0& \mathrm{𝟏}^R\end{array}\right),s^2=\left(\begin{array}{cc}I^R& 0\\ 0& I^R\end{array}\right),s^3=\left(\begin{array}{cc}J^R& 0\\ 0& J^R\end{array}\right),s^4=\left(\begin{array}{cc}K^R& 0\\ 0& K^R\end{array}\right).$$
In a similar way, the group $`Sp(2)Spin(5)`$ is the group of quaternion-valued $`2\times 2`$ matrices with unit determinant. We will view $`Sp(2)`$ as acting by left multiplication on a field $`\mathrm{\Psi }`$ in the defining representation. So an element $`USp(2)`$ acts in the following way:
$$\mathrm{\Psi }U\mathrm{\Psi }.$$
Equivalently, in terms of components
$$\mathrm{\Psi }_aU_{ab}\mathrm{\Psi }_b.$$
Lastly, we can give an explicit form for the gamma matrices (2.3) in terms of quaternions:
$$\gamma ^1=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right),\gamma ^2=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right),\gamma ^3=\left(\begin{array}{cc}0& I\\ I& 0\end{array}\right)$$
$$\gamma ^4=\left(\begin{array}{cc}0& J\\ J& 0\end{array}\right),\gamma ^5=\left(\begin{array}{cc}0& K\\ K& 0\end{array}\right).$$
In turn, $`\{I,J,K\}`$ can be expressed in terms of the Pauli matrices $`\sigma ^i`$
$$\sigma ^1=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right),\sigma ^2=\left(\begin{array}{cc}0& i\\ i& 0\end{array}\right),\sigma ^3=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)$$
as $`4\times 4`$ real anti-symmetric matrices:
$$I=\left(\begin{array}{cc}0& \sigma ^1\\ \sigma ^1& 0\end{array}\right),J=\left(\begin{array}{cc}i\sigma ^2& 0\\ 0& i\sigma ^2\end{array}\right),K=\left(\begin{array}{cc}0& \sigma ^3\\ \sigma ^3& 0\end{array}\right).$$
Appendix B. Forms and Representations of $`Sp(2)`$
Using the complexification (3.7), we obtain a set of Hermitian $`4\times 4`$ matrices from those given in Appendix A:
$$\stackrel{~}{\gamma }^1=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right),\stackrel{~}{\gamma }^2=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right),\stackrel{~}{\gamma }^3=\left(\begin{array}{cc}0& i\sigma ^1\\ i\sigma ^1& 0\end{array}\right)$$
$$\stackrel{~}{\gamma }^4=\left(\begin{array}{cc}0& i\sigma ^3\\ i\sigma ^3& 0\end{array}\right),\stackrel{~}{\gamma }^5=\left(\begin{array}{cc}0& i\sigma ^2\\ i\sigma ^2& 0\end{array}\right).$$
We also need the symplectic metric or charge conjugation matrix,
$$C=\left(\begin{array}{cccc}0& 1& 0& 0\\ 1& 0& 0& 0\\ 0& 0& 0& 1\\ 0& 0& 1& 0\end{array}\right),$$
which implements complex conjugation:
$$C\stackrel{~}{\gamma }C=\stackrel{~}{\gamma }^{}.$$
$`(\text{B.}1)`$
Our forms $`du_a`$ and $`dv_a`$ transform in the $`\mathrm{𝟒}`$ of $`Sp(2)`$. The representation $`^2\mathbf{\hspace{0.17em}4}`$ decomposes into $`\mathrm{𝟓}\mathrm{𝟏}`$. As an example, we can explicitly construct the $`\mathrm{𝟏}`$ from $`(1,1)`$ forms in the following way,
$$duCdv,$$
$`(\text{B.}2)`$
while the $`\mathrm{𝟓}`$ is given by:
$$du\stackrel{~}{\gamma }^\mu Cdv.$$
$`(\text{B.}3)`$
It is not hard to check that these combinations transform correctly. Lastly, we need to consider $`\mathrm{𝟓}\mathrm{𝟓}=\mathrm{𝟏}\mathrm{𝟏𝟎}\mathrm{𝟏𝟒}`$ since the $`\mathrm{𝟏}`$ and the $`\mathrm{𝟏𝟒}`$ appear in (3.14). The $`\mathrm{𝟏}`$ is given by the form,
$$\underset{\mu }{}du\stackrel{~}{\gamma }^\mu Cdudv\stackrel{~}{\gamma }^\mu Cdv,$$
$`(\text{B.}4)`$
while the $`\mathrm{𝟏𝟒}`$ has components:
$$du\stackrel{~}{\gamma }^{(\mu }Cdudv\stackrel{~}{\gamma }^{\nu )}Cdv\frac{1}{5}\delta ^{\mu \nu }\underset{\rho }{}du\stackrel{~}{\gamma }^\rho Cdudv\stackrel{~}{\gamma }^\rho Cdv.$$
$`(\text{B.}5)`$
Appendix C. A Word on Normalizations
To fix the choice of normalizations, we list explicitly the forms appearing in (3.12) through (3.18) at the special point where $`x^10`$ with $`x^\mu =0`$ for $`\mu >1`$ and $`q_10`$ with $`q_j=0`$ for $`j>1`$. All forms act on the canonical vacuum $`|0>`$ which is omitted,
$$f_1|0,0>=f_1(q_1)^2du_1du_2du_3du_4,$$
$`(\text{C.}1)`$
$$\begin{array}{cc}\hfill f_2|0,4>=& f_2(q_1)^2dv_1dv_2dv_3dv_4,\hfill \end{array}$$
$`(\text{C.}2)`$
$$\begin{array}{cc}\hfill f_3|2,2>_\mathrm{𝟏}=& f_3\{(du_1du_2du_3du_4)(dv_1dv_2dv_3dv_4)\hfill \\ & 2(du_1du_4dv_2dv_3+du_2du_3dv_1dv_4)+\hfill \\ & 2(du_1du_3dv_2dv_4+du_2du_4dv_1dv_3)\},\hfill \\ \hfill f_4x^1|2,2>_\mathrm{𝟓}^1=& f_4x^1(du_1du_2du_3du_4)(dv_1dv_2+dv_3dv_4),\hfill \\ \hfill f_5(x^1)^2|2,2>_{\mathrm{𝟏𝟒}}^{11}=& f_5(x^1)^2\{\frac{4}{5}(du_1du_2du_3du_4)(dv_1dv_2dv_3dv_4)+\hfill \\ & \frac{2}{5}(du_1du_4dv_2dv_3+du_2du_3dv_1dv_4)\hfill \\ & \frac{2}{5}(du_1du_3dv_2dv_4+du_2du_4dv_1dv_3)\},\hfill \end{array}$$
$`(\text{C.}3)`$
$$\begin{array}{cc}\hfill f_6|1,3>=& f_6q_1(dv_1dv_2+dv_3dv_4)(du_1dv_2du_2dv_1+du_3dv_4du_4dv_3),\hfill \\ \hfill f_7|1,1>=& f_7q_1(du_1du_2+du_3du_4)(du_1dv_2du_2dv_1+du_3dv_4du_4dv_3),\hfill \end{array}$$
$`(\text{C.}4)`$
$$\begin{array}{cc}\hfill f_8x^1|1,3>^1=& f_8x^1q_1(dv_1dv_2+dv_3dv_4)(du_1dv_2du_2dv_1du_3dv_4+du_4dv_3),\hfill \\ \hfill f_9x^1|1,1>^1=& f_9x^1q_1(du_1du_2+du_3du_4)(du_1dv_2du_2dv_1du_3dv_4+du_4dv_3),\hfill \end{array}$$
$`(\text{C.}5)`$
$$\begin{array}{cc}\hfill f_{10}|0,2>=& f_{10}\frac{1}{2}(q_1)^2(du_1du_2+du_3du_4)(dv_1dv_2+dv_3dv_4),\hfill \end{array}$$
$`(\text{C.}6)`$
$$\begin{array}{cc}\hfill f_{11}x^1|0,2>^1=& f_{11}x^1\frac{1}{2}(q_1)^2(du_1du_2+du_3du_4)(dv_1dv_2dv_3dv_4).\hfill \end{array}$$
$`(\text{C.}7)`$
Appendix D. The Four-Dimensional Radial Harmonic Oscillator
We want to construct eigenstates for the radial four-dimensional simple harmonic oscillator which satisfy,
$$\left\{_y^2\frac{3}{y}_y+r^2y^2\right\}|n>=E_n|n>,$$
$`(\text{D.}1)`$
where $`E_n=4(n+1)r`$. The easiest is the ground state:
$$|0>=e^{ry^2/2}.$$
A general eigenstate takes the form,
$$|n>=\left(1+a_1^{(n)}y^2+\mathrm{}+a_n^{(n)}y^{2n}\right)e^{ry^2/2}.$$
$`(\text{D.}2)`$
It is not hard to check using the nice relation,
$$\left\{_y^2\frac{3}{y}_y+r^2y^2\right\}y^{2n}e^{ry^2/2}=\left\{E_ny^{2n}4n(n+1)y^{2n2}\right\}e^{ry^2/2},$$
$`(\text{D.}3)`$
that each coefficient $`a_m^{(n)}`$ is determined by the recursion relation,
$$a_m^{(n)}=\frac{a_{m1}^{(n)}(mn1)r}{m(m+1)},$$
$`(\text{D.}4)`$
where $`a_0^{(n)}=1`$. Note that these eigenstates are not normalized, but they are orthogonal when integrated with the measure $`y^3dy`$. The norm of these eigenstates is given by the formula,
$$<n|n>=\frac{1}{r^2}\frac{1}{(2+2n)}.$$
$`(\text{D.}5)`$
We will need to evaluate various operators acting on $`|n>`$. The nicest are the three operators $`y^2,y_y,_r`$. These operators raise and lower by at most one unit:
$$\begin{array}{cc}\hfill ry^2|n>& =n|n1>+2(n+1)|n>(n+2)|n+1>,\hfill \\ \hfill y_y|n>& =n|n1>2|n>+(n+2)|n+1>,\hfill \\ \hfill 2r_r|n>& =n|n1>2|n>+(n+2)|n+1>.\hfill \end{array}$$
$`(\text{D.}6)`$
Note that $`y_y`$ is equivalent to $`2r_r`$ when acting on $`|n>`$.
Appendix E. Equations from a Taylor Expansion
We list explicitly the equations that follow from a Taylor expansion of the $`f_i`$ in the $`y`$-direction and which only involve $`t_i^0`$ and $`t_i^2`$,
$$\begin{array}{cc}\hfill (1)& \frac{1}{2}t_1^0\frac{r^2}{5}t_5^2+t_{11}^2+t_3^2=0,\hfill \\ \hfill (2)& r_rt_9^26t_1^2\frac{1}{2}t_7^0+5t_9^2+2t_{10}^2=0,\hfill \\ \hfill (3)& _rt_7^2+2rt_{11}^2+r\left\{t_1^0\frac{1}{2}t_9^0\right\}=0,\hfill \\ \hfill (4)& 2rt_9^2+_rt_1^0+rt_7^0=0,\hfill \\ \hfill (5)& 2t_7^2+r^2t_9^0=0,\hfill \\ \hfill (6)& 4t_{10}^24t_1^2+5t_9^2+r^2f_{11}^0+r_rt_9^2+\frac{1}{2}t_7^0=0,\hfill \\ \hfill (7)& 2r^2t_4^0+3t_9^0r_rt_9^0+4t_3^2+\frac{16}{5}r^2t_5^2=0,\hfill \\ \hfill (8)& \frac{4}{5}r^2t_5^24t_3^2r_rt_9^03t_9^0=0,\hfill \\ \hfill (9)& t_{10}^0+\frac{1}{2}t_9^0+4t_{11}^2+\frac{1}{r}_rt_7^2=0,\hfill \\ \hfill (10)& 2t_3^0+\frac{8}{5}r^2t_5^0+4t_4^2\frac{1}{r}_rt_7^0=0,\hfill \\ \hfill (11)& r\left\{t_4^0+t_7^0\right\}2rt_9^2+_rt_{10}^0=0,\hfill \\ \hfill (12)& \frac{56}{5}rt_5^2rt_7^0+6rt_9^2+2_rt_3^2+\frac{8}{5}r^2_rt_5^2=0,\hfill \\ \hfill (13)& \frac{2}{r}_rt_3^2\frac{2}{5}r_rt_5^2\frac{14}{5}t_5^2t_7^06t_9^2=0,\hfill \\ \hfill (14)& 10t_4^2+6t_7^2r^2t_9^0+2r_rt_4^2=0,\hfill \\ \hfill (15)& t_3^0+\frac{4}{5}r^2t_5^0+r^2t_9^0+t_{11}^02t_7^2+r_rt_{11}^0=0.\hfill \end{array}$$
$`(\text{E.}1)`$
Those equations that only involve $`t_i^0,t_i^2`$ and $`t_i^4`$ are given below:
$$\begin{array}{cc}\hfill (1)& t_3^4+\frac{1}{2}t_1^2\frac{r^2}{5}t_5^4+t_{11}^4=0,\hfill \\ \hfill (2)& r_rt_9^48t_1^4\frac{1}{2}t_7^2+5t_9^4+2t_{10}^4=0,\hfill \\ \hfill (3)& _rt_7^4+2rt_{11}^4+r\left\{t_1^2\frac{1}{2}t_9^2\right\}=0,\hfill \\ \hfill (4)& 4rt_9^4+_rt_1^2+rt_7^2=0,\hfill \\ \hfill (5)& 4t_7^4+r^2t_9^2+\frac{1}{2}t_1^0=0,\hfill \\ \hfill (6)& 6t_{10}^44t_1^4+5t_9^4+r^2f_{11}^2+r_rt_9^4+\frac{1}{2}t_7^2=0,\hfill \\ \hfill (7)& 2r^2t_4^2\frac{1}{2}t_7^0+3t_9^2r_rt_9^2+8t_3^4+\frac{32}{5}r^2t_5^4=0,\hfill \\ \hfill (8)& \frac{8}{5}r^2t_5^48t_3^4r_rt_9^23t_9^2+\frac{1}{2}t_7^0=0,\hfill \\ \hfill (9)& t_{10}^2+\frac{1}{2}t_9^2+6t_{11}^4+\frac{1}{r}_rt_7^4=0,\hfill \\ \hfill (10)& 2t_3^2+\frac{8}{5}r^2t_5^2\frac{1}{2}t_9^0+8t_4^4\frac{1}{r}_rt_7^2=0,\hfill \\ \hfill (11)& r\left\{t_4^2+t_7^2\right\}4rt_9^4+_rt_{10}^2=0,\hfill \\ \hfill (12)& \frac{56}{5}rt_5^4rt_7^2+8rt_9^4+\frac{r}{2}t_{11}^0+2_rt_3^4+\frac{8}{5}r^2_rt_5^4=0,\hfill \\ \hfill (13)& \frac{2}{r}_rt_3^4\frac{2}{5}r_rt_5^4\frac{14}{5}t_5^4t_7^28t_9^4=0,\hfill \\ \hfill (14)& 10t_4^4+8t_7^4r^2t_9^2+\frac{1}{2}t_{10}^0+2r_rt_4^4=0,\hfill \\ \hfill (15)& t_3^2+\frac{4}{5}r^2t_5^2+r^2t_9^2+t_{11}^24t_7^4+r_rt_{11}^2=0.\hfill \end{array}$$
$`(\text{E.}2)`$
References
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warning/0002/astro-ph0002189.html | ar5iv | text | # Spatially-resolved spectra of the accretion disc of the novalike UU Aquarii
## 1 Introduction
The standard picture of a novalike system is that of a close binary in which a late type star fills its Roche lobe and transfers matter to a companion white dwarf via an accretion disc. A bright spot is expected to form where the gas stream from the donor star hits the edge of the accretion disc.
The SW Sex stars (Thorstensen et al. 1991) form a sub-class of the novalikes with orbital periods in the range 3-4 hs that do not seem to fit within the above standard picture, displaying a range of peculiarities: (1) single peaked asymmetric emission lines showing little eclipse, (2) large ($`70\mathrm{°}`$) phase shifts between photometric and spectroscopic conjunction, (3) orbital phase-dependent absorption in the Balmer lines, (4) Doppler tomograms bright in the lower-left quadrant with small or no sign of disc emission, and (5) v-shaped continuum eclipses implying in flat radial temperature profiles in the inner disc (e.g., Warner 1995 and references therein). Earlier proposals to explain the phenomenon include accretion disc winds (Honeycutt, Schlegel & Kaitchuck 1986), magnetic white dwarfs disrupting the inner disc (Williams 1989), and gas stream overflow (Hellier & Robinson 1994). The two most recent models proposed to explain the phenomenology of the SW Sex stars are the disc-anchored magnetic propeller (Horne 1999) and a combination of stream overflow \+ disc winds (Hellier 1999).
UU Aqr is an eclipsing novalike (P$`{}_{\mathrm{orb}}{}^{}=3.9`$ hr) whose spectrum is dominated by single-peaked strong Balmer and He I emission lines (e.g., Downes & Keyes 1988). H$`\alpha `$ spectroscopy revealed that the line profile is highly asymmetric and phase-dependent and that the spectroscopic conjunction lags mid-eclipse by $`0.15`$ cycle (Haefner 1989; Diaz & Steiner 1991). The lack of the rotational disturbance typical of emitting accretion discs during eclipse in H$`\beta `$ led Hessman (1990) to the suggestion that the emission lines have a non-disc origin.
Baptista, Steiner & Cieslinski (1994; hereafter BSC94) derived a photometric model for the binary with $`q=0.30`$, M$`{}_{1}{}^{}=0.67\mathrm{M}_{}`$, an inclination of $`i=78`$ degrees. From the analysis of mid-eclipse fluxes they suggested that the Balmer lines are formed in an extended region only partially occulted during eclipse, possibly in a wind emanating from the inner disc. They also found that UU Aqr presents long-term brightness variations of low amplitude ($`0.3`$ mags) on timescales of years.
The eclipse mapping study of Baptista, Steiner & Horne (1996; thereafter BSH96) indicates that the inner disc of UU Aqr is optically thick, resulting in a distance estimate of 200 pc. Temperatures in the disc range from $`6000`$ K in the outer regions to $`16000`$ K near the white dwarf at disc centre. The radial temperature profiles in the high state follow the T$`R^{3/4}`$ law in the outer and intermediate disc regions but flattens off in the inner disc, leading to mass accretion rates of $`10^{9.2}\mathrm{M}_{}\mathrm{yr}^1`$ at $`R=0.1R_{\mathrm{L1}}`$ and $`10^{8.8}\mathrm{M}_{}\mathrm{yr}^1`$ at $`R=0.3R_{\mathrm{L1}}`$ ($`R_{\mathrm{L1}}`$ is the distance from disc centre to the inner Lagrangian point). Together with other characteristics, this led BSH96 to suggest that UU Aqr was an SW Sex star. The comparison of eclipse maps of the low and high states revealed that the differences are due to changes in the structure of the outer parts of the disc, the most noticeable effect being the appearance of a conspicuous red, bright structure at disc rim, which the authors identified with the bright spot.
According to BSH96, the Ṁ of UU Aqr is barely above the critical limit for disc instability to set in. Warner (1997) noted that the outer disc temperature is only 6000 K and remarked that small variations in Ṁ could lead to dwarf novae type outbursts. Honeycutt, Robertson & Turner (1998) performed a long-term photometric monitoring of UU Aqr which confirmed the high and low brightness states of BSC94 and revealed the existence of small amplitude ($`\stackrel{<}{}1.0`$ mag) brightness variations on timescales of a few days, which they called ‘stunted outbursts’.
The detailed spectroscopic study of Hoard et al. (1998) reinforced the classification of UU Aqr as an SW Sex star. They found evidences for the presence of a bright spot at the impact site of the gas stream with the edge of the disc, and a non-axisymmetric, vertically and azimuthally extended absorbing structure in the disc. They proposed an explanation for the absorbing structure as well as for the other spectroscopic features of UU Aqr in terms of the explosive impact of the accretion stream with the disc. Optical and ultraviolet spectroscopy by Kaitchuck et al. (1998) shows a secondary eclipse at phase 0.4 in the optical and Balmer lines (but not in the UV continuum or lines) which they suggested may be caused by an occultation of the bright spot and stream region by material suspended above the inner disc.
In this paper we report on the analysis of time-resolved spectroscopy of UU Aqr with multi-wavelength eclipse mapping techniques to derive spatially-resolved spectra of the accretion flow in this binary. Section 2 describes the observations and data reduction procedures, while section 3 describes the analysis of the light curves with eclipse mapping techniques. Section 4 presents eclipse maps at selected wavelengths, the radial intensity and brightness temperature distributions, spatially resolved spectra of the accretion disc and gas stream as well as the spectrum of the uneclipsed component. The results are discussed in section 5 and summarized in section 6.
## 2 Observations
Time-resolved spectroscopy covering 5 eclipses of UU Aqr was obtained with the 2.1-m telescope at the Kitt Peak National Observatory (KPNO) on July-August 1993 in the spectral range 3500–6900 Å (spectral resolution of $`\mathrm{\Delta }\lambda =1.5`$ Å pixel<sup>-1</sup>). The observations consist of 5 sets of $`100`$ short exposure ($`\mathrm{\Delta }t=30s`$) spectra at a time resolution of 50 s. A close comparison star (star C1 of BSC94) was included in the slit to allow correction of sky transparency variations and slit losses. The observations (summarized in Table 1) were performed under good (cloud-free) sky conditions and at small to moderate air masses ($`X1.4`$) except for run 1, which started while the object was still at a reasonably high zenith angle ($`X=2.2`$).
The data were bias-subtracted and corrected for flat-field and slit illumination effects using standard iraf procedures. 1-D spectra of both variable and comparison star were extracted with the optimal extraction algorithm of Horne (1986). The individual spectra were checked for the presence of possible cosmic rays and, when appropriate, were corrected by interpolation from the neighboring wavelengths. Arc-lamp observations were used to calibrate the wavelength scale (accuracy of 0.15 Å). Observations of the standard spectrophotometric stars BD+28 4211 and G191 B2B (Massey et al. 1988) were used to derive the instrumental sensitivity function and to flux calibrate the set of extracted spectra on each night. Error bars were computed taking into account the photon count noise and the sensitivity response of the instrument.
The reduced spectra were combined to produce trailed spectrograms of the variable and the comparison star for each night. The display of the trailed spectrograms of the comparison star shows that there were non negligible sky transparency variations and/or time-dependent slit losses along the runs. We defined a reference spectrum of the comparison star by computing an average of 40 spectra on night 5 corresponding to the time for which the star was closest to zenith. We normalized the spectrograms of the comparison star by dividing each spectrum by the reference spectrum. A 2-D cubic spline fit was used to produce a smoothed version of the normalized spectrograms. The sky transparency variations and variable slit losses were corrected by dividing the spectrogram of the variable by the smoothed, normalized spectrogram of the comparison star on each night (a procedure analogous to the flat-field correction). The reference spectrum is consistent with the UBVRI photometry of star C1 (BSC94) at the 1-$`\sigma `$ level. Therefore, the absolute photometric accuracy of these observations should be better than 10 per cent.
Fig. 1 shows average out-of-eclipse and mid-eclipse spectra of UU Aqr on 1995 August 13. The spectra are dominated by strong single-peaked Balmer emission lines but also show He I lines and the blend of C III, N III and He II lines at $`4650`$ Å. The emission lines have asymmetrical shapes, the red side of the line being stronger – in accordance with the results of Hessman (1990) and Diaz & Steiner (1991). The He I lines and the higher energy Balmer lines show a possible double peak structure suggesting either classical double-peaked emission from a highly inclined disc or single peaked emission with a central absorption component. While the continuum is reduced by a factor $`3`$ during eclipse, the emission lines suffer a much smaller reduction in flux suggesting that they possibly arise from a vertically-extended source larger than the accretion disc (responsible for the continuum emission), in accordance with inferences drawn by BSC94.
Fig. 2 shows lightcurves of the 5 runs in the broad band $`50006500`$ Å. The gap in run 5 is due to an interruption of the observations to check the telescope focus. The lightcurves have similar eclipse shapes and out of eclipse flux levels, with variations at the level of $`\stackrel{<}{}20`$ per cent between the runs. Indications that the observations were performed while UU Aqr was in its high brightness state come from the eclipse shape and average out of eclipse flux level. The latter suggests that the object was even slightly brighter than the typical high brightness state of BSC94. These remarks are in agreement with the historical lightcurve of Honeycutt et al. (1998, see their fig. 1), which shows that UU Aqr reached a maximum of its long-term average brightness level during 1993, the epoch of our observations. The spectral range of the lightcurves in Fig. 1 corresponds roughly to the $`V`$ band. The average out of eclipse level of all runs yields an approximate mean magnitude of $`V=13.2\pm 0.2`$ mag, consistent with the value drawn from the lightcurve of Honeycutt et al. (1998), of $`V=13.4\pm 0.6`$ mag.
## 3 Data analysis
### 3.1 Light curve construction
The spectra were divided into 226 passbands of 15 Å in the continuum and fainter lines, and $`500kms^1`$ across the most prominent lines. For each passband a lightcurve was extracted by computing the average flux on the corresponding wavelength range and phase folding the resulting data according to the ephemeris of BSC94. A phase correction of $`0.003`$ cycle was further applied to the data to make the centre of the white dwarf eclipse coincident with phase zero. For those passbands including emission lines the light curves comprise the total flux at the corresponding bin with no subtraction of a possible underlying continuum contribution.
Since the dataset correspond to the same brightness level it was possible to combine the lightcurves of all runs to produce average lightcurves for each passband. This is helpful to increase the signal-to-noise ratio of the lightcurves and to reduce the influence of flickering in the eclipse shape. For each passband, we first normalized the individual lightcurves by fitting a spline function to the phases outside eclipse and dividing the lightcurve by the fitted spline. The normalized lightcurves were combined by separating the data into phase bins of 0.0038 cycle and computing the median for each bin. The median of the absolute deviations with respect to the median is taken as the corresponding uncertainty. The resulting lightcurve is scaled back to flux units by multiplying the combined lightcurve by the median flux of the spline functions at phase zero. This procedure removes orbital variations outside eclipse with only minor effects on the eclipse shape itself.
### 3.2 Eclipse mapping
The eclipse mapping method was used to solve for a map of the disc brightness distribution and for the flux of an additional uneclipsed component in each passband. For the details of the method the reader is referred to Horne (1985, 1993), Baptista & Steiner (1993) and Rutten et al. (1994).
For our analysis we adopted the same eclipse map of BSH96, a $`51\times 51`$ pixel grid centred on the primary star with side $`2R_{L1}`$ where $`R_{L1}`$ is the distance from the disc centre to the inner Lagrangian point. This choice provides maps with a nominal spatial resolution of $`0.039R_{L1}`$, comparable to the expected size of the white dwarf in UU Aqr ($`0.032R_{L1}`$). The eclipse geometry is specified by the mass ratio $`q`$ and the inclination $`i`$. We adopted the parameters of BSC94, $`i=78\mathrm{°}`$ and $`q=0.3`$. The specific intensities in the eclipse map were computed assuming R$`{}_{L1}{}^{}=0.74R_{}`$ (BSC94) and a distance of 200 pc (BSH96).
The statistical uncertainties of the eclipse maps were estimated with a Monte Carlo procedure (e.g., Rutten et al. 1992; Baptista et al. 1995). For a given narrow-band lightcurve a set of 10 artificial lightcurves is generated, in which the data points are independently and randomly varied according to a Gaussian distribution with standard deviation equal to the uncertainty at that point. The lightcurves are fitted with the eclipse mapping algorithm to produce a set of randomized eclipse maps. These are combined to produce an average map and a map of the residuals with respect to the average, which yields the statistical uncertainty at each pixel. The uncertainties obtained with this procedure will be used when estimating the errors in the derived radial temperature and intensity profiles as well as in the spatially-resolved spectra.
Average light curves, fitted models, and eclipse maps at selected passbands are show in Figs. 3 and 4. These will be discussed in detail in section 4.
## 4 Results
### 4.1 Accretion disc structure
In this section we compare eclipse maps at selected passbands in order to study the structure of the accretion disc at different wavelengths.
Fig. 3 shows lightcurves (left panels) and eclipse maps (right panels) of 4 selected continuum passbands close to the Johnson-Cousins UBVR effective wavelengths in order to allow a comparison with the results of BSH96. Dashed horizontal lines depict the uneclipsed component in each case. The continuum lightcurves show a deep eclipse with a slightly asymmetric egress shoulder which is more pronounced for longer wavelengths. This results in eclipse maps with brightness distributions concentrated towards disc centre and asymmetric structures in the trailing quadrant of the disc closest to the secondary star (the upper right quadrant in the eclipse maps of Fig. 3). The uneclipsed component at $`\lambda 3657`$ is perceptibly larger than at $`\lambda 4411`$, suggesting that the Balmer jump is in emission and that the uneclipsed light has an important contribution from optically thin gas. This is in line with previous results by BSC94 and BSH96. The eclipse shapes and out of eclipse levels resemble those of the high brightness state observed by BSH96, although with a less pronounced asymmetry at eclipse egress. Accordingly, the eclipse maps clearly lack the noticeable asymmetric structure at disc edge which was the main characteristic of the high state (BSH96, see their Fig. 3). We will return to this point in section 5.
Fig. 4 shows lightcurves and eclipse maps for the line centre passbands of H$`\alpha `$, H$`\beta `$, H$`\gamma `$ and He I $`\lambda `$5876. We remark that the line lightcurves include the total flux at the corresponding wavelength range with no subtraction of an interpolated continuum. The eclipses are shallow, leading to brightness distributions which are flatter than those of the continuum. Similar to the continuum maps, the asymmetry in the egress shoulder is more pronounced for the lines at longer wavelengths. The uneclipsed components are considerably larger than in the continuum, indicating that the uneclipsed spectrum has strong Balmer and He I emission lines. The large error bars of the H$`\alpha `$ centre lightcurve is due not to low signal-to-noise ratio but to the variability of the eclipse shape at this wavelength. This effect is also seen, although to a lesser extent, in H$`\beta `$ and H$`\gamma `$.
Fig. 5 shows (Doppler) velocity-resolved lightcurves (left) and eclipse maps (right) across the H$`\beta `$ line. There is marginal evidence of rotational disturbance: the minimum of the blue bin lightcurve ($`494kms^1`$) is slightly displaced towards negative phases while that of the red bin lightcurve ($`+494kms^1`$) is correspondingly displaced towards positive phases, suggesting that the line emitting gas rotates in the prograde sense. However, the eclipse maps in the symmetric velocity bins do not show the mirror symmetry (over the line joining both stars) expected for line emission from a Keplerian disc around the white dwarf. Equally remarkable are the facts that the lightcurve in the red bin has a much larger out-of-eclipse flux than its blue counterpart and that the corresponding eclipse map is perceptibly brighter than that of the blue bin anywhere. A similar behaviour is found in the other lines for which velocity-resolved maps were obtained. This cannot be attributed to the underlying continuum since the interpolated continuum has essentially a constant level across each line. It seems clear that most of the line emission does not arise from a disc in Keplerian rotation.
### 4.2 Radial temperature distribution and mass accretion rate estimate
The simplest way of testing theoretical disc models is to convert the intensities in the eclipse maps to blackbody brightness temperatures, which can then be compared to the radial run of the effective temperature predicted by steady state, optically thick disc models. However, as discussed by Baptista et al. (1998), a relation between the effective temperature and a monochromatic brightness temperature is non-trivial, and can only be properly obtained by constructing self-consistent models of the vertical structure of the disc. Therefore, our analysis here is meant as preliminary, and should be complemented by detailed disc spectrum modeling in a future paper.
Fig. 6 shows brightness temperature radial distributions for the continuum maps of Fig. 3 in a logarithmic scale. Each temperature shown is the blackbody brightness temperature that reproduces the observed surface brightness at the corresponding pixel assuming a distance of 200 pc to UU Aqr (BSH96). Steady-state disc models for mass accretion rates of $`10^{8.5}`$, $`10^9`$, $`10^{9.5}`$ and $`10^{10}M_{}yr^1`$ are plotted as dotted lines for comparison. These models assume M$`{}_{1}{}^{}=0.67M_{}`$ and $`R_1=0.012R_{}`$ (BSC94).
The distributions resemble those obtained by BSH96 for the high brightness state of UU Aqr, closely following the $`TR^{3/4}`$ law for steady accretion in the intermediate and outer disc regions ($`R0.2R_{\mathrm{L1}}`$) but displaying a noticeable flattening in the inner disc ($`R<0.1R_{\mathrm{L1}}`$). Temperatures range from $`18000`$ K in the inner disc to 6000 K in the outer disc regions, leading to inferred mass accretion rates of Ṁ= $`10^{9.0\pm 0.3}M_{}yr^1`$ at $`R=0.1R_{\mathrm{L1}}`$ and $`10^{8.7\pm 0.2}M_{}yr^1`$ at $`R=0.3R_{\mathrm{L1}}`$ — in good agreement with the results of BSH96 for the high brightness state. The quoted errors on Ṁ account for the statistical uncertainties in the eclipse maps, obtained from the Monte Carlo procedure described in section 3.2, and the scatter in the temperatures of maps at different wavelengths. The eclipse map at $`\lambda 3657`$ leads to temperatures which are systematically higher than those of the other continuum maps of Fig. 3, in an example of the limitations of using brightness temperatures to estimate the mass accretion rate. This difference reflects the fact that the Balmer jump appears in emission for the intermediate and outer disc regions, as will be seen in section 4.4.
### 4.3 Radial line intensity distributions
Left panel in Fig. 7 shows radial intensity distributions for the most prominent lines (solid) and adjacent continuum (dotted) in a logarithmic scale. The line distributions were obtained from the average of all eclipse maps across the line region, while the continuum distributions were obtained from the average of eclipse maps on both sides of each line. Net line emission distributions were computed by subtracting the distributions of the adjacent continuum from those of the lines, and are shown in the right panel. In the external map regions ($`R\stackrel{>}{}0.7R_{\mathrm{L1}}`$) the intensities of both line and continuum drop by a factor $`10^3`$ with respect to the inner disc regions, making the computation of the net emission quite noisy and unreliable. H$`\alpha `$ is seen in emission (intensities larger than those at the adjacent continuum) at all disc radii and up to $`R0.6R_{\mathrm{L1}}`$. The other lines are in absorption in the inner disc and transition to emission at intermediate ($`R0.2R_{\mathrm{L1}}`$) disc radius. This behaviour is noticeably different from that observed at the low brightness state, where H$`\alpha `$ is seen in emission in the inner disc and disappears into the continuum for $`R0.3R_{\mathrm{L1}}`$ (BSH96). This result suggests that the line emission region increases in size from the low to the high brightness state, possibly in response to changes in mass accretion rate. The transition from absorption to emission occurs at larger disc radii for lines of higher excitation. This can be explained, for the Balmer lines, by the increase in continuum emission at the inner disc for shorter wavelengths.
A set of dotted lines in the right panel indicate the slope of the empirical radial dependency of the line emissivity in accretion discs, $`IR^{1.5}`$, as inferred from Doppler Tomography by assuming a Keplerian distribution of velocities for the emitting gas (Marsh et al. 1990). For H$`\gamma `$ and He I $`\lambda `$5876, the net emission occurs for a narrow range of radii making a comparison with the empirical law difficult. The derived radial distributions for H$`\alpha `$ and H$`\beta `$ are clearly different from the empirical $`IR^{1.5}`$ law; in particular, the H$`\alpha `$ distribution is flat at inner and intermediate disc radii ($`R<0.3R_{\mathrm{L1}}`$). This remark suggests that the line emitting regions on the disc surface are not in Keplerian orbits or that a substantial fraction of the emission lines does not arise from the accretion disc, in line with the inferences drawn by the comparison of velocity-resolved eclipse maps in section 4.1. The latter hypothesis is consistent with the significant uneclipsed components inferred for the Balmer and He I lines (section 4.4.2).
### 4.4 Spatially resolved spectra
Each of the eclipse maps yields spatially-resolved information about the emitting region on a specific wavelength range. By combining all narrow-band eclipse maps we are able to isolate the spectrum of the eclipsed region at any desired position (e.g., Rutten et al. 1994; Baptista et al. 1998).
To investigate the possible influence of the gas stream on the disc emission and motivated by the observed asymmetries in the eclipse maps shown in section 4.1, we divided the disc into two major azimuthal regions to extract spatially-resolved spectra: the gas stream region (upper right quadrant in the eclipse maps of Figs. 3 and 4) and the disc region (the remaining 3/4 of the eclipse map). For each of these regions, we divided the maps into a set of 6 concentric annuli centred on the white dwarf of width $`0.1R_{L1}`$ and with radius increasing in steps of $`0.1R_{L1}`$. Each spectrum is obtained by averaging the intensity of all pixels inside the corresponding annulus and the statistical uncertainties affecting the average intensities are estimated with the Monte Carlo procedure described in section 3.2.
#### 4.4.1 Disc spectra
Fig. 8 shows spatially-resolved spectra of the disc region in a logarithmic scale. The inner annular region is at the top and each spectrum is at its true intensity level. The spectrum of the uneclipsed component is shown in the lower panel and will be discussed in detail in section 4.4.2.
The spectrum of the inner disc is characterized by a blue and bright continuum filled with deep and narrow absorption lines. The continuum emission becomes progressively fainter and redder for increasing disc radius while the lines transition from absorption to emission showing clear P Cygni profiles on all lines mapped at higher spectral resolution. The Balmer jump appears in absorption in the inner disc and weakly in emission in the intermediate and outer disc regions suggesting that the outer disc in UU Aqr is optically thin. The change in the slope and intensity of the continuum with increasing disc radius reflects the temperature gradient in the accretion disc, with the effective temperature decreasing outwards.
The spatially resolved spectra of the disc are plotted in Fig. 9 as a function of velocity for the H$`\alpha `$, H$`\beta `$ and H$`\gamma `$ regions. Vertical dotted lines mark line centre and the maximum blueshift/redshift velocity expected for gas in Keplerian orbits around a $`0.67M_{}`$ white dwarf as seen from an inclination of $`i=78\mathrm{°}`$ ($`v\mathrm{sin}i=3200kms^1`$) \[BSC94\].
The absorption lines at disc centre are perceptibly narrower than expected for emission from either the white dwarf atmosphere or from disc gas in Keplerian orbits around the white dwarf. The discrepancy increases if the larger mass estimates of Diaz & Steiner (1991) and Kaitchuck et al (1998) are assumed for the white dwarf. Moreover, the absorption lines at disc centre are deep, while lines produced in a white dwarf atmosphere or innermost disc regions should be broad and shallow. The width of the lines indicate a velocity dispersion of $`1500kms^1`$ for the line emitting region in the line of sight to the disc centre and higher velocities ($`2000kms^1`$) for the gas in the outer disc at $`R0.5R_{\mathrm{L1}}`$. This is in clear disagreement with the expected behaviour of line emission from gas in a Keplerian disc and provide additional evidence that these lines do not arise from the disc atmosphere. On the other hand, the lines at intermediate and outer disc regions ($`R\stackrel{>}{}0.2R_{\mathrm{L1}}`$) show clear P Cygni profiles indicating origin in an outflowing gas, probably the disc wind.
We note that the H$`\alpha `$ line shows a redshifted ($`v1800kms^1`$) absorption component in spectra of the outer disc regions ($`R>0.3R_{L1}`$). Comparison of disc spectra at different azimuths shows that this absorption is produced in the front side of the disc, but an origin in the gas stream can possibly be ruled out since the absorption component is seen with similar strengths in the leading and trailing (the one containing the gas stream) quadrants. The interpretation of this feature is not straightforward and deserves a bit of caution, since it is not clearly seen in any other line and also because the surface brightness in the corresponding disc region is only a few percent of the intensities in the inner disc.
#### 4.4.2 The uneclipsed spectrum
The spectrum of the uneclipsed light (lower panel of Fig. 8) show prominent Balmer and He I emission lines. The Balmer jump is clearly in emission and the optical continuum rises towards longer wavelengths suggesting that the Paschen jump is also in emission. These results are consistent with the findings of BSH96 and indicate that the uneclipsed light has an important contribution from optically thin gas from outside the orbital plane. The Balmer lines mapped at higher spectral resolution show broad asymmetric profiles, with line peaks displaced to the red side and wings extending up to $`1500kms^1`$. The observed asymmetry is consistent with that previously seen in the integrated spectra of Diaz & Steiner (1991) and Hessman (1990) and is similar to that observed in the resonant ultraviolet lines of UX UMa, where the uneclipsed component was attributed to emission in a vertically-extended disc wind (Baptista et al. 1995; 1998; Knigge & Drew 1997).
The fractional contribution of the uneclipsed component to the total flux was obtained by dividing the flux of the uneclipsed light by the average out of eclipse level at each passband. The result is shown in Fig. 10. The fractional contribution of the uneclipsed light is very significant for the optical emission lines, reaching 40-60 per cent at the Balmer lines and 20-40 per cent at the He I lines, and decreases steadily along the Balmer series. The difference in fractional contribution between the Balmer and He I lines and among the Balmer lines indicates the existence of a vertical temperature gradient in the material above/below the disc, with the He I lines (which require higher excitation energies) being produced closer to the orbital plane. In any case, a substantial fraction of the light at these lines does not arise from the orbital plane and is not occulted during eclipse. The uneclipsed component gives significant contribution also to the continuum emission. About 20 per cent of the flux at the Balmer continuum and similar fraction of the continuum emission at the red end of the spectrum arise from regions outside the orbital plane.
#### 4.4.3 The gas stream region
Fig. 11 shows the ratio between the spectrum of the gas stream region and the disc region at same radius as a function of radius. A dotted line marks the unity level for each panel. The comparison shows that the spectrum of the gas stream is noticeably different from the disc spectrum in the outer disc regions (where one expects a bright spot to form due to the shock between the inflowing stream and the outer disc rim), but also reveals systematic differences between stream and disc spectra in a range of radii. In all cases, the stream emission is stronger than that of the adjacent disc.
This result suggests that the material in the gas stream continues to flow downstream beyond the bright spot position. In this regard, there are three potential cases: (1) stream overflow that arcs high above/below the disc until it hits the disc surface again at a downstream position closer to the white dwarf (hereafter called the “classical stream overflow”); (2) stream overflow that continuously skims the disc surface (hereafter called the “disc-skimming overflow”); and (3) the stream drills into the disc at the impact site (hereafter named the “disc stream penetration”). The latter case seems physically unrealistic and is not supported by hydrodynamic calculations of stream-disc interaction (Armitage & Livio 1996, 1998). The fact that there is enhanced emission in the stream region extending all the way from the outer disc down to $`R0.2R_{L1}`$ argues in favor of an interpretation in terms of disc-skimming overflow instead of the classical stream overflow – as has been suggested to explain the behaviour of SW Sex stars (Hellier & Robinson 1994; Hellier 1996) – since the latter would produce enhanced emission only at the position of the two spots, at the initial impact site in the outer disc edge and at the re-impact site much closer to disc centre (Lubow 1989).
The spectrum of the ratio becomes redder for decreasing disc radius, possibly a combination of the disc emission becoming bluer as one moves inwards and the gas stream emission becoming redder while its energy is continuously lost in the shock with disc material along the inward stream trajectory. This is reminiscent of what was seen in ultraviolet eclipse observations of the dwarf novae IP Peg in quiescence, which revealed a compact blue bright spot with an extended red tail (Baptista et al. 1993).
## 5 Discussion
In this section we present and discuss some possible interpretations for the results of section 4 in the context of the current models for the SW Sex stars.
### 5.1 Where do the lines come from?
In previous sections we have accumulated evidences that the behaviour of the UU Aqr lines in its high state is not consistent with emission in a disc atmosphere, namely: (i) negligible rotational disturbance, (ii) no mirror symmetry between eclipse maps in symmetric velocity bins; (iii) H$`\alpha `$ line emission distribution much flatter than the empirical $`IR^{1.5}`$ law; (iv) significant uneclipsed components, and (v) presence of P Cygni profiles in the disc spectra at intermediate and large disc radii. If the lines do not arise in the disc atmosphere, where do they come from?
The most compelling interpretation is that the lines are produced in a disc chromosphere + wind. This region is hot, dense, opaque and has low expansion velocities close to the orbital plane in order to produce the observed deep, narrow absorption lines in the line of sight to the inner disc. Most of the high excitation lines are produced close to the disc plane. The density and temperature decrease with height above/below the disc as the outflowing gas spreads over an increasing surface area. Optically thin emission from this extended region is probably responsible for the Balmer jump (and lines) in emission observed in the uneclipsed spectrum. Support in favor of this scenario comes from the recent detailed modeling of the C IV wind line of eclipsing nova-likes by Schlosman, Vitello & Mauche (1996) and Knigge & Drew (1997). Their results suggest the existence of a relatively dense ($`n_e4\times 10^{12}`$ cm<sup>-3</sup>) and vertically extended chromosphere between the disc surface and the fast-moving parts of the wind, which could produce significant amounts of optically thin emission. At orbital phases around eclipse, gas outflowing in the direction of the secondary star will be seen along the line of sight to the bright underlying accretion disc as blueshifted absorption features, while gas expelled in the direction away from the secondary star should contribute with redshifted emission.
We tested this scenario by comparing spatially resolved spectra of the disc lune closest to the secondary star (the right hemisphere of the disc in the eclipse maps of Fig. 3, hereafter called the “front” side) and of the disc lune farthest away from the secondary star (the left hemisphere of the disc in Fig. 3, hereafter called the “back” side). For this purpose, we defined two opposite azimuthal disc regions of width $`30\mathrm{°}`$ along the major axis of the binary, and extracted spatially resolved spectra for the same set of annuli as above. These spatially resolved spectra are noisier than those of Figs. 8 and 9 because in this case the average intensity of each annulus is computed from a significantly smaller number of pixels. The results are shown in Fig. 12 for the H$`\beta `$ and H$`\gamma `$ regions and are consistent with our interpretation: The blueshifted absorption component is seen mainly in the front side of the disc while the redshifted emission is generally more prominent in spectra of the back side of the disc. The fact that the blueshifted absorption can still be seen projected along the line of sight at the outer regions of the disc favours a more spherical or equatorial geometry for the outflowing gas instead of a highly collimated, polar jet.
The chromosphere + disc wind interpretation satisfactorily accounts for all the features listed above and also gives a plausible explanation for (1) the distinct semi-amplitude of the radial velocity $`K`$ and systemic velocity $`\gamma `$ as inferred from different emission lines and (2) the time dependent $`K`$ and $`\gamma `$ values (Hoard et al. 1998). The centroid of lines of different excitation level will occur at different locations in the primary lobe and will sample different velocities along the line of sight. With respect to (2), the comparison of the H$`\alpha `$ map of BSH96 (which corresponds to the low brightness state) with that of Fig. 4 (the high state) reveals that in the latter the emission extends over a much larger region of the primary lobe with a pronounced asymmetry in the stream region, suggesting that the wind emission is variable in time and is intimately connected with the mass accretion rate. This remark gives additional support to the suggestion by Hoard et al. (1998) that the observed time dependence of the $`K`$ and $`\gamma `$ velocities might be due to variability of a wind component in UU Aqr.
An alternative possibility is to consider the deep absorption lines seen towards the line of sight to the disc centre as being produced by absorption in a vertically extended disc rim. Although this scenario accounts for the narrow absorption lines, it is not able to explain the large velocities inferred from the line width for intermediate and large disc radius nor the P Cygni profiles. Furthermore, it should result in a perceptible front-back asymmetry in the disc surface brightness (namely, the back side of the disc should be brighter) which is not seen in the eclipse maps.
Recently, Horne (1999) proposed that most of the features of the SW Sex stars could be explained in terms of a disc-anchored magnetic propeller, in which energy and angular momentum are extracted from the magnetic field of the inner disc regions to fling part of the material in the gas stream out of the binary towards the back side of the disc. Although this model is able to explain many of the observed features of UU Aqr, it can only account for the observed P Cygni profiles if the gas trapped by the inner disc magnetic field is expelled in all directions and not only towards the back of the disc. We note that, in this case, there is no significant difference between the propeller and the disc wind models and, in fact, the former could possibly work as the underlying physical mechanism driving the latter.
If disk-skimming overflow does occur, we might expect that dissipation of energy in the collision between the gas stream and the disc material gives rise to a bulge extending along the stream trajectory over and under the disc. This bulge will appear in front of the chromosphere + wind line emitting region at the inner disc when seen along the line of sight at orbital phases 0.5-0.9. This may explain the phase-dependent absorption lines, observed from phases 0.5-0.9 and with maximum at phase $`0.8`$ (Heafner 1989; Hoard et al. 1998). The enhanced line emission along the gas stream (see Fig. 4) is possibly responsible for the phase offset between photometric and spectroscopic conjunction (Diaz & Steiner 1991; Hoard et al. 1998).
In summary, the picture which emerges from our results is consistent with the results from the Doppler tomography and the model proposed for UU Aqr by Hoard et al. (1998).
### 5.2 Where has the bright spot gone?
Although our observations correspond to the high brightness state of UU Aqr, our eclipse maps do not show the conspicuous asymmetric structure seen in the high state eclipse maps of BSH96 and which was interpreted as being the bright spot. The explanation for the ‘disappearance’ of the bright spot may be connected with the stunted outbursts found by Honeycutt et al. (1998).
BSH96 pointed out that the inferred accretion rate of UU Aqr is close to the critical mass accretion rate for disc instability to occur and remarked that the long-term lightcurves of accretion discs with mass transfer rates near their critical limit might display low-amplitude ($`\stackrel{<}{}1.0`$ mag) outbursts caused by thermal instabilities in the outer disc regions (e.g., Lin, Papaloizou & Faulkner 1985). In this case the outburst is restricted to the outer 1/3 of the disc extent while the inner disc remains in a high viscosity, steady state. Honeycutt et al (1998) suggested that such dwarf-nova type instabilities could be an explanation for the stunted outbursts of UU Aqr if a mechanism can be identified to make the amplitudes appear small. We note that the observed low amplitudes can be easily accounted for by the reduced contrast of the light from the outbursting outer regions – where the efficiency in transforming gravitational potential energy in radiation is relatively low – in comparison to the bright, optically thick and steady inner disc.
If the observed stunted outbursts of UU Aqr are caused by thermal instabilities in its outer disc, the disc radius is expected to increase during the outburst and will eventually reach the 3:1 tidal resonance radius leading to an elliptical precessing disc reminiscent of what possibly happens in SU UMa stars in superoutburst (e.g., Warner 1995 and references therein). We suggest that the azimuthally elongated structure seen in the eclipse maps of BSH96 is the signature of such an elliptical disc and not the bright spot. Following this line of reasoning, this structure should not be present when the disc radius is smaller than the tidal resonance radius. Support for this interpretation comes from the comparison of disc radius in the high state eclipse maps of BSH96 and our eclipse maps. From BSH96 data we estimate a disc radius of $`R_d0.7R_{\mathrm{L1}}`$, comparable to the 3:1 tidal resonance radius for a mass ratio of $`q=0.3`$. Our eclipse maps lead to a smaller value of $`R_d=0.65R_{\mathrm{L1}}`$. Therefore, we suggest that UU Aqr was in an occasional superhumper state during the high brightness state observations of BSC94.
In the model of Hoard et al. (1998), after the explosive impact of the high Ṁ accretion stream with the edge of the disc, the incomming gas forms an optically thick absorbing bulge on the disc that either follows roughly the stream trajectory or runs along the rim of the disc, producing the absorption features seen at phases 0.4-0.9. It may alternatively be possible that the structure seen in the eclipse maps of BSH96 is the signature of such post-impact stream material running along the edge of the disc. In this scenario, the azimuthally extended bulge would be present or not depending on the (variable) mass accretion rate and the resulting orbital hump would remain fixed in phase.
It would be interesting (although outside the scope of this paper) to reanalize the data of BSC94 to see if the orbital hump present in the high state precesses in phase in a similar manner as superhumps in superoutbursts (supporting the elliptical disc scenario) or if its maximum occurs always at the same orbital phase range about 0.8 - 0.9 cycle (favouring the post-impact bulge scenario).
## 6 Conclusions
We used time-resolved spectroscopy to study the structure and spectra of the accretion disc and gas stream of the novalike UU Aquarii in the optical range. The main results of this analysis can be summarized as follows:
* The spectrum of the inner disc shows a blue continuum filled with deep, narrow absorption lines which transition to emission with clear P Cygni profiles at intermediate and large radii ($`R\stackrel{>}{}0.2R_{L1}`$).
* The spectrum of the uneclipsed light has strong H I and He I emission lines and a Balmer jump in emission indicating a significant contribution from optically thin regions outside the orbital plane.
* Velocity-resolved eclipse maps and spectra indicate that most of the line emission probably arises in a vertically-extended disc chromosphere + wind.
* Differences in fractional contribution among emission lines suggests a vertical temperature gradient in the material above/below the disc.
* The comparison of the spectrum of the gas stream region and the disc region at the same radius as a function of radius gives evidence of gas stream disc-skimming overflow down to $`R0.2R_{L1}`$. This may explain the phase-dependent absorption in emission lines.
* The comparison of our eclipse maps with those of BSH96 suggests that the asymmetric structure in the outer disc previously identified as the bright spot may be the signature of an elliptical precessing disc similar to those possibly present in SU UMa stars during superoutbursts.
## Acknowledgments
We gratefully acknowledge the director of KPNO for granting telescope time for this project at the Summer Queue Program, Tod Boroson and the team of observers at KPNO for their kind effort in collecting the data, Knox Long and the director of STScI for financial support through the Director Discretionary fund, Susan Keener for helping with the data reduction at STScI, and an anonymous referee for valuable comments and suggestions that helped to improve the presentation of the results. RB acknowledges financial support from CNPq/Brazil through grant no. 300 354/96-7. This work was partially supported by PRONEX grant FAURGS/FINEP 7697.1003.00. |
warning/0002/astro-ph0002525.html | ar5iv | text | # On the particle acceleration near the light surface of radio pulsars
## 1 Introduction
Despite the fact that the structure of the magnetosphere of radio pulsars remains one of the fundamental astrophysical problems, the common view on the key theoretical question – what is the physical nature of the neutron star braking – is absent (Michel 1991, Beskin Gurevich & Istomin 1993, Mestel 1999). Nevertheless, very extensive theoretical studies in the seventies and the eighties allowed to obtain some model-independent results. One of them is the absence of magnetodipole energy loss. This result was first obtained theoretically (Henriksen & Norton 1975, Beskin et al 1983). It was shown that the electric charges filling the magnetosphere screen fully the magnetodipole radiation of a neutron star for an arbitrary inclination angle $`\chi `$ between the rotational and magnetic axes if there are no longitudinal currents flowing in the magnetosphere. Later this result was also confirmed by observations. The direct detections of the interaction of the pulsar wind with a companion star in close binaries (see e.g. Djorgovsky & Evans 1988, Kulkarni & Hester 1988) have shown that it is impossible to explain the heating of the companion by a low–frequency magnetodipole wave.
On the other hand, the detailed mechanism of particle acceleration remains unclear. Indeed, a very high magnetization parameter $`\sigma `$ (Michel 1969) in the pulsar magnetosphere demonstrates that within the light cylinder $`r<R_\mathrm{L}=c/\mathrm{\Omega }`$ the main part of the energy is transported by the Poynting flux. It means that the additional mechanism of particle acceleration must work in the vicinity of the light cylinder. It is necessary to stress that an effective particle acceleration can only take place for small enough longitudinal electric currents $`I<I_{GJ}`$ when the plasma has no possibility to pass smoothly through the fast magnetosonic surface and when the light surface $`|𝑬|=|𝑩|`$ is located at a finite distance. As to the case of the large longitudinal currents $`I>I_{GJ}`$, both analytical (Tomimatsu 1994, Begelman & Li 1994, Beskin et al 1998) and numerical (Bogovalov 1997) considerations demonstrate that the acceleration becomes ineffective outside the fast magnetosonic surface, and the particle-to-Poynting flux ratio remains small: $`\sigma ^{2/3}`$ (Michel 1969, Okamoto 1978).
The acceleration of an electron–positron plasma near the light surface was considered by Beskin Gurevich and Istomin (1983) in the simple $`1D`$ cylindrical geometry for $`II_{GJ}`$. It was shown that in a narrow boundary layer $`\mathrm{\Delta }\varpi /\varpi 1/\lambda `$ almost all electromagnetic energy is actually converted to the particles energy. Nevertheless, cylindrical geometry does not provide the complete picture of particle acceleration. In particular, it was impossible to include self–consistently the disturbance of a poloidal magnetic field and an electric potential, the later playing the main role in the problem of the plasma acceleration (for more details see e.g. Mestel & Shibata 1994). Hence, a more careful $`2D`$ consideration is necessary.
In Sect. 2 we formulate a complete system of $`2D`$ two–fluid MHD equations describing the electron–positron outflow from a magnetized body with a monopole magnetic field. The presence of an exact analytical force–free solution (Michel 1973) allows us to linearize this system which results in the existence of invariants (energy and angular momentum) along unperturbed monopole field lines similar to the ideal one–fluid MHD flow. In Sect. 3 it is shown that for $`\sigma 1`$ and $`\lambda 1`$ ($`\lambda =n/n_{GJ}`$ is the multiplication factor) the one–fluid MHD approximation remains true in the entire region within the light surface. Finally, in Sect. 4 the acceleration of particles near the light surface $`|𝑬|=|𝑩|`$ is considered. It is shown that, as in the case of cylindrical geometry, in a narrow boundary layer $`\mathrm{\Delta }\varpi /\varpi \lambda ^1`$ almost all the electromagnetic energy is converted into the energy of particles.
## 2 Basic Equations
Let us consider a stationary axisymmetric outflow of a two–component plasma in the vicinity of an active object with a monopole magnetic field. It is necessary to stress that, of course, the monopole magnetic field is a rather crude approximation for a pulsar magnetosphere. Nevertheless, even for a dipole magnetic field near the origin, at large distances $`rR_\mathrm{L}`$ in the wind zone the magnetic field can have a monopole–like structure. For this reason the disturbance of a monopole magnetic field can give us an important information concerning particle acceleration far from the neutron star.
The structure of the flow is described by the set of Maxwell‘s equations and the equations of motion
$`𝑬=4\pi \rho _e,\times 𝑬=0,`$
$`𝑩=0,\times 𝑩={\displaystyle \frac{4\pi }{c}}𝒋,`$ (1)
$`(𝒗^\pm )𝒑^\pm =\pm e\left(𝑬+{\displaystyle \frac{𝒗^\pm }{c}}\times 𝑩\right).`$ (2)
Here $`𝑬`$ and $`𝑩`$ are the electric and magnetic fields, $`\rho _e`$ and $`𝐣`$ are the charge and current densities, and $`𝒗^\pm `$ and $`𝒑^\pm `$ are the speed and momentum of particles. In the limit of infinite particle energy
$$\gamma =\mathrm{},v_r^0=c,v_\phi ^0=0,v_\theta ^0=0,$$
(3)
and for charge and current density
$$\rho _e^0=\rho _\mathrm{s}\frac{R_\mathrm{s}^2}{r^2}\mathrm{cos}\theta ,j_r=\rho _\mathrm{s}c\frac{R_\mathrm{s}^2}{r^2}\mathrm{cos}\theta ,$$
(4)
the monopole poloidal magnetic field
$$B_r^0=B_\mathrm{s}\frac{R_\mathrm{s}^2}{r^2},B_\theta ^0=0,$$
(5)
is the exact solution of Maxwell’s equations. In this case
$$B_\phi ^0=E_\theta ^0=B_\mathrm{s}\frac{R_\mathrm{s}\mathrm{\Omega }}{c}\frac{R_\mathrm{s}}{r}\mathrm{sin}\theta ,E_r^0=E_\phi ^0=0,$$
(6)
which coincides with the well–known Michel (1973) solution. Here $`\gamma `$ is the Lorentz-factor of particles, $`B_\mathrm{s}`$ is the magnetic field on the surface of the sphere $`r=R_\mathrm{s}`$, and $`\rho _\mathrm{s}=`$ const. As a result, the angular velocity can be rewritten in a form $`\mathrm{\Omega }=2\pi c|\rho _\mathrm{s}|/B_\mathrm{s}`$. The limit $`\gamma \mathrm{}`$ just corresponds to zero particle mass in the force–free approximation.
It is convenient to introduce the electric field potential $`\mathrm{\Phi }(r,\theta )`$, so that $`𝑬=\mathrm{\Phi }`$ and
$$\mathrm{\Phi }^0=\frac{\mathrm{\Omega }R_\mathrm{s}^2B_\mathrm{s}}{c}\mathrm{cos}\theta ,$$
(7)
and the flux function $`\mathrm{\Psi }(r,\theta )`$, so that
$$𝑩_\mathrm{p}^0=\frac{\mathrm{\Psi }\times 𝒆_\phi }{2\pi r\mathrm{sin}\theta },$$
(8)
and $`\mathrm{\Psi }^0=2\pi B_\mathrm{s}R_\mathrm{s}^2(1\mathrm{cos}\theta )`$. Then one can seek the first–order corrections for the case $`vc`$ in the following form:
$`n^+`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Omega }B_\mathrm{s}}{2\pi ce}}{\displaystyle \frac{R_\mathrm{s}^2}{r^2}}\left[\lambda {\displaystyle \frac{1}{2}}\mathrm{cos}\theta +\eta ^+(r,\theta )\right],`$ (9)
$`n^{}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Omega }B_\mathrm{s}}{2\pi ce}}{\displaystyle \frac{R_\mathrm{s}^2}{r^2}}\left[\lambda +{\displaystyle \frac{1}{2}}\mathrm{cos}\theta +\eta ^{}(r,\theta )\right],`$ (10)
$`\mathrm{\Phi }(r,\theta )`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Omega }R_\mathrm{s}^2B_\mathrm{s}}{c}}\left[\mathrm{cos}\theta +\delta (r,\theta )\right],`$ (11)
$`\mathrm{\Psi }(r,\theta )`$ $`=`$ $`2\pi B_\mathrm{s}R_\mathrm{s}^2\left[1\mathrm{cos}\theta +\epsilon f(r,\theta )\right],`$ (12)
$`v_r^\pm `$ $`=`$ $`c\left[1\xi _r^\pm (r,\theta )\right],v_\theta ^\pm =c\xi _\theta ^\pm (r,\theta ),v_\phi ^\pm =c\xi _\phi ^\pm (r,\theta ),`$ (13)
$`B_r`$ $`=`$ $`B_\mathrm{s}{\displaystyle \frac{R_\mathrm{s}^2}{r^2}}\left(1+{\displaystyle \frac{\epsilon }{\mathrm{sin}\theta }}{\displaystyle \frac{f}{\theta }}\right),`$ (14)
$`B_\theta `$ $`=`$ $`\epsilon {\displaystyle \frac{B_\mathrm{s}R_\mathrm{s}^2}{r\mathrm{sin}\theta }}{\displaystyle \frac{f}{r}},`$ (15)
$`B_\phi `$ $`=`$ $`B_\mathrm{s}{\displaystyle \frac{\mathrm{\Omega }R_\mathrm{s}}{c}}{\displaystyle \frac{R_\mathrm{s}}{r}}\left[\mathrm{sin}\theta \zeta (r,\theta )\right],`$ (16)
$`E_r`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Omega }B_\mathrm{s}R_\mathrm{s}^2}{c}}{\displaystyle \frac{\delta }{r}},`$ (17)
$`E_\theta `$ $`=`$ $`{\displaystyle \frac{\mathrm{\Omega }R_\mathrm{s}^2B_\mathrm{s}}{cr}}\left(\mathrm{sin}\theta {\displaystyle \frac{\delta }{\theta }}\right).`$ (18)
Such a choice corresponds to a constant particle-to-magnetic flux ratio. Here $`\lambda 1`$ is the multiplication parameter ($`\lambda =en_\mathrm{s}/|\rho _\mathrm{s}|`$, where $`n_\mathrm{s}`$ is the concentration of particles on the surface $`r=R_\mathrm{s}`$) which is $`10^310^5`$ for radio pulsars. In what follows we consider for simplicity the case $`\lambda =`$ const.
Substituting (9)–(18) into equations (1)–(2), we obtain to the first-order approximation the following system of equations:
$`{\displaystyle \frac{1}{\mathrm{sin}\theta }}{\displaystyle \frac{}{\theta }}(\zeta \mathrm{sin}\theta )=2(\eta ^+\eta ^{})2\left[\left(\lambda {\displaystyle \frac{1}{2}}\mathrm{cos}\theta \right)\xi _r^+\left(\lambda +{\displaystyle \frac{1}{2}}\mathrm{cos}\theta \right)\xi _r^{}\right],`$ (19)
$`2(\eta ^+\eta ^{})+{\displaystyle \frac{}{r}}\left(r^2{\displaystyle \frac{\delta }{r}}\right)+{\displaystyle \frac{1}{\mathrm{sin}\theta }}{\displaystyle \frac{}{\theta }}\left(\mathrm{sin}\theta {\displaystyle \frac{\delta }{\theta }}\right)=0,`$ (20)
$`{\displaystyle \frac{\zeta }{r}}={\displaystyle \frac{2}{r}}\left[\left(\lambda {\displaystyle \frac{1}{2}}\mathrm{cos}\theta \right)\xi _\theta ^+\left(\lambda +{\displaystyle \frac{1}{2}}\mathrm{cos}\theta \right)\xi _\theta ^{}\right],`$ (21)
$`{\displaystyle \frac{\epsilon }{\mathrm{sin}\theta }}{\displaystyle \frac{^2f}{r^2}}{\displaystyle \frac{\epsilon }{r^2}}{\displaystyle \frac{}{\theta }}\left({\displaystyle \frac{1}{\mathrm{sin}\theta }}{\displaystyle \frac{f}{\theta }}\right)=2{\displaystyle \frac{\mathrm{\Omega }}{rc}}\left[\left(\lambda {\displaystyle \frac{1}{2}}\mathrm{cos}\theta \right)\xi _\phi ^+\left(\lambda +{\displaystyle \frac{1}{2}}\mathrm{cos}\theta \right)\xi _\phi ^{}\right],`$ (22)
$`{\displaystyle \frac{}{r}}\left(\xi _\theta ^+\gamma ^+\right)+{\displaystyle \frac{\xi _\theta ^+\gamma ^+}{r}}=4\lambda \sigma \left({\displaystyle \frac{1}{r}}{\displaystyle \frac{\delta }{\theta }}+{\displaystyle \frac{\zeta }{r}}{\displaystyle \frac{\mathrm{sin}\theta }{r}}\xi _r^++{\displaystyle \frac{c}{\mathrm{\Omega }r^2}}\xi _\phi ^+\right),`$ (23)
$`{\displaystyle \frac{}{r}}\left(\xi _\theta ^{}\gamma ^{}\right)+{\displaystyle \frac{\xi _\theta ^{}\gamma ^{}}{r}}=4\lambda \sigma \left({\displaystyle \frac{1}{r}}{\displaystyle \frac{\delta }{\theta }}+{\displaystyle \frac{\zeta }{r}}{\displaystyle \frac{\mathrm{sin}\theta }{r}}\xi _r^{}+{\displaystyle \frac{c}{\mathrm{\Omega }r^2}}\xi _\phi ^{}\right),`$ (24)
$`{\displaystyle \frac{}{r}}\left(\gamma ^+\right)=4\lambda \sigma \left({\displaystyle \frac{\delta }{r}}{\displaystyle \frac{\mathrm{sin}\theta }{r}}\xi _\theta ^+\right),`$ (25)
$`{\displaystyle \frac{}{r}}\left(\gamma ^{}\right)=4\lambda \sigma \left({\displaystyle \frac{\delta }{r}}{\displaystyle \frac{\mathrm{sin}\theta }{r}}\xi _\theta ^{}\right),`$ (26)
$`{\displaystyle \frac{}{r}}\left(\xi _\phi ^+\gamma ^+\right)+{\displaystyle \frac{\xi _\phi ^+\gamma ^+}{r}}=4\lambda \sigma \left(\epsilon {\displaystyle \frac{c}{\mathrm{\Omega }r\mathrm{sin}\theta }}{\displaystyle \frac{f}{r}}{\displaystyle \frac{c}{\mathrm{\Omega }r^2}}\xi _\theta ^+\right),`$ (27)
$`{\displaystyle \frac{}{r}}\left(\xi _\phi ^{}\gamma ^{}\right)+{\displaystyle \frac{\xi _\phi ^{}\gamma ^{}}{r}}=4\lambda \sigma \left(\epsilon {\displaystyle \frac{c}{\mathrm{\Omega }r\mathrm{sin}\theta }}{\displaystyle \frac{f}{r}}{\displaystyle \frac{c}{\mathrm{\Omega }r^2}}\xi _\theta ^{}\right).`$ (28)
Here
$$\sigma =\frac{\mathrm{\Omega }eB_\mathrm{s}R_\mathrm{s}^2}{4\lambda mc^3}$$
(29)
is the Michel‘s (1969) magnetization parameter, $`m`$ is the electron mass, and all deflecting functions are supposed to be $`1`$. It is necessary to stress that for applications the magnetic field $`B_\mathrm{s}`$ is to be taken near the light cylinder $`R_\mathrm{s}R_\mathrm{L}`$ because in the internal region of the pulsar magnetosphere $`Br^3`$. As it has already been mentioned, only outside the light cylinder the poloidal magnetic field may have quasi monopole structure. As a result,
$$\sigma =\frac{\mathrm{\Omega }^2eB_0R^3}{4\lambda mc^4}10^4B_{12}\lambda _3^1P^2,$$
(30)
where $`B_0`$ – magnetic field on the neutron star surface $`r=R`$. Hence, for ordinary pulsars ($`P1`$s, $`B_010^{12}`$G) we have $`\sigma 10^410^5`$, and only for fast ones ($`P0.10.01`$s, $`B_010^{13}`$G) we have $`\sigma 10^610^7`$.
Formally, this system of equations requires twelve boundary conditions. We consider for simplicity the case $`\mathrm{\Omega }R/c1`$ when the star radius $`R`$ is much smaller than the light cylinder. As a result, one writes down the first six boundary conditions as
$`\xi _\theta ^\pm (R_\mathrm{s},\theta )`$ $`=`$ $`0,`$ (31)
$`\xi _\phi ^\pm (R_\mathrm{s},\theta )`$ $`=`$ $`0,`$ (32)
$`\gamma ^\pm (R_\mathrm{s},\theta )`$ $`=`$ $`\gamma _{\mathrm{in}},`$ (33)
i.e. $`\xi _r^\pm (R_\mathrm{s},\theta )=1/(2\gamma _{\mathrm{in}}^2)`$. According to all theories of particle generation near the neutron star surface (Ruderman Sutherland 1975, Arons Scharlemann 1979), $`\gamma _{\mathrm{in}}10^2`$ for secondary plasma. For this reason in what follows we consider in more details the case
$$\gamma _{\mathrm{in}}^3\sigma ,$$
(34)
when the additional acceleration of particles inside the fast magnetosonic surface takes place (see e.g. Beskin Kuznetsova Rafikov 1998). It is this case that can be realized for fast pulsars. Moreover, it has more general interest because the relation (34) may be true also for AGNs. As to the case $`\gamma _{\mathrm{in}}^3\sigma `$ corresponding to ordinary pulsars, the particle energy remains constant ($`\gamma =\gamma _{\mathrm{in}}`$) at any way up to the fast magnetosonic surface (see Bogovalov 1997 for details).
Further, one can put
$`\delta (R_\mathrm{s},\theta )`$ $`=`$ $`0,`$ (35)
$`\epsilon f(R_\mathrm{s},\theta )`$ $`=`$ $`0,`$ (36)
$`\eta ^+(R_\mathrm{s},\theta )\eta ^{}(R_\mathrm{s},\theta )`$ $`=`$ $`0.`$ (37)
These conditions result from the relation $`c𝑬_\mathrm{s}+\mathrm{\Omega }R_\mathrm{s}𝒆_\phi \times 𝑩_\mathrm{s}=0`$ corresponding rigid rotation and perfect conductivity of the surface of a star. Finally, as will be shown in Sect. 3.2, the derivative $`\delta /r`$ actually determines the phase of plasma oscillations only and plays no role in the global structure. Finally, the determination of the electric current and, say, the derivative $`f/r`$ depend on the problem under consideration. Indeed, as is well–known, the cold one–fluid MHD outflow contains two singular surfaces, Alfvénic and fast magnetosonic ones. It means that for the transonic flow two latter functions are to be determined from critical conditions (Heyvaerts 1996). In particular, the longitudinal electric current within this approach is not a free parameter. On the other hand, if the electric current is restricted by some physical reason, the flow cannot pass smoothly through the fast magnetosonic surface. In this case, which can be realized in the magnetosphere of radio pulsars (Beskin et al 1983, Beskin & Malyshkin 1998), near the light surface $`|𝑬|=|𝑩|`$ an effective particle acceleration may take place. Such an acceleration will be considered in Sect. 4.
## 3 The electron–positron outflow
### 3.1 The MHD Limit
In the general case Eqns. (19) – (28) have several integrals. Firstly, Eqns.(21), (25), and (26) result in
$$\zeta \frac{2}{\mathrm{tan}\theta }\delta +\frac{(\lambda 1/2\mathrm{cos}\theta )\gamma ^++(\lambda +1/2\mathrm{cos}\theta )\gamma ^{}}{2\sigma \lambda \mathrm{sin}\theta }=\frac{1}{\sigma \mathrm{sin}\theta }\gamma _{\mathrm{in}}+\frac{l(\theta )}{\mathrm{sin}\theta },$$
(38)
where $`l(\theta )`$ describe the disturbance of the electric current at the star surface by the equation $`I(R,\theta )=I_A\left[\mathrm{sin}^2\theta +l(\theta )\right]`$. Expression (38) corresponds to conservation of the total energy flux along a magnetic field line. Furthermore, Eqns. (25) – (28) together with the boundary conditions (35), (36) result in
$`\delta `$ $`=`$ $`\epsilon f{\displaystyle \frac{1}{4\lambda \sigma }}\gamma ^+\left(1{\displaystyle \frac{\mathrm{\Omega }r\mathrm{sin}\theta }{c}}\xi _\phi ^+\right)+{\displaystyle \frac{1}{4\lambda \sigma }}\gamma _{\mathrm{in}};`$ (39)
$`\delta `$ $`=`$ $`\epsilon f+{\displaystyle \frac{1}{4\lambda \sigma }}\gamma ^{}\left(1{\displaystyle \frac{\mathrm{\Omega }r\mathrm{sin}\theta }{c}}\xi _\phi ^{}\right){\displaystyle \frac{1}{4\lambda \sigma }}\gamma _{\mathrm{in}}.`$ (40)
They correspond to conservation of the $`z`$–component of the angular momentum for both types of particles. It is necessary to stress that the complete nonlinearized system of equations contains no such simple invariants.
As $`\sigma \lambda 1`$, we can neglect in Eqns. (23)–(28) their left-hand sides. In this approximation we have $`\xi ^+=\xi ^{}`$ i.e. $`\gamma ^{}=\gamma ^+=\gamma `$, so that
$`{\displaystyle \frac{1}{r}}{\displaystyle \frac{\delta }{\theta }}+{\displaystyle \frac{\zeta }{r}}{\displaystyle \frac{\mathrm{sin}\theta }{r}}\xi _r+{\displaystyle \frac{c}{\mathrm{\Omega }r^2}}\xi _\phi =0,`$ (41)
$`\epsilon {\displaystyle \frac{c}{\mathrm{\Omega }r\mathrm{sin}\theta }}{\displaystyle \frac{f}{r}}+{\displaystyle \frac{c}{\mathrm{\Omega }r^2}}\xi _\theta =0,`$ (42)
and
$$\gamma \left(1\frac{\mathrm{\Omega }r\mathrm{sin}\theta }{c}\xi _\phi \right)=\gamma _{\mathrm{in}}.$$
(43)
Hence, within this approximation
$`\delta `$ $`=`$ $`\epsilon f,`$ (44)
$`\zeta `$ $`=`$ $`{\displaystyle \frac{2}{\mathrm{tan}\theta }}\epsilon f+{\displaystyle \frac{l(\theta )}{\mathrm{sin}\theta }}{\displaystyle \frac{1}{\sigma \mathrm{sin}\theta }}(\gamma \gamma _{\mathrm{in}}).`$ (45)
Substituting these expressions into (41) and using Eqns. (19)–(22), we obtain the following equation describing the disturbance of the magnetic surfaces
$`\epsilon (1x^2\mathrm{sin}^2\theta ){\displaystyle \frac{^2f}{x^2}}+\epsilon (1x^2\mathrm{sin}^2\theta ){\displaystyle \frac{\mathrm{sin}\theta }{x^2}}{\displaystyle \frac{}{\theta }}\left({\displaystyle \frac{1}{\mathrm{sin}\theta }}{\displaystyle \frac{f}{\theta }}\right)2\epsilon x\mathrm{sin}^2\theta {\displaystyle \frac{f}{x}}2\epsilon \mathrm{sin}\theta \mathrm{cos}\theta {\displaystyle \frac{f}{\theta }}+2\epsilon (3\mathrm{cos}^2\theta 1)f`$ (46)
$`+{\displaystyle \frac{1}{\mathrm{sin}\theta }}{\displaystyle \frac{\mathrm{d}}{\mathrm{d}\theta }}(l\mathrm{sin}^2\theta )2{\displaystyle \frac{\mathrm{cos}\theta }{\sigma }}\left(\gamma \gamma _{\mathrm{in}}\right){\displaystyle \frac{\mathrm{sin}\theta }{\sigma }}{\displaystyle \frac{\gamma }{\theta }}2\lambda \mathrm{sin}^2\theta (\xi _r^+\xi _r^{})+{\displaystyle \frac{2\lambda }{x}}\mathrm{sin}\theta (\xi _\phi ^+\xi _\phi ^{})=0,`$
where $`x=\mathrm{\Omega }r/c`$. One can see that it actually coincides with the one–fluid MHD Eqns.(32), (52) from Beskin et al (1998), but contains the two last additional nonhydrodynamical terms. Nevertheless, as will be shown in the next subsection, at small distances $`rr_\mathrm{F}`$ where $`r_\mathrm{F}`$ is the radius of the fast magnetosonic surface we have
$$\lambda \mathrm{sin}^2\theta (\xi _r^+\xi _r^{})+\frac{\lambda }{x}\mathrm{sin}\theta (\xi _\phi ^+\xi _\phi ^{})0,$$
(47)
so actually there is perfect agreement with the MHD approximation
$`\epsilon (1x^2\mathrm{sin}^2\theta ){\displaystyle \frac{^2f}{x^2}}+\epsilon (1x^2\mathrm{sin}^2\theta ){\displaystyle \frac{\mathrm{sin}\theta }{x^2}}{\displaystyle \frac{}{\theta }}\left({\displaystyle \frac{1}{\mathrm{sin}\theta }}{\displaystyle \frac{f}{\theta }}\right)2\epsilon x\mathrm{sin}^2\theta {\displaystyle \frac{f}{x}}2\epsilon \mathrm{sin}\theta \mathrm{cos}\theta {\displaystyle \frac{f}{\theta }}`$ (48)
$`+2\epsilon (3\mathrm{cos}^2\theta 1)f+{\displaystyle \frac{1}{\mathrm{sin}\theta }}{\displaystyle \frac{\mathrm{d}}{\mathrm{d}\theta }}(l\mathrm{sin}^2\theta )2{\displaystyle \frac{\mathrm{cos}\theta }{\sigma }}\left(\gamma \gamma _{\mathrm{in}}\right){\displaystyle \frac{\mathrm{sin}\theta }{\sigma }}{\displaystyle \frac{\gamma }{\theta }}=0.`$
As was shown earlier (Beskin et al 1998), to pass through the fast magnetosonic surface it’s necessary to have
$$|l|<\sigma ^{4/3}.$$
(49)
Hence, within the fast magnetosonic surface $`rr_\mathrm{F}`$ one can neglect terms containing $`\delta =\epsilon f`$ and $`\zeta `$. Then, relations (41) and (42) result in
$`\gamma (1x\mathrm{sin}\theta \xi _\phi )`$ $`=`$ $`\gamma _{\mathrm{in}},`$ (50)
$`\xi _r`$ $`=`$ $`{\displaystyle \frac{\xi _\phi }{x\mathrm{sin}\theta }},`$ (51)
$`\xi _\theta `$ $`=`$ $`0.`$ (52)
Finally, using the definition
$$\gamma ^2=\frac{1}{2\xi _r\xi _\phi ^2},$$
(53)
we obtain for $`\sigma \gamma _{\mathrm{in}}^3`$ for $`rr_\mathrm{F}`$
$`\gamma ^2`$ $`=`$ $`\gamma _{\mathrm{in}}^2+x^2\mathrm{sin}^2\theta ,`$ (54)
$`\xi _\phi `$ $`=`$ $`{\displaystyle \frac{\sqrt{\gamma _{\mathrm{in}}^2+x^2\mathrm{sin}^2\theta }\gamma _{\mathrm{in}}}{x\mathrm{sin}\theta \sqrt{\gamma _{\mathrm{in}}^2+x^2\mathrm{sin}^2\theta }}}{\displaystyle \frac{1}{x\mathrm{sin}\theta }},`$ (55)
$`\xi _r`$ $`=`$ $`{\displaystyle \frac{\sqrt{\gamma _{\mathrm{in}}^2+x^2\mathrm{sin}^2\theta }\gamma _{\mathrm{in}}}{x^2\mathrm{sin}^2\theta \sqrt{\gamma _{\mathrm{in}}^2+x^2\mathrm{sin}^2\theta }}}{\displaystyle \frac{1}{x^2\mathrm{sin}^2\theta }},`$ (56)
in full agreement with the MHD results.
Next, to determine the position of the fast magnetosonic surface $`r_\mathrm{F}`$, one can analyze the algebraic equations (38) and (41) which give
$$\frac{\delta }{\theta }+\frac{2}{\mathrm{tan}\theta }\delta \frac{1}{\sigma \mathrm{sin}\theta }\gamma \mathrm{sin}\theta \xi _r+\frac{1}{x}\xi _\phi =0.$$
(57)
Using now expressions (43) and (53), one can find
$$2\gamma ^32\sigma \left[K+\frac{1}{2x^2}\right]\gamma ^2+\sigma \mathrm{sin}^2\theta =0,$$
(58)
where
$$K(r,\theta )=2\mathrm{cos}\theta \delta \mathrm{sin}\theta \frac{\delta }{\theta }.$$
(59)
Equation (58) allows us to determine the position of the fast magnetosonic surface and the energy of particles. Indeed, determining the derivative $`r\gamma /r`$, one can obtain
$$r\frac{\gamma }{r}=\frac{\gamma \sigma \left(rK/rx^2\right)}{3\gamma \sigma \left(2K+x^2\right)}.$$
(60)
As the fast magnetosonic surface is the $`X`$–point, both the nominator and denominator are to be equal to zero here. As a result, evaluating $`rK/r`$ as $`K`$, we obtain
$`\delta `$ $``$ $`\sigma ^{2/3};`$ (61)
$`r_\mathrm{F}`$ $``$ $`\sigma ^{1/3}R_\mathrm{L};`$ (62)
$`\gamma (r_\mathrm{F})`$ $`=`$ $`\sigma ^{1/3}\mathrm{sin}^{2/3}\theta ,`$ (63)
where the last expression is exact. These equations coincide with those obtained within the MHD consideration. It is the self–consistent analysis when $`\delta =\epsilon f`$, and hence $`K`$ depends on the radius $`r`$ that results in the finite value for the fast magnetosonic radius $`r_\mathrm{F}`$. On the other hand, in a given monopole magnetic field, when $`\epsilon f`$ does not depend on the radius, the critical conditions result in $`r_\mathrm{F}\mathrm{}`$ for a cold outflow (Michel, 1969, Li et al 1992).
Near the fast magnetosonic surface $`r\sigma ^{1/3}R_\mathrm{L}`$ the MHD solution gives
$`\gamma `$ $``$ $`\sigma ^{1/3},`$ (64)
$`\epsilon f`$ $``$ $`\sigma ^{2/3}.`$ (65)
Hence, Eqns. (53), (55), and (56) result in
$`\xi _r`$ $``$ $`\sigma ^{2/3},`$ (66)
$`\xi _\theta `$ $``$ $`\sigma ^{2/3},`$ (67)
$`\xi _\phi `$ $``$ $`\sigma ^{1/3}.`$ (68)
As we see, the $`\theta `$–component of the velocity plays no role in the determination of the $`\gamma `$.
However, analyzing the left-hand sides of the Eqns. (23)–(28) one can evaluate the additional (nonhydrodynamic) variations of the velocity components
$`\mathrm{\Delta }\xi _r^\pm `$ $``$ $`\lambda ^1\sigma ^{4/3},`$ (69)
$`\mathrm{\Delta }\xi _\theta ^\pm `$ $``$ $`\lambda ^1\sigma ^{2/3},`$ (70)
$`\mathrm{\Delta }\xi _\phi ^\pm `$ $``$ $`\lambda ^1\sigma ^1.`$ (71)
Hence, for nonhydrodynamic velocities $`\mathrm{\Delta }\xi _r^\pm \xi _r`$ and $`\mathrm{\Delta }\xi _\phi ^\pm \xi _\phi `$ to be small, it is necessary to have a large magnetization parameter $`\sigma 1`$ only. On the other hand, $`\mathrm{\Delta }\xi _\theta ^\pm /\xi _\theta \lambda ^1`$. In other words, for a highly magnetized plasma $`\sigma 1`$ even outside the fast magnetosonic surface the velocity components (and, hence, the particle energy) can be considered hydrodynamically. The difference $`\lambda ^1`$ appears in the $`\theta `$ component only, but it does not affect the particle energy. Finally, one can obtain from (39), (40) that
$$\frac{\delta \epsilon f}{\epsilon f}\lambda ^2\sigma ^{2/3}.$$
(72)
To put it differently, at large distances the nonhydrodynamical terms are much smaller than hydrodynamical ones.
As a result, at large distances where, according to (39)–(40), one can neglect the toroidal component $`\xi _\phi `$, we obtain
$`\delta `$ $`=`$ $`\epsilon f,`$ (73)
$`\zeta `$ $`=`$ $`{\displaystyle \frac{2}{\mathrm{tan}\theta }}\delta \sigma ^1{\displaystyle \frac{1}{\mathrm{sin}\theta }}\gamma .`$ (74)
On the other hand, Eqn. (23) gives
$$\zeta =\frac{\delta }{\theta }+\mathrm{sin}\theta \xi _r.$$
(75)
Together with (21) one can obtain for $`rr_\mathrm{F}`$
$$\gamma =\sigma \left(2\mathrm{cos}\theta \epsilon f\epsilon \mathrm{sin}\theta \frac{f}{\theta }\right),$$
(76)
which coincides with the MHD solution. Finally, using Eqns. (19), (20), and neglecting the nonhydrodynamic term $`4\lambda (\xi _r^+\xi _r^{})`$, one can find
$$\epsilon \frac{}{r}\left(r^2\frac{f}{r}\right)4\mathrm{cos}\theta \xi _r\mathrm{sin}\theta \frac{}{\theta }\xi _r+\frac{1}{x\mathrm{sin}\theta }\frac{}{\theta }(\xi _\phi \mathrm{sin}\theta )=0.$$
(77)
Together with (76) this equation in the limit $`rr_\mathrm{F}`$ coincides with the asymptotic version of the trans–field equation (Tomimatsu 1994, Beskin et al 1998)
$$\epsilon \frac{^2f}{r^2}+2\epsilon r\frac{f}{r}\mathrm{sin}\theta \frac{D+1}{D}\frac{g}{\theta }=0,$$
(78)
where $`g(\theta )=K(\theta )/\mathrm{sin}^2\theta `$, and
$$D+1=\frac{1}{\sigma ^2\mathrm{sin}^4\theta }g^3(\theta )1.$$
(79)
In this limit, none of the terms containing $`\xi _r^\pm `$ and $`\xi _\phi `$ plays role in the asymptotic trans–field equation. Hence, it is not necessary to consider the effect of the nonhydrodynamical term $`4\lambda (\xi _r^+\xi _r^{})`$ either.
### 3.2 Plasma Oscillations
In the intermediate region $`rr_\mathrm{F}`$ Eqn. (77) cannot be used. The point is that in the limit $`\lambda 1`$ the important role in Eqns. (19) and (22) is played by the nonhydrodynamic terms (47) corresponding to different velocities of two components. As a result, the full version of Eqn. (77) has the form
$$\frac{}{r}\left(r^2\frac{\delta }{r}\right)4\mathrm{cos}\theta \xi _r\mathrm{sin}\theta \frac{\xi _r}{\theta }+\frac{1}{x\mathrm{sin}\theta }\frac{}{\theta }(\xi _\phi \mathrm{sin}\theta )+2\lambda (\xi _r^+\xi _r^{})=0.$$
(80)
Indeed, one can see from equations (19) and (20) that near the origin $`x=R_\mathrm{s}`$ in the case $`\gamma _{\mathrm{in}}^+=\gamma _{\mathrm{in}}^{}`$ (and for the small variation of the current $`\zeta \sigma ^{4/3}`$ which is necessary, as was already stressed, to pass through a fast magnetosonic surface) the density variation on the surface is large enough: $`(\eta ^+\eta ^{})\gamma _{\mathrm{in}}^2\zeta `$. Hence, the derivative $`^2\delta /r^2`$ here is of the order of $`\gamma _{\mathrm{in}}^2`$. On the other hand, according to (22), the derivative $`\epsilon ^2f/r^2`$ is $`x^2`$ times smaller. This means that in the two–component system the longitudinal electric field is to appear resulting in a redistribution of the particle energy. Clearly, such a redistribution is impossible for the charge–separated outflow. In other words, for a finite particle energy a one–component plasma cannot maintain simultaneously both the Goldreich charge and Goldreich current density (4). In a two–component system with $`\lambda 1`$ it can be realized by a small redistribution of particle energy (Ruderman & Sutherland 1975, Arons & Scharlemann 1989).
For simplicity, let us consider only small distances $`x1`$. In this case one can neglect the changes of the magnetic surfaces. Using now (25) and (26), we have
$`\gamma ^+`$ $`=`$ $`\gamma _{\mathrm{in}}4\lambda \sigma \delta ;`$ (81)
$`\gamma ^{}`$ $`=`$ $`\gamma _{\mathrm{in}}+4\lambda \sigma \delta .`$ (82)
Finally, taking into account that $`\xi _\theta `$ and $`\xi _\phi `$ are small here, one can obtain from (20)
$$r^2\frac{^2\delta }{r^2}+2r\frac{\delta }{r}+\frac{1}{\mathrm{sin}\theta }\frac{}{\theta }\left(\mathrm{sin}\theta \frac{\delta }{\theta }\right)+A\delta =\frac{\mathrm{cos}\theta }{\gamma _{\mathrm{in}}^2},$$
(83)
where
$$A=16\frac{\lambda ^2\sigma }{\gamma _{\mathrm{in}}^3}1.$$
(84)
Eqn. (83) has a solution
$$\delta =\delta _0+r^{1/2}\left[C_1\mathrm{sin}(\mu \mathrm{ln}r)+C_2\mathrm{cos}(\mu \mathrm{ln}r)\right]\mathrm{cos}\theta ,$$
(85)
where
$$\delta _0\frac{\gamma _{\mathrm{in}}\mathrm{cos}\theta }{16\lambda ^2\sigma },$$
(86)
and $`\mu \sqrt{A}`$. As we see, Eqn. (85) describes plasma oscillations similar to those considered by Shibata (1997) for charge–separated flow. The decrease of oscillations results from a more accurate consideration of the Laplace operator in a 3D space.
One can easily check that the additional potential $`\delta _0`$ is small, and it is not necessary to add it in (38) and (39)–(40). Moreover, the nonhydrodynamic disturbance $`\mathrm{\Delta }\xi _r`$ (as well as $`\mathrm{\Delta }\gamma `$) is also small, $`\mathrm{\Delta }\xi _r/\xi _r\lambda ^1`$. Hence, as was already stressed, the boundary condition $`\delta /r`$ (determining together with (35) the coefficients $`C_1`$ and $`C_2`$) does not affect the general structure of the flow. On the other hand, the presence of an additional electric potential $`\delta _0`$ results in a full compensation of the last term in (19)
$$2\lambda (\xi _r^+\xi _r^{})\mathrm{cos}\theta \xi _r0.$$
(87)
Next, as $`\epsilon f\sigma ^{2/3}`$ for $`rr_\mathrm{F}`$, a similar expression can be written for the $`\phi `$–components as well
$$2\lambda (\xi _\phi ^+\xi _\phi ^{})\mathrm{cos}\theta \xi _\phi 0.$$
(88)
Expressions (87) – (88) must hold for the whole region $`r<r_\mathrm{F}`$. In this case, the final version of Eqn. (80) in the internal region $`rr_\mathrm{F}`$ can be rewritten as
$$\frac{}{r}\left(r^2\frac{\delta }{r}\right)2\mathrm{cos}\theta \xi _r\mathrm{sin}\theta \frac{}{\theta }\xi _r+\frac{1}{x\mathrm{sin}\theta }\frac{}{\theta }(\xi _\phi \mathrm{sin}\theta )=0.$$
(89)
As $`\delta \epsilon f\sigma ^{2/3}`$ for $`rr_\mathrm{F}`$, and $`\xi _r\gamma _0^2\delta `$, the first term in (89) can be omitted. As a result, the solution of Eqn. (89) coincides exactly with the MHD expression, i.e. $`\gamma ^2=\gamma _{\mathrm{in}}^2+x^2\mathrm{sin}^2\theta `$ (54). Finally, using (87), (88), and (55)–(56), one can easily check that the nonhydrodynamical terms (47) in the trans–field equation (48) do actually vanish.
## 4 The Boundary Layer
Let us now consider the case when the longitudinal electric current $`I(R,\theta )`$ in the magnetosphere of radio pulsars is too small (i.e. the disturbance $`l(\theta )`$ is too large) for the flow to pass smoothly through the fast magnetosonic surface. First of all, it can be realized when the electric current is much smaller than the Goldreich one. This possibility was already discussed within the Ruderman–Sutherland model of the internal gap (Beskin et al 1983, Beskin & Malyshkin 1998). But it may take place in the Arons model (Arons & Scharlemann 1979) as well. Indeed, within this model the electric current is determined by the gap structure. Hence, in general case this current does not correspond to the critical condition at the fast magnetosonic surface. In particular, it may be smaller than the critical current (of course, the separate consideration is necessary to check this statement).
For simplicity let us consider the case $`l(\theta )=h\mathrm{sin}^2\theta `$. Neglecting now the last terms $`\sigma ^1`$ in the trans–field equation (48), we obtain
$`\epsilon (1x^2\mathrm{sin}^2\theta ){\displaystyle \frac{^2f}{x^2}}+\epsilon (1x^2\mathrm{sin}^2\theta ){\displaystyle \frac{\mathrm{sin}\theta }{x^2}}{\displaystyle \frac{}{\theta }}\left({\displaystyle \frac{1}{\mathrm{sin}\theta }}{\displaystyle \frac{f}{\theta }}\right)`$
$`2\epsilon x\mathrm{sin}^2\theta {\displaystyle \frac{f}{x}}2\epsilon \mathrm{sin}\theta \mathrm{cos}\theta {\displaystyle \frac{f}{\theta }}+2\epsilon (3\mathrm{cos}^2\theta 1)f+4h\mathrm{sin}^2\theta \mathrm{cos}\theta =0,`$ (90)
which actually coincides with the force–free equation (Beskin et al 1998). This equation has an exact analytical solution
$$\epsilon f=hx^2\mathrm{sin}^2\theta \mathrm{cos}\theta .$$
(91)
For $`h<0`$ (when the electric current is smaller than the Goldreich one) this solution results in the appearance of the light surface $`|𝑬|=|𝑩|`$ at the finite distance
$$\varpi _c=\frac{R_\mathrm{L}}{(2|h|)^{1/4}}.$$
(92)
As we see, for $`l(\theta )=h\mathrm{sin}^2\theta `$ this surface has the form of a cylinder. It is important that the disturbance of magnetic surfaces $`\epsilon f(|h|)^{1/2}`$ remains small here.
Comparing now the position of the light surface (92) with that of the fast magnetosonic surface (62), one can find that the light surface is located inside the fast magnetosonic one if
$$\sigma ^{4/3}|h|1,$$
(93)
which is opposite to (49). One can check that the condition (93) just allows to neglect the non force–free term in Eqn. (48).
Using now the solution (91) and the MHD condition $`\delta =\epsilon f`$, one can find from (58)
$$2\gamma ^32\sigma \left(hx^2\mathrm{sin}^4\theta +\frac{1}{2x^2}\right)\gamma ^2+\sigma \mathrm{sin}^2\theta =0.$$
(94)
This equation shows that near the force–free boundary $`x_{ff}=(2|h|)^{1/4}`$ (92)
$$\gamma =\sigma ^{1/3}\mathrm{sin}^{2/3}\theta \frac{2|2h|^{3/8}}{\sqrt{3}}\sigma ^{1/3}\mathrm{sin}^{4/3}\theta \sqrt{x_0x\mathrm{sin}\theta },$$
(95)
where
$$x_0=\frac{1}{(2|h|)^{1/4}}\left[1\frac{3}{4(2|h|)^{1/2}}\frac{1}{(\sigma \mathrm{sin}^2\theta )^{2/3}}\right],$$
(96)
(see Fig. 1). Hence, the real solution is absent for $`x\mathrm{sin}\theta >x_0`$. Here $`\gamma (x_0)=\sigma ^{1/3}\mathrm{sin}^{2/3}\theta `$, the condition $`\sigma ^{4/3}|h|`$ resulting in $`\varpi _c<r_\mathrm{F}`$, and the last term in (96) being small.
Although the energy of particles at the limiting point is finite, the derivative $`d\gamma /d\varpi `$ moves to infinity. Hence, near the light surface the left-hand sides in the Eqns. (23) – (28) are to be taken into consideration. Since in our case the light surface has the form of a cylinder, one can move to derivatives perpendicular to the boundary layer only by
$`/r\mathrm{sin}\theta /\varpi ;`$ (97)
$`/\theta \varpi \mathrm{cos}\theta /\varpi .`$ (98)
As a result, $`\zeta `$ can be eliminated from (19) and (21). Together with (20) they give the equation for $`\delta `$ (see (104)). Next, the invariants (39 ) and (40) can be used to define $`\xi _\phi ^\pm `$:
$`\xi _\phi ^+`$ $`=`$ $`{\displaystyle \frac{1}{x\mathrm{sin}\theta }}\left[1+{\displaystyle \frac{4\lambda \sigma (\delta \epsilon f)}{\gamma ^+}}\right];`$ (99)
$`\xi _\phi ^{}`$ $`=`$ $`{\displaystyle \frac{1}{x\mathrm{sin}\theta }}\left[1{\displaystyle \frac{4\lambda \sigma (\delta \epsilon f)}{\gamma ^{}}}\right].`$ (100)
Furthermore, one can define
$`2\xi _r^+`$ $`=`$ $`{\displaystyle \frac{1}{(\gamma ^+)^2}}+(\xi _\phi ^+)^2+(\xi _\theta ^+)^2;`$ (101)
$`2\xi _r^{}`$ $`=`$ $`{\displaystyle \frac{1}{(\gamma ^{})^2}}+(\xi _\phi ^{})^2+(\xi _\theta ^{})^2.`$ (102)
As to the energy integral (38), it determines the variation of the current $`\zeta `$. Now it can be rewritten as
$$\zeta =\frac{2}{\mathrm{tan}\theta }\delta \frac{(\gamma ^++\gamma ^{})}{2\sigma \mathrm{sin}\theta }.$$
(103)
Finally, Eqns. (19) – (28) look like
$`\varpi _c^2{\displaystyle \frac{d^2\delta }{d\varpi ^2}}=2\mathrm{sin}\theta \mathrm{cos}\theta \left[\left(\lambda {\displaystyle \frac{1}{2}}\mathrm{cos}\theta \right)\xi _\theta ^+\left(\lambda +{\displaystyle \frac{1}{2}}\mathrm{cos}\theta \right)\xi _\theta ^{}\right]`$
$`2\mathrm{sin}^2\theta \left[\left(\lambda {\displaystyle \frac{1}{2}}\mathrm{cos}\theta \right)\xi _r^+\left(\lambda +{\displaystyle \frac{1}{2}}\mathrm{cos}\theta \right)\xi _r^{}\right],`$ (104)
$`\varpi _cR_\mathrm{L}\epsilon {\displaystyle \frac{d^2f}{d\varpi ^2}}=2\mathrm{sin}^2\theta \left[\left(\lambda {\displaystyle \frac{1}{2}}\mathrm{cos}\theta \right)\xi _\phi ^+\left(\lambda +{\displaystyle \frac{1}{2}}\mathrm{cos}\theta \right)\xi _\phi ^{}\right],`$ (105)
$`\varpi _c{\displaystyle \frac{d}{d\varpi }}\left(\xi _\theta ^+\gamma ^+\right)=4\lambda \sigma \left({\displaystyle \frac{\gamma ^++\gamma ^{}}{2\sigma \mathrm{sin}\theta }}\varpi _c{\displaystyle \frac{\mathrm{cos}\theta }{\mathrm{sin}\theta }}{\displaystyle \frac{d\delta }{d\varpi }}\mathrm{sin}\theta \xi _r^++{\displaystyle \frac{\mathrm{sin}\theta }{x_0}}\xi _\phi ^+\right),`$ (106)
$`\varpi _c{\displaystyle \frac{d}{d\varpi }}\left(\xi _\theta ^{}\gamma ^{}\right)=4\lambda \sigma \left({\displaystyle \frac{\gamma ^++\gamma ^{}}{2\sigma \mathrm{sin}\theta }}\varpi _c{\displaystyle \frac{\mathrm{cos}\theta }{\mathrm{sin}\theta }}{\displaystyle \frac{d\delta }{d\varpi }}\mathrm{sin}\theta \xi _r^{}+{\displaystyle \frac{\mathrm{sin}\theta }{x_0}}\xi _\phi ^{}\right),`$ (107)
$`\varpi _c{\displaystyle \frac{d}{d\varpi }}\gamma ^+=4\lambda \sigma \left(\varpi _c{\displaystyle \frac{d\delta }{d\varpi }}\mathrm{sin}\theta \xi _\theta ^+\right),`$ (108)
$`\varpi _c{\displaystyle \frac{d}{d\varpi }}\gamma ^{}=4\lambda \sigma \left(\varpi _c{\displaystyle \frac{d\delta }{d\varpi }}\mathrm{sin}\theta \xi _\theta ^{}\right),`$ (109)
where we neglected the terms $`\delta /r`$ in (106) and (107).
Comparing the leading terms, we have inside the layer $`\mathrm{\Delta }\varpi /R_\mathrm{L}\lambda ^1`$
$`\gamma ^\pm `$ $``$ $`h_c^{1/2}\sigma ,`$ (110)
$`\xi _\theta ^\pm `$ $``$ $`h_c^{1/4},`$ (111)
$`\xi _r^\pm `$ $``$ $`h_c^{1/2},`$ (112)
$`\mathrm{\Delta }\delta `$ $``$ $`h_c^{3/4}/\lambda ,`$ (113)
where $`h_c=|h|`$. Then the leading terms in (99) – (103) for $`\mathrm{\Delta }\varpi >\lambda ^1R_\mathrm{L}`$ are
$`\xi _\phi ^+`$ $`=`$ $`{\displaystyle \frac{1}{x\mathrm{sin}\theta }}{\displaystyle \frac{1}{x_0}},`$ (114)
$`\xi _\phi ^{}`$ $`=`$ $`{\displaystyle \frac{1}{x\mathrm{sin}\theta }}{\displaystyle \frac{1}{x_0}},`$ (115)
$`2\xi _r^+`$ $`=`$ $`(\xi _\phi ^+)^2+(\xi _\theta ^+)^2,`$ (116)
$`2\xi _r^{}`$ $`=`$ $`(\xi _\phi ^{})^2+(\xi _\theta ^{})^2,`$ (117)
$$\zeta =\frac{(\gamma ^++\gamma ^{})}{2\sigma \mathrm{sin}\theta },$$
(118)
where $`x_0=\varpi _c/R_\mathrm{L}=(2|h|)^{1/4}`$. Hence, one can totally neglect $`\epsilon f`$ and $`\delta `$ in (99) – (100), so Eqns. (104) – (109) in the region $`\mathrm{\Delta }\varpi >\lambda ^1R_\mathrm{L}`$ can be rewritten as
$`\varpi _c^2{\displaystyle \frac{d^2\delta }{d\varpi ^2}}=2\lambda \mathrm{sin}\theta \mathrm{cos}\theta (\xi _\theta ^+\xi _\theta ^{}),`$ (119)
$`\varpi _c{\displaystyle \frac{d}{d\varpi }}\left(\xi _\theta ^+\gamma ^+\right)=4\lambda \sigma \left({\displaystyle \frac{\gamma ^++\gamma ^{}}{2\sigma \mathrm{sin}\theta }}\varpi _c{\displaystyle \frac{\mathrm{cos}\theta }{\mathrm{sin}\theta }}{\displaystyle \frac{d\delta }{d\varpi }}\mathrm{sin}\theta \xi _r^+\right),`$ (120)
$`\varpi _c{\displaystyle \frac{d}{d\varpi }}\left(\xi _\theta ^{}\gamma ^{}\right)=4\lambda \sigma \left({\displaystyle \frac{\gamma ^++\gamma ^{}}{2\sigma \mathrm{sin}\theta }}\varpi _c{\displaystyle \frac{\mathrm{cos}\theta }{\mathrm{sin}\theta }}{\displaystyle \frac{d\delta }{d\varpi }}\mathrm{sin}\theta \xi _r^{}\right),`$ (121)
$`\varpi _c{\displaystyle \frac{d}{d\varpi }}\left(\gamma ^+\right)=4\lambda \sigma \mathrm{sin}\theta \xi _\theta ^+,`$ (122)
$`\varpi _c{\displaystyle \frac{d}{d\varpi }}\left(\gamma ^{}\right)=4\lambda \sigma \mathrm{sin}\theta \xi _\theta ^{},`$ (123)
with all the terms in the right–hand sides of (120) and (121) being of the same order of magnitude.
As a result, the nonlinear equations (119) – (123) and (105) give the following simple asymptotic solution
$`\gamma ^\pm `$ $`=`$ $`4\mathrm{sin}^2\theta \sigma (\lambda l)^2,`$ (124)
$`\xi _\theta ^\pm `$ $`=`$ $`2\mathrm{sin}\theta \lambda l,`$ (125)
$`\mathrm{\Delta }\delta `$ $`=`$ $`{\displaystyle \frac{4}{3}}\mathrm{sin}^2\theta \mathrm{cos}\theta \lambda ^1(\lambda l)^3,`$ (126)
$`\mathrm{\Delta }(\epsilon f)`$ $`=`$ $`\mathrm{sin}^2\theta \mathrm{cos}\theta \lambda ^2(\lambda l)^2,`$ (127)
$`\zeta `$ $`=`$ $`4\mathrm{sin}\theta (\lambda l)^2,`$ (128)
where now $`l=\mathrm{\Delta }\varpi /\varpi _c`$. It is important that the last expressions are correct for arbitrary relations between $`\gamma _{\mathrm{in}}^3`$ and $`\sigma `$. As we can see, in the narrow layer $`\mathrm{\Delta }\varpi =\varpi _c/\lambda `$ the particle energy increases up to the value $`\sigma `$ which corresponds to the full conversion of the electromagnetic energy to the energy of particles. For this reason we have here $`|\zeta |1`$, which just means the diminishing of the toroidal magnetic field determining the flux of the electromagnetic energy. On the other hand, the variation of the electric potential remains small $`\delta \lambda ^1`$, to say nothing about the variation of the magnetic surfaces $`\mathrm{\Delta }\epsilon f\lambda ^2`$. These results coincide exactly with our previous evaluations (Beskin et al 1983) allowing us to neglect variations of the electric potential and the poloidal magnetic structure in the 1D cylindrical case.
The latter result has a simple physical explanation. Indeed, the diminishing of the toroidal magnetic field is connected with the $`\theta `$–component of the electric current which is produced by all the particles moving in opposite directions. On the other hand, the change of the electric potential depends on the small difference between the electron and positron densities. As a result, according to (127) and (128), the change of the toroidal magnetic field is just $`\lambda `$ times larger than the change of the electric potential. Unfortunately, it is impossible to consider this region more thoroughly because for $`\lambda l1`$ we have $`\xi _\theta ^\pm 1`$ and $`\xi _r^\pm 1`$, i.e. the linear approximation (19) – (28) itself becomes incorrect.
It is necessary to stress as well that we do not include into consideration the radiation reaction force
$$F_x^{(\mathrm{rad})}=\frac{2}{3}\frac{e^4}{m^2c^4}\gamma ^2\left[(E_yB_z)^2+(E_zB_y)^2\right],$$
(129)
which can be important for large enough particle energy. Comparing (129) with appropriate terms in (120) – (123) one can conclude that the radiation force can be neglected for $`\sigma <\sigma _{\mathrm{cr}}`$, where
$$\sigma _{\mathrm{cr}}=\left(\frac{c}{\lambda r_e\mathrm{\Omega }}\right)^{1/3}10^6,$$
(130)
and $`r_e=e^2/mc^2`$ – classical electron radius. This relation can be rewritten in the form
$$\frac{\mathrm{\Omega }R}{c}<3\times 10^3B_{12}^{3/7}\lambda _4^{2/7}$$
(131)
which gives
$$P>0.06B_{12}^{3/7}\lambda _4^{2/7}s.$$
(132)
Hence, for most radio pulsars the radiation force indeed can be neglected. As to the pulsars with $`\sigma >\sigma _{\mathrm{cr}}`$, it is clear that for $`\gamma >\sigma _{\mathrm{cr}}`$ the radiation force becomes larger than the electromagnetic one and strongly inhibits any further acceleration. As a result, we can evaluate the maximum gamma–factor which can be reached during the acceleration as
$$\gamma _{\mathrm{max}}\sigma _{\mathrm{cr}}10^6.$$
(133)
## 5 Discussion
Thus, on a simple example it was demonstrated that for real physical parameters of the magnetosphere of radio pulsars ($`\sigma 1`$ and $`\lambda 1`$) the one–fluid MHD approximation remains true in the whole region within the light surface $`|𝑬|=|𝑩|`$. On the other hand, it was shown that in a more realistic 2D case the main properties of the boundary layer near the light surface existing for small enough longitudinal currents $`I<I_{GJ}`$ (effective energy transformation from electromagnetic field to particles, current closure in this region, smallness of the disturbance of electric potential and poloidal magnetic field) remain the same as in the 1D case considered previously (Beskin et al 1983).
It is necessary to stress the main astrophysical consequences of our results. First of all, the presence of such a boundary layer explains the effective energy transformation of electromagnetic energy into the energy of particles. As was already stressed, now the existence of such an acceleration is confirmed by observations of close binaries containing radio pulsars (as to the particle acceleration far from a neutron star, see e.g. Kennel & Coroniti 1984, Hoshino et al 1992, Gallant & Arons 1994). Simultaneously, it allows us to understand the current closure in the pulsar magnetosphere. Finally, particle acceleration results in the additional mechanism of high–energy radiation from the boundary of the magnetosphere (for more details see Beskin et al 1993).
Nevertheless, it is clear that the results obtained do not solve the whole pulsar wind problem. Indeed, as in the cylindrical case, it is impossible to describe the particle motion outside the light surface. The point is that, as one can see directly from Eqn. (126), for a complete conversion of electromagnetic energy into the energy of particles it is enough for them to pass only $`\lambda ^1`$ of the total potential drop between pulsar magnetosphere and infinity. It means that the electron–positron wind propagating to infinity has to pass the potential drop which is much larger than their energy. It is possible only in the presence of electromagnetic waves even in an axisymmetric magnetosphere which is stationary near the origin. Clearly, such a flow cannot be considered even within the two–fluid approximation. In our opinion, it is only a numerical consideration that can solve the problem completely and determine, in particular, the energy spectrum of particles and the structure of the pulsar wind. Unfortunately, up to now such numerical calculations are absent.
## Acknowledgments
The authors are grateful to I. Okamoto and H. Sol for fruitful discussions. VSB thanks National Astronomical Observatory, Japan for hospitality. This work was supported by INTAS Grant 96–154 and by Russian Foundation for Basic Research (Grant 96–02–18203). |
warning/0002/nucl-th0002024.html | ar5iv | text | # 1 Introduction
## 1 Introduction
The connection between hot dense hadronic matter and a plasma of quarks and gluons is receiving increased attention with the advent of the relativistic heavy-ion collider (RHIC) at Brookhaven to complement previous investigations at the CERN SPS <sup>?</sup>. The plasma is expected to reveal itself through modified properties of hadronic reactions and their products. The di-lepton spectra has been of considerable interest <sup>?</sup> as a relatively clean signal of how vector meson correlations and their decay channels and widths might be influenced by a hot and dense environment. The question of how a vector meson strong decay, such as $`\rho \pi \pi `$, might respond as the temperature $`T`$ or chemical potential $`\mu `$ crosses the critical phase boundary for chiral restoration or quark deconfinement requires that this process be studied at the quark-gluon level. To this end it is desirable to be able to describe quark deconfinement and chiral restoration at finite $`T`$ and $`\mu `$ in a manner that can be extended to a variety of hadronic observables in a rapid and transparent way. The present model provides a simple framework for such investigations.
At $`T=\mu =0`$ significant progress has been made within a continuum approach to modeling non-perturbative QCD based on the truncated Dyson-Schwinger equations (DSEs) <sup>?,?,?</sup>. An attractive feature of this approach is that dynamical chiral symmetry breaking and quark confinement can be embodied in the infrared structure of the dressed gluon 2-point function which is constrained by a few chiral meson observables. Recent works have employed the ladder-rainbow truncation of the coupled quark DSE and $`\overline{q}q`$ Bethe-Salpeter equation (BSE) to produce successful descriptions of masses and decay constants of the light pseudoscalars $`\pi `$ and $`K`$ <sup>?</sup> and the vectors $`\rho `$, $`\varphi `$, and $`K^{}`$ <sup>?</sup>. These works have incorporated the one-loop renormalization group evolution of scale characteristic of QCD. With a few exceptions <sup>?</sup>, applications to other hadron observables, including electromagnetic form factors and coupling constants such as $`g_{\rho \pi \pi }`$, have usually required the use of simpler models <sup>?,?,?,?</sup>.
The finite $`T,\mu `$ extension of realistic DSE/BSE models receives extra complications due to the breaking of $`O(4)`$ symmetry and the dynamical coupling between Matsubara modes <sup>?,?,?</sup>. These difficulties are compounded in the subsequent generation of hadronic observables via straightforward adaptation of the approach <sup>?</sup> found to be successful at $`T=\mu =0`$. Studies of hadronic observables at finite $`T,\mu `$ in DSE/BSE models have been restricted to simplifications such as extensions of the infrared-dominant (ID) model <sup>?</sup> in which the effective gluon propagator is restricted to an integrable singularity at zero momentum. Nevertheless such a model has proved capable of yielding qualitatively useful information <sup>?,?,?</sup>.
In this work we explore a separable Ansatz that can implement the essential qualitative features of DSE/BSE models at finite temperature. Separable representations have previously been found capable of an efficient modeling of the effective $`\overline{q}q`$ interaction in the infrared domain for $`T=0`$ meson observables <sup>?,?</sup>. A previous implementation at $`T>0`$ employed an instantaneous separable interaction without a quark confinement mechanism <sup>?</sup>. Here the approach is covariant and the simple separable interaction absolutely confines quarks at $`T=0`$. The few parameters are fixed by $`\pi `$ and $`\rho /\omega `$ properties. The approach is a simplification of one developed earlier <sup>?</sup> that was found to be quite successful for the light meson spectrum at $`T=0`$. The confining mechanism is an infrared enhancement in the quark-quark interaction that is strong enough to remove the possibility of a mass shell pole in the quark propagator for real $`p^2`$. In the simple implementation used here, it is particularly transparent that sufficiently high temperature will necessarily restore a quark mass-shell pole and there will be a deconfinement transition. This model also implements low temperature dynamical chiral symmetry breaking and it preserves the Goldstone theorem in that the generated $`\pi `$ is massless in the chiral limit. The solutions of the BSE for the $`\pi `$ and $`\rho `$ modes are particularly simple and are used to study the $`T`$-dependence of the meson masses and decay constants in the presence of both deconfinement and chiral restoration mechanisms. Some preliminary results have been previously discussed for $`T>0`$ <sup>?</sup> and for $`T,\mu >0`$ <sup>?</sup>.
In Sec. II the $`T=0`$ separable model is introduced, the DSE solutions for the dressed quark propagator are developed and the quark confinement property is described. The BSE solutions for $`\pi `$ and $`\rho `$ are also treated there and the corresponding decay constants are defined. The extension to $`T>0`$ is considered in Sec. III and the chiral symmetry restoration and deconfinement phenomena are discussed. The spatial $`\pi `$ and $`\rho `$ modes are considered there. The Gell-Mann–Oakes–Renner (GMOR) relation and its generalization are used to evaluate the performance of this model in respecting chiral symmetry constraints. The widths of the transverse $`\rho `$ from $`\rho e^+e^{}`$ and $`\rho \pi \pi `$ are also considered in Sec. III. The high $`T`$ behavior of the model is considered in Sec. IV by comparison of masses with results from lattice simulations. A discussion follows in Sec. V. An Appendix details the high $`T`$ behavior of the ID model as a reference.
## 2 Confining separable Dyson-Schwinger equation model
Mesons can be described as $`q\overline{q}`$ bound states using the Bethe-Salpeter equation. In the ladder truncation, this equation readsWe use a Euclidean space formulation, with $`\{\gamma _\mu ,\gamma _\nu \}=2\delta _{\mu \nu }`$, $`\gamma _\mu ^{}=\gamma _\mu `$ and $`ab=_{i=1}^4a_ib_i`$.
$$\lambda (P^2)\mathrm{\Gamma }(p,P)=\frac{4}{3}\frac{d^4q}{(2\pi )^4}D_{\mu \nu }^{\mathrm{eff}}(pq)\gamma _\mu S(q_+)\mathrm{\Gamma }(q,P)S(q_{})\gamma _\nu ,$$
(1)
where $`P`$ is the total momentum, $`q_\pm =q\pm P/2`$, and $`D_{\mu \nu }^{\mathrm{eff}}(k)`$ an “effective gluon propagator”. The meson mass is identified from $`\lambda (P^2=M^2)=1`$. In conjunction with the rainbow truncation for the quark DSE
$`S(p)^1`$ $`=`$ $`Z_2i\gamma p+Z_2m_0+{\displaystyle \frac{4}{3}}{\displaystyle \frac{d^4q}{(2\pi )^4}g^2D_{\mu \nu }^{\mathrm{eff}}(pq)\gamma _\mu S(q)\gamma _\nu }.`$ (2)
this equation forms the basis for the DSE approach to meson physics <sup>?,?</sup>.
Recent studies have employed $`D_{\mu \nu }^{\mathrm{eff}}(k)=𝒢(k^2)D_{\mu \nu }^{\mathrm{free}}(k)`$ where $`D_{\mu \nu }^{\mathrm{free}}(k)`$ is the free gluon propagator in Landau gauge. The effective coupling $`𝒢(k^2)`$ is given by the one-loop perturbative form of the QCD running coupling in the ultraviolet while the phenomenological infrared form is chosen to reproduce the pion and kaon masses and decay constants <sup>?</sup>. This model can succesfully describe the vector meson masses and dominant decays <sup>?</sup>, as well as the pion and kaon electromagnetic form factors <sup>?</sup>, without new parameters.
The direct extension of such an approach to accommodate finite temperature <sup>?</sup> and baryon density entails a significant increase in complexity. In the Matsubara formalism, the number of coupled equations represented by Eqs. (1) and (2) scales up with the number of fermion Matsubara modes included. For studies near and above the transition, $`T100`$ MeV, about $`10`$ such modes appear adequate <sup>?</sup>. The appropriate number can be $`10^3`$ if reasonable continuity with $`T=0`$ results is to be verified. Meson $`\overline{q}q`$ modes at $`T>0`$ have often been studied within the Nambu–Jona-Lasinio model where the contact nature of the effective interaction allows decoupling of the Matsubara modes and also analytic methods of summation, see for example, Ref. <sup>?</sup> and references therein. However the lack of quark confinement in that model leads to unphysical thresholds for $`\overline{q}q`$ dissociation.
We consider here a simple separable interaction that has a finite range, accommodates quark confinement, and facilitates a decoupling of fermion Matsubara modes. We base our approach on a confining separable model <sup>?</sup> previously found to be successful at $`T=0`$ and defined by $`D_{\mu \nu }^{\mathrm{eff}}(pq)\delta _{\mu \nu }D(p^2,q^2,pq)`$ with
$$D(p^2,q^2,pq)=D_0f_0(p^2)f_0(q^2)+D_1f_1(p^2)(pq)f_1(q^2).$$
(3)
Here a Feynman-like gauge is chosen for phenomenological simplicity. This is a rank-2 interaction with two strength parameters $`D_0,D_1`$, and corresponding form factors $`f_i(p^2)`$. The choice for these quantities is constrained by consideration of the resulting solution of the DSE for the quark propagator in the rainbow approximation. For the amplitudes defined by $`S(p)=[i\text{/}pA(p^2)+B(p^2)+m_0]^1`$ this produces<sup>§</sup><sup>§</sup>§We choose the interaction form factors such that they provide sufficient ultraviolet suppression. Therefore no renormalization is needed and $`Z_2=1`$ .
$`B(p^2)`$ $`=`$ $`{\displaystyle \frac{16}{3}}{\displaystyle \frac{d^4q}{(2\pi )^4}D(p^2,q^2,pq)\frac{B(q^2)+m_0}{q^2A^2(q^2)+\left[B(q^2)+m_0\right]^2}},`$ (4)
$`\left[A(p^2)1\right]p^2`$ $`=`$ $`{\displaystyle \frac{8}{3}}{\displaystyle \frac{d^4q}{(2\pi )^4}D(p^2,q^2,pq)\frac{(pq)A(q^2)}{q^2A^2(q^2)+\left[B(q^2)+m_0\right]^2}}.`$ (5)
We note that if terms of higher order in $`pq`$ were to be included in Eq. (3), they would make no contribution to Eqs. (4) and (5) for the DSE solution <sup>?</sup>. The solution for $`B(p^2)`$ is determined only by the $`D_0`$ term, and the solution for $`A(p^2)1`$ is determined only by the $`D_1`$ term. Eq. (3) produces solutions having the form
$$B(p^2)=bf_0(p^2),A(p^2)=1+af_1(p^2),$$
(6)
and Eqs. (4) and (5) reduce to nonlinear equations for the constants $`b`$ and $`a`$.
### 2.1 Confinement and Dynamical Chiral Symmetry Breaking
If there are no poles in the quark propagator $`S(p)`$ for real timelike $`p^2`$ then there is no physical quark mass shell, quarks cannot propagate free of interactions, and the description of hadronic processes will not be hindered by spurious quark production thresholds. This is sufficient but not necessary for quark confinement and it remains a viable possibility for how quark confinement is realized <sup>?,?</sup>. The more general phenomena of confinement of colored multi-quark states and the non-existence of S-matrix elements connecting them to hadronic states are more subtle topics that do not concern us here <sup>?</sup>. A nontrivial solution for $`B(p^2)`$ in the chiral limit ($`m_0=0`$) signals dynamical chiral symmetry breaking. There is a connection between quark confinement as the lack of a quark mass shell and the existence of a strong quark mass function in the infrared through dynamical chiral symmetry breaking. This connection has proved to be empirically successful for the description of ground state light quark mesons and their form factors and decays <sup>?,?</sup>. In the present separable model the strength $`b=B(0)`$, which is generated by solution of Eqs. (4) and (5), controls both confinement and dynamical chiral symmetry breaking.
The propagator is confining if $`m^2(p^2)p^2`$ for real $`p^2`$ where the quark mass function is $`m(p^2)=(B(p^2)+m_0)/A(p^2)`$. With an exponential form factor $`f_0(p^2)=`$ exp $`(p^2/\mathrm{\Lambda }_0^2)`$, this condition is most transparent in the case of a rank-1 separable model where $`D_1=0`$ and $`A(p^2)=1`$, i.e., $`a=0`$. In the chiral limit, the model is confining if $`D_0`$ is strong enough to make $`b/\mathrm{\Lambda }_01/\sqrt{2\mathrm{e}}`$ and this finite $`b`$ also signals dynamical chiral symmetry breaking. Using $`\mathrm{\Lambda }_00.60.8`$ GeV as a typical range for a quark mass function $`m(p^2)`$, both confinement and dynamical chiral symmetry breaking will be compatible with $`m(p=0)0.3`$ GeV, an empirically viable value. The above qualitative properties will also hold in the case of a rank-2 model if the form factor ranges satisfy $`\mathrm{\Lambda }_1>\mathrm{\Lambda }_0`$ and this is compatible with empirical findings that the amplitude $`A(p^2)`$ typically has a larger momentum range than $`B(p^2)`$. At finite temperature, the strength $`b(T)`$ for the quark mass function will decrease with $`T`$ so that this model can be expected to have a deconfinement transition at or before the chiral restoration transition associated with $`b(T)0`$.
It is found that the simple choice $`f_i(p^2)=`$ exp $`(p^2/\mathrm{\Lambda }_i^2)`$ produces numerical solutions that describe infrared properties of $`\pi `$ and $`\omega `$ mesons very well and generate an empirically acceptable chiral condensate. At the same time the produced quark propagator is found to be confining and the infrared strength and shape of the quark amplitudes $`A(p^2)`$ and $`B(p^2)`$ are in qualitative agreement with the results of typical DSE studies <sup>?</sup>. We use the exponential form factors as a minimal way to preserve these properties while realizing that the ultraviolet suppression is much greater than the power law fall-off (with logarithmic corrections) known from asymptotic QCD. Most of our investigation centers on physics below and in the vicinity of the transition. We use the high $`T`$ behavior of masses in comparision with lattice results to discuss improvements appropriate for future work.
### 2.2 $`\pi `$ and $`\rho `$ bound states
With the separable interaction of Eq. (3), the allowed form of the solution of Eq. (1) for the pion BS amplitude $`\stackrel{}{\tau }\mathrm{\Gamma }_\pi (q;P)`$ is <sup>?</sup>
$$\mathrm{\Gamma }_\pi (q;P)=\gamma _5\left(iE_\pi (P^2)+\overline{)}PF_\pi (P^2)\right)f_0(q^2),$$
(7)
which contains the two dominant covariants from the set of four general covariants. The $`q`$ dependence is described only by the first form factor $`f_0(q^2)`$. The second term $`f_1`$ of the interaction can contribute only indirectly via the quark propagators. The $`\pi `$ BSE, Eq. (1), becomes a $`2\times 2`$ matrix eigenvalue problem $`𝒦(P^2)f=\lambda (P^2)f`$ where the eigenvector is $`f=(E_\pi ,F_\pi )`$. The kernel is
$$𝒦_{ij}(P^2)=\frac{4D_0}{3}\mathrm{tr}_\mathrm{s}\frac{d^4q}{(2\pi )^4}f_0^2(q^2)\left[\widehat{t}_iS(q_+)t_jS(q_{})\right],$$
(8)
where the $`\pi `$ covariants are $`t=(i\gamma _5,\gamma _5\overline{)}P)`$ and we have also introduced $`\widehat{t}=(i\gamma _5,`$ $`\gamma _5\overline{)}P/2P^2)`$. We note that the separable model produces the same $`q^2`$ shape for both amplitudes $`F_\pi `$ and $`E_\pi `$; the shape is that of the quark amplitude $`B(q^2)`$. For the general amplitude $`E_\pi (q;P)`$ this is the correct shape in the chiral limit; for physical quark masses it is still a very good approximation <sup>?</sup>. In general, the amplitude $`F_\pi (q;P)`$ does not have that shape; it is in fact linked with $`A(q^2)1`$ through the axial vector Ward-Takahashi identity (AV-WTI) <sup>?</sup>. However Goldstone’s theorem, the key consequence of the AV-WTI, is preserved by the present separable model; in the chiral limit, whenever a nontrivial DSE solution for $`B(p^2)`$ exists, there will be a massless $`\pi `$ solution to Eq. (8).
A simple illustrative truncation is obtained by setting $`F_\pi (P^2)=0`$ for then Eq. (8) reduces to an expression for the eigenvalue which is
$$\lambda _\pi (P^2)=\frac{16D_0}{3}\frac{d^4q}{(2\pi )^4}f_0^2(q^2)\left[\left(q^2\frac{P^2}{4}\right)\sigma _V^+\sigma _V^{}+\sigma _S^+\sigma _S^{}\right],$$
(9)
where the quark propagator amplitudes employed here are defined by $`S(p)=`$ $`i\text{/}p`$ $`\sigma _V(p^2)+`$ $`\sigma _S(p^2)`$. Since $`A=1`$ in rank-1, we omit the amplitude $`F_\pi `$ in that case. Thus it is Eq. (9) that we use for the $`\pi `$ calculation with a rank-1 interaction.
The vector mesons $`\rho `$ and $`\omega `$ are degenerate in the ladder approximation. We shall deal with the isovector $`\rho `$. For the BS amplitude $`\stackrel{}{\tau }\mathrm{\Gamma }_\mu ^\rho (q;P)`$ there are in general eight transverse covariants <sup>?</sup> and the dominant one is $`\gamma _\mu ^T(P)=`$ $`T_{\mu \nu }(P)\gamma _\nu `$ where $`T_{\mu \nu }(P)=\delta _{\mu \nu }P_\mu P_\nu /P^2`$. The solution of the rank-1 model contains only that term, that is, $`\mathrm{\Gamma }_\mu ^\rho (q;P)=`$ $`\gamma _\mu ^T(P)f_0(q^2)F_\rho (P^2)`$. The corresponding eigenvalue is given by
$$\lambda _\rho (P^2)=\frac{8D_0}{3}\frac{d^4q}{(2\pi )^4}f_0^2(q^2)\left[\left(q^2\frac{P^2}{4}\frac{2q^2}{3}(1z^2)\right)\sigma _V^+\sigma _V^{}+\sigma _S^+\sigma _S^{}\right],$$
(10)
where $`z=\widehat{q}\widehat{P}`$. For the rank-2 separable interaction, there are two other covariants besides $`\gamma _\mu ^T`$ that will appear in the solution of the BSE <sup>?</sup>. However it has been found in such a model that these subleading vector covariants make only a few percent contribution to the vector mass <sup>?</sup> and the associated $`g_{\rho \pi \pi }`$ <sup>?</sup>. In this present work we will ignore the subdominant covariants and employ $`\gamma _\mu ^T`$ as the only covariant for the vector meson. The differences in the vector mass obtained from Eq. (10) in rank-1 and rank-2 will be due to differences in the quark propagator in each case.
The normalization condition for the $`\pi `$ BS amplitude can be expressed as
$`2P_\mu ={\displaystyle \frac{}{P_\mu }}\mathrm{\hspace{0.17em}2}N_c\mathrm{tr}_s{\displaystyle \frac{d^4q}{(2\pi )^4}\overline{\mathrm{\Gamma }}_\pi (q;K)S(q_+)\mathrm{\Gamma }_\pi (q;K)S(q_{})}|_{P^2=K^2=M_\pi ^2}.`$ (11)
Here $`\overline{\mathrm{\Gamma }}(q;K)`$ is the charge conjugate amplitude $`[𝒞^1\mathrm{\Gamma }(q,K)𝒞]^\mathrm{t}`$ where $`𝒞=\gamma _2\gamma _4`$ and t denotes matrix transpose. This defines a normalization constant $`N_\pi `$ via $`E_\pi (P^2=M_\pi ^2)f_0(q^2)=`$ $`B(q^2)/N_\pi `$. The pion decay constant $`f_\pi `$ can be expressed as the loop integral
$`f_\pi P_\mu \delta _{ij}`$ $`=`$ $`0|\overline{q}{\displaystyle \frac{\tau _i}{2}}\gamma _\mu \gamma _5q|\pi _j(P)`$ (12)
$`=`$ $`\delta _{ij}N_c\mathrm{tr}_\mathrm{s}{\displaystyle \frac{d^4q}{(2\pi )^4}\gamma _5\gamma _\mu S(q_+)\mathrm{\Gamma }_\pi (q;P)S(q_{})}.`$
Note that Eq. (12) is the exact expression for $`f_\pi `$ except for the absence of the renormalization constant $`Z_2`$ which is unity in the present model. The decay constants thus test the quality of the infrared behavior of the quark propagator and the BS amplitudes. Similar expressions exist for the normalization of the vector meson BS amplitude and the coupling between a vector meson and a photon <sup>?,?</sup>.
In Table 2.2 the results for the ground state $`\pi `$ and $`\rho /\omega `$ mesons as well as related quantities are shown for both rank-1 and rank-2 versions of the model along with the values of the employed parameters. We consider the experimental $`\omega `$ mass to be the appropriate value for comparison with the vector result in ladder approximation in order to allow for subsequent preferential lowering of $`M_\rho `$ due to pion loop dressing <sup>?</sup>. The range parameter $`\mathrm{\Lambda }_0`$ is used to set the mass scale to reproduce the experimental $`M_\omega `$ for rank-1 and $`f_\pi `$ for rank-2. The obtained dynamical quark mass function at $`p^2=0`$ is $`m(p^2=0)=0.685`$ GeV for rank-1. For rank-2 we obtain $`A(p^2=0)=1.94`$ and $`m(p^2=0)=0.405`$ GeV. In both cases the dressed quark propagator is confining. The strengths obtained for the mass function are consistent with results from recent DSE solutions <sup>?,?</sup>. Also shown in Table 2.2 are the results obtained for the electromagnetic ($`g_\rho `$) and strong decays ($`g_{\rho \pi \pi }`$) of the $`\rho `$, which compare reasonably well with experiments.
The difference between $`N_\pi `$ and $`f_\pi `$ shown in Table 2.2 simply reflects the fact that the AV-WTI cannot be exactly satisfied within a separable model. (This is also the case in the Nambu–Jona-Lasinio model due to the required cut-off and in any approach that does not use a translationally invariant interaction or a sufficiently complete set of Dirac covariants for $`\mathrm{\Gamma }_\pi `$ and the vertex amplitudes <sup>?</sup>.) A translationally invariant model that makes the simplifying assumption $`A(p^2)=1`$ can achieve a fortuitious agreement <sup>?</sup> and this dominates the rank-1 result for $`N_\pi /f_\pi `$. Rather than use $`f_\pi `$ to define $`N_\pi `$, we calculate $`N_\pi `$ from its definition in Eq. (11) so that, in subsequent studies of processes such as $`\rho \pi \pi `$, the pion state is physically and consistently normed within the model.
## 3 Finite Temperature Extension
The extension of the separable model studies to $`T0`$ is systematically accomplished by transcription of the Euclidean quark 4-momentum via $`q`$ $`q_n=`$ $`(\omega _n,\stackrel{}{q})`$, where $`\omega _n=(2n+1)\pi T`$ are the discrete Matsubara frequencies. The effective $`\overline{q}q`$ interaction will automatically decrease with increasing $`T`$ without the introduction of an explicit $`T`$-dependence which would require new parameters. We investigate the resulting behavior of the $`\pi `$ and $`\rho `$ meson modes and decays in the presence of deconfinement and chiral restoration.
### 3.1 Chiral symmetry restoration and deconfinement
The result of the DSE solution for the dressed quark propagator now becomes
$$S^1(p_n,T)=i\stackrel{}{\gamma }\stackrel{}{p}A(p_n^2,T)+i\gamma _4\omega _nC(p_n^2,T)+B(p_n^2,T)+m_0,$$
(13)
where $`p_n^2=\omega _n^2+\stackrel{}{p}^{\mathrm{\hspace{0.17em}2}}`$ and there are now three amplitudes due to the loss of $`O(4)`$ symmetry. The solutions have the form $`B=b(T)f_0(p_n^2)`$, $`A=1+a(T)f_1(p_n^2)`$, and $`C=1+c(T)f_1(p_n^2)`$ and the DSE becomes a set of three non-linear equations for $`b(T)`$, $`a(T)`$ and $`c(T)`$. The explicit form is
$$a(T)=\frac{8D_1}{9}T\underset{n}{}\frac{d^3p}{(2\pi )^3}f_1(p_n^2)\stackrel{}{p}^{\mathrm{\hspace{0.17em}2}}[1+a(T)f_1(p_n^2)]d^1(p_n^2,T),$$
(14)
$$c(T)=\frac{8D_1}{3}T\underset{n}{}\frac{d^3p}{(2\pi )^3}f_1(p_n^2)\omega _n^2[1+c(T)f_1(p_n^2)]d^1(p_n^2,T),$$
(15)
$$b(T)=\frac{16D_0}{3}T\underset{n}{}\frac{d^3p}{(2\pi )^3}f_0(p_n^2)[m_0+b(T)f_0(p_n^2)]d^1(p_n^2,T),$$
(16)
where $`d(p_n^2,T)`$ is given by
$$d(p_n^2,T)=\stackrel{}{p}^{\mathrm{\hspace{0.17em}2}}A^2(p_n^2,T)+\omega _n^2C^2(p_n^2,T)+[m_0+B(p_n^2,T)]^2.$$
(17)
It is also useful to introduce the equivalent representation
$$S(p_n,T)=i\stackrel{}{\gamma }\stackrel{}{p}\sigma _A(p_n^2,T)i\gamma _4\omega _n\sigma _C(p_n^2,T)+\sigma _B(p_n^2,T),$$
(18)
where $`\sigma _A=A/d`$, $`\sigma _C=C/d`$ and $`\sigma _B=(B+m_0)/d`$.
The $`T`$-dependence of the solutions for $`A,B`$ and $`C`$ at $`\stackrel{}{p}^2=0`$ for the rank-2 model is displayed in Fig. 1. The results shown are for the lowest Matsubara mode ($`n=0`$) which provides the leading behavior as $`T`$ is increased. The chiral restoration critical temperature $`T_c`$ is identified from the vanishing of the chiral limit amplitude $`B_0(p=0,T)`$ as shown. We find $`T_c=`$ 121 MeV for the rank-2 model. Below $`T_c`$, $`A`$ and $`C`$ are relatively constant, $`O(4)`$ symmetry is approximately manifest, and the main effect is an almost constant quark wave function renormalization via $`1/C`$; the central feature shared by both rank-1 and rank-2 models is a rapidly decreasing mass function. Above $`T_c`$, there remains a significant temperature range where the self-interaction effects are strong, both $`A`$ and $`C`$ are considerably enhanced above their perturbative values, and the breaking of $`O(4)`$ symmetry is manifest. The present model thus captures the qualitative $`T`$-dependence observed for the dressed quark propagator in studies of the quark DSE <sup>?,?</sup>. The rank-1 limit has $`A=C=1`$ for all $`T`$ and the behavior of $`B(p=0,T)`$ is similar to that in Fig. 1 except $`T_c=`$ 146 MeV is obtained.
The order parameter for chiral restoration, the chiral quark condensate, can be obtained from the chiral limit quark propagator as $`\overline{q}q^0=N_c\mathrm{tr}_\mathrm{s}S_0(x,x)`$. Here this produces
$$\overline{q}q^0=4N_cT\underset{n}{}\frac{d^3p}{(2\pi )^3}\frac{b_0(T)f_0(p_n^2)}{d_0(p_n^2,T)},$$
(19)
where $`b_0`$ and $`d_0`$ are obtained from the chiral limit solution of the DSE. Both $`b_0(T)`$ and $`\overline{q}q^0`$ vanish sharply as $`(1T/T_c)^\beta `$ with the critical exponent having the mean field value $`\beta =1/2`$ in agreement with other rainbow DSE studies <sup>?</sup>.
The deconfinement temperature $`T_d`$ is found by a search for a propagator pole (a zero of the function $`d(p_0^2,T)`$) on the real $`\stackrel{}{p}^{\mathrm{\hspace{0.17em}2}}`$ axis. We find $`T_d=T_c=`$ 146 MeV for the rank-1 model and $`T_d=0.9T_c=`$ 105 MeV for the rank-2 model. Only about 15% variation in these transition temperatures can be achieved by variation of the model parameters while retaining a reasonable description of the observables shown in Table 2.2. Since all dynamical information concerning the rank-1 model is contained in one function $`B(p_n^2,T)`$, one may expect to find $`T_d=T_c`$ in that case. For rank-2, the dynamical information is contained in three functions and the result $`T_dT_c`$ is not surprising. It is however possible that the separable form of interaction used here might miss some dynamical correlations between $`A,B`$ and $`C`$ that would otherwise produce $`T_d=T_c`$.
### 3.2 Spatial $`\pi `$ correlations at $`T0`$
At $`T=0`$ the mass-shell condition for a meson as a $`\overline{q}q`$ bound state of the BSE is equivalent to the appearance of a pole in the $`\overline{q}q`$ scattering amplitude as a function of $`P^2`$. At $`T0`$ in the Matsubara formalism, the $`O(4)`$ symmetry is broken by the heat bath and we have $`P(\mathrm{\Omega }_m,\stackrel{}{P})`$ where $`\mathrm{\Omega }_m=2m\pi T`$. Bound states and the poles they generate in propagators may be investigated through polarization tensors, correlators or Bethe-Salpeter eigenvalues. This pole structure is characterized by information at discrete points $`\mathrm{\Omega }_m`$ on the imaginary energy axis and at a continuum of 3-momenta. Analytic continuation for construction of real-time Green’s functions (and related propagation properties) has been well-studied <sup>?</sup>. An unambiguous result is obtained by requiring that the continuation yield a function that is bounded at complex infinity and analytic off the real axis <sup>?</sup>. One may search for poles as a function of $`\stackrel{}{P}^2`$ thus identifying the so-called spatial or screening masses for each Matsubara mode. These serve as one particular characterization of the propagator and the $`T>0`$ bound states.
In the present context the eigenvalues of the meson BSE become $`\lambda (P^2)`$ $`\stackrel{~}{\lambda }(\mathrm{\Omega }_m^2,\stackrel{}{P}^2;T)`$. The temporal meson masses identified by zeros of $`1\stackrel{~}{\lambda }(\mathrm{\Omega }^2,0;T)`$ will be different in general from the spatial masses identified by zeros of $`1\stackrel{~}{\lambda }(0,\stackrel{}{P}^2;T)`$. They are however identical at $`T=0`$ and an approximate degeneracy can be expected to extend over the finite $`T`$ domain where the $`O(4)`$ symmetry is not strongly broken. From Fig. 1 one may reasonably expect this domain to extend up to $`T80100`$ MeV. At and above the transition, temporal and spatial masses can be expected to emphasize different aspects of the bound state modes. In this work we explore the $`T`$-dependence of the lowest spatial masses in the $`\pi `$ and $`\rho `$ channels.
In the $`\pi `$ channel, the correlator $`\mathrm{\Pi }(x)=TJ_{ps}(x)J_{ps}(0)`$ of two currents $`J_{ps}(x)=`$ $`\overline{q}(x)\gamma _5q(x)`$, after transformation to momentum space and extension to $`T>0`$ via the Matsubara formalism, can be expressed as
$$\mathrm{\Pi }(\mathrm{\Omega }^2,\stackrel{}{P}^2)=N_cT\underset{n}{}\mathrm{tr}_s\frac{d^3q}{(2\pi )^3}\gamma _5S(q_{n+})\mathrm{\Gamma }_{ps}(q_n;\mathrm{\Omega },\stackrel{}{P})S(q_n),$$
(20)
where $`q_{n\pm }=q_n\pm P/2`$ and $`\mathrm{\Gamma }_{ps}`$ is the pseudoscalar vertex which satisfies the inhomogeneous version of the pion BSE. The quark propagators appearing here are dressed according to solution of the DSE. When the ladder approximation is employed for $`\mathrm{\Gamma }_{ps}`$, along with the present separable approximation to the BSE kernel, the correlation function can be expressed as
$$\mathrm{\Pi }(\mathrm{\Omega }^2,\stackrel{}{P}^2)=\mathrm{\Pi }^{(0)}(\mathrm{\Omega }^2,\stackrel{}{P}^2)\frac{4D_0}{3}N_cL_i(\mathrm{\Omega }^2,\stackrel{}{P}^2)[1𝒦(\mathrm{\Omega }^2,\stackrel{}{P}^2)]_{ij}^1\widehat{L}_j(\mathrm{\Omega }^2,\stackrel{}{P}^2).$$
(21)
Here $`\mathrm{\Pi }^{(0)}`$ is the contribution to Eq. (20) arising from the zeroth order contribution ($`\gamma _5`$) to the pseudoscalar vertex $`\mathrm{\Gamma }_{ps}`$. The second term of Eq. (21) sums the interaction terms and factorizes due the separability of the effective interaction. The kernel $`𝒦`$ involves the $`T>0`$ extension of the kernel of the $`\pi `$ separable BSE given previously in Eq. (8). The loop integrals for the numerator are given by
$$L_i(\mathrm{\Omega }^2,\stackrel{}{P}^2)=T\underset{n}{}\mathrm{tr}_s\frac{d^3q}{(2\pi )^3}f_0(q_n^2)\gamma _5S(q_{n+})t_iS(q_n),$$
(22)
and
$$\widehat{L}_j(\mathrm{\Omega }^2,\stackrel{}{P}^2)=T\underset{n}{}\mathrm{tr}_s\frac{d^3q}{(2\pi )^3}f_0(q_n^2)\widehat{t}_jS(q_{n+})\gamma _5S(q_n).$$
(23)
With an eigenvector representation of the BSE kernel $`𝒦`$, Eq. (21) develops a denominator $`1\stackrel{~}{\lambda }_\pi (\mathrm{\Omega }^2,\stackrel{}{P}^2;T)`$. There is a pole in the correlator associated with the spatial mode solution to the homogeneous BSE identified from
$$1\stackrel{~}{\lambda }_\pi (0,\stackrel{}{P}^2;T)=Z_\pi ^1(\stackrel{}{P}^2,T)[\stackrel{}{P}^2+M_\pi ^2(T)]=0.$$
(24)
The masses so identified are spatial screening masses of the lowest mode associated with the 3-space asymptotic behavior $`\mathrm{\Pi }(x)\mathrm{exp}(Mx)`$. The identification of spatial masses by location of a pole is equivalent, by Fourier transformation, to the method of large spatial separation of sources used in lattice QCD.
The general form of the finite $`T`$ pion BS amplitude allowed by the separable model is
$$\mathrm{\Gamma }_\pi (q_n;P_m)=\gamma _5\left(iE_\pi (P_m^2)+\gamma _4\mathrm{\Omega }_m\stackrel{~}{F}_\pi (P_m^2)+\stackrel{}{\gamma }\stackrel{}{P}F_\pi (P_m^2)\right)f_0(q_n^2).$$
(25)
The separable BSE becomes a $`3\times 3`$ matrix eigenvalue problem with a kernel that is a generalization of Eq. (8). In the limit $`\mathrm{\Omega }_m0`$, as is required for the spatial mode of interest here, the amplitude $`\widehat{F}_\pi =\mathrm{\Omega }_m\stackrel{~}{F}_\pi `$ is trivially zero. The pseudovector amplitude $`F_\pi `$ is significantly different from zero below $`T_c`$, but decreases rapidly above the transition. In the chiral limit, it vanishes identically at and beyond the transition; above the transition, only $`E_\pi `$ is survives.
The result for the $`\pi `$ mass is displayed in Fig. 2 for rank-2; the results for the rank-1 model are similar. In both models, $`M_\pi (T)`$ is seen to be only weakly $`T`$-dependent until near $`T_c`$ where a sharp rise begins. These qualitative features of the response of the pion mode with $`T`$ agree with the results deduced from the DSE in Ref. <sup>?</sup> and also with the more detailed study of Ref. <sup>?</sup> using a ladder-rainbow truncation of the DSE/BSE system that preserves the one-loop renormalization group properties of QCD. Evidently the detailed character of the effective interaction in the perturbative region does not dominate at the level of the present qualitative investigation.
### 3.3 Chiral symmetry and $`\pi `$ mass relation
To explore the extent to which the model respects the detailed constraints from chiral symmetry, we investigate the exact QCD pseudoscalar mass relation <sup>?</sup> which, after extension to the spatial mode at $`T>0`$, is
$$M_\pi ^2(T)f_\pi (T)=2m_0r_P(T).$$
(26)
Here $`r_P`$, the residue at the pion pole in the pseudoscalar vertex, is given by the pseudoscalar projection of the pion wavefunction onto zero quark-antiquark separation, that is
$$ir_P(T)=N_cT\underset{n}{}\mathrm{tr}_s\frac{d^3q}{(2\pi )^3}\gamma _5S(q_n+\frac{\stackrel{}{P}}{2})\mathrm{\Gamma }_\pi (q_n;\stackrel{}{P})S(q_n\frac{\stackrel{}{P}}{2}).$$
(27)
The generalization of Eq. (12) for $`f_\pi `$ to finite temperature in the case of the spatial pion mode is
$$P_if_\pi (T)=N_cT\underset{n}{}\mathrm{tr}_s\frac{d^3q}{(2\pi )^3}\gamma _5\gamma _iS(q_n+\frac{\stackrel{}{P}}{2})\mathrm{\Gamma }_\pi (q_n;\stackrel{}{P})S(q_n\frac{\stackrel{}{P}}{2}).$$
(28)
Eqs. (27) and (28) are exact expressions for $`r_P(T)`$ and $`f_\pi (T)`$ except for the absence of the renormalization constants which are trivially equal to one in this separable model. The relation in Eq. (26) is a consequence of the pion pole structure of the isovector axial Ward identity which links the quark propagator, the pseudoscalar vertex and the axial vector vertex <sup>?</sup>. In the chiral limit, $`r_P`$ $`\overline{q}q^0/f_\pi ^0`$ and Eq.(26), for small mass, produces the Gell-Mann–Oakes–Renner (GMOR) relation. The exact mass relation, Eq. (26), can only be satisfied approximately when the various quantities are obtained in a manner that does not preserve axial Ward identity. The error can be used to assess the reliability of the present approach to modeling the behavior of the pion spatial mode as the temperature is varied.
Our findings in the case of the rank-2 model are displayed in Fig. 2. There the solid line represents $`r_P(T)`$ calculated from the quark loop integral in Eq. (27); the dotted line represents $`r_P`$ constructed from the other quantities in Eq. (26). The exact mass relation, Eq. (26), is violated by about 25% almost independent of temperature, even above $`T_c`$. The rank-1 model satisfies this relation to within 1%, due to the special but unrealistic case $`A=1=C`$. We have also investigated the (approximate) GMOR relation within the present model. The quantity $`\overline{q}q^0/N_\pi ^0`$, displayed in Fig. 2, is the chiral limit of $`r_P`$ in this model. If all covariants for the pion were retained and the axial vector Ward identity were obeyed, one would have $`N_\pi ^0=f_\pi ^0`$ in the chiral limit <sup>?</sup>. If the GMOR relation were exactly obeyed, the long-dashed line representing $`\overline{q}q^0/N_\pi ^0`$ would coincide with the dotted line representing $`M_\pi ^2f_\pi /2m_0`$. These features are temperature-independent until about $`0.9T_c`$, consistent with an earlier study of low energy theorems at $`T>0`$ within a three-space, non-confining separable interaction model <sup>?</sup>. Close to $`T_c`$, the GMOR relation breaks down, and above $`T_c`$, it is no longer well-defined.
It should be noted that $`f_\pi ^0,N_\pi ^0`$ and $`\overline{q}q^0`$ are equivalent order parameters for the critical behavior near $`T_c`$ and have weak $`T`$-dependence below $`T_c`$. A consequence is that $`M_\pi ^2f_\pi `$, $`r_P`$ and $`\overline{q}q^0/N_\pi ^0`$ are almost $`T`$-independent and so are the estimated errors for the two mass relations linking these quantities. We have employed the physical $`f_\pi `$ defined at non-zero current quark mass, and this does not vanish at $`T_c`$ but continuously decreases.
### 3.4 Spatial $`\rho `$ correlations at $`T0`$
The $`T=0`$ transverse vector meson, that we have described by the covariant $`\gamma _\mu ^T`$, splits for $`T>0`$ into 3-space longitudinal and transverse modes. For the spatial modes characterized by $`P=(0,\stackrel{}{P})`$ the BS amplitudes are
$$\mathrm{\Gamma }_\mu ^{\rho (L)}(q_n;\stackrel{}{P})=\delta _{\mu 4}\gamma _4f_0(q_n^2)F_{\rho (L)}(\stackrel{}{P}^2),$$
(29)
and
$$\mathrm{\Gamma }_i^{\rho (T)}(q_n;\stackrel{}{P})=\left(\gamma _i\frac{P_i\stackrel{}{P}\stackrel{}{\gamma }}{\stackrel{}{P}^2}\right)f_0(q_n^2)F_{\rho (T)}(\stackrel{}{P}^2).$$
(30)
The $`T`$-dependence of the corresponding masses is displayed in Fig. 3 for the rank-2 model. These modes are effectively degenerate and $`T`$-independent until about $`T_c/2`$ where the breaking of $`O(4)`$ invariance becomes significant. The qualitative features $`M_\rho ^L(T)>M_\rho ^T(T)`$ and $`M_\rho ^T(T)\mathrm{const}`$ for $`T<T_c`$ seen here in the present context of a finite range interaction have previously been noted within the limiting case of the zero momentum range ID model <sup>?</sup>. This latter model was not applied for $`T>T_c`$. We discontinue the present study of the longitudinal mode at $`T`$ 180 MeV where it becomes unstable to $`\overline{q}q`$ dissociation. The transverse mode continues to be below the spatial $`\overline{q}q`$ threshold for the temperature range displayed.
The $`T=0`$ expression <sup>?,?</sup> for the electromagnetic coupling constant $`g_\rho `$ has a straightforward extension to $`T>0`$ for a transverse spatial $`\rho ^0`$ mode. Use of the $`\mathrm{\Omega }_m=0`$ solution described above yields
$`{\displaystyle \frac{M_\rho ^T(T)^2}{g_\rho (T)}}`$ $`=`$ $`{\displaystyle \frac{N_cT}{2}}{\displaystyle \underset{n}{}}\mathrm{tr}_s{\displaystyle \frac{d^3q}{(2\pi )^3}\gamma _iS(q_n+\frac{\stackrel{}{P}}{2})\mathrm{\Gamma }_i^{\rho (T)}(q_n;\stackrel{}{P})S(q_n\frac{\stackrel{}{P}}{2})}.`$ (31)
With summation over enough Matsubara modes for convergence, the produced $`g_\rho (T)`$ tends smoothly to the previously determined $`T=0`$ result. The result over a temperature range that extends just beyond $`T_c`$ is displayed in Fig. 4. There it is confirmed that there is very little $`T`$-dependence below about $`0.9T_c`$ where there is approximate $`O(4)`$ symmetry as evident in the quark propagator behavior in Fig. 1.
The impulse approximation for the $`\rho \pi \pi `$ vertex <sup>?,?</sup>, after extension to $`T>0`$ for spatial modes characterized by $`Q=(0,\stackrel{}{Q})`$ for the $`\rho `$ and $`P=(0,\stackrel{}{P})`$ for the relative $`\pi \pi `$ momentum, takes the form
$`\mathrm{\Lambda }_\nu (P,Q)`$ $`=`$ $`P_\nu g_{\rho \pi \pi }(T)`$ (32)
$`=`$ $`2N_cT{\displaystyle \underset{n}{}}\mathrm{tr}_s{\displaystyle \frac{d^3q}{(2\pi )^3}\mathrm{\Gamma }_\pi (k_{n+};\stackrel{}{P}_+)S(q_{n+})\mathrm{\Gamma }_\nu ^{\rho (T/L)}(q_{n+};\stackrel{}{Q})}`$
$`\times S(q_{n++})\mathrm{\Gamma }_\pi (k_n;\stackrel{}{P}_{})S(q_n).`$
Use of Eq. (29) immediately shows that the longitudinal $`\rho `$ mode cannot couple to $`\pi \pi `$. The temperature dependence obtained for the transverse $`\rho `$ coupling constant $`g_{\rho \pi \pi }(T)`$ is displayed in Fig. 4. Again one observes continuity with the $`T=0`$ result and a very weak $`T`$-dependence until about $`0.8T_c`$. Around $`T_c`$, the coupling constant decreases significantly.
For both interactions of the $`\rho `$ mode, and below $`T_c`$, one expects qualitatively similar behavior from spatial modes associated with higher meson Matsubara frequencies $`\mathrm{\Omega }_m0`$ except that the effect of $`O(4)`$ symmetry breaking will be evident earlier. Above $`T_c`$, each $`2\pi T`$ increment to the meson Matsubara frequency adds significantly to the quark effective mass in integrals like Eq. (31). Since dynamical chiral symmetry breaking is now absent, the quark propagators are of Dirac vector character, the meson Matsubara frequencies are the largest mass scale in the system, and perturbative behavior becomes increasingly dominant. The various spatial modes from $`\mathrm{\Omega }_m0`$ are characterized by masses much greater than that of the lowest mode considered here. We anticipate that this lowest mode characterizes the qualitative behavior of the physical decay processes $`\rho ^0e^+e^{}`$ and $`\rho \pi \pi `$. Certainly, a high temperature limit in which some modes vanish and others diverge would indicate a non-analytic behavior in the reconstructed physical amplitude that is physically untenable. We therefore estimate the relevant decay widths by combining the present results for the coupling constants with the relevant phase space factors generated also from the lowest spatial mass modes.
This leads to the electromagnetic decay width
$`\mathrm{\Gamma }_{\rho ^0e^+e^{}}(T)`$ $`=`$ $`{\displaystyle \frac{4\pi \alpha ^2M_\rho ^T(T)}{3g_\rho (T)^2}},`$ (33)
while the corresponding strong decay width is
$$\mathrm{\Gamma }_{\rho \pi \pi }(T)=\frac{g_{\rho \pi \pi }(T)^2}{4\pi }\frac{M_\rho ^T(T)}{12}\left[1\frac{4M_\pi (T)^2}{M_\rho ^T(T)^2}\right]^{3/2}.$$
(34)
The $`T`$-dependence estimated in this way for the decay widths is due to the response of the quark substructure to the heat bath, particularly the restoration of chiral symmetry. The results are displayed in Fig. 5. The contrast between the behavior of the electromagnetic and strong widths near and just above $`T_c`$ should be a more robust finding than the details of the individual processes. The strong width decreases rapidly and vanishes just above $`T_c`$ while the electromagnetic width remains within 20% of the $`T=0`$ value. Part of the strong decrease of the intrinsic $`\pi \pi `$ width of the transverse $`\rho `$ is due to the decrease in the coupling constant, however the dominant effect is the $`T`$-dependence of the last factor in Eq. (34). As displayed in Fig. 3, $`2M_\pi (T)`$ rises faster with $`T`$ than does $`M_\rho ^T(T)`$ until at $`T=1.17T_c`$ we have $`M_\rho ^T=2M_\pi `$. Beyond this point, the phase space factor vanishes and the strong decay $`\rho ^T\pi \pi `$ is blocked. This suggests that the total $`\rho ^T`$ width of 151 MeV at $`T=0`$ decreases by about 50% near $`T=T_c`$ and drops sharply to the electromagnetic value of about 6 keV by $`T=1.17T_c`$.
This narrowing of the intrinsic decay width of the vector meson mode in the heat bath is a mechanism that is distinct from the collisional broadening effect <sup>?</sup> from the many-hadron environment. The present work indicates that there is a non-trivial $`T`$-dependence to intrinsic coupling constants such as $`g_{\rho \pi \pi }`$ and decay phase space. The intrinsic effect tends to significantly decrease the decay width; the many-hadron medium effects have the opposite influence. Phenomenological forms for the $`T`$-dependence of $`M_\rho (T)`$ and $`\mathrm{\Gamma }_{\rho ^0\pi \pi }(T)`$ have often been explored in studies of medium effects in heavy-ion collisions. For example, both an increase in the width of the form $`\mathrm{\Gamma }_\rho ^I/(1T^2/T_c^2)`$, and a $`\rho `$ mass decreasing by 50% at $`T_c`$ have been explored in an effort to understand the heavy-ion dilepton spectrum <sup>?</sup>. A coordinated approach is called for in which hadronic collisional broadening mechanisms are built upon intrinsic coupling constants that respect the temperature and density dependence of the quark-gluon content.
## 4 Behavior at large $`T`$
The quark deconfinement point $`T_d`$ and the chiral restoration point $`T_c`$ are generally expected be to identical or nearly so <sup>?</sup>. (The present separable model produces $`T_d=T_c`$ in rank-1 and $`T_d=0.9T_c`$ in rank-2.) It might be expected therefore that above $`T_c`$ meson modes should dissolve in favor of a gas of essentially massless quarks. However for a significant temperature range above $`T_c`$, the spatial $`\pi `$ and $`\rho `$ modes studied here continue to be stable against $`\overline{q}q`$ dissociation and do not dissolve into a free quark gas. This situation is clearest for the simpler rank-1 separable interaction. The results for $`M_\pi (T)`$ and $`M_\rho ^T(T)`$ obtained from the eigenmass condition $`\lambda (M^2)=1`$ are displayed in Fig. 6 along with the quark dynamical mass function at $`p=0`$. The masses of both spatial meson modes approach the asymptotic behavior $`2\pi T`$ from below. This asymptotic behavior has been discussed previously <sup>?</sup> and is observed in lattice simulations <sup>?,?</sup>.
The manner in which the $`2\pi T`$ behavior at large $`T`$ emerges from the present description is illustrated well by the rank-1 model. For $`T>T_c`$, the dynamically generated mass function of the quarks is essentially negligible, and the quark propagator is dominated by the Dirac vector amplitude $`\sigma _V(q_n^2)1/q_n^2`$. For the spatial $`\pi `$ mode characterized by $`P=(0,\stackrel{}{P})`$, the loop integral for the “polarization” function or BSE eigenvalue $`\lambda _\pi (\stackrel{}{P}^2)`$, given by Eq. (9), yields at large $`T`$
$$\lambda _\pi (\stackrel{}{P}^2)\frac{16D_0}{3}T\frac{d^3q}{(2\pi )^3}f_0^2(\pi ^2T^2+\stackrel{}{q}^{\mathrm{\hspace{0.17em}2}})\frac{\pi ^2T^2+\stackrel{}{q}^{\mathrm{\hspace{0.17em}2}}\frac{\stackrel{}{P}^2}{4}}{[\pi ^2T^2+(\stackrel{}{q}+\stackrel{}{P}/2)^2][\pi ^2T^2+(\stackrel{}{q}\stackrel{}{P}/2)^2]},$$
(35)
where only the dominant zeroth fermion Matsubara mode has been retained. For $`T\mathrm{\Lambda }_0/\pi `$, with $`\mathrm{\Lambda }_0`$ being the range of the interaction form factor $`f_0`$, only small $`q`$ is relevant and the position of the lowest singularity as a function of $`\stackrel{}{P}^2<0`$ approaches $`2\pi T`$. The higher the temperature, the more $`\lambda _\pi (\stackrel{}{P}^2)`$ is suppressed except near the singularity. The value $`\lambda _\pi (M^2)=1`$ must be encountered before the divergence and thus $`M_\pi (T)2\pi T\mathrm{\Delta }_\pi (T)`$ where $`\mathrm{\Delta }_\pi `$ is a positive mass defect that will typically decrease with $`T`$. Thus the spatial meson mass or screening mass will approach the thermal mass of a pair of massless fermions from below. This limit has also been demonstrated from the pseudoscalar correlator within the Nambu–Jona-Lasinio model <sup>?</sup>.
In the general case, the detailed temperature behavior of the mass defect $`\mathrm{\Delta }_\pi (T)`$, or the nature of the approach to the $`2\pi T`$ limit, depends upon the asymptotic behavior of the quark amplitudes $`A(q^2)1`$ and $`B(q^2)`$. In QCD, their leading asymptotic behavior is $`1/q^2`$ apart from slow logarithmic corrections. The form factors chosen in this initial work within a separable model induce an exponential fall-off. It is to be expected therefore that the present model estimate of $`\mathrm{\Delta }_\pi (T)`$ will decrease too rapidly with $`T`$.
In Fig. 7 the $`\pi `$ and transverse $`\rho `$ spatial masses at $`T>T_c`$ from both the rank-1 separable model and the ID model (see the Appendix for more details on the ID model) are presented in comparison with lattice QCD simulations of spatial screening masses <sup>?</sup>. The solid horizontal line marks $`2\pi `$ while the lower horizontal dot-dashed line represents the lattice free limit corrected <sup>?</sup> for the lattice time extent $`N_t`$. From Fig. 7 it is evident that for the $`\pi `$ at $`T>T_c`$, the spatial or screening mass defect $`\mathrm{\Delta }=`$ $`2\pi TM(T)`$ is decreasing more rapidly above $`T_c`$ in the rank-1 separable model than is evident from the lattice simulations. This is consistent with the exponential behavior of the employed form factors. It is also consistent with the absence of quark vector self-energy amplitudes $`A(p,T)1`$ and $`C(p,T)1`$ through which interactions can be quite persistent in the asymptotic region. This can be demonstrated within the chiral limit ID model where those amplitudes are strong and indeed have the power law fall-off. As seen in Fig. 7, the resulting mass defect for both $`\pi `$ and $`\rho `$ is in fact too strong when compared with the lattice QCD simulations. This persistent self-interaction well above $`T_c`$, which slows the approach to free behavior such as Stefan-Boltzmann thermodynamics <sup>?</sup>, may well be what is signaled by the lattice QCD data in Fig. 7.
## 5 Discussion
We have explored $`\pi `$ and $`\rho `$ spatial correlation modes at $`T>0`$ within the rainbow-ladder truncation of the quark Dyson-Schwinger equation and the $`\overline{q}q`$ Bethe-Salpeter equation in the Matsubara formalism. With parameters fitted to $`T=0`$ properties, the model possesses dynamical chiral symmetry breaking and quark confinement. A simple separable form of the effective interaction is employed and this facilitates the use of sufficient Matsubara modes to allow coverage of low temperatures as well as the transition region. Deconfinement and chiral restoration transition temperatures, $`T_d`$ and $`T_c`$, are very similar and in the range 100-150 MeV. The $`T`$-evolution of the $`\pi `$ and $`\rho `$ $`\overline{q}q`$ states in the presence of the deconfinement and chiral restoration mechanisms is studied. The degree to which the model respects the axial vector Ward-Takahashi relation is evaluated in terms of the exact pion mass relation and the related GMOR relation. The $`\rho `$ electromagnetic coupling constant $`g_\rho `$ and the strong coupling constant $`g_{\rho \pi \pi }`$ are also obtained as a function of $`T`$. Estimates are made for the $`T`$-dependence of the widths for $`\rho ^0e^+e^{}`$ and $`\rho \pi \pi `$. Finally, the high $`T`$ behavior of the spatial masses is compared to that of spatial screening masses from lattice QCD simulations <sup>?</sup>.
The masses $`M_\pi (T)`$, $`M_\rho ^T(T)`$ and $`M_\rho ^L(T)`$ are found to be almost $`T`$-independent below $`T_c`$ followed by a strong increase. This behavior is characteristic of lattice QCD simulations <sup>?,?</sup> and DSE studies <sup>?,?</sup>. Our estimate of the $`\rho \pi \pi `$ strong decay width shows a decrease with $`T`$ such that, by $`T_c`$, it has been reduced by 50%. Our arguments suggest a phase-space blocking effect at about 25 MeV above $`T_c`$. (A similar phase-space blocking has been argued before only for the strong $`\pi \pi `$ decay of the scalar-isoscalar partner of the pion near $`T_c`$ <sup>?,?</sup>.) The tendency here is for the transverse $`\rho `$ mode above $`T_c`$ to be left with a narrow total width typical of electromagnetic decay. One would expect the mass of the pseudoscalar $`K`$ correlation to rise with $`T`$ in a similar fashion to $`M_\pi `$, while the masses of the vector $`\varphi `$ and $`K^{}`$ modes should rise like the $`\rho `$. This suggests that the vector modes $`\rho `$, $`K^{}`$ and $`\varphi `$ tend to be trapped with their relatively long electroweak lifetimes and with significantly increased masses for a domain of high temperatures above the transition. Between 0.5-1.0$`T_c`$, the transition probability connecting vector correlations with pairs of pseudoscalars would be reduced. This suggests that within the gas of pions and other pseudoscalars that dominate the hot hadronic product from heavy-ion collisions, the role of vector meson correlations in producing the dilepton spectra could be significantly less than conventional expectations. The present findings follow from the response of the quark-gluon content of the mesons to the heat bath. A different phenomenon is the coupling of the meson modes to the many-hadron environment which introduces a collisional broadening effect <sup>?</sup>. An approach that incorporates both phenomena is clearly called for.
Only the spatial or screening $`\overline{q}q`$ masses have been investigated within this model. The temporal masses (sometimes called dynamical or pole masses) provide a different characterization of the correlations. Lattice simulations indicate that spatial masses become much larger than the temporal masses above $`T_c`$ <sup>?</sup>. A Nambu–Jona-Lasinio model study <sup>?</sup> found that they are significantly different only in the range 150 MeV $`<T<`$ 350 MeV. A possible explanation for this discrepancy lies in our finding that the mass defect $`\mathrm{\Delta }=`$ $`2\pi TM(T)`$ at high $`T`$ is significantly influenced by residual non-perturbative interaction effects in the Dirac vector amplitude $`A(P^2)1`$ of the quark self-energy. Such a term is not present in the Nambu–Jona-Lasinio model which produces a momentum-independent quark constituent mass. The present model allows low temperature confinement and a momentum-dependent quark self-energy in a simple way but at a cost of an exponential asymptotic fall-off instead of a more realistic power law fall-off. Comparison with lattice QCD results illustrates the connection between asymptotic behavior of the interaction and the high $`T`$ behavior of the mass defects $`\mathrm{\Delta }_\pi (T)`$ and $`\mathrm{\Delta }_\rho (T)`$.
The simplicity of a separable representation of the effective quark-quark interaction of the type studied here might be of advantage in the consideration of pion loop effects and bulk thermodynamic properties of the hadron-quark matter phase transitions. Initial thermodynamic considerations have been reported recently <sup>?</sup>. The only $`T`$-dependence of the effective quark-quark interaction implemented by the present separable model is that generated by the Matsubara frequencies that enter through the momentum dependence. The strength and range parameters have been kept $`T`$-independent and no attempt has been made to introduce explicit $`T`$-dependent characteristics such as a Debye mass. Such considerations are more appropriately handled in approaches that have a better connection to a perturbative gluon propagator <sup>?</sup>.
Acknowledgements
We acknowledge fruitful interactions and conversations with B. Van den Bossche, M. Buballa, C.D. Roberts, and S. Schmidt. The work of G.B. and Y.L.K. has been supported by the Max-Planck-Gesellschaft and by the DFG Graduiertenkolleg “Stark korrelierte Vielteilchensysteme”. D.B. and G.B. gratefully acknowledge financial support by the Deutscher Akademischer Austauschdienst (DAAD) for visits to the Center for Nuclear Research at Kent State University where part of this work was conducted. Y.L.K. also acknowledges the Russian Fund for Fundamental Research, under contract number 97-01-01040, and the support of the Heisenberg-Landau program. P.C.T. and P.M. acknowledge support by the National Science Foundation under Grant Nos. INT-9603385 and PHY97-22429 and the hospitality of the University of Rostock where part of this work was conducted during several visits.
Appendix A
The extension to $`T>0`$ of the ID model introduced at $`T=0`$ by Munczek and Nemirovsky <sup>?</sup> provides a semi-analytic perspective on the large $`T`$ behavior of the meson masses. For the Feynman-like gauge used in the present study, the effective interaction of this model is specified by
$$D(pq)(2\pi )^4\frac{3\eta ^2}{16}\delta ^4(pq),$$
(A.1)
which is to be used in the DSEs given in Eqs. (4) and (5) and also in the BSE given in Eq. (1)The results are unchanged when Landau gauge is used if Eq. (A.1) is scaled up by $`4/3`$.. In the chiral limit, closed form expressions exist for the resulting quark propagator amplitudes and for the ladder BSE eigenvalues $`\lambda (P^2)`$ in the pseudoscalar and vector channels of interest here. The chiral limit DSE solution is of the general form given in Eq. (18) with
$`\sigma _A(p^2)=\sigma _C(p^2)`$ $`=`$ $`\{\begin{array}{cc}\frac{2}{\eta ^2},& p^2\frac{\eta ^2}{4}\hfill \\ \frac{2}{p^2}\left[1+\left(1+\frac{2\eta ^2}{p^2}\right)^{\frac{1}{2}}\right]^1,& p^2\frac{\eta ^2}{4}\hfill \end{array}`$ (A.4)
$`\sigma _B(p^2)`$ $`=`$ $`\{\begin{array}{cc}\frac{1}{\eta ^2}(\eta ^24p^2)^{\frac{1}{2}},& p^2\frac{\eta ^2}{4}\hfill \\ 0,& p^2\frac{\eta ^2}{4}\hfill \end{array}.`$ (A.7)
The quark mass function is $`m^2(p^2)=\eta ^2/4p^2`$ for $`p^2\eta ^2/4`$, and $`m^2(p^2)=0`$ otherwise. There is no mass-shell, i.e. $`p^2+m^2(p^2)0`$ for any $`p^2`$, and there is quark confinement. The above solution holds at $`T=0`$ and at $`T>0`$ with $`p^2`$ replaced by $`p_n^2=`$ $`\omega _n^2+\stackrel{}{p}^{\mathrm{\hspace{0.17em}2}}`$. Thus for $`T>T_c=`$ $`\eta /(2\pi )`$ there is chiral restoration ($`\sigma _B=0`$). Results from this model at finite $`T,\mu `$ have been obtained for quark thermodynamics outside the phase boundary of chiral restoration as governed by the quark propagator <sup>?</sup>. The behavior of the masses of both the $`\pi `$ and $`\rho `$ modes for $`T<T_c`$ (and also for chemical potential dependence for $`\mu <\mu _c`$) have been obtained previously <sup>?</sup>.
With Eq. (A.1), the BSE given in Eq. (1) becomes
$$\lambda (P^2)\mathrm{\Gamma }(q;P)=\frac{\eta ^2}{4}\gamma _\mu S(q_+)\mathrm{\Gamma }(q;P)S(q_{})\gamma _\mu ,$$
(A.8)
and solutions are possible if $`q`$ is fixed by $`P`$. At $`T>0`$, the solution that connects smoothly to the $`T=0`$ solution has the normally independent variable $`q_n=(\omega _n,\stackrel{}{q})`$ restricted to $`q_n(\omega _0,\stackrel{}{0})`$. To obtain spatial meson modes to compare with the separable model results, we again set $`P=(0,\stackrel{}{P})`$. Within the temperature domain $`T<T_c=\eta /(2\pi )`$, the results for the pseudoscalar and vector spatial modes are particularly simple and have been discussed previously <sup>?</sup>. One finds $`\lambda _\pi (0)=1`$ for the chiral limit $`\pi `$; thus $`M_\pi =0`$. For the vector meson one finds that $`\lambda _\rho ^T(\eta ^2/2)=1`$; thus $`M_\rho ^T=\eta /\sqrt{2}`$. These are also the correct $`T=0`$ results of the model. They hold over the finite temperature domain for which the quark mass function appropriate to the propagators occurring in the BS Eq. (A.8) is nonzero. For the equation appropriate to a meson of mass $`M`$, the relevant domain for the present model is $`\eta ^24s>0`$ where $`s=\pi ^2T^2M^2/4`$. For $`M_\pi =0`$, this temperature domain corresponds to that for which the model DSE generates a dynamical quark mass function in accord with the Goldstone theorem. For larger $`M`$, this temperature domain will be larger. The vector result $`M_\rho =\eta /\sqrt{2}`$ holds for the larger temperature domain $`T<\sqrt{3/2}T_c`$. We fix the single parameter $`\eta =1.107`$ GeV, so that $`M_\rho =0.783`$ GeV. This produces $`T_c=0.176`$ GeV.
Beyond the temperature domain where the meson mass is constant, one finds $`\lambda _\pi (\stackrel{}{P}^2)=`$ $`\eta ^2s_{}\sigma _A(s_+)^2`$, where $`s_\pm =`$ $`\pi ^2T^2\pm \stackrel{}{P}/4`$, and the $`\rho `$ eigenvalue is simply given by $`\lambda _\rho ^T(\stackrel{}{P}^2)=`$ $`\frac{1}{2}\lambda _\pi (\stackrel{}{P}^2)`$. In general both functions $`\lambda (M^2)`$ decrease with increasing $`T`$ and increase with increasing $`M`$. Thus the spatial eigenmode condition $`\lambda (M^2)=1`$ will tend to maintain $`M_\rho ^T(T)>M_\pi (T)`$ while both masses rise with $`T`$, and the obtained behaviour of the masses $`M_\rho ^T`$ and $`M_\pi `$ is qualitatively the same as what we found in the separable model. The main difference is that the mass defect $`\mathrm{\Delta }(T)=2\pi TM(T)`$ is significantly larger in the ID model. The leading large $`T`$ behavior
$$M_\rho ^T(T)^2(2\pi T)^2\left(1\frac{\eta }{\pi T}\mathrm{}\right),$$
(A.9)
leads to the spatial mass defect of the transverse $`\rho `$ mode having the asymptotic value $`\mathrm{\Delta }_\rho (T)\eta `$. Fig. 7 displays the behavior for both $`\pi `$ and $`\rho `$.
References |
warning/0002/astro-ph0002096.html | ar5iv | text | # The Broad Line Region in Active Galactic Nuclei
## 1 Introduction
The Broad Line Region (BLR) in Active Galactic Nuclei (AGN) is unresolved with present day imaging detectors and it will remain so for the foreseeable future. This is why “one quasar spectrum is really worth a thousand images” as stressed by Gary Ferland at this meeting. In response we would add that a thousand spectra are better than one average spectrum. Understanding the diversity in optical spectroscopic properties of AGN is the key to any realistic AGN modeling (??). The ideal to reconstruct the BLR velocity field from a single profile is not realistic (?).
Determination of BLR structure and kinematics can be approached in two ways. It has been recognized for a long time that strong broad and narrow emission lines coming from both high and low ionization species are present in Seyfert galaxies and quasars. This is considered a defining spectroscopic property of AGN. Restricting attention to broad lines: a) typical (i.e., strongest and most frequently observed) high ionization lines (HIL: ionization potential $`\stackrel{>}{}`$ 50 eV) are Civ$`\lambda `$1549 Heii$`\lambda `$4686 and Heii$`\lambda `$1640 lines; while b) observed low ionization lines (LIL: ionization potential $`\stackrel{<}{}`$ 20 eV) include Hi Balmer lines, Feii multiplets, Mgii$`\lambda `$2800, and the Caii IR triplet.
The first approach involves the study of line variability in response to continuum changes. This approach has been pursued through a number of successful monitoring campaigns using Reverberation Mapping (RM) techniques. RM requires a large amount of telescope time and, consequently, has been achieved only for a handful of sources. RM confirms that photoionization is the main heating process in the BLR and that a large part of the BLR is optically thick to the ionizing continuum (e. g. ?). RM studies have quantified a main difference between HIL and LIL; HIL respond to continuum changes with a time delay of a few days while the LIL respond with a delay of tens of days. This implies that the LIL are emitted at larger distance from the continuum source (??. An exhaustive list of reference can be found in ?). RM applied to line profiles suffers from uncertainty in our knowledge about the physics of continuum and broad line formation, so that conflicting models can still describe the same lag times (??).
The second approach involves statistical analysis of large samples of line profiles which differ because of properties that may affect the BLR structure, for example samples of radio quiet (RQ) and radio loud (RL) AGN. Statistical studies can be done for a single line or by comparing lines sensible to different physical parameters (e.g., strongest LIL and HIL). The statistical approach, on which we will focus here, is more empirical and therefore requires a conceptual framework for interpretation. This approach is often criticized as relying on several assumptions including that the profile variability does not influence profile shapes and that the non-simultaneity of the observations of LIL and HIL are unimportant. Actually as the size of high quality data samples grow these effects become less and less important.
## 2 Baring the Broad Profiles
The collection of moderate resolution ($`\lambda `$/$`\mathrm{\Delta }`$$`\lambda `$$`10^3`$) optical and UV spectra of good quality (S/N $`\stackrel{>}{}`$ 20 in the continuum) has become possible only in recent years thanks to the widespread use of CCD detectors and the unprecedented sensitivity and resolution of the UV spectrographs on board HST. For practical purposes, Hi H$`\beta `$ and Civ$`\lambda `$1549can be considered representative of HIL and LIL respectively. They are also the best lines for statistical studies because they permit comparison in the same sources out to z=1.0. The linearity of response of the currently employed detectors has made possible a reliable correction for emission features contaminating H$`\beta `$ and Civ$`\lambda `$1549, which are: (1) Feii emission; significant Feii<sub>UV</sub> emission contaminates the red wing of Civ$`\lambda `$1549, and has been identified in I Zw 1 by ? and later confirmed by ?; (2) Heii$`\lambda `$4686 emission for H$`\beta `$, and Heii$`\lambda `$1640 + Oiii\]$`\lambda `$1663 for Civ$`\lambda `$1549; (3) \[Oiii\]$`\lambda \lambda `$4959,5007 for H$`\beta `$; (4) narrow component, present in several cases in both H$`\beta `$ and Civ$`\lambda `$1549.
Civ$`\lambda `$1549 often shows a narrow core with FWHM $``$ 1000-2000 km s<sup>-1</sup>, which is systematically broader than the narrow component of H$`\beta `$ (H$`\beta _{\mathrm{NC}}`$). The separation between the broad component Civ$`\lambda `$1549<sub>BC</sub> and the core component is often ambiguous. This core however shows no shift with respect to the AGN rest frame (?), no variations (?), and correlates with \[Oiii\]$`\lambda \lambda `$4959,5007 prominence (?). It can therefore be ascribed to the NLR and considered as the narrow line component of Civ$`\lambda `$1549 (Civ$`\lambda `$1549<sub>NC</sub>). The different width between Civ$`\lambda `$1549<sub>NC</sub> and H$`\beta _{\mathrm{BC}}`$ can be understood in terms of density stratification (?) without invoking an additional emitting region such as the so-called “intermediate line region”. Even if Civ$`\lambda `$1549<sub>NC</sub> fractional intensity is small ($``$ 10%), and in some cases obviously absent, failure to account for Civ$`\lambda `$1549<sub>NC</sub> has led to the erroneous conclusions that FWHM Civ$`\lambda `$1549<sub>BC</sub> $`>`$ FWHM H$`\beta _{\mathrm{BC}}`$ and that the Civ$`\lambda `$1549 peak shows no shift with respect to H$`\beta `$ (?).
## 3 A DIFFERENT BLR STRUCTURE IN RADIO LOUD AND RADIO QUIET AGN?
? made a comparison between Civ$`\lambda `$1549<sub>BC</sub> and H$`\beta _{\mathrm{BC}}`$ for a sample of 52 AGN (31 RL). They presented measures of radial velocity for the blue and red sides of H$`\beta _{\mathrm{BC}}`$ and Civ$`\lambda `$1549<sub>BC</sub> at 5 different values of fractional intensity which provide a quantitative description of the profiles. The reference frame was set by the measured velocity of \[Oiii\]$`\lambda \lambda `$4959,5007$`\lambda `$5007$`\mathrm{\AA }`$ (IZw1 was the only exception). Standard profile parameters like peak shift, FWHM, asymmetry index and curtosis can be extracted from these measures. Representative profiles constructed from the median values of Civ$`\lambda `$1549<sub>BC</sub> and H$`\beta _{\mathrm{BC}}`$ $`v_r`$ are reproduced in the left panels of Fig. 1 and 2 for the radio quiet (RQ) and radio loud (RL) samples respectively.
Civ$`\lambda `$1549<sub>BC</sub> is broader than H$`\beta _{\mathrm{BC}}`$in both RQ and RL samples and it is almost always blueshifted relative to H$`\beta _{\mathrm{BC}}`$. However Fig. 1 and Fig. 2 show significant differences between RQ and RL AGN. In RQ AGN, Civ$`\lambda `$1549<sub>BC</sub> is significantly blueshifted with respect to the source rest frame while H$`\beta _{\mathrm{BC}}`$ is symmetric and unshifted. Contrarily in RL AGN Civ$`\lambda `$1549<sub>BC</sub> is more symmetric, while H$`\beta _{\mathrm{BC}}`$ is shifted to the red at peak intensity and redward asymmetric as well (the median profile corresponds to the type AR,R according to ?). There are two important results which are not displayed in the Figures: (i) in RQ AGN, Civ$`\lambda `$1549<sub>BC</sub> blueshifts are apparently uncorrelated with respect to any H$`\beta _{\mathrm{BC}}`$ line profile parameter and the largest Civ$`\lambda `$1549 blueshifts are associated with the lowest W(Civ$`\lambda `$1549); (ii) in RL AGN, on the contrary, Civ$`\lambda `$1549<sub>BC</sub> and H$`\beta _{\mathrm{BC}}`$ line profile parameters (FWHM and peak shift) appear to be correlated. Asymmetry index of Civ$`\lambda `$1549<sub>BC</sub> and H$`\beta _{\mathrm{BC}}`$, even if not correlated, shows a clear trend toward asymmetries of the same kind (symmetric or redward asymmetric). These findings on Civ$`\lambda `$1549<sub>BC</sub> have been confirmed by other authors (?, save the difference in terminology and line profile decomposition) and especially by an analysis of archival HST/FOS observations which have become publicly available after 1995 (?).
The LIL and HIL emitting regions are apparently de-coupled in at least some RQ sources. The “de-coupling” is well seen in I Zw 1, the prototype Narrow Line Seyfert 1 Galaxy (NLSy1; see Fig. 1). The H$`\beta `$ profile is very narrow, slightly blueward asymmetric and unshifted with respect to the rest frame defined by 21 cm observations, while the Civ$`\lambda `$1549 profile is almost totally blueshifted. At least in the case of I Zw 1 the distinction between LIL and HIL emitting regions appears to be observationally established (it was actually suggested because of the difficulty to explain the relative strengths of LIL and HIL emission using a photoionized “single cloud;” ?). There is a very important zeroth-order result here: since I Zw 1 is a strong Feii<sub>UV</sub> emitter, we have HIL Civ$`\lambda `$1549 and LIL Feii<sub>UV</sub> in the same rest-frame wavelength range. We see that Feii<sub>UV</sub> is obviously unshifted (this can be very well seen by shifting an Feii<sub>UV</sub> template to the peak radial velocity of Civ$`\lambda `$1549). This result disproves models that see an (unknown) wavelength dependent mechanism accounting for the quasar broad line shifts relative to the quasar rest frame.
RL AGN apparently mirror RQ AGN in a curious way: Civ$`\lambda `$1549<sub>BC</sub> is more symmetric, while H$`\beta _{\mathrm{BC}}`$ shows preferentially redshifted profiles and increasingly redward asymmetries. Large peak redshifts ($`v_r\stackrel{>}{}1000`$ km s<sup>-1</sup>, as in the case of OQ 208, Fig. 2) are rarely observed; H$`\beta _{\mathrm{BC}}`$ peak shifts are usually small ($`\mathrm{\Delta }v_r`$/FWHM $``$ 1, median profile of Fig. 2). RL Civ$`\lambda `$1549<sub>BC</sub> and H$`\beta _{\mathrm{BC}}`$ data leave open the possibility that both lines are emitted in the same region. Civ$`\lambda `$1549<sub>BC</sub> shows a red-wing (very evident in the latest, higher S/N spectra analyzed by ?), which cannot be entirely accounted for by Feii<sub>UV</sub>emission. It is interesting to note that superluminal sources with apparent radial velocity $`\beta _{\mathrm{app}}510`$ (whose radio axis is probably oriented close to the line of sight in the sample of ? show very strong Civ$`\lambda `$1549 redward asymmetries, low W(Civ$`\lambda `$1549) and W(H$`\beta _{\mathrm{BC}}`$). This result suggests that redshifts are maximized in RL objects at “face-on” orientation (we assume that any disk is perpendicular to the radio axis).
## 4 NLSy1 NUCLEI ARE NOT A DISJOINT RQ POPULATION
NLSy1 are neither peculiar nor rare. The 8<sup>th</sup> edition of the ? catalogue includes 119 NLSy1 satisfying the defining criterion FWHM Balmer $`\stackrel{<}{}`$ 2000 km s<sup>-1</sup>. They account for $`10`$% of all AGN in the same redshift and absolute magnitude range. Attention toward NLSy1 remained dormant after their identification as a particular class (?) until it was discovered that they may represent $``$1/3–1/2 of all soft X-ray selected Seyfert 1 sources (e. g., ?). NLSy1 are also apparently favored in AGN samples selected on the basis of color. They account for 27% of the RQ ? sample, probably because of an optical continuum that is steeply rising toward the near UV. NLSy1 do not occupy a disjoint region in parameter space. They are at an extremum in the FWHM(H$`\beta `$) vs. R<sub>FeII</sub> (=I(Feii$`\lambda `$4570/I(H$`\beta _{\mathrm{BC}}`$)) and in the “Eigenvector 1” parameter spaces (???). Also, the soft X-ray spectral index $`\mathrm{\Gamma }_{\mathrm{soft}}`$ shows a continuous distribution which includes NLSy1 at the high end (???).
Orientation can easily explain much of the RQ phenomenology observed by ?. I Zw 1 can be considered as an extremum with an accretion disk seen face-on (i$`=0^{}`$), and an outflowing wind observed along the disk axis. We can infer that the opening angle of any Civ$`\lambda `$1549 outflow is probably large (i. e., the wind is quasi spherical) because Civ$`\lambda `$1549 profiles like I Zw 1 are rare. Other NLSy1 show low W(Civ$`\lambda `$1549), and strong Feii<sub>opt</sub>, but do not always show large Civ$`\lambda `$1549 blueshifts (?). Nonetheless, it is still possible that NLSy1 may be structurally different from other RQ AGN. If the soft X-ray excess of NLSy1 is due to high accretion rate, then a slim accretion disk is expected to form (?). Line correlations presented in (?) appear to hold until FWHM(H$`\beta `$)$`\stackrel{<}{}`$ 4000 km s<sup>-1</sup>. For FWHM$`\stackrel{>}{}`$ 4000 km s<sup>-1</sup>line parameters appear to be uncorrelated however it is still unclear it is at present unclear because of the difficulty in measuring weak and broad Feii<sub>opt</sub>sources and/or because of a BLR structural difference. This limit may be related to the possibility of sustaining a particular disk structure and an HIL outflow. A second parameter, independent from orientation is needed to account for the FWHM(H$`\beta _{\mathrm{BC}}`$) vs R<sub>FeII</sub> vs $`\mathrm{\Gamma }_{\mathrm{soft}}`$ sequences (see ?? for a detailed discussion).
## 5 Inferences on BLR Models for RQ AGN
Models developed almost independently of the data through the 80’s and early 90’s. The situation has now changed because of three main developments: (1) the “Eigenvector 1” correlations allow a systematic view of the change in optical emission line properties for different classes of AGN (???); (2) the Civ$`\lambda `$1549 - H$`\beta `$ comparison has yielded direct clues about the structure of the BLR (??) and (3) data collected for RM projects provide a high-sampling description of line variations. For instance, binary black hole scenarios (?) were recently challenged by the failure to detect the radial velocity variations expected from previous observations and model predictions (?).
A model in which Civ$`\lambda `$1549 is emitted by outflowing gas (e.g. a spherical wind) while H$`\beta _{\mathrm{BC}}`$ is emitted in a flattened distribution of gas (observed in a direction that minimizes velocity dispersion) such as an optically thick disk (obscuring the receding half of the Civ$`\lambda `$1549 flow, ?) or at the wind base is immediately consistent with the I Zw 1 data. The big question is whether the results for I Zw 1 can be straightforwardly extended to other RQ AGN.
An accretion disk (AD) provides a high density and high column density medium for Feii production (e. g. ?), and possibly other low ionization lines such as CaII (e. g. ?). AD avoid conflict with the stringent restrictions on line profile smoothness imposed by the first extremely high s/n Balmer line observations (??) of -incidentally- two NLSy1. They place a lower limit of (10<sup>7-8</sup>) on the number of discrete emitters needed to explain the observed profiles.
Winds arise as a natural component of an AD model when the effects of radiative acceleration are properly taken into account (??) or when a hydromagnetic or hydrodynamic treatment is performed (??). A signature of radiative acceleration is provided by observations of “double troughs” in $`\frac{1}{5}`$ of BAL QSOs i. e., of a hump in the absorption profiles of Nv$`\lambda `$1240 and Civ$`\lambda `$1549 at the radial velocity difference between Ly$`\alpha `$ and Nv$`\lambda `$1240, 5900 km s<sup>-1</sup>. Such a feature indicates that Ly$`\alpha `$ photons are accelerating the BAL clouds (?, and references therein). Additional evidence is provided by the radial velocity separation in the narrow absortion components of Ly$`\alpha `$ and Nv$`\lambda `$1240 which show the same $`\mathrm{\Delta }v_r`$ of the two doublet components of Civ$`\lambda `$1549 (e. g. ?).
## 6 The Trouble With Bare Accretion Disks and Bipolar Flows
Relativistic Keplerian disks (??) may explain unusual profile shapes (e.g. double-peaked profiles of Balmer emission lines; ??). Uniform axisymmetric disk models produce double-peaked line profiles with the blue peak stronger than the red peak because of Doppler boosting, a feature that is not always observed in these already rare profiles. To solve this problem, ? and ?, proposed that the lines can originate in an eccentric (i.e. elliptical) disk. Simple disk illumination models can also produce single peaked LIL, provided they are produced at large radii ($`\stackrel{>}{}10^3`$ gravitational radii) or that the disk is observed at small inclination (???). The first of these conditions may be met in all NLSy1 galaxies; both of them seem to be met in I Zw 1.
Aside from NLSy1 sources, there is general disagreement between observations and model predictions for externally illuminated Keplerian disks in a line shift– asymmetry parameter space (?). Only a minority ($`\stackrel{<}{}`$ 10 %) of RL and a handful of RQ AGN show double peaked Balmer lines suggestive of a Keplerian velocity field (??). Double-peakers (e. g. Arp 102B in Fig. 2, FWHM(H$`\beta _{\mathrm{BC}}`$)$`\stackrel{>}{}`$10000 km s<sup>-1</sup>) cannot be like the classical cases because the line widths are much smaller. The peaks often vary out of phase (Arp 102B: ?, 3C 390.3: ?). Double peaks (NGC 1097) (??) or one of the peaks (Pictor A) (?) sometimes appear quite suddenly. Profile variability studies of Balmer lines force us to introduce second order modifications to the basic scheme, such as: eccentric rings and precession (??), inhomogeneities such as orbiting hot spots (?), and warps (?). Not even these epicycles are always capable of explaining the observed variability patterns. Even if elliptical AD models do well in explaining the integrated profiles, they face important difficulties in explaining variability patterns.
A serious problem for AD models of emission lines is emerging from spectropolarimetric observations. If profile broad line shapes are orientation dependent then, in principle, the profile shape in polarized light will depend on the distribution of the scatterers relative to the principal axis and to our line of sight. Early spectropolarimetric results showed a discrepancy with the simple disk + electron-scattering-dominated atmosphere models, which predicted polarization perpendicular to the radio axis. The observed polarization is low, parallel to the disk axis, and shows no statistically significant wavelength dependence (?). Recently ? and ? included double-peakers in their samples, and obtained troublesome results for disk emission models because the polarized H$`\alpha `$ profiles are centrally peaked (?). They investigated the case of disk emission where the scattering particles are located above and below an obscuring torus, along its poles. This “polar scattering model” is successful in explaining the polarized profiles but not the position angle of the polarization vector.
The same polarization studies indicate that the only scenario that can account for both the shapes of the scattered line profiles and the alignment of the optical polarization with the radio jet in wide separation double peakers like Arp 102B (Fig. 2) involves a biconical BLR within an obscuring torus. H$`\alpha `$ photons emitted by clouds participating in a biconical flow are scattered towards the observer by dust or electrons in the inner wall of the surrounding torus. The particular case of biconical outflow was first developed to reproduce observed profiles by ?. This model has been successfully applied to fit observed profiles in double-peaked objects (??).
Double peaked or single blueshifted peak LIL profiles fitted with bi-cone outflow models require that the receding part of the flow is also seen. Self-gravity may be important beyond $`1`$ pc, and the disk may be advection-dominated and optically thin (?). Recent work by ? models the vertical structure of AD and the origin of thermal winds above AD. They not only find that a wind powered by a thermal instability develops in all disks with certain opacity laws but also that in disks dominated by bremsstrahlung radiation, a time-dependent inner hole develops below a critical accretion rate. This scenario provides a natural explanation for transient double-peakers, such as NGC1097 (??), but low accretion rate is a requirement for both advection dominated disks and hole formation.
? explored the idea that the double-peaked emitters represent a geometrical extremum where an outflow is viewed close to pole-on. However, double-peakers are associated with double-lobe radio-sources suggesting that the line of sight has a considerable inclination to the axis of the jets. The problem arises only if the core (pc scale) jets is related to the much larger (100 kpc scale) jets. There is both theoretical (e. g. ?) and observational evidence against this assumption.
## 7 Emission from Clouds Illuminated by an Anisotropic Continuum
Models based on radiative acceleration of optically thick clouds with small volume filling factor gained wide acceptance in the Eighties (? and references therein; see also ?). However, problems with cloud confinements and stability (?) have made them increasingly less frequently invoked to explain observations.
First RM studies on Balmer lines excluded radial, and favored orbital or chaotic motions (e.g. ??). ? applied RM techniques to one of the most extensively monitored objects: NGC 5548. They ruled out radial motions, and found that the Civ$`\lambda `$1549 line variations are broadly consistent with a spherical BLR geometry, in which clouds following randomly inclined circularly Keplerian orbits are illuminated by an anisotropic source of ionizing continuum. A RM result favoring Keplerian motion may be approximately correct also for models in which the emitting gas is not bound, such as a wind, since most of the emission occurs near the base of the flow, when the velocity is still close to the escape velocity which is similar to the Keplerian velocity (?).
## 8 Is the disk + wind model applicable to all AGN?
Emission from a terminal flow can explain the recent observations of Goad et al. 1999 who reported that LIL (Mgii$`\lambda `$2800+ Feii<sub>UV</sub>) in NGC 3516 do not respond to continuum variations which did induce detectable variability in the HIL (Ly$`\alpha `$ and Civ$`\lambda `$1549) lines. Hydromagnetic wind models such as those developed by ? and ? exhibit these basic properties. A two-zone wind provides another scenario for the different origins of LIL and HIL.
Turning to the general population of RL AGN, the predominance of redshifts and redward asymmetric profiles is difficult to explain. Several lines of evidence suggest a significant role of gravitational redshift in RL AGN (?) possibly related to a lower distance (in units of gravitational radii) between BLR and central black hole, which may be systematically more massive in RL than in RQ AGN. If this is the case, then a double zone wind may be present also in RL AGN, since Civ$`\lambda `$1549<sub>BC</sub> is still systematically blueshifted with respect to H$`\beta _{\mathrm{BC}}`$. The “correlation” between Civ$`\lambda `$1549<sub>BC</sub> and H$`\beta _{\mathrm{BC}}`$ parameters could be due to the impossibility of maintaining a radial flow along the disk axis, where a relativistic jet is instead propagating. This will make any HIL outflow possible only at lower latitudes over the disk plane, and therefore will produce more similar H$`\beta _{\mathrm{BC}}`$ and Civ$`\lambda `$1549<sub>BC</sub> profiles.
## 9 CONCLUSION
While observations support emission from an accretion disk and an associated spherical wind in RQ AGN with FWHM(H$`\beta _{\mathrm{BC}}`$) $`\stackrel{<}{}`$ 4000 km s<sup>-1</sup>, there is not enough observational support to warrant the same conclusion for RL AGN (and possibly RQ with FWHM$`>`$4000 km s<sup>-1</sup>), although a disk + wind model is a viable possibility also in this case. Wide separation double peakers (mostly RL) do not provide conclusive evidence in favor of LIL disk emission; rather, there is evidence against disk emission as well as against every other reasonably simple scenario.
###### Acknowledgements.
DD-H acknowledges support through grant IN109698 from PAPIIT-UNAM. PM acknowledges financial support from MURST through grant Cofin 98-02-32, as well as hospitality and support from IA-UNAM. We also acknowledge consistent (cycles5-9 AR and GO) rejection of proposals to continue this work. |
warning/0002/gr-qc0002036.html | ar5iv | text | # Electromagnetic waves in a wormhole geometry
## I Introduction
Since the pioneering article by Morris and Thorne , Lorentzian wormholes have attracted a lot of interest in the literature. Let us recall that a wormhole is a solution to Einstein’s equations that can be understood as a “handle” connecting either two universes or two distant places in the same universe. A number of static and dynamic wormholes have been found both in General Relativity and in alternative theories of gravitation (see and references therein). It has also been shown that under particular circumstances this geometry permits the formation of closed timelike curves (CTC’s) .
One of the most striking features of this field configuration is that it needs matter with negative energy density as a source . Hochberg and Visser, and Ida and Hayward have analyzed in detail the issue of the violation of the energy conditions both in static wormholes and in a completely general (i.e. non-symmetric and time dependent) traversable wormhole . Their analysis shows that the null energy condition (NEC) must be violated in both cases. Since it is believed that all classical matter acting as a source of gravitation satisfies the NEC, only quantum phenomena can, in principle, be responsible for NEC violations. In fact, it has been known for some time that there are states in Quantum Field Theory that may violate the energy conditions . Many examples of systems that can violate these conditions due to quantum effects can be found in ; one such example is the conformally coupled scalar field. Indeed, Hochberg et al. found the first self-consistent wormhole solution in semiclassical General Relativity with quantized conformally coupled scalar matter as a source . Another instance in which quantum effects are essential has been presented recently by Kim and Lee . Using a two-dimensional dilaton-gravity at the one-loop level, they showed that a wormhole can be the final state of an evaporating black hole. This result agrees with previous considerations of Hayward , and might be considered as an example of a process to bring a wormhole into existence. Recently, Krasnikov showed that there exist static wormholes with the vacuum stress-energy tensor of the neutrino, the electromagnetic or the massless sclar field as a source of Einstein’s equations.
In spite of these results, the issue of the existence of large amounts of NEC violating matter has not yet been resolved. However, one can take a complementary view by assuming that wormholes exist and work out possible consequences. This line of reasoning was started by Frolov and Novikov , who studied nontrivial effects arising in a wormhole that interacts with an external electromagnetic or gravitational field. In a subsequent paper , they used wormholes as tools for studying the interior of black holes. Later, Gonzalez-Diaz studied some astrophysical consequences of the existence of these objects . More recently, Cramer et al. considered the unusual features of the lensing of light caused by an object with negative mass (which could be interpreted as one of the mouths of a wormhole). As a natural application of this idea, Torres et al. considered the microlensing caused by wormholes on light coming from distant active galactic nuclei. They showed that these events resemble certain types of Gamma Ray Bursts (GRB’s) and set an upper bound on the negative mass density existing in the universe in the form of wormholes. The GRB’s produced by microlensing by wormholes present a definite feature that differentiate them from those associated with fireballs . Namely, they always appear in pairs called FRED-antiFRED. A subsequent article by Anchordoqui et al investigated profiles of GRB’s from the BATSE database, searching for observable signatures of natural wormholes. They could identify at least one event that may be associated with this mechanism.
From an observational point of view, it is also important to study the propagation of different types of perturbations in the geometry associated with these objects. This study was initiated by Kar and Sahdev , and Kar et al. , who studied the reflection and transmission of massless scalar waves in the presence of an ultrastatic wormhole. In this article we consider the propagation of electromagnetic waves in the same geometry. We will study some properties of the outgoing radiation modified by the gravitational field of the wormhole.
We should mention at this point an article that is at first sight related to ours. In , Clément discussed the problem of the scattering of scalar and electromagnetic waves in an Ellis geometry (which is in fact a wormhole). Following the ideas of Wheeler’s geometrodynamics , his aim was to show that wormholes are particle-like objects. Consequently, the wormhole structure in was to be important only at a microscopic level, while in this work we are considering macroscopic wormholes.
The structure of the paper is as follows. In the next section we transform Maxwell’s equations to a form that is convenient to study the propagation of electromagnetic waves in curved spacetimes. Following this, we introduce the metric of a static wormhole in which we will study the electromagnetic perturbations. Section IV shows how the problem is equivalent to a one-dimensional problem in the presence of a potential barrier. We also present numerical results obtained for the transmission coefficient. Section V investigates the polarization of the outgoing radiation. We conclude with a summary of our results.
## II Maxwell’s equations in a gravitational field
We begin with the equations that govern the propagation of electromagnetic waves in a gravitational background
$`F_{;\nu }^{\mu \nu }=4\pi J^\mu ,F_{\mu \nu ;\sigma }+F_{\nu \sigma ;\mu }+F_{\sigma \mu ;\nu }=0.`$ (1)
where $`J^\mu `$ is the current four-vector. We shall see below that these equations can be rewritten in a more convenient way . In a given coordinate frame such that
$`ds^2=g_{\mu \nu }dx^\mu dx^\nu `$, define
$$H^{\mu \nu }\sqrt{g}F^{\mu \nu },\text{ and }I^\mu \sqrt{g}J^\mu .$$
In terms of these tensors, Maxwell’s equations are given by
$`H_{,\nu }^{\mu \nu }=4\pi I^\mu ,F_{\mu \nu ,\sigma }+F_{\nu \sigma ,\mu }+F_{\sigma \mu ,\nu }=0.`$ (2)
Note that these equations are similar to Maxwell’s equations in flat spacetime. The background geometry has not disappeared, but manifests itself in the constituitive equations $`H^{\mu \nu }=\sqrt{g}g^{\mu \rho }g^{\nu \sigma }F_{\rho \sigma }`$. To be specific, in a Cartesian coordinate system the following decomposition of the tensors is possible:
$$F_{\mu \nu }(\stackrel{}{E},\stackrel{}{B}),H^{\mu \nu }(\stackrel{}{D},\stackrel{}{H}),J^\mu (\rho ,\stackrel{}{J}).$$
(3)
In this notation, Maxwell’s equations explicitly take the form they have in Euclidean spacetime:
$`\stackrel{}{}.\stackrel{}{B}`$ $`=`$ $`0,\stackrel{}{}\stackrel{}{E}={\displaystyle \frac{\stackrel{}{B}}{t}},`$ (4)
$`\stackrel{}{}.\stackrel{}{D}`$ $`=`$ $`0,\stackrel{}{}\stackrel{}{H}={\displaystyle \frac{\stackrel{}{D}}{t}}+4\pi \stackrel{}{j}.`$ (5)
With the decomposition given in Eq.(3) the constituitive relations can be written as
$$D_i=ϵ_{ik}E_k(\stackrel{}{G}\times \stackrel{}{H})_i$$
(6)
and
$$B_i=\mu _{ik}H_k+(\stackrel{}{G}\times \stackrel{}{E})_i,$$
(7)
where
$$ϵ_{ik}=\mu _{ik}=\sqrt{g}\frac{g^{ik}}{g_{00}}$$
(8)
and
$$G_i=\frac{g_{0i}}{g_{00}}.$$
(9)
We can now see that Maxwell’s equations in a gravitational background are formally equivalent to the equations of an electromagnetic field in a flat spacetime in the presence of a medium. Generally, this medium is bi-anisotropic. In the following, we shall restrict our considerations to a medium characterized by diagonal dielectric and magnetic permeabilities :
$$ϵ_{ik}=\mu _{ik}n\delta _{ik}.$$
(10)
In this case, the refraction index $`n`$ corresponds to a static spacetime . Defining the vectors
$$\stackrel{}{F}^\pm \stackrel{}{E}\pm i\stackrel{}{H},\stackrel{}{S}^\pm \stackrel{}{D}\pm i\stackrel{}{B},$$
it is easily shown that Maxwell’s equations can be recast as
$$\stackrel{}{}\stackrel{}{F}^\pm =\pm i\frac{\stackrel{}{S}^\pm }{t},\stackrel{}{}.\stackrel{}{S}^\pm =0.$$
(11)
Using the fact that $`\stackrel{}{D}=ϵ\stackrel{}{E}`$, and $`\stackrel{}{B}=\mu \stackrel{}{H}`$, the first of these equations reduces to
$$\stackrel{}{}\stackrel{}{F}^\pm =\pm in\frac{\stackrel{}{F}^\pm }{t}.$$
(12)
In the next section we describe the background geometry for the wormhole in which the electromagnetic perturbations will be studied.
## III The geometry
As the background metric, we shall adopt that of a static wormhole, given by
$`ds^2=e^{2\varphi (r)}dt^2+{\displaystyle \frac{dr^2}{\left(1\frac{b(r)}{r}\right)}}+r^2d\mathrm{\Omega }^2,`$ (13)
where $`d\mathrm{\Omega }^2`$ is the surface element of $`S^2`$, $`\varphi (r)`$ is the so-called redshift function and $`b(r)`$ is the shape function. Here we shall restrict ourselves to the case $`b(r)=b_0^2/r`$, where $`b_0`$ is the radius of the throat of the wormhole. In order to use the decomposition given by Eq.(3), we need to recast the metric in Cartesian coordinates. First we make the transformation to isotropic coordinates by introducing $`r=f(\rho )`$, where
$`f(\rho )={\displaystyle \frac{4\rho ^2+b_0^2}{4\rho }}.`$ (14)
In these new coordinates, the metric takes the form
$$ds^2=e^{2\varphi (\rho )}dt^2+A^2(\rho )(d\rho ^2+\rho ^2d\mathrm{\Omega }^2),$$
(15)
with
$`A(\rho )={\displaystyle \frac{4\rho ^2+b_0^2}{4\rho ^2}}.`$ (16)
The spatial part of the metric can now be written in Cartesian coordinates by means of the usual definitions $`x^1=\rho \mathrm{sin}\theta \mathrm{cos}\phi `$, $`x^2=\rho \mathrm{sin}\theta \mathrm{sin}\phi `$, $`x^3=\rho \mathrm{cos}\phi `$. With this substitution, the metric becomes
$`ds^2=e^{2\varphi (\rho )}dt^2+A^2(\rho )(\delta _{ij}dx^idx^j).`$ (17)
From Eqs.(8) and (10) it is easy to show that the refraction index is given by
$`n(\rho )={\displaystyle \frac{A(\rho )}{e^{\varphi (\rho )}}}.`$ (18)
It can be seen from the properties of the metric Eq.(13) that far from the wormhole $`n(\rho )`$ tends to 1, in which case we will recover Maxwell’s equations in free space.
## IV The equivalent one-dimensional problem
We demonstrate below how the problem of travelling electromagnetic waves in the geometry given by Eq.(17) can be reduced to a one-dimensional Schrödinger’s equation for a particle with unit mass in a given potential. The following calculations are restricted to the so-called “ultrastatic case”, in which $`\varphi (\rho )0`$.
We begin by developing the Hertz vector $`\stackrel{}{F}^\pm `$ in series of generalized spherical harmonics :
$`\stackrel{}{F}^\pm (\stackrel{}{\rho },t)={\displaystyle \underset{J,M}{}}\stackrel{}{F}_{JM}^\pm (\stackrel{}{\rho },t),`$ (19)
with
$`\stackrel{}{F}_{JM}^\pm (\stackrel{}{\rho },t)={\displaystyle \underset{\lambda =e,m,o}{}}F_{JM}^\pm (\rho ,\omega )\stackrel{}{Y}_{JM}^{(\lambda )}(\widehat{\rho })e^{i\omega t}.`$ (20)
The superindices $`e`$ and $`m`$ refer to transverse modes, while the superindex $`o`$ refers to longitudinal modes. Using the properties of the $`\stackrel{}{Y}_{JM}^{(\lambda )}(\widehat{\rho })`$ , Eq.(11) can be rewritten as
$$\frac{d}{d\rho }\left(\rho F_{JM}^{\pm (m)}\right)=\pm n\omega \rho F_{JM}^{\pm (e)},$$
(21)
$$\frac{d}{d\rho }\left(\rho F_{JM}^{\pm (e)}\right)\sqrt{J(J+1)}F_{JM}^{\pm (o)}=\pm n\omega \rho F_{JM}^{\pm (m)},$$
(22)
and
$$\frac{1}{\rho }\sqrt{J(J+1)}F_{JM}^{\pm (m)}=\pm n\omega F^{\pm (o)}.$$
(23)
At this point, it is convenient to transform $`\rho `$ to a new variable $`x`$ defined by
$`{\displaystyle \frac{dx}{d\rho }}=n(\rho ).`$ (24)
It can be easily seen from the metric given in Eq.(15) that the coordinate $`x`$ is the proper distance, which is given in terms of $`r`$ by the relation
$`x=\pm \sqrt{r^2b_0^2}.`$ (25)
It is also useful to define the functions
$`\chi _{JM}^{\pm (\lambda )}(x,\omega )=\rho (x)F_{JM}^{\pm (\lambda )}(\rho (x),\omega ).`$ (26)
Substituting Eqs.(21) and (23) into Eq.(22) and making the transformation $`z=x/b_0`$ we obtain
$`{\displaystyle \frac{d^2\chi _{JM}^{\pm (m)}}{dz^2}}+\left[k^22U_J(z)\right]\chi _{JM}^{\pm (m)}=0,`$ (27)
with $`k=\omega b_0`$. The potential $`U_J(z)`$, is given by
$$U_J(z)=2J(J+1)\left[\frac{z+\sqrt{1+z^2}}{(z+\sqrt{1+z^2})^2+1}\right]^2.$$
(28)
Note that the potential is asymmetric and tends asymptotically to $`0`$ as $`z\mathrm{}`$. This is illustrated in the following figure.
The potential barrier vanishes for $`J=0`$, and grows like $`J^2`$ for large $`J`$. Notice that $`U_J(0)=J(J+1)/2`$ is a good estimate for the maximum of the barrier. Due to the intricate dependence of the potential on $`z`$, it was not possible to find an exact analytical expression for the transmission coefficient, $`T_J`$. Instead, numerical methods based on the work presented in the appendix of were employed. Our results for the transmission coefficient versus $`k`$, for the values $`J=1,3\text{ and }5`$ are displayed in the following figure.
$`T_J`$ gives the amount of radiation that goes through the wormhole. The lower part of the plot shows a well-known feature of this type of scattering: the heigher the potential barrier, the slower $`T_J`$ grows. However, there is a rapid increment of the transmission coefficient within a small interval of $`k`$ until it reaches its maximum for every value of $`J`$.
## V Polarization of the outgoing waves
Let us now discuss the polarization state of the outgoing waves that passed through the wormhole’s throat. From the equations of motion Eq.(11) it can be seen that the left-circularly-polarized radiation ($`\stackrel{}{F}^+=0`$) is decoupled from the right-circularly-polarized photons, so the helicity is conserved. This remains true even in the case of a stationary spacetime, in which $`\stackrel{}{G}0`$. Thus the polarization state of circularly-polarized photons is not altered by the scattering in gravitational field of a wormhole . Note that this result is valid for any static spacetime (see Eq.(8)) and consequently, helicity is conserved for any static wormhole and not just for the one described by Eq.(17).
The case of linearly polarized waves was studied by Mashhoon for a Schwarzschild black hole . His calculations can be easily adapted to the present case, but are rather lengthy. Instead, we follow the approach developed in . The metric of any stationary spacetime can be written as
$$ds^2=h(dt^2G_idx^i)^2dl^2,$$
(29)
with
$$dl^2=\left(g_{ij}+\frac{g_{0i}g_{0j}}{g_{00}}\right)dx^idx^j.$$
(30)
It was shown in that the rotation in the polarization plane of light along the path between the source and the observer in the above metric (i.e. the so-called Faraday effect) is given by
$$\mathrm{\Omega }=\frac{1}{2}_{\mathrm{sou}.}^{\mathrm{obs}.}\sqrt{h}\stackrel{}{B}_g.\stackrel{}{dl},$$
(31)
where $`\stackrel{}{B}_g`$ is the “gravitomagnetic” vector given by
$$\stackrel{}{B}_g=\stackrel{}{}\stackrel{}{G}.$$
(32)
In the case of a static wormhole, $`\stackrel{}{G}0`$ and consequently there is no change in the linearly polarized light.
## VI Conclusions
We have studied here several aspects of the transit of electromagnetic waves through an ultrastatic wormhole. We have shown that the problem can be rewritten in such a way that the curved background geometry is replaced by a medium whose properties are summarized by the refraction index $`n(\rho )`$. We also showed that the magnetic modes of the electromagnetic field obey a one-dimensional Schrödinger-like equation, with an asymmetric barrier-type potential. The transmission coefficient was calculated numerically, and it exhibits some features common to barrier-like potentials. Namely, in a small interval of $`k`$, the curve rapidly increases until there is no reflection and for higher values of $`J`$, initially the curve rises more slowly.
We have also demonstrated that the interaction of the radiation with the gravitational field of the wormhole cannot change the polarization state of the radiation. Moreover, this was shown for any static wormhole and not only with the particular case described by metric Eq.(17).
It would be interesting to search for special features of the transmission coefficient in more general types of wormholes. In particular, the existence of bound states (i.e.standing waves) for scalar waves was shown in . If resonances with the same characteristics also exist in the case of electromagnetic waves, their position in what could be considered as the emission spectrum of the wormhole may give us information on the size of the throat and its shape .
One can expect a greater richness of observational consequences if the background wormhole geometry is rotating . For instance, differential gravitational scattering of polarized radiation can be caused by the helicity-rotation coupling . Finally, to complete the study of perturbations in these spacetimes, it would be necessary to analyze spinor and gravitational perturbations. We intend to address these problems in future works.
## VII Acknowledgments
We thank S. Joffily, S. Kar, D. Monteoliva, M. Novello, J. Salim, and F. Zyserman for helpful discussions. S.E. Perez Bergliaffa acknowledges financial support from Conicet (Argentina) and K. Hibberd is supported by CNPq (Conselho Nacional de Desenvolvimento Científico e Tecnológico). |
warning/0002/hep-ph0002082.html | ar5iv | text | # 1 Introduction
## 1 Introduction
A substantial part of current particle phenomenology is based upon our ability to study efficiently processes involving a relatively large number of particles. This requires efficient algorithms for matrix element calculation and phase space generation and integration. As far as the matrix element calculation is concerned, the traditional approach utilizes standard Feynman graph representation of the scattering amplitude, resulting to a computational cost that grows asymptotically as $`n!`$, where $`n`$ is the number of particles involved in the process. As an alternative recursive algorithms based on Schwinger-Dyson equations lead asymptotically to a much lower growth of the computational cost, namely $`a^n`$, where $`a3`$. It is the aim of this paper to present such an algorithm as well as the corresponding FORTRAN code which allows the computation of any tree-order electroweak amplitude.
## 2 The algorithm
As advertised, the algorithm is based upon the well-known Dyson-Schwinger equations. These equations are equivalent to the field equations derived from the classical lagrangian and express recursively the $`n`$-point Green’s functions in terms of the $`1,2,\mathrm{},(n1)`$-point functions.
In order to illustrate this point let us consider a simple example of a QED-like theory possessing minimal interactions of a spinor field to a gauge boson. Let $`p_1,p_2,\mathrm{},p_n`$ represent the external momenta involved in the scattering process under consideration, taken to be incoming. For a vector field we define
$$b_\mu (P)=\text{}$$
(1)
a four vector, which describes any sub-amplitude from which a boson $`V`$ with momentum $`P`$ can be constructed. The momentum $`P`$ is given as a sum of momenta of external particles, namely $`P=_{iI}p_i`$ with $`I\{1,\mathrm{},n\}`$. Accordingly, we define by
$$\psi (P)=\text{}$$
(2)
a four-dimensional spinor, which describes any sub-amplitude from which a fermion with momentum $`P`$ can be constructed and by
$$\overline{\psi }(P)=\text{}$$
(3)
a four-dimensional antispinor, which describes any sub-amplitude from which an antifermion with momentum $`P`$ can be constructed.
The Dyson-Schwinger equations take the following simple form. For a boson with momentum $`P^\mu `$,
$`=`$ $`\text{}+\text{}`$
$`b^\mu (P)`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{n}{}}}\delta _{P=p_i}b^\mu (p_i)+{\displaystyle \underset{P=P_1+P_2}{}}(ig)\mathrm{\Pi }_\nu ^\mu \overline{\psi }(P_2)\gamma ^\nu \psi (P_1)ϵ(P_1,P_2)`$ (4)
where the propagator $`\mathrm{\Pi }`$ is given by
$$\mathrm{\Pi }_\nu ^\mu =\frac{i}{P^2m^2}(g_\nu ^\mu +\frac{P^\mu P_\nu }{P^2\xi m^2}(1\xi ).)$$
and $`ϵ(P_1,P_2)`$ is a sign factor, that takes the value $`\pm 1`$ and whose definition will be described below. For a fermion with momentum $`P`$
$`=`$ $`\text{}+\text{}`$
$`\psi (P)`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{n}{}}}\delta _{P=p_i}\psi (p_i)+{\displaystyle \underset{P=P_1+P_2}{}}(ig)𝒫b/(P_2)\psi (P_1)ϵ(P_1,P_2)`$ (5)
where $`𝒫`$ is given by
$$𝒫=\frac{i}{P/m}$$
and for an antifermion with momentum $`P`$,
$`=`$ $`\text{}+\text{}`$
$`\overline{\psi }(P)`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{n}{}}}\delta _{P=p_i}\overline{\psi }(p_i)+{\displaystyle \underset{P=P_1+P_2}{}}(ig)\overline{\psi }(P_1)b/(P_2)\overline{𝒫}ϵ(P_1,P_2)`$ (6)
where
$$\overline{𝒫}=\frac{i}{P/m}.$$
The scattering amplitude can be calculated by any of the following relations,
$$𝒜(p_1,\mathrm{},p_n)=\{\begin{array}{cc}b_0^\mu (P_i)b_\mu (p_i)\hfill & \text{where }i\text{ corresponds to a photon}\hfill \\ \overline{\psi }_0(P_i)\psi (p_i)\hfill & \text{where }i\text{ corresponds to an incoming fermion line}\hfill \\ \overline{\psi }(p_i)\psi _0(P_i)\hfill & \text{where }i\text{ corresponds to an outgoing fermion line}\hfill \end{array}$$
(7)
where
$$P_i=\underset{ji}{}p_j,$$
so that $`P_i+p_i=0`$. The functions with subscript $`0`$ are given by the previous expressions, except for the propagator term. This is because the outgoing momentum $`P_i`$ being on shell the propagator factor is removed by the amputation procedure. The initial conditions are given by
$`b^\mu (p_i)`$ $`=`$ $`ϵ_\lambda ^\mu (p_i),\lambda =\pm 1,0`$
$`\psi (p_i)`$ $`=`$ $`\{\begin{array}{cc}u_\lambda (p_i)\hfill & \text{if }p_i^00\hfill \\ v_\lambda (p_i)\hfill & \text{if }p_i^00\hfill \end{array}`$ (10)
$`\overline{\psi }(p_i)`$ $`=`$ $`\{\begin{array}{cc}\overline{u}_\lambda (p_i)\hfill & \text{if }p_i^00\hfill \\ \overline{v}_\lambda (p_i)\hfill & \text{if }p_i^00\hfill \end{array}`$ (13)
where the explicit form of $`ϵ_\lambda ^\mu ,u_\lambda ,v_\lambda ,\overline{u}_\lambda ,\overline{v}_\lambda `$ are given in the Appendix.
In order to actually solve the recursive equations it is convenient to use a binary representation of the momenta involved . For a process involving $`n`$ external particles with momenta $`p_i^\mu ,i=1\mathrm{},n`$ we assign to the momentum $`P^\mu `$ defined as
$$P^\mu =\underset{iI}{}p_i^\mu $$
where $`I\{1,\mathrm{},n\}`$, a binary vector $`\stackrel{}{m}=(m_1,\mathrm{},m_n)`$, where its components take the values $`0`$ or $`1`$ in such a way that
$$P^\mu =\underset{i=1}{\overset{n}{}}m_ip_i^\mu .$$
Moreover this binary vector can be uniquely represented by the integer
$$m=\underset{i=1}{\overset{n}{}}2^{i1}m_i$$
with $`0m2^{n1}`$. Especially, the external momenta $`p_i`$, $`i=1\mathrm{},n`$ are represented by $`m=2^{i1},i=1,\mathrm{},n`$. All momenta $`P`$ can now be replaced by the corresponding integers
$$b_\mu (P)b_\mu (m).$$
A very convenient ordering of integers in binary representation relies on the notion of level $`l`$, defined simply as
$$l=\underset{i=1}{\overset{n}{}}m_i.$$
As it is easily seen all external momenta are of level $`1`$, whereas the total amplitude corresponds to the unique level $`n`$ integer $`2^{n1}`$. This ordering dictates the natural path of the computation; starting with level-$`1`$ sub-amplitudes, we compute the level-$`2`$ ones using the Dyson-Schwinger equations and so on up to the final expression Eq.(13). As far as the sign factor is concerned
$$ϵ(P_1,P_2)ϵ(m_1,m_2)$$
we define
$$ϵ(m_1,m_2)=(1)^{\chi (m_1,m_2)}$$
(14)
with
$$\chi (m_1,m_2)=\underset{i=n}{\overset{2}{}}\widehat{m}_{1i}\left(\underset{j=1}{\overset{i1}{}}\widehat{m}_{2j}\right)$$
(15)
where hated components are set to $`0`$ if the corresponding external particle is a boson. It is not difficult to see that this sign factor takes properly into account the anti-symmetry of the amplitude with respect to fermionic particles.
As an illuminating example we present here the computation of the amplitude for the process $`e^+e^{}\mu ^+\mu ^{}\gamma `$. The calculation starts with the computation of the level-$`2`$ sub-amplitudes, which are all sub-amplitudes produced from any two of the external particles. These are
$`b_\mu (12)`$ $`=`$ $`(ig)\mathrm{\Pi }_{12\mu }^\nu \overline{\psi }(4)\gamma _\nu \psi (8)`$
$`\overline{\psi }(18)`$ $`=`$ $`(ig)\overline{\psi }(2)b/(16)\overline{𝒫}_{18}`$
$`\overline{\psi }(20)`$ $`=`$ $`(ig)\overline{\psi }(4)b/(16)\overline{𝒫}_{20}`$
$`\psi (24)`$ $`=`$ $`(ig)𝒫_{24}b/(16)\psi (8)`$
Then we compute the level-$`3`$ sub-amplitudes:
$`\overline{\psi }(14)`$ $`=`$ $`(ig)\overline{\psi }(2)b/(12)\overline{𝒫}_{14}`$
$`b_\mu (28)`$ $`=`$ $`(ig)\mathrm{\Pi }_{28\mu }^\nu \left(\overline{\psi }(20)\gamma _\nu \psi (8)+\overline{\psi }(4)\gamma _\nu \psi (24)\right)`$
and then the level-$`4`$ sub-amplitude, which is the final level in this case,
$$\overline{\psi }_0(30)=(ig)\left(\overline{\psi }(2)b/(28)+\overline{\psi }(14)b/(16)+\overline{\psi }(18)b/(12)\right)$$
(16)
Then the amplitude is simply given by
$$𝒜=\overline{\psi }_0(30)\psi (1)$$
Note that we have chosen the particle number $`1`$ as our ending point, so we have computed all sub-amplitudes where the momentum $`p_1`$ does not appear: this excludes all odd integers between $`1`$ and $`2^{n2}`$.
There are two special issues in the computation which go beyond the recursive equations presented above. The first is the generation of all helicity configurations and the second the treatment of the color summation in the case external particles with color are involved.
As far as the helicity configurations are concerned this is done in a rather straightforward way, resulting to an automatic evaluation of all relevant combinations. For any given external particle knowledge of its flavor allows the program to compute all relevant helicity configurations. The total number is given by $`2^{l_2}3^{l_3}`$ where $`l_2`$ is the number of (anti-)fermions and massless gauge bosons and $`l_3`$ the number of massive gauge bosons involved in the scattering process.
The color configurations are taken into account as follows. Since only electroweak processes are considered in this version, the only colored particles that can appear in the amplitude are quarks and antiquarks and let us have $`n`$ pairs of them. Each color amplitude is proportional to the following color structure
$$𝒞_i=\delta _{1,\sigma _i(1)}\delta _{2,\sigma _i(2)}\mathrm{}\delta _{n,\sigma _i(n)}$$
where $`\sigma _i`$ represents the $`i`$-th permutation of the set $`1,2,\mathrm{},n`$. The code computes all non-vanishing color amplitudes as well as the corresponding color matrix
$$_{ij}=C_iC_j$$
and finally performs the color summation.
## 3 The code HELAC
The computation is split in two major phases. During the first phase, which we call initialization phase, the program selects all the relevant sub-amplitudes for the required process. In the simplest version the program accepts as input the following variables:
* n the number of particles involved in the scattering
* ifl(1:n) the array of flavors of the particles
* iflag. If set to 0 sum over all helicity configurations is understood. If set to 1 you must supply also the specific helicity configuration to be computed.
* iunitary. If set to $`0`$ the Feynman gauge is considered whereas if set to $`1`$ the unitary gauge is used.
* ihiggs denotes the inclusion (1) or not (0) of the Higgs particle as an intermediate state.
* iwidth denotes the fixed (0) or complex (1) scheme for the introduction of the width of $`W`$ and $`Z`$<sup>1</sup><sup>1</sup>1For a recent discussion see reference ..
* io(1:n) is used to distinguish among initial (1) and final state particles (-1). By default io(1:2)=1 and io(3:n)=-1.
The first routine called after the input has been read is the routine physics/physics.f located in the homonymous file. In this routine all couplings of the standard electroweak theory are defined .
Then the routine helac\_init/master.f is called. In the beginning the average-over-helicity (avhel), the average-over-color (avcol) and the symmetry(symet) factors are computed. Then, by calling the routine setncc/intpar.f the number of color configurations ncc is set up depending on the number of quarks and anti-quarks involved. For each color configuration the program constructs the skeleton of the amplitude starting at level two and proceeding up to level $`n1`$. All possible vertices are scanned in this phase, as for instance vff/pan1.f which describes the coupling of a fermion anti-fermion to produce a vector boson, in order to select all non-zero sub-amplitudes. A special routine,redo/pan1.f is called at the end just to check whether all the selected sub-amplitudes are indeed contributing to the final amplitude under consideration. The program ends this phase when all color configurations have been considered and the color matrix rmatrix(1:ncc,1:ncc) has been computed.
During the second phase, which we call the computation phase, the code computes numerically the amplitude for each phase space point introduced. The main calling routine is helac\_master/master.f. The first step is the automatic setup of the helicity configurations via the routine sethel/intpar.f. The next step is the computation of the external particle wave functions by the routine iniqq/pan1.f. All relevant routines for this computation can be found in wavef.f. The program proceeds by computing the amplitude for every color configuration via nextq/pan1.f. The vertex and propagator functions needed are included in pan2.f. Moreover the sign factor as described in Eq.(15) is also computed. The routine helac\_master/master.f returns the squared matrix element (smel2)summed and averaged over helicity and color. Finally, a sample main/main.f program to run HELAC is provided. In this part of the code a call to a phase-space generator is also included (rambo/rambo.f), in order to generate appropriate phase space points for the computation.
All floating point computations are performed in real\*8 or double precision. Nevertheless the program is written in such a way that one can choose a higher accuracy if needed. There are two alternatives. The first one is to use real\*16 or quadruple precision. For this task several modules to perform quadruple precision computation with complex numbers are supplied in qprec.f. This is necessary because in most of the existing platforms no FORTRAN quadruple precision complex numbers are available. The most efficient way to use this precision with complex numbers is to exploit the virtues of FORTRAN 90 and define the appropriate derived types and modules. The second alternative make use of a multi-precision library , written also in FORTRAN 90. The precision is now user-defined and is set up by calling the routine zmset in the very beginning of the program (main.f). All these are automatically driven by the appropriate make files included in the package.
## 4 Results
In this section we will try to explain how to read the output of the program and to highlight several aspects of it. First of all let us follow a sample computation of the process
$$e^{}e^+e^{}e^+\gamma .$$
As an input we have to provide
1. The number of particles involved in the process (5).
2. For each particle its flavor (2 -2 2 -2 31). In the following table we provide the correspondence used in HELAC:
| $`\nu _e,e^{},u,d,\nu _\mu ,\mu ^{},c,s,\mathrm{}`$ | $`1,\mathrm{},12`$ |
| --- | --- |
| $`\overline{\nu _e},e^+,\overline{u},\overline{d},\overline{\nu }_\mu ,\mu ^+,\overline{c},\overline{s},\mathrm{}`$ | $`1,\mathrm{},12`$ |
| $`\gamma ,Z,W^+,W^{}`$ | $`31,\mathrm{},34`$ |
| $`H,\chi ,\varphi ^+,\varphi ^{}`$ | $`41,\mathrm{},44`$ |
3. The parameters iflag, iunitary, ihiggs, iwidth (0 1 0 0).
The expected output is as follows
where the first lines show the required input data; then the average-helicity, average-color and symmetry factors are shown; the number of contributing sub-amplitudes for each color configuration is given (44) as well as the total number of Feynman graphs (16). This marks the end of the initialization phase, and afterwards the computation of the squared matrix element for a given phase space point is performed. Results for each helicity configuration are printed and at the end the HELAC amplitude squared is given (1.107657409535622E-04)
As far as the efficiency of the code is concerned extensive comparisons have been made to existing calculations. We have chosen among them two popular tools, namely MADGRAPH and EXCALIBUR . Comparisons with EXCALIBUR have been restricted to four-fermion final states . Running under the same conditions the speed ratio was varying between $`1`$ and $`2`$. The same results have been obtained in comparisons with MADGRAPH taken into account that it can produce results for only up to $`7`$ particles involved in the scattering process. As expected HELAC show an exponential (instead of factorial) CPU-time growth. There is no a priori limitation for the number of particles that HELAC can treat, the only restrictions being that of memory allocation.
In order to have a taste of a multi-precision computation we have computed the squared amplitude for the process
$$e^{}(p_1)e^+(p_2)e^{}(p_3)e^+(p_4)e^{}(p_5)e^+(p_6)$$
at two phase space points. Phase space point (A) is just a randomly generated one by the phase-space generator RAMBO. Phase space point (B) on the other hand is a very special one, where
$`p_1^0/\text{ GeV}=100,\stackrel{}{p_1}+\stackrel{}{p_2}=0,p_3^0/p_1^0=0.9,\theta _3=0,`$
$`(p_5+p_6)^2/(p_4+p_5+p_6)^2=0.1,\theta _4=0,\varphi _4=0,\theta _5=0,\varphi _5=0`$
and $`m_e=0.511\times 10^3\text{ GeV}`$. In this case all particles are co-linear to the beam and therefore important enhancements are expected. In the following table results are provided for these points, by using the real\*8 (DP), real\*16 (QP) and 34-digit multi-precision version of the code.
(A) (B) 1.539728523150595E-008 1.256276706229023E+023 1.53972852315058854156763002825013D-08 3.07162601093710915134136924973089D+22 1.53972852315058854156763002825011853M-8 3.07162601093710915127950109241770808M+22
As one can easily see the numerical stability of the DP computation is spoiled in the co-linear region. This is due to huge gauge cancellations occurring in this region; a way out of this problem will be discussed elsewhere.
We end this presentation by summarizing the main achievements:
* An algorithm based on Dyson-Schwinger equations has been presented that enables the computation of electroweak amplitudes with high efficiency.
* Based on this algorithm a FORTRAN package HELAC has been developed, which includes full massive computation in both the unitary and the Feynman gauge of the standard electroweak theory.
* A quadruple as well as a multi-precision version of HELAC have been incorporated allowing numerically stable results for any phase space point and for arbitrary high energies.
Appendix
In all calculations we are using the light-cone representation of a four-vector $`V^\mu `$, defined as
$$V^A=(V^0+V_z,V^0V_z,V_x+iV_y,V_xiV_y),A=1,\mathrm{},4.$$
(17)
Polarization state-vectors are given by
$`ϵ_{}^A`$ $`=`$ $`({\displaystyle \frac{p_T}{\sqrt{2}|\stackrel{}{p}|}},{\displaystyle \frac{p_T}{\sqrt{2}|\stackrel{}{p}|}},{\displaystyle \frac{(p_x+ip_y)(|\stackrel{}{p}|+p_z)}{\sqrt{2}|\stackrel{}{p}|p_T}},{\displaystyle \frac{(p_xip_y)(|\stackrel{}{p}|+p_z)}{\sqrt{2}|\stackrel{}{p}|p_T}})`$
$`ϵ_+^A`$ $`=`$ $`({\displaystyle \frac{p_T}{\sqrt{2}|\stackrel{}{p}|}},{\displaystyle \frac{p_T}{\sqrt{2}|\stackrel{}{p}|}},{\displaystyle \frac{(p_x+ip_y)(|\stackrel{}{p}|p_z)}{\sqrt{2}|\stackrel{}{p}|p_T}},{\displaystyle \frac{(p_xip_y)(|\stackrel{}{p}|p_z)}{\sqrt{2}|\stackrel{}{p}|p_T}})`$
$`ϵ_0^A`$ $`=`$ $`({\displaystyle \frac{|\stackrel{}{p}|}{\sqrt{p^2}}}+{\displaystyle \frac{p_zp_0}{|\stackrel{}{p}|\sqrt{p^2}}},{\displaystyle \frac{|\stackrel{}{p}|}{\sqrt{p^2}}}{\displaystyle \frac{p_zp_0}{|\stackrel{}{p}|\sqrt{p^2}}},{\displaystyle \frac{(p_x+ip_y)p_0}{|\stackrel{}{p}|\sqrt{p^2}}},{\displaystyle \frac{(p_xip_y)p_0}{|\stackrel{}{p}|\sqrt{p^2}}})`$ (18)
As for the Dirac matrices we are using the chiral representation <sup>2</sup><sup>2</sup>2For conventions see reference .. The wave functions which describe massive spinors are given by:
$`u_+(p)=\left(\begin{array}{c}r/c\\ a(p_x+ip_y)/r\\ mb/r\\ m(p_x+ip_y)/r\end{array}\right)`$ $`\overline{u}_+(p)=\left(\begin{array}{c}mb/r\\ m(p_xip_y)/r\\ r/c\\ a(p_xip_y)/r\end{array}\right)`$ (27)
$`u_{}(p)=\left(\begin{array}{c}m(p_xip_y)/r\\ mb/r\\ a(p_xip_y)/r\\ r/c\end{array}\right)`$ $`\overline{u}_{}(p)=\left(\begin{array}{c}a(p_x+ip_y)/r\\ r/c\\ m(p_x+ip_y)/r\\ mb/r\end{array}\right)`$ (36)
$`v_+(p)=\left(\begin{array}{c}m(p_xip_y)/r\\ mb/r\\ a(p_xip_y)/r\\ r/c\end{array}\right)`$ $`\overline{v}_+(p)=\left(\begin{array}{c}a(p_x+ip_y)/r\\ r/c\\ m(p_x+ip_y)/r\\ mb/r\end{array}\right)`$ (45)
$`v_{}(p)=\left(\begin{array}{c}r/c\\ a(p_x+ip_y)/r\\ mb/r\\ m(p_x+ip_y)/r\end{array}\right)`$ $`\overline{v}_{}(p)=\left(\begin{array}{c}mb/r\\ m(p_xip_y)/r\\ r/c\\ a(p_xip_y)/r\end{array}\right)`$ (54)
where:
$$a=p_0+|\stackrel{}{p}|,b=p_z+|\stackrel{}{p}|,c=2|\stackrel{}{p}|,r=\sqrt{abc}$$
For a massless particle the spinors are
$`u_R(p)=\left(\begin{array}{c}\sqrt{p_0+p_z}\\ (p_x+ip_y)/\sqrt{p_0+p_z}\\ 0\\ 0\end{array}\right)`$ $`\overline{u}_R(p)=\left(\begin{array}{c}0\\ 0\\ \sqrt{p_0+p_z}\\ (p_xip_y)/\sqrt{p_0+p_z}\end{array}\right)`$ (63)
$`u_L(p)=\left(\begin{array}{c}0\\ 0\\ (p_xip_y)/\sqrt{p_0+p_z}\\ \sqrt{p_0+p_z}\end{array}\right)`$ $`\overline{u}_L(p)=\left(\begin{array}{c}(p_x+ip_y)/\sqrt{p_0+p_z}\\ \sqrt{p_0+p_z}\\ 0\\ 0\end{array}\right)`$ (72)
We now proceed to describe the vertex functions. Let us take as an example the Eq.(4). By using the Dirac matrices in the chiral representation and the light-cone expression of four vectors, the reduced form of this equation becomes very simple, namely the four vector
$$V_\mu =\overline{\psi }(P_1)\gamma _\mu \left(g_R\omega _R+g_L\omega _L\right)\psi (P_2)$$
turns out to be
$$V^A=\left(\begin{array}{c}g_R\psi _1\overline{\psi }_3g_L\psi _4\overline{\psi }_2\\ g_R\psi _2\overline{\psi }_4g_L\psi _3\overline{\psi }_1\\ g_R\psi _2\overline{\psi }_3+g_L\psi _4\overline{\psi }_1\\ g_R\psi _1\overline{\psi }_4+g_L\psi _3\overline{\psi }_2\end{array}\right)$$
(73)
where $`\psi _i(\overline{\psi }_i),i=1,\mathrm{},4`$ are the components of the spinor $`\psi (P_2)\left(\overline{\psi }(P_1)\right)`$ and
$$\omega _L=\frac{1}{2}\left(1\gamma _5\right),\omega _R=\frac{1}{2}\left(1+\gamma _5\right).$$
(74)
On the other hand, the spinor
$$u=(P/+m)b/(P_1)\omega _R\psi (P_2)$$
used for instance in Eq.(5), can be reduced to
$$u=\left(\begin{array}{c}(b_2p_1+b_3p_4)\psi _1+(b_4p_1b_1p_4)\psi _2\\ (b_3p_2b_2p_3)\psi _1+(b_1p_2+b_4p_3)\psi _2\\ m(b_2\psi _1b_4\psi _2)\\ m(b_3\psi _1+b_1\psi _2)\end{array}\right)$$
(75)
where $`\psi _i,b_i,p_i,i=1,\mathrm{},4`$ are the components of $`P`$, $`b(P_1)`$ and $`\psi (P_2)`$ respectively. In a similar way, for the standard electroweak theory in both the unitary and the Feynman gauge, twenty eight different vertex functions have been implemented in vertices/pan2.f.
Acknowledgements
Helpful discussions with F. A. Berends, A. P. Chapovsky and R. Pittau are kindly acknowledged. |
warning/0002/math0002145.html | ar5iv | text | # Absolutely singular dynamical foliations
## Introduction
Let $`A_2`$ be the automorphism of the $`2`$-torus, $`𝐓^2=𝐑^2/𝐙^2,`$ given by $`\left(\begin{array}{cc}2& 1\\ 1& 1\end{array}\right).`$ Let $`A_3`$ be the automorphism of the $`3`$-torus $`𝐓^3=𝐑^3/𝐙^3`$ given by $`\left(\begin{array}{cc}A_2& 0\\ 0& 1\end{array}\right).`$ Let $`\text{Diff}_\mu ^2(𝐓^3)`$ be the set of $`C^2`$ diffeomorphisms of $`𝐓^3`$ that preserve Lebesgue-Haar measure $`\mu `$.
In \[SW1\], M. Shub and A. Wilkinson prove the following theorem.
Theorem: Arbitrarily close to $`A_3`$ there is a $`C^1`$-open set $`U\text{Diff}_\mu ^2(𝐓^3)`$ such that for each $`gU`$,
1. $`g`$ is ergodic.
2. There is an equivariant fibration $`\pi :𝐓^3𝐓^2`$ such that $`\pi g=A_2\pi `$ The fibers of $`\pi `$ are the leaves of a foliation $`𝒲_g^c`$ of $`𝐓^3`$ by $`C^2`$ circles. In particular, the set of periodic leaves is dense in $`𝐓^3`$.
3. There exists $`\lambda ^c>0`$ such that, for $`\mu `$-almost every $`w𝐓^3`$, if $`vT_w𝐓^3`$ is tangent to the leaf of $`𝒲_g^c`$ containing $`w`$, then
$$\underset{n\mathrm{}}{lim}\frac{1}{n}\mathrm{log}T_wg^nv=\lambda ^c.$$
4. Consequently, there exists a set $`S𝐓^3`$ of full $`\mu `$-measure that meets every leaf of $`𝒲_g^c`$ in a set of leaf-measure $`0`$. The foliation $`𝒲_g^c`$ is not absolutely continuous.
Additionally, it is shown that the diffeomorphisms in $`U`$ are nonuniformly hyperbolic and Bernoullian. In this note, we prove:
Theorem I: Let $`g`$ satisfy conclusions 1.–3. of the previous theorem. Then there exist $`S𝐓^3`$ of full $`\mu `$-measure and $`k𝐍`$ such that $`S`$ meets every leaf of $`𝒲_g^c`$ in exactly $`k`$ points. The foliation $`𝒲_g^c`$ is absolutely singular.
Remark: In A. Katok’s example of an absolutely singular foliation in \[Mi\], the leaves of the foliation meet the set of full measure in one point. In the \[SW1\] examples, the set $`S`$ may necessarily meet leaves of $`𝒲_g^c`$ in more than one point, as the following argument of Katok’s shows.
It follows from Theorem II in \[SW2\] that for $`k𝐙_+`$ and for small $`a,b>0`$, the map $`g=j_{a,k}h_b`$ satisfies the hypotheses of Theorem I, where
$$h_b(x,y,z)=(2x+y,x+y,x+y+z+b\mathrm{sin}2\pi y),\text{and}$$
$$j_{a,k}(x,y,z)=(x,y,z)+a\mathrm{cos}(2\pi kz)(1+\sqrt{5},2,0).$$
For $`k𝐍`$, let $`\rho _k`$ be the vertical translation that sends $`(x,y,z)`$ to $`(x,y,z+\frac{1}{k})`$. Note that $`h_b\rho _k=\rho _kh_b`$ and $`j_{a,k}\rho _k=\rho _kj_{a,k}`$. Thus $`g\rho _k=\rho _kg`$.
The fibration $`\pi :𝐓^3𝐓^2`$ was obtained in \[SW1\] by using the persistence of normally hyperbolic submanifolds under perturbations. In the present case the symmetries $`\rho _k`$ preserve the fibers of the trivial fibration $`P:𝐓^3𝐓^2`$ from which one starts, and also the maps $`g`$. Therefore the fibers of $`\pi :𝐓^3𝐓^2`$ (i.e., the leaves of center foliation $`𝒲_g^c`$) are invariant under the action of the finite group $`<\rho _k>`$.
Let $`S`$ be the (full measure) set of points in $`𝐓^3`$ for which the center direction is a positive Lyapunov direction (i.e. for which conclusion 3 holds). Since $`\rho _k(𝒲_g^c)=𝒲_g^c`$, it follows that $`\rho _kS=S`$. If $`pS𝒲^c(p)`$, then $`\rho _k(p)\rho _k(S)\rho _k(𝒲^c(p))=S𝒲^c(p)`$; that is, $`S𝒲^c(p)`$ contains at least $`k`$ points.
Thus Theorem I is “sharp” in the sense that we cannot say more about the value of $`k`$ in general. We see no reason why $`k=1`$ should hold even for a residual set in $`U`$.
Theorem I has an interesting interpretation. Recall that a $`G`$-extension of a dynamical system $`f:XX`$ is a map $`f_\phi :X\times GX\times G`$, where $`G`$ is a compact group, of the form $`(x,y)(g(x),\phi (x)y)`$. If $`f`$ preserves $`\nu `$, and $`\phi :XG`$ is measurable, then $`f_\phi `$ preserves the product of $`\nu `$ with Lebesgue-Haar measure on $`G`$. A $`𝐙/k𝐙`$-extension is also called a $`k`$-point extension.
Let $`\lambda `$ be an invariant probability measure for a $`k`$-point extension of $`f:XX`$, and $`\{\lambda _x\}`$ the family of conditional measures associated with the partition $`\{\{x\}\times G\}`$. We remark that if $`\lambda `$ is ergodic, then each atom of $`\lambda _x`$ must have the same weight $`1/k`$ (up to a set of $`\lambda `$-measure $`0`$).
Now take $`gU`$. Choose a coherent orientation on the leaves of $`\{\pi ^1(x)\}_{xT^2}`$. Take $`h:𝐓^3𝐓^2\times 𝐓`$ to be any continuous change of coordinates such that $`h`$ restricted to $`\pi ^1(x)`$ is smooth and orientation preserving to $`\{x\}\times 𝐓`$. We may then write $`F=hgh^1:𝐓^2\times 𝐓𝐓^2\times 𝐓`$ in the form
$$F(x,p)=(A_2x,\phi _x(p))$$
where $`\phi _x:𝐓𝐓`$ is smooth and orientation preserving. If $`P:𝐓^2\times 𝐓𝐓^2`$ is the projection on the first factor of the product, we have $`Ph=\pi `$. Therefore, writing $`\lambda =h^{}\mu `$, we have $`P^{}\lambda =\pi ^{}\mu `$. Let $`\{\lambda _x\}`$ be the disintegration of the measure $`\lambda `$ along the fibers $`\{x\}\times 𝐓`$. By a further measurable change of coordinates, smooth along each $`\{x\}\times 𝐓`$ fiber, we may assume that $`\lambda `$-almost everywhere, the atoms of $`\lambda _x`$ are at $`l/k`$, for $`l=0,\mathrm{},k1`$. But then $`\phi _x`$ permutes the atoms cyclically, and we obtain the following corollary.
Corollary: For every $`gU`$ there exists $`k𝐍`$ such that $`(𝐓^3,\mu ,g)`$ is isomorphic to an (ergodic) $`k`$-point extension of $`(𝐓^2,\pi ^{}\mu ,A_2)`$.
M. Shub has observed that if $`g=j_{a,k}h_b`$, then $`\pi ^{}\mu `$ is actually Lebesgue measure on $`𝐓^2`$.
## 1 Proof of Theorem I
The proof of Theorem I follows from a more general result about fibered diffeomorphisms. Before stating this result, we describe the underlying setup and assumptions.
Let $`X`$ be a compact metric space with Borel probability measure $`\nu `$, and let $`f:XX`$ be invertible and ergodic with respect to $`\nu `$. Let $`M`$ be a closed Riemannian manifold and $`\phi :X\text{Diff}^{1+\alpha }(M)`$ a measurable map. Consider the skew-product transformation $`F:X\times MX\times M`$ given by
$$F(x,p)=(f(x),\phi _x(p)).$$
Assume further that there is an $`F`$-invariant ergodic probability measure $`\mu `$ on $`X\times M`$ such that $`\pi _{}\mu =\nu `$, where $`\pi :X\times MX`$ is the projection onto the first factor.
For $`xX`$, let $`\phi _x^{(0)}`$ be the identity map on $`M`$ and for $`k𝐙`$, define $`\phi _x^{(k)}`$ by
$$\phi _x^{(k+1)}=\phi _{f^k(x)}\phi _x^{(k)}.$$
Since the tangent bundle to $`M`$ is measurably trivial, the derivative map of $`\phi `$ along the $`M`$ direction gives a cocycle $`D\phi :X\times M\times 𝐙GL(n,𝐑)`$, where $`n=\text{dim}(M)`$:
$$(x,p,k)D_p\phi _x^{(k)}.$$
Assume that $`\mathrm{log}^+D\phi _\alpha L^1(X\times M,\mu )`$, where $`_\alpha `$ is the $`\alpha `$-Hölder norm. Let $`\lambda _1<\lambda _2\mathrm{}<\lambda _l`$ be the Lyapunov exponents of this cocycle; they exist for $`\mu `$-a.e. $`(x,p)`$ by Oseledec’s Theorem and are constant by ergodicity. We call these the fiberwise exponents of $`F`$. Under the assumptions just described, we have the following result.
Theorem II: Suppose that $`\lambda _l<0`$. Then there exists a set $`SX\times M`$ and an integer $`k1`$ such that
* $`\mu (S)=1`$
* For every $`(x,p)S`$, we have $`\mathrm{\#}(S\{x\}\times M)=k`$.
This has the immediate corollary:
Corollary: Let $`f\text{Diff}^{1+\alpha }(M)`$. If $`\mu `$ is an ergodic measure with all of its exponents negative, then it is concentrated on the orbit of a periodic sink.
The corollary has a simple proof using regular neighborhoods. Our proof is a fibered version. Theorem I is also a corollary of Theorem II. For this, the argument is actually applied to the inverse of $`g`$, which has negative fiberwise exponents, rather than to $`g`$ itself, whose fiberwise exponents are positive. As we described in the previous remarks, there is a measurable change of coordinates, smooth along the leaves of $`𝒲_g^c`$ in which $`g^1`$ is expressed as a skew product of $`𝐓^2\times 𝐓`$.
Remark: Without the assumption that $`f`$ is invertible, Theorem II is false. An example is described by Y. Kifer \[Ki\], which we recall here. Let $`f:𝐓𝐓`$ be a $`C^{1+\alpha }`$ diffeomorphism with exactly two fixed points, one attracting and one repelling. Consider the following random diffeomophism of $`𝐓`$: with probability $`p(0,1)`$, apply $`f`$, and with probability $`1p`$, rotate by an angle chosen randomly from the interval $`[ϵ,ϵ]`$.
Let $`X=(\{0,1\}\times 𝐓)^𝐍`$. To generate a sequence of diffeomorphisms $`f_0,f_1,\mathrm{}`$ according to the above rule, we first define $`\phi :X\text{Diff}^{1+\alpha }(𝐓)`$ by
$$\phi (\omega )=\{\begin{array}{cc}f\hfill & \text{if }\omega (0)=(0,\theta )\text{,}\hfill \\ R_\theta \hfill & \text{if }\omega (0)=(1,\theta )\text{,}\hfill \end{array}$$
where $`R_\theta `$ is rotation through angle $`\theta `$. Next, we let $`\nu _ϵ`$ be the product of $`p,1p`$-measure on $`\{0,1\}`$ with the measure on $`𝐓`$ that is uniformly distributed on $`[ϵ,ϵ]`$. Then corresponding to $`\nu _ϵ^𝐍`$-almost every element $`\omega X`$ is the sequence $`\{f_k=\phi (\sigma ^k(\omega ))\}_{k=0}^{\mathrm{}},`$ where $`\sigma :XX`$ is the one-sided shift $`\sigma (\omega )(n)=\omega (n+1)`$.
Put another way, the random diffeomorphism is generated by the (noninvertible) skew product $`\tau :X\times 𝐓X\times 𝐓`$, where $`\tau (\omega ,x)=(\sigma (\omega ),\phi (\omega )(x))`$. An ergodic $`\nu _ϵ`$-stationary measure for this random diffeomorphism is a measure $`\mu _ϵ`$ on $`𝐓`$ such that $`\mu _ϵ\times \nu _ϵ^𝐍`$ is $`\tau `$-invariant and ergodic. Such measures always exist (\[Ki\], Lemma I.2.2), but, for this example, there is an ergodic stationary measure with additional special properties.
Specifically, for every $`ϵ>0`$, there exists an ergodic $`\nu _ϵ`$-stationary measure $`\mu _ϵ`$ on $`𝐓`$ such that, as $`ϵ0`$, $`\mu _ϵ\delta _{x_0}`$, in the weak topology, where $`\delta _{x_0}`$ is Dirac measure concentrated on the sink $`x_0`$ for $`f`$. From this, it follows that, as $`ϵ0`$, the fiberwise Lyapunov exponent for $`\mu _ϵ`$ approaches $`\mathrm{log}|f^{}(x_0)|<0`$, which is the Lyapunov exponent of $`\delta _{x_0}`$. Thus, for $`ϵ`$ sufficiently small, the fiberwise exponent for $`\tau `$ with respect to $`\mu _ϵ`$ is negative. Nonetheless, it is easy to see that $`\mu _ϵ`$ for $`ϵ>0`$ cannot be uniformly distributed on $`k`$ atoms; if $`\mu _ϵ`$ were atomic, then $`\tau `$-invariance of $`\mu _ϵ\times \nu _ϵ^𝐍`$ would imply that, for every $`x𝐓`$,
$`\mu _ϵ(\{x\})`$ $`=`$ $`p\mu _ϵ(\{f^1(x)\})+(1p){\displaystyle _ϵ^ϵ}\mu _ϵ(\{R_\theta (x)\})𝑑\theta `$
$`=`$ $`p\mu _ϵ(\{f^1(x)\}),`$
which is impossible if $`\mu _ϵ`$ has finitely many atoms. In fact, $`\mu _ϵ`$ can be shown to be absolutely continuous with respect to Lebesgue measure (see \[Ki\], p. 173ff and the references cited therein). Hence invertibility is essential, and we indicate in the proof of Theorem II where it is used.
Proof of Theorem II: We first establish the existence of fiberwise “stable manifolds” for the skew product $`F`$. A general theory of stable manifolds for random dynamical systems is worked out in (\[Ki\], Theorem V.1.6; see also \[BL\]); since we are assuming that all of the fiberwise exponents for $`F`$ are negative, we are faced with the simpler task of constructing fiberwise regular neighborhoods for $`F`$ (see the Appendix by Katok and Mendoza in \[KH\]). We outline a proof, following closely \[KH\].
###### Theorem 1.1
(Existence of Regular Neighborhoods) There exists a set $`\mathrm{\Lambda }_0X\times M`$ of full measure such that for $`ϵ>0`$:
* There exists a measurable function $`r:\mathrm{\Lambda }_0(0,1]`$ and a collection of embeddings $`\mathrm{\Psi }_{(x,p)}:B(0,q(x,p))M`$ such that $`\mathrm{\Psi }_{(x,p)}(0)=p`$ and $`\text{exp}(ϵ)<r(F(x,p))/r(x,p)<\text{exp}(ϵ)`$.
* If $`\phi _{(x,p)}=\mathrm{\Psi }_{F(x,p)}^1\phi _x\mathrm{\Psi }_{(x,p)}:B(0,r(x,p))𝐑^n`$, then $`D_0\phi _{(x,p)}`$ satisfies
$$\text{exp}(\lambda _1ϵ)D_0\phi _{(x,p)}^1^1,D_0\phi _{(x,p)}\text{exp}(\lambda _l+ϵ).$$
* The $`C^1`$ distance $`d_{C^1}(\phi _{(x,p)},D_0\phi _{(x,p)})<ϵ`$ in $`B(0,r(x,p))`$.
* There exist a constant $`K>0`$ and a measurable function $`A:\mathrm{\Lambda }_0𝐑`$ such that for $`y,zB(0,r(x,p))`$,
$$K^1d(\mathrm{\Psi }_{(x,p)}(y),\mathrm{\Psi }_{(x,p)}(z))yzA(x)d(\mathrm{\Psi }_{(x,p)}(y),\mathrm{\Psi }_{(x,p)}(z)),$$
with $`\text{exp}(ϵ)<A(F(x,p))/A(x,p)<\text{exp}(ϵ)`$.
Proof: See the proof of Theorem S.3.1 in \[KH\]. $`\mathrm{}`$
Decompose $`\mu `$ into a system of fiberwise measures $`d\mu (x,p)=d\mu _x(p)d\nu (x)`$. Invariance of $`\mu `$ with respect to $`F`$ implies that, for $`\nu `$-a.e. $`xX`$,
$$\phi _{x}^{}{}_{}{}^{}\mu _x=\mu _{f(x)}.$$
###### Corollary 1.2
There exists a set $`\mathrm{\Lambda }X\times M`$, and real numbers $`R>0`$, $`C>0`$, and $`c<1`$ such that
* $`\mu (\mathrm{\Lambda })>.5`$, and, if $`(x,p)\mathrm{\Lambda }`$, then $`\mu _x(\mathrm{\Lambda }_x)>.5`$, where $`\mathrm{\Lambda }_x=\{pM|(x,p)\mathrm{\Lambda }\}`$,
* If $`(x,p)\mathrm{\Lambda }`$ and $`d_M(p,q)R`$, then
$$d_M(\phi _x^{(m)}(p),\phi _x^{(m)}(q))Cc^md_M(p,q),$$
for all $`m0`$.
Proof: This follows in a standard way from the Mean Value Theorem and Lusin’s Theorem.$`\mathrm{}`$
To prove Theorem II, it suffices to show that there is a positive $`\nu `$-measure set $`BX`$, such that for $`xB`$, the measure $`\mu _x`$ has an atom, as the following argument shows. For $`xX`$, let $`d(x)=sup_{pM}\mu _x(p)`$. Clearly $`d`$ is measurable, $`f`$-invariant, and positive on $`B`$. Ergodicity of $`f`$ implies that $`d(x)=d>0`$ is positive and constant for almost all $`xX`$. Let $`S=\{(x,p)X\times M|\mu _x(p)d\}`$. Observe that $`S`$ is $`F`$-invariant, has measure at least $`d`$, and hence has measure $`1`$. The conclusions of Theorem II follow immediately.
Let $`\mathrm{\Lambda }`$, $`R>0`$, $`C>0`$, and $`c<1`$ be given by Corollary 1.2, and let $`B=\pi (\mathrm{\Lambda })`$. Let $`N`$ be the number of $`R/10`$-balls needed to cover $`M`$. We now show that for $`\nu `$-almost every $`xB`$, the measure $`\mu _x`$ has at least one atom.
For $`xX`$, let
$$m(x)=inf\text{diam }(U_j),$$
where the infimum is taken over all collections of closed balls $`U_1,\mathrm{},U_k`$ in $`M`$ such that $`kN`$ and $`\mu _x(_{j=1}^kU_j).5`$. Let $`m=\text{ess sup }_{xB}m(x)`$.
We now show that $`m=0`$. If $`m>0`$, then there exists an integer $`J`$ such that
$`C\mathrm{\Delta }c^JN`$ $`<`$ $`m/2,`$ (1)
where $`\mathrm{\Delta }`$ is the diameter of $`M`$. Let $`𝒰`$ be a cover of $`M`$ by $`N`$ closed balls of radius $`R/10`$. For $`xB`$, let $`U_1(x),\mathrm{},U_{k(x)}(x)`$ be those balls in $`𝒰`$ that meet $`\mathrm{\Lambda }_x`$. Since these balls cover $`\mathrm{\Lambda }_x`$, and $`\mu _x(\mathrm{\Lambda }_x)>.5`$, it follows that $`\mu _x(_{j=1}^{k(x)}U_j(x)).5`$. But $`\phi _{x}^{(i)}{}_{}{}^{}\mu _x=\mu _{f^i(x)},`$ and so it’s also true that
$`\mu _{f^i(x)}({\displaystyle \underset{j=1}{\overset{k(x)}{}}}\phi _x^{(i)}(U_j(x)))`$ $``$ $`.5,`$ (2)
for all $`i`$.
We now use the fact that $`\phi _x^{(i)}`$ contracts regular neighborhoods to derive a contradiction. The balls $`U_j(x)`$ meet $`\mathrm{\Lambda }_x`$ and have diameter less than $`R/10`$, and so by Corollary 1.2, (2), we have
$`\text{diam }(\phi _x^{(i)}(U_j(x)))`$ $``$ $`C\mathrm{\Delta }c^i.`$ (3)
Let $`\tau :B𝐍`$ be the first-return time of $`f^J`$ to $`B`$, so that $`f^{J\tau (x)}(x)B`$, and $`f^{Ji}(x)B`$, for $`i\{1,\mathrm{},\tau (x)1\}`$. Decompose the set $`B`$ according to these first return times:
$$B=\underset{i=1}{\overset{\mathrm{}}{}}B_i(\text{mod }0),$$
where $`B_i=\tau ^1(i)`$. Because $`f`$ is invertible and $`f^1`$ preserves measure, we also have the mod 0 equivalence:
$$B^{}:=\underset{i=1}{\overset{\mathrm{}}{}}f^{Ji}(B_i)=B(\text{mod }0).$$
Let $`yB^{}`$. Then $`y=f^{Ji}(x)`$, where $`xB_iB`$, for some $`i1`$. It follows from the definition of $`m(y)`$ and inequalities (2), (3) and (1) that
$`m(y)`$ $``$ $`{\displaystyle \underset{j=1}{\overset{k(x)}{}}}\text{diam }(\phi _x^{(Ji)}(U_j(x)))`$
$``$ $`Ck(x)\mathrm{\Delta }c^{Ji}`$
$``$ $`CN\mathrm{\Delta }c^J`$
$`<`$ $`m/2.`$
But then
$`m`$ $`=`$ $`\text{ess sup }_{xB}m(x)`$
$`=`$ $`\text{ess sup }_{yB^{}}m(y)`$
$`<`$ $`m/2,`$
contradicting the assumption $`m>0`$.
Thus $`m=0`$, and, for $`\nu `$-almost every $`xB`$, we have $`m(x)=0`$. If $`m(x)=0`$, then there is a sequence of closed balls $`U^1(x),U^2(x),\mathrm{}`$ with $`lim_i\mathrm{}\text{diam }(U^i(x))=0`$ and $`\mu _x(U^i(x)).5/N`$, for all $`i`$. Take $`p_iU^i(x)`$; any accumulation point of $`\{p_i\}`$ is an atom for $`\mu _x`$. Since we have shown that $`\mu _x`$ has an atom, for $`\nu `$-a.e. $`xB`$, the proof of Theorem II is complete.$`\mathrm{}`$
We thank Michael Shub and Anatole Katok for useful conversations. |
warning/0002/hep-th0002195.html | ar5iv | text | # 1 What is the chiral anomaly?
## 1 What is the chiral anomaly?
The chiral abelian anomaly has been discovered, in the past century, by Adler, Bell and Jackiw, after earlier work on $`\pi ^0`$-decay starting with Steinberger and Schwinger; see e.g. and references given there. It has been rederived in many different ways of varying degree of mathematical rigor by many people. Diverse physical implications, especially in particle physics, have been discussed. It is hard to imagine that one may still be able to find new, interesting implications of the chiral anomaly that specialists have not been aware of, for many years. Yet, until very recently — in the past century, but only two to three years ago — this turned out to be possible, and we suspect that further applications may turn up in the future! This little review is devoted to a discussion of physical implications of the chiral anomaly that have been discovered recently.
Before we turn to physics, we recall what is meant by “chiral (abelian) anomaly”. In general terms, one speaks of an anomaly if some quantum theory violates a symmetry present at the classical level, (i.e., in the limit where $`\mathrm{}0`$). By “violating a symmetry” one means that it is impossible to construct a unitary representation of the symmetry transformations of the classical system underlying some quantum theory on the Hilbert space of pure state vectors of the quantum theory. (By “violating a dynamical symmetry” is meant that it is impossible to construct such a representation that commutes with the unitary time evolution of the quantum theory.)
It is quite clear that understanding anomalies can be viewed as a problem in group cohomology theory. A fundamental example of an anomalous symmetry group is the group of all symplectic transformations of the phase space of a classical Hamiltonian system underlying some quantum theory.
The anomalies considered in this review are ones connected with infinite-dimensional groups of gauge transformations which are symmetries of some classical Lagrangian systems with infinitely many degrees of freedom (Lagrangian field theories). Thus, we consider a theory of charged, massless fermions coupled to an external electromagnetic field in Minkowski space-time of even dimension $`2n`$. Let $`\gamma ^0,\gamma ^1,\mathrm{},\gamma ^{2n1}`$ denote the usual Dirac matrices, and define
$$\gamma :=i\gamma ^0\gamma ^1\mathrm{}\gamma ^{2n1}.$$
(1.1)
Then $`\gamma `$ anti-commutes with the covariant Dirac operator
$$D:=i\gamma ^\mu \left(_\mu iA_\mu \right),$$
(1.2)
where $`A`$ is the vector potential of the external electromagnetic field. Let $`\psi (x)`$ denote the Dirac spinor field and $`\overline{\psi }(x)`$ the conjugate field. We define the vector current, $`𝒥^\mu `$, and the axial vector current $`\stackrel{~}{𝒥}^\mu `$, by
$$𝒥^\mu :=\overline{\psi }\gamma ^\mu \psi ,\stackrel{~}{𝒥}^\mu :=\overline{\psi }\gamma ^\mu \gamma \psi .$$
(1.3)
At the classical level, these currents are conserved,
$$_\mu 𝒥^\mu =\mathrm{\hspace{0.33em}0},_\mu \stackrel{~}{𝒥}^\mu =\mathrm{\hspace{0.33em}0},$$
(1.4)
on solutions of the equations of motion, $`(D\psi =0)`$. The conservation of the vector current is intimately connected with the electromagnetic gauge invariance of the theory,
$`\psi (x)e^{i\chi (x)}\psi (x),\overline{\psi }(x)\overline{\psi }(x)e^{i\chi (x)}`$
$`A_\mu (x)A_\mu (x)+_\mu \chi (x),`$ (1.5)
where $`\chi (x)`$ is a test function on space-time. When $`\chi `$ is constant in $`x`$ the transformations (1.5) are a global symmetry of the classical theory corresponding to the conserved quantity
$$Q=𝑑\underset{¯}{x}𝒥^0(x^0,\underset{¯}{x})$$
(1.6)
which is the electric charge. The conservation of $`Q`$ (independence of $`x^0`$) follows, of course, from the fact that the Noether current $`𝒥^\mu `$ associated with (1.5) satisfies the continuity equation (1.4).
The conservation of the axial vector current $`\stackrel{~}{𝒥}^\mu `$, in the classical theory, is connected with the invariance of the theory under local chiral rotations
$`\psi (x)e^{i\alpha (x)\gamma }\psi (x),\overline{\psi }(x)\overline{\psi }(x)e^{i\alpha (x)\gamma }`$
$`A_\mu (x)A_\mu (x)+\gamma _\mu \alpha (x),`$ (1.7)
where $`\alpha (x)`$ is a test function on space-time. In particular, when $`\alpha `$ is a constant the transformations (1.7) are a global symmetry of the classical theory corresponding to the conserved charge
$$\stackrel{~}{Q}=𝑑\underset{¯}{x}\stackrel{~}{𝒥}^0(x^0,\underset{¯}{x})$$
(1.8)
(which, according to (1.4), is independent of $`x^0`$).
It turns out that, in the quantum theory, the local chiral rotations (1.7) do not leave quantum-mechanical transition amplitudes invariant, and the axial vector current $`\stackrel{~}{𝒥}^\mu `$ is not a conserved current, for arbitrary external electromagnetic fields. This phenomenon is called chiral (abelian) anomaly.
Let us see where the chiral anomaly comes from, for theories in two and four space-time dimensions. We start with the discussion of two-dimensional theories. We consider a quantum theory which has a conserved vector current $`𝒥^\mu `$ and — if the external electromagnetic field vanishes — a conserved axial vector current $`\stackrel{~}{𝒥}^\mu `$, i.e.,
$$_\mu 𝒥^\mu =\mathrm{\hspace{0.33em}0},_\mu \stackrel{~}{𝒥}^\mu =\mathrm{\hspace{0.33em}0}.$$
(1.9)
In two space-time dimensions, $`𝒥^\mu `$ and $`\stackrel{~}{𝒥}^\mu `$ are related to each other by
$$\stackrel{~}{𝒥}^\mu =\epsilon ^{\mu \nu }𝒥_\nu $$
(1.10)
where $`\epsilon ^{00}=\epsilon ^{11}=0,\epsilon ^{01}=\epsilon ^{10}=1`$. The continuity equation
$$_\mu 𝒥^\mu =\mathrm{\hspace{0.33em}0}$$
has the general solution
$$𝒥^\mu (x)=\frac{q}{2\pi }\epsilon ^{\mu \nu }\left(_\nu \phi \right)(x),$$
(1.11)
where $`\phi (x)`$ is an arbitrary scalar field on space-time, and $`q`$ denotes the electric charge. Using eqs. (1.11) and (1.10) and the continuity equation,
$$_\mu \stackrel{~}{𝒥}^\mu =\mathrm{\hspace{0.33em}0},$$
for the axial vector current, we find that the field $`\phi `$ must obey the equation of motion
$$\mathrm{}\phi (x)=\mathrm{\hspace{0.33em}0}.$$
(1.12)
Thus, if the vector- and axial vector currents are conserved then the potential $`\phi `$ of the vector current is a massless free field. The theory of the massless free field is an example of a Lagrangian field theory. It has an action functional, $`S`$, given by
$$S(\phi )=\frac{1}{4\pi }d^2x\left(_\mu \phi \right)(x)\left(^\mu \phi \right)(x).$$
(1.13)
The “momentum”, $`\pi (x)`$, canonically conjugate to $`\phi (x)`$ is defined, as usual, by
$$\pi (x)=\delta S(\phi )/\delta \left(_0\phi (x)\right)=\frac{1}{2\pi }\frac{\phi (x)}{t},$$
(1.14)
where $`t=x^0`$ denotes time; (the “velocity of light” $`c=1`$). In the quantum theory, $`\phi `$ and $`\pi `$ are operator-valued distributions on Fock space satisfying the equal-time canonical commutation relations
$$[\pi (t,\underset{¯}{x}),\phi (t,\underset{¯}{y})]=i\delta \left(\underset{¯}{x}\underset{¯}{y}\right).$$
(1.15)
Since
$$𝒥^\mu (x)=\frac{q}{2\pi }\epsilon ^{\mu \nu }\left(_\nu \phi \right)(x),$$
and
$$\stackrel{~}{𝒥}^\mu (x)=\epsilon ^{\mu \nu }𝒥_\nu (x)=\frac{q}{2\pi }\left(^\mu \phi \right)(x),$$
eq. (1.15) yields the well known anomalous commutator
$$[𝒥^0(t,\underset{¯}{x}),\stackrel{~}{𝒥}^0(t,\underset{¯}{y})]=i\frac{q^2}{2\pi }\delta ^{}\left(\underset{¯}{x}\underset{¯}{y}\right).$$
(1.16)
Next, we imagine that the system is coupled to a classical external electric field $`E(x)`$. In two space-time dimensions, the electric field is given in terms of the electromagnetic vector potential $`A_\mu `$ by
$$E(x)=\epsilon ^{\mu \nu }\left(_\mu A_\nu \right)(x).$$
(1.17)
The action functional for the theory in an external electric field is given by
$`S(\phi ,A)`$ $`=`$ $`{\displaystyle \frac{1}{4\pi }}{\displaystyle d^2x\left(_\mu \phi \right)(x)\left(^\mu \phi \right)(x)}+{\displaystyle \frac{1}{q}}{\displaystyle d^2x𝒥^\mu (x)A_\mu (x)}`$ (1.18)
$`=`$ $`{\displaystyle \frac{1}{4\pi }}{\displaystyle d^2x\left\{\left(_\mu \phi \right)(x)\left(^\mu \phi \right)(x)+2\epsilon ^{\mu \nu }_\nu \phi (x)A_\mu (x)\right\}}.`$
The equation of motion (Euler-Lagrange equation) obtained from the action function (1.18) is
$$\mathrm{}\phi (x)=E(x).$$
(1.19)
Using (1.10) and (1.11), we see that equation (1.19) is equivalent to
$$_\mu \stackrel{~}{𝒥}^\mu (x)=\frac{q}{2\pi }E(x),$$
(1.20)
i.e., the axial vector current fails to be conserved in a non-vanishing external electric field $`E`$. Equation (1.20) is the standard expression of the chiral anomaly in two space-time dimensions.
From the currents $`𝒥^\mu `$ and $`\stackrel{~}{𝒥}^\mu `$ one can construct chiral currents, $`𝒥_{\mathrm{}}^\mu `$ and $`𝒥_r^\mu `$, for left-moving and right-moving modes by setting
$$𝒥_{\mathrm{}}^\mu =𝒥^\mu \stackrel{~}{𝒥}^\mu ,𝒥_r^\mu =𝒥^\mu +\stackrel{~}{𝒥}^\mu .$$
(1.21)
They satisfy the equations
$$_\mu 𝒥_{\mathrm{}/r}^\mu =\frac{q}{2\pi }E(x).$$
(1.22)
From eqs. (1.17) and (1.22) we infer that one can define modified chiral currents, $`\widehat{𝒥}_{\mathrm{}/r}^\mu `$, which are conserved:
$$\widehat{𝒥}_{\mathrm{}/r}^\mu (x):=𝒥_{\mathrm{}/r}^\mu \pm \frac{q}{2\pi }\epsilon ^{\mu \nu }A_\nu (x).$$
(1.23)
Then
$$_\mu \widehat{𝒥}_{\mathrm{}/r}^\mu (x)=\mathrm{\hspace{0.33em}0},$$
but $`\widehat{𝒥}_{\mathrm{}/r}^\mu `$ fail to be gauge-invariant. Nevertheless the conserved charges,
$$N_{\mathrm{}}:=𝑑\underset{¯}{x}\widehat{𝒥}_{\mathrm{}}^0(t,\underset{¯}{x}),N_r:=𝑑\underset{¯}{x}\widehat{𝒥}_r^0(t,\underset{¯}{x}),$$
(1.24)
are gauge-invariant. They count the total electric charge of left-moving and of right-moving modes, respectively, present in a physical state of the system.
The anomalous commutators are given by
$$[𝒥^0(t,\underset{¯}{x}),\widehat{𝒥}_{\mathrm{}/r}^0(t,\underset{¯}{y})]=i\frac{q^2}{2\pi }\delta ^{}\left(\underset{¯}{x}\underset{¯}{y}\right).$$
The left-moving / right-moving charged fields of the theory can be expressed as normal-ordered exponentials of spatial integrals of $`\frac{1}{q}\widehat{𝒥}_{\mathrm{}/r}^0(x)`$, i.e., as vertex operators; they transform correctly under gauge transformations.
This completes our review of the chiral anomaly and of anomalous commutators in two dimensions, and we now turn to four\- (or higher-) dimensional systems.
We consider charged, massless Dirac fermions described by a Dirac spinor field $`\psi (x)`$ and its conjugate field $`\overline{\psi }(x)=\psi ^{}(x)\gamma _0`$. We study the effect of coupling these fields to external vector- and axial-vector potentials, $`A_\mu `$ and $`Z_\mu `$, respectively. The theory of these fields provides an example of Lagrangian field theory, the action functional being given by
$$S(\overline{\psi },\psi ;A,Z):=d^{2n}x\overline{\psi }(x)D_{A,Z}\psi (x),$$
(1.25)
where the covariant Dirac operator is
$$D_{A,Z}=i\gamma ^\mu \left(_\mu iA_\mu iZ_\mu \gamma \right),$$
(1.26)
with $`\gamma (=`$$`\gamma ^5`$”) as in eq. (1.1). The fields $`A_\mu `$ and $`Z_\mu `$ are arbitrary external fields (i.e., they are not quantized, for the time being). We define the effective action, $`S_{\mathrm{eff}}(A,Z)`$, by
$$e^{iS_{\mathrm{eff}}(A,Z)}:=\mathrm{const}𝒟\psi 𝒟\overline{\psi }e^{iS(\overline{\psi },\psi ;A,Z)},$$
(1.27)
where the constant is chosen such that $`S(A=0,Z=0)=0`$, and $`\mathrm{}`$ and $`c`$ have been set to 1. After Wick rotation,
$$t=x^0ix^0,A_0iA_0,Z_0iZ_0,\gamma ^0i\gamma ^0,$$
(1.28)
eq. (1.27) reads
$$e^{S_{\mathrm{eff}}^E(A,Z)}=\left[𝒟\psi 𝒟\overline{\psi }e^{S^E(\overline{\psi },\psi ;A,Z)}\right]_{\mathrm{ren}}$$
(1.29)
where the integral on the R.S. is interpreted as a renormalized Gaussian Berezin integral. Thus
$$e^{S_{\mathrm{eff}}^E(A,Z)}=\mathrm{det}_{\mathrm{ren}}\left(D_{A,Z}\right),$$
(1.30)
where, after Wick rotation,
$$D_{A,Z}=i\gamma ^\mu \left(_\mu iA_\mu iZ_\mu \gamma \right)$$
is an anti-hermitian elliptic operator, and the subscripts “ren” indicate that (for $`n2`$) a multiplicative renormalization must be made.
The effective action $`S_{\mathrm{eff}}^E(A,Z)`$ is the generating function for the Euclidian Green functions of the vector- and axial vector currents. At non-coinciding arguments,
$`𝒥^{\mu _1}(x_1)\mathrm{}\stackrel{~}{𝒥}^{\nu _1}(y_1)\mathrm{}_{A,Z}^c`$
$`=(iq){\displaystyle \frac{\delta }{\delta A_{\mu _1}(x_1)}}\mathrm{}(iq){\displaystyle \frac{\delta }{\delta Z_{\nu _1}(y_1)}}\mathrm{}S_{\mathrm{eff}}^E(A,Z),`$ (1.31)
where $`q`$ is the electric charge, and $`()_{A,Z}^c`$ denotes a connected expectation value.
We should like to understand how $`S_{\mathrm{eff}}^E(A,Z)`$ changes under the gauge transformations
$$A_\mu A_\mu +_\mu \chi ,Z_\mu Z_\mu +_\mu \alpha .$$
(1.32)
Following Fujikawa , we perform a phase transformation and a chiral rotation of $`\psi `$ and $`\overline{\psi }`$ under the integral on the R.S. of eq. (1.29). We set
$$\psi ^{}(x)=e^{i\left(\chi (x)+\alpha (x)\gamma \right)}\psi (x),\overline{\psi }^{}(x)=\overline{\psi }(x)e^{i\left(\chi (x)\alpha (x)\gamma \right)}.$$
(1.33)
Then
$$S^E(\overline{\psi }^{},\psi ^{};A+d\chi ,Z+d\alpha )=S^E(\overline{\psi },\psi ;A,Z),$$
(1.34)
where $`d\chi `$ denotes the gradient, $`(_\mu \chi )`$, of $`\chi `$. Next, we must determine the Jacobian, $`J`$, of the transformation (1.33),
$$𝒟\overline{\psi }^{}𝒟\psi ^{}=:J𝒟\overline{\psi }𝒟\psi .$$
(1.35)
Obviously, phase transformations,
$$\psi ^{}=e^{i\chi }\psi ,\overline{\psi }^{}=\overline{\psi }e^{i\chi }$$
have Jacobian $`J=1`$. However, this may not be so for chiral rotations. Formally, under chiral rotations, the Jacobian turns out to be
$$J=\mathrm{exp}\left[2i\mathrm{Tr}\left(\alpha \gamma \right)\right].$$
(1.36)
The problem with the R.S. of (1.36) is that, a priori, it is ill-defined. Let us assume that non-compact Euclidian space-time is replaced by a $`2n`$-dimensional sphere. Then $`D_{A,Z}`$ has discrete spectrum, with eigenvalues $`i\lambda _m`$ corresponding to eigenspinors $`\psi _m(x),m`$. Formally,
$$\mathrm{Tr}(\alpha \gamma )=\underset{m}{}d^{2n}x\alpha (x)\psi _m^{}(x)\gamma \psi _m(x).$$
We regularize the R.S. by replacing it by
$$\underset{m}{}e^{\left(\lambda _m^2/M^2\right)}d^{2n}x\alpha (x)\psi _m^{}(x)\gamma \psi _m(x)$$
(1.37)
and, afterwards, letting $`M\mathrm{}`$. Expression (1.37) is nothing but
$$\mathrm{Tr}\left(\alpha \gamma e^{\left(D_{A,Z}^2/M^2\right)}\right).$$
(1.38)
From Alvarez-Gaumé’s calculations concerning the index theorem, for example, we infer that
$$\underset{M\mathrm{}}{lim}\mathrm{Tr}\left(\alpha \gamma e^{\left(D_{A,Z}^2/M^2\right)}\right)=d^{2n}x\alpha (x)𝒜(x),$$
(1.39)
where $`𝒜(x)`$ is the index density described more explicitly below. From (1.39) and (1.36) we obtain that
$$J=\mathrm{exp}\left[2id^{2n}x\alpha (x)𝒜(x)\right].$$
(1.40)
With (1.34), (1.35) and (1.29), eq. (1.40) yields
$$S_{\mathrm{eff}}^E(A+d\chi ,Z+d\alpha )=S_{\mathrm{eff}}^E(A,Z)\mathrm{\hspace{0.17em}2}id^{2n}x\alpha (x)𝒜(x).$$
(1.41)
When combined with (1.31) eq. (1.41) is seen to yield
$`\left[\delta S_{\mathrm{eff}}^E(A+d\chi ,0)/\delta \chi (x)\right]_{\chi =0}`$
$`=_\mu \left(\delta S_{\mathrm{eff}}^E(A,0)/\delta A_\mu (x)\right)`$
$`={\displaystyle \frac{i}{q}}_\mu 𝒥^\mu (x)_A=\mathrm{\hspace{0.33em}0},`$ (1.42)
and
$`\left[\delta S_{\mathrm{eff}}^E(A,Z+d\alpha )/\delta \alpha (x)\right]_{Z=\alpha =0}`$
$`=_\mu \left(\left[\delta S_{\mathrm{eff}}(A,Z)/\delta Z_\mu (x)\right]_{Z=0}\right)`$
$`={\displaystyle \frac{i}{q}}_\mu \stackrel{~}{𝒥}^\mu (x)_A=2i𝒜(x),`$ (1.43)
i.e.,
$$_\mu 𝒥^\mu (x)_A=\mathrm{\hspace{0.33em}0},_\mu \stackrel{~}{𝒥}^\mu (x)_A=2q𝒜(x).$$
(1.44)
Introducing the chiral currents
$$𝒥_{\mathrm{}}^\mu :=𝒥^\mu \stackrel{~}{𝒥}^\mu ,𝒥_r^\mu :=𝒥^\mu +\stackrel{~}{𝒥}^\mu ,$$
(1.45)
where $`𝒥_{\mathrm{}/r}^\mu `$ is the current of left-handed/right-handed fermions, we see that (1.44) is equivalent to
$$_\mu 𝒥_{\mathrm{}}^\mu (x)_A=\mathrm{\hspace{0.33em}2}q𝒜(x),_\mu 𝒥_r^\mu (x)_A=2q𝒜(x).$$
(1.46)
Locally, we can solve the equation
$$\delta \omega (x;A)=𝒜(x),$$
(1.47)
where $`\delta `$, the co-differential, is the dual of exterior differentiation $`d`$, the solution $`\omega (;A)`$ being a 1-form. The 1-form $`\omega (;A)`$ is, however, not gauge-invariant. We may now define modified currents,
$$\widehat{𝒥}_{\mathrm{}/r}^\mu (x)=𝒥_{\mathrm{}/r}^\mu (x)\mathrm{\hspace{0.33em}2}q\omega ^\mu (x;A).$$
(1.48)
They are not gauge-invariant, but, according to eqs. (1.47), (1.48), they are conserved, i.e.,
$$_\mu \widehat{𝒥}_{\mathrm{}/r}^\mu (x)=\mathrm{\hspace{0.33em}0}.$$
(1.49)
Passing to the operator formulation of quantum field theory (i.e., undoing the Wick rotation (1.28), which amounts to Osterwalder-Schrader reconstruction), the conserved currents $`\widehat{𝒥}_{\mathrm{}/r}^\mu `$ give rise to conserved charges,
$$N_{\mathrm{}/r}:=𝑑\underset{¯}{x}\widehat{𝒥}_{\mathrm{}/r}^0(t,\underset{¯}{x})$$
(1.50)
which (for gauge-transformations continuous at infinity) are gauge-invariant.
We should like to determine the equal-time commutators of the (gauge-invariant) currents $`𝒥_{\mathrm{}/r}^\mu (x)`$. Let $`𝒱`$ denote the affine space of configurations of external electromagnetic vector potentials, $`A`$, corresponding to static electromagnetic fields. We consider the Hilbert bundle, $``$, over $`𝒱`$ whose fibre, $`_A`$, at a point $`A𝒱`$ is the Fock space of state vectors of free, chiral (e.g., left-handed) fermions coupled to the vector potential $`A`$. Then $``$ carries a projective representation, $`U`$, of the group $`𝒢`$ of time-independent electromagnetic gauge transformations,
$$g=\left(g^\chi (x)\right),g^\chi (x)=e^{i\chi (x)},\chi (x)=\chi (\underset{¯}{x})(\mathrm{indep}.\mathrm{of}t),$$
with the following properties:
$`U(g):_A_{A+d\chi }`$ ,
and, on the fibre $`_{A+d\chi }`$ ,
$`U(g)\psi (x;A)|__AU(g)^1=e^{i\chi (x)}\psi (x;A+d\chi )`$ ,
where $`\psi (x;A)`$ is the Dirac spinor field acting on $`_A`$; (and similarly for $`\overline{\psi }(x;A)`$). The generator, $`G(\chi )`$, of the gauge transformation $`U(g^\chi ())`$ is given by
$$G(\chi )=𝑑\underset{¯}{x}\chi (\underset{¯}{x})G(x),$$
where
$$G(x)=i\underset{¯}{}\frac{\delta }{\delta \underset{¯}{A}(x)}+\frac{1}{q}𝒥_{\mathrm{}}^0(x;A).$$
(1.51)
Here
$$\underset{¯}{}\frac{\delta }{\delta \underset{¯}{A}}=\underset{j=1}{\overset{2n1}{}}_j\frac{\delta }{\delta A_j}.$$
Locally, the (phase) factor of the projective representation $`U`$ of $`𝒢`$ can be made trivial by redefining the generators $`G`$:
$$G(x)\widehat{G}(x):=i\underset{¯}{}\frac{\delta }{\delta \underset{¯}{A}(x)}+\frac{1}{q}\widehat{𝒥}_{\mathrm{}}^0(x;A).$$
(1.52)
The operators $`\widehat{G}(x)`$ generate a representation of the group $`𝒢`$ of gauge transformations on $``$ iff
$$[\widehat{G}(t,\underset{¯}{x}),\widehat{G}(t,\underset{¯}{y})]=\mathrm{\hspace{0.33em}0}$$
(1.53)
for all times $`t`$. That (1.52) is the right choice of generators compatible with (1.53) follows, heuristically, from the fact that $`\widehat{𝒥}_{\mathrm{}}^\mu (x;A)`$ is a conserved current.
Because the current $`𝒥_{\mathrm{}}^\mu (x;A)`$ is gauge-invariant, we have that
$$[\underset{¯}{}\frac{\delta }{\delta \underset{¯}{A}(x)},𝒥_{\mathrm{}}^0(y;A)]=\mathrm{\hspace{0.33em}0}.$$
(1.54)
Thus, using (1.48), (1.54) and (1.53), we find that
$`0`$ $`\stackrel{!}{=}`$ $`[\widehat{G}(t,\underset{¯}{x}),\widehat{G}(t,\underset{¯}{y})]`$ (1.55)
$`=`$ $`{\displaystyle \frac{1}{q^2}}[𝒥_{\mathrm{}}^0(t,\underset{¯}{x}),𝒥_{\mathrm{}}^0(t,\underset{¯}{y})]2i\underset{¯}{}{\displaystyle \frac{\delta }{\delta A(t,\underset{¯}{x})}}\omega ^0(t,\underset{¯}{y};A)`$
$`+\mathrm{\hspace{0.33em}2}i\underset{¯}{}{\displaystyle \frac{\delta }{\delta A(t,\underset{¯}{y})}}\omega ^0(t,\underset{¯}{x};A).`$
This equation determines the anomalous commutator
$$[𝒥_{\mathrm{}}^0(t,\underset{¯}{x}),𝒥_{\mathrm{}}^0(t,\underset{¯}{y})]=[\widehat{𝒥}_{\mathrm{}}^0(t,\underset{¯}{x}),\widehat{𝒥}_{\mathrm{}}^0(t,\underset{¯}{y})].$$
(1.56)
Of course, our arguments are heuristic, but, hopefully, provide a reasonably clear idea about the origin of anomalous commutators. A more erudite, mathematically clean derivation of (1.55) can be based on an analysis of the cohomology of $`𝒢`$; see e.g. .
In order to arrive at explicit versions of eqs. (1.46), (1.47) and (1.55) in various even dimensions, we must know the explicit expressions for the index density $`𝒜(x)`$ and the one-form $`\omega (x;A)`$. We shall not have any occasion to consider systems coupled to a non-trivial chiral gauge field $`Z`$. We therefore set $`Z=0`$. Then, in two space-time dimensions,
$$𝒜(x)=\frac{1}{4\pi }E(x),$$
(1.57)
by comparison of (1.45) with (1.20), and, by (1.48) and (1.57),
$$\omega ^\mu (x;A)=\frac{1}{4\pi }\epsilon ^{\mu \nu }A_\nu (x),$$
(1.58)
see also (1.23) and (1.17). In four space-time dimensions
$`𝒜(x)`$ $`=`$ $`{\displaystyle \frac{1}{32\pi ^2}}\left(FF\right)(x)`$ (1.59)
$`=`$ $`{\displaystyle \frac{1}{32\pi ^2}}F_{\mu \nu }(x)\stackrel{~}{F}^{\mu \nu }(x),`$
where $``$ denotes the exterior product and $``$ the “Hodge dual”. By eq. (1.47),
$$\omega ^\mu (x;A)=\frac{1}{32\pi ^2}\epsilon ^{\mu \nu \lambda \rho }A_\nu (x)F_{\lambda \rho }(x).$$
(1.60)
Thus eqs. (1.47) read
$$_\mu 𝒥_{\mathrm{}/r}^\mu (x)_A=\frac{q}{32\pi ^2}\left(FF\right)(x),$$
(1.61)
and, from eqs. (1.55), (1.56) and (1.60), we conclude that
$`[𝒥_{\mathrm{}/r}^0(t,\underset{¯}{x}),𝒥_{\mathrm{}/r}^0(t,\underset{¯}{y})]=[\widehat{𝒥}_{\mathrm{}/r}^0(t,\underset{¯}{x}),\widehat{𝒥}_{\mathrm{}/r}^0(t,\underset{¯}{y})]`$
$`=\pm i{\displaystyle \frac{q^2}{4\pi ^2}}\left(\underset{¯}{B}(\underset{¯}{x},t)\underset{¯}{}\right)\delta \left(\underset{¯}{x}\underset{¯}{y}\right),`$ (1.62)
a well known result; see .
The key fact reviewed in this section, from which all other results can be derived, is eq. (1.41), i.e.,
$$S_{\mathrm{eff}}^E(A+d\chi ,Z+d\alpha )=S_{\mathrm{eff}}^E(A,Z)2id^{2n}x\alpha (x)𝒜(x).$$
(1.63)
In order to describe a system in which only the left-handed fermions are charged, while the right-handed fermions are neutral, one may just set
$$A=Z=a$$
(1.64)
in eq. (1.63), where $`a`$ is the electromagnetic vector potential to which the left-handed fermions are coupled; see (1.25), (1.26) and (1.46). Denoting the effective action of this system by $`W_{\mathrm{}}(a)`$, we find from (1.63) and (1.64) that
$$W_{\mathrm{}}\left(a+d\chi \right)=W_{\mathrm{}}(a)+\mathrm{\hspace{0.17em}2}id^{2n}x\chi (x)𝒜(x).$$
(1.65)
Similarly,
$$W_r\left(a+d\chi \right)=W_r(a)2id^{2n}x\chi (x)𝒜(x),$$
(1.66)
for charged right-handed fermions.
Eqs. (1.65) and (1.66) show that a theory of massless chiral fermions coupled to an external electromagnetic field is anomalous, in the sense that it fails to be gauge-invariant. But let us imagine that space-time, $`^{2n}`$, is the boundary of a $`(2n+1)`$-dimensional half-space $`M`$, (i.e., $`M`$ = physical space-time $`^{2n}`$). Let $`A`$ denote a smooth U(1)-gauge potential on $`M`$ which is continuous on $`M`$, with
$$A|_M=a.$$
(1.67)
Let $`\omega ^{2n+1}(;A)`$ denote the usual Chern-Simons $`(2n+1)`$-form on $`M`$. The Chern-Simons action on $`M`$ is defined by
$$S_{CS}(A):=i_M\omega ^{2n+1}(\xi ;A),$$
(1.68)
where $`\xi `$ denotes a point in $`M`$. It should be recalled that
$$\omega ^{2n+1}(;A+d\chi )=\omega ^{2n+1}(;A)+d\chi (𝒜).$$
(1.69)
Since $`d(𝒜)=0`$, $`d\chi (𝒜)=d\left(\chi (𝒜)\right)`$, and hence, by Stokes’ theorem,
$`S_{CS}\left(A+d\chi \right)`$ $`=`$ $`S_{CS}(A)+i{\displaystyle _M}\chi (x)(𝒜)(x)`$ (1.70)
$`=`$ $`S_{CS}(A)+i{\displaystyle _M}d^{2n}x\chi (x)𝒜(x).`$
It follows that
$$W_{\mathrm{}/r}(a)S_{CS}(A)\mathrm{is}gaugeinvariant.$$
(1.71)
This result has a $`(2n+1)`$-dimensional interpretation (see ): Consider a (parity-violating) theory of massive, charged fermions described by $`2^n`$-component Dirac spinors on a $`(2n+1)`$-dimensional space-time $`M`$ with non-empty boundary $`M`$. These fermions are minimally coupled to an external electromagnetic vector potential $`A`$. We impose some anti-selfadjoint spectral boundary conditions on the $`(2n+1)`$-dimensional, covariant Dirac operator $`D_A`$. The action of the system is given by
$$S(\overline{\psi },\psi ;A):=_Md^{2n+1}\xi \overline{\psi }(\xi )\left(D_A+m\right)\psi (\xi ),$$
(1.72)
where $`m`$ is the bare mass of the fermions. The effective action of the system is defined by
$`e^{S_{\mathrm{eff}}^E(A)}`$ $`:=`$ $`\left[{\displaystyle 𝒟\psi 𝒟\overline{\psi }e^{S^E(\overline{\psi },\psi ;A)}}\right]_{\mathrm{ren}}`$ (1.73)
$`=`$ $`\mathrm{det}_{\mathrm{ren}}\left(D_A+m\right),`$
where the subscript “ren” indicates that renormalization may be necessary to define the R.S. of (1.73). Actually, for $`n=1`$, no renormalization is necessary; but, for $`n=2`$, e.g. an infinite charge renormalization must be made. It turns out that, for $`n=1`$ and $`n=2`$ (after renormalization),
$$S_{\mathrm{eff}}^E(A)=W_{\mathrm{}/r}\left(A|_M\right)S_{CS}(A)+𝒪\left(\frac{1}{m}\right),$$
(1.74)
up to a Maxwell term depending on renormalization conditions, where the correction terms are manifestly gauge-invariant; see . (Whether the R.S. of (1.74) involves $`W_{\mathrm{}}`$ or $`W_r`$ depends on the definition of $`D_A`$).
The physical reason underlying the result claimed in eq. (1.74) is that, in a system of massive fermions described by $`2^n`$-component Dirac spinors confined to a space-time $`M`$ with a non-empty, $`2n`$-dimensional boundary $`M`$, there are massless, chiral fermionic surface modes propagating along $`M`$.
This completes our heuristic review of aspects of the chiral abelian anomaly that are relevant for the physical applications to be discussed in subsequent sections. The abelian anomaly is, of course, but a special case of the general theory of anomalies involving also non-abelian, gravitational, global, …anomalies. In recent years, this theory has turned out to be important in connection with the theory of branes in string theory and with understanding aspects of $`M`$-theory. But, in this review, such applications will not be described.
In Sect. 2, we describe physical systems, important features of which can be understood as consequences of the two-dimensional chiral anomaly: incompressible (quantum) Hall fluids and ballistic wires.
In Sect. 3, we describe degrees of freedom in four dimensions which may play an important rôle in the generation of seeds for cosmic magnetic fields in the very early universe. This will turn out to be connected with the four-dimensional chiral anomaly.
In Sect. 4, a brief review of the theory of “transport in thermal equilibrium through gapless modes” developed in is presented.
In Sects. 5 and 6, we combine the results of this section with those in Sect. 4 to derive physical implications of the chiral anomaly for the systems introduced in Sects. 2 and 3.
Some conclusions and open problems are described in Sect. 7.
## 2 Quantized conductances
The original motivation for the work described in this review has been to provide simple and conceptually clear explanations of various formulae for quantized conductances, which have been encountered in the analysis of experimental data. Here are some typical examples.
Example 1. Consider a quantum Hall device with, e.g., an annular (Corbino) geometry. Let $`V`$ denote the voltage drop in the radial direction, between the inner and the outer edge, and let $`I_H`$ denote the total Hall current in the azimuthal direction. The Hall conductance, $`G_H`$, is defined by
$$G_H=I_H/V.$$
(2.1)
One finds that if the longitudinal resistance vanishes (i.e., if the two-dimensional electron gas in the device is “incompressible”) then $`G_H`$ is a rational multiple of $`\frac{e^2}{h}`$, i.e.,
$$G_H=\frac{n}{d}\frac{e^2}{h},n=\mathrm{\hspace{0.33em}0},1,2,\mathrm{},d=\mathrm{\hspace{0.33em}1},2,3,\mathrm{}.$$
(2.2)
In (2.2), $`e`$ denotes the elementary electric charge and $`h`$ denotes Planck’s constant. Well established Hall fractions, $`f_H:=\frac{n}{d}`$, in the range $`0<f_H1`$ are listed in Fig. 1; (see ; and for general background).
Example 2. In a ballistic (quantum) wire, i.e., in a pure, very thin wire without back scattering centers, one finds that the conductance $`G_W=I/V`$ ($`I`$: current through the wire, $`V`$: voltage drop between the two ends of the wire) is given by
$$G_W=\mathrm{\hspace{0.33em}2}N\frac{e^2}{h},N=\mathrm{\hspace{0.33em}0},1,2,\mathrm{},$$
(2.3)
under suitable experimental conditions (“small” $`V`$, temperature not “very small”, “adiabatic gates”); see .
Example 3. In measurements of heat conduction in quantum wires, one finds that the heat current is an integer multiple of a “fundamental” current which depends on the temperatures of the two heat reservoirs at the ends of the wire.
If electromagnetic waves are sent through an “adiabatic hole” connecting two half-spaces one approximately finds an “integer quantization” of electromagnetic energy flux.
Our task is to attempt to provide a theoretical explanation of these remarkable experimental discoveries; hopefully one that enables us to predict further related effects.
Conductance quantization is observed in a rather wide temperature range. It appears that it is only found in systems without dissipative processes. When it is observed it is insensitive to small changes in the parameters specifying the system and to details of sample preparation; i.e., it has universality properties. — It will turn out that the key feature of systems exhibiting conductance quantization is that they have conserved chiral charges; (such conservation laws will only hold approximately, i.e., in slightly idealized systems). Once one has understood this point, the right formulae follow almost automatically, and one arrives at natural generalizations.
In order to give a first indication how the effects described here might be related to the two-dimensional chiral anomaly, we consider Example 1, the quantum Hall effect, in more detail. For readers not familiar with this remarkable effect , we summarize some of its key features.
A quantum Hall fluid (QHF) is an interacting electron gas confined to some domain in a two-dimensional plane (an interface between a semiconductor and an insulator, with compensating background charge) subject to a constant magnetic field $`\stackrel{}{B}^{(0)}`$ transversal to the confinement plane. Among experimental control parameters is the filling factor, $`\nu `$, defined by
$$\nu =\frac{n^{(0)}}{B^{(0)}/\left(\frac{hc}{e}\right)}$$
where $`n^{(0)}`$ is the (constant) electron density, $`B^{(0)}`$ is the component of the magnetic field $`\stackrel{}{B}^{(0)}`$ perpendicular to the plane of the fluid, and $`\frac{hc}{e}`$ is the quantum of magnetic flux. The filling factor $`\nu `$ is dimensionless.
Transport properties of a QHF in an external electric field (of small frequency) are described by the equation
$$\underset{¯}{J}(t,\underset{¯}{x})=\left(\genfrac{}{}{0pt}{}{\sigma _L\sigma _H}{\sigma _H\sigma _L}\right)\underset{¯}{E}(t,\underset{¯}{x}),$$
(2.4)
where $`\underset{¯}{x}`$ is a point in the sample, $`\underset{¯}{J}`$ is the bulk electric current parallel to the sample plane and $`\underset{¯}{E}`$ is the component of the external electric field parallel to the sample plane. Furthermore, $`\sigma _L`$ denotes the longitudinal conductivity, and $`\sigma _H`$ is the transverse – or Hall conductivity. In two dimensions, conductances and conductivities have the same dimension of \[(charge)<sup>2</sup>/action\], and it is not difficult to see that
$$G_H=\sigma _H.$$
(2.5)
Experimentally, one observes that the longitudinal conductivity, $`\sigma _L`$, vanishes when the filling factor $`\nu `$ belongs to certain small intervals , a sign that there are no dissipative processes in the fluid. Such a QHF is called “incompressible”, for reasons explained below. Furthermore, on every interval of $`\nu `$ where $`\sigma _L`$ vanishes, the Hall conductivity $`\sigma _H`$ is a rational multiple of $`\frac{e^2}{h}`$ , as claimed in (2.2).
Next, we recall the basic equations of the electrodynamics of an incompressible QHF; see . It is useful to combine the two-dimensional space of the fluid and time to a three-dimensional space-time. The electromagnetic field tensor of the system is given by
$$F_{\mu \nu }=\left(\begin{array}{ccc}0\hfill & E_1\hfill & E_2\hfill \\ E_1\hfill & 0\hfill & B\hfill \\ E_2\hfill & B\hfill & 0\hfill \end{array}\right),$$
(2.6)
where $`E_1`$ and $`E_2`$ are the components of an external electric field in the plane of the sample, and $`B`$ is the component of an external magnetic field, $`\stackrel{}{B}`$, perturbing the constant field $`\stackrel{}{B}^{(0)}`$ perpendicular to the sample plane; $`\left(\stackrel{}{B}_{\mathrm{total}}=\stackrel{}{B}^{(0)}+\stackrel{}{B}\right)`$.
We define $`J^0(x)`$ to denote the sum of the electron charge density in the space-time point $`x=(t,\underset{¯}{x})`$ and the uniform background charge density $`en^{(0)}`$. We set $`J^\mu =(J^0,\underset{¯}{J})`$.
From the three-dimensional homogeneous Maxwell equations (Faraday’s law),
$$_\mu F_{\nu \lambda }+_\nu F_{\lambda \mu }+_\lambda F_{\mu \nu }=\mathrm{\hspace{0.33em}0},$$
(2.7)
the continuity equation for the electric current density (conservation of electric charge),
$$_\mu J^\mu =\mathrm{\hspace{0.33em}0},$$
(2.8)
and from the transport equation (2.4) with $`\sigma _L=0`$, it follows that
$$J^0=\sigma _HB.$$
(2.9)
Equations (2.4), for $`\sigma _L=0`$, and (2.9) can be combined to the equation
$$J^\mu =\sigma _H\epsilon ^{\mu \nu \lambda }F_{\nu \lambda }$$
(2.10)
of Chern-Simons electrodynamics, . Eqs. (2.10) describe the response of an incompressible QHF to an external electromagnetic field (perturbing the constant magnetic field $`\stackrel{}{B}^{(0)}`$).
Unfortunately, eqs. (2.10) are compatible with the continuity equation (2.8) for $`J^\mu `$ only if $`\sigma _H`$ is constant throughout space-time. But realistic samples have a finite extension.
The finite extension of the sample, confined to a space-time region $`\mathrm{\Omega }=D\times `$, where $`D`$ is e.g. a disk or an annulus, is taken into account by setting the Hall conductivity $`\sigma _H()`$ to zero outside $`\mathrm{\Omega }`$, i.e.,
$$\sigma _H(\xi )=\sigma _H\chi _\mathrm{\Omega }(\xi ),$$
(2.11)
for $`\xi ^3`$, where $`\sigma _H`$ is the (constant) value of the Hall conductivity inside the sample, and $`\chi _\mathrm{\Omega }`$ is the characteristic function of $`\mathrm{\Omega }`$. Taking the divergence of eq. (2.10), we get that
$$_\mu J^\mu =\sigma _H\epsilon ^{\mu \nu \lambda }\left(_\mu \chi _\mathrm{\Omega }\right)F_{\nu \lambda },$$
(2.12)
i.e., $`_\mu J^\mu `$ fails to vanish on the boundary, $`D`$, of the sample. However, conservation of electric charge is a fundamental law of nature for closed systems. Thus, there must be an electric current, $`J_\mathrm{\Omega }`$, localized on the boundary $`\mathrm{\Omega }`$ of the sample space-time such that the total electric current
$$J_{\mathrm{total}}^\mu =J^\mu +J_\mathrm{\Omega }^\mu $$
(2.13)
satisfies the continuity equation. The boundary current $`J_\mathrm{\Omega }^\mu `$ must be tangential to the boundary $`\mathrm{\Omega }`$ of the sample space-time. Hence it determines a current density, $`I^\alpha `$, on the (1+1)-dimensional space-time $`\mathrm{\Omega }`$, where the index $`\alpha `$ refers to a choice of coordinates on $`\mathrm{\Omega }`$. Eq. (2.12) and the continuity equation for $`J_{\mathrm{total}}^\mu `$ then imply that
$$_\alpha I^\alpha =\sigma _H\epsilon ^{\alpha \beta }F_{\alpha \beta }.$$
(2.14)
This equation identifies $`I^\alpha `$ as an anomalous current. Thus, there must be chiral modes (left-movers or right-movers, depending on the orientation of $`\mathrm{\Omega }`$ and the direction of the external magnetic field) propagating along the boundary. They carry the well known diamagnetic edge currents. If $`𝒥_{\mathrm{}}^\alpha `$ (or $`𝒥_r^\alpha `$) denotes the corresponding quantum-mechanical current operator then the edge current $`I^\alpha `$ is given by the quantum-mechanical expectation value, $`𝒥_{\mathrm{}/r}^\alpha _A`$, of $`𝒥_{\mathrm{}}^\alpha `$ (or $`𝒥_r`$). The currents $`𝒥_{\mathrm{}}^\alpha `$ have the anomalous commutators
$$[𝒥_{\mathrm{}}^0(t,\underset{¯}{x}),𝒥_{\mathrm{}}^0(t,\underset{¯}{y})]=\frac{i\sigma _H}{2\pi }\delta ^{}\left(\underset{¯}{x}\underset{¯}{y}\right),$$
(2.15)
see eqs. (1.21) and (1.16), and hence generate a chiral $`\widehat{u}(1)`$-current algebra with central charge given by $`\sigma _H`$.
We now return to the physics of the bulk of an incompressible QHF. The absence of dissipation $`(\sigma _L=0)`$ in the transport of electric charge through the bulk can be explained by the existence of a mobility gap in the energy spectrum between the ground state energy of the QHF and the energies of extended, excited bulk states. This property motivates the term “incompressible”: It is not possible to add an additional electron to, or subtract one from the fluid by injecting only an arbitrarily small amount of energy. An important consequence of incompressibility is that the total electric charge is a good quantum number to label different sectors of physical states of an incompressible QHF (at zero temperature).
We propose to study the bulk physics of incompressible QHF’s in the scaling limit, in order to describe the universal transport laws of such fluids. For this purpose, we consider a QHF confined to a sample of diameter $`\theta `$, where $`\theta `$ is a dimensionless scale factor. The scaling limit is the limit where $`\theta \mathrm{}`$, with distances and time rescaled by a factor $`\theta ^1`$. In rescaled coordinates, the fluid is thus confined to a sample of constant finite diameter.
The presence of a positive mobility gap in the system implies that, in the scaling limit, the effective theory describing an incompressible QHF must be a “topological field theory”. The states of a topological field theory are indexed by static, pointlike sources localized in the bulk and labelled by certain charge quantum numbers which generate a fusion ring; see .
It is not difficult to find the effective action, $`S_{\mathrm{eff}}(A)`$, in the scaling limit, where $`A`$ is the electromagnetic vector potential of the external electromagnetic field $`F_{\mu \nu }`$, see eq. (2.6). A possible starting point is eq. (2.10), relating the expectation value of the electric current to the external electromagnetic field:
$$J^\mu (\xi )=\delta S_{\mathrm{eff}}(A)/\delta A_\mu (\xi )=\sigma _H\epsilon ^{\mu \nu \lambda }F_{\nu \lambda }(\xi ).$$
(2.16)
The solution of eq. (2.16) is
$$S_{\mathrm{eff}}(A)=\sigma _HS_{CS}(A)=\frac{\sigma _H}{2}_\mathrm{\Omega }d^3\xi \epsilon ^{\mu \nu \lambda }A_\mu (\xi )_\nu A_\lambda (\xi ),$$
(2.17)
i.e., $`S_{\mathrm{eff}}`$ is proportional to the Chern-Simons action $`S_{CS}`$. The Chern-Simons action is not invariant under gauge transformations of $`A`$ that do not vanish on the boundary $`\mathrm{\Omega }`$ of the sample. Since electromagnetic gauge invariance is a fundamental property of quantum-mechanical systems, eq. (2.17) for $`S_{\mathrm{eff}}(A)`$ must be corrected by a boundary term. Let $`a`$ denote the restriction of $`A`$ to the boundary $`\mathrm{\Omega }`$ of the sample. Then, as pointed out in eq. (1.71), the expression $`W_{\mathrm{}/r}(a)S_{CS}(A)`$ is gauge-invariant, where $`W_{\mathrm{}/r}(a)`$ is the effective action of charged chiral modes propagating along $`\mathrm{\Omega }`$. Thus, in the scaling limit,
$$S_{\mathrm{eff}}(A)=\sigma _H\left[W_{\mathrm{}/r}(a)+S_{CS}(A)\right],$$
(2.18)
(depending on the sign of $`\sigma _H`$). It is well known that the action $`W_{\mathrm{}/r}(a)`$ is the generating function for the connected Green functions of the chiral current operators, $`𝒥_{\mathrm{}/r}^\alpha `$, on $`\mathrm{\Omega }`$, which generate a $`\widehat{u}(1)`$-current algebra. Formula (2.18) plays an important rôle in understanding the physics of incompressible quantum Hall fluids.
In the next section, we consider systems of massless chiral modes in four-dimensional space-time, with physical properties some of which are related to the four-dimensional chiral anomaly, and which may play a significant rôle in the physics of the early universe.
## 3 Branes, axions and charged fermions
The very early universe is filled with a hot plasma of charged leptons, quarks, gluons, photons, … . At a time after the big bang when the temperature $`T`$ is of the order of 80 $`TeV`$ chirality flips of light charged leptons, in particular of right-handed electrons, constitute a dynamical process slower than the expansion rate of the universe. Thus, for $`T`$ 80 $`TeV`$, the chiral charges, $`N_{\mathrm{}}`$ and $`N_r`$, defined in eq. (1.50) of Sect. 1, are approximately conserved for electrons. They are related to an approximate chiral symmetry of the electronic sector of the standard model. Among other results, we shall attempt to show that if, in the very early universe, the chemical potentials of left-handed and right-handed electrons are different from each other, this may give rise to the generation of large, cosmic magnetic fields, ; (see also for a similar, independent suggestion). This effect is, in a sense explained in Sects. 4 and 6, an effect in equilibrium statistical mechanics. However, this is precisely what may make it appear quite unnatural and implausible: The chiral charges, $`N_{\mathrm{}}`$ and $`N_r`$, are not really conserved; leptons are massive. The very early universe is not really in an equilibrium state, and the chemical potentials of left-handed and right-handed electrons neither have an unambiguous meaning, nor would they be space- and time-independent. It may then be wrong, or, at least, misleading, to invoke results from equilibrium statistical mechanics to explore effects in the physics of the very early universe.
A way out from these difficulties can be found by seeking inspiration from an analogy with the quantum Hall effect: Consider a quantum Hall fluid (QHF), confined to a strip of macroscopic width $`\mathrm{}`$ in the plane. If the QHF is incompressible then there are no light (gapless) modes propagating through the bulk of the sample; but, as shown in the last section, there are gapless, chiral modes propagating along the boundaries of the sample. Let $`\mathrm{\Omega }`$ denote the space-time of the fluid; it is a slab of width $`\mathrm{}`$ in three-dimensional Minkowski space. The two components of the boundary, $`\mathrm{\Omega }`$, of $`\mathrm{\Omega }`$ are denoted by $`_+\mathrm{\Omega }`$, $`_{}\mathrm{\Omega }`$, respectively. As shown in the last section, eq. (2.18), (see also for more details) the effective action of such an incompressible QHF (in the scaling limit) is given by
$$S_{\mathrm{eff}}(A)=\sigma _H\left[W_{\mathrm{}}\left(a_+\right)+W_r\left(a_{}\right)S_{CS}(A)\right],$$
(3.1)
(if the direction of the external magnetic field $`\stackrel{}{B}^{(0)}`$ is chosen appropriately, given an orientation of $`\mathrm{\Omega }`$). In (3.1), $`A`$ is an external electromagnetic vector potential on $`\mathrm{\Omega }`$, and
$$a_\pm :=A|_{_\pm \mathrm{\Omega }},$$
(3.2)
is the restriction of the 1-form $`A`$ to a component, $`_\pm \mathrm{\Omega }`$, of the boundary of $`\mathrm{\Omega }`$; $`W_{\mathrm{}/r}()`$ is the two-dimensional, anomalous effective action for charged, chiral (left-moving, or right-moving, respectively) surface modes propagating along $`_+\mathrm{\Omega },_{}\mathrm{\Omega }`$, respectively; and $`S_{CS}()`$ is the three-dimensional topological Chern-Simons action, see (2.17). Many universal features of the quantum Hall effect can be derived directly from eq. (3.1).
Suppose, in analogy to what we have just discussed, that the world, as known to us, is a movie showing the dynamics of light modes propagating along two parallel 3-branes in a five-dimensional space-time, $`M`$. More precisely, we imagine that $`M`$ is a slab of width $`\mathrm{}`$ in five-dimensional space-time, $`^5`$, the two components, $`_+M`$ and $`_{}M`$, of the boundary of $`M`$ being identified with the two parallel 3-branes. Let us imagine that, through the five-dimensional bulk $`M`$ of the system, a massive, charged, four-component spinor field $`\psi `$ propagates. We consider the response of this system to coupling the charged fermions described by $`\psi `$ to a five-dimensional, external electromagnetic vector potential, $`\widehat{A}`$. By $`A_\pm `$ we denote the four-dimensional vector potentials on $`_\pm M`$ obtained by restricting $`\widehat{A}`$ to $`_\pm M`$. As discussed at the end of Sect. 1, there are chiral, left-handed or right-handed, charged, fermionic surface modes propagating along $`_+M`$, $`_{}M`$, which are coupled to $`A_+,A_{}`$, respectively; see . In eq. (1.74), the effective action of this system has been reported. It is given by
$`S_{\mathrm{eff}}^E(\widehat{A})=W_{\mathrm{}}(A_+)+W_r(A_{})S_{CS}(\widehat{A})`$
$`+(4\mathrm{}e^2)^1{\displaystyle _M}d^5\xi F_{\widehat{A}}(\xi )^2+\mathrm{},`$ (3.3)
where the dots stand for terms $`O\left(\frac{1}{m}\right)`$ , and the renormalization conditions have been chosen in such a way that the constant $`e^2`$ in front of the five-dimensional Maxwell term is the four-dimensional feinstructure constant. The components, $`\widehat{A}_K`$, of $`\widehat{A}`$ are denoted by
$$\widehat{A}_\mu =:A_\mu ,\mu =\mathrm{\hspace{0.17em}0},1,2,3,\widehat{A}_4=:\phi ,$$
(3.4)
i.e., $`(\widehat{A}_K)=(A,\phi ),K=0,1,2,3,4.`$
In order to make contact with the laws of physics in four space-time dimensions, we should insist on the requirement that left-handed and right-handed fermions propagating along $`_+M`$ and $`_{}M`$, respectively, couple to the same electromagnetic vector potential, i.e., that
$$A_+\left(x,x^4=\mathrm{}\right)=A_{}\left(x,x^4=0\right)A(x).$$
(3.5)
This requirement is met if we assume that
$$\widehat{A}(x,x^4)\mathrm{is}independent\mathrm{of}x^4.$$
(3.6)
In this case,
$`S_{CS}(\widehat{A})`$ $`=`$ $`{\displaystyle \frac{i\mathrm{}}{32\pi ^2}}{\displaystyle _N}\phi \left(F_AF_A\right)`$ (3.7)
$`=`$ $`{\displaystyle \frac{i\mathrm{}}{32\pi ^2}}{\displaystyle _N}d^4x\phi (x)\epsilon ^{\mu \nu \lambda \rho }F_{\mu \nu }(x)F_{\lambda \rho }(x)`$
where $`N^4`$ is a slice through $`M`$ parallel to $`_\pm M`$, $`\mu ,\nu =0,1,2,3,`$ and $`F_A=(F_{\mu \nu })`$ is the four-dimensional field tensor; (the trivial integration over $`x^4`$ has produced the factor $`\mathrm{}`$). Furthermore, the Maxwell term on the R.S. of (3.3) reduces to
$$\frac{1}{4e^2}\left\{_Nd^4xF_{\mu \nu }(x)F^{\mu \nu }(x)+\mathrm{\hspace{0.17em}2}_Nd^4x\left(_\mu \phi \right)(x)\left(^\mu \phi \right)(x)\right\}.$$
(3.8)
Finally,
$$W_{\mathrm{}}\left(A_+=A\right)+W_r\left(A_{}=A\right)=S_{\mathrm{eff}}^E(A),$$
(3.9)
with $`S_{\mathrm{eff}}^E(A)=S_{\mathrm{eff}}^E(A,Z=0)`$ as in eqs. (1.29), (1.30). Thus, the complete effective action of the system is given by
$`S_{\mathrm{eff}}^E(\phi ;A)`$ $`=`$ $`S_{\mathrm{eff}}^E(A)+{\displaystyle \frac{i\mathrm{}}{32\pi ^2}}{\displaystyle _N}\phi \left(F_AF_A\right)`$ (3.10)
$`+`$ $`{\displaystyle \frac{1}{4e^2}}\left\{{\displaystyle _N}d^4xF_A^2(x)+\mathrm{\hspace{0.17em}2}{\displaystyle _N}d^4x\left(\phi \right)^2(x)\right\}.`$
Clearly, there is something quite unnatural about this approach: It is conditions (3.5) and (3.6)! If $`A_+`$ were different from $`A_{}`$ then the fermionic effective action $`S_{\mathrm{eff}}^E(A)=S_{\mathrm{eff}}^E(A,Z=0)`$ would be replaced by $`S_{\mathrm{eff}}^E(A,Z)`$, where $`A=\frac{1}{2}(A_++A_{})`$ and $`Z=\frac{1}{2}(A_++A_{})`$. Thus the surface modes would not only couple to the electromagnetic field, but also to a chiral gauge field $`Z`$ for which there is no experimental evidence, and the gauge fields would sample a five-dimensional space-time.
These unnatural features can be avoided by following Connes’ formulation of gauge theories with fermions . Then the effective action displayed in eq. (3.10) can be reproduced as follows: One sets $`M=N\times _2`$, $`N^4`$ and treats the discrete “fifth dimension”, $`_2`$, by using elementary tools from non-commutative geometry . By adding a “non-commutative”, five-dimensional Chern-Simons action, as constructed in , to Connes’ version of the Yang-Mills action (for a U(1)-gauge field) and to the standard fermionic effective action, one can reproduce actions like the one in eq. (3.10); see . There is no room, here, to review the details of these constructions.
In analogy to what we have discussed above, one may argue that string theories arise as effective theories of surface modes propagating along 9-branes in an “eleven-dimensional” space-time, starting from eleven-dimensional $`M`$-theory, (with anomalies of the surface theories cancelled by certain eleven-dimensional Chern-Simons actions). One realization of this idea appears in . But we shall not pursue these ideas any further, in this review.
Instead, we ask whether the effective action in (3.10) ought to look familiar to people holding a conventional point of view that physical space-time is four-dimensional. The answer is “yes”! The scalar field $`\phi `$ appearing in the effective action on the R.S. of (3.10) can be interpreted as the axion. The axion field was originally introduced by Peccei and Quinn to solve the strong CP problem. There are various reasons, including, primarily, experimental ones, to feel unhappy about introducing an axion into the standard model. But there is also a good reason to do so: String theory predicts the existence of an axion, the “model-independent axion” first described by Witten .
The argument in favor of the model-independent axion goes as follows: String theory tells us that there must exist a second-rank antisymmetric tensor field, i.e., a two-form, $`B_{\mu \nu }`$. The gauge-invariant field strength, $`H`$, a three-form, corresponding to $`B`$ is given by
$$H=dB\omega _{3YM}+\omega _{3G},$$
(3.11)
where $`d`$ denotes exterior differentiation, and $`\omega _{3YM}`$ and $`\omega _{3G}`$ are the gauge-field (“Yang-Mills”) and gravitational (Lorentz) Chern-Simons three-forms. (The coefficients in front of these Chern-Simons forms are proportional to the number, $`N_f`$, of species of fermions coupled to the gauge- and gravitational fields. In the following we shall set $`N_f=1`$.) The field strength $`H`$ is invariant under the gauge transformations $`BB+d\lambda `$, where $`\lambda `$ is an arbitrary one-form, and under gauge- and local Lorentz transformations accompanied by shifts of $`B`$. The equation of motion of $`H`$ is
$$^\mu H_{\mu \nu \lambda }=\mathrm{\hspace{0.33em}0},$$
(3.12)
or $`\delta H=0`$, where $`\delta `$ is the co-differential. We consider the components of $`B_{\mu \nu }`$ with $`\mu ,\nu =0,\mathrm{},3`$ and assume that $`B`$ is independent of coordinates of internal dimensions (of the string theory target). Then, in four-dimensional (non-compact) space-time, the three-form $`H`$ is dual to a one-form, $`Z`$, and the equation of motion (3.12) becomes
$$_\mu Z_\nu _\nu Z_\mu =\mathrm{\hspace{0.33em}0},\mathrm{or}dZ=\mathrm{\hspace{0.33em}0}.$$
(3.13)
By Poincaré’s lemma,
$$Z_\mu =_\mu \alpha ,\mathrm{or}Z=d\alpha ,$$
(3.14)
where $`\alpha `$ is a scalar field. By (3.11), the scaling dimension of $`\alpha `$ is two. Introducing a constant, $`\mathrm{}`$, with the dimension of length, we set
$$\alpha =\frac{1}{\mathrm{}e^2}\phi ,$$
(3.15)
where $`\phi `$ has scaling dimension = 1; ($`e^2`$ is the feinstructure constant).
From $`d^2=0`$ and (3.11) we obtain the equation
$$dH(x)=𝒜(x)+\mathrm{const}.\mathrm{tr}(R(x)R(x))$$
(3.16)
where
$$𝒜(x)=\frac{i}{32\pi ^2}(F(x)F(x))$$
(3.17)
is the index density, see eq. (1.59), ($``$ denotes the Hodge dual), and $`R(x)`$ is the Riemann curvature tensor. Assuming that space-time is flat, hence $`R=0`$, and considering the special case, where the electromagnetic field is the only gauge field in the system, we obtain
$$dH=\frac{i}{32\pi ^2}\left(F_AF_A\right).$$
(3.18)
Recalling that
$$H=\left(\frac{1}{\mathrm{}e^2}d\phi \right),$$
see (3.13)–(3.15), we find that (3.18) yields the following equation of motion for $`\phi `$:
$$\mathrm{}\phi =\frac{i\mathrm{}e^2}{32\pi ^2}\left(F_AF_A\right).$$
(3.19)
This equation is the Euler-Lagrange equation corresponding to the action functional
$$\frac{1}{2e^2}d^4x\left(\phi \right)^2(x)+\frac{i\mathrm{}}{32\pi ^2}\phi \left(F_AF_A\right),$$
(3.20)
which reproduces the R.S. of (3.10), up to the fermionic effective action and the Maxwell term! The second term in (3.20) can be understood as arising from coupling fermions to the axion. The term in the bare action of the fermions describing their coupling to the axion is given by
$$\frac{\mathrm{}^2}{2}d^4xH_{\mu \nu \lambda }\overline{\psi }\gamma ^\mu \gamma ^\nu \gamma ^\lambda \psi =\frac{\mathrm{}}{2}d^4x_\mu \phi \overline{\psi }\gamma ^\mu \gamma \psi ,$$
(3.21)
where $`\gamma =\gamma ^5`$. Carrying out the Berezin integral over the fermionic degrees of freedom — see eq. (1.29) — we find an effective action for the fermions given by
$`S_{\mathrm{eff}}^E\left(A,Z={\displaystyle \frac{\mathrm{}}{2}}d\phi \right)`$ $`=`$ $`S_{\mathrm{eff}}^E\left(A,Z=0\right)i\mathrm{}{\displaystyle d^4x\phi (x)𝒜(x)}`$ (3.22)
$`=`$ $`S_{\mathrm{eff}}^E(A){\displaystyle \frac{i\mathrm{}}{32\pi ^2}}{\displaystyle \phi \left(F_AF_A\right)},`$
in accordance with (3.20). The first equation in (3.22) is eq. (1.41), the second follows from (1.59).
Thus, coupling charged Dirac fermions to an external electromagnetic vector potential $`A`$ and an axion $`\phi `$ yields the effective action (3.22). Adding to it the Maxwell term and the kinetic energy term for $`\phi `$, we again obtain the action (3.10)!
One may argue that, in any case, the presence of an axion in the theory may be an indication that there must exist extra (classical or, perhaps more plausibly, discrete or “non-commutative”) dimensions. But, for our applications in Sect. 6, this point is not important. What will matter is that the time derivative of the axion field will play the rôle of a, generally speaking, space-time dependent “chemical potential” for right-handed leptons.
But, quite independently of the properties of fermions (which, for example, may acquire masses through a Higgs-Kibble mechanism), the axion, $`\phi `$, will turn out to be the driving force for a possible generation of large cosmic magnetic fields.
As our discussion at the beginning of this section, up to eq. (3.10), has shown it is legitimate to view a four-dimensional system of fermions in an external electromagnetic and an external axion field as the four-dimensional analogue of the edge degrees of freedom of an incompressible quantum Hall fluid. It supports electric currents analogous to the diamagnetic edge currents of a quantum Hall fluid.
## 4 Transport in thermal equilibrium through gapless <br>modes
In this section we prepare the ground for a theoretical explanation of effects such as the ones described in Sects. 2 (Examples 1 through 3) and 3. We consider a quantum-mechanical system $`𝒮`$ whose dynamics is determined by a Hamiltonian $`H`$, which is a selfadjoint operator on the Hilbert space $``$ of pure state vectors of $`𝒮`$ with discrete energy spectrum. It is assumed that the system obeys conservation laws described by some conserved “charges” $`N_1,\mathrm{},N_L`$ commuting with all observables of the system. Hence
$$[H,N_{\mathrm{}}]=\mathrm{\hspace{0.33em}0},[N_{\mathrm{}},N_k]=\mathrm{\hspace{0.33em}0},\mathrm{},k=\mathrm{\hspace{0.33em}1},\mathrm{},L,$$
(4.1)
(e.g. in the sense that the spectral projections of $`H`$ and of $`N_{\mathrm{}},N_k`$ commute with one another, for all $`k`$ and $`\mathrm{}`$.) The system $`𝒮`$ is coupled to $`L`$ reservoirs, $`_1\mathrm{},_L`$, with the property that the expectation value of the conserved charge $`N_{\mathrm{}}`$ in a stationary state of $`𝒮`$ can be tuned to some fixed value through exchange of “quasi-particles” between $`𝒮`$ and $`_{\mathrm{}}`$, i.e., through a current between $`𝒮`$ and $`_{\mathrm{}}`$ that carries “$`N_{\mathrm{}}`$-charge”, for all $`\mathrm{}=1,\mathrm{},L`$ .
We are interested in describing a thermal equilibrium state of $`𝒮`$ coupled to $`_1,\mathrm{},_L`$, at a temperature $`T=(k_B\beta )^1`$. According to Gibbs, we should work in the grand-canonical ensemble. The reservoirs $`_1,\mathrm{}_L`$ then enter the description of the thermal equilibrium of $`𝒮`$ only through their chemical potentials $`\mu _1,\mathrm{},\mu _L`$. The chemical potential $`\mu _{\mathrm{}}`$, is a thermodynamic parameter canonically conjugate to the charge $`N_{\mathrm{}}`$; in particular, the dimension of $`\mu _{\mathrm{}}N_{\mathrm{}}`$ is that of an energy. According to Landau and von Neumann, the thermal equilibrium state of $`𝒮`$ at temperature $`(k_B\beta )^1`$ in the grand-canonical ensemble, with fixed values of $`\mu _1,\mathrm{},\mu _L`$, is given by the density matrix
$$\rho _{\beta ,\underset{¯}{\mu }}=\mathrm{\Xi }_{\beta ,\underset{¯}{\mu }}^1\mathrm{exp}\left[\beta \left(H\underset{\mathrm{}=1}{\overset{L}{}}\mu _{\mathrm{}}N_{\mathrm{}}\right)\right],$$
(4.2)
where the grand partition function $`\mathrm{\Xi }_{\beta ,\underset{¯}{\mu }}`$ is determined by the requirement that
$$\mathrm{Tr}\rho _{\beta ,\underset{¯}{\mu }}=\mathrm{\hspace{0.33em}1}.$$
(4.3)
(It is assumed here that $`\mathrm{exp}\left[\beta \left(H\mu _{\mathrm{}}N_{\mathrm{}}\right)\right]`$ is a trace-class operator on $``$, for all $`\beta >0`$; we are studying a system in a compact region of physical space.) The equilibrium expectation of a bounded operator, $`a`$, on $``$ is defined by
$$a_{\beta ,\underset{¯}{\mu }}:=\mathrm{Tr}\left(\rho _{\beta ,\underset{¯}{\mu }}a\right).$$
(4.4)
Let $`𝒥(x)=(𝒥^0(x),\underset{¯}{𝒥}(x))`$ be a conserved quantum-mechanical current density of $`𝒮`$, where $`x=(\underset{¯}{x},t)`$, $`t`$ is time and $`\underset{¯}{x}`$ is a point of physical space contained inside $`𝒮`$. We are interested in calculating the expectation values of products of components of $`𝒥`$ in the state $`\rho _{\beta ,\underset{¯}{\mu }}`$; in particular, we should like to calculate $`\underset{¯}{𝒥}(x)_{\beta ,\underset{¯}{\mu }}`$. Of course, if the dimension of space is larger than one, $`\underset{¯}{𝒥}(x)_{\beta ,\underset{¯}{\mu }}`$ vanishes unless rotation invariance is broken by some external field. If $`\underset{¯}{𝒥}(x)`$ is a vector current then $`\underset{¯}{𝒥}(x)_{\beta ,\underset{¯}{\mu }}`$ vanishes unless the state $`\rho _{\beta ,\underset{¯}{\mu }}`$ is not invariant under space-reflection and time reversal. This happens if some of the charges $`N_1,\mathrm{},N_L`$ are not invariant under space-reflection and time reversal, i.e., if they are chiral.
To say that $`𝒥`$ is conserved means that it satisfies the continuity equation
$$_\mu 𝒥^\mu =\mathrm{\hspace{0.33em}0},$$
(4.5)
where $`x^0=t`$ denotes time, and $`_\mu =/x^\mu `$. If the space-time of the system $`𝒮`$ is topologically trivial (“star-shaped”) then eq. (4.5) implies that there is a globally defined vector field $`\underset{¯}{\phi }(x)`$ such that
$$𝒥^0(x)=\frac{q}{2\pi }\mathrm{div}\underset{¯}{\phi }(x),\underset{¯}{𝒥}(x)=\frac{q}{2\pi }\frac{}{t}\underset{¯}{\phi }(x),$$
(4.6)
with $`q`$ the electric charge.
Let us suppose that $`\underset{¯}{\phi }(x)`$ is an operator-valued distribution on $``$, whose time-dependence is determined by the formal Heisenberg equation
$$\frac{}{t}\underset{¯}{\phi }(x)=\frac{i}{\mathrm{}}[H,\underset{¯}{\phi }(x)].$$
(4.7)
\[Technically, we are treading on somewhat slippery ground here; but we shall proceed formally, in order to explain the key ideas on a few pages.\] From (4.6) and (4.7) we derive that
$$\underset{¯}{𝒥}(x)_{\beta ,\underset{¯}{\mu }}=\frac{iq}{h}[H,\underset{¯}{\phi }(x)]_{\beta ,\underset{¯}{\mu }}.$$
(4.8)
Formally, the R.S. of (4.8) vanishes, because $`()_{\beta ,\underset{¯}{\mu }}`$ is a time-translation invariant state. However, the field $`\underset{¯}{\phi }`$ turns out to have ill-defined zero-modes, and it is not legitimate to pretend that $`[H,\underset{¯}{\phi }(x)]=H\underset{¯}{\phi }(x)\underset{¯}{\phi }(x)H`$, because both terms on the R.S. are divergent, due to the zero-modes of $`\underset{¯}{\phi }`$. What is legitimate is to claim that
$$\frac{}{t}\underset{¯}{\phi }(x)=\frac{i}{h}[H\underset{\mathrm{}=1}{\overset{L}{}}\mu _{\mathrm{}}N_{\mathrm{}},\underset{¯}{\phi }(x)]+\frac{i}{h}\underset{\mathrm{}=1}{\overset{L}{}}\mu _{\mathrm{}}[N_{\mathrm{}},\underset{¯}{\phi }(x)],$$
(4.9)
and that the expectation value
$$[H\underset{\mathrm{}=1}{\overset{L}{}}\mu _{\mathrm{}}N_{\mathrm{}},\underset{¯}{\phi }(x)]_{\beta ,\underset{¯}{\mu }}$$
vanishes. This can be seen by replacing the Hamiltonian $`H`$ by a regularized Hamiltonian $`H^{(\epsilon )}`$ generating a dynamics that eliminates the zero-modes of $`\underset{¯}{\phi }`$. One replaces the state $`\rho _{\beta ,\underset{¯}{\mu }}`$ by a regularized state $`\rho _{\beta ,\underset{¯}{\mu }}^{(\epsilon )}`$ proportional to $`\mathrm{exp}\left[\beta \left(H^{(\epsilon )}\mu _{\mathrm{}}N_{\mathrm{}}\right)\right]`$, and we set
$$a_{\beta ,\underset{¯}{\mu }}^{(\epsilon )}:=\mathrm{tr}\left(\rho _{\beta ,\underset{¯}{\mu }}^{(\epsilon )}a\right),$$
for any bounded operator $`a`$ on $``$. Then
$`\underset{¯}{𝒥}(x)_{\beta ,\underset{¯}{\mu }}`$ $`=`$ $`\underset{\epsilon 0}{lim}{\displaystyle \frac{iq}{h}}[H^{(\epsilon )}{\displaystyle \underset{\mathrm{}=1}{\overset{L}{}}}\mu _{\mathrm{}}N_{\mathrm{}},\underset{¯}{\phi }(x)]_{\beta ,\underset{¯}{\mu }}^{(\epsilon )}`$ (4.10)
$`+`$ $`\underset{\epsilon 0}{lim}{\displaystyle \underset{\mathrm{}=1}{\overset{L}{}}}{\displaystyle \frac{iq\mu _{\mathrm{}}}{h}}[N_{\mathrm{}},\underset{¯}{\phi }(x)]_{\beta ,\underset{¯}{\mu }}^{(\epsilon )}.`$
Obviously
$$[H^{(\epsilon )}\underset{\mathrm{}=1}{\overset{L}{}}\mu _{\mathrm{}}N_{\mathrm{}},\underset{¯}{\phi }(x)]_{\beta ,\underset{¯}{\mu }}^{(\epsilon )}=\mathrm{\hspace{0.33em}0},$$
(4.11)
and one might be tempted to expect that $`\underset{\epsilon 0}{lim}[N_{\mathrm{}},\underset{¯}{\phi }(x)]_{\beta ,\underset{¯}{\mu }}^{(\epsilon )}`$ vanishes, for all $`\mathrm{}`$, because the charges $`N_{\mathrm{}}`$ are conserved. However, as long as the regularization is present $`(\epsilon 0)`$, these charges are not conserved, and there is no guarantee that the second term on the R.S. of (4.10) vanishes!
We conclude that
$`\underset{¯}{𝒥}(x)_{\beta ,\underset{¯}{\mu }}`$ $`=`$ $`\underset{\epsilon 0}{lim}{\displaystyle \underset{\mathrm{}=1}{\overset{L}{}}}{\displaystyle \frac{iq\mu _{\mathrm{}}}{h}}[N_{\mathrm{}},\underset{¯}{\phi }(x)]_{\beta ,\underset{¯}{\mu }}^{(\epsilon )}`$ (4.12)
$`=:`$ $`{\displaystyle \underset{\mathrm{}=1}{\overset{L}{}}}{\displaystyle \frac{iq\mu _{\mathrm{}}}{h}}[N_{\mathrm{}},\underset{¯}{\phi }(x)]_{\beta ,\underset{¯}{\mu }}.`$
Eq. (4.12) might be called a current sum rule.
Let us assume that the conserved charges $`N_{\mathrm{}},\mathrm{}=1,2,\mathrm{},`$ are given as integrals of the 0-components of conserved currents over space. Then the current sum rule (4.12) implies that if $`\underset{¯}{𝒥}(x)_{\beta ,\underset{¯}{\mu }}0`$ there must be gapless modes in the system. The proof, see , is analogous to the proof of the Goldstone theorem in the theory of broken continuous symmetries.
The sum rule (4.12) is the main result of this section. A careful derivation of equation (4.12) and of our analogue of the Goldstone theorem could be given by using the operator-algebra approach to quantum statistical mechanics . But, in order to reach our punch line on a reasonable number of pages, we refrain from entering into a careful technical discussion.
## 5 Conductance quantization in ballistic wires and in incompressible quantum Hall fluids
In this section, we combine the results of Sects. 2 and 4, in order to gain insight into the phenomena of conductance quantization, as discussed at the beginning of Sect. 2. We first study a ballistic wire, i.e., a very thin, long, clean conductor without back scattering centers (impurities). The ends of the wire are connected to two reservoirs filled with electrons at chemical potentials $`\mu _{\mathrm{}},\mu _r`$, respectively, with
$$\mu _{\mathrm{}}\mu _r=V,$$
(5.1)
where $`V`$ is the voltage drop through the wire.
A ballistic wire is a three-dimensional, elongated metallic object with a tiny cross section in the plane perpendicular to its principal axis. Thus, at low temperature, the three-dimensional nature of the wire merely implies that there are several, say $`N`$, species of electrons labelled by discrete quantum numbers that originate from the motion in the plane perpendicular to the axis of the wire. Every species of electrons forms a one-dimensional Luttinger liquid , and these Luttinger liquids may interact with each other. Every Luttinger liquid has two conserved vector current operators, $`𝒥^{(i,s)\mu }`$, and conserved chiral current operators, $`\widehat{𝒥}_{\mathrm{}/r}^{(i,s)\mu }`$, where $`s=,`$ denotes the magnetic quantum number of the electrons in the i<sup>th</sup> Luttinger liquid (“spin up” and “spin down”), and $`i=1,\mathrm{},N`$. The chiral current operators $`\widehat{𝒥}_{\mathrm{}}^{(i,s)\mu }`$ are as in eqs. (1.21)–(1.23). The total electric current operator and the total chiral current operators are given by
$$𝒥^\mu =\underset{\genfrac{}{}{0pt}{}{i=1}{s=,}}{\overset{N}{}}𝒥^{(i,s)\mu },\widehat{𝒥}_{\mathrm{}/r}^\mu =\underset{\genfrac{}{}{0pt}{}{i=1}{s=,}}{\overset{N}{}}\widehat{𝒥}_{\mathrm{},r}^{(i,s)\mu }.$$
(5.2)
They are conserved. The total electric charge operators counting the electric charges of chiral (left-moving and right-moving) modes in the wire are the operators $`N_{\mathrm{}}`$ and $`N_r`$ defined in eq. (1.24). Their expectation values in a thermal equilibrium state of the wire are tuned by the chemical potentials, $`\mu _{\mathrm{}},\mu _r`$, respectively, of the reservoirs at the right and left end of the wire.
Imagine that the wire is kept at a constant temperature $`\beta ^1`$. Our description of the electron gas in the wire in terms of a finite number of Luttinger liquids correctly captures electric transport properties of the wire only if $`\beta ^1`$ and $`eV`$, with $`e`$ the elementary electric charge, are tiny as compared to the energy scale of the motion in the plane perpendicular to the axis of the wire. (However, $`\beta ^1`$ and $`eV`$ should be large as compared to the energy scale of weak back scattering centers.) We shall assume that these conditions are met. Then we may apply the current sum rule (4.12) derived in the last section, and the formulae for the anomalous commutators derived in Sect. 1, see (1.16) and the equation after (1.24), in order to calculate the electric current, $`I`$, in the wire corresponding to a voltage drop $`V`$. The current sum rule (4.12) yields
$`I`$ $`=`$ $`\underset{¯}{𝒥}(x)_{\beta ,\underset{¯}{\mu }}`$ (5.3)
$`=`$ $`{\displaystyle \frac{iq}{h}}\left\{\mu _{\mathrm{}}[N_{\mathrm{}},\underset{¯}{\phi }(x)]_{\beta ,\underset{¯}{\mu }}+\mu _r[N_r,\underset{¯}{\phi }(x)]_{\beta ,\underset{¯}{\mu }}\right\},`$
where $`\underset{¯}{\phi }`$ is the potential of the current $`𝒥^\mu `$. Since the currents $`𝒥^{(i,s)\mu }`$ of all the Luttinger liquids are conserved, every one of them can be derived from a potential, $`\phi ^{(i,s)}`$,
$$𝒥^{(i,s)\mu }(x)=\frac{q}{2\pi }\epsilon ^{\mu \nu }\left(_\nu \phi ^{(i,s)}\right)(x),$$
(5.4)
see eq. (1.11), and $`q=e`$, because the electric charge of an electron is equal to minus the elementary electric charge.
Plugging (5.4) and (5.2) into eq. (5.3) and recalling eq. (1.24) and the anomalous commutator
$`[\widehat{𝒥}_{\mathrm{}/r}^{(i,s)0}(\underset{¯}{y},t),\phi ^{(i^{},s^{})}(\underset{¯}{x},t)]`$
$`=\pm i{\displaystyle \frac{e}{2\pi }}\delta _{ii^{}}\delta _{ss^{}}\delta \left(\underset{¯}{x}\underset{¯}{y}\right),`$ (5.5)
see eqs. (1.11), (1.15), (1.16), we find that
$`I`$ $`=`$ $`{\displaystyle \frac{ie}{h}}{\displaystyle \underset{\genfrac{}{}{0pt}{}{i=1}{s=,}}{\overset{N}{}}}\{\mu _{\mathrm{}}[N_{\mathrm{}}^{(i,s)},\phi ^{(i,s)}(\underset{¯}{x},t)]_{\beta ,\underset{¯}{\mu }}`$ (5.6)
$`+\mu _r[N_r^{(i,s)},\phi ^{(i,s)}(\underset{¯}{x},t)]_{\beta ,\underset{¯}{\mu }}\}`$
$`=`$ $`{\displaystyle \frac{e^2}{h}}\times 2N\times \left(\mu _{\mathrm{}}\mu _r\right)`$
$`=`$ $`2N{\displaystyle \frac{e^2}{h}}V.`$
Thus, we have derived the formula
$$G_W=\frac{I}{V}=\mathrm{\hspace{0.33em}2}N\frac{e^2}{h},$$
(5.7)
as claimed in Example 2 at the beginning of Sect. 2.
Of course, the number, $`N`$, of Luttinger liquids of electrons in the wire depends on the mean Fermi energy of the wire (at zero temperature) and hence on the electron density in the wire and can be tuned.
The quantization of the Hall conductance of an incompressible Hall fluid in a Hall sample with e.g. an annular (Corbino) geometry (see Example 1) can be understood by using very similar arguments as in the example of quantum wires. Let $`V`$ denote the voltage drop between the outer and the inner edge of the sample. We assume that $`eV`$ and the temperature $`\beta ^1`$ are tiny, as compared to the mobility gap in the bulk of the fluid. Let us also assume, temporarily, that the electric field created by connecting the outer and inner edge to the two leads of a battery with voltage drop $`V`$ does not penetrate into the bulk of the sample (i.e., that, in the bulk, it is screened completely). If this assumption (which will actually turn out to be irrelevant, later) is made then the entire Hall current, $`I_H`$, in the sample is carried by the chiral modes propagating along the edges of the sample, i.e., $`I_H`$ is given by the expectation value of the sum, $`\widehat{𝒥}_{\mathrm{}}^1+\widehat{𝒥}_r^1`$, of the edge currents, $`\widehat{𝒥}_{\mathrm{}}^\mu ,\widehat{𝒥}_r^\mu `$. For an appropriate choice of orientation, $`\widehat{𝒥}_{\mathrm{}}^\mu `$ is the current at the outer edge and $`\widehat{𝒥}_r^\mu `$ is the current at the inner edge of the sample. The two edges are separated by the bulk, and, for a macroscopic sample, tunnelling of quasi-particles from one edge to the other one can be neglected for all practical purposes. This implies that the currents, $`\widehat{𝒥}_{\mathrm{}}^\mu `$ and $`\widehat{𝒥}_r^\mu `$, and hence the charge operators $`N_{\mathrm{}}`$ and $`N_r`$ defined in eq. (1.24), are conserved to very high accuracy. The anomalous commutators of $`\widehat{𝒥}_{\mathrm{}}^\mu `$ and $`\widehat{𝒥}_r^\mu `$ are given in eq. (2.15), and the analogue of Eqs. (1.11) and (5.4) is
$$𝒥^\mu (x)=e\frac{\sqrt{f_H}}{2\pi }\epsilon ^{\mu \nu }\left(_\nu \phi \right)(x).$$
(5.8)
Inserting these equations into the current sum rule (4.12), one finds that
$`I_H`$ $`=`$ $`{\displaystyle \frac{e}{h}}\sqrt{f_H}\left\{\mu _{\mathrm{}}[N_{\mathrm{}},\phi (\underset{¯}{x},t)]_{\beta ,\underset{¯}{\mu }}+\mu _r[N_r,\phi (\underset{¯}{x},t)]_{\beta ,\underset{¯}{\mu }}\right\}`$ (5.9)
$`=`$ $`{\displaystyle \frac{e^2}{h}}f_H\left(\mu _{\mathrm{}}\mu _r\right)`$
$`=`$ $`\sigma _HV.`$
These arguments do not make it clear why the Hall fraction $`f_H=\left(e^2/h\right)^1\sigma _H`$ is a rational number, and we have no clue, so far, which rational numbers may turn up in physical samples. Understanding the rational quantization of $`f_H`$ is not quite an easy matter; see . Here we can only sketch some key ideas. Let $`\psi (\underset{¯}{x},t)`$ denote the field (a “chiral vertex operator”) creating an electron or a hole propagating along the inner (or along the outer) edge of the sample. This field has the form
$$\psi (\underset{¯}{x},t)=:e^{iq\phi _{\mathrm{}/r}(\underset{¯}{x},t)}:\epsilon (\underset{¯}{x},t),$$
(5.10)
where $`q`$ is a real number to be determined, $`\phi _{\mathrm{}/r}(\underset{¯}{x},t)`$ is the potential of the conserved chiral edge current, i.e., it is a massless, chiral free field, and $`\epsilon (\underset{¯}{x},t)`$ is an electrically neutral so-called simple current of a rational chiral conformal field theory describing chiral modes of zero charge propagating along the edge. The field $`\psi (\underset{¯}{x},t)`$ must carry electric charge $`\pm e`$. Using formula (5.8) and recalling that $`\epsilon `$ has zero electric charge, we find that
$$q=\mathrm{\hspace{0.33em}1}/\sqrt{f_H}.$$
(5.11)
Furthermore, the field $`\psi (\underset{¯}{x},t)`$ must obey Fermi statistics (because electrons and holes are fermions). Hence it must have half-integer “conformal spin”, i.e.,
$$s_\psi \mathrm{\hspace{0.33em}1}/2\mathrm{mod}.1.$$
(5.12)
By eq. (5.10), the conformal spin of $`\psi `$ is given by
$$s_\psi =\frac{q^2}{2}+s_\epsilon =\frac{1}{2f_H}+s_\epsilon ,$$
(5.13)
where $`s_\epsilon `$ is the conformal spin of $`\epsilon `$. Because $`\epsilon `$ is a simple current of a rational chiral conformal field theory, $`s_\epsilon `$ is a rational number, i.e., $`s_\epsilon =\frac{k}{\mathrm{}}`$, with $`k`$ and $`\mathrm{}`$ two relatively prime integers. Thus (5.12) and (5.13) imply that
$$\frac{1}{2f_H}+\frac{k}{\mathrm{}}1/2\mathrm{mod}.1.$$
(5.14)
It follows that $`f_H`$ is a rational number. For more details see and, especially, . Properties of the rational chiral conformal field theories that may appear in the context of the quantum Hall effect are discussed in . One noteworthy result is that, unless $`f_H`$ is an integer, there must be chiral modes (quasi-particles) of fractional electric charge and fractional statistics, sometimes called Laughlin vortices, propagating along the edges of the sample.
Let us see what happens if the electric field $`\underset{¯}{E}`$ can penetrate into the bulk of an incompressible quantum Hall fluid. Electric transport in such Hall fluids can be understood by combining the arguments outlined above with Hall’s law in the bulk. The total Hall current, $`I_H`$, is given by
$$I_H=I_H^{\mathrm{edge}}+I_H^{\mathrm{bulk}},$$
(5.15)
where $`I_H^{\mathrm{edge}}`$ is the edge current studied above, and $`I_H^{\mathrm{bulk}}`$ is a current carried by extended bulk states. Let $`\gamma `$ denote an arbitrary smooth oriented curve connecting a point on the inner edge to a point on the outer edge of the sample. Then
$$I_H^{\mathrm{bulk}}=\underset{k,\mathrm{}}{}_\gamma J^k(\underset{¯}{x},t)\epsilon _k\mathrm{}𝑑s^{\mathrm{}}\left(\underset{¯}{x}\right),$$
(5.16)
where $`J^k`$ is the $`k`$-component of the bulk current; see eq. (2.4). As usual,
$$J^k(\underset{¯}{x},t)=𝒥^k(\underset{¯}{x},t)_A=\delta S_{\mathrm{eff}}(A)/\delta A_k(\underset{¯}{x},t).$$
(5.17)
By eqs. (2.17), (2.18), the R.S. of (5.17) is given by
$$\delta S_{\mathrm{eff}}(A)/\delta A_k(\underset{¯}{x},t)=\sigma _H\epsilon _k\mathrm{}E^{\mathrm{}}(\underset{¯}{x},t),$$
(5.18)
see also (2.4) (with $`\sigma _L=0`$). Thus
$$I_H^{\mathrm{bulk}}=\sigma _H_\gamma \underset{¯}{E}(\underset{¯}{x},t)𝑑\underset{¯}{s}(\underset{¯}{x}).$$
(5.19)
We have shown in eq. (5.9) that
$$I_H^{\mathrm{edge}}=\sigma _H\left(\mu _{\mathrm{}}\mu _r\right).$$
(5.20)
Thus, combining (5.15), (5.19) and (5.20), we conclude that
$$I_H=I_H^{\mathrm{edge}}+I_H^{\mathrm{bulk}}=\sigma _H\left(\mu _{\mathrm{}}\mu _r+_\gamma \underset{¯}{E}(\underset{¯}{x},t)𝑑\underset{¯}{s}\left(\underset{¯}{x}\right)\right).$$
(5.21)
But the expression in the parenthesis on the R.S. of (5.21) is nothing but the total voltage drop $`V`$ between the outer and the inner edge. Hence (5.21) implies that
$$I_H=\sigma _HV,$$
(5.22)
as desired.
Transport phenomena such as heat conduction through a quantum wire or a Hall sample (see Example 3 at the beginning of Sect. 2) can be studied along similar lines: In a physical system where modes of different chirality do not interact with each other (such as the modes at the inner and at the outer edge of the sample containing an incompressible Quantum Hall fluid) the left-moving and the right moving modes can be coupled to different reservoirs at different temperatures $`\beta _{\mathrm{}}^1`$ and $`\beta _r^1`$. This results in a non-zero heat current given by an expectation value of the component $`T^{01}`$ of the energy-momentum tensor of the conformal field theory describing the chiral modes in an equilibrium state where the left-movers are at temperature $`\beta _{\mathrm{}}^1`$ and the right-movers at temperature $`\beta _r^1\beta _{\mathrm{}}^1`$. (Such expectation values can be calculated from Virasoro characters.) These ideas lead to a conceptually clean understanding of the effects described in Example 3 at the beginning of Sect. 2.
## 6 A four-dimensional analogue of the Hall effect, and the generation of large, cosmic magnetic fields in the early universe
In this section, we further explore the four-dimensional analogue of the Hall effect described in Sect. 3. We shall apply our findings to exhibit effects that may play an important rôle in early-universe cosmology. Our results represent an elaboration upon those in .
We start our analysis by studying a system of massless Dirac fermions coupled to an external electromagnetic field in four-dimensional Minkowski space. Using results derived in Sects. 1 and 4, we derive equations analogous to eqs. (5.3)–(5.6) for the conductance of a quantum wire.
From Sect. 1 we recall the expression for the anomalous commutators between vector- and axial-vector — or chiral currents.
$$[\widehat{𝒥}_{\mathrm{}/r}^0(t,\underset{¯}{x}),\widehat{𝒥}_{\mathrm{}/r}^0(t,\underset{¯}{y})]=\pm i\frac{q^2}{4\pi ^2}\left(\underset{¯}{B}(\underset{¯}{x},t)\underset{¯}{}\right)\delta \left(\underset{¯}{x}\underset{¯}{y}\right),$$
(6.1)
where $`q`$ is the charge of the fermions — see eq. (1.62) — and
$$[\widehat{𝒥}_{\mathrm{}}^0(t,\underset{¯}{x}),\widehat{𝒥}_r^0(t,\underset{¯}{y})]=\mathrm{\hspace{0.33em}0}.$$
(6.2)
With (1.45) and (1.48), these equations yield
$$[\widehat{𝒥}_{\mathrm{}/r}^0(t,\underset{¯}{y}),𝒥^0(t,\underset{¯}{x})]=\pm i\frac{q^2}{8\pi ^2}\left(\underset{¯}{B}(\underset{¯}{y},t)\underset{¯}{}_{\underset{¯}{x}}\right)\delta \left(\underset{¯}{x}\underset{¯}{y}\right),$$
(6.3)
where $`𝒥^\mu `$ is the $`\mu `$-component of the conserved vector current. In Sect. 4, we have introduced the vector potential, $`\underset{¯}{\phi }`$, of $`𝒥^\mu `$:
$$𝒥^0(x)=\frac{q}{2\pi }\mathrm{div}\underset{¯}{\phi }(x),\underset{¯}{𝒥}(x)=\frac{q}{2\pi }\frac{}{t}\underset{¯}{\phi }(x).$$
(6.4)
Eqs. (6.3) and (6.4) imply that
$`[\widehat{𝒥}_{\mathrm{}/r}^0(\underset{¯}{y},t),\underset{¯}{\phi }(\underset{¯}{x},t)]`$ $`=`$ $`\pm i{\displaystyle \frac{q}{4\pi }}\underset{¯}{B}(\underset{¯}{y},t)\delta \left(\underset{¯}{x}\underset{¯}{y}\right)`$ (6.5)
$`\pm `$ $`\mathrm{curl}\underset{¯}{\mathrm{\Pi }}(\underset{¯}{x}\underset{¯}{y},t)`$
where $`\underset{¯}{\mathrm{\Pi }}`$ is some vector-valued distribution.
Next, we recall that the operators
$$N_{\mathrm{}/r}:=𝑑\underset{¯}{y}\widehat{𝒥}_{\mathrm{}/r}^0(\underset{¯}{y},t)$$
(6.6)
are conserved. They are interpreted as the electric charge operators for left-handed/right-handed fermionic modes. The chemical potentials conjugate to $`N_{\mathrm{}/r}`$ are denoted by $`\mu _{\mathrm{}/r}`$. Let us imagine that, at very early times in the evolution of our universe (or others), there was an asymmetry in the population of left-handed and right-handed fermionic modes, (as argued in for the example of electrons before the electroweak phase transition). Then
$$\mu _{\mathrm{}}\mu _r,$$
(6.7)
in the state of the universe at those very early times. Let us furthermore imagine that the state of the universe at those early times was, to a good approximation, a thermal equilibrium state at an inverse temperature $`\beta `$ ( $`\left(80TeV\right)^1`$, as argued in ) and with chemical potentials $`\mu _{\mathrm{}}`$ and $`\mu _r`$. (It may well be that this is an unrealistic assumption. — It will subsequently turn out that it is unimportant!)
Under these assumptions, we may apply the current sum rule (4.12) derived in Sect. 4. Combining eqs. (6.5), (6.6) and (4.12), and using that $`\underset{^3}{}𝑑\underset{¯}{y}\mathrm{curl}\underset{¯}{\mathrm{\Pi }}(\underset{¯}{x}\underset{¯}{y},t)=0`$, for all $`\underset{¯}{x},t`$, we find that
$`\underset{¯}{𝒥}(x)_{\beta ,\underset{¯}{\mu }}`$ $`=`$ $`{\displaystyle \frac{iq}{h}}\left\{\mu _{\mathrm{}}[N_{\mathrm{}},\underset{¯}{\phi }(x)]_{\beta ,\underset{¯}{\mu }}+\mu _r[N_r,\underset{¯}{\phi }(x)]_{\beta ,\underset{¯}{\mu }}\right\}`$ (6.8)
$`=`$ $`{\displaystyle \frac{q^2}{4\pi h}}\left(\mu _{\mathrm{}}\mu _r\right)\underset{¯}{B}(x),`$
as claimed in . This equation is the analogue of (5.6).
Treating the electromagnetic field as a classical, but dynamical field, its dynamics is governed by Maxwell’s equations,
$$\underset{¯}{}\underset{¯}{B}=\mathrm{\hspace{0.33em}0},\underset{¯}{}\underset{¯}{E}+_t\underset{¯}{B}=\mathrm{\hspace{0.33em}0},$$
and
$$\underset{¯}{}\underset{¯}{E}=𝒥^0_{\beta ,\underset{¯}{\mu }},\underset{¯}{}\underset{¯}{B}_t\underset{¯}{E}=\underset{¯}{𝒥}_{\beta ,\underset{¯}{\mu }}.$$
(6.9)
There is no reason to imagine that the charge density, $`𝒥^0_{\beta ,\underset{¯}{\mu }}`$, in the very early universe is different from zero. In the last equation of (6.9), the current on the R.S. is given by eq. (6.8). Actually, assuming that there are some dissipative processes evolving in the early universe, an equation for the current,
$$\underset{¯}{J}(x):=\underset{¯}{𝒥}(x)_{\beta ,\underset{¯}{\mu }},$$
more realistic than (6.8) may be
$$\underset{¯}{J}(x)=\sigma _L\underset{¯}{E}(x)+\sigma _TV\underset{¯}{B}(x),$$
(6.10)
where $`\sigma _L`$ is an Ohmic longitudinal conductivity, and
$$\sigma _T:=\frac{q^2}{4\pi h}$$
(6.11)
is the analogue of the “transverse” or Hall conductivity; furthermore,
$$V:=\mu _{\mathrm{}}\mu _r$$
(6.12)
is the analogue of the voltage drop considered in the Hall effect. The quantity $`\sigma _T`$ is “quantized”, just like the Hall conductivity: If there are $`N>1`$ species of charged, massless fermions, with electric charges $`q_1,\mathrm{},q_N,`$ then
$$\sigma _T=\frac{1}{4\pi h}\left(\underset{j=1}{\overset{N}{}}q_j^2\right),$$
(6.13)
which is the precise analogue of a formula for the quantization of the Hall conductivity derived in , and, for $`q_j=\pm e,j=1,\mathrm{},N,`$ of eq. (5.6).
Let us temporarily assume that $`\sigma _L=0`$, (i.e., we neglect dissipative processes). Then Maxwell’s equations, together with eq. (6.10) (for $`\sigma _L=0`$) and the assumption that the charge density vanishes, yield the following system of linear equations:
$`\underset{¯}{}\underset{¯}{B}=\mathrm{\hspace{0.33em}0},\underset{¯}{}\underset{¯}{E}+_t\underset{¯}{B}=\mathrm{\hspace{0.33em}0},`$
$`\underset{¯}{}\underset{¯}{E}=\mathrm{\hspace{0.33em}0},\underset{¯}{}\underset{¯}{B}_t\underset{¯}{E}=\sigma _TV\underset{¯}{B}.`$ (6.14)
Because all coefficients are constant, these equations can be solved by Fourier transformation, and it is enough to construct propagating wave solutions corresponding to an arbitrary, but fixed wave vector $`\underset{¯}{k}`$. The equations $`\underset{¯}{}\underset{¯}{B}=\underset{¯}{}\underset{¯}{E}=0`$ imply that
$$\underset{¯}{k}\underset{¯}{\overset{^}{B}}=\underset{¯}{k}\underset{¯}{\overset{^}{E}}=\mathrm{\hspace{0.33em}0},$$
(6.15)
i.e., that only the components of the Fourier transforms $`\underset{¯}{\overset{^}{B}}`$ and $`\underset{¯}{\overset{^}{E}}`$ of $`\underset{¯}{B}`$ and $`\underset{¯}{E}`$ (evaluated at the wave vector $`\underset{¯}{k}`$) perpendicular to $`\underset{¯}{k}`$ can be non-zero. Denoting the components of $`\underset{¯}{\overset{^}{B}}`$ and $`\underset{¯}{\overset{^}{E}}`$ perpendicular to $`\underset{¯}{k}`$ by $`\underset{¯}{\overset{^}{B}}^T,\underset{¯}{\overset{^}{E}}^T`$, respectively, the remaining equations in (6.14) yield
$$_t\left(\genfrac{}{}{0pt}{}{\underset{¯}{\overset{^}{E}}^T}{\underset{¯}{\overset{^}{B}}^T}\right)=K(\underset{¯}{k})\left(\genfrac{}{}{0pt}{}{\underset{¯}{\overset{^}{E}}^T}{\underset{¯}{\overset{^}{B}}^T}\right),$$
(6.16)
where (in an orthonormal basis chosen in the plane perpendicular to $`\underset{¯}{k}`$) the matrix $`K(\underset{¯}{k})`$ is given by
$$K(\underset{¯}{k})=\left(\begin{array}{cccc}0& 0& \sigma _TV& ik\\ 0& 0& ik& \sigma _TV\\ 0& ik& 0& 0\\ ik& 0& 0& 0\end{array}\right)$$
(6.17)
with $`k=|\underset{¯}{k}|`$. The circular frequency of a propagating wave solution of (6.14) with wave vector $`\underset{¯}{k}`$ is given by $`\omega (\underset{¯}{k})`$, where $`i\omega (\underset{¯}{k})`$ is an eigenvalue of $`K(\underset{¯}{k})`$. By (6.17),
$$\omega (\underset{¯}{k})^2=k^2\pm k\sigma _TV,$$
(6.18)
as one readily checks. Thus, if
$$|\underset{¯}{k}|=k<\sigma _TV$$
(6.19)
there are two purely imaginary frequencies, and eqs. (6.14) have solutions $`(\underset{¯}{B}(\underset{¯}{x},t),\underset{¯}{E}(\underset{¯}{x},t))`$ growing exponentially fast in time and with the property that
$$\underset{¯}{B}(\underset{¯}{x},t)\underset{¯}{E}(\underset{¯}{x},t)\mathrm{\hspace{0.33em}0}.$$
(6.20)
It is almost as easy to solve Maxwell’s equations (6.9), with $`\underset{¯}{J}`$ given by (6.10), for $`\sigma _L0`$. For wave vectors $`\underset{¯}{k}`$ satisfying
$$\sigma _L^2<|\underset{¯}{k}|\sigma _TV<\left(\sigma _TV\right)^2,$$
(6.21)
one again finds exponentially growing electromagnetic fields; (perturbation theory). Dissipative processes will subsequently damp out electric fields.
In , calculations similar to those just presented are used to argue that, in the very early universe, large, cosmic electromagnetic fields may have been generated as a consequence of an asymmetric population of left-handed and right-handed electron modes $`(q=e)`$. However, these arguments rest on rather shaky hypotheses; (the state of the early universe is assumed to be a thermal equilibrium state, and the charges $`N_{\mathrm{}}`$ and $`N_r`$, see eq. (6.6), are assumed to be approximately conserved). We propose to reconsider these arguments in the light of the analogy between the (2+1)-dimensional (bulk) description of the Hall effect and the (4+1)-dimensional description of chiral fermions discussed at the beginning of Sect. 3, eqs (3.3) through (3.10). What we have described, so far, in this section are calculations analogous to those reported in eqs. (5.6), (5.8) and (5.9). Next, we generalize our analysis in a way analogous to that followed in eqs. (5.15) through (5.22), starting from the effective action given in (3.10); (see also (3.20)).
We integrate out all degrees of freedom (quarks, gluons, leptons, the weak gauge fields — $`W,Z`$ — etc.), except for the electromagnetic and the axion field. We have seen, at the beginning of Sect. 3, eqs. (3.4), (3.10), that the axion could be viewed as the four-component of a five-dimensional electromagnetic vector potential, $`\widehat{A}`$, which does not depend on the coordinate, $`x^4`$, in the direction perpendicular to the four-dimensional branes on which we live; see (3.6). We could pursue a five- (or higher-) dimensional approach to early-universe cosmology (as presently popular), — but let’s not! We propose to view the axion as the “model-independent (invisible)” axion first described in . It has a geometrical origin (in superstring theory). It couples to all gauge fields present in the system through a term
$$\frac{i\mathrm{}}{32\pi ^2}\phi \left(F_WF_W\right),$$
(6.22)
where $`F_W`$ is the field strength of a gauge field $`W`$ appearing in our theoretical description, and to the curvature tensor $`R`$; see (3.16). All gauge fields, except for the electromagnetic vector potential $`A`$, shall be integrated out. The (Euclidian-region-) functional integrals have the form
$$d\mu (W)\mathrm{exp}[\frac{i\mathrm{}}{32\pi ^2}\phi (F_WF_W)]=:e^{U(\phi )}.$$
(6.23)
Since $`\frac{i}{32\pi ^2}\left(F_WF_W\right)`$ is the index density, the integrand in $`U(\phi )`$ can be shown to be periodic in $`\phi `$, for $`\phi `$ independent of $`x`$, with period $`\frac{2\pi }{\mathrm{}}`$. It is known that (somewhat loosely speaking) $`d\mu `$ is a positive measure and that it is invariant under space reflection, which changes the sign of $`F_WF_W`$. It follows that $`\mathrm{exp}\left(U(\phi )\right)`$ is real and has its maxima at $`\phi =\frac{2\pi }{\mathrm{}}n,n=0,\pm 1,\pm 2,\mathrm{}.`$ (See e.g. for more details.)
A transition amplitude from a configuration $`(A_{\mathrm{in}},\phi _{\mathrm{in}})`$ of the electromagnetic — and the axion field at a very early time, $`t_1`$, to a configuration $`(A_{\mathrm{out}},\phi _{\mathrm{out}})`$ at a much later time, $`t_2`$, can be computed from the Feynman path integral
$$𝒟A𝒟\phi e^{iS_{\mathrm{eff}}(A,\phi )/\mathrm{}},$$
(6.24)
with boundary conditions $`(A(t_1),\phi (t_1))=(A_{\mathrm{in}},\phi _{\mathrm{in}})`$ and $`(A(t_2),\phi (t_2))=(A_{\mathrm{out}},\phi _{\mathrm{out}})`$. In (6.24), $`S_{\mathrm{eff}}(A,\phi )`$ denotes the total effective action over Minkowski space. It is obtained from $`S_{\mathrm{eff}}^E(A,\phi )`$, the effective action in the Euclidian region, by undoing the Wick rotation described in eq. (1.28). By eqs. (3.10) or (3.20) and (6.23), $`S_{\mathrm{eff}}(A,\phi )`$ has the general form
$`S_{\mathrm{eff}}(A,\phi )`$ $`=`$ $`{\displaystyle \frac{1}{4e^2}}{\displaystyle d^4x\left\{F_{\mu \nu }(x)F^{\mu \nu }(x)+\mathrm{\hspace{0.17em}2}\left(_\mu \phi \right)(x)\left(^\mu \phi \right)(x)\right\}}`$ (6.25)
$`+`$ $`{\displaystyle \frac{\mathrm{}}{32\pi ^2}}{\displaystyle \phi (x)\left(FF\right)(x)}U(\phi )+W(A),`$
where $`W(A)`$ is of higher than second order in $`A`$ and arises from integrating out all charged fields in the theory<sup>1</sup><sup>1</sup>1$`W`$ depends on the boundary conditions, at times $`t_1,t_2`$, imposed on the fields that have been integrated out.; furthermore, $`e^2`$ is the effective (one-loop renormalized) feinstructure constant. It is not necessary, in this approach, to assume that all the fermions in the theory be massless. They can acquire masses through the Higgs–Kibble mechanism. (The arguments of complex chiral Higgs fields then contain a term proportional to the axion field $`\phi `$ which, however, can be absorbed in a change of variables.) Furthermore, calculating transition amplitudes with the help of Feynman path integrals does not presuppose that the system is in or close to thermal equilibrium.
We now insert expression (6.25) into the functional integral (6.24) and try to evaluate the latter by using a semi-classical expansion based on the stationary-phase method. The equations for the saddle point are
$$\delta S_{\mathrm{eff}}(A,\phi )/\delta A_\mu (x)=\mathrm{\hspace{0.33em}0},\delta S_{\mathrm{eff}}(A,\phi )/\delta \phi (x)=\mathrm{\hspace{0.33em}0}.$$
(6.26)
To simplify matters, we consider solutions of these equations describing fairly small electromagnetic fields and an axion field that varies only slowly in space-time. Then we can neglect the term $`W(A)`$ in (6.25) and we may omit all contributions to $`U(\phi )`$ involving derivatives, $`_\mu \phi `$, of the axion field $`\phi `$. The saddle point equations (6.26) then yield the following coupled Maxwell–Dirac-axion equations:
$`_\mu F^{\mu \nu }`$ $`=`$ $`{\displaystyle \frac{\mathrm{}e^2}{8\pi ^2}}_\mu \left(\phi \stackrel{~}{F}^{\mu \nu }\right),`$
$`\mathrm{}\phi `$ $`=`$ $`{\displaystyle \frac{\mathrm{}e^2}{32\pi ^2}}\left(FF\right)U^{}(\phi ),`$ (6.27)
(and we have set $`c=1`$ and $`\mathrm{}=1`$). Let $`J_M^\mu `$ denote the magnetic current that could be present if there were magnetic monopoles moving through the early universe. Then the full set of Maxwell–Dirac-axion equations reads
$`_\mu \stackrel{~}{F}^{\mu \nu }`$ $`=`$ $`J_M^\nu ,_\mu F^{\mu \nu }={\displaystyle \frac{\mathrm{}e^2}{8\pi ^2}}\left\{\left(_\mu \phi \right)\stackrel{~}{F}^{\mu \nu }+\phi J_M^\nu \right\},`$
$`\mathrm{}\phi `$ $`=`$ $`{\displaystyle \frac{\mathrm{}e^2}{32\pi ^2}}\left(FF\right)U^{}(\phi ).`$ (6.28)
The first equation in (6.28) replaces the homogeneous Maxwell equations,
$`\left(_\mu \stackrel{~}{F}^{\mu \nu }=0,\mathrm{for}J_M^\nu =0\right).`$ In vector notation, the system of equations (6.28) reads
$`\underset{¯}{}\underset{¯}{B}`$ $`=`$ $`J_M^0,\underset{¯}{}\underset{¯}{E}+\underset{¯}{\overset{\dot{}}{B}}=\underset{¯}{J}_M,`$
$`\underset{¯}{}\underset{¯}{E}`$ $`=`$ $`{\displaystyle \frac{\mathrm{}e^2}{8\pi ^2}}\left\{\left(\underset{¯}{}\phi \right)\underset{¯}{B}+\phi J_M^0\right\},`$
$`\underset{¯}{B}`$ $``$ $`\underset{¯}{\overset{\dot{}}{E}}={\displaystyle \frac{\mathrm{}e^2}{8\pi ^2}}\left\{\dot{\phi }\underset{¯}{B}+\underset{¯}{}\phi \underset{¯}{E}+\phi \underset{¯}{J}_M\right\},`$
$`\mathrm{}\phi `$ $`=`$ $`{\displaystyle \frac{\mathrm{}e^2}{8\pi ^2}}\underset{¯}{E}\underset{¯}{B}U^{}(\phi ).`$ (6.29)
In order to gain some insight into properties of solutions of these highly non-linear equations, we study their linearization around various special solutions. Already this part of the analysis, let alone a study of the full, non-linear equations, is quite lengthy; see for a beginning. Here we just sketch results in a few interesting special situations.
We shall first assume that $`J_M^\mu 0`$, i.e., that there aren’t any magnetic monopoles around.
(i) We set $`U(\phi )=0`$ and consider the following special solution of eqs. (6.29).
$`\underset{¯}{E}=\underset{¯}{B}\mathrm{\hspace{0.33em}0},`$
$`\phi (\underset{¯}{x},t)={\displaystyle \frac{V}{\mathrm{}}}t,`$ (6.30)
where $`V`$ is a constant. Linearizing (6.29) around (6.30), we obtain the equations
$`\underset{¯}{}\underset{¯}{B}`$ $`=`$ $`0,\underset{¯}{}\underset{¯}{E}+\underset{¯}{\overset{\dot{}}{B}}=\mathrm{\hspace{0.33em}0},`$
$`\underset{¯}{}\underset{¯}{E}`$ $`=`$ $`0,\underset{¯}{}\underset{¯}{B}\underset{¯}{\overset{\dot{}}{E}}={\displaystyle \frac{e^2}{8\pi ^2}}V\underset{¯}{B},`$
$`\mathrm{}\phi `$ $`=`$ $`0.`$ (6.31)
With the exception of the wave equation for the axion field $`\phi `$, these equations are identical to eqs. (6.14), with $`\sigma _T=\frac{e^2}{8\pi ^2}.`$ Had we not set $`\mathrm{}=1`$, the equation for $`\sigma _T`$ would read
$$\sigma _T=\frac{e^2}{4\pi h},$$
which is precisely eq. (6.11), with $`q=e`$! Recall that, in the analysis presented at the beginning of this section,
$$V=\mu _{\mathrm{}}\mu _r.$$
This equation and (6.30) tell us that, apparently, the field $`\mathrm{}\dot{\phi }`$ has the interpretation of the difference of chemical potentials of left- and right-handed fermions! This interpretation magically fits with the five-dimensional interpretation of the axion field $`\phi `$ as the four-component, $`\widehat{A}_4`$, of an electromagnetic vector potential $`\widehat{A}`$ defined over a slab of height $`\mathrm{}`$ in five-dimensional Minkowski space; see eqs. (3.4), (3.6) and (3.10). Then
$$\dot{\phi }=\dot{\widehat{A}}_4=E_4$$
is the four-component of the electric field. Integrating $`E`$ along an oriented curve, $`\gamma `$, joining a point on the lower face of the slab to a point on the upper face, at fixed time, we obtain
$$\underset{\gamma }{}\underset{K=1}{\overset{4}{}}E_Kds^K=\underset{0}{\overset{\mathrm{}}{}}𝑑x^4E_4(\xi )=\underset{0}{\overset{\mathrm{}}{}}𝑑x^4\dot{\phi }(\xi ),$$
(6.32)
where $`\xi =(x,x^4)=(t,\underset{¯}{x},x^4)`$, and we have assumed in the first equality that $`E`$ does not depend on $`x^4`$ (see assumption (3.6)) and $`E_4`$ does not depend on $`\underset{¯}{x}`$. Since, for solution (6.30),
$$E_4(\xi )=\dot{\phi }(\xi )=\frac{V}{\mathrm{}},$$
eq. (6.32) yields
$$\underset{\gamma }{}\underset{K=1}{\overset{4}{}}E_Kds^K=V.$$
(6.33)
This shows that, in the five-dimensional interpretation of the axion, $`V`$ is the “voltage drop” between the two four-dimensional branes corresponding to the lower and upper face of the five-dimensional slab. This observation makes the analogy between the effects studied here and the Hall effect yet a little more precise.
Solutions of eqs. (6.31) have been studied earlier in this section; see (6.16) through (6.20). They have unstable modes growing exponentially in time, with $`\underset{¯}{B}(\underset{¯}{x},t)\underset{¯}{E}(\underset{¯}{x},t)0`$ .
(ii) Now $`U(\phi )0`$; $`U^{}(\phi )(x):=\delta U(\phi )/\delta \phi (x)`$ is a periodic function with minima at $`\frac{2\pi }{\mathrm{}}n`$, $`n=0,\pm 1,\pm 2,\mathrm{}`$ . We linearize equations (6.29) around the solution $`\underset{¯}{E}=\underset{¯}{B}0`$, $`\phi =\phi _c(t)`$ , where $`\phi _c(t)`$ solves the equation
$$\ddot{\phi }(t)=U^{}\left(\phi \left(t\right)\right).$$
(6.34)
This is the equation of motion of a planar pendulum in a force field with potential $`U`$. We have learnt in our courses on elementary mechanics how to solve (6.34), using energy conservation. For “small energy”, a solution, $`\phi _c(t)`$, of (6.34) is a periodic function of $`t`$; for “large energy”, $`\phi _c(t)`$ grows linearly in $`t`$, with periodic modulations superimposed; and $`\dot{\phi }_c(t)`$ is periodic in $`t`$.
Eqs. (6.29), with $`J_M^\mu 0`$, linearized around $`\underset{¯}{E}=\underset{¯}{B}=0,`$ $`\phi =\phi _c(t),`$ yield the equations
$`\underset{¯}{}\underset{¯}{B}=\mathrm{\hspace{0.33em}0},\underset{¯}{}\underset{¯}{E}+\underset{¯}{\overset{\dot{}}{B}}=\mathrm{\hspace{0.33em}0},`$
$`\underset{¯}{}\underset{¯}{E}=\mathrm{\hspace{0.33em}0},\underset{¯}{}\underset{¯}{B}\underset{¯}{\overset{\dot{}}{E}}={\displaystyle \frac{\mathrm{}e^2}{8\pi ^2}}\dot{\phi }_c\underset{¯}{B},`$ (6.35)
which can be solved by Fourier transformation in the space variables. The equations for the components, $`\underset{¯}{\overset{^}{B}}^T`$ and $`\underset{¯}{\overset{^}{E}}^T`$, of the Fourier components of $`\underset{¯}{B}`$ and $`\underset{¯}{E}`$ perpendicular to the wave vector $`\underset{¯}{k}`$ are two Mathieu equations of the form
$$\left(\genfrac{}{}{0pt}{}{\dot{\xi }}{\dot{\eta }}\right)=\left(\begin{array}{cc}0\hfill & 1\hfill \\ h_k(t)\hfill & 0\hfill \end{array}\right)\left(\genfrac{}{}{0pt}{}{\xi }{\eta }\right),$$
where $`k=|\underset{¯}{k}|`$ and $`h_k(t)`$ depends on $`k`$ and is linear in $`\dot{\phi }_c(t)`$; see . These equations yield
$$\ddot{\xi }(t)=h_k(t)\xi (t).$$
(6.36)
In solving this equation one encounters the phenomenon of the parametric resonance, i.e., for $`k`$ in a family of intervals, eq. (6.36) has a solution growing exponentially in time. Hence the electromagnetic field has unstable modes growing exponentially in time and with $`\underset{¯}{B}\underset{¯}{E}0`$.
The parametric resonance has appeared in cosmology in other contexts. In our analysis it plays an entirely natural and essentially model-independent rôle and may help to explain where large, cosmic (electro) magnetic fields might come from.
Of course, eqs. (6.29) are Lagrangian equations of motion. They are derived from the action functional (6.25), (with $`W=0`$ and $`U`$ independent of derivatives of $`\phi `$). The Lagrangian density does not depend on time explicitly. Therefore, there is a conserved energy functional, $`(A,\phi )`$. The special solutions considered in (6.30) and (6.34) have infinite (axionic) energy. The instabilities in the time evolution of the electromagnetic field are due to a reshuffling of energy from axionic to electromagnetic degrees of freedom.
Clearly, it would be interesting to construct finite-energy solutions of eqs. (6.29), with an initial axion field depending not only on time but also on space. Of particular significance is situation (ii), with $`U0`$. Interpreting $`\mathrm{}\dot{\phi }`$ as a difference of chemical potentials for left- and right-handed fermions, we are thus considering states of the universe with spatially varying, time-dependent chemical potentials triggering an asymmetric population of left-handed and right-handed fermionic modes. This asymmetry gradually disappears, due to chirality-changing processes, and the field energy stored in axionic degrees of freedom is reshuffled into certain electromagnetic field modes triggering the growth of cosmic electromagnetic fields. Large electric fields rapidly die out because of dissipative processes; (the energy loss from the electric field into matter degrees of freedom is described by $`\underset{¯}{E}\underset{¯}{J}\sigma _L|\underset{¯}{E}|^2+\mathrm{}.`$) But large magnetic fields may survive for a comparatively long time.
Describing these phenomena within the approximation of linearizing eqs. (6.29) (possibly supplemented by a dissipative Ohmic term) around special solutions, including space-dependent ones, of infinite or finite energy, is feasible; . But our understanding of the effects of the non-linearities in eqs. (6.29) remains, not surprisingly, very rudimentary.
Some speculations on the rôle played by magnetic monopoles in the effects described here are contained in the last section; see also .
## 7 Conclusions and outlook
In this review we have shown how the chiral, abelian anomaly helps to explain important features of the (quantum) Hall effect, such as the existence of edge currents and aspects of the quantization of the Hall conductivity, and of its four-dimensional cousin, which may play a significant rôle in explaining the origin of large, cosmic magnetic fields. Our analysis is essentially model-independent, a fact that makes it quite trustworthy. How significant the four-dimensional variant of the Hall effect is in early-universe cosmology remains to be understood in more detail. This will require a better understanding of orders of magnitude of various physical quantities and of the properties of solutions of the non-linear Maxwell–Dirac-axion equations (6.29). A beginning has been made in . — There is no doubt that the following equations
$$J_{\mathrm{bulk}}=\sigma _TF,\delta J_{\mathrm{edge}}=\sigma _TE,$$
(7.1)
with $`\sigma _T=\sigma _H`$, for bulk- and edge-currents of an incompressible Hall fluid (see eqs. (2.10) and (2.14)), and
$$J^\nu =\sigma _T\mathrm{}\left\{\left(_\mu \phi \right)\stackrel{~}{F}^{\mu \nu }+\phi J_M^\nu \right\},$$
(7.2)
where $`\sigma _T=\frac{1}{4\pi h}\left(_{j=1}^Nq_j^2\right),`$ with $`N`$ the number of species of charged fermions with electric charges $`q_1,\mathrm{},q_N`$, (see eqs. (6.13) and (6.28)) are significant laws of nature connected with the chiral anomaly.
For the future, it would be important to gain a better understanding of the contents of equations (6.29), (possibly corrected by dissipative terms and/or ones coming from $`\delta W(A)/\delta A_\mu (x)`$, which have been neglected), including the rôle played by magnetic monopoles and dyons $`\left(J_M^\mu 0\right)`$. (Eqs. (6.29) and their fully quantized counterparts appear to offer some clue for understanding (axion-driven) monopole–anti-monopole annihilation, triggering the growth of certain modes of the electromagnetic field.) Some understanding of these issues has been gained in ; but much work remains to be done. We have also studied the influence of gravitational fields on the processes described in Sect. 6 (in analogy to the “geometric” (or gravitational) Hall effect in 2+1 dimensions described in the third paper quoted under and to the phenomenon of “quantized” heat currents in quantum wires mentioned in Sects. 3 and 5). But there is no room here to describe our results in detail. Our findings will have to be combined with cosmic evolution equations.
————————
In this review, we have only quoted literature that we used in carrying out the calculations described here. Many further references may be found in .
## Acknowledgments
The results described in Sects. 2, 4 and 5 have been obtained in collaboration (of J.F.) with A. Alekseev and V. Cheianov , in continuation of earlier work with T. Kerler, U. Studer and E. Thiran. We thank these colleagues, Chr. Schweigert and Ph. Werner for many useful discussions. We are grateful to R. Durrer, E. Seiler and D. Wyler for drawing our attention to some useful earlier work in the literature and for encouragement. |
warning/0002/nlin0002046.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Perceptual alternation phenomena of ambiguous figures have been studied for a long time. Figure-ground, perspective (depth) and semantic ambiguities are well known (As an overview, for example, see and ). Actually, when we view the Necker cube which is a classic example of perspective alternation, a part of the figure is perceived either as front or back of a cube and our perception switches between the two different interpretations as shown in Fig.1. In this circumstance the external stimulus is kept constant, but perception undergoes involuntary and random-like change. The measurements have been quantified in psychophysical experiments and it becomes evident that the times between such changes are approximately Gamma distributed .
Theoretical model approaches to explaining the facts have been made mainly from three situations based on the synergetics , the BSB (brain-state-in-a-box) neural network model , and the PDP (parallel distributed processing) schema model . Common to these approaches is that top-down designs are applied so that the model can be manipulable by a few parameters and upon this basis fluctuating sources are brought in. The major interests seem to be not in the relation between the whole function and its element (neuron), but in the model building at the phenomenological level.
Until now diverse types of chaos have been confirmed at several hierarchical levels in the real neural systems from single cells to cortical networks (e.g. ionic channels, spike trains from cells, EEG) . This suggests that artificial neural networks based on the McCulloch-Pitts neuron model should be re-examined and re-developed. Chaos may play an essential role in the extended frame of the Hopfield neural network beyond the only equilibrium point attractors. To make this point clear, following the model of chaotic neural network , the dynamic learning and retrieving features of the associative memory have been studied . In this paper, we present a perception model of ambiguous patterns based on the chaotic neural network from the viewpoint of bottom-up approach , aiming at the functioning of chaos in dynamic perceptual processes.
## 2 Model and Method
The chaotic neural network (CNN) composed of N chaotic neurons is described as
$`X_i(t+1)`$ $`=`$ $`f(\eta _i(t+1)+\zeta _i(t+1)),`$ (1)
$`\eta _i(t+1)`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{N}{}}}w_{ij}{\displaystyle \underset{d=0}{\overset{t}{}}}k_f^dX_j(td),`$ (2)
$`\zeta _i(t+1)`$ $`=`$ $`\alpha {\displaystyle \underset{d=0}{\overset{t}{}}}k_r^dX_i(td)\theta _i,`$ (3)
where $`X_i`$ : output of neuron $`i(1X_i1),w_{ij}`$ : synaptic weight from neuron $`j`$ to neuron $`i,\theta _i`$ : threshold of neuron $`i,k_f(k_r`$) : decay factor for the feedback(refractoriness) $`(0k_f,k_r<1),\alpha `$ : refractory scaling parameter, $`f`$ : output function defined by $`f(y)=tanh(y/2\epsilon )`$ with the steepness parameter $`\epsilon `$. Owing to the exponentially decaying form of the past influence, Eqs.(2) and (3) can be reduced to
$`\eta _i(t+1)`$ $`=`$ $`k_f\eta _i(t)+{\displaystyle \underset{j=1}{\overset{N}{}}}w_{ij}X_j(t),`$ (4)
$`\zeta _i(t+1)`$ $`=`$ $`k_r\zeta _i(t)\alpha X_i(t)+a,`$ (5)
where $`a`$ is temporally constant $`a\theta _i(1k_r)`$. All neurons are updated in parallel, that is, synchronously. The network corresponds to the conventional discrete-time Hopfield network :
$`X_i(t+1)=f({\displaystyle \underset{j=1}{\overset{N}{}}}w_{ij}X_j(t)\theta _i)`$ (6)
when $`\alpha =k_f=k_r=0`$ (Hopfield network point (HNP)). The asymptotical stability and chaos in discrete-time neural networks are theoretically investigated in Refs. .
Under external stimuli, Eq.(1) is influenced as
$`X_i(t+1)=f\left(\eta _i(t+1)+\zeta _i(t+1)+\sigma _i\right),`$ (7)
where $`\{\sigma _i\}`$ is the effective term by external stimuli. This is a simple and unartificial incorporation of stimuli as the changes of neural active potentials.
The two competitive interpretations are embedded in the network as minima of the energy map :
$`E={\displaystyle \frac{1}{2}}{\displaystyle \underset{ij}{}}w_{ij}X_iX_j`$ (8)
at HNP. This is done by using a iterative perception learning rule for $`p(<N)`$ patterns $`\{\xi _i^\mu \}(\xi _1^\mu ,\mathrm{},\xi _N^\mu ),(\mu =1,\mathrm{},p;\xi _i^\mu =+1or1)`$ in the form :
$`w_{ij}^{new}=w_{ij}^{old}+{\displaystyle \underset{\mu }{}}\delta w_{ij}^\mu `$ (9)
with
$`\delta w_{ij}^\mu ={\displaystyle \frac{1}{N}}\theta (1\gamma _i^\mu )\xi _i^\mu \xi _j^\mu ,`$ (10)
where $`\gamma _i^\mu \xi _i^\mu _{j=1}^Nw_{ij}\xi _j^\mu `$ and $`\theta (h)`$ is the unit step function. The learning mode is separated from the performance mode by Eq.(7).
The conceptual picture of our model is shown in Fig.2. Under the external stimulus $`\{\sigma _i\}`$, chaotic activities arise on the neural network and cause the transitions between stable states of HNP. This situation corresponds to the dynamic multistable perception. Note that $`\theta (1\gamma _i^\mu )`$ turns off learning for overlearned patterns. This will be empirically shown in the next section to have dynamics characteristic of a chaotic dynamical system.
## 3 Simulations and Results
To carry out computational experiments, we use the $`12\times 13(N=156)`$ non-orthogonal 10 random patterns $`\{\xi _i^\nu \}(\nu =1,\mathrm{},10,;i=1,\mathrm{},N)`$ as a set of ambiguous figure stimuli: $`\{\sigma _i\}=s\{\xi _i^\nu \}`$. $`s`$ is the strength factor of stimulation. For each of them, two interpretation patterns $`\{\xi _i^{\nu 1}\}`$ and $`\{\xi _i^{\nu 2}\}`$ are prepared by changing 15 white ($`\xi _i=1`$) pixels to black ($`\xi _i=+1`$) ones which do not overlap between $`\nu 1`$ and $`\nu 2`$ as shown by shaded and dotted in Fig.3, and are memorized following the above learning rule $`(p=20)`$.
Figure 4 shows a time series evolution of CNN $`(k_f=0.5,k_r=0.8,\alpha =0.34,a=0,\epsilon =0.015)`$ under the stimulus $`\{\sigma _i\}=0.7\{\xi _i^1\}`$. Here,
$`m^{11}(t)={\displaystyle \frac{1}{N}}{\displaystyle \underset{i=1}{\overset{N}{}}}\xi _i^{11}X_i(t)`$ (11)
and is called the overlap of the network state $`\{X_i\}`$ and the interpretation pattern $`\{\xi _i^{11}\}`$. A switching phenomenon between $`\{\xi _i^{11}\}`$ ($`m^{11}=1.0`$) and $`\{\xi _i^{12}\}`$ ($`m^{11}=0.62`$) can be observed. Bursts of switching are interspersed with prolonged periods during which $`\{X_i\}`$ trembles near $`\{\xi _i^{11}\}`$ or $`\{\xi _i^{12}\}`$. Evaluating the maximum Lyapunov exponent to be positive ($`\lambda _1=0.26`$), we find that the network is dynamically in chaos. In the cases $`\lambda _1<0`$, such switching phenomena do not arise.
From the $`2\times 10^5`$ iteration data (until $`t=2\times 10^5`$) of Fig.4, we get 1257 events staying near one of the two interpretations, $`\{\xi _i^{12}\}`$. As can be seen from Fig.5 magnified for t-axis, they have various persistent durations $`T(1)T(1257)`$ which seem to have a random time course by the return map in $`(T(n),T(n+1))`$ shown in Fig.6. From the evaluation of the autocorrelation function for $`T(n)`$, $`C(k)=<T(n+k)T(n)><T(n+k)><T(n)>`$ (here, $`<>`$ means an average over time), we get $`0.06<C(k)/C(0)<0.06`$ against $`k=1100`$. This suggests successive durations $`T(n)`$ are independent. The frequency of occurrence of $`T`$ is plotted for 1257 events in Fig.7. The distribution is well fitted by Gamma distribution
$`G(\stackrel{~}{T})={\displaystyle \frac{b^n\stackrel{~}{T}^{n1}e^{b\stackrel{~}{T}}}{\mathrm{\Gamma }(n)}}`$ (12)
with $`b=0.918,n=4.68(\chi ^2=0.0033,r=0.98)`$, where $`\mathrm{\Gamma }(n)`$ is the Euler-Gamma function. $`\stackrel{~}{T}`$ is the normalized duration $`T/15`$ and here 15 step interval is applied to determine the relative frequencies.
The results are in good agreement with the characteristics of psychophysical experiments . Similar results to the above example are obtained in the appropriate parameter regions where the network may induce chaotic activities under external stimuli. It is found that aperiodic spontaneous switching does not necessitate some stochastic description as in the synergetic model .
These results can not be easily explained through the use of a standard Hopfield network with a stochastic fluctuation forcing function. We look into the case that the stochastic fluctuation $`\{F_i\}`$ is attached to Eq.(6) of HNP together with the external stimulus $`\{\sigma _i\}`$ :
$`X_i(t+1)=f({\displaystyle \underset{j=1}{\overset{N}{}}}w_{ij}X_j(t)+\sigma _i+F_i(t)),`$ (13)
where
$`\{\begin{array}{cc}<F_i(t)>=0\hfill & \\ <F_i(t)F_j(t^{})>=D^2\delta _{tt^{}}\delta _{ij}.\hfill & \end{array}`$ (16)
Using this equation in the same framework, we examined many cases with different values of the noise strength $`D`$, but couldn’t find successive alternation phenomena. Figure 8 is a typical result in $`s=0.5,D=0.65`$ and is far from the realization of sudden perceptual changes. In such a simple scheme, the noise does not drive a quick motion of $`\{X_i\}`$ for the energy barrier. In Fig.9 the frequency of occurrence $`T`$ is plotted for 1387 events obtained until $`t=2\times 10^5`$. The distribution becomes the exponential-like, not Gamma distribution.
## 4 Conclusion
We have shown that the neural chaos leads to perceptual alternations as responses to ambiguous stimuli in the chaotic neural network. Its emergence is based on the simple process in a realistic bottom-up framework. In the same stage, similar results can not be obtained by the stochastic activity. In order to compare the simulation results with experimental ones in a concrete form, there remain problems which are to analyze the relationship between the iteration (step) time and the visual real time, and to study the switching dependence on the cube’s orientation or size.
Finally, our demonstration suggests functional usefulness of the chaotic activity in perceptual systems even at higher cognitive levels. The perceptual alternation appears to be an inherent feature built in the chaotic neuron assembly. It may be interesting to study the brain with the experimental technique (e.g., fMRI) under the circumstance where the perceptual alternation is running. |
warning/0002/quant-ph0002065.html | ar5iv | text | # Unitary relation between a harmonic oscillator of time-dependent frequency and a simple harmonic oscillator with and without an inverse-square potential
## Abstract
The unitary operator which transforms a harmonic oscillator system of time-dependent frequency into that of a simple harmonic oscillator of different time-scale is found, with and without an inverse-square potential. It is shown that for both cases, this operator can be used in finding complete sets of wave functions of a generalized harmonic oscillator system from the well-known sets of the simple harmonic oscillator. Exact invariants of the time-dependent systems can also be obtained from the constant Hamiltonians of unit mass and frequency by making use of this unitary transformation. The geometric phases for the wave functions of a generalized harmonic oscillator with an inverse-square potential are given.
It is certainly of importance to find a complete set of wave functions for a system of the time-dependent Hamiltonian \[1-17\]. It has long been known that a harmonic oscillator of time-dependent frequency with or without an inverse-square potential is the system of practical applications (see e.g. Ref. ), where the wave functions are described in terms of solutions of classical equation of motion of the oscillator without the inverse-square potential . In Ref. it has been shown that, for a generalized harmonic oscillator system, the kernel of the system is determined by the classical action. This is one of the basic reasons of the fact that the wave functions are described by the classical solutions. On the other hand, it has long been noticed that there exist classical canonical transformations which relate the (driven) harmonic oscillators of different parameters (see e.g. Refs. ). Recently, in Ref. , it has been shown that a driven harmonic oscillator of time-dependent frequency is related, through canonical transformations, to the simple harmonic oscillator of unit mass and unit frequency but with a different time-scale . This fact has been used to find the wave functions of a driven system which exactly agree with the known results .
In this Rapid Communication, we will show that for both oscillators with and without an inverse-square potential, there is a unitary operator which transforms the harmonic oscillator systems of time-dependent frequency into those of the unit-mass and unit-frequency oscillators with different time-scales. This unitary operator can be used to find complete sets of wave functions of the systems with time-dependent parameters from the well-known sets of wave functions of the simple harmonic oscillator with or without an inverse-square potential. It has been known that there exist exact invariants in the systems of time-dependent parameters which have long been used to find the wave functions . As might have been implied by the classical treatments through canonical transformations, it will also be shown that, the exact quantum-mechanical invariants in oscillator systems of time-dependent parameters can be obtained from the constant Hamiltonians of unit mass and frequency (which, certainly, are invariants in their systems respectively), through the unitary transformation given here and those in Refs. .
The unit-mass harmonic oscillator of time-dependent frequency $`w_0(t)`$ is described by the Hamiltonian
$$H_0(x,p,t)=\frac{p^2}{2}+\frac{w_0^2(t)}{2}x^2,$$
(1)
with the classical equation of motion
$$\ddot{x}+w_0^2(t)x=0.$$
(2)
If we denote the two linearly independent solutions of Eq. (2) as $`u_0(t)`$ and $`v_0(t)`$, the $`\rho _0(t)`$ defined by $`\rho _0(t)=\sqrt{u_0^2+v_0^2}`$ should satisfy
$$\frac{d^2}{dt^2}\rho _0+w_0^2(t)\rho _0\frac{\mathrm{\Omega }_0^2}{\rho _0^3}=0,$$
(3)
with a time-constant $`\mathrm{\Omega }_0`$ ($`\dot{v}_0u_0\dot{u}_0v_0`$). Without losing generality we assume that $`\mathrm{\Omega }_0`$ is positive. The wave functions $`\psi _n^0(x,t)`$ of the system should satisfy the Schrödinger equation
$$O_0(t)\psi _n^0(x,t)=0,$$
(4)
where $`O_0(t)=ih\frac{}{t}+H_0(x,p,t)`$. For the simple harmonic oscillator system of unit mass and frequency whose time is $`\tau `$, the wave functions $`\psi _n^s(x,\tau )`$ should satisfy
$$O_s(\tau )\psi _n^s(x,\tau )=0,$$
(5)
where $`O_s(\tau )=ih\frac{}{\tau }+H^s`$ with $`H^s=\frac{1}{2}(p^2+x^2).`$ If $`t`$ and $`\tau `$ is related through the relation
$$d\tau =\frac{\mathrm{\Omega }_0}{\rho _0^2}dt,$$
(6)
by defining the unitary operator
$$U_{w0}(\rho _0,\mathrm{\Omega }_0)=\mathrm{exp}(\frac{i\dot{\rho _0}}{2\mathrm{}\rho _0}x^2)\mathrm{exp}[\frac{i}{2\mathrm{}}(\mathrm{ln}\frac{\rho _0}{\sqrt{\mathrm{\Omega }_0}})(xp+px)],$$
(7)
one may find the relation
$$\frac{\mathrm{\Omega }_0}{\rho _0^2}U_{w0}O_s(\tau )U_{w0}^{}_{\tau =\tau (t)}=O_0(t).$$
(8)
In Eq. (7), the overdot denotes the differentiation with respect to time $`t`$, while in Eq. (8) the notation $`\mathrm{"}_{\tau =\tau (t)}\mathrm{"}`$ is to mention that $`\tau `$ should be replaced by the function of $`t`$ satisfying the relation (6). In a different vein, the relation (8) has also been noticed in Ref. . Eqs. (5,8) imply the following relation in wave functions;
$$\psi _n^0(x,t)=U_{w0}\psi _n^s|_{\tau =\tau (t)}.$$
(9)
As is well-known , the simplest choice of $`\{\psi _n^s|n=0,1,2,\mathrm{}\}`$ may be given as
$`\psi _n^s(x,\tau )|_{\tau =\tau (t)}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2^nn!\sqrt{\pi \mathrm{}}}}}e^{i(n+\frac{1}{2})\tau }`$ (11)
$`\times \mathrm{exp}[{\displaystyle \frac{x^2}{2\mathrm{}}}]H_n({\displaystyle \frac{1}{\sqrt{\mathrm{}}}}x)|_{\tau =\tau (t)}`$
$`=`$ $`{\displaystyle \frac{1}{\sqrt{2^nn!\sqrt{\pi \mathrm{}}}}}({\displaystyle \frac{u_0(t)iv_0(t)}{\rho _0(t)}})^{n+1/2}`$ (13)
$`\times \mathrm{exp}[{\displaystyle \frac{x^2}{2\mathrm{}}}+ic_0]H_n({\displaystyle \frac{1}{\sqrt{\mathrm{}}}}x),`$
where $`c_0`$ is an arbitrary real number which will be set to zero from now on. In obtaining Eq. (11), we make use of the fact:
$$d\tau =\frac{\mathrm{\Omega }_0}{\rho _0^2}dt=i(\frac{\dot{u}_0i\dot{v}_0}{u_0iv_0}\frac{\dot{\rho }_0}{\rho _0})dt.$$
(14)
In order to find a general expression of $`\psi _n^0(x,t)`$, we consider another unitary transformation. By defining $`\delta _{u_1}(t)`$ through the relations
$$\dot{\delta }_{u_1}=\frac{1}{2}w_0^2u_1^2\frac{1}{2}\dot{u}_1^2$$
(15)
where $`u_1`$ is a linear combination of $`u_0(t)`$ and $`v_0(t)`$, one may find that the unitary operator $`U_f`$ given as
$$U_f=\mathrm{exp}[\frac{i}{\mathrm{}}(\dot{u}_1x+\delta _{u_1}(t)]\mathrm{exp}(\frac{i}{\mathrm{}}u_1p)$$
(16)
satisfies the following relation
$$U_fO_0U_f^{}=O_0.$$
(17)
Therefore, the wave functions $`\psi _n^0`$ satisfying Schrödinger equation of Eq. (4) may in general be written as
$`\psi _n^0(x,t)`$ $`=`$ $`U_fU_{w0}\psi _n^s(x,\tau )_{\tau =\tau (t)}`$ (18)
$`=`$ $`{\displaystyle \frac{1}{\sqrt{2^nn!\rho _0(t)}}}({\displaystyle \frac{\mathrm{\Omega }_0}{\pi \mathrm{}}})^{1/4}({\displaystyle \frac{u_0(t)iv_0(t)}{\rho _0(t)}})^{n+1/2}`$ (22)
$`\times \mathrm{exp}[{\displaystyle \frac{i}{\mathrm{}}}(\dot{u}_1(t)x+\delta _{u_1}(t))]`$
$`\times \mathrm{exp}[{\displaystyle \frac{(xu_1(t))^2}{2\mathrm{}}}({\displaystyle \frac{\mathrm{\Omega }_0}{\rho _0^2(t)}}+i{\displaystyle \frac{\dot{\rho }_0}{\rho _0}})]`$
$`\times H_n(\sqrt{{\displaystyle \frac{\mathrm{\Omega }_0}{\mathrm{}}}}{\displaystyle \frac{xu_1(t)}{\rho _0(t)}}).`$
This wave function, of course, agrees with the known one if we consider $`u_1`$ as a (fictitious) particular solution. If $`u_1=0`$, the wave function given in Eq. (16) also agrees with that in Refs. .
It may be interesting to find that how many free parameters are in the wave function $`\psi _n^0(x,t)`$. First, there are two parameters in determining $`u_1(t)`$. In the case of $`u_1=0`$, one may think that there are four parameters which come from determining $`u_0(t),v_0(t)`$. However, one of them is not a free parameter, since the wave functions are invariant under the multiplication of $`u_0(t)`$ and $`v_0(t)`$ with same constant factor. For the simple harmonic oscillator of time-translational symmetry, one of the remaining three parameters of $`u_1=0`$ is simply related to a time-shifting of the wave function. This can be seen from the fact that, for the unit frequency case, the $`u_0`$ and $`v_0`$ can be taken as $`\mathrm{cos}(t+t_0)`$ and $`C\mathrm{sin}(t+\beta +t_0)`$, respectively, with real constants $`C,\beta ,t_0`$.
If one considers $`\rho _s(\tau )`$ satisfying
$$\frac{d^2}{d\tau ^2}\rho _s+\rho _s\frac{\mathrm{\Omega }_s^2}{\rho _s^3}=0,$$
(23)
and a simple harmonic oscillator of unit mass and frequency and with time $`\tau ^{}`$ which is related to $`\tau `$ as
$`d\tau ^{}={\displaystyle \frac{\mathrm{\Omega }_s}{\rho _s^2}}d\tau ,`$
by defining
$$U_s=\mathrm{exp}(\frac{id\rho _s/d\tau }{2\mathrm{}\rho _s}x^2)\mathrm{exp}[\frac{i}{2\mathrm{}}(\mathrm{ln}\frac{\rho _s}{\sqrt{\mathrm{\Omega }_s}})(xp+px)],$$
(24)
one may find that
$$\frac{\mathrm{\Omega }_s}{\rho _s^2}U_sO_s(\tau ^{})U_s^{}_{\tau ^{}=\tau ^{}(\tau )}=O_s(\tau ).$$
(25)
The wave functions $`\stackrel{~}{\psi }_n^s(\tau )`$ defined by
$`\stackrel{~}{\psi }_n^s(\tau )U_s\psi _n^s(x,\tau ^{})_{\tau ^{}=\tau ^{}(\tau )}`$
then satisfy the Schrödinger equation $`O_s(\tau )\stackrel{~}{\psi }_n^s=0`$; In fact, $`\stackrel{~}{\psi }_n^s(\tau )`$ is closely related to the wave functions of the squeezed states .
One may think that a more general expression of the unitary operator, $`U_{w0}`$, may be obtained by combining use of $`U_{w0}`$ and $`U_s`$. This, however, is not the case as can be seen from the relation
$$U_{w0}(\rho _0,\mathrm{\Omega }_0)U_s_{\tau =\tau (t)}=U_{w0}(\rho _0\rho _s,\mathrm{\Omega }_0\mathrm{\Omega }_s)_{\tau =\tau (t)},$$
(26)
which is in accordance with the number counting of free parameters in $`\psi _n^0(x,t)`$.
The harmonic oscillator of unit mass and frequency with an inverse-square potential is described by the Hamiltonian
$$H_{in}^s=\frac{p^2}{2}+\frac{x^2}{2}+\frac{g}{x^2}.$$
(27)
We only consider the case of $`g>\mathrm{}^2/8`$, and the region of $`x>0`$. By defining $`\alpha =\frac{1}{2}(1+8g/\mathrm{}^2)^{1/2}`$ and
$$O_s^{in}(\tau )=i\mathrm{}\frac{}{\tau }+H_{in}^s,$$
(28)
the wave function $`\varphi _n^s`$ satisfying the Schrödinger equation $`O_s^{in}(\tau )\varphi _n^s=0`$ is given as
$`\varphi _n^s`$ $``$ $`x|\varphi _n^s`$ (29)
$`=`$ $`({\displaystyle \frac{4}{\mathrm{}}})^{1/4}({\displaystyle \frac{\mathrm{\Gamma }(n+1)}{\mathrm{\Gamma }(n+\alpha +1}})^{1/2}e^{i(2n+\alpha +1)\tau }`$ (31)
$`\times ({\displaystyle \frac{x^2}{\mathrm{}}})^{(2\alpha +1)/4}\mathrm{exp}({\displaystyle \frac{x^2}{2\mathrm{}}})L_n^\alpha ({\displaystyle \frac{x^2}{\mathrm{}}}).`$
By defining $`O_0^{in}`$ as
$$O_0^{in}(t)=i\mathrm{}\frac{}{t}+\frac{p^2}{2}+w_0^2(t)\frac{x^2}{2}+\frac{g}{x^2},$$
(32)
as in the case without the inverse-square term, one may find the relation
$$\frac{\mathrm{\Omega }_0}{\rho _0^2}U_{w0}O_s^{in}(\tau )U_{w0}^{}_{\tau =\tau (t)}=O_0^{in}(t).$$
(33)
In deriving Eq. (25), we make use of the commutator relation
$$[xp+px,\frac{1}{x^2}]=4i\mathrm{}\frac{1}{x^2}.$$
(34)
For a further generalization, we define a unitary operator
$$U_g=\mathrm{exp}[\frac{i}{\mathrm{}}(Max^2\frac{\dot{M}}{4})x^2]\mathrm{exp}[i\frac{\mathrm{ln}M}{4\mathrm{}}(xp+px)],$$
(35)
where $`M`$ is a positive function of $`t`$, and $`a(t)`$ is a real function. One may then easily find the relation
$`U_gO_0^{in}U_g^{}`$ $``$ $`O^{in}`$ (36)
$`=`$ $`i\mathrm{}{\displaystyle \frac{}{t}}+H_{in},`$ (37)
where (see Ref. )
$`H_{in}`$ $`=`$ $`{\displaystyle \frac{p^2}{2M(t)}}a(t)(xp+px)`$ (39)
$`+{\displaystyle \frac{1}{2}}M(t)c(t)x^2+{\displaystyle \frac{g}{M(t)}}{\displaystyle \frac{1}{x^2}}`$
with
$`c(t)=w_0^2(t)+{\displaystyle \frac{1}{\sqrt{M}}}{\displaystyle \frac{d^2\sqrt{M}}{dt^2}}+4a^22{\displaystyle \frac{1}{M}}{\displaystyle \frac{d}{dt}}(Ma).`$
For convenience , we consider the equation
$$\frac{d}{dt}(M\dot{x})+w^2(t)x=0,$$
(40)
where $`w^2(t)=w_0^2(t)+\frac{1}{\sqrt{M}}\frac{d^2\sqrt{M}}{dt^2}.`$ The two linearly independent solutions $`u(t),v(t)`$ of Eq. (31) can be given from $`u_0(t),v_0(t)`$ as $`u(t)=\frac{u_0}{\sqrt{M}},v(t)=\frac{v_0}{\sqrt{M}}`$, so that one may find the relation $`\mathrm{\Omega }_0=M(\dot{v}u\dot{u}v)`$. We also define the $`\rho (t)`$ as $`\rho (t)=\frac{\rho _0}{\sqrt{M}}`$. The wave function $`\varphi _n`$ of the system described by the Hamiltonian $`H_{in}(x,p,t)`$ can then be obtained as
$`\varphi _n`$ $`=`$ $`U_G\varphi _n^s(\tau )_{\tau =\tau (t)}`$ (41)
$`=`$ $`({\displaystyle \frac{4\mathrm{\Omega }_0}{\mathrm{}\rho ^2}})^{1/4}({\displaystyle \frac{\mathrm{\Gamma }(n+1)}{\mathrm{\Gamma }(n+\alpha +1}})^{1/2}`$ (44)
$`\times ({\displaystyle \frac{uiv}{\rho }})^{(2n+\alpha +1)}({\displaystyle \frac{\mathrm{\Omega }_0x^2}{\mathrm{}\rho ^2}})^{(2\alpha +1)/4}`$
$`\times \mathrm{exp}[{\displaystyle \frac{x^2}{2\mathrm{}}}({\displaystyle \frac{\mathrm{\Omega }_0}{\rho ^2}}iM{\displaystyle \frac{\dot{\rho }}{\rho }}2iMa)]L_n^\alpha ({\displaystyle \frac{\mathrm{\Omega }_0x^2}{\mathrm{}\rho ^2}}),`$
where
$$U_G=U_gU_{w0}.$$
(45)
For $`a=0`$, the wave functions $`\varphi _n`$ agree with those in Refs. . As in Ref. , by considering the kernel of the system , it may be easy to see that the wave functions $`\varphi _n(x,t)`$ form a complete set. The form of $`\varphi _n`$ in Eq. (33) indicates that, even for the system described by the constant Hamiltonian $`H_{in}^s`$ given in Eq. (21), there are wave functions whose probability density distributions pulsate as in those of the squeezed states.
For the system of the Hamiltonian $`H_{in}`$, if $`M(t),w_0^2(t),`$ and $`a(t)`$ are periodic with a period $`T`$, one may study the non-adiabatic geometric phases . The wave function $`\varphi _n`$ is (quasi)periodic, only if $`\rho (t)`$ is periodic. The condition for periodic $`\rho (t)`$ with the period $`T^{}`$ $`(=T`$ or $`2T)`$ has been analyzed in Ref. . Here, we only consider the case of such a periodic $`\rho (t)`$. The overall phase change of $`\varphi _n`$ under the $`T^{}`$ evolution is given as
$`\chi _n=(2n+\alpha +1){\displaystyle _0^T^{}}{\displaystyle \frac{\mathrm{\Omega }_0}{\rho _0^2(t)}}𝑑t.`$
The expectation value of the $`H_{in}`$ can be evaluated by making use of the relation
$$H_{in}\varphi _n=(i\mathrm{}\frac{}{t}U^G)\varphi _n^s+U_Gi\mathrm{}\frac{}{t}\varphi _n^s.$$
(46)
From the fact that
$$i\mathrm{}\frac{}{t}\varphi _n^s=i\mathrm{}\frac{d\tau }{dt}\frac{}{\tau }\varphi _n^s=(2n+\alpha +1)\mathrm{}\frac{\mathrm{\Omega }_0}{\rho _0^2(t)}\varphi _n^s,$$
(47)
one may find that the geometric phase $`\gamma _n`$ for the wave function $`\varphi _n`$ under the $`T^{}`$ evolution is written as
$`\gamma _n`$ $`=`$ $`\chi _n+{\displaystyle \frac{1}{\mathrm{}}}{\displaystyle _0^T^{}}\varphi _n|H_{in}|\varphi _n𝑑t`$ (48)
$`=`$ $`{\displaystyle \frac{1}{\mathrm{}\mathrm{\Omega }_0}}{\displaystyle _0^T^{}}(M\dot{\rho }^2+2Ma\rho \dot{\rho })𝑑t\varphi _n^s|x^2|\varphi _n^s`$ (49)
$`=`$ $`(2n+\alpha +1){\displaystyle \frac{1}{\mathrm{\Omega }_0}}{\displaystyle _0^T^{}}(M\dot{\rho }^2+2Ma\rho \dot{\rho })𝑑t.`$ (50)
The unitary operators can be used in finding the exact invariants for the cases without and with the inverse-square potential from $`H^s`$ and $`H_{in}^s`$, respectively. First of all, it is clear that $`H^s`$ and $`H_{in}^s`$ are invariants in the systems they describe, respectively. For the system described by $`H_0(x,p,t)`$, if we only consider the case of $`u_1=0`$, the invariant $`I_0`$ is obtained by applying the unitary transformation to the invariant $`H^s`$
$`I_0`$ $`=`$ $`U_{w0}H^sU_{w0}^{}`$ (51)
$`=`$ $`{\displaystyle \frac{1}{2\mathrm{\Omega }_0}}[({\displaystyle \frac{\mathrm{\Omega }_0x}{\rho _0}})^2+(\rho _0p\dot{\rho }x)^2],`$ (52)
which agrees with those in Refs. . For the system described by $`H_{in}(x,p,t)`$, the invariant is again given from the invariant $`H_{in}^s`$ as
$`I_{in}`$ $`=`$ $`U_GH_{in}^sU_G^{}`$ (53)
$`=`$ $`{\displaystyle \frac{1}{2\mathrm{\Omega }_0}}[({\displaystyle \frac{\mathrm{\Omega }_0x}{\rho }})^2`$ (55)
$`+\{\rho p(M\dot{\rho }+2Ma\rho )x\}^2+2\rho ^2{\displaystyle \frac{g}{x^2}}].`$
For the case of $`a=0`$, the invariant $`I_{in}`$ reduces to the known one . One can explicitly check that the invariant $`I_{in}`$ indeed satisfies the relation
$$i\mathrm{}\frac{I_{in}}{t}+[I_{in},H_{in}]=0.$$
(56)
Alternatively, making uses of Eqs. (35,36) and relying on the completeness of the set $`\{\varphi _n^s|n=0,1,2,\mathrm{}\}`$, a simple proof of Eq. (40) may also be possible.
In summary, we have found a unitary operator which transforms a harmonic oscillator system of time-dependent frequency into that of a simple harmonic oscillator of different time-scale, with and without the inverse-square potential. Making use of the unitary operator, the exact invariants and wave functions of the time-dependent systems have been evaluated from the well-known results in the corresponding system of constant Hamiltonians. It should be mentioned, however, that the classical solutions of the time-dependent harmonic oscillator system must be found for actual applications, while the classical equation (see Eq. (2)) is formally equivalent to a one-dimensional time-independent Schrödinger equation (of arbitrary potential). The classical correspondent of unitary transformation is the canonical transformation which has been studied in the model . It would be interesting if the relationship could be used in finding relations among the quantities in classical and quantum mechanics such as that between the geometric phases and the Hannay’s angle (see Ref. ). |
warning/0002/astro-ph0002019.html | ar5iv | text | # Magnetized Atmospheres around Neutron Stars Accreting at Low Rates
## 1 Introduction
The problem of investigating the properties of radiation emitted by neutron stars (NSs) accreting at low rates, $`\dot{M}10^{10}10^{14}`$ g/s, became of interest after it was realized that the Galaxy may contain a large population of low luminosity magnetic accretors (see e.g. Nelson, Wang, Salpeter, & Wasserman nel95:1995 (1995)). The Galaxy should harbor more than $`10^3`$ Be/X–ray binaries with an accreting NS shining at $`10^{32}10^{34}`$ erg/s (Rappaport, & van den Heuvel rvdh82:1982 (1982); van den Heuvel, & Rappaport vdhr86:1986 (1986)). Moreover, assuming a supernova birth rate of 10–100 yr<sup>-1</sup>, $`10^810^9`$ old, isolated NSs (ONSs) should be present in the Galaxy. Accretion onto a strongly magnetized, moving neutron star may be severely hindered for different reasons, but there is the possibility that a small, albeit non–negligible, fraction of ONSs may be accreting directly from the interstellar medium and some of them might be above the sensitivity threshold of ROSAT (e.g. Treves, & Colpi tc91:1991 (1991); Blaes, & Madau bm93:1993 (1993); see Treves, Turolla, Zane, & Colpi review:1999 (1999) for a review).
At variance with neutron stars accreting at high rates, e.g. in X–ray pulsators, in low–luminosity sources the interaction of the escaping radiation with the inflowing material is of little importance, so they provide a much simpler case for investigating the physics of accretion in a strongly magnetized environment. For luminosities far below the Eddington limit, the accretion problem becomes germane to that of calculating the spectrum emerging from static atmospheres around cooling NSs. Spectra from cooling NSs have been widely investigated by a number of authors in connection with the X–ray emission from young, millisecond pulsars and isolated NSs, both for low and high magnetic fields and for different chemical compositions (see e.g. Romani rom87:1987 (1987); Shibanov, Zavlin, Pavlov & Ventura sh92:1992 (1992); Rajagopal, & Romani rr96:1996 (1996); Pavlov, Zavlin, Trümper, & Neuhäuser pztn96:1996 (1996)). Emerging spectra are not very different from a blackbody at the star effective temperature, the distinctive hardening present at low fields ($`B10^9`$ G) becoming less pronounced when the magnetic field is $`10^{12}10^{13}`$ G. Similar conclusions were reached for the spectrum emitted by low–luminosity, low–field, accreting NSs by Zampieri, Turolla, Zane, & Treves (ztzt95:1995 (1995), hereafter ZTZT) for a pure hydrogen atmospheric composition.
The search of isolated neutron stars with ROSAT produced in recent years half a dozen promising candidates (Walter, Wolk, & Neühauser wwn96:1996 (1996); Haberl et al. ha97:1997 (1997); Haberl, Motch & Pietsch hmp98:1998 (1998); Schwope et al. sch99:1999 (1999); Motch et al. mo99:1999 (1999); Haberl, Pietsch & Motch hpm99:1999 (1999)). All of them show a soft, thermal X–ray spectrum, with typical energies $`100`$ eV, and have an exceedingly large X–ray to optical flux ratio, $`10^4`$. Although their association with isolated neutron stars is firmly established, their interpretation in terms of an accreting or a cooling object is still a matter of lively debate. Present models predict rather similar spectral distributions in both cases, especially in the X–ray band. It is therefore of particular importance to improve our theoretical understanding of these two classes of sources, looking, in particular, for spectral signatures which can enable us to discriminate between them.
In this paper we present a first detailed calculation of spectra emitted by strongly magnetized ($`B10^{12}`$ G), accreting neutron stars, focusing our attention on low luminosities, $`L10^{30}`$$`10^{33}`$ erg/s, such as those expected from old neutrons stars accreting the interstellar medium. Spectral distributions are computed solving the transfer equations for the normal modes coupled to the hydrostatic equilibrium and the energy balance for different values of the accretion rate and the magnetic field.
We find that spectra emerging from magnetized, accretion atmospheres are blackbody–like in the X–ray band, in close agreement with the known results for cooling, magnetized atmospheres. However, accretion spectra show a new and distinctive feature at low energies, being characterized by an excess over the Raleigh–Jeans tail of the X–ray continuum below $`10`$ eV. The same behaviour is found in accretion atmospheres around unmagnetized neutron stars, but, as already pointed out by ZTZT, the X–ray spectrum is sensibly harder in this case. This result may be relevant in connection with the isolated neutron star candidate RX J18563.5-3754. Multiwavelength observations of this source indicate that, while ROSAT data are well fitted by a blackbody at $`T_{eff}60`$ eV, HST points lie above the extrapolation of the fit in the optical (Walter, & Matthews wm97:1997 (1997)).
The plan of the paper is as follows. The input physics relevant to our model is presented in §2: radiative transfer in a magnetized plasma is discussed in §2.1 and the structure of an accreting, magnetized atmosphere in §2.2. Computed spectra are presented in §3. Discussion and conclusions follow in §4.
## 2 The Model
### 2.1 Radiation Transfer
In this paper we consider a magnetized, nondegenerate, pure hydrogen, cold plasma, in which the main radiative processes are free–free emission/absorption and Thomson scattering. The plasma is in local thermal equilibrium (LTE) at temperature $`T`$. We consider a plane–parallel geometry with normal $`𝐧`$ parallel to the magnetic field $`𝐁`$ and to the $`z`$–direction. The stratification of the atmosphere is described by using as a parameter the scattering depth $`\tau `$, as defined for an unmagnetized medium
$$\tau =\kappa _{es}_z^{\mathrm{}}\rho 𝑑z$$
(1)
where $`z`$ is the coordinate variable, $`\rho `$ is the plasma density, $`\kappa _{es}=\sigma _T/m_p`$ is the Thomson opacity and $`\sigma _T`$ is the Thomson cross section.
In the following, we neglect collective plasma effects and consider only the limit $`\omega _p^2/\omega ^21`$, where $`\omega _p=(4\pi n_ee^2/m_e)^{1/2}`$ is the plasma frequency and $`n_e`$ electron density. We also consider only frequencies lower than the electron cyclotron frequency, $`\omega _{c,e}=eB/m_ec`$, so the semitransverse approximation can be assumed to hold. Since, for $`\tau 0.01`$, the temperature in the atmosphere is always $`10^7`$ K (see §2.2 and ZTZT) and scattering dominates over true absorption only for $`\tau 1`$, Comptonization is negligible. For this reason, similarly to what is done for cooling, magnetized atmospheres (see e.g. Shibanov, et al. sh92:1992 (1992)), only conservative scattering is accounted for in the transfer equations (see, however, the discussion in §2.2 for the role of Compton heating/cooling in the energy balance of the external atmospheric layers).
Under these assumptions, the coupled equations for the transfer of the two normal modes take the form (see e.g. Gnedin & Pavlov gp74:1974 (1974); Yahel y80:1980 (1980))
$`y_G\mu {\displaystyle \frac{df^1}{d\tau }}`$ $`=`$ $`{\displaystyle K_s^{11}f^1^{}𝑑\mu ^{}}+{\displaystyle K_s^{21}f^2^{}𝑑\mu ^{}}k_s^1f^1+k_{ab}^1\left({\displaystyle \frac{1}{2}}b_\nu f^1\right)`$
$`y_G\mu {\displaystyle \frac{df^2}{d\tau }}`$ $`=`$ $`{\displaystyle K_s^{22}f^2^{}𝑑\mu ^{}}+{\displaystyle K_s^{12}f^1^{}𝑑\mu ^{}}k_s^2f^2+k_{ab}^2\left({\displaystyle \frac{1}{2}}b_\nu f^2\right)`$ (2)
where $`y_G=\sqrt{12GM/Rc^2}`$ is the gravitational redshift factor, $`R`$ and $`M`$ are the star mass and radius, $`\mu =𝐧𝐬=\mathrm{cos}\theta `$, $`f^i(\tau ,\mu ,\nu )=c^2I^i/2h^4\nu ^3`$ denotes the photon occupation number for the ordinary ($`i=1`$) and extraordinary ($`i=2`$) mode, $`I^i`$ is the specific intensity, $`b_\nu =c^2B_\nu /2h^4\nu ^3`$, and $`B_\nu `$ is the Planck function. In equations (2.1) $`k_{ab}^i`$ is the total free–free opacity,
$$k_s^i(\mu ,\nu )=\underset{j=1}{\overset{2}{}}K_s^{ij}𝑑\mu ^{}$$
(3)
and $`K_s^{ij}(\mu ,\mu ^{},\nu )`$ is the probability that an incident photon, which has polarization $`\widehat{e}^i`$ and propagates in the direction $`\mu `$, scatters into a direction $`\mu ^{}`$ and polarization $`\widehat{e}^j`$. All the quantities appearing in equations (2.1) are referred to the local observer, at rest on the star surface; the photon energy measured by an observer at infinity is given by $`h\nu _{\mathrm{}}=y_Gh\nu `$. The integral terms appearing into equations (2.1) account for the scattering emissivities, and all opacity/emissivity coefficients are normalized to $`\kappa _{es}`$. The expressions for the opacities relevant to the present calculation are reported in appendix A.
### 2.2 Atmospheric Structure
Accretion atmosphere models are constructed by solving the transfer equations (2.1) coupled to the hydrostatic equilibrium and the energy equation. The hydrostatic balance is simply expressed as
$$\frac{dP}{d\tau }=\frac{GM}{y_G^2R^2\kappa _{es}},$$
(4)
where $`P=k\rho T/\mu _em_p`$ ($`\mu _e1/2`$ for completely ionized hydrogen) is the gas pressure and we consider only the case $`L/L_{Edd}1`$, where $`L=L(\tau )`$ is the luminosity measured by the local observer. Since in all our models the ram pressure of the accreting material, $`1/2\rho v^2`$, turns out to be much smaller than the thermal pressure, it has been neglected, together with the radiative force. In this limit equation (4) is immediately integrated, and gives the density as a function of depth
$$\rho =\frac{GMm_p}{2y_G^2R^2\kappa _{es}}\frac{\tau }{kT(\tau )}.$$
(5)
The energy balance just states that the net radiative cooling must equate the heating $`W_H`$ supplied by accretion. The radiative energy exchange is obtained adding equations (2.1) together, after multiplying them by $`(\mathrm{}\omega )^3`$, and integrating over angles and energies. Since we assumed conservative scattering, its contribution to the energy balance clearly vanishes. However, as discussed in previous investigations (Alme, & Wilson aw73:1973 (1973); ZTZT), Comptonization is ineffective in modifying the spectrum, but plays a crucial role in determining the temperature in the external, optically thin layers. Contrary to what happens in cooling atmospheres, the temperature profile shows a sudden rise (or “jump”) in the external, low–density layers where the heating produced by the incoming protons is mainly balanced by Compton cooling. Including Compton heating/cooling the energy equation becomes
$$k_P\left(\frac{aT^4}{2}\frac{k_{am}^1}{k_P}U^1\frac{k_{am}^2}{k_P}U^2\right)+\left(\mathrm{\Gamma }\mathrm{\Lambda }\right)_C=\frac{W_H}{c\kappa _{es}}$$
(6)
where $`U^i`$ is the radiation energy density of mode $`i`$ and $`k_P`$, $`k_{am}^i`$ are defined in strict analogy with the Planck and absorption mean opacities in the unmagnetized case. In evaluating the previous expression the approximated formula by Arons, Klein, & Lea (akl87:1987 (1987)) for the Compton rate in a magnetized plasma, $`(\mathrm{\Gamma }\mathrm{\Lambda })_C`$, has been used.
The detailed expressions for the heating rate $`W_H`$ and the stopping depth $`\tau _B`$ in a magnetized atmosphere are presented in appendix B.
## 3 Results
### 3.1 Numerical Method
The numerical calculation was performed adapting to radiative transfer in a magnetized medium the tangent–ray code developed by Zane, Turolla, Nobili, & Erna (crm:1996 (1996)) for one–dimensional, general–relativistic radiation transfer. The method performs an ordinary $`\mathrm{\Lambda }`$–iteration for computing the scattering integrals. Schematically, the calculation proceeds as follows. First an initial temperature profile is specified (usually that of an unmagnetized model with similar parameters calculated by ZTZT) and the zero–th order approximation for $`f^1`$, $`f^2`$ is computed solving equations (2.1) with no scattering emissivity. The boundary conditions for ingoing ($`\mu <0`$) trajectories are $`f^i=0`$ at $`\tau =\tau _{min}`$ while diffusive boundary conditions at $`\tau =\tau _{max}`$ were used for outgoing ($`\mu >0`$) trajectories. The computed intensities are then used to evaluate the scattering integrals and the whole procedure is repeated, keeping the temperature profile unchanged. As soon as corrections on the intensities are small enough, new temperature and density profiles are obtained solving equations (6) and (5). The whole scheme is then iterated to convergence. Each model is completely characterized by the magnetic field strength, the total luminosity and the luminosity at $`\tau _B`$, or, equivalently, by $`B`$ and the heating rate $`W_H`$ (see appendix B). Since the code solves the full transfer problem, it allows for the complete determination of the radiation field, including its angular dependence. This ensures a more accurate treatment of non–anisotropic radiative process with respect to angle–averaged, diffusion approximations. Owing to the gravitational redshift, the total accretion luminosity measured at infinity is related to the local luminosity at the top of the atmosphere by $`L_{\mathrm{}}=y_G^2L(0)`$.
Models presented below were computed using a logarithmic grid with 300 equally–spaced depth points, 20 equally–spaced angular points and 48 energies. We explored a wide range of luminosities and considered two representative values of the magnetic field, $`B=10^{12}`$ and $`10^{13}`$ G. The model parameters are reported in Table 1, together with the accretion rate
$$\dot{M}=L(0)\frac{\left[1L(\tau _B)/L(0)\right]}{\eta c^2};$$
(7)
here $`\eta =1y_G`$ is the relativistic efficiency. Since the values of $`\tau _B`$ and $`\omega _{c,e,p}`$ depend on $`B`$, the adopted boundaries in depth and energy vary from model to model. Typical values are $`\tau _{max}10^2`$ and $`\tau _{min}10^610^8`$; the energy range goes from 0.16 eV to 5.45 keV. The angle–averaged effective depth is always $`100`$ at $`\tau _{max}`$ and $`10^5`$ at $`\tau _{min}`$.
Convergence was generally achieved after 20–30 iterations with a fractional accuracy $`0.01`$ both in the hydrodynamical variables and in the radiation field. As a further check on the accuracy of our solutions, the luminosity evaluated numerically was compared with
$$L(\tau )L(0)\left[L(0)L(\tau _B)\right]\frac{1[1(1v_{th}^4/v_{ff}^4)(\tau /\tau _B)]^{1/2}}{1v_{th}^4/v_{ff}^4}$$
(8)
which follows integrating the first gray moment equation
$$\frac{dL}{d\tau }=\frac{4\pi R^2f_AW_H}{y_G\kappa _{es}}$$
(9)
and provides an analytical expression for $`L`$ at $`\tau <\tau _B`$ (see appendix B for notation). Within the range of validity of equation (8), the two values of $`L`$ differ by less than 4%.
In all models $`M=1M_{}`$, $`R=6GM/c^20.89\times 10^6`$ cm which correspond to $`\tau _s3.3`$ (see equation \[B2\]). Only model A5, the unmagnetized one, was computed with $`\tau _s8`$. One of the largest complication introduced by the presence of the magnetic field is the large–scale pattern of the accretion flow. In particular, when the accretion rate is small the way the spherically symmetric infalling plasma enters the magnetosphere is not fully understood as yet (see e.g. Blaes, & Madau bm93:1993 (1993), Arons & Lea al80:1980 (1980)). In order to bracket uncertainties, we assume a fiducial value for the fraction of the star surface covered by accretion, $`f_A=0.01`$. We stress that this choice has no effect on the spectral properties of the emitted radiation and affects only the total luminosity, which scales linearly with $`f_A`$.
### 3.2 Emerging Spectra
The emergent spectra for models A1–A2 and A3–A4 are shown in figures 1 and 2, together with the blackbody at the NS effective temperature, $`T_{eff}`$. An unmagnetized model (A5) with similar luminosity is shown in figure 3. The corresponding temperature profiles are plotted in figure 4. Note that in all the plots the photon energy is already corrected for the gravitational redshift, so the spectral distribution is shown as a function of the energy as observed at Earth. As it is apparent comparing the different curves (see also the solutions computed by ZTZT) the thermal stratification of the atmosphere shows the same general features (inner layers in LTE, outer region dominated by Comptonization) independently of $`B`$.
The sudden growth of $`T`$ (up to $`10^710^8`$ K) that appears in the external layers is basically due to the fact that free–free cooling can not balance the heating produced by accretion at low densities. The temperature then rises until it reaches a value at which Compton cooling becomes efficient. This effect has been discussed by ZTZT and Zane, Turolla, & Treves (1998) in connection with unmagnetized atmospheres and they have shown that the temperature jump is located at the depth where the free–free and Compton thermal timescales become comparable, $`t_{ff}t_C`$. In the magnetized case the situation is very similar, but now there are two relevant free–free timescales, $`t_{ff}^{(1)}`$ and $`t_{ff}^{(2)}`$, one for each mode. Comptonization becomes the dominant cooling process for $`t_C\mathrm{min}(t_{ff}^{(1)},t_{ff}^{(2)})`$ and gives rise to the large jump present in all accretion models. For some values of the model parameters, the region where $`\rho =\rho _{vac}`$ coincides with the photospheric region for both modes (see appendix A.3). In this case, vacuum effects can produce the peculiar “double jump” structure present in model A1, with the first (small) jump located at a depth where $`t_{ff}^{(1)}/t_{ff}^{(2)}1`$.
In order to compare our results with models available in the literature, some cooling atmospheres have been also computed setting $`W_H=0`$ in equation (6); accordingly, the luminosity is now a constant. The emergent spectra for two such models (C1 and C2, see Table 1) show a good agreement with those computed by Shibanov et al. (sh92:1992 (1992)) using the diffusion approximation.
We find that the spectral hardening at low luminosities, typical of unmagnetized atmospheres, is far less pronounced but still present up to field strengths $`10^{13}`$ G and tends to disappear at large enough luminosities. For comparison, with $`L4\times 10^{33}`$ erg/s, the magnetized ($`B=10^{12}`$ G) spectrum has negligible hardening while the hardening ratio is still $`1.6`$ in unmagnetized models of similar luminosity (ZTZT). The overall dependence of the continuum is not particularly sensitive to the value of the magnetic field, although the absorption feature at the proton cyclotron energy becomes more prominent with increasing $`B`$.
### 3.3 The Optical Excess
The most striking result emerging from our computations is that, although spectra from cooling and accreting H atmospheres are rather similar in the X–rays, they differ substantially at low energies. Below $`10`$ eV spectra from accreting atmospheres exhibit a soft excess with respect to the blackbody spectrum which is not shared by the cooling models. This feature, which is present also in the unmagnetized case, can be viewed as a distinctive spectral signature of a low–luminosity, accreting neutron star. The fact that it was not reported by ZTZT (and by previous investigators, see e.g. Alme, & Wilson aw73:1973 (1973)) is because they were mainly interested in the shape of the X–ray continuum and their energy range was not large enough to cover the optical band; besides, some numerical problems, related to the moment formalism used to solve the transfer, prevented ZTZT to reach very low frequencies. The evidence of the soft excess is even more apparent in figures 5, 6 and 7 where synthetic spectra are plotted together with the best fitting blackbody in the X–ray band. The excess at two selected optical wavelengths ($`\lambda =3000,6060`$ A, see discussion below), together with the temperature $`T_{fit}`$ of the best–fitting blackbody in the X–ray band, are reported in Table 2.
The appearance of an optical excess in accreting models is related to the behaviour of the temperature, which is different from that of cooling models. The external layers are now hotter because Comptonization dominates the thermal balance there. The low energy tail of the spectrum decouples at a depth that corresponds to the (relatively) high temperatures near the jump, and emerges at infinity as a planckian at a temperature higher than $`T_{eff}`$. By decreasing the luminosity, the temperature jump moves at lower scattering depth and the frequency below which the spectrum exceeds the blackbody at $`T_{eff}`$ becomes lower.
The presence of an optical excess has been reported in the spectrum of a few isolated, nearby pulsars (Pavlov, Stringfellow, & Córdova psc96:1996 (1996)) and, at lower luminosities, in the spectrum of the ONSs candidate RX J18563.5-3754. RX J18563.5-3754 was observed in the X–ray band with ROSAT (Walter, Wolk, & Neuhäuser wwn96:1996 (1996)) and by HST at $`\lambda =3000`$ and 6060 A (Walter, & Matthews wm97:1997 (1997)). These multiwavelength observations made evident that the spectrum of RX J18563.5-3754 is more complex than a simple blackbody. The blackbody fit to PSPC data underpredicts the optical fluxes $`f_{3000}`$ and $`f_{6060}`$ by a factor 2.4 and 3.7 respectively (Walter, & Matthews wm97:1997 (1997); see also Pavlov et al. pztn96:1996 (1996)). Models of cooling atmospheres based on different chemical compositions also underestimate the optical fluxes (Pavlov et al. psc96:1996 (1996)), while models with two blackbody components or with a surface temperature variation may fit the $`f_{3000}`$ flux. Recent spectra from non–magnetic atmospheres with Fe or Si–ash compositions (see Walter, & An wa98:1998 (1998)) may also provide a fit of both the X–ray and optical data, although these spectral models agree with those of Rajagopal, & Romani (rr96:1996 (1996)) but not with those of Pavlov et al. (pztn96:1996 (1996)). Given the considerable latitude of the unknown parameters $`f_A`$, $`L`$, $`B`$, however, results presented here indicate that the full spectral energy distribution may be consistent with the picture of an accreting NS. We want also to note that all models computed here have a blackbody temperature higher than that required to fit the X–ray spectrum of RX J18563.5-3754 (see Table 2) and, although we are far from having explored the model parameter space, present results indicate that the excess decreases for decresing luminosities. Although the optical identification of another isolated NS candidate, RX J0720.4-3125, still lacks a definite confirmation, it is interesting to note that the counterpart proposed by Kulkarni & van Kerkwijk (kvk98:1998 (1998)) also shows a similar excess.
### 3.4 Fraction of Polarization
The fraction of polarization strongly depends on the energy band and, in the presence of a temperature and density gradient, shows a variety of different behaviours (see figure 8). Its sign is determined by the competition between plasma and vacuum properties in the photospheric layers. However, independently of the model parameters, the degree of polarization crosses zero at the very vicinity of the proton cyclotron energy (see eq. \[A10\]), where the mode absorption coefficients cross each other. The bulk of the thermal emission from low–luminosity accreting NSs falls in the extreme UV/soft X–ray band, which is subject to strong interstellar absorption, making difficult their identification. It has been suggested that the detection of the non–thermal cyclotron emission feature at highest energies may be a distinguishing signature for most of these low–luminosity sources (Nelson, et al. nel95:1995 (1995)). Our results show that observations of the proton cyclotron line combined with measures of polarization may also provide a powerful tool to determine the magnetic field of the source, even in the absence of pulsations.
## 4 Discussion and Conclusions
We have discussed the spectral distribution of the radiation emitted by a static, plane–parallel atmosphere around a strongly magnetized neutron star which is heated by accretion. Synthetic spectra have been computed solving the full transfer problem in a magnetoactive plasma for several values of the accretion luminosity and of the star magnetic field. In particular, we explored the low luminosities ($`L10^{30}`$$`10^{33}`$ erg/s), typical e.g. of isolated accreting NSs, and found that model spectra show a distinctive excess at low energies over the blackbody distribution which best–fits the X–rays. The energy $`\mathrm{}\omega _{ex}`$ at which the spectrum rises depends both on the field strength and the luminosity, but for $`B10^{12}`$$`10^{13}`$ G $`\mathrm{}\omega _{ex}10`$ eV, so that the optical emission of an accreting, NS is enhanced with respect to what is expected extrapolating the X–ray spectrum to optical wavelengths. The presence of an optical/UV bump is due to the fact that the low–frequency radiation decouples at very low values of the scattering depth in layers where the gas is kept at larger temperatures by Compton heating. Since a hot outer zone is exhibited by accreting and not by cooling atmospheres, the optical excess becomes a distinctive spectral feature of NSs accreting at low rates.
Despite present results are useful in shedding some light on the emission properties of accreting, magnetized neutron stars, many points still need further clarification before the problem is fully understood and some of the assumptions on which our investigation was based deserve a further discussion. In this paper we have considered only fully ionized, pure hydrogen atmospheres. In the case of cooling NSs, the emitted spectra are strongly influenced by the chemical composition of the surface layers which results from the supernova explosion and the subsequent envelope fallback. Present uncertainties motivated several authors to compute cooling spectra for different abundances (see e.g. Miller, & Neuhäuser mn91:1991 (1991); Miller mil92:1992 (1992); Pavlov, et al. pszm95:1995 (1995); Rajagopal, Romani, & Miller rrm97:1997 (1997)) and led to the suggestion that the comparison between observed and synthetic X–ray spectra may probe the chemistry of the NS crust (Pavlov, et al. pztn96:1996 (1996)). The assumption of a pure hydrogen atmosphere, although crude, is not unreasonable for an accretion atmosphere. In this case, in fact, incoming protons and spallation by energetic particles in the magnetosphere may enrich the NS surface with light elements (mainly H) and, owing to the rapid sedimentation, these elements should dominate the photospheric layers (Bildsten, Salpeter, & Wasserman bsw92:1992 (1992)).
Even in this simplified picture, however, the assumption of complete ionization may not provide an entirely realistic description. In a strongly magnetized plasma atomic binding is greatly enhanced, mainly for light elements. For a typical field of $`10^{12}`$ G, the ionization potentials for the hydrogen ground state are $`100300`$ eV, moving the photoionization thresholds into the soft X–rays (see e.g. Ruderman r72:1972 (1972); Shibanov et al. sh92:1992 (1992)). The role of photoionization and of pressure ionization in a magnetized hydrogen atmosphere has been investigated by a number of authors (e.g Ventura, Herold, Ruder, & Geyer vhrg92:1992 (1992); Potekhin, & Pavlov pp97:1997 (1997); Potekhin, Shibanov, & Ventura psv98:1998 (1998)). The neutral fraction, $`f_H`$, reaches a peak at densities $`1\mathrm{g}/\mathrm{cm}^3`$ then decreases due to pressure ionization and turns out to be highly dependent on the temperature. While $`f_H`$ never exceeds a few percent at $`T=10^6`$ K, it becomes as large as $`80\%`$ for $`T=10^{5.5}`$ K and $`B=10^{13}`$ G. Although in the models we presented here such low temperatures are only reached in relatively low–density layers (see figure 4), hydrogen ionization equilibrium should be properly included in a more detailed analysis.
The effects of different orientations of the magnetic field through the atmosphere have been also neglected by keeping $`𝐁||𝐧`$, a key simplifying assumption which allowed us to solve the transfer problem in one spatial dimension. Clearly, such an approximate description is valid only if the size of the emitting caps is small. Due to the intrinsic anisotropy of a magnetized medium, the emerging flux is expected to depend on the angle $`\theta _B`$ between the magnetic field and the normal to the surface. Shibanov et al. (sh92:1992 (1992)) estimated that in a cooling atmosphere the flux at 1 keV from a surface element perpendicular to the field may exceed that of an element parallel to $`𝐁`$ by nearly 50%. This result shows that a non–uniform magnetic field may lead to a significant distortion in the emerging spectra. Moreover, if the orientation of $`B`$ varies along the atmosphere, tangential components of the radiative flux may induce some meridional circulation to maintain the heat balance (Kaminker, Pavlov, & Shibanov kps82:1982 (1982)).
When the finite size of the emitting regions is accounted for, the problem can not reduced to a plane–parallel geometry and its solution necessarily demands for multidimensional transfer algorithms (e.g. Burnard, Klein, & Arons bka88:1988 (1988), bka90:1990 (1990); Hsu, Arons, & Klein hak97:1997 (1997)). Calculations performed so far were aimed to the solution of the frequency–dependent radiative problem on a fixed background or of the full radiation hydrodynamical problem in the frequency–integrated case. Numerical codes for the solution of the full transfer problem in axially symmetric media under general conditions are now available, see e.g. ZEUS (Stone, & Norman sn92:1992 (1992)), ALTAIR (Dykema, Castor, & Klein dck96:1996 (1996)) and RADICAL (Dullemond, & Turolla dt99:1999 (1999)). Their application to the transfer of radiation in accretion atmospheres around magnetized NSs can add new insights on the properties of the emitted spectra.
We are grateful to G.G. Pavlov and to V.E. Zavlin for several very helpful discussions, and to A.Y. Potekhin for providing us with the expressions for the ion opacities. We also thank S. Rappaport for calling our attention to some useful references.
## Appendix A Radiative Processes in a Magnetized Medium
### A.1 Electron Scattering
The electron contribution to the scattering source terms has been evaluated using the expression of the differential cross section, $`d\sigma ^{ij}/d\mathrm{\Omega }`$, discussed by Ventura (v79:1979 (1979)) and Kaminker, Pavlov, & Shibanov (kps82:1982 (1982)). The electron contribution to $`K_s^{ij}`$ can be simply written as
$$K_{s,e}^{ij}=\frac{1}{m_p\kappa _{es}}𝑑\varphi ^{}\frac{d\sigma ^{ij}}{d\mathrm{\Omega }}\frac{3}{4}\underset{\alpha =1}{\overset{1}{}}\left|e_\alpha ^j^{}\right|^2\left|e_\alpha ^i\right|^2\frac{1}{\left(1+\alpha u^{1/2}\right)^2+\gamma _r^2}.$$
(A1)
where $`u=\omega _{c,e}^2/\omega ^2`$, the $`e_\alpha ^i`$ are the components of the normal mode unit polarization vector in a coordinate frame with the $`z`$–axis along $`𝐁`$ and $`\gamma _r=(2/3)(e^2/m_ec^3)\omega `$ is the radiation damping. Further integration over $`\theta ^{}`$ gives the total opacities
$$k_{s,e}^i=\underset{j=1}{\overset{2}{}}K_{s,e}^{ij}𝑑\mu ^{}\underset{\alpha =1}{\overset{1}{}}\left|e_\alpha ^i\right|^2\frac{1}{\left(1+\alpha u^{1/2}\right)^2+\gamma _r^2}.$$
(A2)
For energies near the proton cyclotron frequency $`\omega _{c,p}=(m_e/m_p)\omega _{c,e}`$, Thomson scattering on ions becomes important. The corresponding opacity has been derived by Pavlov et al. (pszm95:1995 (1995)) in a relaxation–time approximation and it is
$$K_{s,p}^{ij}\frac{3}{4}\mu _m^2\underset{\alpha =1}{\overset{1}{}}\left|e_\alpha ^j^{}\right|^2\left|e_\alpha ^i\right|^2\frac{1}{\left(1\alpha u_p^{1/2}\right)^2+\mu _m^2\gamma _r^2}.$$
(A3)
$$k_{s,p}^i\mu _m^2\underset{\alpha =1}{\overset{1}{}}\left|e_\alpha ^i\right|^2\frac{1}{\left(1\alpha u_p^{1/2}\right)^2+\mu _m^2\gamma _r^2},$$
(A4)
where $`\mu _m=m_e/m_p`$ and $`u_p=\omega _{c,p}^2/\omega ^2`$. The total opacities $`K_s^{ij}`$ and $`k_s^i`$ appearing in equations (2.1) are then evaluated by adding the contributions of the two species.
### A.2 Bremsstrahlung
The electron and proton contributions to free–free opacity have a structure similar to that discussed for scattering. In a pure hydrogen plasma, they are given by (see Pavlov, & Panov pp76:1976 (1976); Mészáros mes92:1992 (1992); Pavlov et al. pszm95:1995 (1995))
$$k_{ab,e}^i\frac{\kappa _{ff}}{\kappa _{es}}\underset{\alpha =1}{\overset{1}{}}\left|e_\alpha ^i\right|^2\frac{g_\alpha }{\left(1+\alpha u^{1/2}\right)^2+\gamma _r^2},$$
(A5)
$$k_{ab,p}^i\frac{\kappa _{ff}}{\kappa _{es}}\mu _m^2\underset{\alpha =1}{\overset{1}{}}\left|e_\alpha ^i\right|^2\frac{g_\alpha }{\left(1\alpha u^{1/2}\right)^2+\mu _m^2\gamma _r^2}$$
(A6)
where
$$\kappa _{ff}=4\pi ^2\alpha _F^3\frac{\mathrm{}^2c^2}{m_e^2}\frac{n_e^2}{v_T\omega ^3}\left[1\mathrm{exp}\left(\mathrm{}\omega /kT\right)\right],$$
(A7)
$`g_0=g_{||}`$, $`g_1=g_{+1}=g_{}`$, and $`g_{||}`$, $`g_{}`$ are the modified Gaunt factors, which account for the anisotropy induced by the magnetic field. The quantity $`\kappa _{ff}`$ is the free–free opacity of a non–magnetic plasma apart from a factor $`(4\pi /3\sqrt{3})g`$, where $`g`$ is the unmagnetized Gaunt factor. The total absorption opacity can be then evaluated by summing over the two species.
The modified Gaunt factors were computed evaluating numerically their integral form as given by Pavlov, & Panov (pp76:1976 (1976)). At low temperatures and small frequencies ($`u1`$), where direct numerical quadrature becomes troublesome, the Gaunt factors have been obtained from the simpler formulas by Nagel (n80:1980 (1980); see also Mészáros mes92:1992 (1992); Rajagopal, Romani, & Miller rrm97:1997 (1997) and references therein). In order to decrease the computational time, the Gaunt factors have been evaluated once for all over a sufficiently large grid of temperatures and frequencies. In the transfer calculation they are then obtained at the required values of $`T`$ and $`\omega `$ by polynomial interpolation.
### A.3 Vacuum Effects and Mode Switching
The opacities of a real plasma start to change, due to the vacuum corrections in the polarization eigenmodes, when the field approaches the critical value $`B_c=m_e^2c^3/\mathrm{}e4.41\times 10^{13}`$ G. The vacuum contribution has ben included modifying the expressions for the $`e_\alpha ^i`$ as discussed by Kaminker, Pavlov, & Shibanov (kps82:1982 (1982)), and is controlled by the vacuum parameter $`W`$
$$W=\left(\frac{3\times 10^{28}\mathrm{cm}^3}{n_e}\right)\left(\frac{B}{B_c}\right)^4.$$
(A8)
The inclusion of vacuum and of the protons produces the breakdown of the NM approximation near the mode collapse points (MCPs; see e.g. Pavlov, & Shibanov ps79:1979 (1979); Mészáros mes92:1992 (1992) and references therein). For $`W>4`$, or $`\rho <\rho _{vac}=3.3\times 10^3(B/10^{12}\mathrm{G})^4\mathrm{g}/\mathrm{cm}^3`$, the MCPs appear at the two critical frequencies
$$\omega _{c1,2}^2=\frac{1}{2}\omega _{c,e}^2\left[1\pm \left(1\frac{4}{W}\right)^{1/2}\right].$$
(A9)
MCPs play an important role in the transfer of radiation through a magnetized medium, since the absorption coefficients of the two modes either cross each other or have a close approach, depending on the angle. Following the discussion by Pavlov, & Shibanov (ps79:1979 (1979)), under the typical conditions at hand mode switching is likely to occur at nearly all values of $`\mu `$. For this reason and for the sake of simplicity, in the present calculation we assumed mode switching for any value of $`\mu `$ at the two vacuum critical frequencies.
As shown by Bulik, & Pavlov (bp96:1996 (1996)) in a fully ionized hydrogen plasma the presence of protons introduces (even in the absence of vacuum) a new MCP at
$$\omega _{c,3}=\frac{\omega _{c,p}}{\sqrt{1\omega _{c,p}/\omega _{c,e}+\omega _{c,p}^2/\omega _{c,e}^2}}\omega _{c,p}\left(1+\frac{m_e}{2m_p}\right).$$
(A10)
At $`\omega =\omega _{c,3}`$ the opacity coefficients cross (Zavlin, private communication), and the MCP related to the proton contribution is again a mode switching point. Although we are aware of no detailed calculation of the polarization modes in a “protons + electrons + vacuum” plasma, it seems natural to assume that, at least if $`\omega _{c,2}>\omega _{c,3}`$, the proton contribution to the “vacuum” term may be safely neglected, being a function of the ratio $`m_e/m_p`$.
## Appendix B Stopping Depth and Heating Rate in a Magnetized Atmosphere
In the unmagnetized case, under the assumption that all the proton stopping power is converted into electromagnetic radiation within the atmosphere, $`W_H`$ can be approximated as (Alme, & Wilson aw73:1973 (1973); ZTZT)
$$W_H\{\begin{array}{cc}\frac{y_G\kappa _{es}}{8\pi R^2\tau _sf_A}\left[L(0)L(\tau _s)\right]\frac{f(x_e)}{[1(1v_{th}^4/v_{ff}^4)(\tau /\tau _s)]^{1/2}}\hfill & \tau <\tau _s\hfill \\ & \\ 0\hfill & \tau \tau _s\hfill \end{array}$$
(B1)
where $`f_A`$ is the fraction of the star surface covered by accretion, $`\tau _s`$ is the proton stopping depth and $`v_{th}^2/v_{ff}^2=3kT(\tau _s)R/(2m_pGM)`$ is the squared ratio of the proton thermal velocity to the free–fall velocity. For cold atmospheres in which the proton kinetic energy is much larger than the thermal energy, $`f(x_e)`$ can be safely taken to be unity for all practical purposes.
Infalling protons loose their energy to electrons via Coulomb collisions and the generation of collective plasma oscillations and $`\tau _s`$ can be approximated as (e.g. Zel’dovich, & Shakura zs69:1969 (1969); Nelson, Salpeter, & Wasserman nel93:1993 (1993))
$$\tau _s\frac{1}{6}\frac{m_p}{m_e}\frac{v_{ff}^4}{c^4\mathrm{ln}\mathrm{\Lambda }_c}2.6\left(\frac{M}{M_{}}\right)^2\left(\frac{R}{10^6\mathrm{cm}}\right)^2\left(\frac{10}{\mathrm{ln}\mathrm{\Lambda }_c}\right)$$
(B2)
where $`\mathrm{ln}\mathrm{\Lambda }_c`$ is a (constant) Coulomb logarithm.
The proton stopping process in strongly magnetized atmospheres presents a few differences, and has been discussed in detail by Nelson, Salpeter, & Wasserman (nel93:1993 (1993)). Now the proton stopping depth, $`\tau _B`$, is larger than $`\tau _s`$, essentially because the magnetic field reduces the effective Coulomb logarithm perpendicular to the field. The expression for the proton stopping power retains, nevertheless, the same form as in the unmagnetized case. In the present calculation we used their approximate expression which relates $`\tau _B`$ to $`\tau _s`$ at different field strengths
$$\frac{\tau _B}{\tau _s}\{\begin{array}{cc}\frac{\mathrm{ln}\mathrm{\Lambda }_c}{\mathrm{ln}(2n_{max})}\hfill & n_{max}1\hfill \\ & \\ 2\mathrm{ln}(m_p/m_e)15\hfill & n_{max}1\hfill \end{array}$$
(B3)
$`n_{max}=m_ev_{ff}^2/(2\mathrm{}\omega _{c,e})`$ $`6.4(M/M_{})(R/10^6\mathrm{cm})^1(B/10^{12}\mathrm{G})^1`$. Since proton–proton interactions limit the stopping length to the mean–free path for nuclear collisions, which corresponds to $`\tau _{pp}22`$ (Mészáros mes92:1992 (1992)), if the value of $`\tau _B`$ which follows form equation (B3) exceeds $`\tau _{pp}`$, the latter is used instead.
In the very strong field limit ($`n_{max}1`$), expression (B1) for the heating rate is still valid, provided that $`\tau _s`$ is replaced by $`\tau _B`$. The situation is more complicated in the moderate field limit ($`n_{max}1`$). In this case, not only $`\tau _s`$ must be replaced by $`\tau _B`$, but a further effect must be accounted for. Now a sizeable fraction, $`11/\mathrm{ln}(2n_{max}`$), of the initial proton energy goes into excitations of electrons Landau levels. The cyclotron photons produced by radiative deexcitation will be partly thermalized by absorption and Compton scattering while the remaining ones escape forming a broad cyclotron emission feature (Nelson, et al. nel95:1995 (1995)). To account for this we take in the moderate field limit,
$$W_{MF}=(1f)W_H+fW_H(\mathrm{\Delta }ϵ/ϵ)$$
(B4)
where $`f`$ is the fraction of the initial proton energy which goes into Landau excitations and
$$\frac{\mathrm{\Delta }ϵ}{ϵ}\frac{(\mathrm{}\omega _{c,e}/m_ec^2)\tau ^2}{1+(\mathrm{}\omega _{c,e}/m_ec^2)\tau ^2}$$
(B5)
is the electron fractional energy gain due to repeated scatterings. Strictly speaking, equation (B5) is valid only for photons produced below $`\tau _C(\mathrm{}\omega _{c,e}/m_ec^2)^{1/2}7(B/10^{12}\mathrm{G})^{1/2}`$, because only these photons loose a significant fraction of their energy in Compton collisions with electrons. We have also to take into account that cyclotron photons produced at depths greater than the thermalization depth, $`\tau _{th}16(B/10^{12}\mathrm{G})^{7/6}(M/M_{})^{1/3}(R/10^6\mathrm{cm})^{2/3}(T/10^7\mathrm{K})^{1/3}`$, are absorbed. In our models, however, $`\tau _B`$ is always less than $`\tau _{th}`$, so we do not need to worry about free–free absorption in (B4). A parabolic fit to the sum of first two curves ($`01`$ and $`02`$ Landau transitions) in figure 4 of Nelson, Salpeter, & Wasserman (nel93:1993 (1993)) is very accurate and gives
$$f(\tau )=6.1\times 10^4+6.5\times 10^2(\tau /\tau _s)+8.7\times 10^3(\tau /\tau _s)^2.$$
(B6) |
warning/0002/cond-mat0002106.html | ar5iv | text | # Quantum Criticality at the Metal Insulator Transition
## I Introduction
The Metal-Insulator (M-I) transition has been understood within the seminal paper in 1979. Focusing on noninteracting electrons the authors demonstrated that in two dimension (2D) even weak disorder is sufficient to localize the electrons at $`T=0`$. Few years later it has been realized by Finkelstein that the particle-hole interaction in the triplet channel might enhance the conductivity. However a detailed analysis revealed that at long scale the interaction term diverges making difficult to determine what will happen at long scales. Recently a remarkable experiment has been performed on a 2D electron gas in zero magnetic field strongly points towards a M-I transition in two dimensions. The characteristic of this experiment performed on a 2DES silicon ( $`n_s10^{11}cm^2`$) the mean free path “$`\mathrm{}`$” is large, the electron-electron interaction was $`5mev`$, while the Fermi energy is only $`0.6mev`$. The lowest temperature in the experiment was $`0.2K`$. These experimental condition might suggest that the non-linear sigma model introduced in ref. might not be applicable since it ignores the interaction effects at length scales shorter than the mean free path. Since the mean free path is large quantum effects in the momentum range $`2\pi /\mathrm{}q\mathrm{\Lambda }`$ ($`\mathrm{\Lambda }^1a`$ particle separation) might be important for weak disorder, $`\mathrm{}\mathrm{}`$. This suggests that a phase transition due to a collective many body interaction might occur before the diffusive limit is reached. One might have a phase transition from a superconductor to insulator , Wigner crystal , or quantum Hall-insulator transition . In one dimension it is known that attractive interaction or ferromagnetic spin fluctuations can suppress the $`2k_F`$ backscattering leading to a delocalization transition . We investigate the problem in the presence of interaction and large mean free paths. In order to clarify the situation in 2D we propose to use the Renormalization Group (RG) analysis. Motivated by the fact that the mean free path “$`\mathrm{}`$” can be large with respect to the particle separation $`a\mathrm{\Lambda }^1`$ (standard transport theories start at the scale “$`\mathrm{}`$” and investigate only processes at larger scales governed by diffusion) we investigate at finite temperatures the competition between localization and interaction. The competition between multiple scattering (due to disorder) and the interactions is investigated within a RG theory. The method used here is different from the procedure used in ref.. In ref. one emphasizes the disorder by replacing the multiple elastic scattering by a diffusion theory and in the second step the interactions are treated perturbatively. We consider a situation where the elastic mean free path is much larger than any microscopic length. Therefore we might have a situation that before entering the diffusive region we have to stop scaling. This can happen if the thermal wave length is shorter than the elastic mean free path or that the Cooper channel diverges giving rise to superconductivity. In the quantum region the single particle excitations are well-defined and the Fermi surface is parametrized in terms of $`N_o=\frac{\pi k_F}{\mathrm{\Lambda }}`$ channels. When the cutoff $`\mathrm{\Lambda }`$ is reduced $`\mathrm{\Lambda }\mathrm{\Lambda }/b`$, one finds that the interactions scale like $`\mathrm{\Gamma }\mathrm{\Gamma }b^{1d}`$ and the number of channels, increases like $`N=N_ob`$ . The disorder scales like $`DDb^{2d}`$. Due to the fact that the number of channels increase under scaling, we find that the interaction is marginal and the disorder is relevant. The quantum region gives rise to a set of scaling equations for the interaction term $`\mathrm{\Gamma }`$: $`\mathrm{\Gamma }_2^{(c)}`$–particle-hole singlet, $`\mathrm{\Gamma }_2^{(s)}`$–particle-hole triplet, $`\mathrm{\Gamma }_3^{(s)}`$–particle-particle singlet and disorder $`D`$ ($`d_3^{(s)}`$–the Cooperon). Our results show that due to disorder $`\mathrm{\Gamma }_3^{(s)}`$ might becomes negative resulting in a superconducting instability at $`T0`$. This might give rise to an Insulator-Superconductor transition similar to what one has for superconducting films where a phase transition is expected . In the absence of an instability the standard method at length scale $`b>b_{Dif}`$, $`b_{Dif}=\frac{\mathrm{\Lambda }}{2\pi /\mathrm{}}`$ is the diffusion theory developed by Finkelstein. Here we consider the situation where the system is in the clean limit such that the microscopic mean free path $`\mathrm{}_o=\mathrm{}(b=1)`$ is large. Due to interaction we obtain that the mean free path $`\mathrm{}(b>1)`$ increases, $`\mathrm{}(b)>\mathrm{}_o`$.
In this paper we will work at finite temperatures such the the thermal wavelength is shorter than the mean free path $`\mathrm{}`$. We introduce a thermal length scale $`b_T=\frac{v_F\mathrm{\Lambda }}{T}`$ and consider the situation where $`b_{Dif}>b_T`$. Since we have to stop the scaling scaling at $`b=b_T`$ we are allowed to ignore the diffusive region. In the recent transport experiment $`E_F/T5`$ and $`K_F\mathrm{}5`$, therefore the condition $`b_{Dif}>b_T`$ is realized. The presence of the cutoff $`b_T`$ prevent the number of channels to scale to infinity, instead we have $`N_o<N(b)N(b_T)=\overline{N}=\frac{E_F}{T}`$. We solve the model under the condition $`b_{Dif}>b_T`$ and find that the physics is controlled by the disorder “$`D`$” and the particle-hole triplet $`\mathrm{\Gamma }_2^{(s)}`$. We find that when the number of channels does not scale (This might be the case at finite temperature or for non-spherical Fermi surfaces, which obeys $`N(b)Const.`$), a fixed point in the plane $`\mathrm{\Gamma }_2^{(s)}`$ and $`D`$ is obtained. This fixed point separates a metallic phase from a localized one. The metallic phase is caused by the fact that the particle-hole triplet flows to a stable fixed point causing a shift in the critical dimension from $`d=2`$ to $`d<2`$. The presence of the stable fixed point in the triplet channel causes power law behavior of the spin-spin correlations. The resistivity is expected to obey the scaling behavior: $`\rho (D,\mathrm{\Gamma }_2^{(s)},T)=\rho (D(b),\mathrm{\Gamma }_2^{(s)}(b),Tb^z);z1`$ where $`\mathrm{\Gamma }_2^{(s)}(b)=\mathrm{\Gamma }_2^{}+(\mathrm{\Gamma }_2^{(s)}\mathrm{\Gamma }_2^{})b^{1/\nu _2}`$ and $`D(b)=D^{}+(DD^{})b^{1/\nu _1}`$. Choosing $`Tb^z=T_o`$ we obtain: $`\rho (D,\mathrm{\Gamma }_2^{(s)},T)\rho (D^{},\mathrm{\Gamma }_2^{},T_o)+const.(\frac{DD^{}}{D^{}})(\frac{T_o}{T})^{1/z\nu _1}`$. In agreement with the experimental results given in ref. the resistivity increases for $`D>D^{}`$ and decreases for $`D<D^{}`$. In the literature alternative theories have been proposed already: ref. (phenomenological), ref. (within the Finkelstein theory), as well as models which focus on the insulating side ref..
The plan of this paper is: We introduce in Chapter II our microscopic model. We consider a two dimensional gas in the presence of a screened two-body potential and a static random potential. We follow a standard method for treating disorder. We use the “replica” method and perform the statistical average over the disorder. In the second step we parametrize the Fermi Surface (FS) in terms of $`N`$ channels. Using this parametrization we identify in Appendix A all the possible interaction and disorder terms. We find that the interaction and disorder is best described in terms of chiral currents carrying indices of charge, spin, replica, and channel. In Chapter III the method of the Renormalization Group (RG) based on the Operator Product Expansion (OPE) is introduced. We compute the OPE rules for the different interaction terms, particle-hole (p-h) singlet, p-h triplet, particle-particle (p-p) and the Cooperon (the effective interaction induced by the disorder). Chapter IV is devoted to the derivation of the RG equations based on the OPE results obtained in Chapter III. In Chapter V we consider the scaling equations in the quantum limit. Chapter VI is devoted to the possible superconducting instability which might occur in the quantum region. In Chapter VII we investigate the scaling equations at finite temperatures. Here we observe that the physics is determined by the effective number of channels $`\overline{N}`$. In Chapter VIII we solve the RG equations and compute the resistivity. Chapter IX is limited to discussions and conclusions.
## II The Microscopic Model
We introduce the screened two-body potential and perform a statistical average over the disorder using the replica method. We parametrize the FS in terms of $`N`$ Fermions. Using these Fermions we replace the interaction terms and the Cooperon by chiral currents. The starting point of our investigation is the averaged disorder partition function, $`\overline{Z^\alpha },\alpha =1,\mathrm{},\alpha 0`$,
$$\overline{Z^\alpha }=D[\overline{\psi },\psi ]e^S,\alpha =1,\mathrm{},\alpha 0$$
(1)
$$S_o=d^dx𝑑t\{\underset{\sigma }{}\underset{\alpha }{}[\overline{\psi }_{\sigma ,\alpha }_t\psi _{\sigma ,\alpha }\overline{\psi }_{\sigma ,\alpha }(\frac{^2}{2m}+E_F)\psi _{\sigma ,\alpha }]\}$$
(2)
$$S_{int}=d^dxd^dy𝑑t\underset{\sigma ,\sigma ^{}}{}\underset{\alpha }{}\{\overline{\psi }_{\sigma ,\alpha }(x)\overline{\psi }_{\sigma ^{},\alpha }(y)v(xy)\psi _{\sigma ^{},\alpha }(y)\psi _{\sigma ,\alpha }(x)\}$$
(3)
$$S_D=𝑑t_1𝑑t_2d^dxd^dy\underset{\sigma ,\sigma ^{}}{}\underset{\alpha ,\beta }{}\{\overline{V(x)V(y)}\overline{\psi }_{\sigma ,\alpha }(x,t_1)\overline{\psi }_{\sigma ^{},\beta }(y,t_2)\psi _{\sigma ^{},\beta }(y,t_2)\psi _{\sigma ,\alpha }(x,t_1)\}$$
(4)
$`v(xy)`$” is the two body screened potential and $`\overline{V(x)V(y)}=D\delta (xy)`$ where $`D=\frac{v_F^2}{K_F\mathrm{}}`$ is the disorder parameter controlled by the elastic scattering time $`\tau =\mathrm{}/v_F`$. Next we parametrize the Fermi surface (FS) in terms of $`N`$ Fermions or $`N/2`$ pairs of right and left movers ( see ref. ):
$$\psi _{\sigma ,\alpha }(\stackrel{}{x})=\underset{n=1}{\overset{N/2}{}}(e^{ik_F\widehat{n}\stackrel{}{x}}R_{n,\sigma ,\alpha }(\stackrel{}{x})+e^{ik_F\widehat{n}\stackrel{}{x}}L_{n,\sigma ,\alpha }(\stackrel{}{x}))$$
(5)
$`R_{n,\sigma ,\alpha }(\stackrel{}{x})`$ and $`L_{n,\sigma ,\alpha }(\stackrel{}{x})`$ are right and left movers defined by momenta $`q_{}<\mathrm{\Lambda }`$, $`q_{}<\mathrm{\Lambda }`$ around each Fermi point $`k_F=k_F\widehat{n}`$. The Fermi momentum is determined by the renormalized Fermi energy $`\overline{E}_F`$ which is related to the non-interacting Fermi energy $`E_F`$ by the relation $`\overline{E}_F=E_F+\delta \mu _F`$, such that $`\overline{E}_F=\frac{k_F^2}{2m^{}}`$. The value of $`\delta \mu _F`$ is obtained from the interaction. The two dimensional Fermions are expressed in terms of the one dimensional Fermions $`\widehat{R}_{n,\sigma ,\alpha }(x_{})`$ and $`\widehat{L}_{n,\sigma ,\alpha }(x_{})`$:
$`R_{n,\sigma ,\alpha }(\stackrel{}{x})=\widehat{R}_{n,\sigma ,\alpha }(x_{})Z_n(x_{}),L_{n,\sigma ,\alpha }(\stackrel{}{x})=\widehat{L}_{n,\sigma ,\alpha }(x_{})Z_n(x_{})`$
$`Z_n(x_{})`$ is scalar function which ensures the conservation of momentum in the transversal direction. The number of channels (Fermions) is related to $`k_F`$ and cutoff $`\mathrm{\Lambda }<k_F`$, $`N_o=\frac{\pi k_F}{\mathrm{\Lambda }}`$. Using the representation given in Eq.5, we introduce the normal order currents $`J_{n,\alpha ,\sigma }^R(Z)`$ (right mover) and $`J_{n,\alpha ,\sigma }^L(\overline{Z})`$ (left mover) with $`Z`$ and $`\overline{Z}`$ given by $`Z=(Z_{},Z_{})`$, $`\overline{Z}=(\overline{Z}_{},\overline{Z}_{})`$, $`Z_{}=v_Ftix_{}`$, $`\overline{Z}_{}=v_Ft+ix_{}`$, and $`Z_{}=\overline{Z}_{}=x_{}`$,
$$J_{n,\alpha ,\sigma }^R(Z)=:R_{n,\alpha ,\sigma }^{}(Z)R_{n,\alpha ,\sigma }(Z):R_{n,\alpha ,\sigma }^{}(Z+ϵ)R_{n,\alpha ,\sigma }(Z)R_{n,\alpha ,\sigma }^{}(Z+ϵ)R_{n,\alpha ,\sigma }(Z)_o$$
(6)
with $`ϵ=\epsilon _xi\delta `$, $`ϵ0`$ and the expectation value:
$$R_{n,\alpha ,\sigma _1}^{}(\stackrel{}{x},t_1)R_{m,\beta ,\sigma _2}(\stackrel{}{y},t_2)_o\delta _{n,m}\delta _{\alpha ,\beta }\delta _{\sigma _1,\sigma _2}\delta _\mathrm{\Lambda }^{d1}(x_{}y_{})[v_F(t_1t_2)i(x_{}y_{})]^1$$
(7)
Similarly we introduce for the left movers:
$$J_{n,\alpha ,\sigma }^L(\overline{Z})=:L_{n,\alpha ,\sigma }^{}(\overline{Z})L_{n,\alpha ,\sigma }(\overline{Z}):L_{n,\alpha ,\sigma }^{}(\overline{Z}+\overline{ϵ})L_{n,\alpha ,\sigma }(\overline{Z})L_{n,\alpha ,\sigma }^{}(\overline{Z}+\overline{ϵ})L_{n,\alpha ,\sigma }(\overline{Z})_o$$
(8)
We write the interaction and the disorder parts in the normal order form. From the disorder part we obtain the elastic scattering term $`\frac{1}{2\tau }D`$ (see ref.). From the disorder part (Eq.4) we obtain the normal order form $`\stackrel{~}{S}_D`$.
From the interaction part we find the normal order representation $`\stackrel{~}{S}_{int}`$ plus a shift of the Fermi energy: $`\delta \mu _{int}(J_{n,\alpha ,\sigma }^R(\stackrel{}{x},t)+J_{n,\alpha ,\sigma }^L(\stackrel{}{x},t))`$. We choose $`\delta \mu _F`$ such that it cancels the interaction shift, $`\delta \mu _F+\delta \mu _{int}=0`$. As a result $`S_0`$ becomes:
$$\stackrel{~}{S}_o=\underset{n=1}{\overset{N/2}{}}\underset{\sigma }{}\underset{\alpha }{}d^dx𝑑t\{\overline{R}_{n,\alpha ,\sigma }[_tv_F\widehat{n}\stackrel{}{}]R_{n,\alpha ,\sigma }+\overline{L}_{n,\alpha ,\sigma }[_t+v_F\widehat{n}\stackrel{}{}]L_{n,\alpha ,\sigma }\}$$
(9)
Using the representation given in Eq.5 we replace the interaction term and disorder in terms of the currents ( see appendix A ). The interaction part is decomposed in terms of forward scattering $`Q_{n,m}^{(F)}(t,\stackrel{}{x},\stackrel{}{y})`$ (charge part) $`H_{n,m}^{(F)}(t,\stackrel{}{x},\stackrel{}{y})`$ (spin part), $`Q_{n,m}^{(B)}(t,\stackrel{}{x},\stackrel{}{y})`$ (particle-hole in the singlet channel), $`H_{n,m}^{(B)}(t,\stackrel{}{x},\stackrel{}{y})`$ (particle-hole in the triplet channel), $`O_{n,m}^{(s)}(t,\stackrel{}{x},\stackrel{}{y})`$ (particle-particle in the singlet channel), $`O_{n,m}^{(t)}(t,\stackrel{}{x},\stackrel{}{y})`$ (particle-particle in the triplet channel). From the screened two-body potential $`v(\stackrel{}{q})`$ we obtain the scattering matrix elements for the different processes, $`\mathrm{\Gamma }^{(c)}(\stackrel{}{n},\stackrel{}{m})`$,$`\mathrm{\Gamma }^{(s)}(\stackrel{}{n},\stackrel{}{m})`$, $`\mathrm{\Gamma }_2^{(c)}(\stackrel{}{n},\stackrel{}{m})`$,$`\mathrm{\Gamma }_2^{(s)}(\stackrel{}{n},\stackrel{}{m})`$, $`\mathrm{\Gamma }_3^{(s)}(\stackrel{}{n},\stackrel{}{m})`$,$`\mathrm{\Gamma }_3^{(t)}(\stackrel{}{n},\stackrel{}{m})`$. For the screened case the matrix elements $`\mathrm{\Gamma }(\stackrel{}{n},\stackrel{}{m})`$ depend only on the angles “$`\theta `$” on the FS. For example, if $`\kappa `$ is the inverse of the screening length we have $`\mathrm{\Gamma }_2^{(s)}(\stackrel{}{n},\stackrel{}{m})=2\kappa [1+\frac{2k_F}{\kappa }\mathrm{cos}\theta /2]^1,\mathrm{\hspace{0.33em}\hspace{0.33em}0}\theta \pi `$ (“$`\theta `$” is the angle between the unit vectors $`\stackrel{}{n}`$ and $`\stackrel{}{m}`$). The particle-particle matrix is, $`\mathrm{\Gamma }_3^{(s)}(\stackrel{}{n},\stackrel{}{m})=\frac{\kappa }{2}[(1+\frac{2k_F}{\kappa }\mathrm{sin}\theta /2)^1+(1+\frac{2k_F}{\kappa }\mathrm{cos}\theta /2)^1],\mathrm{\hspace{0.33em}\hspace{0.33em}0}\theta \pi `$. We introduce the left and right currents and obtain the representations for the interaction and disorder terms:
$`J_{n,\alpha ,\sigma _1;m,\beta ,\sigma _2}^R(\stackrel{}{x},t_1,t_2)=:R_{n,\alpha ,\sigma _1}^{}(\stackrel{}{x},t_1)R_{m,\beta ,\sigma _2}(\stackrel{}{x},t_2):`$
$$J_{n,\alpha ,\sigma _1;m,\beta ,\sigma _2}^L(\stackrel{}{x},t_1,t_2)=:L_{n,\alpha ,\sigma _1}^{}(\stackrel{}{x},t_1)L_{m,\beta ,\sigma _2}(\stackrel{}{x},t_2):$$
(10)
For the interaction term we have $`t_1=t_2`$ and $`\alpha =\beta `$. We obtain that the interaction part for a screened two-body potential takes the form:
$`\stackrel{~}{S}_{int}={\displaystyle \frac{\mathrm{\Lambda }^{1d}}{2N_o}}{\displaystyle \underset{n}{}}{\displaystyle \underset{m}{}}{\displaystyle }d^dx{\displaystyle }dt{\displaystyle \underset{\alpha }{}}\{\mathrm{\Gamma }^{(c)}(\stackrel{}{n},\stackrel{}{m})Q_{n,m;\alpha }^{(F)}(\stackrel{}{x},t)\mathrm{\Gamma }^{(s)}(\stackrel{}{n},\stackrel{}{m})H_{n,m;\alpha }^{(F)}(\stackrel{}{x},t)`$
$$+\mathrm{\Gamma }_2^{(c)}(\stackrel{}{n},\stackrel{}{m})Q_{n,m;\alpha }^{(B)}(\stackrel{}{x},t)\mathrm{\Gamma }_2^{(s)}(\stackrel{}{n},\stackrel{}{m})H_{n,m;\alpha }^{(B)}(\stackrel{}{x},t)+\mathrm{\Gamma }_3^{(s)}(\stackrel{}{n},\stackrel{}{m})O_{n,m;\alpha }^{(s)}(\stackrel{}{x},t)+\mathrm{\Gamma }_3^{(t)}(\stackrel{}{n},\stackrel{}{m})O_{n,m;\alpha }^{(t)}(\stackrel{}{x},t)\}$$
(11)
In Eq.11 we have to restrict $`\mathrm{\Gamma }_3^{(s)}(\stackrel{}{n},\stackrel{}{m})`$ and $`\mathrm{\Gamma }_3^{(t)}(\stackrel{}{n},\stackrel{}{m})`$ to $`\stackrel{}{n}\stackrel{}{m}`$ in order to avoid double counting. If we ignore the angle dependence of $`\mathrm{\Gamma }_3^{(t)}`$ we have $`\mathrm{\Gamma }_3^{(t)}0`$. For $`\stackrel{}{n}=\stackrel{}{m}`$ we have the relation $`\mathrm{\Gamma }_3^{(s)}(\stackrel{}{n},\stackrel{}{n})=\frac{1}{2}\mathrm{\Gamma }_2^{(s)}(\stackrel{}{n},\stackrel{}{n})`$. Based on dimensional analysis we obtain that the interaction term has the dimension of $`\mathrm{\Lambda }^{1d}`$. Due to the fact that the interaction is defined at the scale $`\mathrm{\Lambda }<k_F`$ we have the relation $`k_F^{1d}\mathrm{\Gamma }(k_F)=\mathrm{\Lambda }^{1d}(\frac{k_F}{\mathrm{\Lambda }})^{d1}\mathrm{\Gamma }(\mathrm{\Lambda })\frac{\mathrm{\Lambda }^{1d}}{N_o}\mathrm{\Gamma }(\mathrm{\Lambda })`$ where $`N_o=\pi (\frac{k_F}{\mathrm{\Lambda }})^{d1}`$ is the number of channels. The operators $`Q_{n,m;\alpha }^{(F)}`$, $`H_{n,m;\alpha }^{(F)}`$, $`Q_{n,m;\alpha }^{(B)}`$, $`H_{n,m;\alpha }^{(B)}`$, $`O_{n,m;\alpha }^{(s)}`$, and $`O_{n,m;\alpha }^{(t)}`$ are given by:
$`Q_{n,m;\alpha }^{(F)}(\stackrel{}{x},t)=J_{n,\alpha }^R(\stackrel{}{x},t)J_{m,\alpha }^R(\stackrel{}{x},t)+J_{n,\alpha }^L(\stackrel{}{x},t)J_{m,\alpha }^L(\stackrel{}{x},t),`$
$`H_{n,m;\alpha }^{(F)}(\stackrel{}{x},t)=\stackrel{}{J}_{n,\alpha }^R(\stackrel{}{x},t)\stackrel{}{J}_{m,\alpha }^R(\stackrel{}{x},t)+\stackrel{}{J}_{n,\alpha }^L(\stackrel{}{x},t)\stackrel{}{J}_{m,\alpha }^L(\stackrel{}{x},t),`$
$`Q_{n,m;\alpha }^{(B)}(\stackrel{}{x},t)=J_{n,\alpha }^R(\stackrel{}{x},t)J_{m,\alpha }^L(\stackrel{}{x},t)+J_{n,\alpha }^L(\stackrel{}{x},t)J_{m,\alpha }^R(\stackrel{}{x},t),`$
$`H_{n,m;\alpha }^{(B)}(\stackrel{}{x},t)=\stackrel{}{J}_{n,\alpha }^R(\stackrel{}{x},t)\stackrel{}{J}_{m,\alpha }^L(\stackrel{}{x},t)+\stackrel{}{J}_{n,\alpha }^L(\stackrel{}{x},t)\stackrel{}{J}_{m,\alpha }^R(\stackrel{}{x},t),`$
$`O_{n,m;\alpha }^{(s)}(\stackrel{}{x},t)=O_{n,m;\alpha }^{()}(\stackrel{}{x},t)O_{n,m;\alpha }^{()}(\stackrel{}{x},t),O_{n,m;\alpha }^{(t)}(\stackrel{}{x},t)=O_{n,m;\alpha }^{()}(\stackrel{}{x},t)+O_{n,m;\alpha }^{()}(\stackrel{}{x},t),`$
$`O_{n,m;\alpha }^{()}(\stackrel{}{x},t)={\displaystyle \underset{\sigma }{}}:R_{n,\alpha ,\sigma }^{}(\stackrel{}{x},t)R_{m,\alpha ,\sigma }(\stackrel{}{x},t)::L_{n,\alpha ,\sigma }^{}(\stackrel{}{x},t)L_{m,\alpha ,\sigma }(\stackrel{}{x},t):`$
$`+:L_{n,\alpha ,\sigma }^{}(\stackrel{}{x},t)L_{m,\alpha ,\sigma }(\stackrel{}{x},t)::R_{n,\alpha ,\sigma }^{}(\stackrel{}{x},t)R_{m,\alpha ,\sigma }(\stackrel{}{x},t):,`$
$`O_{n,m;\alpha }^{()}(\stackrel{}{x},t)={\displaystyle \underset{\sigma }{}}:R_{n,\alpha ,\sigma }^{}(\stackrel{}{x},t)R_{m,\alpha ,\sigma }(\stackrel{}{x},t)::L_{n,\alpha ,\sigma }^{}(\stackrel{}{x},t)L_{m,\alpha ,\sigma }(\stackrel{}{x},t):`$
$$+:L_{n,\alpha ,\sigma }^{}(\stackrel{}{x},t)L_{m,\alpha ,\sigma }(\stackrel{}{x},t)::R_{n,\alpha ,\sigma }^{}(\stackrel{}{x},t)R_{m,\alpha ,\sigma }(\stackrel{}{x},t):,$$
(12)
where
$$J_{n,\alpha }^R(\stackrel{}{x},t)=\underset{\sigma }{}:R_{n,\alpha ,\sigma }^{}(\stackrel{}{x},t)R_{n,\alpha ,\sigma }(\stackrel{}{x},t):,\stackrel{}{J}_{n,\alpha }^R(\stackrel{}{x},t)=\frac{1}{2}:R_{n,\alpha ,\sigma _1}^{}(\stackrel{}{x},t)\stackrel{}{\sigma }_{\sigma _1,\sigma _2}R_{n,\alpha ,\sigma _2}(\stackrel{}{x},t):$$
(13)
with similar expressions for the left movers. $`\mathrm{`}\mathrm{`}\mathrm{\Lambda }\mathrm{"}<k_F`$ is the cutoff of the theory and we find that the naive dimension of the interaction field is $`\mathrm{\Lambda }^{1d}(d=2)`$. This follows from the fact that Eq.9 is invariant under the scaling $`\mathrm{\Lambda }\mathrm{\Lambda }/b`$, $`x=x^{}b`$, $`t=t^{}b`$, $`R_n(\stackrel{}{x},t)=b^{d/2}R_n(\stackrel{}{x}^{},t^{})`$, $`L_n(\stackrel{}{x},t)=b^{d/2}L_n(\stackrel{}{x}^{},t^{})`$, and $`N(b)=bN_o`$. Following the same procedure as for the interaction we express the disorder part using again the respective part:
$`\stackrel{~}{S}_D={\displaystyle \frac{\mathrm{\Lambda }^{2d}}{N_o}}{\displaystyle \underset{n}{}}{\displaystyle \underset{m}{}}{\displaystyle \underset{\alpha ,\beta }{}}{\displaystyle }dt_1{\displaystyle }dt_2{\displaystyle }d^dx\{d_2^{(d)}\rho _{n,m,\alpha ,\beta }(\stackrel{}{x};t_1,t_2)d_2^{(c)}q_{n,m,\alpha ,\beta }(\stackrel{}{x};t_1,t_2)`$
$$d_2^{(s)}h_{n,m,\alpha ,\beta }^{(B)}(\stackrel{}{x};t_1,t_2)+d_3^{(s)}c_{n,m,\alpha ,\beta }^{(s)}(\stackrel{}{x};t_1,t_2)+d_3^{(t)}c_{n,m,\alpha ,\beta }^{(t)}(\stackrel{}{x};t_1,t_2)$$
(14)
The operators in Eq.14 are in complete analogy with the ones in Eq.11, except that they are at different times and have double replica index:
$`\rho _{n,m,\alpha ,\beta }(\stackrel{}{x};t_1,t_2)Q_{n,m,\alpha }^{(F)}(\stackrel{}{x},t);`$
$`q_{n,m,\alpha ,\beta }(\stackrel{}{x};t_1,t_2)Q_{n,m,\alpha }^{(B)}(\stackrel{}{x},t);`$
$`h_{n,m,\alpha ,\beta }^{(B)}(\stackrel{}{x};t_1,t_2)H_{n,m,\alpha }^{(B)}(\stackrel{}{x},t);`$
$`C_{n,m,\alpha ,\beta }^{(s)}(\stackrel{}{x};t_1,t_2)O_{n,m,\alpha }^{(s)}(\stackrel{}{x},t);`$
$`C_{n,m,\alpha ,\beta }^{(t)}(\stackrel{}{x};t_1,t_2)O_{n,m,\alpha }^{(t)}(\stackrel{}{x},t)`$
The corresponding constants in Eq.14 have the initial values: $`d_3^{(s)}=d_2^{(c)}=\frac{1}{2}d_2^{(s)}D`$, $`d_3^{(t)}=0`$. We will find that only the Cooperon term, $`d_3^{(s)}c_{n,m,\alpha ,\beta }^{(s)}(\stackrel{}{x};t_1,t_2)`$ is important. For the rest part of this paper we will ignore the rest of the terms and consider only the Cooperon part.
## III The Renormalization Group Method
In the first part of this chapter we will introduce the RG method based on the OPE. This method is needed in order to analyze the possible phase diagram of our problem. The real space method based on the Operator Product Expansion (OPE) introduced in ref. is in particular advantageous. In order to explain how this works we express the action in Eq.11 by a formal expression $`S\mathrm{\Gamma }_iA_i`$ where $`A_i`$ are the operators and $`\mathrm{\Gamma }_i`$ are the coupling constants. Using the fact that the time ordered product of the single particle operator is given by,
$`R_{n,\alpha ,\sigma }(\stackrel{}{x},t_1)R_{m,\beta ,\sigma _1}^{}(\stackrel{}{y},t_2){\displaystyle \frac{1}{2\pi }}\delta _{n,m}\delta _{\alpha ,\beta }\delta _{\sigma ,\sigma _1}\delta _\mathrm{\Lambda }^{d1}(x_{}y_{})\theta (t_1t_2)`$
$$[v_F(t_1t_2)i(x_{}y_{})]^1$$
(15)
where $`x_{}=\widehat{n}\stackrel{}{x}`$, $`x_{}=\stackrel{}{x}\widehat{n}\stackrel{}{x}`$. We find for any two operators given in Eq.11 the OPE:
$$A_i(\stackrel{}{x},t_1)A_j(\stackrel{}{x}+a,t_2)\underset{K}{}\frac{C_{ij}^KF_K(t_1t_2)A_K(\stackrel{}{x},\frac{t_1+t_2}{2})}{[a^2+v_F^2(t_1t_2)^2]^{x_i+x_jx_K}}$$
(16)
with $`C_{ij}^K`$ the structure constant and $`F_K(t_1t_2)1`$. As a result the product of any number of operators can be reduced to a sum of operators. This implies that once the cutoff $`\mathrm{\Lambda }`$ is reduced to $`\mathrm{\Lambda }/b`$, one can obtain the scaling equations for coupling constants $`\mathrm{\Gamma }_i`$. For $`\mathrm{\Gamma }_i`$ with the scaling dimension $`\mathrm{\Gamma }_i\mathrm{\Gamma }_ib^{(x_id)}`$, one obtains:
$$\frac{d\mathrm{\Gamma }_K}{d\mathrm{ln}b}=(dx_K)\mathrm{\Gamma }_K\frac{1}{2}\underset{i,j}{}\stackrel{~}{C}_{i,j}^K\mathrm{\Gamma }_i\mathrm{\Gamma }_j+\frac{1}{3!}\underset{i,j}{}\underset{p,q}{}\stackrel{~}{C}_{i,j}^p\stackrel{~}{C}_{p,q}^K\mathrm{\Gamma }_i\mathrm{\Gamma }_j\mathrm{\Gamma }_q$$
(17)
where the $`\stackrel{~}{C}_{i,j}^K`$ are proportional to the structure constants $`C_{i,j}^K`$. In order to be able to complete the RG equation given in Eq.17 we have to compute the operator product expansion of the operators which appear in Eqs.11 and 14. The second part of this chapter will be devoted to the calculation of the OPE for the interaction and disorder operators. Using current algebra of the chiral currents given in ref. we will establish the OPE rules for our problem. The calculation is based on the Wick theorem which replaces the time order product by the normal ordered form plus all the possible ways of contracting pairs of Fermion fields. This calculation is standard and lengthy therefore we will present only the results. We start with the results for the p-p singlet:
$`O_{n,m,\alpha }^{(s)}(\stackrel{}{x},t)O_{k,l,\beta }^{(s)}(\stackrel{}{x}+\stackrel{}{a},t+\tau )={\displaystyle \frac{1}{(2\pi )^2}}({\displaystyle \frac{\mathrm{\Lambda }}{2\pi }})^{d1}\delta _{\alpha ,\beta }[a^2+(v_F\tau )^2]^1\{O_{k,m,\alpha }^{(s)}(\stackrel{}{x},t)\delta _{n,l}`$
$$+O_{n,l,\alpha }^{(s)}(\stackrel{}{x},t)\delta _{k,m}\delta _{n,l}\delta _{k,m}(Q_{n,m,\alpha }^{(B)}(\stackrel{}{x},t)+Q_{m,n,\alpha }^{(B)}(\stackrel{}{x},t))\}+[\mathrm{`}\mathrm{`}c\mathrm{"}number].$$
(18)
From Eq.18 we learn that the OPE generates p-p and p-h singlets.
For the p-h in the triplet channel no new terms are generated:
$`H_{n,m,\alpha }^{(B)}(\stackrel{}{x},t)H_{k,l,\beta }^{(B)}(\stackrel{}{x}+\stackrel{}{a},t+\tau )={\displaystyle \frac{1}{(2\pi )^2}}({\displaystyle \frac{\mathrm{\Lambda }}{2\pi }})^{d1}\delta _{\alpha ,\beta }[a^2+(v_F\tau )^2]^1`$
$$\{2H_{n,m,\alpha }^{(B)}(\stackrel{}{x},t)[\delta _{n,l}\delta _{k,m}+\delta _{n,k}\delta _{l,m}]\}+[\mathrm{`}\mathrm{`}c\mathrm{"}number]$$
(19)
The p-h singlet generates only a “c” number:
$$Q_{n,m,\alpha }^{(B)}(\stackrel{}{x},t)Q_{k,l,\beta }^{(B)}(\stackrel{}{x}+\stackrel{}{a},t+\tau )=[\mathrm{`}\mathrm{`}c\mathrm{"}number]$$
(20)
The OPE for the Cooperon do not generate new terms:
$`C_{n,m;\alpha ,\beta }^{(s)}(\stackrel{}{x};t_1,t_2)C_{k,l;\alpha ^{}\beta ^{}}^{(s)}(\stackrel{}{x}+\stackrel{}{a};t_1+\tau _1,t_2+\tau _2)={\displaystyle \frac{1}{(2\pi )^2}}({\displaystyle \frac{\mathrm{\Lambda }}{2\pi }})^{d1}[{\displaystyle \frac{1/2}{(v_F\tau _1ia)(v_F\tau _2+ia)}}`$
$`+{\displaystyle \frac{1/2}{(v_F\tau _1+ia)(v_F\tau _2ia)}}]\{2\delta _{\alpha ,\beta ^{}}\delta _{\alpha ^{},\beta }[\delta _{m,k}C_{n,l;\alpha ,\beta }^{(s)}(\stackrel{}{x};t_1,t_2)+\delta _{n,l}C_{k,m;\alpha ,\beta }^{(s)}(\stackrel{}{x};t_1,t_2)]`$
$$\frac{1}{2}\delta _{\alpha ,\alpha ^{}}\delta _{\beta ,\beta ^{}}\delta _{m,k}\delta _{n,l}[C_{n,m;\alpha ,\beta }^{(s)}(\stackrel{}{x};t_1,t_2)+C_{m,n;\alpha ,\beta }^{(s)}(\stackrel{}{x};t_1,t_2)]\}+[\mathrm{`}\mathrm{`}c\mathrm{"}number]$$
(21)
The OPE between the p-p and p-h triplet generates the p-p operator and the p-h singlet:
$`O_{n,m;\alpha }^{(s)}(\stackrel{}{x},t)H_{k,l;\beta }^{(B)}(\stackrel{}{x}+\stackrel{}{a},t+\tau )={\displaystyle \frac{1}{(2\pi )^2}}({\displaystyle \frac{\mathrm{\Lambda }}{2\pi }})^{d1}\delta _{\alpha ,\beta }[a+(v_F\tau )^2]^1\{{\displaystyle \frac{3}{4}}\delta _{l,k}\delta _{n,l}\delta _{m,k}Q_{n,m}^{(B)})\stackrel{}{x},t)]\}`$
$$\frac{\delta _{l,k}}{8}[\delta _{n,l}O_{l,m;\alpha }^{(s)}(\stackrel{}{x},t)+\delta _{m,l}O_{n,l;\alpha }^{(s)}(\stackrel{}{x},t)]\}+[\mathrm{`}\mathrm{`}c\mathrm{"}number]$$
(22)
The OPE between the p-p term and the p-h singlet generates a “c” number:
$$O_{n,m;\alpha }^{(B)}(\stackrel{}{x},t)Q_{k,l;\beta }^{(B)}(\stackrel{}{x}+\stackrel{}{a},t+\tau )=[\mathrm{`}\mathrm{`}c\mathrm{"}number]$$
(23)
The product for the product p-h triplet and p-h singlet gives a “c” number:
$$H_{n,m;\alpha }^{(B)}(\stackrel{}{x},t)Q_{k,l;\beta }^{(B)}(\stackrel{}{x}+\stackrel{}{a},t+\tau )=[\mathrm{`}\mathrm{`}c\mathrm{"}number]$$
(24)
In the remaining part we present the OPE between the Cooperon and the interaction operators. For the p-p case we generate the Cooperon and p-p operator:
$`O_{n,m;\gamma }^{(s)}(\stackrel{}{x},t)C_{k,l;\alpha ,\beta }^{(s)}(\stackrel{}{x}+\stackrel{}{a},t_1+\tau _1,t_2+\tau _2)={\displaystyle \frac{1}{(2\pi )^2}}({\displaystyle \frac{\mathrm{\Lambda }}{2\pi }})^{d1}\{[{\displaystyle \frac{1/2}{(v_F\tau _1ia)(v_F\tau _2+ia)}}+`$
$`{\displaystyle \frac{1/2}{(v_F\tau _1+ia)(v_F\tau _2ia)}}]\delta _{\gamma ,\alpha }\delta _{\gamma ,\beta }[O_{k,m;\gamma }^{(s)}(\stackrel{}{x},t)\delta _{n,l}+O_{n,l;\gamma }^{(s)}(\stackrel{}{x},t)\delta _{k,m}]+{\displaystyle \frac{1/2}{(v_F\tau _1)^2+a^2}}\delta _{k,m}\delta _{n,l}\delta _{n,m}(\delta _{\gamma ,\alpha }+\delta _{\gamma ,\beta })`$
$$C_{n,m;\alpha ,\beta }^{(s)}(\stackrel{}{x},t,t+\tau _2)+\frac{1/2}{(v_F\tau _2)^2+a^2}\delta _{k,m}\delta _{n,l}\delta _{n,m}(\delta _{\gamma ,\alpha }+\delta _{\gamma ,\beta })C_{n,m;\alpha ,\beta }^{(s)}(\stackrel{}{x},t+\tau _1,t)\}+[\mathrm{`}\mathrm{`}c\mathrm{"}number]$$
(25)
For the p-h triplet one obtains the p-p and Cooperon terms:
$`H_{n,m;\gamma }^{(B)}(\stackrel{}{x},t)C_{k,l;\alpha ,\beta }^{(s)}(\stackrel{}{x}+\stackrel{}{a},t_1+\tau _1,t_2+\tau _2)={\displaystyle \frac{1}{(2\pi )^2}}({\displaystyle \frac{\mathrm{\Lambda }}{2\pi }})^{d1}\{{\displaystyle \frac{1}{2}}[{\displaystyle \frac{1}{a^2+(v_F\tau _1)^2}}C_{n,m;\alpha ,\beta }^{(s)}(\stackrel{}{x},t,t+\tau _2)`$
$`+{\displaystyle \frac{1}{a^2+(v_F\tau _2)^2}}C_{n,m;\alpha ,\beta }^{(s)}(\stackrel{}{x},t+\tau _1,t)][\delta _{k,m}\delta _{n,l}\delta _{n,m}(\delta _{\gamma ,\alpha }+\delta _{\gamma ,\beta }){\displaystyle \frac{3}{4}}][{\displaystyle \frac{1/2}{(v_F\tau _1ia)(v_F\tau _2+ia)}}`$
$$+\frac{1/2}{(v_F\tau _1+ia)(v_F\tau _2ia)}][\delta _{m,k}\delta _{m,l}+\delta _{n,k}\delta _{n,l}][\delta _{\gamma ,\alpha }\delta _{\gamma ,\beta }][\frac{1}{4}O_{n,m;\gamma }^{(s)}(\stackrel{}{x},t)+\frac{1}{4}O_{n,m;\gamma }^{(t)}(\stackrel{}{x},t)]\}+[\mathrm{`}\mathrm{`}c\mathrm{"}number]$$
(26)
When we consider the p-h singlet we generate the p-p and Cooperon terms.
$`Q_{n,m;\gamma }^{(B)}(\stackrel{}{x},t)C_{k,l;\alpha ,\beta }^{(s)}(\stackrel{}{x}+\stackrel{}{a},t_1+\tau _1,t_2+\tau _2)={\displaystyle \frac{1}{(2\pi )^2}}({\displaystyle \frac{\mathrm{\Lambda }}{2\pi }})^{d1}\{{\displaystyle \frac{1}{2}}[{\displaystyle \frac{1}{a^2+(v_F\tau _1)^2}}(C_{n,m;\alpha ,\beta }^{(s)}(\stackrel{}{x},t,t+\tau _2)`$
$`C_{n,m;\alpha ,\beta }^{(t)}(\stackrel{}{x},t,t+\tau _2)]))+{\displaystyle \frac{1}{a^2+(v_F\tau _2)^2}}(C_{n,m;\alpha ,\beta }^{(s)}(\stackrel{}{x},t+\tau _1,t)C_{n,m;\alpha ,\beta }^{(t)}(\stackrel{}{x},t+\tau _1,t))]`$
$`{\displaystyle \frac{1}{2}}(\delta _{n,k}\delta _{m,l}+\delta _{m,k}\delta _{n,l})(\delta \gamma ,\alpha +\delta _{\gamma ,\beta })+[{\displaystyle \frac{1/2}{(v_F\tau _1ia)(v_F\tau _2+ia)}}+{\displaystyle \frac{1/2}{(v_F\tau _1+ia)(v_F\tau _2ia)}}]`$
$$\delta _{\gamma ,\alpha }\delta _{\gamma ,\beta }(\delta _{m,k}\delta _{m,l}+\delta _{n,k}\delta _{n,l})(O_{n,m;\gamma }^{(s)}(\stackrel{}{x},t)O_{n,m;\gamma }^{(t)}(\stackrel{}{x},t))\}+[\mathrm{`}\mathrm{`}c\mathrm{"}number]$$
(27)
The Eqs.18-27 have been obtained using the free Fermion action given in Eq.9. We will work at a finite temperature, therefore Eqs.15 is an approximation of the exact propagator $`\{\frac{v_F\beta }{\pi }\mathrm{sin}[\frac{\pi }{v_F\beta }(v_F(t_1t_2)i(x_{}y_{}))]\}^1`$. This means that the “$`t`$” range of integration in Eqs.17 and 27 is restricted to $`t<\beta `$.
## IV Derivation of the RG equations
This chapter is designated to the computation of the RG equations. This will be done by expanding the partition function $`Z`$ in terms of the interaction and disorder operators. Using the OPE rules derived in Eqs.18-27 will allow to replace the product of operators in terms of a sum of operators. When rescaling the minimal distance “a” to “ba” will allow to find the scaling equations. It is important to remark that the method used here is different from the standard method used for problems with disorder. The traditional method starts from the diffusion theory and includes the interaction terms as a perturbation. Here we start from the Fermion theory and include simultaneously on equal footing the effects of interaction and disorder. In the standard approach the quantum diffusion theory ignores completely the effects of interactions at short distances (distances shorter than the mean free path). We will see that considering the disorder and interaction on equal footing new terms will appear in the RG equations. The scaling equations will contain terms which are controlled by the number of channels.
Following the analysis given in section II we have:
$$\stackrel{~}{S}_o=\underset{n=1}{\overset{N/2}{}}\underset{\sigma }{}\underset{\alpha }{}d^dx𝑑t\{\overline{R}_{n,\alpha ,\sigma }[_tv_F\widehat{n}\stackrel{}{}]R_{n,\alpha ,\sigma }+\overline{L}_{n,\alpha ,\sigma }[_t+v_F\widehat{n}\stackrel{}{}]L_{n,\alpha ,\sigma }\}.$$
(28)
The action $`\stackrel{~}{S}_o`$ determines the partition function $`Z_o`$,
$`Z_o={\displaystyle D[\overline{\psi },\psi ]e^{\stackrel{~}{S}_o}}.`$
We perturb the partition function $`Z_o`$ by the interaction $`\stackrel{~}{S}_{int}`$ and disorder $`\stackrel{~}{S}_D`$:
$`\stackrel{~}{S}_{int}={\displaystyle \frac{\mathrm{\Lambda }^{1d}}{2N_o}}{\displaystyle \underset{n}{}}{\displaystyle \underset{m}{}}{\displaystyle \underset{\alpha }{}}{\displaystyle }d^dx{\displaystyle }dt\{\mathrm{\Gamma }^{(c)}(\stackrel{}{n},\stackrel{}{m})Q_{n,m;\alpha }^{(F)}(\stackrel{}{x},t)\mathrm{\Gamma }^{(s)}(\stackrel{}{n},\stackrel{}{m})H_{n,m;\alpha }^{(F)}(\stackrel{}{x},t)+\mathrm{\Gamma }_2^{(c)}(\stackrel{}{n},\stackrel{}{m})Q_{n,m;\alpha }^{(B)}(\stackrel{}{x},t)`$
$`\mathrm{\Gamma }_2^{(s)}(\stackrel{}{n},\stackrel{}{m})H_{n,m;\alpha }^{(B)}(\stackrel{}{x},t)+(1\delta _{n,m})\mathrm{\Gamma }_3^{(s)}(\stackrel{}{n},\stackrel{}{m})O_{n,m;\alpha }^{(s)}(\stackrel{}{x},t)\}`$
$`={\displaystyle \frac{\mathrm{\Lambda }^{1d}}{2N_o}}{\displaystyle \underset{n}{}}{\displaystyle \underset{m}{}}{\displaystyle \underset{\alpha }{}}{\displaystyle }d^dx{\displaystyle }dt\{\mathrm{\Gamma }^{(c)}(\stackrel{}{n},\stackrel{}{m})Q_{n,m;\alpha }^{(F)}(\stackrel{}{x},t)\mathrm{\Gamma }^{(s)}(\stackrel{}{n},\stackrel{}{m})H_{n,m;\alpha }^{(F)}(\stackrel{}{x},t)+e_2^{(c)}(\stackrel{}{n},\stackrel{}{m})Q_{n,m;\alpha }^{(B)}(\stackrel{}{x},t)`$
$$e_2^{(s)}(\stackrel{}{n},\stackrel{}{m})H_{n,m;\alpha }^{(B)}(\stackrel{}{x},t)+e_3^{(s)}(\stackrel{}{n},\stackrel{}{m})O_{n,m;\alpha }^{(s)}(\stackrel{}{x},t)\}$$
(29)
In Eq.29 we have ignored for simplicity the particle-particle triplet and consider only the particle-particle singlet.
Due to the relation between the particle-particle and the particle-hole triplets we remove the term $`(1\delta _{n,m})`$ by defining new coupling constants:
$`e_2^{(c)}(\stackrel{}{n},\stackrel{}{m})=\mathrm{\Gamma }_2^{(c)}(\stackrel{}{n},\stackrel{}{m}){\displaystyle \frac{1}{2}}\delta _{n,m}\mathrm{\Gamma }_3^{(s)}(\stackrel{}{n},\stackrel{}{m}),`$
$`e_2^{(s)}(\stackrel{}{n},\stackrel{}{m})=\mathrm{\Gamma }_2^{(s)}(\stackrel{}{n},\stackrel{}{m})2\delta _{n,m}\mathrm{\Gamma }_3^{(s)}(\stackrel{}{n},\stackrel{}{m}),`$
$$e_3^{(s)}(\stackrel{}{n},\stackrel{}{m})=\mathrm{\Gamma }_3^{(s)}(\stackrel{}{n},\stackrel{}{m}).$$
(30)
The results in Eqs.30 follows from the operator identity
$$O_{n,m;\alpha }^{(s)}(\stackrel{}{x},t)=(1\delta _{n,m})O_{n,m;\alpha }^{(s)}(\stackrel{}{x},t)+\delta _{n,m}[\frac{1}{2}Q_{n,m;\alpha }^{(B)}(\stackrel{}{x},t)2H_{n,m;\alpha }^{(B)}(\stackrel{}{x},t)]$$
(31)
Eq.31 follows directly from the definitions of the particle-particle singlet for $`\stackrel{}{n}=\stackrel{}{m}`$ in terms of the currents ( see Eqs.12-13 ). For $`\stackrel{}{n}=\stackrel{}{m}`$ the particle-hole triplet $`\mathrm{\Gamma }_2^{(s)}(\stackrel{}{n},\stackrel{}{n})`$ is related to the particle-particle singlet $`\mathrm{\Gamma }_3^{(s)}(\stackrel{}{n},\stackrel{}{n})`$ ( see Eq.A4 )
$$\frac{1}{2}\mathrm{\Gamma }_2^{(s)}(\stackrel{}{n},\stackrel{}{n})=\mathrm{\Gamma }_3^{(s)}(\stackrel{}{n},\stackrel{}{n}).$$
(32)
For the disorder part we will consider only the dominant Cooperon term $`d_3^{(s)}C_{n,m;\alpha ,\beta }^{(s)}(\stackrel{}{x};t_1,t_2)`$ and we will ignore the effect of forward disorder
$$\stackrel{~}{S}_D=\frac{\mathrm{\Lambda }^{2d}}{N_o}\underset{n}{}\underset{m}{}\underset{\alpha }{}\underset{\beta }{}d^dx𝑑t_1𝑑t_2\{d_3^{(s)}C_{n,m;\alpha ,\beta }^{(s)}(\stackrel{}{x};t_1,t_2)\}.$$
(33)
Following ref. we compute the partition function $`Z`$ of the action $`\stackrel{~}{S}_o+\stackrel{~}{S}_{int}+\stackrel{~}{S}_D`$ by expanding up to the third order in $`\stackrel{~}{S}_{int}+\stackrel{~}{S}_D`$. Using $`Z_o`$ we obtain:
$`Z=Z_o\{1[\stackrel{~}{S}_{int}_a+\stackrel{~}{S}_D_a{\displaystyle \frac{1}{2}}\stackrel{~}{S}_{int}^2_a\stackrel{~}{S}_{int}\stackrel{~}{S}_D_a{\displaystyle \frac{1}{2}}\stackrel{~}{S}_D^2_a`$
$$+\frac{1}{3!}\stackrel{~}{S}_{int}^3_a+\frac{1}{3!}\stackrel{~}{S}_D^3_a+\frac{1}{2}\stackrel{~}{S}_{int}^2\stackrel{~}{S}_D_a+\frac{1}{2}\stackrel{~}{S}_{int}\stackrel{~}{S}_D^2_a]\}.$$
(34)
The meaning of $`\mathrm{}_a`$ is to take the expectation value with respect to $`\stackrel{~}{S}_o`$ defined in Eq.28. Since we want to perform a RG analysis we will take the expectation value only in the interval $`(\mathrm{\Lambda },\mathrm{\Lambda }/b)`$, $`b1`$. In real space this means to integrate from the microscopic distance $`a`$ to $`ba`$.
Next we will compute the first term in Eq.34
$$\stackrel{~}{S}_{int}_{ba}=b^{2d}\frac{N_o}{N(b)}\stackrel{~}{S}_{int}_a,N(b)=N_ob$$
(35)
where $`\mathrm{}_{ba}`$ represents the expectation value with respect to Eq.28 with the new cutoff $`\mathrm{\Lambda }/b=2\pi /ba`$. The expectation value of $`\stackrel{~}{S}_D_a`$ is different from Eq.35. The difference is due to the two times $`t_1`$ and $`t_2`$. For times $`t_1t_2a/v_F`$, $`C_{n,m;\alpha ,\beta }^{(s)}(\stackrel{}{x};t_1,t_2)`$ is replaced by the singlet particle-particle interaction.
$$\stackrel{~}{S}_D_{ba}=b^{3d}\frac{N_o}{N(b)}\stackrel{~}{S}_D_a$$
(36)
$$\mathrm{\Delta }\stackrel{~}{S}_{int}_{ba}=\frac{2a}{v_F}b^{2d}\frac{N_o}{N(b)}\stackrel{~}{S}_D(t_1=t_2)_a.$$
(37)
Eq.37 represents the contribution from the disorder Cooperon to the singlet particle-particle term when $`t_1t_2a/v_F`$.
In order to compute the higher order term we have to use the rule of the operator product expansion defined in Eqs.18-27, and have to perform the time integration. We introduce the notation $`\mathrm{}_{da}`$ which stands for the expectation value in the domain ($`a`$,$`ba`$)
$`{\displaystyle \frac{1}{2}}\stackrel{~}{S}_{int}^2_{da}={\displaystyle \frac{da}{a}}\{{\displaystyle \frac{\mathrm{\Lambda }^{1d}}{2N}}{\displaystyle \frac{A^1}{4N}}(1){\displaystyle \underset{n}{}}{\displaystyle \underset{m}{}}{\displaystyle \underset{\alpha }{}}{\displaystyle }d^dx{\displaystyle }dt\{T(\stackrel{}{n},\stackrel{}{m})Q_{n,m;\alpha }^{(B)}(\stackrel{}{x},t)`$
$$S(\stackrel{}{n},\stackrel{}{m})H_{n,m;\alpha }^{(B)}(\stackrel{}{x},t)+R(\stackrel{}{n},\stackrel{}{m})O_{n,m;\alpha }^{(B)}(\stackrel{}{x},t)\}$$
(38)
where $`\frac{da}{a}=\frac{baa}{a}d\mathrm{ln}b`$. $`R(\stackrel{}{n},\stackrel{}{m})`$, $`S(\stackrel{}{n},\stackrel{}{m})`$, and $`T(\stackrel{}{n},\stackrel{}{m})`$ are a set of polynomials defined by the rules of the OPE given by Eqs.18-27.
$`T(\stackrel{}{n},\stackrel{}{m})=[2(e_3^{(s)}(\stackrel{}{n},\stackrel{}{m}))^2+{\displaystyle \frac{3}{4}}\delta _{n,m}e_2^{(s)}(\stackrel{}{n},\stackrel{}{m})e_3^{(s)}(\stackrel{}{n},\stackrel{}{m})],`$
$`R(\stackrel{}{n},\stackrel{}{m})=2{\displaystyle \underset{\stackrel{}{l}}{}}e_3^{(s)}(\stackrel{}{n},\stackrel{}{l})e_3^{(s)}(\stackrel{}{l},\stackrel{}{m})+{\displaystyle \frac{1}{4}}e_3^{(s)}(\stackrel{}{n},\stackrel{}{m})e_2^{(s)}(0),`$
$$S(\stackrel{}{n},\stackrel{}{m})=4(e_2^{(s)}(\stackrel{}{n},\stackrel{}{m}))^2$$
(39)
where $`e_2^{(s)}(0)e_2^{(s)}(\stackrel{}{n},\stackrel{}{n})`$. $`A^1`$ is determined by the time integration
$$A^1=\frac{I_1(\widehat{\beta })}{(2\pi )^{d1}2\pi v_F};\widehat{\beta }\frac{\beta }{\pi a}$$
(40)
$$I_1(\widehat{\beta })=\frac{2}{\pi }_0^{\mathrm{}}𝑑x\frac{\mathrm{cos}x/\widehat{\beta }}{x^2+1}$$
(41)
where $`\widehat{\beta }`$ is the dimensionless inverse temperature. The function $`I_1(\widehat{\beta })`$ originates at $`T0`$. In the limit $`\beta 1`$, $`I_1(\widehat{\beta })1`$. In the limit $`\beta 1`$, $`I_1(\widehat{\beta })1`$ and the time integration can be neglected.
$$\frac{1}{2}\stackrel{~}{S}_D^2_{da}=\frac{da}{a}\{\frac{\mathrm{\Lambda }^{2d}}{N}\frac{B^1}{2N}(1)\underset{n}{}\underset{m}{}\underset{\alpha }{}\underset{\beta }{}d^dx𝑑t_1𝑑t_2[2(d_3^{(s)})^2(1\frac{1}{8N})NC_{n,m;\alpha ,\beta }^{(s)}(\stackrel{}{x};t_1,t_2)]\}$$
(42)
and
$`\stackrel{~}{S}_{int}\stackrel{~}{S}_D^2_{da}={\displaystyle \frac{da}{a}}\{{\displaystyle \frac{\mathrm{\Lambda }^{1d}}{2N}}{\displaystyle \frac{B^1}{N}}{\displaystyle \underset{n}{}}{\displaystyle \underset{m}{}}{\displaystyle \underset{\alpha }{}}{\displaystyle }d^dx{\displaystyle }dt[d_3^{(s)}\widehat{L}(\stackrel{}{n},\stackrel{}{m})O_{n,m;\alpha }^{(s)}(\stackrel{}{x},t)]`$
$$\frac{\mathrm{\Lambda }^{2d}}{N}\frac{A^1}{2N}\underset{n}{}\underset{m}{}\underset{\alpha }{}\underset{\beta }{}d^dxdt_1dt_2[d_3^{(s)}\widehat{M}(\stackrel{}{n},\stackrel{}{m})C_{n,m;\alpha ,\beta }^{(s)}(\stackrel{}{x};t_1,t_2)]\}$$
(43)
where
$$B^1=\frac{I_2(2\widehat{\beta })}{(2\pi )^{d1}2v_F^2}$$
(44)
$$I_2(\widehat{\beta })=[I_1(\widehat{\beta })]^2$$
(45)
The term $`\widehat{L}(\stackrel{}{n},\stackrel{}{m})`$ and $`\widehat{M}(\stackrel{}{n},\stackrel{}{m})`$ are given by:
$$\widehat{L}(\stackrel{}{n},\stackrel{}{m})=2\underset{\stackrel{}{l}}{}e_3^{(s)}(\stackrel{}{l},\stackrel{}{m})+\frac{1}{2}e_2^{(s)}(\stackrel{}{n},\stackrel{}{m})2e_2^{(c)}(\stackrel{}{n},\stackrel{}{m})$$
(46)
and
$$\widehat{M}(\stackrel{}{n},\stackrel{}{m})=\frac{3}{2}e_2^{(s)}(\stackrel{}{n},\stackrel{}{m})+2\gamma _3^{(s)}(\stackrel{}{n},\stackrel{}{m})\delta _{n,m}2e_2^{(c)}(\stackrel{}{n},\stackrel{}{m})$$
(47)
The set of Eqs.43-47 concludes the RG calculation to second order.
The presence of the elastic mean free path introduces a cutoff in the time domain and allows us to apply the method of OPE to higher order. To third order in the interaction parameters $`e_2^{(s)}`$, $`e_2^{(c)}`$, $`e_3^{(s)}`$, and disorder $`d_3^{(s)}`$ we obtain:
$`\stackrel{~}{S}_{int}^2\stackrel{~}{S}_D_{da}={\displaystyle \frac{da}{a}}\{{\displaystyle \frac{\mathrm{\Lambda }^{2d}}{N}}{\displaystyle \frac{A^1}{2N^2}}{\displaystyle \underset{n}{}}{\displaystyle \underset{m}{}}{\displaystyle \underset{\alpha }{}}{\displaystyle \underset{\beta }{}}{\displaystyle }d^dx{\displaystyle }dt_1{\displaystyle }dt_2[{\displaystyle \frac{A^1}{2}}{\displaystyle \frac{J_1(\widehat{\beta })}{I_1(\widehat{\beta })}}G_1(\stackrel{}{n},\stackrel{}{m})`$
$`+B^1{\displaystyle \frac{J_2(\widehat{2\beta })}{I_2(2\widehat{\beta })}}\widehat{\mathrm{}}G_2(\stackrel{}{n},\stackrel{}{m})]C_{n,m;\alpha ,\beta }^{(s)}(\stackrel{}{x};t_1,t_2)\}+{\displaystyle \frac{da}{a}}\{{\displaystyle \frac{\mathrm{\Lambda }^{1d}}{2N}}{\displaystyle \frac{A^1B^1}{N^2}}{\displaystyle \underset{n}{}}{\displaystyle \underset{m}{}}{\displaystyle \underset{\alpha }{}}{\displaystyle }d^dx{\displaystyle }dt`$
$`[{\displaystyle \frac{J_2(\widehat{\beta })}{2I_2(\widehat{\beta })}}K_3(\stackrel{}{n},\stackrel{}{m})+{\displaystyle \frac{J_1(\widehat{\beta })}{I_1(\widehat{\beta })}}K_2(\stackrel{}{n},\stackrel{}{m})+{\displaystyle \frac{B}{A}}{\displaystyle \frac{J_1(\widehat{\beta })}{I_1(\widehat{\beta })}}K_3(\stackrel{}{n},\stackrel{}{m})]O_{n,m;\alpha }^{(s)}(\stackrel{}{x},t)\}`$
$$+\frac{da}{a}\{\frac{\mathrm{\Lambda }^{1d}}{2N}\frac{A^1B^1}{N^2}\frac{J_1(\widehat{\beta })}{I_1(\widehat{\beta })}\underset{n}{}\underset{m}{}\underset{\alpha }{}d^dx𝑑tF(\stackrel{}{n},\stackrel{}{m})Q_{n,m;\alpha }^{(B)}(\stackrel{}{x},t)\}$$
(48)
In Eq.48 the time integration introduces:
$$J_1(\widehat{\beta })I_1(\widehat{\beta }),J_2(\widehat{\beta })I_2(\widehat{\beta }).$$
(49)
The integral in Eq.49 depends explicitly on the dimensionless $`\widehat{\beta }`$. At the scale $`b=1`$ we have $`\widehat{\beta }(b=1)=\widehat{\beta }1`$ and for $`b=b_T=\beta /a`$ we have $`\widehat{\beta }(b)=1`$ and have to stop scaling. By Using the OPE rules we generate the polynomials $`G_1`$, $`G_2`$, $`K_1`$, $`K_2`$, $`K_3`$, and $`F`$. These polynomials are obtained from the microscopic couplings and the OPE results obtained at second order ( the polynomials $`R`$, $`S`$, $`T`$, $`L`$, and $`M`$).
$$G_1(\stackrel{}{n},\stackrel{}{m})=d_3^{(s)}[\frac{3}{2}S(\stackrel{}{n},\stackrel{}{m})+2R(0)\delta _{n,m}T(\stackrel{}{n},\stackrel{}{m})],$$
(50)
$$G_2(\stackrel{}{n},\stackrel{}{m})=d_3^{(s)}[2\widehat{M}(0)e_3^{(s)}(0)\delta _{n,m}+\frac{3}{2}\widehat{M}(\stackrel{}{n},\stackrel{}{m})e_2^{(s)}(\stackrel{}{n},\stackrel{}{m})2\widehat{M}(\stackrel{}{n},\stackrel{}{m})e_2^{(c)}(\stackrel{}{n},\stackrel{}{m})],$$
(51)
$$K_1(\stackrel{}{n},\stackrel{}{m})=d_3^{(s)}[2\underset{\stackrel{}{l}}{}R(\stackrel{}{l},\stackrel{}{m})+\frac{1}{2}S(\stackrel{}{n},\stackrel{}{m})T(\stackrel{}{n},\stackrel{}{m})],$$
(52)
$$K_2(\stackrel{}{n},\stackrel{}{m})=d_3^{(s)}[2\underset{\stackrel{}{l}}{}\stackrel{}{L}(\stackrel{}{n},\stackrel{}{l})e_3^{(s)}(\stackrel{}{l},\stackrel{}{m})+\frac{1}{4}\stackrel{}{L}(\stackrel{}{n},\stackrel{}{m})e_2^{(s)}(0)],$$
(53)
$$K_3(\stackrel{}{n},\stackrel{}{m})=d_3^{(s)}[2\underset{\stackrel{}{l}}{}\stackrel{}{M}(\stackrel{}{n},\stackrel{}{l})e_3^{(s)}(\stackrel{}{l},\stackrel{}{m})+\frac{1}{2}\stackrel{}{M}(\stackrel{}{n},\stackrel{}{m})e_2^{(s)}(\stackrel{}{n},\stackrel{}{m})+2\stackrel{}{M}(\stackrel{}{n},\stackrel{}{m})e_2^{(c)}(\stackrel{}{n},\stackrel{}{m})],$$
(54)
$$F(\stackrel{}{n},\stackrel{}{m})=d_3^{(s)}[2\stackrel{}{L}(\stackrel{}{n},\stackrel{}{m})e_3^{(s)}(\stackrel{}{n},\stackrel{}{m})\frac{3}{4}\delta _{n,m}e_2^{(s)}(0)\widehat{L}(\stackrel{}{n},\stackrel{}{m})].$$
(55)
Next we compute:
$`\stackrel{~}{S}_{int}\stackrel{~}{S}_D^2_{da}={\displaystyle \frac{da}{a}}\{{\displaystyle \frac{\mathrm{\Lambda }^{2d}}{N}}{\displaystyle \frac{A^1B^1}{N^2}}{\displaystyle \underset{n}{}}{\displaystyle \underset{m}{}}{\displaystyle \underset{\alpha }{}}{\displaystyle \underset{\beta }{}}{\displaystyle }d^dx{\displaystyle }dt_1{\displaystyle }dt_2[{\displaystyle \frac{J_1(\widehat{\beta })}{2I_1(\widehat{\beta })}}(d_3^{(s)})^2\widehat{M}(\stackrel{}{n},\stackrel{}{m})(4N1)`$
$$+\frac{J_2(2\widehat{\beta })}{I_2(2\widehat{\beta })}(d_3^{(s)})^2(2\widehat{L}(0)\delta _{n,m}+4\underset{\widehat{l}}{}\widehat{M}(\stackrel{}{l},\stackrel{}{m})\widehat{M}(\stackrel{}{n},\stackrel{}{m}))]C_{n,m;\alpha ,\beta }^{(s)}(\stackrel{}{x};t_1,t_2)\}.$$
(56)
The OPE in Eq.56 determines the behavior of the Cooperon as a function of the polynomials $`\widehat{M}`$ (see Eq.47) and Cooperon coupling $`d_3^{(s)}`$.
$`{\displaystyle \frac{1}{3!}}\stackrel{~}{S}_{int}^3_{da}={\displaystyle \frac{da}{a}}\{{\displaystyle \frac{\mathrm{\Lambda }^{1d}}{2N}}{\displaystyle \frac{(A^1)^2}{3!4N^2}}{\displaystyle \frac{J_1(\widehat{\beta })}{I_1(\widehat{\beta })}}{\displaystyle \underset{n}{}}{\displaystyle \underset{m}{}}{\displaystyle \underset{\alpha }{}}{\displaystyle }d^dx{\displaystyle }dt`$
$$[W(\stackrel{}{n},\stackrel{}{m})Q_{n,m;\alpha }^{(B)}(\stackrel{}{x},t)V(\stackrel{}{n},\stackrel{}{m})H_{n,m;\alpha }^{(B)}(\stackrel{}{x},t)+U(\stackrel{}{n},\stackrel{}{m})O_{n,m;\alpha }^{(s)}(\stackrel{}{x},t)]\}$$
(57)
where the functions $`U`$, $`V`$, and $`W`$ are defined in terms of the microscopic couplings and the second order functions $`R`$ and $`S`$ defined in Eq.39.
$`W(\stackrel{}{n},\stackrel{}{m})=[2R(\stackrel{}{n},\stackrel{}{m})e_3^{(s)}(\stackrel{}{n},\stackrel{}{m})+{\displaystyle \frac{3}{4}}\delta _{n,m}(R(0)e_2^{(s)}(0)+S(0)e_3^{(s)}(0))]`$
$`V(\stackrel{}{n},\stackrel{}{m})=4S(\stackrel{}{n},\stackrel{}{m})e_2^{(s)}(\stackrel{}{n},\stackrel{}{m})`$
$$U(\stackrel{}{n},\stackrel{}{m})=2\underset{\stackrel{}{l}}{}R(\stackrel{}{n},\stackrel{}{l})e_3^{(s)}(\stackrel{}{l},\stackrel{}{m})+\frac{1}{4}(R(\stackrel{}{n},\stackrel{}{m})e_2^{(s)}(0)+S(0)e_3^{(s)}(\stackrel{}{n},\stackrel{}{m}))$$
(58)
and
$`{\displaystyle \frac{1}{3!}}\stackrel{~}{S}_D^3_{da}={\displaystyle \frac{da}{a}}\{{\displaystyle \frac{\mathrm{\Lambda }^{2d}}{N}}(B^1)^2{\displaystyle \frac{J_2(2\widehat{\beta })}{I_2(2\widehat{\beta })}}{\displaystyle \frac{16}{3!}}{\displaystyle \underset{n}{}}{\displaystyle \underset{m}{}}{\displaystyle \underset{\alpha }{}}{\displaystyle \underset{\beta }{}}{\displaystyle }d^dx{\displaystyle }dt_1{\displaystyle }dt_2`$
$$(d_3^{(s)})^3(1\frac{1}{4N})^3C_{n,m;\alpha ,\beta }^{(s)}(\stackrel{}{x};t_1,t_2)\}$$
(59)
Using the results given in Eqs.29-59 we will obtain the RG equations.
## V The RG equations in the Quantum limit
The quantum region is defined by $`\mathrm{\Lambda }/b_T<|q|<\mathrm{\Lambda }`$ where $`b_T=\frac{v_F\mathrm{\Lambda }}{T}`$. In principle it is possible that before the scale $`b_T`$ has been reached, one of the coupling constants has reached values of order one. If this happens at a scale $`b_o<b_T`$ we have to stop at $`b_o`$ and for the interval $`\frac{\mathrm{\Lambda }}{b_T}|q|<\frac{\mathrm{\Lambda }}{b_o}`$ we have a different theory. If the Cooperon coupling constant $`\frac{\widehat{t}}{N}(k_F\mathrm{})^1`$ reaches values of order one at $`b_o<b_T`$ we must crossover to the Finkelstein diffusion theory. From the other hand if one of the two-body interactions reaches large values we have to construct a new theory. If the two-body interaction which grows under scaling is the Cooper coupling constant we have to construct a theory based on a superconductivity with disorder. We will consider here the situation where the effects of interactions are such that the value of $`b_ob_{Dif}`$ obeys $`b_{Dif}>b_T`$ or $`b_ob_{SC}`$, $`b_{SC}<b_T`$ ($`b_{SC}`$ is the length scale where the Cooper coupling constant diverges.). Therefore we will ignore the diffusive region.
We introduce the following rescaled coupling constants:
$`d_3^{(s)}=\widehat{t}B;e_2^{(c)}(\stackrel{}{n},\stackrel{}{m})=\widehat{e}_2^{(c)}(\stackrel{}{n},\stackrel{}{m})A;`$
$`\mathrm{\Gamma }_2^{(c)}(\stackrel{}{n},\stackrel{}{m})=\widehat{\gamma }_2^{(c)}(\stackrel{}{n},\stackrel{}{m})A;e_2^{(s)}(\stackrel{}{n},\stackrel{}{m})=\widehat{e}_2^{(s)}(\stackrel{}{n},\stackrel{}{m})A;`$
$`\mathrm{\Gamma }_2^{(s)}(\stackrel{}{n},\stackrel{}{m})=\widehat{\gamma }_2^{(s)}(\stackrel{}{n},\stackrel{}{m})A;e_3^{(s)}(\stackrel{}{n},\stackrel{}{m})=\widehat{e}_3^{(s)}(\stackrel{}{n},\stackrel{}{m})A;`$
$$\mathrm{\Gamma }_3^{(s)}(\stackrel{}{n},\stackrel{}{m})=\widehat{\gamma }_3^{(s)}(\stackrel{}{n},\stackrel{}{m})A.$$
(60)
where the constants $`A`$ and $`B`$ are defined in Eqs.40 and 44. In the quantum regime the number of channels obeys $`N_0N(b)=\pi (\frac{k_F}{\mathrm{\Lambda }/b})=N_0b`$. Due to the fact that when the cutoff $`\mathrm{\Lambda }`$ is reduced to $`\mathrm{\Lambda }/b`$ the number of channels scales like $`N(b)=N_0b`$, it follows that the naive scaling dimension of the interaction and disorder will be
$`{\displaystyle \frac{\gamma }{N}}\stackrel{\mathrm{\Lambda }/b}{}{\displaystyle \frac{\gamma }{N}}b^{2d}and{\displaystyle \frac{\widehat{t}}{N}}\stackrel{\mathrm{\Lambda }/b}{}{\displaystyle \frac{\widehat{t}}{N}}b^{3d}.`$
We observe that the interaction becomes marginal while the disorder is relevant. In the opposite situation where the number of channels does not scale, we have: $`\gamma \gamma b^{1d}`$ and $`\widehat{t}\widehat{t}b^{2d}`$.
For the disorder Cooperon coupling constant $`\widehat{t}`$ we have the scaling equation:
$`{\displaystyle \frac{d\widehat{t}}{d\mathrm{ln}b}}=\widehat{t}[1{\displaystyle \frac{1}{N}}({\displaystyle \frac{3}{4}}\widehat{\gamma }_2^{(s)}(\stackrel{}{n},\stackrel{}{m})\widehat{\gamma }_2^{(c)}(\stackrel{}{n},\stackrel{}{m})+{\displaystyle \frac{J_1(\widehat{\beta })}{I_1(\widehat{\beta })}}\delta _{n,m}[\widehat{\gamma }_3^{(s)}]_{n,m}^2)]`$
$$+2\widehat{t}^2[1\frac{1}{N}\frac{J_1(\widehat{\beta })}{I_1(\widehat{\beta })}(\frac{3}{2}\widehat{\gamma }_2^{(s)}(\stackrel{}{n},\stackrel{}{m})2\widehat{\gamma }_2^{(c)}(\stackrel{}{n},\stackrel{}{m}))\frac{1}{N}\frac{J_2(\widehat{\beta })}{I_2(\widehat{\beta })}(3\widehat{\gamma }_2^{(s)}2\widehat{\gamma }_2^{(c)})].$$
(61)
In Eq.61 we use the notation:
$$[\widehat{\gamma }_3^{(s)}]_{n,m}^2\frac{1}{N}\underset{\stackrel{}{l}}{}\widehat{\gamma }_3^{(s)}(\stackrel{}{n},\stackrel{}{l})\widehat{\gamma }_3^{(s)}(\stackrel{}{l},\stackrel{}{m})$$
(62)
$$\widehat{\gamma }_2^{(s)}=\frac{1}{N}\underset{\stackrel{}{l}}{}\gamma _2^{(s)}(\stackrel{}{l},\stackrel{}{n}),\widehat{\gamma }_2^{(c)}=\frac{1}{N}\underset{\stackrel{}{l}}{}\gamma _2^{(c)}(\stackrel{}{l},\stackrel{}{n}).$$
(63)
From Eq.61 we see that we can have a M-I transition in two dimensions when the p-p interaction $`\gamma _3^{(s)}`$ and the p-h $`\gamma _2^{(s)}`$ increases such that the linear term in “$`\widehat{t}`$” becomes negative (see Eq.61). We observe in Eq.61 that the effect of the p-h singlet $`\gamma _2^{(c)}`$ is opposite to the p-h triplet $`\gamma _2^{(s)}`$. $`\gamma _2^{(c)}`$ enhances the localization while $`\gamma _2^{(s)}`$ drives the system metallic. This is consistent with the known fact that a “Hartree” term ($`\gamma _2^{(c)}`$) favors localization while the “Fock” exchange term ($`\gamma _2^{(s)}`$) drives the system metallic. From dimensional analysis it follows that Eq.61 must be linear in $`\widehat{t}`$. In addition we have that the number of channels obey the scaling law, $`N=N(b)=N_ob`$.
The scaling equation for the particle-hole singlet is:
$`{\displaystyle \frac{d\widehat{\gamma }_2^{(c)}(\stackrel{}{n},\stackrel{}{m})}{d\mathrm{ln}b}}={\displaystyle \frac{1}{N}}\{(\widehat{\gamma }_3^{(s)}(\stackrel{}{n},\stackrel{}{m}))^2+{\displaystyle \frac{J_1(\widehat{\beta })}{I_1(\widehat{\beta })}}[\widehat{t}\widehat{\gamma }_3^{(s)}(\widehat{\gamma }_3^{(s)}(\stackrel{}{n},\stackrel{}{m})(13\delta _{n,m})+{\displaystyle \frac{3}{2}}\delta _{n,m}\widehat{\gamma }_2^{(s)}(0))`$
$$\frac{1}{24}[\widehat{\gamma }_3^{(s)}]_{n,m}^2(\widehat{\gamma }_3^{(s)}(\stackrel{}{n},\stackrel{}{m})(43\delta _{n,m})+\frac{3}{2}\delta _{n,m}\widehat{\gamma }_3^{(s)}(0))]\}+\frac{1}{2}\delta _{n,m}\frac{d\widehat{\gamma }_3^{(s)}(\stackrel{}{n},\stackrel{}{m})}{d\mathrm{ln}b}$$
(64)
and the particle-hole triplet $`\widehat{\gamma }_2^{(s)}(\stackrel{}{n},\stackrel{}{m})`$ is given by:
$$\frac{d\widehat{\gamma }_2^{(s)}(\stackrel{}{n},\stackrel{}{m})}{d\mathrm{ln}b}=\frac{1}{N}\{(\widehat{\gamma }_2^{(s)}(\stackrel{}{n},\stackrel{}{m}))^2+\frac{1}{6N}\frac{J_1(\widehat{\beta })}{I_1(\widehat{\beta })}(\widehat{\gamma }_2^{(s)}(\stackrel{}{n},\stackrel{}{m}))^3\})+2\delta _{n,m}\frac{d\widehat{\gamma }_3^{(s)}(\stackrel{}{n},\stackrel{}{m})}{d\mathrm{ln}b}$$
(65)
From Eqs.64 and 65 we see that the particle-particle channel affects the particle-hole singlet. In addition for $`\stackrel{}{n}=\stackrel{}{m}`$ the particle-particle channel $`\widehat{\gamma }_3^{(s)}(\stackrel{}{n},\stackrel{}{n})`$ is identicle to the particle-hole triplet $`\frac{1}{2}\widehat{\gamma }_2^{(s)}(\stackrel{}{n},\stackrel{}{n})`$, $`\widehat{\gamma }_3^{(s)}(\stackrel{}{n},\stackrel{}{n})=\frac{1}{2}\widehat{\gamma }_2^{(s)}(\stackrel{}{n},\stackrel{}{n})`$.
The particle-particle singlet term obeys the scaling equation:
$`{\displaystyle \frac{d\widehat{\gamma }_3^{(s)}(\stackrel{}{n},\stackrel{}{m})}{d\mathrm{ln}b}}={\displaystyle \frac{1}{2}}[\widehat{\gamma }_3^{(s)}]_{n,m}^2+\widehat{t}\widehat{\gamma }_3^{(s)}+{\displaystyle \frac{1}{3!}}{\displaystyle \frac{J_1(\widehat{\beta })}{I_1(\widehat{\beta })}}[\widehat{\gamma }_3^{(s)}]_{n,m}^3`$
$`2{\displaystyle \frac{J_2(\widehat{\beta })}{I_2(\widehat{\beta })}}\widehat{t}(\widehat{\gamma }_3^{(s)})^24{\displaystyle \frac{J_1(\widehat{\beta })}{I_1(\widehat{\beta })}}\widehat{t}(\widehat{\gamma }_3^{(s)})^2+8{\displaystyle \frac{J_2(\widehat{\beta })}{I_2(\widehat{\beta })}}\widehat{t}^2(\widehat{\gamma }_3^{(s)})+8{\displaystyle \frac{J_1(\widehat{\beta })}{I_1(\widehat{\beta })}}\widehat{t}^2(\widehat{\gamma }_3^{(s)})`$
$`+{\displaystyle \frac{1}{N}}\{[\widehat{t}+4\widehat{t}^2({\displaystyle \frac{J_1(\widehat{\beta })}{I_1(\widehat{\beta })}}+{\displaystyle \frac{A}{B}}{\displaystyle \frac{J_2(\widehat{\beta })}{I_2(\widehat{\beta })}})][{\displaystyle \frac{1}{2}}\widehat{\gamma }_2^{(s)}(\stackrel{}{n},\stackrel{}{m})2\widehat{\gamma }_2^{(c)}(\stackrel{}{n},\stackrel{}{m})]\}+{\displaystyle \frac{1}{N}}\{{\displaystyle \frac{1}{8}}\widehat{\gamma }_3^{(s)}(\stackrel{}{n},\stackrel{}{m})\widehat{\gamma }_3^{(s)}(0){\displaystyle \frac{1}{16}}\widehat{\gamma }_3^{(s)}(\stackrel{}{n},\stackrel{}{m})\widehat{\gamma }_2^{(s)}(0)`$
$$\frac{1}{3!}\frac{J_1(\widehat{\beta })}{I_1(\widehat{\beta })}\delta _{n,m}[\widehat{\gamma }_3^{(s)}]_{n,m}^2\widehat{\gamma }_3^{(s)}(0)+\frac{1}{2!3!}\frac{J_1(\widehat{\beta })}{I_1(\widehat{\beta })}[\widehat{\gamma }_3^{(s)}]_{n,m}^2\widehat{\gamma }_2^{(s)}(0)\frac{J_2(2\widehat{\beta })}{I_2(2\widehat{\beta })}(\widehat{\gamma }_3^{(s)})(\frac{1}{2}\widehat{\gamma }_2^{(s)}(0)\widehat{\gamma }_3^{(s)}(0)\delta _{n,m})\}$$
(66)
where
$`[\widehat{\gamma }_3^{(s)}]_{n,m}^2={\displaystyle \frac{1}{N}}{\displaystyle \underset{\stackrel{}{l}}{}}\widehat{\gamma }_3^{(s)}(\stackrel{}{n},\stackrel{}{l})\widehat{\gamma }_3^{(s)}(\stackrel{}{l},\stackrel{}{m})`$
$`(\widehat{\gamma }_3^{(s)})^2={\displaystyle \frac{1}{N}}{\displaystyle \underset{\stackrel{}{l}}{}}(\widehat{\gamma }_3^{(s)}(\stackrel{}{n},\stackrel{}{l}))^2`$
$`(\widehat{\gamma }_3^{(s)})={\displaystyle \frac{1}{N}}{\displaystyle \underset{\stackrel{}{l}}{}}\widehat{\gamma }_3^{(s)}(\stackrel{}{n},\stackrel{}{l})`$
$$[\widehat{\gamma }_3^{(s)}]_{n,m}^3=\frac{1}{N^2}\underset{\stackrel{}{l}}{}\underset{\stackrel{}{l}^{}}{}\widehat{\gamma }_3^{(s)}(\stackrel{}{n},\stackrel{}{l})\widehat{\gamma }_3^{(s)}(\stackrel{}{l},\stackrel{}{l}^{})\widehat{\gamma }_3^{(s)}(\stackrel{}{l}^{},\stackrel{}{m}).$$
(67)
The scaling relation for the forward part are trivial:
$$\frac{d\mathrm{\Gamma }^{(c)}}{d\mathrm{ln}b}=\frac{d\mathrm{\Gamma }^{(s)}}{d\mathrm{ln}b}=0$$
(68)
The set of Eqs.64-66 show that in the limit of $`N\mathrm{}`$ the interaction is controlled only by the particle-particle singlet $`\widehat{\gamma }_3^{(s)}(\stackrel{}{n},\stackrel{}{m})`$. In addition we observe that the disorder renormalizes the $`\widehat{\gamma }_3^{(s)}(\stackrel{}{n},\stackrel{}{m})`$. We observe that the scaling equation for $`\widehat{\gamma }_3^{(s)}(\stackrel{}{n},\stackrel{}{m})`$ can be negative at $`b=1`$. The origin of the negative term is given by Eq.37, where it has been shown that at short times the Cooperon behaves like a Cooper p-p singlet. As a result the initial values of the particle-particle singlet $`\widehat{\gamma }_3^{(s)}(\stackrel{}{n},\stackrel{}{m};b=1)`$ are replaced by $`\widehat{\gamma }_3^{(s)}(\stackrel{}{n},\stackrel{}{m};b=1)2v_F(\frac{B}{A})\widehat{t}`$. In Eqs.61, 64, 65, and 66 the scaling of the number of the channels is stopped when diffusive region is reached. At finite temperature we stop scaling at the scale $`b=b_T=E_F/T`$. This will fix the number of channels to $`\overline{N}N_T=E_F/T`$ (see ref.). It might be possible that in two dimensions the decoherency introduced by the temperature might be stronger than $`T`$. This might be the case if we have in mind dephasing effects in two dimensions which can define an effective temperature $`T_{eff}(T)>T`$ replacing $`\overline{N}`$ by $`E_F/T_{eff}`$.
## VI The conducting phase due to the superconducting instability in the quantum region
In the low temperature limit we can ignore all the many body effect except the particle-particle singlet $`\widehat{\gamma }_3^{(s)}(\stackrel{}{n},\stackrel{}{m})`$. The reason being the $`1/N`$ factor which appears in Eqs.64 and 65 and is missing for the particle-particle singlet in Eq.66. The growth of the number of channels $`N(b)`$ is determined by the topology of the Fermi surface. In particular this is the case for spherical Fermi surface where $`N(b)=N_ob`$. For $`T0`$ we obtain $`N(b=b_T)=\frac{E_F}{T}`$. (In chapter VIII we will consider non-spherical Fermi surface with repulsive interaction which might lead to a Ferromagnetic instability.)
Due to the fact that the $`1/N`$ factor is only absent for the particle-particle singlet, we will investigate the problem in the parameter space $`(\widehat{\gamma }_3^{(s)},\widehat{t})`$ using the angular momentum representation:
$`\gamma _3^{(s)}(r)\gamma _r={\displaystyle _0^\pi }{\displaystyle \frac{d\theta }{\pi }}\gamma _3^{(s)}(\theta )\mathrm{cos}(r\theta ),r=0,2,4,\mathrm{}.`$
For the singlet case $`r=0`$ we have $`\gamma _{r=0}=\gamma _o`$:
$$\gamma _o\frac{1}{N}\underset{\stackrel{}{l}}{}\widehat{\gamma }_3^{(s)}(\stackrel{}{l},\stackrel{}{n})$$
(69)
From Eq.66 we obtain to leading order in $`1/N`$ the following equation for particle-particle singlet:
$$\frac{d\gamma _o}{d\mathrm{ln}b}=\frac{1}{2}\gamma _o^2+\gamma _o\widehat{t}6\gamma _o^2\widehat{t}+16\gamma _o\widehat{t}^2+\frac{1}{3!}\gamma _o^3,$$
(70)
with
$`\gamma _o(b=1)\gamma _o(b=1)2\stackrel{~}{v_F}\widehat{t}`$
In Eq.70 we have used $`I_1I_2J_11`$ and $`\stackrel{~}{v_F}=v_FB/A`$ where the constants $`A`$ and $`B`$ have been defined in Eqs.40 and 44.
We investigate Eq.70 in the limit of weak disorder $`\widehat{t}0`$. We find that even for positive value of $`\gamma _o`$ the effect of disorder is to drive $`\gamma _o(b)`$ to negative values. The reason for this is the fact that the negative linear term in $`\widehat{t}`$ can cause an initial negative value for $`\gamma _o(b=1)`$. As a result the term $`\frac{1}{2}\gamma _o^2`$ (for negative value of $`\gamma _o`$, $`\gamma _o(b=1)<0`$) might drive the particle-particle interaction towards a superconducting instability. This behavior can be seen in the following way. In the limit of $`\widehat{t}0`$ we keep in Eq.70 only the two first order terms and obtain the solution for $`\gamma _o(b)`$:
$$\gamma _o(b=e^{\mathrm{}})=\gamma _o(b=1)exp(_0^{\mathrm{}}\widehat{t}(x)𝑑x)[1+\frac{1}{2}\gamma _o(b=1)_0^{\mathrm{}}𝑑yexp(_0^y\widehat{t}(x)𝑑x)]^1$$
(71)
For $`\gamma _o(b=1)<0`$, $`\gamma _o(b)`$ diverges at a length scale $`b=b_{SC}\frac{v_F\mathrm{\Lambda }}{T_{SC}}`$ where $`T_{SC}`$ represents the superconducting instability temperature,
$`T_{SC}=v_F\mathrm{\Lambda }(1+{\displaystyle \frac{2\widehat{t}}{|\gamma _o(b=1)|}})^{\frac{1}{t}}\stackrel{\widehat{t}0}{}v_F\mathrm{\Lambda }exp({\displaystyle \frac{2}{|\gamma _o(b=1)|}}).`$
Next we consider the RG equation for the Cooperon (see Eq.61 with $`J_1(\widehat{\beta })I_1(\widehat{\beta })I_2(\widehat{\beta })1`$). From Eq.61 we observe that in the limit of vanishing interactions the Cooperon coupling constant scales like $`\frac{\widehat{t}(b)}{N(b)}\frac{\widehat{t}(b=1)}{N_o}[12\widehat{t}(b=1)\mathrm{log}b]^1`$ and diverges at $`bb_{Loc}\frac{v_F\mathrm{\Lambda }}{T_{Loc}}`$ ($`b_{Loc}b_{Dif}`$, $`T_{Dif}T_{Loc}`$), $`T_{Loc}v_F\mathrm{\Lambda }exp[\frac{1}{2\widehat{t}(b=1)}]`$. In order to understand the physics of the system we have to compare the physical temperature $`T`$ with the other two, $`T_{SC}`$ and $`T_{Loc}`$. We have to consider separately the cases: a) $`T<T_{SC}<T_{Loc}`$; b) $`T<T_{Loc}<T_{SC}`$; c) $`T_{SC}<T_{Loc}<T`$; d) $`T_{SC}<T<T_{Loc}`$; e) $`T_{Loc}<T<T_{SC}`$; f) $`T_{Loc}<T_{SC}<T`$.
a) $`T<T_{SC}<T_{Loc}`$
This is the localized case where the mean free path “$`\mathrm{}`$” is the shortest length scale in the problem. This case will not be analyzed here. Most of the work in the past has been concentrated towards this case, in particular the Finkelstein theory which has investigated the interactions within the diffusion theory.
b) $`T<T_{Loc}<T_{SC}`$
Here the shortest length scale is the Cooper coherence length. Physically one can describe this region by a system of disorder bosons (the bosons describe the pairs). The critical theory might correspond to a disorder X-Y model.
c) $`T_{SC}<T_{Loc}<T`$
This is a region where interactions are not important. The physics is controlled by classical hopping transport.
d) $`T_{SC}<T<T_{Loc}`$
As in case a) here the system is localized. This case will not be considered here. (See the Finkelstein theory.)
e) $`T_{Loc}<T<T_{SC}`$
Again a bosonic X-Y theory with disorder is applicable here as in case b).
f) $`T_{Loc}<T_{SC}<T`$
In this region we will have transport controlled by pair breaking.
In the rest part of this section we will investigate the RG equation for the negative particle-particle singlet $`\widehat{\gamma }_3^{(s)}(\stackrel{}{n},\stackrel{}{m})`$ and the Cooperon coupling constant $`\widehat{t}`$. In agreement with Eq.69 we introduce the angular momentum representation for the Cooper and Cooperon channels: $`\gamma _o=\frac{1}{N}_\stackrel{}{ł}\widehat{\gamma }_3^{(s)}(\stackrel{}{l},\stackrel{}{n})`$, $`t_o=\frac{1}{N}_\stackrel{}{ł}\widehat{t}(\stackrel{}{l},\stackrel{}{n})`$. We obtain from Eq.61 and Eq.70 the following RG equations:
$`{\displaystyle \frac{d\lambda }{d\mathrm{ln}b}}={\displaystyle \frac{1}{2}}\lambda ^2+\lambda t_o,\lambda \gamma _o`$
$`{\displaystyle \frac{dt_o}{d\mathrm{ln}b}}=t_o[1({\displaystyle \frac{\lambda }{N}})^2]+2t_o^2`$
$$\rho (b)\frac{t_o(b)}{N(b)}\overline{t}_o(b),N(b)=N_ob$$
(72)
$`\rho (b)`$ is the resistance with $`b`$ restricted to $`1<b\frac{v_F\mathrm{\Lambda }}{T}`$. From Eq.72 we observe that in the limit $`b\mathrm{}`$ ($`T0`$) the parameter $`\lambda `$ diverges. In particular we observe that the ratio $`\frac{\lambda (b)}{N(b)}\stackrel{b\mathrm{}}{}\mathrm{}`$. As a result the RG equation behaves like $`\frac{dt_o}{d\mathrm{ln}b}=t_o(\frac{\lambda }{N})^2`$. Due to the large value of $`(\frac{\lambda }{N})^2`$ it follows that $`t_o(b)\stackrel{b\mathrm{}}{}0`$. As a result we obtain a superconducting ground state. At finite temperature we consider the case $`b_T<b_{SC}<b_{Loc}`$. We substitute the solution of $`\lambda (b)`$ into $`t_o(b)`$ and obtain
$$\overline{t}_o(b_T)=\frac{t_o}{N_o}exp\{_0^{\mathrm{log}b_T}(\frac{\lambda (x)}{N(x)})^2𝑑x\}\frac{t_o}{N_o}exp\{\frac{(4/N_o^2)T_{SC}}{|TT_{SC}|}\},T>T_{SC}$$
(73)
$`N_o\frac{\pi k_F}{\mathrm{\Lambda }}1`$ and $`T_{SC}`$ is given by Eq.71. As a result we obtain that the resistance obeys $`\rho (T)\stackrel{TT_{SC}}{}0`$, $`\rho (T)Const.exp\{\frac{(4/N_o^2)T_{SC}}{|TT_{SC}|}\}`$. To conclude this section ($`\gamma _o<0`$) we remark that the transport data show some similarity with the one reported for disorder bosons in ref.. This might suggest that the correct starting point might be a disordered bosonic system instead of a diffusion theory .
## VII The RG equation at finite temperature
At a temperature $`T`$ the scaling is restricted to $`\frac{\mathrm{\Lambda }}{b_T}<|q|<\mathrm{\Lambda }`$ where $`b_T=\frac{v_F\mathrm{\Lambda }}{T}`$. In this interval the number of channels is restricted to $`\overline{N}=N(b_T)=\frac{E_F}{T}`$, with $`N(b)`$ obeying the condition $`N_o<N(b)\overline{N}`$. We replace in Eqs.61-68 $`J_1(\widehat{\beta })J_2(\widehat{\beta })I_1(\widehat{\beta })I_2(\widehat{\beta })1`$ and find a simplified form
$`{\displaystyle \frac{d\widehat{t}}{d\mathrm{ln}b}}=ϵ(b)\widehat{t}{\displaystyle \frac{\widehat{t}}{N}}({\displaystyle \frac{3}{4}}\widehat{\gamma }_2^{(s)}(\stackrel{}{n},\stackrel{}{m})\widehat{\gamma }_2^{(c)}(\stackrel{}{n},\stackrel{}{m})+\delta _{n,m}[\widehat{\gamma }_3^{(s)}]_{n,m}^2)`$
$$+2\widehat{t}^2[1\frac{1}{N}(\frac{3}{2}\widehat{\gamma }_2^{(s)}(\stackrel{}{n},\stackrel{}{m})2\widehat{\gamma }_2^{(c)}(\stackrel{}{n},\stackrel{}{m}))\frac{1}{N}(3\widehat{\gamma }_2^{(s)}2\widehat{\gamma }_2^{(c)})]$$
(74)
The parameter $`ϵ(b)`$ controls the crossover at finite temperatures. $`ϵ(b)`$ is given by, $`ϵ(b)=1`$ for $`b<b_T`$ and $`ϵ(b)0`$ for $`b>b_T`$. Eq.74 replaces the scaling Eq.61 for the disorder coupling constant $`\widehat{t}`$. In Eq.74 we observe that the interaction has produced a shift in the critical dimensionality. The disorder parameter $`\widehat{t}`$ has accumulated a finite anomalous dimension, $`\frac{1}{N}(\frac{3}{4}\widehat{\gamma }_2^{(s)}\mathrm{})`$, which will control the M-I transition. (In the limit $`T0`$, $`N\mathrm{}`$ causing this term to disappear.)
At finite temperatures the scaling Eqs.64 and 65 for the interactions $`\widehat{\gamma }_2^{(s)}`$ and $`\widehat{\gamma }_2^{(c)}`$ are the same except that linear terms of the form $`[ϵ(b)1]\widehat{\gamma }_2^{(s)}`$ and $`[ϵ(b)1]\widehat{\gamma }_2^{(c)}`$ are added to the Eqs.64 and 65, respectively. For the particle-particle singlet $`\widehat{\gamma }_3^{(s)}`$ we have
$`{\displaystyle \frac{d\widehat{\gamma }_3^{(s)}(\stackrel{}{n},\stackrel{}{m})}{d\mathrm{ln}b}}=[ϵ(b)1]\widehat{\gamma }_3^{(s)}(\stackrel{}{n},\stackrel{}{m}){\displaystyle \frac{1}{2}}[\widehat{\gamma }_3^{(s)}]_{n,m}^2+\widehat{t}\widehat{\gamma }_3^{(s)}+{\displaystyle \frac{1}{3!}}[\widehat{\gamma }_3^{(s)}]_{n,m}^36\widehat{t}(\widehat{\gamma }_3^{(s)})^2+16\widehat{t}\widehat{\gamma }_3^{(s)}`$
$`+{\displaystyle \frac{1}{N}}\{(\widehat{t}+8\widehat{t}^2)({\displaystyle \frac{1}{2}}\widehat{\gamma }_2^{(s)}(\stackrel{}{n},\stackrel{}{m})2\widehat{\gamma }_2^{(c)}(\stackrel{}{n},\stackrel{}{m}))\}+{\displaystyle \frac{1}{N}}\{{\displaystyle \frac{1}{8}}\widehat{\gamma }_3^{(s)}(\stackrel{}{n},\stackrel{}{m})\widehat{\gamma }_3^{(s)}(0){\displaystyle \frac{1}{16}}\widehat{\gamma }_3^{(s)}(\stackrel{}{n},\stackrel{}{m})\widehat{\gamma }_2^{(s)}(0)`$
$$\frac{1}{3!}\delta _{n,m}[\widehat{\gamma }_3^{(s)}]_{n,m}^2\widehat{\gamma }_3^{(s)}(\stackrel{}{n},\stackrel{}{m})+\frac{1}{2!3!}[\widehat{\gamma }_3^{(s)}]_{n,m}^2\widehat{\gamma }_2^{(s)}(0)\widehat{\gamma }_3^{(s)}(\frac{1}{2}\widehat{\gamma }_2^{(s)}(0)\widehat{\gamma }_3^{(s)}(0)\delta _{n,m})\}$$
(75)
In Eq.75 we use the same definitions as given in Eq.67. Eq.75 must be supplemented by the condition $`\frac{1}{2}\widehat{\gamma }_2^{(s)}(\stackrel{}{n},\stackrel{}{n})=\widehat{\gamma }_3^{(s)}(\stackrel{}{n},\stackrel{}{n})`$ plus Eqs.64 and 65.
## VIII The scaling equations for the resistivity at finite temperatures and strong repulsive interactions
We restrict ourselves to finite temperatures or/and cases where the scaling of the number of channels is different from $`N(b)=N_ob`$ (spherical Fermi surface). For flat Fermi surface the number of the channels does not scale. We have $`N(b)N(b=1)N_o`$. At finite temperature for spherical Fermi surface the number of channels is finite and is restricted by the temperature $`N_o<N(b)<N(b_T)\frac{E_F}{T}`$. Since the coupling constants depend on the number of channels (finite), we will normalize the coupling constant by $`N`$, the number of channels
$$\overline{\gamma }_2^{(c)}\frac{\widehat{\gamma }_2^{(c)}}{N},\overline{\gamma }_2^{(s)}\frac{\widehat{\gamma }_2^{(s)}}{N},\overline{\gamma }_3^{(s)}\frac{\widehat{\gamma }_3^{(s)}}{N},\overline{t}\frac{\widehat{t}}{N}.$$
(76)
As a result the new RG equations are given in terms of the original Eqs. 75, 65, and 66:
$`{\displaystyle \frac{d\overline{\gamma }_2^{(c)}}{d\mathrm{ln}b}}={\displaystyle \frac{1}{N}}({\displaystyle \frac{d\widehat{\gamma }_2^{(c)}}{d\mathrm{ln}b}})ϵ_T\overline{\gamma }_2^{(c)};`$
$`{\displaystyle \frac{d\overline{\gamma }_2^{(s)}}{d\mathrm{ln}b}}={\displaystyle \frac{1}{N}}({\displaystyle \frac{d\widehat{\gamma }_2^{(s)}}{d\mathrm{ln}b}})ϵ_T\overline{\gamma }_2^{(s)};`$
$$\frac{d\overline{\gamma }_3^{(s)}}{d\mathrm{ln}b}=\frac{1}{N}(\frac{d\widehat{\gamma }_3^{(s)}}{d\mathrm{ln}b})ϵ_T\overline{\gamma }_3^{(s)}$$
(77)
The parameter $`ϵ_T`$ depends on the topology of the Fermi surface and temperature
$$ϵ_T|\frac{d\mathrm{ln}N(b)}{d\mathrm{ln}b}|,N_oN(b)N(b_T).$$
(78)
The parameter $`ϵ_T`$ takes values of $`0ϵ_T1`$. The value of $`ϵ_T=1`$ is obtained for spherical Fermi surface $`N(b)=N_ob`$ and $`ϵ_T=0`$ is obtained for flat Fermi surface or high temperatures, $`N(b)\overline{N}\frac{E_F}{T}`$.
Here we consider a special case of repulsive interactions such that the particle-particle singlet and particle-hole triplet are strong in the backward direction. This means that the most relevant interactions are those with $`\stackrel{}{n}=\stackrel{}{m}`$. In order to be specific we will consider a special model for which the terms $`\widehat{\gamma }_2^{(c)}(\stackrel{}{n},\stackrel{}{m})`$, $`\widehat{\gamma }_2^{(s)}(\stackrel{}{n},\stackrel{}{m})`$, and $`\widehat{\gamma }_3^{(s)}(\stackrel{}{n},\stackrel{}{m})`$ are zero for $`\stackrel{}{n}\stackrel{}{m}`$. We keep only terms with $`\stackrel{}{n}=\stackrel{}{m}`$ and introduce the definition:
$`\widehat{\gamma }_2^{(c)}\widehat{\gamma }_2^{(c)}(\stackrel{}{n},\stackrel{}{n}),\widehat{\gamma }_2^{(s)}\widehat{\gamma }_2^{(s)}(\stackrel{}{n},\stackrel{}{n})=2\widehat{\gamma }_3^{(s)}(\stackrel{}{n},\stackrel{}{n})`$
$$\widehat{\gamma }_2^{(c)}(\stackrel{}{n}\stackrel{}{m})\widehat{\gamma }_2^{(s)}(\stackrel{}{n}\stackrel{}{m})\widehat{\gamma }_3^{(s)}(\stackrel{}{n}\stackrel{}{m})0$$
(79)
Using Eqs. 76, 77, and 78 we obtain:
$$\frac{\overline{\gamma }_2^{(c)}}{d\mathrm{ln}b}=\overline{\gamma }_2^{(c)}(ϵ_T+\widehat{t})+\frac{1}{2}\widehat{t}\overline{\gamma }_2^{(s)}\frac{3}{4}(\overline{\gamma }_2^{(s)})^2(\widehat{t}\frac{1}{4})$$
(80)
$$\frac{\overline{\gamma }_2^{(s)}}{d\mathrm{ln}b}=\overline{\gamma }_2^{(s)}(2\widehat{t}+8\widehat{t}^2ϵ_T)(\overline{\gamma }_2^{(s)})^2(\frac{5}{4}+3\widehat{t})+\frac{5}{24}(\overline{\gamma }_2^{(s)})^34\overline{\gamma }_2^{(c)}(\widehat{t}+8\widehat{t}^2)$$
(81)
$$\frac{d\widehat{t}}{d\mathrm{ln}b}=\widehat{t}[1\frac{3}{4}\overline{\gamma }_2^{(s)}+\overline{\gamma }_2^{(c)}\frac{1}{4}(\overline{\gamma }_2^{(s)})^2]+2\widehat{t}^2[1\frac{3}{2}\overline{\gamma }_2^{(s)}+2\overline{\gamma }_2^{(c)}]$$
(82)
$$\overline{t}=\frac{\widehat{t}}{N},N=N(b),\mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}1}bb_T=\frac{E_F}{T}$$
(83)
From Eq.80 we conclude that the particle-hole singlet $`\overline{\gamma }_c`$ is irrelevant. Therefore we will take $`\overline{\gamma }_c=0`$ and ignore Eq.80. We will solve the RG equation in the space of $`\overline{\gamma }_2^{(s)}`$ and $`\widehat{t}`$ (Eqs. 81 and 82). In the parameter space $`(\overline{\gamma }_2^{(s)},\widehat{t})`$ we find a non-trivial fixed point. In the limit $`ϵ_T0`$ we find $`(\overline{\gamma }_2^{(s)})^{}=\frac{8}{5}(\widehat{t})^{}\frac{4}{25}`$.
We linearize the equations around this fixed point and find: $`\widehat{t}(b)=\widehat{t}^{}+(\widehat{t}\widehat{t}^{})b^{1/\nu _1}`$, $`\nu _11+\frac{2}{25}`$ and $`\overline{\gamma }_2^{(s)}(b)=\overline{\gamma }_2^{}+(\overline{\gamma }_2^{(s)}\overline{\gamma }_2^{})b^{1/\nu _2}`$. These equations show that for $`\widehat{t}<\widehat{t}^{}`$ the disorder decreases and in the same time $`\overline{\gamma }_2^{(s)}`$ flows to $`\overline{\gamma }_2^{}`$. For large value of disorder we obtain that $`\widehat{t}`$ increase and $`\overline{\gamma }_2^{(s)}`$ flows to $`\overline{\gamma }_2^{}`$. Experimentally the presence of the stable fixed point $`\overline{\gamma }_2^{}`$ might be identified by a power law behavior in the spin-spin correlation. This is similar to what one has in one dimension and might corresponds to a spin-liquid phase. For the transport properties, we believe that our predictions are in a qualitative agreement with the experiments , we find for the resistivity $`\rho (\overline{t},\overline{\gamma }_2^{(s)},T)`$ at a finite temperature and the dynamical exponent $`z1`$: $`\rho (\overline{t},\overline{\gamma }_2^{(s)},T)=\rho (\overline{t}(b),\overline{\gamma }_2^{(s)}(b),Tb^z)=\rho (\overline{t}^{}+(\overline{t}\overline{t}^{})b^{1/\nu _1}`$, $`\overline{\gamma }_2^{}+(\overline{\gamma }_2^{(s)}\overline{\gamma }_2^{})b^{1/\nu _2},Tb^z)`$. We introduce $`Tb^zT_ob(\frac{T_o}{T})^{1/z}`$ and use the definitions $`\overline{t}\frac{\widehat{t}}{\overline{N}}`$. As a result we find:
$$\rho (\overline{t},\overline{\gamma }_2^{(s)},T)\rho (\overline{t}^{},\overline{\gamma }_2^{},T_o)+const.(\frac{\overline{t}\overline{t}^{}}{\overline{t}^{}})(\frac{T_o}{T})^{1/z\nu _1}$$
(84)
Eq.84 shows that for $`\overline{t}<\overline{t}^{}`$ the resistivity $`\rho `$ decreases as we lower the temperature and increases when $`\overline{t}>\overline{t}^{}`$. In order to make contact with the experiments we replace: $`\overline{t}(k_F\mathrm{})^1`$, $`k_Fn_c^{1/2}`$, $`\overline{t}^{}(n_c^{})^{1/2}`$ ($`n_c^{}`$ is the critical density) and identify$`\frac{\overline{t}\overline{t}^{}}{\overline{t}^{}}\frac{n_c^{}n_c}{n_c^{}}\delta `$. As a result we find: $`\rho (n_c,\overline{\gamma }_2^{(s)},T)\rho ^{}f(\frac{\delta }{T^{1/z\nu _1}})`$, $`\rho ^{}\rho (n_c^{},\gamma _2^{},T_o)`$ which is the result observed in ref.. We hope that more accurate experiments will confirm the existence of the suggested fixed point.
## IX Conclusion
A new method for studying many-body systems and disorder has been introduced. The method is based on the extension of the OPE to two dimensional systems. Using a real space version of RG we have derived a set of RG equations for disorder and interaction. We have constructed an alternative theory to the one constructed by Finkelstein . The basic assumption in ref. is that the elastic mean free path is the shortest length in the problem. As a result the multiple elastic scatterings are replaced by a diffusion theory (the non-linear $`\sigma `$-model) and the interactions are considered as a perturbation of the diffusion theory. The method used here is based on a RG analysis which studies the competitions between the multiple elastic scattering and the interaction. We identified the following regions:
1) The multiple elastic scattering is the shortest length scale and diverges first. For this case we agree with the results given in ref. and do not have anything to add.
2) The particle-particle singlet is negative and a superconducting instability occurs for $`TT_{SC}`$ where $`T_{SC}>T_{Loc}`$. As a result one has to treat first the interaction within an effective Ginzburg-Landau theory. We reproduced a bosonic model (X-Y) which is perturbed by disorder.
3) The interactions are positive and the particle-hole is dominant in the backward direction. At finite temperature and non-spherical Fermi surfaces which obey $`|\frac{d\mathrm{ln}N(b)}{d\mathrm{ln}b}|1`$ one obtains a non-trivial fixed point in the plane $`(\overline{\gamma }_2^{(s)},\widehat{t})`$ which separates the conducting from the insulating phase. This fixed point is characterized by a stable fixed point in the $`\overline{\gamma }_2^{(s)}`$ direction. No divergence in the particle-hole triplet occurs expect the infinite correlation length for the spin-spin ferromagnetic correlations when $`\overline{\gamma }_2^{(s)}\overline{\gamma }_2^{}`$.
ACKNOWLEDGMENTS
D. Schmeltzer would like to thank professor A.M. Finkelstein for his valuable comment concerning the differences between his theory and the one presented here.
## A
The Fermion field $`\psi _{\sigma ,\alpha }(\stackrel{}{x})`$ is decomposed into $`N`$ Fermions, $`\psi _{\sigma ,\alpha }(\stackrel{}{x})=_{\stackrel{}{\omega }=1}^Ne^{ik_F\stackrel{}{\omega }\stackrel{}{x}}\psi _{\stackrel{}{\omega },\sigma ,\alpha }`$. Using this representation we obtain from Eq.3 the result:
$`S_{int}{\displaystyle d^dx𝑑t\underset{\sigma ,\sigma ^{}}{}\underset{\alpha }{}\underset{\stackrel{}{\omega }_1}{}\underset{\stackrel{}{\omega }_2}{}\underset{\stackrel{}{\omega }_3}{}\underset{\stackrel{}{\omega }_4}{}\delta _{\stackrel{}{\omega }_1+\stackrel{}{\omega }_2=\stackrel{}{\omega }_3+\stackrel{}{\omega }_4}}`$
$$v(\stackrel{}{\omega }_1,\stackrel{}{\omega }_2,\stackrel{}{\omega }_3,\stackrel{}{\omega }_4)\psi _{\stackrel{}{\omega }_1,\sigma ,\alpha }^{}(\stackrel{}{x})\psi _{\stackrel{}{\omega }_2,\sigma ^{},\alpha }^{}(\stackrel{}{x})\psi _{\stackrel{}{\omega }_3,\sigma ^{},\alpha }(\stackrel{}{x})\psi _{\stackrel{}{\omega }_4,\sigma ,\alpha }(\stackrel{}{x})$$
(A1)
where $`v(\stackrel{}{\omega }_1,\stackrel{}{\omega }_2,\stackrel{}{\omega }_3,\stackrel{}{\omega }_4)`$ represents the projection of the screened two-body potential on the Fermi surface. The presence of the Kroneker-delta function imposes the condition $`\stackrel{}{\omega }_1+\stackrel{}{\omega }_2=\stackrel{}{\omega }_3+\stackrel{}{\omega }_4`$. As a result we separate the interaction term into three processes: 1) direct, 2) exchange, and 3) Cooperon channel:
1. The direct process is realized when $`\stackrel{}{\omega }_1=\stackrel{}{\omega }_4`$, $`\stackrel{}{\omega }_2=\stackrel{}{\omega }_3`$.
2. The exchange process: $`\stackrel{}{\omega }_1=\stackrel{}{\omega }_3\stackrel{}{\omega }`$, $`\stackrel{}{\omega }_2=\stackrel{}{\omega }_4\stackrel{}{\omega }^{}`$
3. The Cooperon channel: $`\stackrel{}{\omega }\stackrel{}{\omega }_1=\stackrel{}{\omega }_2`$, $`\stackrel{}{\omega }^{}\stackrel{}{\omega }_3=\stackrel{}{\omega }_4`$
As a result Eq.A1 becomes
$`S_{int}{\displaystyle }d^dx{\displaystyle }dt{\displaystyle \underset{\sigma ,\sigma ^{}}{}}{\displaystyle \underset{\alpha }{}}{\displaystyle \underset{\stackrel{}{\omega }}{}}{\displaystyle \underset{\stackrel{}{\omega }^{}}{}}\{v(0)\psi _{\stackrel{}{\omega },\sigma ,\alpha }^{}(\stackrel{}{x})\psi _{\stackrel{}{\omega },\sigma ,\alpha }(\stackrel{}{x})\psi _{\stackrel{}{\omega }^{},\sigma ^{},\alpha }^{}(\stackrel{}{x})\psi _{\stackrel{}{\omega }^{},\sigma ^{},\alpha }(\stackrel{}{x})`$
$`v(\stackrel{}{\omega },\stackrel{}{\omega }^{},\stackrel{}{\omega },\stackrel{}{\omega }^{})\psi _{\stackrel{}{\omega },\sigma ,\alpha }^{}(\stackrel{}{x})\psi _{\stackrel{}{\omega },\sigma ^{},\alpha }(\stackrel{}{x})\psi _{\stackrel{}{\omega }^{},\sigma ^{},\alpha }^{}(\stackrel{}{x})\psi _{\stackrel{}{\omega }^{},\sigma ,\alpha }(\stackrel{}{x})`$
$$+v(\stackrel{}{\omega },\stackrel{}{\omega },\stackrel{}{\omega }^{},\stackrel{}{\omega }^{})\psi _{\stackrel{}{\omega },\sigma ,\alpha }^{}(\stackrel{}{x})\psi _{\stackrel{}{\omega },\sigma ^{},\alpha }(\stackrel{}{x})\psi _{\stackrel{}{\omega }^{},\sigma ^{},\alpha }^{}(\stackrel{}{x})\psi _{\stackrel{}{\omega }^{},\sigma ,\alpha }(\stackrel{}{x})\}$$
(A2)
In Eq.A2 we observe that the Cooperon channel is identical to the exchange one if we substitute in the exchange term $`\stackrel{}{\omega }^{}=\stackrel{}{\omega }`$. This means that we have to take into consideration this identity in order to avoid double counting.
We replace the “$`N`$” fermions by $`N/2`$ pairs of chiral fermions (see Eq.5). In the second step we replace Eq.A2 by the current representation:
$`S_{int}{\displaystyle }d^dx{\displaystyle }dt{\displaystyle \underset{\alpha }{}}{\displaystyle \underset{\stackrel{}{n},\stackrel{}{m}}{}}\{(v(0){\displaystyle \frac{1}{2}}v(\stackrel{}{n},\stackrel{}{m}))(J_{n,\alpha }^R(\stackrel{}{x})J_{m,\alpha }^R(\stackrel{}{x})+J_{n,\alpha }^L(\stackrel{}{x})J_{m,\alpha }^L(\stackrel{}{x}))`$
$`2v(\stackrel{}{n},\stackrel{}{m})(J_{n,\alpha }^R(\stackrel{}{x})J_{m,\alpha }^R(\stackrel{}{x})+J_{n,\alpha }^L(\stackrel{}{x})J_{m,\alpha }^L(\stackrel{}{x}))+(v(0){\displaystyle \frac{1}{2}}v(\stackrel{}{n},\stackrel{}{m}+\pi ))(J_{n,\alpha }^R(\stackrel{}{x})J_{m,\alpha }^L(\stackrel{}{x})`$
$`+J_{n,\alpha }^L(\stackrel{}{x})J_{m,\alpha }^R(\stackrel{}{x}))2v(\stackrel{}{n},\stackrel{}{m}+\pi )(J_{n,\alpha }^R(\stackrel{}{x})J_{m,\alpha }^L(\stackrel{}{x})+J_{n,\alpha }^L(\stackrel{}{x})J_{m,\alpha }^R(\stackrel{}{x}))`$
$`+(1\delta _{n,m})[v(\stackrel{}{n},\stackrel{}{m}){\displaystyle \underset{\sigma =,}{}}(J_{n,\sigma ,\alpha ;m,\sigma ,\alpha }^R(\stackrel{}{x})J_{n,\sigma ,\alpha ;m,\sigma ,\alpha }^L(\stackrel{}{x})+J_{n,\sigma ,\alpha ;m,\sigma ,\alpha }^L(\stackrel{}{x})J_{n,\sigma ,\alpha ;m,\sigma ,\alpha }^R(\stackrel{}{x}))`$
$$v(\stackrel{}{n},\stackrel{}{m}+\pi )\underset{\sigma =,}{}(J_{n,\sigma ,\alpha ;m,\sigma ,\alpha }^R(\stackrel{}{x})J_{n,\sigma ,\alpha ;m,\sigma ,\alpha }^L(\stackrel{}{x})+J_{n,\sigma ,\alpha ;m,\sigma ,\alpha }^L(\stackrel{}{x})J_{n,\sigma ,\alpha ;m,\sigma ,\alpha }^R(\stackrel{}{x}))]\}$$
(A3)
We introduce the following definitions:
$`\stackrel{~}{\mathrm{\Gamma }}^{(c)}(\stackrel{}{n},\stackrel{}{m})v(0){\displaystyle \frac{1}{2}}v(\stackrel{}{n},\stackrel{}{m}),\stackrel{~}{\mathrm{\Gamma }}^{(s)}(\stackrel{}{n},\stackrel{}{m})2v(\stackrel{}{n},\stackrel{}{m}),`$
$`\mathrm{\Gamma }_2^{(c)}(\stackrel{}{n},\stackrel{}{m})v(0){\displaystyle \frac{1}{2}}v(\stackrel{}{n},\stackrel{}{m}+\pi ),\mathrm{\Gamma }_2^{(s)}(\stackrel{}{n},\stackrel{}{m})2v(\stackrel{}{n},\stackrel{}{m}+\pi ),`$
$$\mathrm{\Gamma }_3^{(s)}(\stackrel{}{n},\stackrel{}{m})\frac{1}{2}(v(\stackrel{}{n},\stackrel{}{m})+v(\stackrel{}{n},\stackrel{}{m}+\pi )),\mathrm{\Gamma }_3^{(t)}(\stackrel{}{n},\stackrel{}{m})\frac{1}{2}(v(\stackrel{}{n},\stackrel{}{m})v(\stackrel{}{n},\stackrel{}{m}+\pi )).$$
(A4)
Using the definitions of the interaction operators given in Eqs.LABEL:012 and 13 we obtain the result given in Eq.11. |
warning/0002/astro-ph0002473.html | ar5iv | text | # 1 Introduction
## 1 Introduction
IRC +10 216 (CW Leo) is the nearest and best–studied carbon star and one of the brightest infrared sources. It experiences strong mass-loss rates of $`\dot{M}25\times 10^5`$M$`_{}`$yr<sup>-1</sup> (Loup et al. 1993). The central star of IRC +10 216 is a long–period variable with a period of $`649`$ days (Le Bertre 1992). Recent distance estimates range from 110 pc to 150 pc (Groenewegen 1997, Crosas & Menten 1997). IRC +10 216’s initial mass can be expected to be close to 4 M (Guelin et al. 1995, Weigelt et al. 1998). The bipolar appearance of the nebula around this object was already reported, e.g., by Kastner & Weintraub . The non-spherical structure is consistent with the conjecture that IRC +10 216 is in a phase immediately before entering the protoplanetary nebula stage. The most recent high–resolution observations of this object and its circumstellar dust shell were reported by Weigelt et al. , Haniff & Buscher , and Tuthill et al. . The results of Dyck et al. and Haniff & Buscher showed that the dust-shell structure of IRC +10 216 is changing within some years.
## 2 Observations and bispectrum speckle interferometry results
The IRC +10 216 speckle interferograms were obtained with the 6 m telescope at the Special Astrophysical Observatory in Russia. At all continuum epochs data within the $`K`$–band were obtained (date / center wavelength of the filter in $`\mu `$m / FWHM bandwidth of the filter in $`\mu `$m: 8.10.95/2.12/0.02, 3.4.96/2.17/0.02, 23.1.97/2.19/0.41, 14.6.98/2.17/0.33, 3.11.98/2.19/0.19). $`J`$\- and $`H`$-band data were recorded at one epoch (2.4.96/1.24/0.28, 23.1.97/1.64/0.31) each.
Figures 1 and 2 show the $`K`$, $`H`$, and $`J`$ images of the central region of IRC +10 216 for all epochs. The high–resolution images were reconstructed from the speckle interferograms using the bispectrum speckle interferometry method (Weigelt 1977, Lohmann et al. 1983, Weigelt 1991). We denote the resolved components in the center of the nebula as A, B, C, and D (see Fig. 2b) in the order of decreasing peak intensity (based on the $`K`$ band results from 1996). Figure 2b shows in addition three fainter components denoted with E, F, and G. The faint extended feature at position angle PA $``$340 in the $`J`$ image corresponds quite well to the faint component E visible in all images in Figs. 1 and 2 (assuming that the brightest component in the $`J`$ image is coinciding with component A in the $`H`$ and $`K`$ images, see Figure 3a).
Figure 3 shows the results of polarimetric observations with the HST Nicmos camera at a wavelength of 1.1 $`\mu `$m (raw data retrieved from the Hubble Data Archive, STScI). The data were obtained at a photometric phase of $`\mathrm{\Phi }=0.76`$. From the data the total intensity (Figs. 3a and b), the polarized intensity (Fig. 3c), the degree and position angle of the polarization (Fig. 3b) have been derived.
## 3 Nebula structures
The bright inner region of IRC +10 216 is surrounded by a larger faint nebula. The bipolar shape of this nebula is most prominently present in the $`J`$–band image (Fig. 2c) and in the HST images at shorter wavelengths (Fig. 3a). However, the fact that even in the polarized intensity the nebula is very faint on the southeastern side suggests the main axis of the nebula to be at PA$`20^{}`$ to $`30^{}`$ along the direction from component A to B. This fits well to the main axis of the H<sup>13</sup>CN($`J=10`$) emission (Dayal & Bieging 1995) which is weakly elongated on a scale of about 10<sup>′′</sup>.
We determined the separations of the components A and B for the 5 epochs from 1995 to 1998 shown in Fig. 1. The separations are: 191 mas, 201 mas, 214 mas, 245 mas, and 265 mas. A linear regression fit gives a value of 23 mas/yr for the average increase in the apparent separation of the components. Interpreting this increase as a real motion would lead to 14 km/s within the plane of the sky (for a distance of $`130`$ pc). The result comprises data from more than one pulsation period ($``$649 days, see Le Bertre 1992) and it is thus obvious that the apparent relative motion of the nebula components is not simply related to the stellar variability.
Besides the motion of the components, Fig. 1 shows that these components change their shapes and relative fluxes. The brightest component A becomes narrower along the axis A–B ($``$20). The peak–to–peak intensity ratio of B and A is almost constant from 1995 to 1997. Later component B is fading. At the same time the other components become brighter and detached from A. Note that the photometric phases and the integral $`K`$ magnitudes of IRC +10 216 in January 1997 and November 1998 are almost identical (Osterbart et al. 1999). Again we find that the time scale for the changes seen in our images is significantly different from the period of the stellar pulsation.
## 4 Discussion
In the following we want to address the question where, behind all the dust, the central star in the system of IRC +10 216 is located. This question is of specific interest to understand the physical properties of the nebula. At short wavelengths the nebula shows a bipolar structure. A comparison of the observed structure to other bipolar nebulae like the Red Rectangle (Men’shchikov et al. 1998) suggests that the X–like arms originate mainly from scattering of stellar light on the surfaces of cavities. The star then is at least partially obscured by an optically thick dust shell or torus. Convincingly Haniff & Buscher argued that the main axis of the object is tilted with its southern side towards the observer.
Is component A the star? At first glance it seems reasonable to assume that the star is at the position of the brightest component A. However, the synthetic polarization maps of Fischer et al. show that a significant polarization at the position of the star is only present for nearly edge–on configurations. The high degree of polarization of A ($`P14\%`$) is thus not in agreement with A being the star because the structure of the bipolar nebula at short wavelengths suggests that we are looking at an intermediate viewing angle (e.g. 50 to 60). At larger separations from A the polarization pattern is centrosymmetric. The center of such a pattern is thought to be at the position of the illuminating source (Fischer et al. 1996). Since in the map in Fig. 3b this center does not coincide with A, but is located significantly north of A, we are led to the conclusion that A is not the star.
Is the star at the position of B? In the following we will show that it is consistent with the observation to assume that the star is at the position of B.
(a) The cometary shapes of A (Fig. 2) as well as the presented polarization data strongly suggest that A is part of a scattering lobe within a bipolar structure. Consistently, component A and its southern tails are relatively blue ($`HK`$ in the range from 2 to 3.2; see Osterbart et al. 1999) compared to the integral color ($`HK=3.2`$).
(b) The northern components B, C, and D are, on the other hand, rather red ($`HK4.2`$) in comparison with the integral color. This suggests that these structures are strongly obscured and reddened by circumstellar dust.
(c) The brightest northern component in the $`J`$ image (Fig. 2c) can hardly be seen in the $`H`$ image and is thus very blue. This component and A can be considered as opposite lobes of an almost bipolar structure. Component B is almost in the center between these counterlobes and thus approx. at the position where the star would be expected.
(d) The polarization map (Fig. 3b) fits well to the picture that the star is at B. The center of the centrosymmetric polarization pattern at larger separations from A is located between the two $`J`$–band lobes. This is consistent with the assumption that the illuminating source is at or near the postion of B (cf. Fischer et al. 1996). Unfortunately, the polarization at the position of B itself is contaminated by contributions from the diffraction pattern associated with the peak A.
Changes in the mass loss rate. The change of the shape of component A and the fading of B can be attributed to an increasing mass loss which is accompanied by a gradual increase of the optical depth of the dust shell. This is most obvious for the later observation epochs, suggesting an enhanced mass loss since 1997. A strongly variable mass loss has, in fact, been predicted by theoretical models treating the dust formation mechanism in the envelopes of long–period variable carbon stars (Winters et al. 1995). Periods of this mechanism may be significantly longer than the stellar pulsation period.
Alternative models: the star between A and B or near B. It is not possible to exclude that the star may be located between A and B, close to B, or in the center of A, B, C, and D. The precision of the polarization map is not sufficient to conclude whether the star is at the position of B or only close to it. Two–dimensional radiative transfer calculations in progess show, however, that the observed intensity ratio of A and B as well as the components’ shapes clearly require the star to be at B.
## 5 Stellar evolution and bipolar structure
IRC +10 216 is without doubt in a very advanced stage of its AGB evolution due to its long pulsational period, high mass-loss rate, and carbon-rich dust-shell chemistry indicating that already a significant number of thermal pulses did take place. The star’s initial mass can be estimated to be $`4`$M$`{}_{}{}^{}\pm 1`$M due to the observed isotopic ratios of C, N and O in the dust shell (Guelin et al. 1995) and the luminosity of the central star (Weigelt et al. 1998). Accordingly, the core mass should be $`0.7`$ to $`0.8`$M with corresponding thermal-pulse cycle times of $`1\mathrm{3\hspace{0.17em}10}^4`$yr (Blöcker 1995). Introducing the mean observed mass-loss rate to these thermal-pulse periods shows that the present stellar wind leads to a very effective erosion of the envelope per thermal pulse cycle, possibly as high as $`1`$M/cycle. Consequently, the whole envelope may be lost during the next few thermal pulses leading to the termination of the AGB evolution. Thus, it is not unlikely to assume that IRC +10 216 has entered a phase immediately before moving off the AGB. This is strongly supported by the non-spherical appearence of its dust shell showing even bipolar structures. Unlike AGB stars, post-AGB objects as protoplanetary nebulae often expose prominent features of asphericities, in particular in axisymmetric geometry (e.g. Olofsson 1996). Accordingly, IRC +10 216 can be thought to be an object in transition. It is noteworthy that the establishment of bipolar structures, i.e. the metamorphosis into a protoplanetary nebulae, obviously already begins during the (very end of) AGB evolution. The clumpiness within the bipolar shape is probably due to small scale fluctuations of the dust condensation radius which, in turn, might be influenced by, e.g., giant surface convection cells (Schwarzschild 1975). The formation of giant convection cells can be assumed to be a common phenomenon in red giants.
The shaping of planetary nebulae can successfully be described by interacting stellar wind models (Kwok et al. 1978) Within this scenario a fast (spherical) wind from the central star interacts with the slow wind of the preceding AGB evolution. The slow AGB wind is asssumed to be non-spherical (axisymmetric) which leads to the observed morphology of planetary nebulae (Mellema 1996). Different mechanisms to provide the required equatorial density enhancements are discussed (cf. Livio 1993). Among these, binarity is one channel including common envelope evolution and spin up of the AGB star due to the interaction with its companion (Morris 1981). Not only stellar companions are found to be able to spin up the AGB star but also substellar ones as brown dwarfs and planets, most effectively by evaporation in the AGB star’s envelope (Soker 1997). Currently there is no observational evidence for a possible binary nature of IRC +10 216. The fact that the polarization pattern in the southeastern part of the nebula at 1.1 $`\mu `$m has a different orientation than in the rest of the nebula might be an indication for a second illuminating source.
Mechanisms inherent to the star include rotation, non-radial pulsations, and magnetic fields (see e.g. Dorfi & Höfner 1996, Soker & Harpaz 1992, Garcia-Segura et al. 1999). Both non-radial $`p`$-modes and magnetic fields appear to be only important for significant rotation rates. Often spin-up agents due to binarity are assumed. For instance, Groenewegen favours non-radial pulsation or a yet unidentified companion which spun up the central star as the most likely explanation for the non-spherical shape of the dust shell of IRC +10 216.
AGB stars are known to be slow rotators. Stars with initial masses below $`1.3`$M can be expected to lose almost their entire angular momentum during the main sequence phase due to magnetic braking operating in their convective envelopes. Consequently they are not believed to develop non-spherical mass-loss due to rotation. Stars with larger initial masses are spun down due to mass loss but may achieve sufficiently high rotation rates at the end of AGB evolution (Garcia-Segura et al. 1999). Already small rotation rates influence dust-driven winds considerably yielding a mass loss preferentially driven in the equatorial plane (Dorfi & Höfner 1996). For supergiants leaving the Hayashi line Heger & Langer found that significant spin up of the surface layers may take place. Thus, at second glance, rotation might be able to support axisymmetric mass loss during the transition to the proto-planetary nebula phase for AGB stars as IRC +10 216.
## 6 Conclusions
We have presented high-resolution $`J`$–, $`H`$–, and $`K`$–band observations of IRC +10 216 with the highest resolution so far at $`H`$ of 70 mas. A series of $`K`$–band images from five epochs between October 1995 and November 1998 shows that the inner nebula is non-stationary. The separations of the four dominant resolved components increased within the 3 years by almost $`40\%`$. For the two brightest components a relative velocity within the plane of the sky of about 23 mas/yr or 14 km/s was found. Within these 3 years the rather faint components C and D become brighter whereas component B is fading. The general geometry of the nebula seems to be bipolar.
We find that the most promising model to explain the structures and changes in the inner nebula is to assume that the star is at the position of component B. The star then is strongly but not totally obscured at $`H`$ and $`K`$. Consistently component B is very red in the $`HK`$ color while A and the northern $`J`$–band components are relatively blue. Similarly the polarization pattern with strong polarization in the northern arms and still a significant polarization in the peak supports this model. The inner nebula and the apparent motions seem to be rather symmetric around this position and the observed changes are consistent with an enhanced mass loss since 1997.
IRC +10 216 is without doubt in a very advanced stage of its AGB evolution. The observed bipolarity of its dust shell even reveals that it has possibly entered the phase of transformation into a protoplanetary nebula. |
warning/0002/nucl-th0002014.html | ar5iv | text | # KUNS-1636 Magnetic instability of quark matter
## 1 Introduction
Pulsars are rotating neutron stars emitting radio waves, X-rays or gamma rays. Ordinary radio pulsars have a magnetic field of $`O(10^{1213})`$G, which causes various radiations. The origin of such strong magnetic field is still an open problem. Recently a new type of neutron stars, called magnetars, has been proposed to explain the observational data on pulsars, which should have an extraordinary magnetic field of $`O(10^{15})`$ G $`^\mathrm{?}`$. There are reported several magnetar candidates so far for anomalous X-ray pulsars (AXP) and pulsars associated with soft-gamma-ray repeaters (SGR).
There has been a naive working hypothesis to understand the magnetic field in neutron stars; if the magnetic flux of a main sequence star is conserved during its evolution, the decrease in radius leads to an increase in the magnetic field. For example, the sun, a typical main sequence star, has a magnetic field of $`O(10^3)`$G with the radius $`R10^{1011}`$cm. By squeezing the radius to $`10^6`$cm for neutron stars we have $`O(10^{1113})`$G, which is consistent with observations for radio pulsars. However, if this argument is extrapolated to explain the intense of the magnetic field for magnetars, their radius should be $`O(10^4)`$cm, which is much less than the Schwartzschild radius of neutron stars with the canonical mass $`M=1.4M_{}`$, $`R_{Sch}=2GM/c^2=4\times 10^5`$cm.
These observations seems to enforce our reconsideration of the origin of the magnetic field in neutron stars. Since there is a bulk hadronic matter beyond the nuclear density ($`n_B0.16`$fm<sup>-3</sup>) inside neutron stars, it should be interesting to consider the hadronic origin of the magnetic field; ferromagnetism or spin-polarization of hadronic matter may give such magnetic field. Unfortunately there has been little suggestion about the possibility of spontaneous magnetization of hadronic matter. We consider here the possibility of ferromagnetism of quark liquid interacting with the one-gluon-exchange (OGE) interaction $`^\mathrm{?}`$.
One believes that there are deconfinement transition and chiral symmetry restoration at several times the nuclear density, while their critical densities have not been fixed yet. One interesting suggestion is that three-flavor symmetric quark matter (strange quark matter) around or above the nuclear density may be the true ground state of matter $`^\mathrm{?}`$. If this is the case, strange quark stars, where quarks occupy almost whole the inner region of stars, can exist in a different branch from the neutron-star branch in the mass-radius plane. Otherwise quark matter may exist in the small core region of neutron stars. We shall see our results should give an origin of the strong magnetic field in the context of strange quark-star scenario.
## 2 Ferromagnetism of quark liquid
Quark liquid should be totally color singlet (neutral), which means that only the exchange interaction between quarks is relevant there. This may remind us of electron system with the Coulomb interaction in a neutralizing positive charge background. In 1929 Bloch first suggested a possibility of ferromagnetism of electron system $`^\mathrm{?}`$. He has shown that there is a trade off between the kinetic and the exchange energies as a function of density, the latter of which favors the spin alignment due to the Pauli principle. This was a beginning of of the concept of itinerant magnetism. In the following we discuss the possibility of ferromagnetism of quark liquid on the analogy with electron gas.
It is to be noted that there is one big difference between them; quarks should be treated in a relativistic way. The concept of the direction of spins is not well defined in relativistic theories, while each quark has two polarization degrees of freedom. Here we define the spin-up and -down states in the rest frame of each quark. Then the projector onto states of definite polarization is given by
$$P(a)=\frac{1}{2}(1+\gamma _5a/)$$
(1)
with the 4-peudovector $`a`$,
$$𝐚=𝜻+\frac{𝐤(𝜻𝐤)}{m_q(E_k+m_q)},a^0=\frac{𝜻𝐤}{m_q}$$
(2)
for a quark moving with the momentum $`k=(E_k,𝐤)`$ $`^\mathrm{?}`$. The 4-peudovector $`a`$ is reduced into the axial vector $`𝜻`$ ($`|𝜻|=1`$) in the rest frame, which is twice the mean spin vector in the rest frame. Actually if we choose $`𝜻`$ along the $`z`$ axis, $`𝜻=(0,0,\pm 1)`$, we can see each value corresponds to the spin-up or -down state. The mean value of the spin is given by
$$\overline{𝐬}=\frac{1}{2}\frac{m_q}{E_k}\left(𝜻+\frac{𝐤(𝜻𝐤)}{m_q(E_k+m_q)}\right).$$
(3)
Finally the projection operator $`P(a)`$ gives the polarization density matrix $`\rho `$,
$$\rho (k,\zeta )=\frac{1}{2m_q}(k/+m_q)P(a),P(a)=\frac{1}{2}(1+\gamma _5a/).$$
(4)
The exchange interaction between two quarks with momenta $`𝐤`$ and $`𝐪`$ is given by
$$f_{𝐤\zeta ,𝐪\zeta ^{}}=\frac{m_q}{E_k}\frac{m_q}{E_q}_{𝐤\zeta ,𝐪\zeta ^{}}.$$
(5)
$`_{𝐤\zeta ,𝐪\zeta ^{}}`$ is the usual Lorentz invariant matrix element, and is evaluated with the help of the polarization density matrix (4)
$`_{𝐤\zeta ,𝐪\zeta ^{}}`$ $`=`$ $`g^2{\displaystyle \frac{1}{9}}\mathrm{tr}(\lambda _a/2\lambda _a/2){\displaystyle \frac{1}{4}}\mathrm{tr}\left[\gamma _\mu \rho (k,\zeta )\gamma ^\mu \rho (q,\zeta ^{})\right]{\displaystyle \frac{1}{(kq)^2}}`$ (6)
$`=`$ $`g^2{\displaystyle \frac{2}{9}}[2m_q^2kqm_q^2ab]{\displaystyle \frac{1}{(kq)^2}},`$ (7)
where the 4-pseudovector $`b`$ is given by the same form as in Eq. (2) for the momentum $`𝐪`$.
Although the vector $`𝜻`$ of each quark may point in a different direction on the two dimentional sphere $`S^2`$, we assume here it along the same direction, say $`z`$ axis. The exchange energy for quark liquid is then given by the integration of the interaction (5) over the two Fermi seas with the spin-up and -down states; eventually, it consists of two contributions,
$$ϵ_{ex}=ϵ_{ex}^{nonflip}+ϵ_{ex}^{flip}.$$
(8)
The first one arises from the interaction between quarks with the same polarization, while the second one with the opposite polarization. The non-flip contribution is the similar one as in electron gas, while the flip contribution is a genuine relativistic effect and never appears in electron gas. We shall see that this relativistic effect leads to a novel mechanism of ferromagnetism of quark liquid.
## 3 Examples
We show some results about the total energy of quark liquid, $`ϵ_{tot}=ϵ_{kin}+ϵ_{ex}`$, by adding the kinetic term $`ϵ_{kin}`$. Since gluons have not the flavor quantum numbers, we can consider one flavor quark matter without loss of generality. Then quark number density directly corresponds to baryon number density, if we assume the three flavor symmetric quark matter as mentioned in §1.
There are two parameters in our theory: the quark mass $`m_q`$ and the quark-gluon coupling constant $`\alpha _c`$. These values are not well determined so far. In particular, the value of quark mass involves subtle issues; it depends on the current or constituent quark picture and may be also related to the existence of chiral phase transition. Here we allow some range for these parameters and take, for example, a fiducial set, $`m_q=300`$MeV for strange quark and $`\alpha _c=2.2`$, given by the MIT bag model $`^\mathrm{?}`$. In Fig.1 two results are presented as functions of the polarization parameter $`p`$ defined by the difference of the number of the spin-up and -down quarks, $`n_q^+n_q^{}pn_q`$. The results clearly show that the ground state should be ferromagnetic for lower density, while it is in the paramagnetic phase for higher density. The phase transition is of first order and its critical density is around $`n_q^c0.16`$fm<sup>-3</sup> in this case, which corresponds to the nuclear density for flavor symmetric quark matter. Note that there is a metastable ferromagnetic state (the local minimum) even above the critical density. This ferromagnetic phase is a spontaneously symmetry broken state with respect to the rotational symmetry: the order parameter is the mean value of $`𝜻`$, $`𝜻`$, and symmetry is broken from $`G=O(3)`$ to $`H=O(2)`$ once $`𝜻`$ takes a special direction on $`S^2`$.
Magnetic properties of quark liquid are characterized by three quantities, $`\delta ϵ,\chi `$ and $`\eta `$; $`\delta ϵϵ_{tot}(p=1)ϵ_{tot}(p=0)`$, which is the measure for ferromagnetism to appear in the ground state. For small $`p1`$,
$$ϵ_{tot}ϵ_{tot}(p=0)=\chi ^1p^2+O(p^4).$$
(9)
$`\chi `$ is proportional to the magnetic susceptibility and plays an important role if the phase transition is of second order. In our case it is less relevant since the phase transition is of first order. Finally, $`\eta ϵ_{tot}/p|_{p=1}`$, which is the measure for metastability to to exist. In Fig.2 the density dependence of three quantities are given for a fiducial set of parameters. We can see that ferromagnetic phase is the ground state below $`n_q^c`$, while the metastable state is possible up to rather high densities.
Finally we present a phase diagram in the $`m_q\alpha _c`$ plane for $`n_q=0.3`$fm<sup>-3</sup>, which corresponds to about twice the nuclear density for flavor symmetric quark matter. The region above the solid line shows the ferromagnetic phase and that above the dashed line indicates the existence of the matastable state. For massive quarks with the large mass, which may correspond to the current $`s`$ quarks or the constituent quarks before chiral symmetry restoration, the ferromagnetic state is favored for small coupling constant due to the same mechanism as in electron gas. The ferromagnetic state is favored again for light quarks with small mass, which may correspond to the current $`u,d`$ quarks, while the nonrelativistic calculation does not show such tendency. Hence this is due to a genuine relativistic effect, where the spin-flip interaction plays an essential role.
## 4 Summary and Concluding remarks
We have seen that the ferromagnetic phase is realized at low densities and the metastable state is plausible up to rather high densities for a reasonable range of the QCD parameters. If a ferromagnetic quark liquid exists stably or metastably around or above nuclear density, it has some implications on the properties of strange quark stars and strange quark nuggets. They should be magnetized in a macroscopic scale. For quark stars with the quark core of $`r_q`$, simply assuming the dipolar magnetic field, we can estimate its strength at the surface $`R`$,
$$B_{max}=\frac{8\pi }{3}\left(\frac{r_q}{R}\right)^3\mu _qn_q,$$
(10)
amounts to order of $`O(10^{1517})`$G for $`r_qO(R)`$ and $`n_q=O(0.1)`$fm<sup>-3</sup> , which should be large enough for magnetars, using the quark magnetic moment $`\mu _q\mu _N`$($`\mu _N:`$nuclear magneton$`5\times 10^{24}\mathrm{erg}\mathrm{gauss}^1`$) for massive quarks and $`10^2\mu _N`$ for light quarks. Hence it might be interesting to model SGR or AXP using our idea.
We have found that ferromagnetic instability is feasible not only in the massive quark system but also in the light quark system: the spin-nonflip contribution is dominant in the nonrelativistic case as in electron gas, while a novel mechanism appears as a result of the large spin-flip contribution in the relativistic case.
Our calculation is basically a perturbative one and the Fermi sea remains in a spherical shape. However, if we get more insight about the ferromagnetic phase, we must solve the Hartree-Fock equation and thereby derive a self-consistent mean-field for quark liquid. Moreover, we need to examine the long range correlation among quarks by looking into the ring diagrams, which has been known to be important in the calculation of the susceptibility of electron gas.
## References |
warning/0002/math-ph0002049.html | ar5iv | text | # Classical and Quantum Probability
## 1 Introduction
The are few mathematical topics that are as badly taught to physicists as probability theory. Maxwell, Boltzmann and Gibbs were using probabilistic methods long before the subject was properly established as mathematics. Their language, of ensembles, complexions, fluctuations and most probable state, are still used. When quantum theory came along, the same notions were fitted into the new theory, sometimes leading to confusion. We review the mathematical development of probability, emphasising that quantum theory is a generalisation. The approach to history is in the same spirit as used by Milligan in
There are three ‘philosophies’ concerning probability. In the easy case, when there are finitely many possible outcomes to the experiment being considered, Laplace’s principle of equal ignorance tells us that the probability of each of the outcomes is the same. In the case of a die with six sides, experiments suggest that the probabilities are not all exactly equal. Nevertheless, there is not much error if we assume that the probability of each number is $`1/6`$. An objection to Laplace’s principle in general is that it is not always clear that the outcome of a particular experiment is a matter of chance, even when we do not know which outcome will turn up; it could even be that a particular outcome is inevitable. Thus a more robust version of Laplace’s principle might be that in events governed by chance, the probability of each possible outcome is the same. This still leaves open the meaning of the phrase, ‘governed by chance’. The difficulty of defining the ‘uniform’ distribution when variates take continuous values is illustrated by Bertrand’s paradox (, p. 246). This demolished Laplace’s principle for continuous variables.
The philosophy of Laplace applied to probability theory might be described as Platonic. A real die is the shadow of the ideal die, which has perfect sides and exact probabilities of $`1/6`$ for each outcome. This has a modern form of expression: we model the real die by the sample space $`\mathrm{\Omega }=\{1,2,\mathrm{},6\}`$ whose elements $`\omega `$ are called outcomes, and assign the probability $`1/6`$ to each. The value of a random variable $`f`$ is known if we know the outcome $`\omega `$; $`f`$ is therefore a real-valued function on $`\mathrm{\Omega }`$. More generally, if the sample space is a finite set $`\mathrm{\Omega }`$, an event $`E`$ is a subset of $`\mathrm{\Omega }`$; we say that the event has occurred if the $`\omega `$ that occurs lies in $`E`$. The probability that $`E`$ occurs is the sum of the probabilities of the points in E:
$$p(E)=\underset{\omega E}{}p(\omega ).$$
(1)
We say that two events, $`E`$, $`F`$, are independent if $`p(EF)=p(E)p(F)`$. In this way the binomial distribution can be derived for the total shown by $`n`$ dice thrown independently, and all of Laplace’s probability theory can be derived. It can tell us what bets to lay on an event $`E`$, even when only one trial is going to occur.
Laplace’s method has been successfully applied to statistical mechanics; the space of states is discretised, thus avoiding Bertrand’s paradox (the choice of bins being suggested by quantum mechanics). Each bin is said to be equally probable, and some hypotheses about independence is postulated. Then it is shown that the complexion (macroscopic state) given by the Gibbs distribution is not just the most probable, but is overwhelmingly the most probable. The chance of any complexion minutely different is put at $`10^{170}`$. The Gibbs distribution is, of course, the equilibrium state; if it is so probable, how come systems manage to be out of equilibrium, and remain so for years at a time? This remark is not aimed at Tolman , who made it clear that the assumption of equal probabilities applies only to equilibrium, and is to be tested against experiment; it passes the test well, but he then spoils it by adding as a further justification, ‘without this postulate there would be nothing to correspond to the circumstance that nature does not have any tendency to present us with systems in conditions which we regard as mechanically entirely possible but statistically improbable’. The word ‘improbable’ is itself based on Laplace’s assumption!
The second ‘philosophy’ of probability can be described as Aristotelian; it had taken hold by 1920, and is known as the ‘frequentist’ approach. It is essential that we can reproduce a long run of independent trials each conducted under exactly the same experimental conditions. In this respect, the theory makes sense only within a scientific culture. Suppose that we have one ‘variate’, which may take continuous or discrete values. The result of a measurement of the variate is assigned to one of a preassigned set of ‘bins’, which are intervals on the real axis. We repeat a number of times, to find the histogram, that is, the number $`n_i`$ of events (out of N trials) in the $`i^{\mathrm{th}}`$ bin. If the histogram settles down to a stable shape as we increase $`N`$, we declare that the value of the variate is random (or, random enough). We then define the probability of the event $`i`$ to be
$$p_i=\underset{N\mathrm{}}{lim}\frac{n_i}{N}.$$
(2)
This approach avoids the above problems that beset the Laplace philosophy. However, it is completely useless as mathematics; a ‘definition’ should not depend on an infinite number of future experimental results. There is not one theorem that can be proved from this definition. Feller points out that we must avoid confusion between a definition, and a method of measurement. There is great heuristic value to the frequentist approach. It is easy to teach ; we do not prejudge the possible values that the variate can have, or the probability of a given value; we can introduce another variate $`Y`$, and observe its distribution, and its joint probability distribution with $`X`$; we can by extending this idea get access to the joint probability distribution of any finite number of variates; we can get some idea as to whether the variates are random by examining a sequence of independent trials. We can even cover situations in which two variates are not simultaneously observable, as in quantum mechanics, by listing only the joint distributions of compatible observables, and omitting those we cannot measure. If we measure a variate $`X`$ with $`n`$ different values $`x_i`$ with relative frequency $`p_i`$, we can construct a sample space $`x_1,\mathrm{},x_n`$, and assign the probability $`p_i`$ to the occurrence of the outcome $`x_i`$. Similarly, we can construct a sample space and probability for any finite set of compatible variates if each measurement records their values. The observed probabilities are more reliable than assuming all points are equally probable.
However, there is one grave disadvantage of the approach, apart from not being mathematics: it is simply a description of data, and has much less predictive power than Laplace’s method. In particular, the method takes no position on the question as to what are the possible variates. If $`\mathrm{\Omega }`$ has $`|\mathrm{\Omega }|=n`$ points, then the random variables form a vector space, denoted $`𝒜(\mathrm{\Omega })`$, of dimension $`n`$, so that at most $`n`$ random variables can be linearly independent. No similar constraint holds in the frequentist point of view. Thus a variate is not the same as a random variable. In fact, it has no definition, other than the statement that its values are random.
The frequentist approach is the safest one to use in studies involving humans; social or financial matters are so complicated that it is not likely that a sample space, $`\mathrm{\Omega }_1`$ say, chosen to accommodate the data observed so far, can describe all the possible new variates and the values available to them. In the frequentist approach, faced with a new variate, $`Y`$, one simply takes the set of possible values of $`Y`$, say $`\mathrm{\Omega }_2=\{y_1,\mathrm{},y_m\}`$, and uses $`\mathrm{\Omega }_1\times \mathrm{\Omega }_2`$ as the sample space of the enhanced problem.
In a classical system in physics or chemistry, treated by classical statistical mechanics, we want to follow the scientific method: we model the system, do experiments, and reject the model if forced to. In that case, we make another model, estimate its parameters, and suggest more testing experiments. We want and expect to be able to make predictions about variates not measured yet. So we must reject the frequentist approach.
The third philosophy of probability was made clear by Kolmogorov, and combines something of the first two; it is to regard a probability theory as a model, to be tested against experiment. It is like Plato’s ideal, in that it is based on a specified sample space $`\mathrm{\Omega }`$; but now the probability $`p`$ is not determined by pure thought; any $`p`$ satisfying the axioms below provides us with a model.
###### Definition 1.1
Let $`\mathrm{\Omega }`$ be a countable space. A map $`p:\mathrm{\Omega }[0,1]`$ is a probability if $`p(\omega )0`$ and $`_\omega p(\omega )=1`$.
The probability of an event $`E\mathrm{\Omega }`$, and the concept of independence of two events, are then as in Laplace’s theory and clearly depend on the choice of $`p`$.
A random variable $`f:\mathrm{\Omega }𝐑`$ is chosen to represent the variate being observed, the particular choice being part of the interpretation of the model. A theoretical idea, or else the first few experiments on the variate $`f`$, allow us to get some guide-lines for $`\mathrm{\Omega }`$ and the values of $`p`$. This is the subject of estimation theory. We can judge the validity of the model (the choices we have made for $`(\mathrm{\Omega },p,f)`$) by comparing the predictions of the model with the observed frequencies $`n_i`$ using the theory of significance tests. Both estimation theory and significance were developed before Kolmogorov’s book. The founders of these techniques were often frequentists; they realised that one could not use an extreme frequentist point of view: in estimation, they often postulated that the data had Gaussian distributions, but with unknown parameters. In significance testing, to make a start, they assumed a probability distribution for the variate being measured; this is called the ‘hypothesis H’, which is part of the model; it can be rejected if the data are significantly unlikely. This has a version within Kolmogorov’s formulation, in which we are given a probability space, the pair $`(p,\mathrm{\Omega })`$, and model the variate with a random variable, $`f`$. To make contact with the well-established theory of estimation and significance, we must relate the probability distribution of $`f`$ to the probability $`p`$. We now remind the reader how this is done.
Given a finite probability space $`(p,\mathrm{\Omega })`$ and a random variable $`f:\mathrm{\Omega }𝐑`$, the probability distribution of $`f`$ is denoted $`p_f(i)`$, and is determined as follows: let $`x_i,i=1,2,\mathrm{},n`$ be the values that $`f`$ takes, and let $`p_f(i)`$ be the probability that the event $`\{\omega :f(\omega )=x_i\}`$. That is
$$p_f(i):=\underset{\omega :f(\omega )=x_i}{}p(\omega ).$$
(3)
This is what is accessible to experiments when we measure $`f`$. The mean of $`f`$ is determined by $`x_i`$ and $`p_f`$:
$$E_p[f]:=\underset{\omega }{}p(\omega )f(\omega )=\underset{i}{}x_ip_f(i),\text{ also written }p.f.$$
(4)
Given two random variables on $`(\mathrm{\Omega },p)`$, $`f,g`$ we define the joint distribution, denoted $`p_{f,g}(i,j)`$ to be
$$p_{f,g}(i,j):=p\{\omega :f(\omega )=x_i\text{ and }g(\omega )=y_j\}.$$
(5)
We say two r. v. are independent if the events $`\{f(\omega )=x_i\}`$ and $`\{g(\omega )=y_j\}`$ are independent for all $`i,j`$. This is equivalent to the frequentists’ version: $`p_{f,g}(i,j)=p_f(i)p_g(j)`$. The joint distribution determines $`p_f`$ and $`p_g`$ as its marginals, and also all moments, e. g. the cross-moment $`E_p[fg]`$ can be shown to be $`_{ij}x_iy_jp_{f,g}(i,j).`$
A probability $`p`$ defines a linear functional on the set $`𝒜(\mathrm{\Omega })`$ by the expectation, (4): $`fE_p[f]`$. We shall call any such functional a state: it is linear and positive, taking the value $`1`$ on the sure function $`I`$. The dual space $`𝒜^d`$ is the set of all linear functionals, so the states form a subset of $`𝒜^d`$; it is denoted $`\mathrm{\Sigma }(\mathrm{\Omega })`$. Given two states $`p_1`$ and $`p_2`$ and $`0<\lambda <1`$, their mixture with probabilities $`\lambda `$ and $`1\lambda `$, $`p=\lambda p_1+(1\lambda )p_2`$, is again a state. So the states form a convex set.
Whether $`\mathrm{\Omega }`$ is countable or not, for a random variable $`f`$ on $`(\mathrm{\Omega },p)`$ the probability of the occurrence of a single value $`f_0`$ might be zero, even when there is an $`\omega _0\mathrm{\Omega }`$ with $`f(\omega _0)=f_0`$; for, $`p(\omega _0)`$ might be zero. This often happens when $`\mathrm{\Omega }`$ is not countable, and $`f`$ takes continuous values. Then, more information about the probability measure is provided by the ‘cumulative’ distribution function
$$P_f(x)=p\{\omega :f(\omega )<x\}.$$
(6)
This is an increasing function of $`x`$, going from $`0`$ at $`x=\mathrm{}`$ to $`1`$ at $`x=\mathrm{}`$. We say that $`f`$ possesses a density $`\rho _f`$ if $`P_f(x)`$ is differentiable, and we write
$$\rho _f(x)=\frac{dP_f(x)}{dx}.$$
(7)
It is clear that we cannot cope with this subject without a certain amount of real analysis.
A cumulative probability distribution $`P_f(x)`$ is determined by its characteristic function
$$C_f(\lambda ):=e^{i\lambda x}𝑑P_f(x)$$
(8)
Here we use the Stieltjes integral. Any characteristic function satisfies
1. $`C(\lambda )`$ is continuous;
2. $`C(0)=1`$;
3. $`C`$ is of positive type:
$$\underset{ij}{}\overline{z}_iz_jC(\lambda _j\lambda _i)0.$$
Conversely, any function $`C`$ obeying $`1,2`$ and $`3`$ is the characteristic function of a probability distribution; this is Bochner’s theorem. In terms of the original $`(\mathrm{\Omega },p)`$ and random variable $`f`$, the characteristic function is
$$C_f(\lambda ):=E_p[e^{i\lambda f}].$$
(9)
If $`C_f`$ is analytic in $`\lambda `$ around $`\lambda =0`$, we can easily justify the expansion
$$C_f(\lambda )=\underset{n}{}(i\lambda )^nE_p[f^n]/n!=\underset{n}{}(i\lambda )^nM_n/n!.$$
here, $`M_n`$ are the $`n^{\mathrm{th}}`$ moments of $`f`$; for this reason, $`C_f`$ acts as a moment generating function for the r. v. $`f`$. An important variant of this is the cumulant generating function
$$\mathrm{log}C_f(\lambda )=\underset{n}{}(i\lambda )^n\kappa _n/n!.$$
We prefer to keep the imaginary unit in these formulas, since if we drop it the mean $`C_f(\lambda )`$ might not be finite. The cumulants $`\kappa _n`$ are determined by induction from the system
$$M_n=\underset{k}{}\underset{_k}{}\kappa _{n_1}\mathrm{}\kappa _{n_k}.$$
(10)
Here, $`_k=_1_2\mathrm{}_k`$ is an arbitrary partition of $`\{1,2,\mathrm{},n\}`$ into $`k`$ parts, including the identity partition, and $`n_j=|_j|,j=1,\mathrm{},k`$. The condition for independence, $`p_{f,g}(i,j)=p_f(i)p_g(j)`$ for all $`i,j`$ is equivalent to $`C_{f+g}=C_fC_g`$; it follows that then the cumulants of $`f+g`$ are the sums of those of $`f`$ and $`g`$.
For a Gaussian distribution, all the cumulants beyond the second are zero. There are results of the following kind: if all the cumulants $`\kappa _n`$ of a distribution are zero for $`nN`$, then they are zero beyond $`n=2`$, and so the distribution is Gaussian. These results use the positivity of the mean of a positive polynomial in $`f`$:
$$\underset{ij}{}\overline{z}_iz_jM_{i+j}=E[\underset{ij}{}\overline{z}_iz_jf^{i+j}]=E[|\underset{j}{}z_if^j|^2]0.$$
(11)
Given a set of real numbers $`\{M_n\}`$ satisfying the positivity condition in (11), it is not obvious that $`M_n`$ is the $`n^{\mathrm{th}}`$ moment of a random variable $`f`$, or that if so, $`f`$ is unique. This has led to a body of work called the moment problem.
The distribution of a random variable $`f`$ determines that of any differentiable function $`g(f)`$ of $`f`$; this is also a random variable; the density of the distribution of $`g`$ is determined by the usual rule: if $`g`$ is bijective, so that $`f`$ is a function of $`g`$, the probability that $`g`$ lies between $`y`$ and $`y+dy`$ is $`\rho _g(y)dy`$, and this occurs if and only if $`f`$ lies between $`x`$ and $`x+dx`$, where $`y=g(x)`$. Therefore $`\rho _gdy=\rho _fdx`$, giving the relation
$$\rho _g=|(df/dg)|\rho _f.$$
(12)
If $`g`$ is not bijective, but has a local inverse with various branches, $`f_i`$, then we have to sum over the contribution $`|(df_i/dg)|\rho _{f_i}`$ of each branch.
The remarkable thing is that the methods of probability theory give good results in many cases that are not governed by chance, such as the distribution of digits in $`\pi `$. Another example is the configuration of a chaotic system at the time $`t`$, where $`t`$ is large, given the initial configuration at time zero. If the initial state is not specified sufficiently accurately, then the configuration at time $`t`$ seems to be governed by chance, although it is not. It was suggested by Krylov that statistical physics is a successful method exactly in the cases when the underlying dynamics is chaotic. This will occur when nearby initial points become exponentially far apart as time progresses, and this is signalled by a positive real part to the dominant eigenvalue of the linearised dynamics. This largest real part is called the Lyapunov index. We are talking here about a chaotic theory; actual experimental measurements will always have further uncertainty, influenced by small effects omitted from the theory. In a non-chaotic system small forces can be omitted in the first few approximations. However, in a chaotic system, the inclusion of one such small force can change the outcome of the calculation at the large time $`t`$, making it appear to be random. This is well modelled by omitting any attempt to include all the actual forces, replacing those omitted by a ‘noise’, that is, a random term. Thus, we expect chaos to be well-modelled by a system with increasing uncertainty, as measured by entropy. Kolmogorov, and then Sinai took up Krylov’s cause, and were able to relate the rate of ‘entropy’ production to the Lyapunov exponent of the dynamics. However, Ruelle interprets this as an increase in information, available as time goes by.
Laplace’s problem, of whether to assign equal probabilities to each energy-level of a system, arises in quantum theory. Krylov takes von Neumann to task for assuming that the density matrix for the state of a particle with spin produced by a quantum process should, in the absence of any theory or experiment, be taken to be totally unpolarised. Krylov says that this is not true for most known processes, as the polarisation is found to be nonzero, small for some and large for others. Krylov’s view is that it should be assigned a general density matrix; we can then estimate this matrix in the light of experiments. This leads to the subject of quantum estimation, for which there is a body of theory. Krylov believed that physics is not in the gambling business; we do not second guess the state of the system and follow a strategy of hedging against wrong guesses; rather, in physics we predict what will happen (with various probabilities) at a later time, when the initial state is known.
Estimation theory has received an impetus from a modern development, information theory. Shannon introduced the entropy of the random variable $`f`$ taking values $`x_i`$ as
$$S_f:=\underset{i}{}p_f(x_i)\mathrm{log}p_f(x_i).$$
(13)
Note that $`S_f`$ does not depend on the actual values that $`f`$ takes. The distribution with the maximum possible entropy is easily proved to be the uniform distribution. The school of probability known as Bayesian therefore argues that if we know nothing whatever about $`f`$ it must be assigned the uniform distribution, called the prior. Thus, Laplace’s intuition gets very respectable support. There is one big problem with this: the uniform distribution for $`f`$ is not in general consistent with the uniform distribution for say $`g=f^3`$, as we see from eq. (12); so the prior depends on the random variable we choose to name as the one we know nothing about. This echoes Bertrand’s paradox.
A quantum version of entropy was earlier given by von Neumann. For the classical case $`(\mathrm{\Omega },p)`$ with $`\mathrm{\Omega }`$ countable it reduces to
$$S(p):=\underset{\omega }{}p(\omega )\mathrm{log}p(\omega ).$$
(14)
It does not make any reference to a random variable. We may obtain Shannon’s entropy of a random variable $`f`$ as the von Neumann entropy of $`p_f`$ regarded as a probability on the space of values that $`f`$ takes. Note that $`S_f=S(p)`$ if $`f`$ takes different values at different points of $`\mathrm{\Omega }`$, that is, if $`f`$ separates the points of $`\mathrm{\Omega }`$. We then say that $`f`$ is a sufficient statistic. $`S_f`$ is in general less than $`S(p)`$, and it reduces to zero when $`f`$ takes only one value. The entropies of Shannon and von Neumann are not the same concepts, and this difference reflects their different interpretations; the point $`\omega `$ is the message, and $`S_f`$ is the information about the message that is on average conveyed by measuring $`f`$; it cannot exceed $`S(p)`$, which is the entropy (missing information) in the original probability space. Naturally, if $`f`$ is sure it conveys no info at all. Since $`S_f`$ depends on the random variable $`f`$ only through its distribution, it has a meaning in the frequentist approach. To compute $`S(p)`$, the model $`(\mathrm{\Omega },p)`$ must be given, and it does not depend on $`f`$. More generally, we can define the Shannon entropy of a set $`(f_1,\mathrm{},f_n)`$ as the von Neumann entropy of their joint distribution on the sample space of their values. Some authors regard the Shannon entropy as the physical entropy of a reduced description of a physical model. The trouble with this idea is that the introduction of noise in the measurement of $`f`$ causes the Shannon entropy to decrease, instead of to increase as we would want.
A simple example of noise is that caused by a mapping $`T:\mathrm{\Omega }\mathrm{\Omega }`$. This defines a co-action on the set of random variables: $`fT^{}f:=fT`$, which in fact in an endomorphism of $`𝒜`$. If $`T`$ is not bijective, there might be points that can be distinguished by measuring $`f`$, but not by measuring $`T^{}f`$; thus $`S_{T^{}f}S_f`$. This also holds more generally, when $`T^{}`$ is a convex linear sum of such maps, thus: $`T^{}=\lambda _iT_i^{}`$. It can be shown that this is the most general stochastic map on $`𝒜`$, that is, linear map, taking $`I`$ to $`I`$ and non-negative functions to non-negative functions. The reduction in the information carried by $`f`$ in the presence of noise is natural in telephony. The von Neumann entropy, on the other hand, increases if we add noise. This is achieved by a bistochastic map $`T`$, (a stochastic map whose adjoint is also stochastic). We write it as a right action, thus: $`ppT`$. By the deep theorem of Birkhoff , a bistochastic matrix is a mixture of permutations. Since a permutation $`\mathrm{\Omega }`$ does not alter $`S(p)`$, and $`p\mathrm{log}p`$ is concave, we see that $`S(pT)S(p)`$. Moreover, the von Neumann entropy is not decreased by a reduced description, unlike the Shannon version. Thus $`S(p)`$ is the correct concept to represent physical entropy .
If we are given information about which $`\omega `$ has occurred, the probability on $`\mathrm{\Omega }`$, called the prior, changes. Suppose that $`p`$ is the prior. If the information is that the sample lies in a known subset (event) $`\mathrm{\Omega }_0\mathrm{\Omega }`$, then Bayes’s theorem on conditional probability is used; the conditional probability is
$$p(E|\mathrm{\Omega }_0):=\frac{p(E\mathrm{\Omega }_0)}{p(\mathrm{\Omega }_0)}.$$
(15)
This is called the posterior probability, and correctly describes the probability among outcomes all of which lie in $`\mathrm{\Omega }_0`$. A conditional probability satisfies the axioms of probability, def. (1.1).
This use of information to modify the probability $`p`$ should not be confused with estimation theory. There, we do not change $`p`$, since after the measurement of independent samples, we continue to assume that new samples are governed by the original $`p`$. The method of estimation using the principle of maximum entropy proceeds as follows. Suppose that we know $`\mathrm{\Omega }`$, and $`f`$, with $`|\mathrm{\Omega }|<\mathrm{}`$ and we are also told the average of $`f`$ over a number of independent trials. We can vary $`p`$ over the simplex $`\mathrm{\Sigma }(\mathrm{\Omega })`$ to find the point that maximises the entropy of the probability, given the observed mean value, $`\eta `$ say. Thus we use the method of Lagrange multipliers to maximise
$$\underset{\omega }{}p(\omega )\mathrm{log}(p(\omega ))\text{ subject to }E_p[f]=\eta .$$
Gibbs knew that the solution to this is
$$p(\omega )=Z^1\mathrm{exp}(\beta f(\omega )),$$
(16)
where $`Z=_\omega e^{\beta f(\omega )}`$, is the Lagrange multiplier for the normalisation condition $`p(\omega )=1`$, and is called the partition function. The parameter $`\beta `$ is the Lagrange multiplier for the condition $`p.f:=E_p[f]=\eta `$, and is determined by it. Then (16) is the least prejudiced estimate for the probability, given the mean .
The method of maximum entropy solves an important problem in the theory of estimation. Let $`X`$ be a variate of which the distribution is known to be one of a family, $`=\{p_\eta (i)\}_{\eta 𝐑}`$; we hope to estimate $`\eta `$ by measuring $`X`$ independently $`m`$ times. An estimator $`f`$ is a function of the data $`x_1,x_2,\mathrm{},x_m`$ that is used for this estimate. Thus $`f`$ is a function of $`X`$, and so is a random variable. Since we do not know $`\eta `$, to be useful, the estimator must be independent of $`\eta `$. We say an estimator is unbiased if its mean is the desired parameter, thus:
$$p_\eta .f:=\underset{i}{}p_\eta (i)f(x_i)=\eta .$$
(17)
Apart from being unbiased, a good estimator should have a small chance of being far from the mean; so we are interested in estimators of minimum variance, $`V=p_\eta .[(f\eta )^2]`$. To any probability $`p`$ define the Fisher information as
$$G=p_\eta .\left(\frac{\mathrm{log}p_\eta }{\eta }\right)^2$$
(18)
We recognise this as the variance of the random variable $`Y=p/\eta `$. The Cramer-Rao theorem puts limits on the smallness of the variance $`V`$ of an estimator $`f`$:
###### Theorem 1.2
$$VG^1.$$
(19)
For the proof, differentiate (17) with respect to $`\eta `$, to get
$$\underset{i}{}\frac{p_\eta (i)}{\eta }f(x_i)=1.$$
Now use $`/\eta [_ip_\eta (i)]=0,`$ and rearrange, to get
$$\underset{i}{}p_\eta (i)\left(\frac{\mathrm{log}p_\eta (i)}{\eta }(f(x_i)\eta )\right)=1.$$
(20)
This is the correlation between the random variables $`Y`$ and $`f`$; the (positive-semidefinite) covariance matrix is therefore
$$\left(\begin{array}{cc}G& 1\\ 1& V\end{array}\right)$$
(21)
Schwarz’s inequality then gives (19).
The minimum variance allowed by (19) occurs when the Cauchy inequality is equality, which occurs when the factors in (20) are proportional (with ratio dependent on $`\eta `$). Calling this factor $`\xi /\eta `$, we see that the distribution of minimum variance must satisfy
$$\mathrm{log}p_\eta (i)=^\eta \xi /\eta (f(x_i)\eta )d\eta =\xi f(x_i)\psi ,$$
(22)
showing that a necessary and sufficient condition is that $`\{p_\eta \}`$ be the exponential family.
In the case of several parameters $`\eta _1,\mathrm{}\eta _n`$, we have estimators $`f_1,\mathrm{},f_n`$, which can be taken to be linearly independent, but need not be functionally independent. The state of maximum entropy, given the means $`E_p[f_i]=\eta _i,i=1,\mathrm{},n`$ is easily shown to be of the form
$$p(\omega )=Z^1\mathrm{exp}\{\xi _1f_1(\omega )+\mathrm{}\xi _nf_n(\omega )\}=Z^1\mathrm{exp}(f)\text{ say},$$
(23)
where the Lagrange multipliers $`\xi _i`$ are determined by the given conditions on the means. The set of probabilities of the form eq. (23) form the set $``$ called the info manifold, or the exponential family, determined by Span$`\{f_i\}`$. We can regard $`\{\xi _i\}`$ or indeed $`f`$ as coordinates, called canonical; or we can regard $`\{\eta _i\}`$ or indeed $`p`$ as coordinates, called the expectation coordinates. In this case the Fisher information matrix is defined to be
$$G^{ij}:=p.\left(\frac{\mathrm{log}p}{\eta _i}\frac{\mathrm{log}p}{\eta _j}\right).$$
(24)
Then the Cramer-Rao inequality (19) becomes a matrix inequality, where $`V`$ is the covariance matrix $`V_{ij}:=p.[(f_i\eta _i)(f_j\eta _j)]`$. Equality holds only if $`G=V^1`$, which leads to the exponential family.
Rao showed that $`G`$ defines a Riemannian metric on the tangent spaces of $``$ ; as such, its components depend on the coordinates chosen for the tangent space and it transforms as a tensor under changes in variables. At the point $`p`$, a vector in the tangent space is given in in canonical coordinates by a random variable $`f`$ in the span of the ‘score variables’ $`\widehat{f}_j:=f_i\eta _j`$. Writing $`f=_k\xi ^k\widehat{f}_k`$ introduces contravariant components $`\xi ^k`$. These are dual to the $`\eta _j`$, which are covariant components. The covariant metric is the covariance matrix
$$G_{ij}=G(\widehat{f}_i,\widehat{f}_j)=E_p[\widehat{f}_i\widehat{f}_j].$$
(25)
It is the inverse of the contravariant $`G^{ij}`$, which explains why we get equality in (19).
The Massieu function $`\psi :=\mathrm{log}Z`$, where $`Z`$ is the partition function, is related to the free energy; it is the generating function for the cumulants; so we have
$`\eta _j`$ $`=`$ $`{\displaystyle \frac{\psi }{\xi ^j}}`$ (26)
$`G_{ij}:=V_{ij}`$ $`=`$ $`{\displaystyle \frac{^2\psi }{\xi ^i\xi ^j}}.`$ (27)
The entropy is the Legendre transform of $`\psi `$, and its second variation is the Fisher information matrix, $`G^{ij}`$, the metric in the coordinates $`\eta `$.
Amari showed that $``$ is furnished with a pair of affine flat connections, for which the global affine coordinates are $`\xi ^i`$ and $`\eta _i`$ . These connections are not metric connections, but are dual relative to $`G`$. An important role in information geometry is played by the relative information $`S(p|p^{}):=_\omega p(\omega )(\mathrm{log}p(\omega )\mathrm{log}p^{}(\omega ))`$. This distinguishes between the points $`p`$ and $`p^{}`$ in $``$, in that $`S(p|p^{})0`$ and vanishes only when $`p=p^{}`$. For a modern version, see .
The observables form the algebra $`𝒜(\mathrm{\Omega })`$ in which multiplication is pointwise: $`(fg)(\omega ):=f(\omega )g(\omega )`$; the states lie in its dual. Thus, states and observables are not the same kind of thing, and they transform as duals under stochastic maps. However, states like observables are functions of $`\omega `$; to distinguish them we can write $`p(\omega )`$ for a state and $`(\omega )f`$ for an observable. If $`|\mathrm{\Omega }|<\mathrm{}`$, either can be identified with an element of the formal vector space spanned by $`\mathrm{\Omega }`$, thus: $`_\omega \alpha (\omega )\omega \alpha `$, whether $`\alpha `$ is regarded as an observable or a state. Then $``$ is the interior of the convex hull of $`\mathrm{\Omega }`$. The permutation group of $`\mathrm{\Omega }`$ acts by right action $`\omega \omega T`$. Its inverse $`\omega \omega T^1`$ is a co-action of the group (its product law is the opposite of that of the group) and so can be written as a left action: $`\omega T\omega :=\omega T^1`$. These induce a right action on probabilities, and a left action on observables, by $`pT(\omega ):=p(T\omega )`$ and $`(\omega )Tf:=(\omega T)f`$, the latter written without the dual symbol . These express associativity, as does the dual relation $`pT.f=p.Tf`$.
These definitions can be extended to any map $`T:\mathrm{\Omega }\mathrm{\Omega }`$, whether invertible or not: we define the action on probabilities using $`T(\omega ):=(\{\omega \}T^1)`$, the inverse image of the point-set $`\{\omega \}`$. Every algebraic endomorphism of $`𝒜`$ is of the form $`fTf`$ for some map $`T:\mathrm{\Omega }\mathrm{\Omega }`$, and these make up exactly the extreme points of the convex set of stochastic maps.
In infinite dimensions, there is more than one useful topology on the states and observables. The modern view is that the state $`p`$ and the observable $`\mathrm{log}p`$ are merely alternative coordinates for a point in the info manifold. The natural class of charts are related by monotone, convex functions, of which the stochastic maps, , as well as the non-linear maps $`p\mathrm{log}p`$ and $`pp^\alpha `$, $`0<\alpha <1`$ are examples.
An active field of research is to set up quantum analogues of all this
## 2 From Bachelier to Wiener
In 1900, Bachelier proposed a random model of the stock market ; the idea was that the decision to buy or sell a stock is randomly taken by independent investors. Let us suppose that the chance $`\lambda `$ that the price goes up one unit $`dx`$ is the same as that for going down, during any unit trading period $`dt`$. Let $`X𝐙`$ be the random price, and $`p(x,t)`$ be the probability that the price is $`x`$ at time $`t`$; then the new probability $`p(x,t+dt)`$ can be unchanged, or can change due to a movement down from $`x+dx`$ or a movement up from $`xdx`$. The probabilities of these are, respectively, $`12\lambda `$, $`\lambda `$ and $`\lambda `$. Thus we get the relation
$$p(x,t+dt)=(12\lambda )p(x,t)+\lambda p(x+dx,t)+\lambda p(xdx,t).$$
(28)
Let $`T`$ be the tridiagonal infinite matrix $`\{\lambda ,12\lambda ,\lambda \}`$. Then the row $`T_{xy},y𝐙`$ is the conditional probability, that the price will be $`y`$ at time $`(N+1)dt`$, given that it is $`x`$ at time $`Ndt`$. In fact, $`T`$ is a stochastic matrix, which happens to be symmetric.
Suppose that at $`t=0`$ the price is $`x_0`$; then in time $`Ndt`$, the price will follow the path $`\gamma :=x_0x_1\mathrm{}x_N`$ with probability
$$p(\gamma )=T(x_0,x_1)T(x_1,x_2)\mathrm{}T(x_{N1},x_N).$$
(29)
This is called the random walk on $`𝐙`$ determined by $`T`$, starting at $`x_0`$. The set of allowed paths starting at $`x_0`$ is a finite subset of $`\mathrm{\Omega }=𝐙^N`$. $`p(\gamma )`$ is a probability on $`\mathrm{\Omega }`$, and the structure is called a Markov chain. An alternative point of view is to start with $`p_0\mathrm{\Sigma }(𝐙)`$, and to follow the path in $`\mathrm{\Sigma }(𝐙)`$ given by the time evolution. By Bayes’s law, the probability that at time $`t=1`$ the particle is at $`x_1`$ whatever its initial position, is $`p(x_1,1)=_{x_0}p(x_0,0)T(x_0,x_1)`$; this can be written as the matrix product $`p_0T`$, where $`p_t`$ is a row vector made from the components of $`p(x_t,t)`$. By induction, the probability that at time $`N`$ the particle is at $`x`$ is $`p_0T^N`$. In this way, a Markov chain is described by a semi-group of stochastic maps $`T(t):=T^t`$ acting on $`\mathrm{\Sigma }(\mathrm{\Lambda })`$. Obviously
$`T(0)`$ $`=`$ $`1`$ (30)
$`T(s)T(t)`$ $`=`$ $`T(s+t),s,t𝐍.`$ (31)
One of the themes of probability theory is the relationship between a semi-group of stochastic maps and a probabiliy on the corresponding path space. The latter is called a dilation of the former. Since $`T`$ is independent of time, the chain is said to be stationary; if we limit the allowed space to be a finite set $`\mathrm{\Lambda }𝐙`$, we get a finite Markov chain, in which case there is at least one stationary distribution $`p^{}`$; this means that $`p^{}T=p^{}`$, so that 1 is a left eigenvalue of $`T`$. If some power of $`T`$ has all its matrix elements positive, then the Perron-Frobenius theorem tells us that $`1`$ is a simple eigen-value, and all the others have modulus less than $`1`$. One can then show that $`pT^np^{}`$ as $`n\mathrm{}`$; the system converges exponentially to equilibrium. We then say that the dynamics is mixing. There are similar results in infinite dimensions, but to get exponential convergence we need to show that there is a spectral gap. This means that $`1`$ is simple and lies a finite distance from the next eigenvalue of $`TT^{}`$. To prove this in the case at hand is usually the key to the study of the long-time behaviour. The Markov property is that the probability of getting to $`x`$ at time $`t+1`$ depends only on where the particle was at time $`t`$, and not on the previous path. The study of Markov chains was started in the $`19^{\mathrm{th}}`$ century, and is a huge subject.
Fick obtained an equation similar to (28) for the diffusion of particles in one dimension. If $`dx`$ and $`dt`$ become small such that $`(dx)^2/dta`$, a finite limit, we say the system is following the diffusion limit. Rearranging, and taking the diffusion limit, Fick obtained the heat equation for the probability density, which we call $`\rho `$:
$$\frac{\rho }{t}=\kappa \frac{^2\rho }{x^2}.$$
(32)
Here, $`\kappa =a\lambda `$. This is not a very good model of the market; apart from the omission of drift, the gains in price should grow with the overall price. As it is, negative prices are possible.
The heat equation (32) can be written in the form of a conservation law:
$$\frac{\rho }{x}+\text{div}j(x,t)=0$$
(33)
where $`j(x,t)=\kappa \rho `$. At this stage, mathematicians did not have the continuous version of the sample space $`\mathrm{\Omega }`$; this was to be Wiener’s great construction.
In his celebrated work of 1905 , Einstein also used (32) to describe the Brownian motion of small particles in a warm liquid. He was mindful of Stokes’s law of diffusion; this says that in a viscous liquid a small particle under a constant force, such as gravity, will increase its speed towards a terminal velocity $`v`$ say, which is proportional to the force. Einstein required that in equilibrium the current $`v\rho `$ due to this flow should balance the diffusion due to the density gradient, so that steady state should obey
$$\kappa \rho +v\rho =0.$$
(34)
The solution to this in the case of gravity, where $`v=|v|`$ in the $`z`$-direction, is
$$\rho (x,y,z)=\text{const.}e^{|v|z/\kappa },$$
(35)
and this should be the Maxwell-Boltzmann law at the temperature $`\mathrm{\Theta }`$ of the liquid,
$$\rho (x)=Z^1e^{mgz/(k_B\mathrm{\Theta })}.$$
(36)
Einstein thus obtained the famous Einstein relation
$$F=k_B\mathrm{\Theta }v/\kappa .$$
(37)
His treatment is not complete, since he omitted the drift term in the diffusion equation! See §4, (1) in . In a detailed study, Smoluchowski wrote down the diffusion equation with drift
$$\frac{\rho }{t}=\kappa \frac{\rho }{x^2}v\frac{\rho }{x}.$$
(38)
now known as the Smoluchowski equation, and is a special case of the Fokker-Planck or backward Kolmogorov equation. He solved this by using the method of images for several systems with boundaries, such as the mass of air above the ground, and obtained the approach to the stationary state expected by Einstein.
It was known that one can solve eq. (38) exactly, to fit a more or less arbitrary initial function $`\rho (x,0)=f(x)`$, by using the Green function (in one dimension)
$$G(x,t):=[4\pi \kappa t]^{(1/2)}e^{(xvt)^2/(4\kappa t)}.$$
(39)
This satisfies eq. (38), and converges in the sense of distributions to the Dirac $`\delta `$-function as $`t0`$. Then
$$\rho (x,t)=_{\mathrm{}}^{\mathrm{}}G(xy,t)f(y)𝑑y$$
(40)
satisfies eq. (38) and the boundary condition. The operator whose kernel is $`G`$ is the continuum analogue of the matrix $`T^n`$ of the Markov chain,
When the force and temperature are slowly varying, we get the coupled system
$`{\displaystyle \frac{\rho }{t}}`$ $`+`$ $`\text{div}J=0;`$ (41)
$`{\displaystyle \frac{\mathrm{\Theta }}{t}}`$ $`=`$ $`\kappa ^{}\text{div}\mathrm{\Theta }+\kappa 𝐅.𝐅/k_B\mathrm{\Theta }.`$ (42)
Here, $`J(x,t)=\kappa (\rho +V\rho /\mathrm{\Theta })`$, where $`V`$ is the potential giving rise to the force $`F`$. The source term in the heat equation is $`F.J`$, the power of the external force supplied to the particle, all of which is converted into heat. This system obeys the first and second laws of thermodynamics .
Consider now the solution (40) to (38). Because $`G`$ is positive, the density remains positive for all time, and the conservation law shows that the integral of $`\rho `$ over space is constant. So we get a flow through the space of probabilities. The question arises, is there a process in continuous time associated with the Smoluchowski equation? The answer is yes, and this was the result of the work of Wiener, and later, Ito. An alternative idea was introduced by Langevin, who considered Newton’s laws, in which a part of the external force, denoted $`F`$, is random; friction enters as a damping force proportional to the velocity, parametrised by $`\gamma >0`$. Thus his equation is
$$\frac{d^2x}{dt^2}=\frac{V}{x}\gamma \dot{x}+F(t).$$
(43)
This is the equation for a single particle, but as $`F`$ is random, the position $`x(t)`$ becomes random as time goes by, even if its initial condition is given. Statistical properties of $`x`$ are determined by those of $`F`$; the relation of these to the Smoluchowski equation were studied by Fokker and Planck, but were fully understood only in terms of stochastic calculus. One might assume that $`F`$ is Gaussian distributed, and is of mean zero, with independent values at different times. This would now be described as white noise. Langevin’s work started the enormous field of stochastic differential equations.
In 1904 Lebesgue tried to set up a general theory in which every subset of $`[0,1]`$ is assigned a measure. . The very next year, G. Vitale showed that the scheme was inconsistent . Hausdorff and Banach and Tarski, showed that the measure could not be additive . The point is that some sets are so bad they cannot be assigned a measure, even a finitely additive one. This led to the concept of measurable set. Let us start with the Borel measurable sets on $`[0,1]`$.
Let $`\mathrm{\Omega }=[0,1]`$; let us say that a collection $``$ of subsets of $`\mathrm{\Omega }`$ form a tribe if
1. $`\mathrm{\Omega }`$
2. whenever $`B`$, we have $`B^c:=\mathrm{\Omega }B`$;
3. whenever $`A`$ and $`B`$, we have $`AB`$.
Such a collection of subsets is also called a Boolean ring, or a Boolean algebra. The collection $``$ is actually a ring, with multiplication given by intersection, and addition given by symmetric difference, that is $`A+B:=ABAB`$. It is also an algebra in the technical modern sense, but trivially in that any ring is an algebra over the field consisting of two numbers, $`0`$ and $`1`$. Since this ring structure plays no role in the theory, we prefer not to furnish $``$ with the extra structure ‘+’, and will use the word ‘tribe’ instead.
We define a $`\sigma `$-tribe to be a collection $``$ of sets $`B_i\mathrm{\Omega }`$ such that 3. above is replaced by
$`\mathrm{\hspace{0.33em}\hspace{0.33em}3}_{\mathrm{}}.\text{ if }B_i\text{ is a countable family of disjoint sets, then }B_i.`$
The set of all subsets of a set $`\mathrm{\Omega }=[0,1]`$ is obviously a $`\sigma `$-tribe, and indeed satisfies uncountable additivity as well. This $`\sigma `$-tribe is called the power set of $`\mathrm{\Omega }`$. But, as we saw, there are no useful definitions of measure on the power set. Another easy case is the collection of all countable subsets of $`\mathrm{\Omega }`$: the union of a countable collection of countable subsets is countable. However, any countable set has length zero, since it can be covered by a sequence of intervals of length $`ϵ/2,ϵ/4ϵ/8,\mathrm{}`$, of total length $`ϵ`$. Since $`ϵ`$ can be anything, the set has length zero. To get some sets of non-zero length, let us consider the tribe $`_0`$ of all finite disjoint unions of open, closed and half-open intervals. We could add to $`_0`$ all countable unions of sets in $`_0`$, and all complements in $`\mathrm{\Omega }`$ of sets in the tribe so obtained. Call this $`_1`$. Then we would need to consider the collection of countable unions of sets in $`_1`$, and their complements, to get a new tribe $`_2`$, and so on. Does this end up with a well-defined $`\sigma `$-tribe? The following argument does the trick. Let $`𝒢`$ be any $`\sigma `$-tribe containing all sets in $`_0`$, and let $`C`$ be the set of all such $`\sigma `$-tribes. Then $`C`$ is non-empty, as it contains the power set at least. Then form
$$=\underset{𝒢C}{}𝒢.$$
(44)
That is, $``$ contains those subsets of $`\mathrm{\Omega }`$ that lie in all $`\sigma `$-tribes $`𝒢`$, and no other subsets. In particular, $``$ contains all subsets in $`_0`$, $`_1`$ etc. In fact, by using the techniques of set theory, one can prove that $``$ is smallest $`\sigma `$-tribe containing all the open intervals in $`\mathrm{\Omega }=[0,1]`$; it is called the Borel tribe. One can ask whether we have arrived at the power set after all, or have something without the pathological sets. That $``$ contains only nice sets follows the construction of a countable measure on its sets, namely the Lebesgue measure.
A finitely additive measure on a tribe $``$ is a map $`\mu :𝐑_+\{+\mathrm{}\}`$ such that
$$\mu (AB)=\mu (A)+\mu (B)\text{ for all disjoint }A,B.$$
If $`\mu (\mathrm{\Omega })=1`$, it is a finitely additive probability measure. To do analysis, we must be able to take some limits, and so we now assume that $``$ is a $`\sigma `$-tribe.
A probability measure on $`(\mathrm{\Omega },)`$ is a map $`\mu :𝐑^+`$ such that
1. $`\mu (B)0`$ for all $`B`$;
2. $`\mu (\mathrm{\Omega })=1`$;
3. if $`B_i`$ is a countable collection of disjoint sets in $``$, then
$$\mu (B_i)=\underset{i}{}\mu (B_i).$$
Considering the tribe $`_0`$ of finite unions of disjoint open, closed and half-open intervals, we can define the Lebesgue measure of $`B_0`$ to be the sum of the usual lengths of the intervals involved. It is then proved that there is a countably additive measure on the Borel $`\sigma `$-tribe, which agrees with the length on the intervals. This measure is called the Lebesgue measure.
It is sometimes useful to extend the concept of measure to unbounded sets such as $`𝐑`$, whose total length is infinite. For this, we just drop axiom 2. above.
So much for the measure; integration theory needs a remark as well. Suppose that we have a function $`y=f(x)`$, where $`x[0,1]`$ and $`y`$ is real-valued and bounded, and we seek a way of finding the area under the graph of $`y`$ against $`x`$. In Riemann’s method of integration we divide the $`x`$-axis into a large number of small intervals, $`[0,x_1],(x_1,x_2],\mathrm{},(x_N,1]`$, and define $`y_i`$ to be the smallest value of $`y`$ in the interval $`(x_i,x_{i+1}]`$ and $`Y_i`$ to be the largest value. Now define the two approximations to the area, known as the upper sum and the lower sum, $`R^+=_iY_i(x_{i+1}x_i)`$ and $`R^{}=_iy_i(x_{i+1}x_i)`$. As we refine the subdivision, $`R^+`$ decreases and $`R^{}`$ increases. If the limits of these are equal, we say that the function is Riemann-integrable, and take their common value as the area under the curve $`y=f(x),\mathrm{\hspace{0.33em}\hspace{0.33em}0}x1`$. One shows that continuous functions are integrable, and can establish the fundamental theorem of the calculus; a generalisation, called the Riemann-Stieltjes integral, can be defined, if we replace $`x_{i+1}x_i`$ by $`P(x_{i+1})P(x_i)`$, where $`P`$ is an increasing function of bounded variation, continuous from the left. We write the integral as $`y(x)𝑑P(x)`$. To define the integral of unbounded functions, various limiting methods were invented. The theory is not really satisfactory.
Lebesgue introduced a new form of integration: compared with Riemann’s method, it is done the other way round. As the first step, only positive functions are considered. Then, we divide the $`y`$-axis into intervals $`([0,y_1],(y_1,y_2],\mathrm{},(y_N,\mathrm{}))`$, and for each interval, look for the inverse image of each interval under the map $`f`$. That is, we consider the subset of the $`x`$-axis consisting of $`x`$ such that $`f(x)(y_i,y_{i+1}]`$. This set, denoted by $`f^1(y_i,y_{i+1}]:=B_i`$, may consist of many pieces, and so will not always be an interval. We require, however that it should be a set in the Borel $`\sigma `$-tribe $``$; if this holds for every subdivision of the $`y`$-axis into intervals, we say that the function $`f`$ is $``$-measurable. The set $`B_i`$ will have a ‘length’, namely, its Lebesgue measure, $`\mu (B_i)`$. We approximate the area under the graph of $`f`$ by the sum
$$L(f):=\underset{i}{}y_i\mu (B_i).$$
This is positive and increases as we refine the partition of the $`y`$-axis. If its supremum over all partitions is finite, we say that $`f`$ is Lebesgue-integrable, and write
$$_0^1f(x)𝑑x=supL(f).$$
(45)
We can integrate functions that are not positive, provided that the positive and negative parts are separately integrable, and we integrate complex functions by treating the real and imaginary parts separately. This generalises the Riemann integral in that any Riemann-integrable function is Lebesgue integrable, and then both versions give the same answer.
Lebesgue integration has the following easy generalisation, which is important for probability. Suppose that $`\mathrm{\Omega }`$ is any set, provided with a $`\sigma `$-tribe $``$; the pair $`(\mathrm{\Omega },)`$ is called a measurable space. A real-valued function is said to be $``$-measurable if the inverse image of every open interval lies in $``$:
$$f^1(y_1,y_2):=\{\omega \mathrm{\Omega }:y_1<f(\omega )<y_2\}.$$
A random variable is then simply a real-valued $``$-measurable function on $`\mathrm{\Omega }`$. Given a measure $`\mu `$ on $`(\mathrm{\Omega },)`$, not necessarily of finite total measure, we can regard as the same random variable $`f`$ two that differ only on a set of $`\mu `$-measure zero; they are called versions of $`f`$. The set of all bounded random variables forms a commutative algebra $`𝒜(\mathrm{\Omega })`$ with norm $`f_{\mathrm{}}:=infsup_\omega |f(\omega )|`$; here, $`inf`$ is taken over all versions of $`f`$. The sets in $``$ are called events. The integral of a positive measurable function (with respect to the measure $`\mu `$) is defined similarly to the case when $`\mathrm{\Omega }=𝐑`$. If $`\mu `$ is a probability measure, this integral is called the mean $`\mu .f`$ of $`f`$ in the state $`\mu `$. and if $`\mu .|f|<\mathrm{}`$ we write $`fL^1(\mathrm{\Omega },,\mu )`$. More generally, we write $`fL^p(\mathrm{\Omega },,\mu )`$, $`1p<\mathrm{}`$ if $`f`$ is $``$-measurable and $`|f|^p`$ is integrable. These are Banach spaces with norm $`f_p:=(|f(\omega )|^p𝑑\mu )^{1/p}`$. The probability of an event $`B`$ is taken to be $`\mu (B)`$. Each measure $`\mu `$ defines an element of the dual space of $`𝒜`$ by the linear form $`ff𝑑\mu `$.
We have remarked that the original motivation for introducing the $`\sigma `$-tribe was to avoid pathology. However, the concept has been very useful in a heuristic way, to describe the information carried by events and observables in a random theory based on a measure space $`(\mathrm{\Omega },,\mu )`$; in particular, it is useful to consider a sub-tribe or sub-$`\sigma `$-tribe, of $``$. Suppose that $`B`$ is an event; it is determined by its indicator function, $`\chi __B(\omega )`$ which is $`1`$ if $`\omega B`$ and zero outside $`B`$. If $`\mu (B)0,1`$ and $`f`$ is measurable, we can define the conditional expectation
$$E[f|B]=\underset{\omega }{}f(\omega )\mu (\omega |B).$$
(46)
We may also find the conditional probability of $`A`$, given that $`B`$ did not happen: $`\mu (A|B^c)=\mu (AB^c)/\mu (B^c)`$, and the corresponding conditional expectation
$$E[f|B^c]=\underset{\omega }{}f(\omega )\mu (\omega |B^c).$$
(47)
We may regard the pair of numbers, $`\{E[f|B],E[f|B^c]\}`$ as defining a simple measurable function on $`\mathrm{\Omega }`$, equal to $`E[f|B]`$ if $`\omega B`$ and to $`E[f|B^c]`$ if $`\omega B`$. Let us now generalise this idea. Let $`B_1,\mathrm{},B_n`$ be disjoint measurable sets such that $`\mu (B_j)0`$ for all $`j`$, and $`\mu (_jB_j)=1`$. These sets generate a tribe, say $`_0`$ (by various unions; there are $`2^n`$ such unions). If $`f`$ is measurable, the functions on $`\mathrm{\Omega }`$ defined by
$$F_f(\omega )=E[f|B_j]\text{ if }\omega B_j$$
(48)
are measurable relative to $`_0`$. They take constant values, $`E[f|B_j]`$ on each $`B_j`$ and so can be written
$$F_f(\omega )=\underset{j}{}\chi _j(\omega )c_j\text{ where }c_j=E[f|B_j].$$
(49)
Conversely, every function $`F`$, measurable relative to $`_0`$, has this form for some $`\{c_j\}`$. The map, $`fF_f`$, is linear and is called the conditional expectation of $`f`$ given $`_0`$. This map leaves invariant the vector space of $`_0`$-measurable functions, and indeed is the orthogonal projection of $`L^2(\mathrm{\Omega },,\mu )`$ onto $`L^2(\mathrm{\Omega },_0,\mu )`$.
The tribe $`_0`$ tells how fine was the division into the sets $`B_j`$, and determines how much detail can be obtained from the functions that are $`_0`$-measurable. From the fact the $`F`$ is the orthogonal projection, we see that $`E[f|_0]`$ is the best approximation (in the $`L^2`$-sense) to $`f`$ by functions that are $`_0`$-measurable.
Consider for example the price of a stock at time $`t`$, where $`t`$ is a non-negative integer; $`S(t)_{t0}`$ are then a family of random variables on $`\mathrm{\Omega }`$, and while we can find out the prices up to the present time, we cannot know the future. Suppose that $`t=N`$ is the present. The information contained in the knowledge of the prices at $`N+1`$ previous times, namely $`S(t=0)=s_0,S(t=1)=s_1,\mathrm{},S(t=N)=s_N`$, selects in $`\mathrm{\Omega }`$ a particular level set of these functions: this is the event
$$\{\omega \mathrm{\Omega }:S(0)(\omega )=s_0\mathrm{}S(N)(\omega )=s_N\}$$
Since we assume that $`S(t)`$ are $``$-measurable, this set lies in $``$, by the intersection property. The same for any other possible set of values of these observations. There is a smallest $`\sigma `$-tribe with respect to which all these functions are measurable, and in fact, this $`\sigma `$-tribe is generated by all the level sets described above. Call this $`_N`$. The set of all random variables that are $`_N`$-measurable is exactly the set of functions of the data $`S(0),\mathrm{},S(N)`$, measurable in the Lebesgue sense; they can therefore be computed from the data we have access to.
The increasing family $`\{_n\}`$ is called the filtration generated by the process. It provides a neat formulation of the Markov condition for a process $`X_n`$; let $`_n`$ be the $`\sigma `$-tribe generated by the r. v. $`X_n`$. Then a process is called Markovian if
$$E[X_n|_m]=E[X_n|_m]\text{ if }nm.$$
(50)
The idea is that the information contained in $`X_m`$, the present value, tells us as much about the future as the whole previous history. Consider again the semigroup $`\{T^n\}`$ of stochastic maps, acting on $`\mathrm{\Sigma }(𝐙)`$ in one time-step as in eq. (28). One can check that if $`p_0`$ is the initial probability distribution of the initial point of the path, then
$$p_0T^n=E_{p_n}[x_n|_0].$$
(51)
Here, $`\gamma =(x_0,\mathrm{},x_n)`$ and $`p_n(\gamma )=p_0(x_0)p(\gamma )`$ where $`p(\gamma )`$ is as in (29).
Wiener was able to put the Bachelier-Einstein diffusion theory on a rigorous footing. He has to define, first, the sample space $`\mathrm{\Omega }`$; then he needs a $`\sigma `$-tribe $``$ and a measure on it; he also needs a family of $``$-measurable functions $`X_t(\omega )`$ whose distribution has density of probability equal to $`\rho (x,t)`$ obeying the diffusion equation. Finally, he needs to get the continuum version of eq. (51).
Let $`\mathrm{\Omega }`$ be the set of all continuous functions $`\omega `$ of $`t0`$ with $`\omega (0)=0`$; these are called ‘Brownian paths’. Let $`(x_1,y_1)`$ be an interval of the real line, which we call a gate; we now consider the subset of paths which pass through the gate at time $`t_1`$. This set is called the cylinder set based on $`(x_1,y_1)`$. In symbols, it is
$$\{\omega \mathrm{\Omega }:x_1<\omega (t_1)<y_1\}$$
The $`\omega (t)`$ for various $`t`$ are coordinates of the point $`\omega `$; we have a condition on only one of the coordinates; the rest run over the real line. Consider another cylinder set, similarly constructed at time $`t_2>t_1`$, based on another open interval $`(x_2,y_2)`$. The intersection of these sets is a cylinder set based on rectangle $`(x_1,y_1)\times (x_2,y_2)`$ in the plane made by the coordinates $`\omega (t_1),\omega (t_2)`$. The path $`\omega (t)`$ passes through the first gate at time $`t_1`$ and the second at time $`t_2`$; it is a slalom. Consider the collection of subsets of $`\mathrm{\Omega }`$ consisting of all these cylinder sets defined by slaloms with any finite number of gates, at any selection of different positive times. The finite unions of these form a tribe. The smallest $`\sigma `$-tribe $``$ containing all these is the one we choose, so obtaining the measurable space $`(\mathrm{\Omega },)`$.
We first define a finitely additive measure on the tribe of cylinder sets. It is enough to give the measure of a general cylinder set, and to use the finite additivity. Starting at $`x=0`$, the probability density that a diffusing particle reaches $`x_1`$ at time $`t_1`$ is taken to be the Gaussian given by the Green function; thus the probability of lying in the interval $`x_1,y_1`$ is
$`\text{Prob}\{\omega (t_1)(x_1,y_1)|\omega (0)=0\}`$ $`=`$ $`{\displaystyle \frac{1}{(4\pi \kappa t_1)^{1/2}}}{\displaystyle _{x_1}^{y_1}}e^{x^2/(4\kappa t_1)}𝑑x`$ (52)
$`=`$ $`{\displaystyle _{x_1}^{y_1}}G(x,t_1)𝑑x.`$
The probability that the path goes through two gates, $`(x_1,y_1)`$ at $`t_1`$ and $`(x_2,y_2)`$ at $`t_2`$ is defined to be
$`\text{Prob}\{\omega (t_1)`$ $``$ $`(x_1,y_1)\text{ and }\omega (t_2)(x_2,y_2)|\omega (0)=0\}`$ (53)
$`=`$ $`{\displaystyle _{x_1}^{y_1}}𝑑x{\displaystyle _{x_2}^{y_2}}𝑑x^{}G(x,t_1)G(x^{}x,t_2t_1).`$
This can be interpreted as Bayes’s theorem, in which $`G`$ is the conditional probability density. Similarly, the probability of any cylinder set, based on a finite set of gates, can be given. The probability is the same, whether the gates are open, closed or half-open. We would like the measure we are constructing to be at least finitely additive. Thus we take the measure of the union of two disjoint cylinder sets to be the sum of the measures we have just given them individually. A possible problem arises if we add together infinitely many gates at time $`t_1`$ to make up the whole line; for, we would like our measure to be countably additive, and we need the consistency condition between the two ways to define the probability of reaching the gate $`(x_2,y_2)`$: from $`0`$ directly, with no gate at $`t_1`$, as given by eq. (52), or as the sum over all paths going through any complete set of disjoint gates at $`t_1`$, as got by summing eq. (53). Indeed, we do get the same answer, because of the propagating property of $`G`$:
$$_{\mathrm{}}^{\mathrm{}}𝑑x^{}G(xx^{},t_1)G(x^{}y,t_2t_1)=G(xy,t_2).$$
(54)
This is a continuous version of the obvious property of the stochastic matrices $`T^n`$ of a Markov chain, namely $`T^mT^n=T^{m+n}`$, in which the matrix product, expressed as the sum over an intermediate index, is replaced by the integral over the point $`x^{}`$. Thus our equation just expresses the semi-group property of the time-evolution of a first-order equation, here the heat equation. It is seen here as the main point which establishes the additivity of the finitely additive measure we have constructed on the tribe of cylinder sets.
Let us define $`_{[s,t]}`$ as the $`\sigma `$-tribe generated by the cylinder sets labelled by times in the interval $`[s,t]`$, and $``$ that generated by all of these. Then Wiener proved that there exists a unique measure on the measurable space $`(\mathrm{\Omega },)`$ that coincides with the measure above on the tribe of unions of such cylinder sets. This measure is now called Wiener measure. The Wiener process starting at $`0`$ is then the family of random variables defines by $`W_t(\omega ):=\omega (t),t0`$. The process has the following properties:
1. $`W_tW_s`$ is Gaussian with mean zero and variance $`ts`$, for $`t>s`$.
2. $`W_tW_s`$ is independent of $`W_vW_u`$ if $`0uvst`$.
3. $`W_0=0`$.
These properties characterise the process. By requiring that $`W_0^x=x`$ we get the Wiener process $`W_t^x`$ starting at $`x𝐑`$.
We now need the concept of the symmetric Fock space $`\mathrm{\Gamma }()`$ of a Hilbert space $``$. The $`n`$-fold tensor product $`^n`$ is the completed span of the symbols $`_{i=1}^n\psi _i=\psi _1\mathrm{}\psi _n`$ with the scalar product
$$\psi _i,\varphi _i:=\underset{i=1}{\overset{n}{}}\psi _i,\varphi _i.$$
The symmetric tensor product $`^n:=_S^n`$ is the subset of symmetric tensors, called the $`n`$-particle space; the zeroth tensor power is taken to be $`𝐂`$. The Fock space $`\mathrm{\Gamma }()`$ is the direct sum $`_{n=0}^{\mathrm{}}^n`$. This has the functorial property
$$\mathrm{\Gamma }(_1)\mathrm{\Gamma }(_2)=\mathrm{\Gamma }(_1)\mathrm{\Gamma }(_2).$$
As a special case, $`\mathrm{\Gamma }(𝐂)=𝐂𝐂\mathrm{}`$.
There is a unitary map $`L^2(\mathrm{\Omega },,\mu )\mathrm{\Gamma }(L^2([0,\mathrm{}),dt)`$, in such a way that $`_{j=0}^n^j`$ is identified with the $`L^2`$-completed span of the polynomials in $`W_t`$ of degree $`n`$ . The $`n`$-particle space is then identified successively by Gram-Schmidt orthogonalisation with the part orthogonal to the $`k`$-particle spaces, $`k<n`$. This is Wiener’s chaos expansion . In particular, the one-particle space is spanned by $`\{W_t\}_{t0}`$.
For each fixed $`t`$, the space $`L^2(\mathrm{\Omega },,\mu )`$ contains the random variables $`I`$, $`W_t`$, …$`W_t^n`$…They can act as multiplication operators successively on the vector $`1`$, to get $`n`$ vectors. Suppose we orthogonalise them by the Gram-Schmidt procedure. Since $`W_t`$ is Gaussian, we get the Hermite polynomials in successive spaces, and any $`L^2`$ function of $`W_t`$ has a convergent expansion as a sum of its components in these spaces. The subspace we get can be identified as the Fock space over the one-dimensional space spanned by $`W_t`$. We shall see that these polynomials are Wick-ordered powers of $`W_t`$, and that they are martingales.
Now suppose that $`u>s`$; since $`W_u`$ is independent of $`W_uW_s`$, and they are Gaussian, they are orthogonal in the one-particle space, which can thus be written as the direct sum $`L^2([0,\mathrm{}),dt)=L^2([0,s),dt)L^2([s,\mathrm{}),dt)`$. By the functorial property of Fock space, we therefore can write
$$L^2(\mathrm{\Omega },,\mu )=L^2(\mathrm{\Omega },_{[0,s]},\mu )L^2(\mathrm{\Omega },_s,\mu ).$$
(55)
We can similarly split Fock space into arbitrarily many factors, corresponding to any partition of the time axis into intervals: it has the property of a continuous tensor product.
The continuous analogue of the semigroup $`(\gamma T^n)_m:=\gamma _{m+n},m,n=0,1,2\mathrm{}`$ of the random walk is the left-shift of the paths: $`(\omega T_s)(t)=\omega (s+t)`$. This induces the dual action on the observables:
$$T_s^{}:L^2(\mathrm{\Omega },,\mu )L^2(\mathrm{\Omega },_s,\mu ),s0.$$
(56)
This operator is isometric but not invertible. We can also embed $`L^2(\mathrm{\Omega },,\mu )`$ in the two-sided space $`\mathrm{\Gamma }(L^2(\mathrm{},\mathrm{}))`$, on which the left shift is unitary, and induces the action of the group $`𝐑`$ rather than the semigroup $`𝐑^+`$. In that case the paths are not conditioned to pass through the origin, and only the differences, $`W_tW_s`$ make sense as vectors or operators.
## 3 The Quantum Leap
The remarkable discovery of matrix mechanics by Heisenberg in 1925 is comparable to that of the theory of relativity in 1917. Clifford had speculated that the world might have chosen a geometry other than Euclidean. It was agreed that it was an experimental question, and that the data agreed with Einstein’s theory. Though the classical axioms were yet to be written down by Kolmogorov, Heisenberg, with help of the Copenhagen interpretation, invented a generalisation of the concept of probability, and physicists showed that this was the model of probability chosen by atoms and molecules.
According to Einstein et al. a concept is deemed to be an element of reality within a specified theory if there is a mathematical object in the theory which is assigned to the concept, and which takes a definite value (when the state of the system is given). This is now called an observable. For example, the choice of the zero-level of a potential function, is not an observable since it is not determined by the state of the system. They are not here discussing random samples, which at the time would have been described as an ensemble. In that case, they might have conceded that a concept could be regarded as an element of reality if, in a random selection of the system from an ensemble, there is a definite random variable assigned to the physical concept. The interpretation of a theory is not complete unless it is specified at the outset which mathematical objects arising in the theory correspond to observables. Thus in a theory with randomness in classical physics, there is a space $`(\mathrm{\Omega },)`$ and an observable is a random variable, and an ensemble is a probability measure on $`\mathrm{\Omega }`$. A non-random state is given by a point-measure. In this state any r. v. has zero variance, thus satisfying EPR.
In quantum mechanics, this is not the case; an observable is a Hermitian matrix $`A`$, or in modern terms, a self-adjoint operator on a given Hilbert space $``$; the possible values one can find in a measurement are the eigenvalues of $`A`$. A wave-function is determined by a vector $`\psi `$; but only unit vectors are used, and $`e^{i\theta }\psi `$ represents the same state as $`\psi `$. Thus the state is the equivalence class $`\{\psi \}=\{e^{i\alpha }\psi ,\alpha 𝐑\}`$. If $`dim=n<\mathrm{}`$, such equivalence classes make up the projective space $`CP^{n1}`$. An element of $`CP^{n1}`$ determines the expectation value of any observable $`A`$ by $`\psi ,A\psi `$, which according to the Copenhagen interpretation, is the mean value of $`A`$ if measured many times in the state $`\{\psi \}`$. It is seen to be independent of the representative vector $`\psi \{\psi \}`$. Such a state is called a vector state. The concept of state was generalised by von Neumann to include random mixtures of vector states . Let $`()`$ denote the set of bounded operators on $``$; this is a complex vector space, and also -algebra, where conjugation is given by the adjoint and multiplication is the usual product of operators. A state is given by a positive operator $`\rho `$ of trace 1, called a density operator, and the expectation of an observable $`A`$ is taken to be $`m_1(A):=\text{Tr}(\rho A)`$. Any density operator determines an element of the dual space to $`()`$ by the map $`Am_1(A)`$. We also can define $`m_n(A):=\text{Tr}\rho A^n`$ to be the $`n^{\mathrm{th}}`$ moment of $`A`$, and $`\kappa _2(A):=m_2(A)m_1(A)^2`$ to be the second cumulant, the variance, uncertainty or dispersion of $`A`$ in the state $`\rho `$. von Neumann showed that there are no dispersion-free states. Thus, quantum mechanics is intrinsically random. Heisenberg’s uncertainty relation, which is a theorem, not a postulate, is the best-known facet of this:
###### Theorem 3.1
Let $`A,B,C()`$ be such that $`[A,B]:=ABBA=C`$; then in any state $`\rho `$, we have $`\kappa _2(A)\kappa _2(B)m_1(C)^2/4`$.
There is no uncertainty relation for commuting operators $`A,B`$, and such observables are said to be compatible. If $`[A,B]0`$, we say that $`A`$ and $`B`$ are complementary.
Segal has emphasised that the bounded observables in any quantum theory should form the Hermitian part of a $`C^{}`$-algebra with identity. This is a complex vector space $`𝒜`$ with
1. a product $`AB`$ is defined for all $`A,B𝒜`$, which is distributive and associative, but not necessarily commutative;
2. a conjugation $`AA^{}`$, which is complex-antilinear, is specified;
3. $`𝒜`$ is provided with a norm $``$ which obeys Gelfand’s condition
$$A^{}A=A^2;$$
(57)
4. $`𝒜`$ is complete in the topology given by this norm.
This concept includes all the examples we have seen so far; the set $`_n(𝐂)`$, denoting $`n\times n`$ matrices, with matrix addition and product, is a $`C^{}`$-algebra. The operation is Hermitian conjugate, and the norm $`A`$ is the maximum eigenvalue of $`|A|=(A^{}A)^{1/2}`$. For any Hilbert space, $`()`$ is also a $`C^{}`$-algebra, and more generally, so is any von Neumann algebra, which can be defined as any weakly closed -subalgebra of $`()`$ containing the identity. Another notable example is the subset $`𝒞(n)`$ of $`(n)`$ consisting of real diagonal matrices $`A=\text{diag}(a_1,\mathrm{},a_n)`$. This is clearly commutative, and the diagonal elements are the eigenvalues. Thus, each $`A𝒞`$ determines uniquely a function $`ia_i,\mathrm{\hspace{0.33em}1}in`$ from the set $`\mathrm{\Omega }_n=(1,2,\mathrm{},n)`$ to $`𝐑`$. Conversely, any random variable $`f`$ on $`\mathrm{\Omega }_n`$ defines a unique diagonal matrix diag$`(f(1),\mathrm{},f(n))`$. So the classical observables on $`\mathrm{\Omega }_n`$ can be described as a special type of quantum mechanics, namely, the diagonal matrices. Moreover, the interpretation in classical theory, of the values of the random variables $`f`$ as possible observed values, coincides with the quantum interpretation of the eigenvalues. Also, each $`n\times n`$ density matrix $`\rho `$ defines a unique probability measure $`p`$ on $`\mathrm{\Omega }_n`$, by using the diagonal elements: $`p(i)=\rho _{ii},\mathrm{\hspace{0.33em}1}in`$. Clearly, a probability $`p`$ can define a density matrix by the same formula, but there are other, non-diagonal density matrices giving the same $`p`$. If all the observables are contained in $`𝒞`$, then the off-diagonal elements of the density matrix are of no relevance, and all the information on the state of the system is contained in $`p`$. A concept that captures the essentials of this idea, removing redundant description, is due to Segal. Given the algebra of observables, $`𝒜`$, we say a state on $`𝒜`$ is a positive, normalised linear map $`\rho :𝒜𝐂`$. Thus
1. $`\rho `$ is complex linear;
2. $`\rho (I)=1`$;
3. $`\rho (A^{}A)0`$ for all $`A𝒜`$.
Naturally, two constructions that lead to the same map are said to define the same state. We should note that we only need the expectations, i.e. the first moments, of the observables, because $`𝒜`$ itself contains all powers of $`A`$, and (as it is complete), also elements such as $`e^{iA}`$; so if we know the state we know the characteristic function of every observable, and so its distribution too.
More generally, the classical measure theory $`(\mathrm{\Omega },,\mu )`$, where $`\mu `$ is a positive measure, can be written as a (commutative) quantum theory by using the von Neumann algebra $`L^{\mathrm{}}(\mathrm{\Omega },\mu )`$ acting as multiplication operators on $`L^2(\mathrm{\Omega },,\mu )`$; its normal states correspond to (countably additive) probability measures, which vanish on $`\mu `$-null sets. Indeed, given a state $`\rho `$ we can define the corresponding measure of a set $`B`$ as $`\rho (\chi __B)`$. In this, sets $`B`$ and $`B^{}`$ are indistinguishable if the differ by a $`\mu `$-null set; we do not really need $`\mathrm{\Omega }`$ itself, but only the $`\sigma `$-tribe $``$, modulo this equivalence.
The set of states of a $`C^{}`$-algebra $`𝒜`$ forms a convex set, which we shall call $`\mathrm{\Sigma }(𝒜)`$ or just $`\mathrm{\Sigma }`$. The convex sum
$$\rho =\lambda \rho _1+(1\lambda )\rho _2,\text{ where }0<\lambda <1$$
(58)
represents the random mixing of the states $`\rho _1`$ and $`\rho _2`$ with weights $`\lambda `$ and $`1\lambda `$. All expectations in the state $`\rho `$ are then the same mixtures of the expectations in the states $`\rho _1`$ and $`\rho _2`$. If $`\rho _1\rho _2`$ we say that $`\rho `$ is a mixed state. If $`\rho `$ cannot be written as a mixed state (so that in any relation such as eq. (58) we must have $`\rho _1=\rho _2`$), the we say that $`\rho `$ is a pure state. Every $`C^{}`$-algebra possesses many pure states. For the full matrix algebra $`𝐌_n`$, every pure state is given by a unit ray $`\{\psi \}`$ in the Hilbert space $`𝐂^n`$, using the usual quantum-mechanical expression; every density operator is a mixture of such. This is an example of the Krein-Milman theorem, which says that a weak-compact convex set in the dual of a Banach space is generated by its extreme points. The representation of a mixed state as eq. (58) is, in general, not unique. For example, if $`=𝐂^2`$, the fully unpolarised state is $`(1/2)I`$, and this the equal mixture of the pure states, the eigen-vectors of $`J_3`$, the spin operator in the direction of quantisation, as well as the equal mixture of the eigenstates of $`J_1`$, or any other spin direction. This means that all statistical properties of the observables are the same however the state was made up. We express this by saying that the state-space in quantum probability is in general not a simplex: in a simplex, any mixed state has only one decomposition into pure states. In classical probability, in contrast, the state space $`\mathrm{\Sigma }(\mathrm{\Omega })`$ is a simplex. This is true in quantum probability only if $`𝒜`$ is abelian. The density matrix contains all the information there is. Our inability to distinguish the history of how the state was made is due to the quantum phenomenon of coherent sums of wave-functions.
There is an important connection between states and representations of a $`C^{}`$-algebra. A representation $`\pi `$ of $`𝒜`$ is a -homomorphism from $`𝒜`$ into $`()`$ for some Hilbert space $``$. Thus, $`\pi (A)`$ is an operator on $``$ and the map $`\pi `$ satisfies, for all $`A,B𝒜`$,
1. $`\pi (\lambda A+B)=\lambda \pi (A)+\pi (B),`$ for all $`\lambda 𝐂`$ (linearity);
2. $`\pi (A^{})=(\pi (A))^{}`$(hermiticity).
A representation is said to be faithful if $`\pi (A)`$ is non-zero if $`A0`$. A state $`\rho `$ is said to be faithful if $`\rho (A^{}A)=0`$ only for $`A=0`$. To each state $`\rho `$ there is a representation $`\pi _\rho `$, on a Hilbert space $`_\rho `$, and a unit vector $`\psi _\rho _\rho `$, such that $`\rho `$ is vector state $`\psi _\rho `$; that is,
$$\rho (A)=\psi _\rho ,\pi _\rho (A)\psi _\rho ,A𝒜.$$
(59)
If the state $`\rho `$ is faithful, then so is the corresponding representation $`\pi _\rho `$. Moreover, $`\pi `$ is irreducible if and only if $`\rho `$ is pure.
The proof of this theorem, which asserts the existence of $`_\rho `$ and the homomorphism $`\pi _\rho `$, follows the common mathematical trick: we construct these objects out of the material at hand. Let us do this when $`𝒜`$ has an identity and $`\rho `$ is faithful. We start with the vector space $`𝒜`$ and provide it with the scalar product
$$A,B:=\rho (A^{}B).$$
The completion of this space is then taken to be $`_\rho `$. The operator $`\pi _\rho (A)`$ is taken to be left-multiplication of $`𝒜`$ by $`A`$, thus: $`\pi _\rho (A)B:=AB`$. This defines $`\pi _\rho (A)`$ on the dense set $`𝒜_\rho `$, and can be shown to be bounded. We take $`\psi _\rho =I`$, the identity in the algebra. One can then verify that $`(_\rho ,\pi _\rho ,\psi _\rho )`$ satisfy eq. (59). A slightly more elaborate construction can be given if there is no identity or the state is not faithful. This realisation of the algebra is called the GNS construction, based on $`\rho `$.
It took some time before it was understood that quantum theory is a generalisation of probability, rather than a modification of the laws of mechanics. This was not helped by the term quantum mechanics; more, the Copenhagen interpretation is given in terms of probability, meaning as understood at the time. Bohr has said that the interpretation of microscopic measurements must be done in classical terms, because the measuring instruments are large, and are therefore described by classical laws. It is true, that the springs and cogs making up a measuring instrument themselves obey classical laws; but this does not mean that the information held on the instrument, in the numbers indicated by the dials, obey classical statistics. If the instrument faithfully measures an atomic observable, then the numbers indicated by the dials should be analysed by quantum probability, however large the instrument is.
We now present Gelfand’s theorem, which shows that any commutative quantum theory can be viewed as a classical probability theory. We give a proof in finite dimensions.
###### Theorem 3.2
Given a commutative -algebra $`𝒞`$ of finite dimension, there exists a (finite) space $`\mathrm{\Omega }`$ and an algebraic -isomorphism $`J`$ from $`𝒞`$ onto $`𝒜(\mathrm{\Omega })`$, such that for any state $`\rho `$ on $`𝒞`$ there exists a probability $`p`$ on $`\mathrm{\Omega }`$, such that for any element $`A𝒞`$ we have
$$\rho (A)=E_p[J(A)].$$
(60)
Proof
Since dim$`𝒞=n<\mathrm{}`$, the dimension of the dual space is the same. There is a faithful state $`\omega `$ on $`𝒞`$; this could be for example a mixture of a basis of the state-space with non-zero coefficients. We can therefore construct a faithful realisation of $`𝒞`$ as a matrix algebra. In this, the GNS construction, the Hilbert space is built out of $`𝒞`$ and so is of dimension $`n`$. A commutative collection of normal matrices can be simultaneously diagonalised, so there is a basis in the Hilbert space such that each element of $`𝒞`$ is a diagonal $`n\times n`$ matrix. Since exactly $`n`$ of these diagonal matrices make up a linearly independent set, every diagonal matrix appears. Every element of $`𝒞`$ is therefore a sum of multiples of units $`\{e_j\}`$ of the algebra, satisfying $`e_j^2=e_j`$ and $`e_ie_j=0,ij`$. In the above matrix realisation, $`e_j`$ is the matrix with $`1`$ on the diagonal in position $`j`$, and zero elsewhere. Thus $`A=a_je_j`$. So let $`\mathrm{\Omega }=\{e_j\}_{j=1,\mathrm{}n}`$, and let $`JA`$ be the function $`JA(e_j)=a_j`$. Then one verifies that $`J`$ is an algebraic -isomorphism. To the state $`\rho `$ we associate the probability $`p(e_j)=\rho (e_j)`$, and see easily that eq. (60) holds.$`\mathrm{}`$
In this proof, instead of identifying $`\mathrm{\Omega }`$ with the collection of elements $`e_j`$ in the algebra, we could have taken the dual, and identified $`\mathrm{\Omega }`$ with the set of characters on $`𝒞`$. This is the set of multiplicative states, that is, states $`\omega `$ obeying $`\omega (AB)=\omega (A)\omega (B)`$ for all $`A,B𝒞`$. The set of characters of a $`C^{}`$-algebra is called its spectrum. Our proof shows that there are exactly $`n`$ of these, defined by $`\omega _j(e_k)=\delta _{jk}`$. Putting $`A=B`$ we see that any character is dispersion-free. This is why the spectrum is taken by Gelfand to be the definition of $`\mathrm{\Omega }`$ in the infinite-dimensional case:
###### Theorem 3.3
Let $`𝒞`$ be a commutative $`C^{}`$-algebra with identity. Then the set of characters can be given a topology so as to form a compact Hausdorff space $`\mathrm{\Omega }`$ such that $`𝒞`$ is $`C^{}`$-isomorphic to $`C(\mathrm{\Omega })`$, and every state on $`𝒞`$ corresponds to a finitely additive measure on $`\mathrm{\Omega }`$ (with the Borel tribe).
Bohm asked whether the observed statistics, agreeing with experiment, can be obtained from a larger, more complicated classical theory. This is the idea behind the attempts to introduce hidden variables. This is certainly true of the statistics of any fixed complete commuting set of observables; for they form an abelian algebra, and so can be represented by the classical statistics of multiplication operators on a sample space (the spectrum of the algebra). Obviously the full non-abelian algebra cannot be a subalgebra of an abelian algebra, so the way hidden variables are introduced must be more elaborate than extending the algebra by adding them. However, the deep result of J. S. Bell shows (if the dimension is 4 or higher) that the full set of statistics predicted by quantum theory cannot be got from any underlying classical theory. In the quantum model of two spin-half systems, Bell constructs a sum of four correlations which in a certain state is equal to $`2\sqrt{2}`$, a factor $`\sqrt{2}`$ larger than the greatest value allowed in any classical theory.
Let us follow . Let $`P,Q`$ be complementary projections, and also let $`P^{},Q^{}`$ be complementary projections, while $`P`$ is compatible with $`P^{}`$ and with $`Q^{}`$, and $`Q`$ is compatible with $`P^{}`$ and $`Q^{}`$. Define $`a=2PI`$, $`b=2QI`$, and similarly for $`a^{}`$ and $`b^{}`$. For any state $`\rho `$ define $`R`$ by
$$R:=\rho (aa^{}+ab^{}+bb^{}ba^{})=\rho (C)$$
where $`C=a(a^{}+b^{})+b(b^{}a^{})`$. Then $`a^2=b^2=a^2=b^2=1`$, so
$$C^2=4+[a,b][a^{},b^{}]=4+16[P,Q][P^{},Q^{}].$$
(61)
Since $`a=b=a^{}=b^{}=1`$, it follows that
$$[a,b][a^{},b^{}]4,$$
so $`C^28`$ and $`|R|^2=|\rho (C)|^2\rho (C^2)8`$. So in quantum theory, $`|R|2\sqrt{2}`$. If there is a joint probability space on which we can describe $`a,\mathrm{},b^{}`$ by the r. v. $`f,\mathrm{}g^{}`$ taking the values $`\pm 1`$, and a measure $`p`$ on it, then $`R=E_p[h]`$ where
$$h=f(f^{}+g^{})+g(g^{}f^{}).$$
Then these r. v. commute, so eq. (61) becomes $`h^2=4`$, and
$$|R^2|=E_p[h]^2E_p[h^2]=4.$$
So $`|R|2`$, (Bell’s inequality). Bell showed that the entangled states of the Bohm-EPR set-up give a $`\rho `$ such that $`R=2\sqrt{2}`$, violating this. Thus no description by classical probability is possible.
The famous Aspect experiment tested Bell’s inequalities. This involves observing a system (in a pure entangled state) in a long run of measurements; the correlations singled out by Bell, between several compatible pairs of spin observables, were measured. The experiments showed that $`R`$ was just less than $`2\sqrt{2}`$, in agreement with the quantum predictions.
The upshot is that in quantum probability there is no sample space; we have the $`C^{}`$-algebra $`𝒜`$, and this plays the rôle of the space of bounded functions.
Let us now examine Bohm’s claim that there is a hidden assumption in Bell’s proof, that of ‘locality’. It is now generally agreed that the term ‘local’, referring to the space-localisation, is not the best, and that ‘non-contextual’ is a better term; namely that the choice of random variable $`f`$ assigned to represent a certain observable $`a`$ which is being measured, does not depend on which of the other observables, $`a^{}`$ or $`b^{}`$, is being measured at the same time. This is now called a non-contextual assignment. Bohm suggested that we should allow a contextual choice of assignment of random variable, so that the r.v. representing the observable $`a`$ when $`a^{}`$ is also measured is not the same as the choice of r.v. for $`a`$ when it is measured with $`b^{}`$. The two choices will, however, have the same distribution. It should be said straight away that this idea is contrary to the practice of probabilists, who would expect there to be a unique random variable representing an observable. It also goes against the definition of ‘element of reality’ of EPR as extended by us to the random case. The quantum version does not suffer from this unreality, since the mathematical object assigned to the observable, the Hermitian matrix, does not depend on the context, i.e. is local in Bohm’s language.
Bohm’s idea leads to a theory with very few rules. However there are some restrictions, since the choice must be done so that all statistical measurements of compatible observables (means, correlations, third moments etc) of the model can be arranged to give the same answers as the quantum theory. This is achieved as follows. Let $`a,a^{}`$ be compatible, generating a commutative $`C^{}`$-algebra, $`𝒞`$ and let $`\rho `$ be a state on the full algebra $`𝒜`$. By restriction, $`\rho `$ defines a state on $`𝒞`$. By Gelfand’s isomorphism, we can construct a space $`\mathrm{\Omega }`$, the spectrum of $`𝒞`$, and a measure $`\mu `$ on it, such that $`a,a^{}`$ can be represented as multiplication operators on $`𝒞(\mathrm{\Omega })`$, so they are random variables, $`f,g`$. The joint probability distribution of $`f,g`$ is the same as that of the (diagonal) matrices $`a,a^{}`$ in the state $`\rho `$. On the other hand, $`\mathrm{\Omega },\mu `$ depends on the set $`a,a^{}`$. Let us record this by denoting this Gelfand representation by $`\mathrm{\Omega }_{a,a^{}},\mu _{a,a^{}}`$. If we measure $`a`$ and $`b^{}`$, and proceed as Bohm suggests, then we get a different space $`\mathrm{\Omega }_{a,b^{}}`$, the spectrum of a different algebra $`𝒞_{a,b^{}}`$, say. The state $`\rho `$ leads to a different measure $`\mu _{a,b^{}}`$. The r.v. assigned to $`a`$ cannot be $`f`$ this time; it must a function on $`\mathrm{\Omega }_{a,b^{}}`$, a different space; it has the same distribution in $`\mu _{a,b^{}}`$ as the $`f`$ had in $`\mu _{a,a^{}}`$. In this set-up, there is no obvious definition of $`a^{}+b^{}`$, as they are not functions on the same space. This problem does not arise in the quantum formulation: there is an underlying $`C^{}`$-algebra, in which we can add the operators.
Bohm’s suggestion might be said to be an interpretation of quantum mechanics in terms of classical probability . However, the construction is not a probability theory in the sense of Kolmogorov, as there is no single sample space; the theory is preKolmogorovian, in the tradition of the frequentist school. One can generalise the frequentist point of view, and specify that certain collections of observations are compatible, and others are not; then we can by observation construct the joint probabilities of each compatible set, and have no need of the sample space (the space of joint values). A different compatible set need have no analytic relation to the first, even though it contain common observables. Bell’s inequality need not hold, but then neither need the quantum version, which is $`\sqrt{2}`$ times more generous. It is a feeble theory, not much more that data collection, and has no predictive power. Mere data give us no more than mere data.
Another variant of quantum mechanics, a new form of algebra called ‘quantum logic’, was developed in . New rules by which propositions can be manipulated are given. This was worked on later by Jauch and coworkers , culminating in Piron’s thesis. This says that the propositions form a lattice isomorphic to the lattice of subspaces in a Hilbert space (but not necessarily over the complex field). Apart from this result, quantum logic has not been very successful, and it is more productive to keep to classical logic, but to generalise the concept of probability algebra from commutative to non-commutative. Another alternative to quantum probability is stochastic mechanics, founded but now abandonned by Nelson as not being correct physics. Thus Segal’s approach is the one we adopt here. It is well explained in .
Quantum theory has its version of estimation theory . In finite dimensions, the method of maximum likelihood is to find the density matrix $`\rho `$ that maximises the entropy, subject to given values for the means, $`\{\eta _i\}`$ of observables in the subspace of hermitian operators spanned by a named list $`\{X_1,\mathrm{},X_n\}`$ of slow variables. So $`\eta _i=\text{Tr}[\rho X_i]`$. It is well known that the answer is the Gibbs state
$$\rho =Z^1\mathrm{exp}(H)=Z^1\mathrm{exp}[\xi ^1X_1+\mathrm{}+\xi ^nX_n],Z=\text{Tr}[\mathrm{exp}(H)].$$
(62)
Again, $`\mathrm{log}Z`$ is strictly convex, and its Hessian gives a Riemannian metric on the manifold $``$ of all faithful density operators . In this case we get the Kubo-Mori-Bogoliubov metric; in terms of the centred variables $`\widehat{X}_i:=X_i\eta _i`$, the metric is
$$g(\widehat{X}_i,\widehat{X}_j)=\text{Tr}\left[_0^1\rho ^\lambda \widehat{X}_i\rho ^{1\lambda }\widehat{X}_j𝑑\lambda \right]$$
(63)
This is the closest point on $``$ to any state with the given means, where ‘distance’ is measured by the relative entropy $`S(\rho |\rho ^{}):=\text{Tr}\rho [\mathrm{log}\rho \mathrm{log}\rho ]`$. Again, the $`\xi ^j`$ are uniquely determined by the measured means $`\eta _i`$.
## 4 Kolmogorov and Ito
A stochastic process over a set $`T`$ is a family $`\{X_t,tT\}`$ of random variables on a measure space $`(\mathrm{\Omega },,\mu )`$. We might have $`T=\{0,1,2\mathrm{}\}`$, or $`T=𝐑^+`$, when we interpret $`t`$ as time. From the frequentist point of view, we can observe $`X_{t_1},X_{t_2},\mathrm{}X_{t_N}`$ at finitely many points of time. In this way, we can test any a model as to what the joint distribution of these r.v. is.
Kolmogorov’s existence theorem says that a family of joint (cumulative) distributions $`F_{1,2\mathrm{}n}(x_1,\mathrm{},x_n)`$, given for all finite subsets of $`T`$, is the set of joint distributions of a stochastic process over $`T`$ if and only if the consistency conditions hold. Thus, (the hatted variable is omitted):
1. For any permutation $`\pi `$ of $`(1,2,\mathrm{},n)`$, we have
$$F_{1,\mathrm{}n}(x_1,\mathrm{},x_n)=F_{\pi (1),\mathrm{},\pi (n)}(x_{\pi (1)},\mathrm{},x_{\pi (n)})$$
2. For any $`j`$, we have
$$F_{1,\mathrm{},n}(x_1,\mathrm{},x_j=\mathrm{},x_{j+1},\mathrm{},x_n)=F_{1,\mathrm{},\widehat{j},\mathrm{},n}(x_1,\mathrm{},\widehat{x_j},\mathrm{},x_n).$$
If these hold, he shows that the sample space $`\mathrm{\Omega }`$ may be taken to be $`𝐑^T`$, an enormous space (of all functions $`\omega `$ of $`t`$); the r. v. $`X_t`$ is then the function $`X_t(\omega )=\omega (t)`$. He proved the existence of a measure $`\mu `$, which reproduces the given joint distributions; the $`\sigma `$-tribe $``$ has the following structure. Let $`_t`$ be the smallest $`\sigma `$-tribe such that all $`X_s`$, for $`st`$, are measurable; then this is an increasing family of $`\sigma `$-tribes, called a filtration. Then $``$ is the smallest $`\sigma `$-tribe containing all the $`_t`$.
Apply this to the Brownian paths, and the measures defined by a finite set of gates as in the last section; this proves that there is a probability theory underlying the finite joint distributions. However, it does not prove Wiener’s theorem, in that the sample space obtained by the Kolmogorov construction is the huge set of all functions of time. It is then a hard problem to show that the subset of continuous functions has measure 1. This fact is very important for specialists in Brownian motion, but is not a general feature of processes covered by Kolmogorov’s theorem, and is not needed to construct the usual $`L^p`$ spaces of functional analysis. Without Wiener’s version we lose the power of the path-wise methods, and also lots of intuition. The modern method is to get the cow off the ice using Kolmogorov, and supplement it with further estimates, on tightness and radonifying maps, if we need to find smaller carrier spaces for the measure . After Kolmogorov’s treatise, the subject could develop ‘in the usual professional mathematical way’, to use Segal’s phrase. That is, theorems could be stated and proved, and then sharpened. The most important of these were the laws of large numbers, the zero-one laws, the central limit theorems, the theory of large deviations, the classification of all processes with independent increments, martingales, and stochastic integration.
The conditional expectation $`E_t:=E[|_t]`$ takes a random variable in $`L^2(\mathrm{\Omega },,\mu )`$ into one in $`L^2(\mathrm{\Omega },_t,\mu )`$; since it is the identity on the latter space, and is Hermitian, it must be the orthogonal projection onto $`L^2(\mathrm{\Omega },_t,\mu )`$. None of these ideas depends on which version of the sample space we have.
The concept of conditional expectation can be extended to integrable r.v., thus:
###### Definition 4.1
Let $`(\mathrm{\Omega },,\mu )`$ be a probability space, and let $`_0`$ be a sub-$`\sigma `$-tribe of $``$. Let $`X`$ be a random variable with $`E[X]<\mathrm{}`$. Then there exists a $`_0`$-measurable r.v. $`Y`$, written $`E[X|_0]`$, such that for each set $`B_0`$ we have
$$_BY𝑑\mu =_BX𝑑\mu $$
(64)
Further, if $`\widehat{Y}`$ is another r.v. with these properties, then $`\widehat{Y}=Y`$ almost everywhere.
See for a proof, and other things.
A martingale is a stochastic process $`X_t`$ on $`(\mathrm{\Omega },_{t0},\mu )`$ such that $`X_t`$ is integrable, and
$$E[X_t|X_s]=E[X_s]\text{ for all }ts.$$
(65)
A martingale is a fair game. For example, consider the independent tosses of a fair coin, and let $`X_n=H_nT_n`$, where $`H_n`$ is the number of heads and $`T_n`$ is the number of tails at the $`n^{\mathrm{th}}`$ toss. Let $`S_N=_{j=1}^NX_j`$. Then $`S_N`$ is a martingale , p 202.
There are four concepts of convergence of a sequence $`\{X_n\}`$ to $`X`$ in the space of random variables on a probability space $`(\mathrm{\Omega },,\mu )`$.
1. We say $`X_nX`$ almost surely (or, almost everywhere) if
$$\mu \{\omega :X_n(\omega )X(\omega )\}=1.$$
2. We say $`X_nX\text{ in }_r`$ if
$$X_nX_r0\text{ as }n\mathrm{}.$$
3. We say that $`X_nX`$ in probability if
$$\mu \{\omega :|X_n(\omega )X(\omega )|>ϵ\}0\text{ as }n\mathrm{}\text{for all }ϵ>0.$$
4. We say $`X_nX`$ in law if
$$\mu \{\omega :X_nx\}\mu \{\omega :Xx\}\text{ for all }x\text{ at which }F_X(x):=\mu \{Xx\}$$
is continuous.
These concepts are not equivalent; (1) and (2) are not comparable, but (1) or (2) imply (3), which implies (4) . Convergence in law can be related to convergence of the characteristic functions of $`X_n`$ to that of $`X`$; we see that if $`X_n`$ converges to $`X`$ in law implies that $`X_n`$ also converges to $`Y`$ in law if $`Y`$ has the same distribution as $`X`$. This shows that convergence in law in a very feeble concept. The four concepts of convergence do not depend on the version of sample space adopted, and so are the same whether we use Wiener space or Kolmogorov’s abstract construction.
For a given $`\mu `$, we can complete the $`\sigma `$-tribe $``$ to include all subsets of sets of $`\mu `$-measure zero; then the events that can happen are described by the quotient $`\sigma `$-tribe, in which events which differ by a set of measure zero are identified. This idea is not wise when we are interested in measures with different sets of zero measure, as happens when we condition a Wiener path to pass through a given point. The Dirac measure on $`𝐑`$ is a simple example of the trouble we get into. If two measures $`\mu _1,\mu _2`$ have the same sets of zero measure in $`(\mathrm{\Omega },)`$, we say that they are equivalent. If $`\mu _1(B)=0`$ whenever $`\mu _2(B)=0`$ we say that $`\mu _1`$ is absolutely continuous relative to $`\mu _2`$; in that case there exists a function $`wL^1(\mathrm{\Omega },,\mu _2)`$ such that $`\mu _1(B)=_B\rho (\omega )𝑑\mu _2`$ for all $`B`$. We write $`w=d\mu _1/d\mu _2`$, the Radon-Nikodym derivative. This is the abstract version of eq. (12). We shall be interested in other measures, singular relative to a given one. Then the best formalism is to start with an abelian $`C^{}`$-algebra $`𝒜`$ and consider its states.
Estimation is assisted by the law of large numbers. Let $`X`$ be a random variable on a probability space, whose mean $`\eta `$ we wish to find, making use of a random experiment which is believed to be well modelled by $`X`$. We set up a sequence of independent copies $`X_n`$ of $`X`$, and consider the stochastic process $`\{X_n\}`$ on e.g. the probability space constructed by the theorem of Kolmogorov. The strong law of large numbers says that if $`E(X)=\eta `$ and $`E(X^2)<\mathrm{}`$, then putting $`S_n=_{j=1}^NS_j`$ we have
$$S_N/N\eta I\text{ in }_2.$$
If $`E|X|<\mathrm{}`$, we get almost sure convergence. These are necessary and sufficient conditions. Weaker conditions ensure that the sum converges in probability . This is called the weak law. Note that the meaning of convergence uses the measure on the Kolmogorov space, so for almost all sequences, randomly chosen, we get convergence to the mean. It does not say how fast the convergence is. For example if $`X_n`$ is the number of heads minus the number of tails, at the $`n^{\mathrm{th}}`$ toss of a fair coin, then $`S_N`$ is the number of heads in $`N`$ tosses minus the number of tails, and $`S_N/N`$ converges almost surely to zero. If we know that $`S_m0`$ after $`m`$ results, the law does not say that the bias evens out in the long run. $`S_n`$ is a martingale, and its expected value for $`N>m`$ is its present value $`S_m`$. It is $`S_N/N`$, which converges; the bias at time $`m`$ gets divided by $`N`$, and so goes away for large $`N`$.
Another famous limit law is the central limit theorem; if the standard deviation $`X`$ is 1 and the mean is zero, then one can show that
$$S_n/(\sqrt{n})N(0,1)\text{ in law}.$$
Versions of this were known to Bernouilli and Gauss, if we assume that the moment-generating function exists. It explains the ubiquity of the normal distribution; many random processes are the sums of small and independent random things, and so tend to be Gaussian. The theory of large deviations tells us something about the rate of convergence of $`S_n`$; this stuff is deeper . There is also a large body of work on sums of nearly independent random variables, and also on the cases where the variances are not all equal.
Doob proved that martingales often converge; e. g.,
###### Theorem 4.2
If $`\{S_n\}`$ is a martingale with $`E(S_n^2)<M<\mathrm{}`$ for some $`M`$ and all $`n`$, then there exists a random variable $`S`$ such that $`S_nS`$ almost surely.
Consider now a process $`(X_t,\mathrm{\Omega },_{t0},\mu )`$ in continuous time with independent increments; that is, $`X_tX_s`$ is independent of $`X_r`$ for $`r<s<t`$. Since we can write $`X_tX_s`$ as the sum of more and more independent differences, we might expect that $`X_tX_s`$ must be Gaussian, by the central limit theorem. However, this is not the case since the distributions of the difference $`X_tX_s`$ might change as the interval is made smaller. This question led Levy to characterise all processes that are stationary and have independent increments. This can be done by showing that the characteristic function
$$C(\lambda ):=E[\mathrm{exp}i(X_tX_s)\lambda ]$$
should not only be of positive type, but so should any fractional power. Such a function is called infinitely divisible, and so is the corresponding random variable. The necessity is easy to see; we can write $`X_tX_s`$ as the sum of $`N`$ identical and independent random variables, namely, the increments for time intervals $`(ts)/N`$; then the characteristic function of this sum is the product of the $`N`$ characteristic functions (which are all equal, by stationarity) of these increments, and so C has an $`N^{\mathrm{th}}`$ root that is of positive type. This condition is also sufficient, to which we shall return. The characteristic function of the Gaussian is infinitely divisible, and so is that for the Poisson distribution. This means that Gaussian and Poisson processes with independent increments exist. Levy found that by mixing these he got some new processes (Levy processes), and he found the most general form of the characteristic function, which is
$$\mathrm{log}C(\lambda )=a\lambda ^2+ib\lambda +𝑑\sigma (\alpha )[e^{i\alpha \lambda }(1+i\alpha \lambda )Z(\alpha )]$$
(66)
where $`a0`$, $`b`$ is real, $`d\sigma (\alpha )0`$ obeys $`_1^1\alpha ^2𝑑\sigma (\alpha )<\mathrm{}`$. There are some further conditions on the weight $`d\sigma `$ at infinity . If $`\sigma =0`$ we get the Gaussian, and if $`\sigma `$ has a discontinuity, we get a Poisson process. These can be understood in terms of Hilbert space cohomology of $`𝐑`$, as in §(5).
During this period, physicists and engineers studied stochastic differential equations, similar to the Langevin equation. Often the random force was chosen to be the derivative of Brownian motion, called white noise. Since $`B_t`$ is at best continuous, this work lacked rigour, and remained poorly defined even after appeal to Dirac’s generalised concept of function. This sorry state of affairs was cleared up by Ito.
Let $`W(t)`$ be Brownian motion starting at zero. At first sight, an equation for an unknown $`X(t)`$ similar to the Langevin equation, of the form eq. (43)
$$\frac{dX_t}{dt}=a(t)+b\frac{dW_t}{dt},\text{ for almost all }\omega $$
makes no sense, since for almost all $`\omega \mathrm{\Omega }`$, $`W(t)`$ is not differentiable. The equation does make sense if written in the form: find a family of random variables, $`\{X(t)\}`$, such that for a given initial random variable $`X(0)`$, the r. v. $`X(t)X(0)_0^ta(s)𝑑s`$ is the known r.v. $`W_t`$ for almost all $`\omega `$. This does not prove that there is such a process, but it is does make sense. For the more general case when $`a,b`$ depend on the unknown $`X(t)`$, the integral form is
$$X(t)X(0)=_0^ta(s,X(s))𝑑s+_0^tb(s,X(s))𝑑W(s).$$
(67)
The last expression, called a stochastic integral, looks like a Stieltjes integral, but the needed condition of bounded variation on $`W(s)`$ do not hold. Solve the equation by iteration (Picard’s method); we see that at each stage, the approximation to $`X(t)`$ is a function of $`W(s)`$ only for $`st`$. So it would appear that we need only give a meaning to the stochastic integral for the cases where $`X(t)`$ is a function of $`W(s)`$ for $`st`$, and so the same holds for $`b(t,X(t))`$. This can be neatly put in terms of the filtration $`_t`$ generated by the Wiener process: for all $`t0`$, $`X(t)`$ and so $`b(t,X(t))`$ is measurable relative to the $`\sigma `$-tribe $`_t`$. This makes sense physically; it says that we can know the present configuration $`X(t)`$ if we know the initial configuration $`X(0)`$ and the outcomes of all the randomness, $`W_s,st`$, so far. A random function of time, $`f`$ is said to be adapted to the filtration $`_t`$ if $`f(t)`$ is $`_t`$-measurable for all $`t0`$.
Let $`f(t)`$ be an adapted process in the time interval $`0tT`$, which is simple: that is there is a finite partition $`0=t_0,t_1,\mathrm{}<t_n=T`$ such that $`f(t)=f_j`$ for $`t_{j1}t<t_j`$ for all integers $`j(1,\mathrm{},N)`$. Here, $`f_j`$ is a random variable independent of time, and equality of random variables means almost everywhere. Following Ito, we can define the stochastic integral of an adapted simple function $`f`$ to be the random variable
$$_0^Tf(t)𝑑W_t:=\underset{j}{}f_j(W_{t_{j+1}}W_{t_j}).$$
(68)
Note that the increment $`dW`$ is in the future of the random variable $`f_j`$ that multiplies it. The mapping, $`f_0^Tf(t)𝑑W_t`$ takes the linear space of simple adapted functions into the space of random variables, and is clearly a linear map. The brilliant remark of Ito is then that the following identity, called Ito’s isometry, holds:
$$E[|_0^Tf(t)𝑑W_t|^2]=_0^TE[|f(t)|^2]𝑑t$$
(69)
Proof:
$`E[|{\displaystyle _0^T}f(t)𝑑W_t|^2]`$ $`=`$ $`{\displaystyle \underset{i}{}}{\displaystyle \underset{j}{}}E[f_i(W_{t_{i+1}}W_{t_i})f_j(W_{t_{j+1}}W_{t_j})]`$
$`=`$ $`{\displaystyle \underset{i}{}}E[|f_i|^2(W_{t_{i+1}}W_{t_i})^2`$
$`+`$ $`2{\displaystyle \underset{i<j}{}}f_if_j(W_{t_{i+1}}W_{t_i})(W_{t_{j+1}}W_{t_j})].`$
Now, the future increment $`W_{t_{i+1}}W_{t_i}`$ is independent of $`f_i`$, which is adapted, i.e. a function of earlier $`W(s)`$. So the expectation value in the first term factorises:
$`E[|f_i|^2(W_{t_{i+1}}W_{t_i})^2]`$ $`=`$ $`E[|f_i|^2]E[(W_{ti+1}W_{t_i})^2]`$
$`=`$ $`E[f_i|^2](t_{i+1}t_i),`$
by the property of Brownian motion. This gives the desired term in eq. (69). It remains to show that the remaining double sum vanishes. This is true, because the factor $`(W_{t_{j+1}}W_{t_j})`$ for $`j>i`$ is independent of the remaining factors $`f_if_j(W_{t_{i+1}}W_{t_i})`$ and so the expectation of the product is the product of the expectations; but the expectation of the future increment of $`W_t`$ is zero.$`\mathrm{}`$
Ito’s isometry is a mapping from the set of simple adapted processes to random variables; by a simple theorem of normed spaces, it can be extended by continuity to a linear isometry (unitary transformation) between the completions of both sides in the norms given. The completion of simple functions in the norm
$$f^2=_0^TE[|f(t)|^2]𝑑t$$
(70)
is the space of processes such that $`E[|f|^2]`$ is Lebesgue integrable; so Ito can define the stochastic, or Ito integral, of all processes with this property; it is the limit in this norm of simple adapted processes approximating it. Naturally, we must prove that the adapted simple processes are $`L^2`$-dense in the square-integrable adapted processes; this is not difficult, since the projection $`E_t`$ is a bounded operator and maps onto the space of $`_t`$-adapted square-integrable processes.
We can now give a meaning to the question, do there exist solutions to the stochastic differential equation
$$\frac{dX_t}{dt}=a(X_t,t)+b(X_t,t)\frac{dW_t}{dt}\mathrm{?}$$
(71)
We say the a process $`X_t`$ satisfies this equation if, on substituting $`X_s`$ in the integrals in eq. (67) we get back $`X_tX_0`$.
For a wide class of functions $`a`$ and $`b`$ of two variables we can then get a convergent iterated approximation, the Picard series, which converges to a process $`X_t`$ obeying the (integral form of) the stochastic differential equation. This holds for example if $`a(x,y)`$ is uniformly Lipschitz in $`y`$ in a region, and $`b(x,y)`$ is uniformly elliptic in $`y`$ and measurable in $`x,y`$. This result can be improved and generalised, so that vector-valued stochastic processes can be studied, and the noise can be of a much more general martingale than $`W_t`$. This can be reworded as a ‘martingale problem’ .
The converse to Ito integration should be a form of differentiation: it is called (Ito) stochastic differentiation; we may say that the process $`f(W_t,t)`$ is the stochastic derivative of $`_0^tf(W_s,s)𝑑W_s`$. The Ito integral is always a martingale, and every martingale is a stochastic integral, and so has a stochastic derivative, namely the integrand in its representation as an Ito integral. One can show that this is unique. It is interesting to form the repeated stochastic integrals
$`W_t`$ $`=`$ $`{\displaystyle _0^t}𝑑W_s`$
$`:W_t^2:=W_t^2t`$ $`=`$ $`2{\displaystyle _0^t}W_s𝑑W_s`$
$`:W_t^3:`$ $`=`$ $`3{\displaystyle _0^t}:W_s^2:dW_s`$
$`\mathrm{}`$ $`\mathrm{}`$
in which the Wick ordered (Hermite polynomials) occurring in the Wiener chaos are the successive stochastic integrals. They are all contained in the exponential martingale $`e^{\lambda X_t\frac{1}{2}\lambda ^2t}`$. The second one illustrates the Doob-Meyer decomposition: $`W_t^2`$ is a submartingale, and is written as the sum of a martingale, $`:W_t^2:`$, and an increasing function, $`t`$ of bounded variation.
Manipulation of stochastic integration can be summarised by the Ito multiplication table: keep all differentials in $`dt`$ up to first order, using $`dt.dW=0`$ and $`dW.dW=dt`$. From this, we can get the important relation between a certain parabolic partial differential equations known as Kolmogorov’s forward equation, and the corresponding stochastic differential equation. Suppose that $`X_t`$ satisfies the stochastic differential equation eq. (71), with initial r.v. equal to $`X_0`$, which has law $`p(x)`$. Let $`p(x,t)`$ be the law of $`X_t`$; then it can be shown that $`p(x,t)`$ is smooth and satisfies the parabolic equation
$$\frac{p}{t}=\frac{1}{2}\frac{}{x}\left(b(x,t)^2\frac{p}{x}\right)+\frac{}{x}(a(x,t)p),$$
(72)
with initial condition $`p(x,0)=p(x)`$. To see why this is, we note that if $`f(x)`$ is any smooth function, we can apply Ito’s lemma to the random process $`f(X_t)`$. We recover $`p(x,t)f(x)𝑑x`$ as $`E[f(X_t)|X_t=x]`$. We now expand $`f(X_{t+dt})`$ in a Taylor series about $`X_t`$ up to second order in $`dW`$:
$$f(X_{t+dt})=f(X_t)+\frac{f}{x}dX+\frac{1}{2}\frac{^2f}{x^2}(dX)^2.$$
(73)
Eq. (71) tells us that $`(dX_t)^2=b^2dt`$ and $`dX=adt+bdW`$. Here, $`dW`$ is the forward difference. Then the expectation vanishes: $`E[f^{}bdW|X_tx]=E[f^{}b|X_t=x]E[dW|X_t=x]=0`$, since $`dW`$ is independent of $`f^{}b`$ at time $`t`$, and has zero expectation. So, taking the conditional expectation of eq. (73),
$$E[f(X_{t+dt}f(X_t)|X_t=x]=E[\frac{f}{x}a]dt+\frac{1}{2}E[b^2f^{\prime \prime }|X_t=x]dt.$$
(74)
Since $`a,b,f,f^{},f^{\prime \prime }`$ are functions of $`X_t,t`$ they become sure functions, evaluated at $`x`$ under the conditioning; thus we get the equation for the increment $`f(x,t+dt):=E[f(X_{t+dt}|X_t=x]`$:
$$(f(x,t+dt)f(x))/dt:=f=(1/2)b(x,t)^2f^{\prime \prime }+a(x,t)f^{}.$$
This is Kolmogorov’s backward equation, which applies to the dynamics of the process. To get the dynamics of the probability density, we take the dual operator $`^{}`$, defined by
$$p(x,t)f(x,t)𝑑x=^{}p(x,t)f(x,t)𝑑x$$
which on integration by parts, and discarding the boundary term at $`\mathrm{}`$ gives
$$^{}f:=\frac{1}{2}\frac{}{x}\left(b(x,t)^2\frac{}{x}f\right)+\frac{}{x}\left(a(x,t)f\right).$$
Since $`f`$ was arbitrary, we see that $`p(x,t)`$ satisfies the forward equation in the weak sense (after smoothing with a test-function $`f`$). It is known from the theory of elliptic regularity, that any weak solution is a strong solution. If $`a`$ and $`b`$ are constants, we arrive at the Smoluchowski equation, and the continuum version of (51):
$$p(x,t)=E[p(X_t,0)|X(0)=x].$$
This representation for the solution of the pde gives an immediate proof that the solution remains non-negative if the initial condition is non-negative, since $`p(X_t,0)0`$; also, one sees that the time-evolution must be a contraction in the $`L^{\mathrm{}}`$-norm, and the $`L^2`$-norm as the conditional expectation is a projection.
Sometimes, we can rewrite the solution $`X_t`$ in terms of time-translation $`\omega \omega T_t`$ if we modify the measure . Suppose $`\mu ^{}`$ is absolutely continuous relative to $`\mu `$. Then there exists an adapted process $`u(t)`$ in $`(\mathrm{\Omega },,\mu )`$ such that
$$dX_t=dW_t+u(t)dt,X_0=0,$$
(75)
has a weak solution $`X_t`$ whose law is the same as $`Y_t(\omega ):=\omega (t)`$ as a r. v. on $`(\mathrm{\Omega },,\mu ^{})`$. Then the Radon-Nikodym derivative is
$$d\mu ^{}/d\mu ^=\mathrm{exp}\left[_0^tu(s)𝑑X_s(1/2)_0^tu(s)^2𝑑s\right].$$
(76)
Conversely, if $`u`$ is such that the r. h. s. of (76) has Wiener expectation 1, (as will happen if $`u`$ is bounded), then there exists an absolutely continuous measure $`\mu ^{}`$ given by (76), such that $`T_t^{}`$ on $`(\mathrm{\Omega },,\mu ^{})`$ produces a weak solution to (75). This is the Girsanov-Cameron-Martin theorem.
This change of measure is closely linked to the change of ground state in the corresponding quantum theory, when an interaction is introduced. We see this in the Feynman-Kac formula, below.
One can, using similar methods, integrate adapted functions relative to $`dM`$, where $`M`$ is any martingale. The stochastic integral has other variants, such as the Stratonovitch version ; one can also integrate non-adapted processes, subject to other conditions (Skorokhod), or use another noise which is not quite a martingale . The Ito version has an interesting interpretation in mathematical finance. Suppose that the price of an asset is a random process $`S_t`$, and it obeys the Ito equation
$$dS_t=a(S_t,t)dt+b(S_t,t)dW_t.$$
If we choose to hold $`\phi (t)`$ units of this asset, our portfolio at time $`t`$ is worth $`\phi (t)S_t`$. The change in the value of our portfolio in time $`dt`$ is $`d(\phi (t)S_t)`$, and we evaluate this as $`\phi (t)dS_t`$, because we do not change our holding $`\phi (t)`$ until after we have seen the change in the asset price. Here $`dS_t=S_{t+dt}S_t`$, so the total change in the asset over the time-interval $`[0,T]`$ is the Ito integral $`_0^T\phi (s)𝑑S_s`$, in which $`\phi `$ is adapted and the stochastic increment is the forward difference.
We now give a brief account of the Feynman-Kac formula . Feynman related the quantum transition amplitude $`\psi ,e^{iHt}\varphi `$ to the integral over histories of $`\psi ,e^{i{\scriptscriptstyle L(s)𝑑s}}\varphi `$ where $`L`$ is the Lagrangian . The trouble is, the Feynman ‘integral’ over histories is not based on measures, but on oscillatory integrals, and these rarely converge. In quantum physics, the spectrum of the energy is bounded below (at least at zero temperature). This expresses the stability of the theory. It follows that the unitary time-evolution group $`e^{iHt}`$ has an analytic continuation to complex times with negative imaginary part. In particular this is true of all the matrix elements of this operator. This is the underlying fact used in Euclidean quantum field theory, but also holds for quantum systems without any large symmetry group; only invariance under time-evolution is needed. In particular, we can consider the group for negative imaginary times, giving a semigroup $`e^{Ht}`$. The large-time behaviour of this is very good. This was used by Nelson to study certain perturbations of the free Hamiltonian: it is easier to study perturbations of a contraction semigroup than a unitary group.
###### Theorem 4.3
Let $`H_0=\frac{1}{2}\frac{^2}{x^2}`$ and $`V`$ be a real-valued $`C^{\mathrm{}}`$-function of $`x𝐑`$, vanishing at $`\mathrm{}`$. Then $`H_0+V`$ is self-adjoint on Dom$`H_0`$ and
$$\psi ,e^{(H_0+V)t}\phi =\overline{\psi (\omega (0))}\phi (\omega (t))\mathrm{exp}\left(_0^tV(\omega (s))𝑑s\right)𝑑\mu .$$
(77)
For the proof, see or . For a version within quantum probability, see . In this way, we construct an interacting theory in terms of a path integral using the Wiener measure $`\mu `$, weighted with an exponential function. The similarity with the Gibbs state of a system of paths in a potential $`V`$ is noteworthy. Suppose that $`V=0`$ outside a region $`\mathrm{\Lambda }`$, and converges to $`+\mathrm{}`$ inside $`\mathrm{\Lambda }`$. Then we see from Feynman-Kac formula that the measure vanishes on all paths that enter the region $`\mathrm{\Lambda }`$. After a normalisation, the weighted measure thus becomes the conditional Wiener measure, $`\mu (.|\omega (t)\mathrm{\Lambda }\text{ for all }t)`$. The formula then solves the heat equation subject to the condition of no-flow through the boundary $`\mathrm{\Lambda }`$. We do not need to find this conditioned measure to use the formula; we can, for example, use the Monte Carlo method, and sample paths by computer, rejecting any that enter $`\mathrm{\Lambda }`$; we can also use the conditioned measure to get results on monotonicity, since e. g. if the region $`\mathrm{\Lambda }`$ is enlarged, obviously more paths are allowed, and so the integral of a positive integrand is increased. This relation with pde’s has developed into the subject called potential theory , and is one of the tools used in constructive quantum field theory .
Dyson saw the usefulness of using imaginary time in quantum field theory . Schwinger had introduced the idea of the Euclidean quantum field as a way of avoiding the difficulties of Lorentz invariance; these are replaced by invariance under $`O(4)`$, the orthogonal group; since we analytically continue all the time-ordered functions to imaginary time, time $`t`$ gets replaced by $`it`$, often attributed to Minkowski. In fact, Minkowski did not know about the consequences of positive energy; he did not analytically continue anything, but simply replaced time by $`ix_4`$, where $`x_4=it`$. This means that he considered the complex $`O(4)`$, and the invariance group was a particular subgroup $`L`$ of it consisting of matrices some of whose entries were complex. In fact, $`L`$ is isomorphic to the real Lorentz group, and is thus non-compact. Nothing has been gained by Minkowski’s trick. Indeed, lots of confusion arose in electromagnetic texts up until recently, where other four-vectors such as $`A^\mu `$ were regarded as having a complex zero<sup>th</sup> component. Schwinger’s programme of Euclidean field theory is a special case of a theory developed by Wightman , in which the expectation values of the field are proved to have an analytic continuation in all the space-time components, into a domain that includes real position variables and purely imaginary time.
Symanzik started the mathematical programme of Euclidean quantum field theory. Glimm and Jaffe developed constructive quantum field theory using their theory of the perturbation of contraction semigroups. This is almost a Euclidean point of view. A beautiful probabilistic version of the subject resulted from Nelson’s rewrite of Symanzik’s programme. Let us outline this for the quantum mechanics of an oscillator.
We start with the self-adjoint Hamiltonian
$$H=H_0+V=\frac{1}{2}(\frac{^2}{q_j^2}+q^21)$$
(78)
Then the lowest eigenvalue, say $`0`$, is simple; let $`U(t)=e^{iHt}`$ and let $`\psi _0`$ be the eigenfunction of the eigenvalue $`0`$. Then $`\psi _0>0`$ holds. That is, there are no nodes in the ground state, a kind of Perron-Frobenius theorem. It is then convenient to replace the Hilbert space of the theory, $`=L^2(𝐑,dq)`$ by the unitarily equivalent space $`^{}=L^2(𝐑,|\psi _0(q)|^2dq)`$. The unitary map $`W:^{}`$ is given by $`(W\psi )(q)=\psi (q)/\psi _0(q)`$. This is obviously organised so that $`W\psi _0=1`$, the unit constant function in $`^{}`$. An observable $`A`$, acting on $``$, is converted to $`A^{}=WAW^1`$. The operator $`q`$ commutes with $`W`$, so is unchanged; but its canonical conjugate, $`p`$ does not commute with $`W`$, and neither does $`q(t):=U(t)qU(t)`$, so these operators do not take the usual Schrödinger form on $`^{}`$.
The positivity of the energy ensures that the Wightman function $`1,q(t_1)\mathrm{}q(t_n)1`$ has an analytic continuation to purely imaginary times,
$$t_j=is_j,\text{ such that }s_js_{j+1}>0,s_j𝐑,j=1,\mathrm{}n1.$$
(79)
Define the Schwinger function
$$S_n(s_1,\mathrm{},s_n)=W_n(is_1,\mathrm{},is_n)$$
(80)
at points given by eq. (79); we take $`S_n`$ to be defined by symmetry in the other regions; since the $`w_n`$ are symmetric at real points, the $`n!`$ analytic functions coincide at a common boundary of real dimension $`n`$. So by the edge-of-the-wedge theorem there is one common analytic function coinciding with these Schwinger functions. Obviously, $`S_n`$ determines $`W_n`$, by the uniqueness of analytic continuation.
Then two properties hold: there is a stochastic process $`X(t)`$ such that $`S_n`$ is the $`n^{\mathrm{th}}`$ moment:
$$S_n(s_1,\mathrm{},s_n)=E[X(s_1)\mathrm{}X(s_n)];$$
Moreover, the process is stationary and Markovian; that is
$$E[X_t|_s]=E[X_t|_s],\text{ for }ts.$$
(81)
Here, $`_s`$ is the $`\sigma `$-tribe generated by $`X_r,rs`$, and $`_s`$ that generated by $`X_s`$. Neither of these properties is true for a general Hamiltonian theory, so they reflect somehow the Lagrangian origins of the theory.
We can recover the physical Hilbert space as the initial space, $`L^2(\mathrm{\Omega },_0,\mu )`$ generated by powers $`X(0)`$ acting on the vacuum, $`\psi _0`$ which is the function 1. Also $`q`$ is then multiplication by $`X(0)`$. The Hamiltonian can be recovered by the identity (c.f. (51))
$$e^{Ht}P(q)\psi _0=E[P(X(t))|_0]$$
(82)
for any polynomial $`P`$. This is the continuous version of the fact that the transition matrix of a Markov chain can be recovered as the conditional probability of one time-step. We find
$$\psi _0,q(t_1)q(t_2)\psi _0=(1/2)\mathrm{exp}\{i(t_1t_2)\}.$$
(83)
This leads by analytic continuation to
$$S(s_1,s_2)=(1/2)\mathrm{exp}\{|s_1s_2|\}=E[X(s_1)X(s_2)]$$
(84)
where $`X(t)`$ is the Ornstein-Uhlenbeck process.
Nelson was able to follow this programme for the free quantised field, and so rewrite the problem of finding solutions to relativistic quantum fields in terms of generalised random fields. A selection of good reading on this subject is .
## 5 Quantum Processes
Is friction a classical concept? ‘There is no friction in quantum systems: the ground state of the atom does not grind to a halt. The introduction of friction, e. g. the term $`\gamma \dot{x}`$ in Newton’s laws, is to account for atomic phenomena such as radiation of moving charges, in a very crude way. Such effects are treated exactly in quantum mechanics, and therefore frictional terms do not appear’. The view is still widespread but not universal among physicists. Friction does not appear in classical mechanics either if it is not put in.
A quantum process is, in a general way, a Hilbert space $``$ and a family of self-adjoint operators $`\{A(t)\}_{t0}`$ on $``$. A quantum field used as noise appeared in . Senitzky obtained the approximate dynamics of a quantum oscillator by reduction from the dynamics of a larger conservative system. He arrived at the following quantum Langevin equation with a Gaussian positive-energy quantum driving term $`(\phi (t),\pi (t))`$ (the noise):
$$\frac{dQ(t)}{dt}=\omega P(t)\gamma Q(t)+\phi (t)\frac{dP(t)}{dt}=\omega Q(t)\gamma P(t)+\pi (t).$$
(85)
He noticed that without the ‘noise’, the Heisenberg commutation relations fade with time: $`[Q(t),P(t)]=ie^{2\gamma t}`$; he considered this to be inconsistent with quantum mechanics. With the noise, the solutions obey $`[Q(t),P(t)]i`$ for all time. The noise was a free quantum field with constant energy spectrum from $`0`$ to $`\mathrm{}`$. This does not quite satisfy the requirement that the Heisenberg cummutation relations should hold for all time. In we found the general exact solution to this problem. A special case is
$$\phi (t)=2^{1/2}(a(t)+a^{}(t)),\pi (t)=i2^{1/2}(a(t)a^{}(t)),$$
where
$$a(t)=(2\gamma /\pi )^{1/2}_\omega ^{\mathrm{}}e^{ikt}a(k)𝑑k,[a(k),a^{}(k^{})]=\delta (kk^{}).$$
This has a constant energy spectrum from $`\omega `$ to $`\mathrm{}`$. The feature of this solution, and Senitzky’s approximate solution, is the relationship between the dissipation $`\gamma `$ and the correlation of the quantum noise, which at zero temperature is
$$a(s)a^{}(t)=\frac{2\gamma }{\pi }e^{i\omega (ts)}\frac{1}{ts+iϵ}.$$
This is called the fluctuation-dissipation theorem.
Lax used noise with all frequencies, with two-point function
$$a(s)a^{}(t)=\frac{\gamma }{\pi }\delta (ts).$$
This is closer to the classical white noise, in that the increments to the process are independent, and the field obeys a quantum version of the Markov property. It was to be used later by Hudson and Parthasarathy in a rigorous body of theory . As physics, it was criticised by Kubo and others, as violating the KMS condition, which comes from the axiom of positive energy . The correct treatment (at non-zero temperature) was obtained by Ford et al., by taking the limit of one oscillator coupled to a large system of oscillators (or a string ). This was truly the quantum Langevin equation, in that the noise is added only to the equation for $`P`$ and not to $`Q`$. This can also be obtained as a singular limit of the asymmetric solution given in . The quantum noises in are not martingales, and have not got independent increments. They do fit in to the axiomatic scheme offered in . In , Ford emphasizes the role played by causality. Instead of eq. (43), he considers the equation with memory
$$m\stackrel{..}{x}+_{\mathrm{}}^t\mu (ts)\dot{x}(s)𝑑s+V^{}(x)=F(t).$$
(86)
The fact that the dissipation due to the future must be zero leads us to consider only those $`\mu `$ which vanish for negative argument. Perhaps this is a lesson for those who like to work on Lax’s version.
The first work to use the words ‘continuous tensor products’ (CPT) was . The notable conclusion was that the theory can always be embedded in a boson Fock space; the Wiener chaos is an example of this. We start with a definition of current algebra, or better, current group. Let $`G`$ be a Lie group, with Lie algebra $`𝒢`$, and denote by $`𝒟(G)`$ the set of $`C^{\mathrm{}}`$-maps from $`𝐑^n`$ into $`G`$, being the identity outside a compact set. We can furnish $`𝒟(G)`$, the current group, with a group law by pointwise multiplication: $`fg(x):=f(x)g(x)`$. This group has a Lie algebra, denoted $`𝒟(𝒢)`$, which is the set of all $`C^{\mathrm{}}`$-maps $`F:𝐑^n𝒢`$, of compact support, under the pointwise bracket
$$[F(f),G(g)]:=[F,G](fg).$$
The problem is to find representations of the current groups and the current algebras, by unitary or self-adjoint operators respectively.
Guichardet proposed a construction for the tensor product of Banach spaces or algebras, labelled by a continuous index. The first thing is to define, if possible, the continuous product of $`f(x)`$ over $`x𝐑^n`$, when $`f`$ has compact support. He tries
$$\underset{x}{}f(x):=\mathrm{exp}\left(\mathrm{log}f(x)𝑑x\right).$$
(87)
For Hilbert spaces, we wish to define the scalar product between two fields of vectors $`\psi (x)`$ and $`\varphi (x)`$. We put $`f(x)=\psi (x),\varphi (x)`$ and use eq. (87), provided that $`f(x)=1`$ outside a compact set and we take $`\mathrm{log}1=0`$ (the principal branch). We then need to be able to extend the scalar product to linear combinations of product vectors. In , we give an example of a non-existent Hilbert continuous product, in that the positivity fails on linear combinations. Guichardet presents a class of Hilbert spaces for which the construction works, and writes the Fock representation of the free field in these terms. To explain his examples, let $``$ be a Hilbert space, and $`\mathrm{\Gamma }()`$ the Fock space over $``$. We define the map $`\mathrm{exp}\mathrm{\Gamma }()`$ by
$$\mathrm{exp}\varphi :=1\varphi 2^{1/2}\varphi \varphi \mathrm{}(n!)^{1/2}^n\varphi \mathrm{}$$
(88)
The $`\mathrm{exp}\varphi \mathrm{\Gamma }()`$ is called the coherent state determined by the one-particle state $`\varphi `$. One shows that they form a total set (their span is dense) in $`\mathrm{\Gamma }()`$; clearly,
$$\mathrm{exp}\varphi ,\mathrm{exp}\psi =\mathrm{exp}\varphi ,\psi .$$
(89)
In , the Hilbert spaces $`_x`$ at each point is itself the Fock space $`\mathrm{\Gamma }(H)`$ of a Hilbert space $`H`$, and the family $``$ consists of coherent states at each point. This is a special case of the construction given below.
The case of current groups was treated in . We give here a special case when the continuous label is $`𝐑`$, interpreted as time; we start with $`(,U,\psi )`$, where $``$ is a Hilbert space, $`\psi `$, and $`U`$ is a representation of $`G`$ on $``$ such $`\{U(g)\psi ,gG\}`$ has dense span. The triple $`(,U,\psi )`$ is called a cyclic representation of $`G`$.
We say that $`(_1,U_1,\psi _1)`$ and $`(_1,U_2,\psi _2)`$ are cyclic equivalent if there exists a unitary isomorphism $`W:_1_2`$ such that for all $`gG`$,
$$WU_1(g)W^1=U_2(g);W\psi _1=\psi _2.$$
A cyclic representation gives us a function on the group, analogous to the characteristic function of a random variable. Indeed, it reduces to the characteristic function when the group is $`𝐑`$. Thus
$$C(g):=\psi ,U(g)\psi .$$
(90)
Let Span$`G`$ denote the complex vector space of finite formal sums of elements of $`G`$. Then $`C`$ is continuous and of positive type on Span$`G`$, which determines $`(,U,\psi )`$ up to cyclic equivalence. Conversely, a continuous function $`C`$ of positive type on $`G`$ determines a cyclic representation $`(,\pi ,\psi )`$ related to $`C`$ by eq. (90). The construction is very similar to the proof of the GNS representation. First, we construct the vector space, Span$`G`$, and furnish it with the scalar product, determined by its values on the linearly independent elements $`g_1,g_2,\mathrm{}`$, by
$$g_i,g_j=C(g_i^1g_j);$$
we complete Span$`G`$ in the norm (or, if a semi-norm with kernel $`K`$, we complete the quotient Span$`G/K`$), giving the space $``$. Then we choose $`\psi `$ to be the identity of the group. The operator $`U(g)`$ can be defined on Span$`G`$ as left multiplication; this is easily shown to be unitary, and so can be extended to the whole space to get the representation $`U`$ of $`G`$.
In an infinite tensor product over a discrete index, von Neumann was able to end up with a separable Hilbert space only by labelling a special vector, say $`\psi _x`$ in each factor $`_x`$, and then considering products $`\varphi _x`$ of vectors in a subset $`\mathrm{\Delta }`$ that at infinity are close to $`\psi _x`$. Only then does the infinite product $`\varphi _x,\psi _x`$ converge. The tensor product then carries the labels $`\{\psi (x),\mathrm{\Delta }\}`$. Guichardet used a similar idea for the continuous product. We are less ambitious, in that we ask for the tensor product of a cyclic representation $`(,U,\psi )`$ of a group. We use the same representation at each point of the time axis, because we want to get a stationary quantum process. We then define the function $`C:𝒟(G)𝐂`$ as
$$C(g(.):=\underset{x}{}\psi ,U(g(x))\psi ,$$
(91)
which is well defined if we choose at each $`x`$ one branch of the logarithm. To get a representation of the current group, it is necessary and sufficient that this be of positive type on the current group, in which case we say that the CTP exists. We also want the function to be extendable to step functions, constant in an interval $`[s,t]`$ and the identity outside. For such a $`g(.)`$, we divide an interval $`[s,t]`$ into an arbitrary number, $`N`$, of equal intervals; then $`C(g)`$ is the product of $`N`$ equal factors, each a characteristic function on $`G`$. Thus $`C`$ has the property that it has an $`N^{\mathrm{th}}`$ root that is also a characteristic function. Such $`C`$ is called infinitely divisible. By the relation of characteristic functions to cyclic representations, we are able to transfer the concept of $`\mathrm{}`$-divisibility to cyclic representations:
###### Definition 5.1
Let $`(,U,\psi )`$ be a cyclic representation of a group $`G`$. We say that it is $`\mathrm{}`$-divisible if, for any integer $`N>0`$, there is another cyclic representation $`(^{1/n},U^{1/n},\psi ^{1/n})`$, called the $`n^{\mathrm{th}}`$-root, such that $`(,U,\psi )`$ is cyclically equivalent to
$$(^{1/n},U^{1/n},\psi ^{1/n})$$
where the tensor product is over $`N`$ factors, and the resulting representation is restricted to the cyclic subspace spanned by the group acting on the product vector $`\psi ^{1/n}`$.
We see immediately that if for some $`n`$ the $`n^{\mathrm{th}}`$ root of the representation exists, then it is unique (up to cyclic equivalence). For, the characteristic function of two $`n^{\mathrm{th}}`$-roots, $`C_1,C_2`$ say, both satisfy $`C_i^n=C`$, and so their ratio is $`\omega _n`$, an $`n^{\mathrm{th}}`$-root of unity. But this violates positivity unless $`\omega _n=1`$. The converse also holds: if $`C`$ is the product of $`n`$ functions of positive type, then $`C`$ itself is of positive type. In we assumed that $`C(g)`$ never vanishes; we prove this later.
Following we can now give the criterion for the positivity of the scalar product in a continuous tensor product $`^{\psi ,\mathrm{\Delta }}_x`$ of cyclic group representations, relative to the cyclic vector $`\psi `$ and the set of states $`\mathrm{\Delta }:=\{U(g)\psi :gG\}`$.
###### Theorem 5.2
The following are equivalent.
1. The function $`C(g)`$ is a continuous function of positive type on $`G`$ with $`C(e)=1`$, and is $`\mathrm{}`$-divisible.
2. There exists an $`\mathrm{}`$-divisible cyclic representation $`(,U,\psi )`$ of $`G`$ such that $`C(g)=\psi ,U(g)\psi `$.
3. $`^{\psi ,\mathrm{\Delta }}`$ exists.
4. $`C(e)=1`$ and a branch of $`\mathrm{log}C(g)`$ is a conditionally positive function on $`G`$.
In (3) and (4) the branch of the logarithm is determined by which root of $`C`$ is of positive type. Only the item (4) needs explanation. A function $`F(g)`$ on a group is said to be conditionally positive if
$$\underset{ij}{}\overline{z}_iz_jF(g_i^1g_j)0$$
for all $`n`$-tuples $`(g_1,\mathrm{},g_n)`$ of group elements and all complex $`n`$-tuples
$`(z_1,\mathrm{},z_n)`$ summing to zero: $`_iz_i=0`$.
To sketch the proof, if $`C`$ is $`\mathrm{}`$-divisible, and $`C=e^F`$, then $`C^s`$ is also of positive type, for all small $`s>0`$. Then
$$\underset{ij}{}\overline{z}_iz_j(1+sF_{ij}+\mathrm{})0,$$
(92)
and so if $`_iz_i=0`$, we get that $`F`$ is conditionally positive semidefinite. For the converse, if $`F`$ is conditionally positive definite, then $`e^F`$ is of positive type for all $`s>0`$, see , page 280.
The following result is called an Araki-Woods embedding theorem , because of the similarity with , (but with different hypotheses). We remark that under the above conditions $`F`$ is conditionally positive semidefinite; then the function
$$g,h:=F(g^1h)F(g)F(h^1)$$
(93)
is of positive type, and so can be used to define a semi-definite form on Span$`G`$ by sesquilinearity.
Let $`𝒦`$ be the (separated, completed) Hilbert space formed using this as scalar product on Span$`G`$. Let $`G_0`$ be the subgroup of $`G`$ such that $`U(g)\psi =e^{i\lambda }\psi `$ for some real $`\lambda `$. We see that $`g,h`$ vanishes on Span$`G_0`$, and defines a scalar product on $`\text{Span}G/(\text{Span}G_0)`$, (perhaps after identifying vectors of zero norm with zero). We then complete this to give a Hilbert space, $`𝒦`$. We see that the equivalence class of the identity $`eG`$ is the zero vector in $`𝒦`$. The original cyclic representation $`(,U,\psi )`$ can then be embedded in the Fock space over $`𝒦`$, as follows: define the map $`W`$ from $``$ to $`\mathrm{\Gamma }(𝒦)`$ by its action on the total set $`U(G)\psi `$:
$$W(U(g)\psi )=C(g)\mathrm{exp}[g],gG.$$
(94)
One easily sees that this preserves the scalar product, using (93). Thus it can be extended by linearity and continuity to $``$. We see that the cyclic vector $`\psi `$ is mapped to the ‘vacuum’ vector $`\psi _0`$ of the Fock space. As for the group action, we use the fact that $`G/G_0`$ is a $`g`$-space, with left multiplication $`\tau _g[h]=[gh]`$. This defines an action $`\mathrm{exp}\{\tau _g\}`$ on the Fock space as usual, by its actions on the coherent vectors:
$$\mathrm{exp}\{\tau _g\}\mathrm{exp}[h]:=\mathrm{exp}[gh].$$
Define an operator $`U^{}`$ closely related to $`\mathrm{exp}\{\tau _g\}`$:
$$U^{}(g)C(h)\mathrm{exp}[h]:=C(gh)\mathrm{exp}[gh]$$
(95)
Then by calculation one sees that $`(,U,\psi )`$ is cyclically equivalent to the cyclic subspace of $`(𝒦,U^{},\psi _0)`$; $`W`$ intertwines $`U`$ and $`U^{}`$ and maps $`\psi `$ to $`\psi _0=\mathrm{exp}[e]`$. From the unitarity of $`U^{}`$ we see that $`|C(g)|^2=e^{[g],[g]}0`$.
The Gaussian measure is $`\mathrm{}`$-divisible, and the representation of the translation group, $`U(\lambda )`$, with Gaussian cyclic vector $`\psi (x)=(2\pi )^{1/4}e^{x^2/4}`$, is $`\mathrm{}`$-divisible. The corresponding CTP contains Brownian motion §2; the continuous product $`_0^tU(\lambda )`$ is the exponential martingale. A representation of the oscillator group is $`\mathrm{}`$-divisible, and the CTP of this is the free non-relativistic quantised fields .
H. Araki independently obtained similar results . Instead of $`\mathrm{}`$-divisible cyclic representations of groups, Araki started with a factorizable representation of current algebra. He remarked that, putting $`[g]=\varphi _g`$ the map $`V(g)\varphi _h:=\varphi _{gh}\varphi _g`$ is a unitary representation of $`G`$; this is proved on the vectors $`\varphi _h`$, $`\varphi _k`$ by use of (93). The equation expresses that the map $`g\varphi _g𝒦`$ is a one-cocycle of the group, with values in $`𝒦`$. We briefly explain this.
So, let $`G`$ be a group, and let $`𝒦`$ be a Hilbert space on which $`G`$ acts by unitary operators $`gV(g)`$. We shall write the left action $`\varphi V(g)\varphi `$ as left multiplication, $`\varphi g\varphi `$. The right action, which appears in the general theory of group cohomology, is taken to be trivial: $`\varphi g=\varphi `$. An $`n`$-cochain with values in $`𝒦`$ is a map from $`G^n`$ into $`𝒦`$, that is, it is a function of $`n`$ group elements with values in $`𝒦`$, thus: $`\varphi (g_1,\mathrm{},g_n)`$. We shall need only the $`0`$-cochains, which make up the space $`C^0:=𝒦`$ of vectors independent of $`g`$, and the 1-cochains, which are vector fields $`\varphi (g)𝒦`$ defined on the group. These make up the vector space $`C^1`$. We shall also need the $`2`$-co-chains, when $`𝒦=𝐂`$; these are complex-valued functions of two group elements. We see that the cochains of any degree $`k`$ form a vector space $`C^k`$. Fundamental to any cohomology theory is the coboundary operator, which is a linear map, $`\delta :C^kC^{k+1}`$, so increasing the degree of the cochain. It obeys $`\delta ^2=0`$. In the case of a group $`G`$ and a left and right action of $`G`$ on $`𝒦`$, $`\delta `$ is the linear map defined on $`C^0`$ by
$$(\delta \varphi _0)(g)=g\varphi _0\varphi _0g.$$
On $`C^1`$, $`\delta `$ is the linear map defined by
$$(\delta \varphi _1)(g_1,g_2)=g_1\varphi _1(g_2)\varphi _1(g_1g_2)_+\varphi _1(g_1)g_2.$$
On $`C^2`$, $`\delta `$ is the linear map defined by
$$(\delta \varphi _2)(g_1,g_2,g_3)=g_1\varphi _2(g_2,g_3)\varphi _2(g_1g_2,g_3)+\varphi _2(g_1,g_2g_3)\varphi _2(g_1,g_2)g_3.$$
The vector space of cocycles of degree $`k`$ in a vector space $`𝒦`$, with left and right actions $`\tau _1,\tau _2`$, is denoted $`Z^k(G,𝒦,\tau _1,\tau _2)`$. One checks that $`\delta ^2=0`$. A coboundary of degree $`k`$ is a vector function of the form $`\delta \psi `$, where $`\varphi `$ is a cochain of degree $`k1`$. The coboundaries of degree $`k`$ form the vector space $`B^k(G,V,\tau _1,\tau _2)`$. Since $`\delta ^2=0`$, we see that every coboundary is a cocycle. If the converse holds, the cohomology group $`H^k:=Z^k/B^k`$, is trivial. One sees that if $`\varphi `$ is a one-cocycle in $`C^1(G,𝒦,V)`$, then $`\varphi (g_1^1),\varphi (g_2)`$ is a two-cocycle in $`C^2(G,𝐂,I)`$.
A $`2`$-cocycle $`\sigma (g,h)`$ with values in the unit circle is also called a multiplier for the group. A multiplier representation of a group $`G`$ is a map $`gU(g),gG`$, such that $`U(g)U(h)=\sigma (g,h)U(gh)`$ for all $`g,hG`$. Although Wigner’s analysis of symmetry in quantum mechanics leads naturally to multiplier representations, their occurrence is sometimes called an ‘anomaly’ by physicists. When the CTP exists, we can represent the element $`g(.)`$ of the current group by the operator $`(U)_g`$, defined on the product vectors $`_xU(h(x))\psi _x`$ by
$$(U)_g(_xU(h(x))\psi _x):=_xU(g(x)h(x))\psi _x,$$
(96)
The space of the CTP is then $`\mathrm{\Gamma }(_{}\mathrm{exp}𝒦dx)`$ and $`\mathrm{\Delta }`$ consists of coherent states of the form $`\mathrm{exp}\varphi _{g(x)}`$. So we obtain a local representation of the current algebra. We get a multiplier when the branch of the logarithm in (91) obtained by the group law differs from the one needed to give a function of positive type on the group. This gives rise to an anomaly.
Araki showed that if $`\varphi `$ is the cocycle defined by the $`\mathrm{}`$-divisible representation $`U`$, then it is necessary that $`\mathrm{Im}\varphi (g_1^1),\varphi (g_2)`$ be a coboundary. Conversely, given a cocycle $`\varphi `$ with this property, it comes from an $`\mathrm{}`$-divisible representation. He proved that if $`G`$ is compact, then any cocycle is a coboundary, i. e. of the form $`\varphi _g=(V(g)I)\chi `$ for some $`\chi 𝒦`$. Use of a coboundary leads to a CTP of the form assumed by Guichardet . Araki was able to obtain analogues of the Levy formula (66) for various groups; for the group $`𝐑`$ this takes on a new meaning, as the decomposition of a cocycle into its parts coming from primitive cocycles, some algebraic and some topological. The topological cocycles are of the form $`(V(g)I)\chi `$; it is not a coboundary because $`\chi `$ is not in $`𝒦`$, but lies in a larger space that admits an extension of $`V`$; the $`V(g)I`$ brings the vector back into $`𝒦`$. Some groups, e. g. $`𝐑`$, also have cocycles called algebraic by Araki. For example, in the case $`G=𝐑`$, take $`𝒦=𝐂`$, and $`V(a)=I`$ for all $`aG`$. The cocycle is $`\varphi (a)=a`$. Then $`\varphi _a,\varphi _b=ab`$ is real, and $`C(\lambda )=\mathrm{exp}\{\frac{1}{2}\lambda ^2\}`$, the characteristic function of the Gaussian distribution. The Poisson part of the Levy formula comes from the coboundaries, and the Levy processes from the topological cocycles.
The question arises, given $`𝒦,V`$ and a cocycle $`g\varphi _g`$, can we construct a CTP? We can construct $`(,U(g),\psi )`$ from $`C`$, which can be regarded as a function such that $`C(e)=1`$ and the map $`C(h)\mathrm{exp}\varphi _hC(gh)\mathrm{exp}\varphi _{gh}`$ is unitary. The next big step was by Parthasarathy and Schmidt , who showed that given a cocycle there is indeed an $`\mathrm{}`$-divisible representation associated with it, but that it is a multiplier representation, with an $`\mathrm{}`$-divisible multiplier $`\sigma `$. The corresponding function $`C(g)`$ is $`\sigma `$-positive. This means that
$$\underset{ij}{}\overline{z}_iz_j\sigma (g_i^1,g_j)C(g_i^1g_j)0.$$
(97)
Naturally, this gives to a multiplier representation of the current group in general, and they found the multiplier; this leads to a tidier theory than , since the condition for the absence of multiplier can be dropped. Since the physical interpretation of a symmetry group leads (according to Wigner) to the ambiguity of the induced unitary representation up to a coboundary, the projective theory is certainly the right setting. Holevo has presented some similar concepts at the level of the algebra of observables, and found applications in quantum theory . Notable in the development was the work of Gelfand, et al. who used a cocycle of $`SL(2,𝐑)`$ to construct a factorisable representation of the corresponding current group. The whole theory is well explained in .
A theory of processes with independent increments and values in a Lie algebra $`𝒢`$ was developed in , extended to multiplier representations by Mathon and to Clifford algebras in . Corresponding central limit theorems were proved by Hudson, and Cushen and Hudson . A Lie process can be obtained by differentiation of the corresponding object for a Lie group. For example, near the identity any group element $`g`$ lies on a one-parameter subgroup generated by an $`X𝒢`$, and we write (Exp means the exponential map from $`𝒢`$ to $`G`$, not the Fock map) $`g(t)=\text{Exp}tX,g(0)=e,g(1)=g`$; given a representation $`U(g)`$ we get a representation of $`𝒢`$ by $`\pi (X)=d/dt[U(g(t)]_{t=0}`$. By Stone’s theorem, $`X`$ is self-sdjoint. However, given a cyclic vector $`\psi `$ for $`U`$ it does not follow that $`\psi `$ is cyclic for $`\pi `$, because of domain questions. Let $``$ be the universal enveloping algebra of $`𝒢`$. This is the nonabelian polynomial algebra, modulo the ideal generated by the commutators $`XYYX[X,Y]`$. Here, $`[X,Y]𝒢`$ is the Lie product, a polynomial of degree 1. A cyclic representation $`(,\pi ,\psi )`$ is determined (up to equivalence) by a positive linear functional, or state, on $``$:
$$X_1X_2\mathrm{}X_n\psi ,\pi (X_1)\pi (X_2)\mathrm{}\pi (X_n)\psi =W_n(X_1\mathrm{}X_n).$$
These are the noncommutative moments, or Wightman functions; they determine a representation, by the Wightman reconstruction theorem . They are generated by the characteristic function
$$C(\lambda )=\psi ,U(\text{Exp}\lambda _1X_1)\mathrm{}U(\text{Exp}\lambda _nX_n)\psi ,\lambda 𝐑^n.$$
(98)
Here, $`\{X_j\}`$ is a basis of the Lie algebra, and any moment out of order is determined by a derivative of $`C`$ and use of the commutation relations. The truncated functions $`W_T`$ are generated by $`\mathrm{log}C`$, and are related to $`W`$ by a formula similar to eq. (10), relating cumulants to the moments. Two cyclic representations with the same $`W`$, or the same $`W_T`$, are cyclic equivalent. The cumulants of $`\mathrm{exp}U`$ (the Fock construction) are the same as the moments of $`U`$; this follows from $`\mathrm{exp}U(g)\mathrm{exp}U(h)\psi )=\mathrm{exp}(U(gh)\psi )`$ and (98).
Given two representations $`U_1`$, $`U_2`$ of $`G`$, their tensor product $`U_1U_2`$, restricted to the diagonal subgroup of $`G\times G`$, gives the representation $`\pi _1I+I\pi _2`$ of $`𝒢`$. This led to the use of a coproduct, though it was not recognised as such until . Whereas a product on an algebra $`𝒜`$ is a linear map $`𝒜𝒜𝒜`$, a coproduct is a map $`𝒜𝒜𝒜`$. For Lie algebras the coproduct is $`XXI+IX`$. Then we say that a cyclic representation $`(,\pi ,\psi )`$ is $`\mathrm{}`$-divisible if for each $`N`$ there is another, $`(^{1/N},\pi ^{1/n},\psi ^{1/N})`$ such that $`(,\pi ,\psi )`$ is cyclically equivalent to $`\pi ^{1/N}I+I\pi ^{11/N}`$. Starting at $`N=2`$ this gives the concept of rational powers of $`\pi `$.
The differentiation of a CTP representation $`_tU_t(g(t))`$ of the current group leads to an ultralocal field . These are such that the truncated Wightman functions have the form
$$W_T(X_1(f_1)\mathrm{}X_n(f_n))=\kappa _n(X_1\mathrm{}X_n)f_1(t)\mathrm{}f_n(t)𝑑t.$$
(99)
Here, $`\{\kappa _n\}`$ are the cumulants of $`\pi =dU`$. The commutative analogue was analysed in . For Lie algebras, we found :
###### Theorem 5.3
The following are equivalent;
1) Eq. (99) defines a representation of $`𝒟(𝒢)`$.
2) The $`\kappa _n`$ are the cumulants of some $`\mathrm{}`$-divisible cyclic representation of $`𝒢`$.
3) The $`\kappa _n`$ are positive semi-definite on $`_1`$, the subalgebra of $``$ with identity omitted.
We note that (3) is the expression of conditional positivity at the algebraic level. Since the cumulants of $`\mathrm{exp}U`$ are the moments of $`U`$, we can get a set of $`\kappa _n`$ that obey the positivity (3) by using the moments of $`\mathrm{exp}U`$. These happen to have a positive extension to $``$: any conditionally positive functional is positive. Th. (5.3) has a cohomological version, which we outline.
Let $``$ be an associative algebra with identity, $`𝒦`$ a linear space and $`\tau `$ a representation of $``$ on $`𝒦`$. The $`p`$-cochain group $`C^p(,𝒦,\tau )`$ is the linear space of $`p`$-multilinear maps $`\varphi :\times \mathrm{}𝒦`$. The coboundary operator $`\delta :C^pC^{p+1}`$ is given by
$$(\delta \varphi )(X_1,\mathrm{},X_{p+1})=\tau (X_1)\varphi (X_2,\mathrm{},X_{p+1})+(1)^j\varphi (X_1,\mathrm{},X_jX_{j+1},\mathrm{},X_{p+1}).$$
Then $`\delta ^2=0`$ and we define as usual the cocycle group $`Z^{}:=\text{ker}\delta `$ and the coboundary group $`B^{}:=\text{Ran}\delta `$, and the cohomology as $`H^{}:=Z^{}/B^{}`$. ( means for any $`p`$). We see that a 1-cocycle is a map $`\varphi :𝒦`$ that satisfies $`\varphi (XY)=\tau (X)\varphi (Y)`$, and a 1-coboundary is a cocycle of the form $`\varphi (X)=\tau (X)\varphi _0`$ for some $`\varphi _0𝒦`$.
The states on $`_1`$ are positive elements of $`B^2(_1,𝐂,0)`$. Thus if $`(,\pi ,\psi )`$ is $`\mathrm{}`$-divisible, then its cumulants $`W_T`$ define a state on $`_1`$, and thus a scalar product: $`X,Y:=W_T(X^{}Y)`$. Here we define $`X^{}=X`$, since we want $`\pi `$ to represent the generators $`iX`$ of one-parameter subgroups by hermitian operators. Define $`𝒦`$ as the separated prehilbert space obtained from $`_1`$ as usual. Let $`\varphi :_1𝒦`$ be the embedding obtained from this, and define a \*-action $`\tau `$ of $`𝒢`$ on monomials by
$$\tau (X)\varphi (X_1\mathrm{}X_n):=\varphi (XX_1\mathrm{}X_n).$$
This states that $`\varphi `$ is a 1-cocycle. We then show that there is a bijection between the set of $`\mathrm{}`$-divisible cyclic representations $`(,\pi ,\psi )`$ of $`𝒢`$ and the triples $`(\tau ,\varphi ,\chi )`$, where $`\tau `$ is a hermitian representation of $`𝒢`$ on a prehilbert space $`𝒦`$, $`\chi `$ is a real character, and $`\varphi Z^1(_1,𝒦,\tau )`$ such that
$$\gamma :=\text{Im}\varphi (X),\varphi (Y)B^2(_1,𝐑,0).$$
(100)
In this bijection, $``$ is embedded in $`\mathrm{\Gamma }(𝒦)`$, $`\psi `$ is mapped to the Fock vacuum, and $`\pi `$ is related to $`\mathrm{exp}\tau `$ . So this is the Araki-Woods embedding theorem in this case. If (100) fails then we get a projective representation of $`𝒢`$, with multiplier $`\sigma `$ related to the cocycle $`\gamma `$ . We see that a cocycle for $`𝐑`$ is defined by a function $`\chi L^1(𝐑)`$ such that $`x\chi L^2(𝐑)`$. We thus see the origin of the condition near $`\alpha =0`$ in (66).
In we show that for Clifford algebras, the only possible $`\mathrm{}`$-divisible states are Gaussian (all cumulants above the second vanish). Here the coproduct is that of Chevalley, $`AAI+(1)^FIA`$ where $`F`$ is the degree of $`A`$, for elements of even or odd degree.
The algebraic theory was extended to associative algebras (that were not enveloping algebras of Lie algebras) by Hegerfeldt, who applied it to classify $`\mathrm{}`$-divisible quantum fields .
Goldin et al. have, independently of this work, constructed representations of a vector form of charge-current algebra, starting with the Fock space creation-annihilation operators ; they have been able to identify the representations in terms of the general anlysis of semi-direct products.
Schürmann introduced the concept of infinite divisibility for a representation of a Hopf algebra, and obtained essentially all the results of in this more general setting. Stochastic integrals for these processes were also constructed. For a clear account, see .
Voiculecsu developed the algebraic side into a subject called ‘free probability’ , as it lives in full Fock space, without symmetry or antsymmetry.
Albeverio and Hoegh-Krohn have constructed representations of current groups, and been able to replace the independence at every point by a covariance similar to the Nelson free field.
## 6 Quantum Stochastic Semigroups
These models of non-commutative noise, or quantum noise, are possible driving random terms for noisy quantum dynamics. What should we be looking for in a nonequilibrium stochastic quantum dynamics? From 1970, E. B. Davies made progress in formulating stochastic quantum dynamics . Suppose that the $`C^{}`$-algebra of observables is $`𝒜`$. We look at the Fokker-Planck equation in the classical case, and we see that we might expect a quantum stochastic process to be determined by a semigroup (in continuous or discrete time) of maps $`T_t`$ from the state space $`\mathrm{\Sigma }(𝒜)`$ to itself. It must map positive operators, the density matrices, to positive operators, and preserve the trace. We also do not want it to map a normal state to one of the finitely additive ones, so we require a stochastic map to obey
1. $`T`$ maps $`\mathrm{\Sigma }`$ to itself;
2. $`T`$ is linear;
3. In continuous time, $`(T_tI)A_10`$ as $`t0`$.
We can throw the action onto to algebra, to get the dual action $`T^{}:𝒜𝒜`$, by the requirement that for $`A𝒜`$,
$$T\rho ,A=\rho ,T^{}A\text{ for all }\rho \mathrm{\Sigma }.$$
$`T^{}`$ is automatically normal. We see that if $`𝒜`$ is abelian, then our conditions reduce to the properties needed for a classical stochastic process. It is obvious that a unitary time-evolution gives us a one-parameter family of stochastic maps, which can be extended to a group by including the inverses. We can get a large class of stochastic maps by forming mixtures of unitary groups; thus if $`\tau _j`$ is a family of invertible dynamics, then $`T=_j\lambda _j\tau _j`$ is stochastic if $`\lambda _j0`$ and $`\lambda _j=1`$. Any stochastic map is non-invertible if it is not unitary, and so is in this sense dissipative , p 25. In addition, in the quantum case, Kraus has argued that to get a satisfactory interpretation of the semigroup, $`T`$ must be completely positive. We say that a map $`T:𝒜𝒜`$ is $`n`$-positive if $`TI_n`$ is positive on the algebra $`𝒜𝐌^n`$. This is needed, since if our quantum system is described by the algebra $`𝒜`$, and there is an $`n`$-state quantum system far away, then the combined system will be described by $`𝒜𝐌^n`$, and the dynamics on the combined system could be $`TI_n`$. This must be positivity preserving, or else some state of the combined system will evolve to give negative probabilities. Since we want to avoid this for all $`n`$, we want $`T`$ to be $`n`$-positive for all $`n=1,2\mathrm{}`$. Such a condition is called complete positivity. It should be said that any positive map on an abelian algebra is always completely positive, so this concept only seriously arises in quantum probability.
Kraus showed that a map $`T`$ is completely positive if and only if $`T(A)`$ is a sum of maps of the form $`S_n^{}AS_n`$, where the $`S_n`$ are bounded; , p. 140.
The great result in the subject is the classification of continuous semi-groups of completely positive maps. In finite dimensions this was achieved in , and independently, by Lindblad, whose result holds for norm-continuous dynamics on $`C^{}`$ algebras. Their result is the quantum analogue of the heat equation, i. e. it is a dynamical equation for the density matrix. For a simple derivation, see . The result is:
###### Theorem 6.1
Let $`T_t`$ be a semigroup of completely positive stochastic maps on $`𝐌^n`$. Then there exists a Hermitian matrix $`H`$ and matrices $`S_j`$ such that the generator of the semigroup has the form
$$Z(A)=i[H,A]\frac{1}{2}(RA+AR)+\underset{j}{}S_j^{}AS_j,\text{ where }R=\underset{j}{}S_j^{}S_j.$$
(101)
This can be thrown onto the density matrices by duality. The first term $`i[H,A]`$ is non-dissipative, and is called the hamiltonian term. The second term is the dissipation.
It is very interesting that the first two terms of the Heisenberg expansion of the dynamics are of this form. Thus,
$`\left(e^{iHt}Ae^{iHt}A\right)t^1`$ $`=`$ $`i[H,A]{\displaystyle \frac{1}{2}}[H,[H,A]]t+O(t^2)`$
$`=`$ $`i[H,A]{\displaystyle \frac{1}{2}}(AS^2+S^2A)+SAS`$
$`\text{where }S=Ht^{1/2},`$
up to O(t), so it is of the form eq. (101) with $`R=S^2`$. In the anti-van Hove limit we replace $`S`$ by $`\lambda H`$.
It has been remarked that the commutator $`Ai[H,A]`$ is a derivation of the operator algebra, and so has many of the properties of a derivative. The double commutator has many of the properties of the second derivative, including some positivity, which mimics the positive spectrum of $`\mathrm{\Delta }`$ and the positivity improving properties of $`e^{\mathrm{\Delta }t}`$. Lindblad has analysed continuous semigroups of cp maps, with generator $``$, in terms of the ‘dissipation operator’, being minus the coboundary of $`L`$:
$$D(A,B):=\delta (A,B)=(AB)(A)BA(B).$$
(102)
He proves that $`T_t:=\mathrm{exp}(it)`$ is a continuous semigroup of cp maps if and only if $`D`$ is positive in the sense that
$$\underset{ij}{}C_i^{}D(A_i^{},A_j)C_j0\text{ for all }A_i,C_j𝒜.$$
(103)
Note the formal similarity with . Fannes and Quaegebeur have defined the concept of $`\mathrm{}`$-divisible completely positive mappings on groups, in which the function $`C(g)`$ is replaced by a cp operator. They prove an Araki-Woods embedding theorem for such structures.
Recall that for Markov chains, Brownian motion and Euclidean field theory, we can express the given semigroup as an isometric time-translation, followed by the conditional expectation onto the initial space. By using two-sided time, the isometries can be replaced by a unitary group. The finding of the appropriate unitary group is called the dilation of the semi-group. It is not unique, but there is a unique minimal one. . It would be nice to interpret the dilated system as representing the full physics of system plus environment, with a unitary evolution; the projection onto a subspace represents our loss of information due to incomplete knowledge. The ambiguity of the dilation then shows that several different models give the same (crude) coarse-grained dynamics. However, it will rarely be the case that a dilation has the good properties, such as positivity of the energy, needed for this interpretation.
This is illustrated in the quantum case, which in finite dimensions takes the form
###### Theorem 6.2
Let $`T_t`$ be a semigroup of cp stochastic maps on $`𝐌_n`$ acting on $``$. Then there exists a Hilbert space $`𝒦`$, a pure state $`\rho `$ on $`𝒦`$ and a one parameter unitary group $`V_t`$ on $`𝒦`$ such that
$$T_t(A)=E_\rho [V_t^{}(AI)V_t]$$
for all $`A`$ and all $`t𝐑`$.
This is proved by putting together Theorem (4.2) and §7.2 of . Note that the Hilbert space $`𝒦`$ is constructed by adding Wiener noise, and so is not finite-dimensional. The semi-group has been dilated to a unitary group on the Wiener space with two-sided time; the generator of time-evolution is not bounded below, since it has white spectrum. This does not represent an environment at any finite temperature. A special case is the dilation of the semigroup given by the anti-van Hove limit. In that case the process is given by
$$X(t)=(2\pi t\lambda ^2)^{1/2}e^{s^2/(2\lambda ^2t)}U(s+t)XU(st)𝑑s.$$
(104)
This has the interpretation as the Heisenberg evolution, but with the time $`t`$ slightly uncertain, and getting more uncertain in the future. This interpretation is only a slight variation on the methods used in the justification of the microcanonical state by ergodic theory. There, it is said that the atomic times are so small that we never measure an observable at a particular time; rather, we measure the average over the time $`0st`$ of the measurement thus: $`\overline{A}=t^1_0^tA(s)𝑑s`$. Since $`t`$ is so large compared with the atomic processes, we take the limit $`t\mathrm{}`$. This idea is a non-starter for non-equilibrium statistical mechanics, since if the limit exists it is time-independent. Instead, we may say that we cannot measure an observable at an exact time, but form the weighted average, with Gaussian weight, around the desired time $`t`$. The uncertainty in the Gaussian is $`\lambda ^2t`$, growing with time. $`\lambda `$ is the dissipation parameter. In models it turns out to be the hopping parameter of the atomic system.
Some authors limit the concept of quantum stochastic process to the case where the possible observed path of measurements themselves make up a classical process. The grounds for this is that the observations (in a set of repeated experiments) have actually been seen; these form the quantum record; take them to form a sample space. However, this is not true. The process $`X(t)`$ at different times might not commute, so the measurement of $`X(t)`$ alters the state (by collapse), and subsequent measurements are not those predicted by $`X(t+s)`$, $`s>0`$, as computed using the given initial state. It needs conditioning to the new information, and quantum conditional expectations only commute on abelian subalgebras. Moreover, one can measure $`X(t)`$ in one sampling and $`Y(t)`$ in another, where $`X`$ and $`Y`$ do not commute. No classical model would predict the statistics of the process; the classical theorist is liable to be hit by the EPR paradox in acute form. We regard $`X(t)`$ as the observable seen at time $`t`$ when no measurement has been made in $`\{s:0<s<t\}`$. So we cannot agree with the idea that the randomness itself is caused by the reduction of the wave-function due to continuous measurement; it might be due to interaction with a large other body, but not one designed to measure any particular observable.
Davies’s dilation of the Lindblad semigroup uses a number of independent Wiener processes to provide the set-up. The question arises whether there is a relation between quantum dynamical semigroups and a class of quantum stochastic differential equations, similar to the relation between the Fokker-Planck equation (72) and the sde (71). For this, we need a quantum version of Ito’s integral. In 1956, Umegaki defined the concept of conditional expectation in non-commutative integration theory . Let $`𝒜`$ be a von Neumann algebra with a semi-finite trace, and say an operator $`A`$ is integrable if $`\text{Tr}|A|<\mathrm{}`$. The vector space of integrable operators can be completed to form the space $`L^1(𝒜)`$. Segal and Nelson showed that there is a closed operator representing an element of the completion. Let $`𝒜_t`$ be an increasing family of subalgebras which generate $`𝒜`$ and are right continuous , such that the trace, restricted to each $`𝒜_t`$ is semi-finite. Then a conditional expectation relative to the trace is a linear map $`M:L^1(𝒜)L^1(𝒜_t),t0`$, such that
$$\text{Tr}(XA)=\text{Tr}(M_t(X)A)\text{ for all }A𝒜_t,XL^1(𝒜).$$
A martingale is a process $`X_t`$ of integrable operators such that
$$M_sX_t=X_s$$
for all $`0st`$. This concept can be generalised to a filtration of algebra with specified state, rather than trace.
Cuculescu proved a martingale convergence theorem for discrete time. Barnett obtained a martingale theorem for continuous time. This work persuaded us to look for examples of noncommuting martingales. Soon we found plenty within the theory of continuous tensor products . Let $`(,U,\psi )`$ be an $`\mathrm{}`$-divisible representation of a Lie group G, and consider $`_{t=0}^{\mathrm{}}_t`$ relative to the vector $`\psi _t`$ and the set $`\mathrm{\Delta }`$ of coherent vectors. Here, all factors are the same. To $`gG`$ we associate the family of unitary operators
$$V_t(g):=_0^tU(g)_t^{\mathrm{}}I.$$
(105)
We call such an operator simple, localised in $`[0,t]`$. Let $`𝒜_t`$ be the algebra generated by $`\{V_s(g)\}`$ with $`0st`$ and $`gG`$. Then for $`s<t`$ define the map $`M_s:𝒜_t𝒜_s`$ by continuous linear extension of $`M_s_{r=0}^tV_r(g)=_{r=0}^sV_r(g)`$. Then $`M_s`$ is a conditional expectation, and relative to $`M_s`$, the family $`V_t`$ is a martingale. Applied to $`G=𝐑`$ with $`\psi `$ a Gaussian state, $`V_t`$ is the exponential martingale of Brownian motion. When $`G`$ is the oscillator group, the lie algebra is spanned by $`P,Q,H`$ and a central element $`I`$. There is a representation by self-adjoint operators on $`L^2(𝐑)`$, with the ground state of the harmonic oscillator as cyclic vector. This is infinitely divisible, and the unitary operators in (105) are copies of the exponential martingale $`e^{iW_t}`$ for the subgroups generated by $`P`$ and $`Q`$, and is the Poisson exponential martingale for the subgroup generated by $`H`$ . This became known as the gauge process . All these martingales are defined on the total set of coherent states. Since they are unitary, they can be extended to an everywhere-defined unitary group, the generators of which are self-adjoint operators. This is the main technique of the Hudson-Parthasarathy calculus
Examples of martingales with trace were given in . Consider the Fock Fermi operators $`b(f),b(g)`$ with anticommutation relations $`[b(f),b^{}(g)]=f,g`$ for $`f,gL^2(𝐑_+)`$. The algebra generated by these and the Fock condition $`b(f)|0=0`$ is represented on antisymmetric Fock space over $`L^2(𝐑_+)`$ as the $`W^{}`$-algebra generated by the Fermi field $`\psi (f)=b(f)+b^{}(\overline{f})`$ acting on the Fock vacuum $`|0`$. The Clifford process is the set of operators
$$\mathrm{\Psi }(t):=\psi (\xi _{[0,t]}).$$
(106)
The non-commutative integration theory , taking the place of measure theory, is that based on the hyperfinite von Neumann factor of type $`II_1`$, which is furnished with a faithful trace $`\phi (A)=0|A|0`$. The completion of $`𝒜`$ in the norm $`A=\phi (A^{}A)^{1/2}`$ is denoted $`L^2(𝒜,\phi )`$. The projection $`M_t`$ from $`L^2(𝒜,\phi )`$ onto $`L^2(𝒜_t,\phi )`$ is the same as the projection from $`\mathrm{\Gamma }(L^2[0,\mathrm{}])`$ onto $`\mathrm{\Gamma }(L^2[0,t])`$; it obeys the laws for a conditional expectation, and $`\mathrm{\Psi }(t)`$ is a martingale.
The increments of $`\mathrm{\Psi }(t)`$ are independent, but anti-commute. Otherwise, all the properties are analogous to Brownian motion. The isometric time-evolution analogous to the left shift of the classical theory is that given by the map $`U_s:\mathrm{\Psi }(t)\mathrm{\Psi }(s+t)`$. The antisymmetric Fock space over $`L^2(𝐑)`$ carries a unitary extension of $`U_s`$, namely the second quantisation of translation in $`𝐑`$. We define an adapted process $`h(t)`$ to be a family of operators such that $`h(t)𝒜_t`$; it is simple if it can be expressed as
$$h=\underset{k=1}{\overset{n}{}}h_{k1}\chi _{[t_{k1},t_k)}\text{ on }[0,t).$$
(107)
We then define the stochastic integral of any simple adapted process, relative to $`\mathrm{\Psi }`$, to be that constructed in the manner of Ito, with the forward difference in $`d\mathrm{\Psi }`$:
$$_0^tf(s)𝑑\mathrm{\Psi }(s):=\underset{k=1}{\overset{n}{}}h_{k1}\left(\mathrm{\Psi }(t_k)\mathrm{\Psi }(t_{k1})\right).$$
(108)
As in Ito’s theory, what make it work is an isometry property:
###### Theorem 6.3
If $`h(t)`$ is a simple process made up of $`L^2`$ operators, then $`_0^th(s)𝑑\mathrm{\Psi }(s)L^2`$, and
$$_0^th(s)𝑑\mathrm{\Psi }(s)_2^2=_0^th(s)_2^2𝑑s.$$
The proof is similar to Ito’s. We use this to construct the integral of square-integrable adapted processes, and some $`L^p`$ processes, by extension to the completion of the space of simple adapted processes. The stochastic integral is the quantised field $`\mathrm{\Psi }`$, smeared with an operator $`h`$ rather than a test-function. There is a Doob-Meyer theorem: $`M_t^2`$ is the sum of a martingale, denoted by $`[M_t,M_t]`$ in classical theory, (NOT the commutator!) and an increasing process of bounded variation, denoted $`M_t,M_t`$. Any stochastic integral is a martingale, and we show the converse, that any $`L^2`$ martingale of mean zero is a stochastic integral. We also define the stochastic integral $`N(t)=_0^th(s)𝑑M(s)`$, where $`h`$ is adapted and square-integrable relative to $`M_t,M_t`$. Here, $`M`$ is an $`L^2`$-martingale. This representation of $`N`$ is unique; we then write $`h`$ as the stochastic derivative: $`h=N/M`$. We show that we can change variables in the integral: the stochastic Radon-Nikodym theorem .
We are able to show that the quantum sde
$$dX_t=F(X_t,t)dM_t+dM_tG(X_t,t)+H(X_t,t)dt$$
(109)
has a solution in $`L^2(𝒜,\phi )`$ for $`F,G,H`$ continuous, adapted and locally uniformly Lipschitz, for any martingale $`M_t`$ of degree $`n`$, and that the solution obeys the Markov property .
Manipulations of differentials are similar to the Ito calculus: $`(dt)^2=0=(dt)(d\mathrm{\Psi })`$; $`(d\mathrm{\Psi })^2=dt`$. Pisier and Xu have obtained ‘Burkholder-Gundy’ inequalities within this theory .
The central state $`\phi `$ of the Clifford algebra corresponds physically to an infinite temperature. For the CCR and CAR algebras, we constructed the stochastic integrals starting with quasifree states with no Fock part, using the non-central state in place of the trace . This theory is somewhat technical (‘unreadable’ ).
The general Lindblad semigroup can be dilated using the flow defined by a solution to a quantum stochastic equation in the sense of Hudson and Parthasarathy . It was extended to some unbounded cases by Belavkin . We now give a brief account of this, following Frigerio .
Let $`T_t=\mathrm{exp}(t)`$ be a semigroup of completely positive normal stochastic maps on the algebra $`()`$.
###### Theorem 6.4
There exists a Hilbert space $``$, a group $`\{\alpha _t:t𝐑\}`$ of -automorphisms of $`()`$ and a conditional expectation $`E_0`$ of $`()`$ onto $`()I_{}`$ such that
$$T_t(X)I_{}=E_0[\alpha _t(XI_{}0],X(),t𝐑.$$
(110)
The evolution $`\alpha _t`$ is a perturbation of the ‘free evolution’ $`\alpha _t^0`$ on $`()`$, of the form
$$\alpha _t(.)=U(t)\alpha _t^0(.)U(t)^{},$$
(111)
where $`\{U(t):t𝐑\}`$ satisfies the cocycle condition
$$U(t)\alpha _t^0(U(s))=U(s+t),t,s𝐫,$$
(112)
is unitary and is the solution of a qsde in the sense of . We give the details in the simplest case, eq. (101) with only one term $`S`$ in the sum. We take $`=\mathrm{\Gamma }(L^2𝐑)`$, with total the set of coherent vectors $`\mathrm{exp}\varphi :\varphi L^2(𝐑)L^1(𝐑)`$. We define the annihilation process, creation process and gauge process on this total set by
$`A(t)\mathrm{exp}\varphi `$ $`=`$ $`({\displaystyle _0^t}\varphi (s)𝑑s)\mathrm{exp}\varphi `$ (113)
$`A^{}(t)\mathrm{exp}\varphi `$ $`=`$ $`{\displaystyle \frac{d}{dϵ}}\mathrm{exp}\left(\varphi +ϵ\chi _{[0,t]}\right)|_{ϵ=0}`$ (114)
$`\mathrm{\Lambda }(t)\mathrm{exp}\varphi `$ $`=`$ $`{\displaystyle \frac{d}{dϵ}}\mathrm{exp}\left(e^{ϵ\chi _{[0,t]}}\varphi \right)|_{ϵ=0}.`$ (115)
The conditional expectation $`M_t`$ is as for the CTP, $`_{s=0}^t(_s)`$, based on the Fock vacuum, and $`A(t),A^{}(t)`$ are the creators and annihilators defined by the generators $`P,Q`$ of the Heisenberg subgroup of the oscillator group; $`\mathrm{\Lambda }`$ is the number operator.
We identify any operator $`X`$ in $`()`$ with its ampliation $`XI_{}`$, and any operator $`Y`$ with domain $`𝒟`$ with the algebraic tensor product $`I_{}Y`$. A family $`U(t)`$ is found by solving the qsde
$$dU(t)=U(t)\left[iS^{}dA(t)+iSdA^{}(t)+(iHS^{}S/2)dt\right],$$
(116)
with the initial condition $`U(0)=I`$. The structure of the equation is designed to ensure that the solution, defined on the set of coherent states, is continuous, unitary, and adapted. The term $`S^{}S/2`$ arises as the Ito correction, or as due to the Wick ordering . To ensure that $`\alpha _t`$ obeys the group law, the usual free evolution $`\alpha ^0`$ on $``$, the second quantisation of the translation group on $`L^2(𝐑)`$, is chosen. It is then proved that $`\alpha _s^0[U^{}(s)U(s+t)]`$ satisfies the same qsde as $`U(t)`$, and so, by uniqueness, must be $`U(t)`$. So $`U`$ satisfies the cocycle condition. On multiplying out, we see that $`\alpha _y`$ is a group.
The theorem for a semigroup with a finite number of operators $`S_j`$ follows a similar line.$`\mathrm{}`$
There is a fermionic version of this dilation .
Quantum stochastic calculus has become a mature field of mathematics. The approach of , rather than , has the disadvantage that the stochastic integrals are defined as operators only on a dense set. It is not always clear that they have a unique closed extension. This is overcome in by limiting the class of equations to those with unitary solutions. Another help in the analysis is by the use of Maassen kernels . Alternatively, one may give a meaning to these objects as maps between test-functions and distributions, using white-noise analysis.
One problem with this work, and this includes as well, is that the spectrum of the noise is white, so that random negative energy is added as well as positive energy. We saw that positive energy seems to exclude martingales . In fact, the KMS condition excludes the existence of a conditional expectation except in trivial cases. It has been remarked that it also excludes the Markov property and the regression theorem . Lindblad has remarked that for the oscillator, the KMS condition is not compatible with the axioms of dynamical semigroups. So to model random external forces in a real system, coupled to a heat-bath, the white noise sde is an approximation, that might be good if the time-interval is large compared with the memory time. These ideas are used to describe quantum systems like lasers, which are subject to external forces; this was the original intention of Senitzky and Lax. The modern version is described in . Since external forces introduce energy and entropy into a system, such models have two drawbacks:
1. The first law of thermodynamics is not obeyed.
2. The second law of thermodynamics is not obeyed.
This is the starting point of . One step of the linear dynamics is given by a bistochastic map $`\rho \rho T`$, so that entropy increases. We require that $`T^{}`$ maps any spectral projection of the energy to itself; this will preserve energy. To reduce the description, we then project the new state $`\rho T`$ onto the information manifold $``$ defined by the set of slow variables, to get the state $`\rho TQ`$. To preserve mean energy, the energy must be a slow variable. The map $`Q`$ is nonlinear and is interpreted as the thermalisation of the fast variables. Thus, after the map $`T`$, the system itself decides to find the best estimate $`\rho TQ`$ to $`\rho T`$ within $``$. The resulting map gives a nonlinear motion through the manifold, obeying the first and second laws of thermodynamics. This theory, called statistical dynamics, is still being explored . |
warning/0002/cond-mat0002296.html | ar5iv | text | # Delocalization of two-particle ring near the Fermi level of 2d Anderson model
\[
## Abstract
We study analytically and numerically the problem of two particles with a long range attractive interaction on a two-dimensional (2d) lattice with disorder. It is shown that below some critical disorder the interaction creates delocalized coupled states near the Fermi level. These states appear inside well localized noninteracting phase and have a form of two-particle ring which diffusively propagates over the lattice.
\]
Recently a great deal of attention has been attracted to the problem of interaction effects in disordered systems with Anderson localization . From the theoretical point of view the problem is rather nontrivial. Indeed, even if a great progress has been reached in the theoretical investigation of the properties of localized eigenstates still the analytical expressions for interaction matrix elements between localized states are lacking. In spite of these theoretical difficulties it has been shown recently that a repulsive or attractive interaction between particles can destroy localization and lead to a propagation of pairs in the noninteracting localized phase. This two interacting particles (TIP) effect has been studied recently by different groups and it has been understood that the delocalization of TIP pairs is related to the enhancement of interaction in systems with complex, chaotic eigenstates. Such an enhancement had been already known for parity violation induced by the weak interaction in heavy nuclei where the interaction is typically increased by a factor of thousand. However, since there the two-body interaction is really weak the final result still remains small. On the contrary, for TIP pairs in the localized phase the enhancement of interaction qualitatively changes the dynamics leading to a coherent propagation of TIP on a distance $`l_c`$ being much larger than the pair size and one-particle localization length $`l_1`$. The enhancement factor $`\kappa `$ is determined by the density of two-particle states $`\rho _2`$, coupled by interaction, and the interaction induced transition rate $`\mathrm{\Gamma }_2`$ between noninteracting eigenstates, so that $`\kappa =\mathrm{\Gamma }_2\rho _2`$. At $`\kappa 1`$ the interaction matrix element becomes comparable with two-particle level spacing and the Anderson localization starts to be destroyed by interaction. For excited states the TIP density $`\rho _2`$ is significantly larger than the one-particle density $`\rho `$ and the delocalization can be reached for relatively weak interaction if $`l_1`$ is large. However, when the excitation energy $`ϵ`$ above the Fermi level decreases then $`\rho _2`$ becomes smaller and it approaches the one-particle density $`\rho `$ at low energy: $`\rho _2ϵ\rho ^2`$ . As a result the value of $`\kappa `$ also drops with $`ϵ`$ so that the delocalization of TIP pairs practically disappears near the Fermi energy. This result has been found in in the approximation of the frozen Fermi sea created by fermions. Recent numerical studies of TIP pairs with short range interaction near the Fermi level confirmed these theoretical expectations.
In this paper we discuss another type of situation in which TIP delocalization (see Fig. 1) takes place mainly due to geometrical reasons and not due to the relation $`\rho _2\rho `$. As a result the TIP pair can be delocalized in a close vicinity to the Fermi level that opens new interesting possibilities for interaction induced delocalization in the localized noninteracting phase. To study this new situation we investigate a model with a long range attractive interaction between particles in the 2d Anderson model. In this case the particles can rotate around their center of mass, being far from each other and keeping the same energy, while the center can move randomly in two dimensions. As a result the system has effectively three degrees of freedom that makes it rather similar to the case of one particle in the 3d Anderson model where delocalization takes place at sufficiently weak disorder. A somewhat similar situation has been studied recently for particles with Coulomb interaction but only excited states were considered there and the delocalization was attributed to the large ratio $`\rho _2/\rho `$ . Here we show that in fact the conditions for delocalization are much less restrictive.
To illustrate the above ideas let us first discuss the case of only two particles with attractive interaction $`U(r)<0`$ in the 2d Anderson model described by the Schrödinger equation
$$\begin{array}{c}(E_{𝐧_1}+E_{𝐧_2}+U(𝐧_1n_2))\psi _{𝐧_1,n_2}+V(\psi _{𝐧_1+1,n_2}\\ +\psi _{𝐧_11,n_2}+\psi _{𝐧_1,n_2+1}+\psi _{𝐧_1,n_21})=E\psi _{𝐧_1,n_2}.\end{array}$$
(1)
Here $`𝐧_{1,2}`$ are the indices of the two particles on the 2d lattice with $`L^2`$ sites and periodic boundary conditions, $`V`$ is the hopping between nearby sites and the random on-site one-particle energies $`E_{𝐧_{1,2}}`$ are homogeneously distributed in the interval $`[W/2,W/2]`$. The long range attractive interaction depends only on the distance between particles $`r=𝐧_\mathrm{𝟏}𝐧_\mathrm{𝟐}`$ and is equal to a constant $`U<0`$ if $`|rR|\mathrm{\Delta }R`$ and zero otherwise. The value of $`r`$ is determined as the minimal inter-particle distance on the periodic lattice. Thus the interaction takes place only inside a ring of radius $`R`$ and width $`\mathrm{\Delta }R`$, and we assume that $`R\mathrm{\Delta }R1`$. For $`U=0`$ the eigenstates are given by the product of two one-particle (noninteracting) eigenstates which are always localized in 2d in a presence of disorder .
In the limit of very strong attractive interaction $`|U|V`$ the TIP coupled states form the energy band of width $`16V`$ around $`E|U|`$ (we consider only the states symmetric in respect to particle interchange). For states in this band the particles are located always inside the ring which center can move over the 2d lattice. Since $`|U|V`$ these states are decoupled from all other states with particles outside the ring. In the ring the Schrödinger equation is in fact rather similar to the case of 3d Anderson model of one particle. In this analogy the number of sites inside the ring $`M_R2\pi R\mathrm{\Delta }R`$ determines the effective number of 2d planes placed one over another in the third $`z`$-dimension (length size $`L_z=M_R`$). In this 3d model the effective strength of disorder is approximately $`2W`$ since the diagonal term is now the sum of two $`E_𝐧`$ values. Also one site is coupled with $`Z=8`$ neighbours contrary to $`Z=6`$ for 3d case (assuming $`\mathrm{\Delta }R1`$). Since in 3d the Anderson transition at the band center takes place at $`W_c=2.75ZV=16.5V`$ , we expect that TIP states inside the ring will be delocalized in the middle of the band when $`2W/ZV=2.75`$ that gives the transition at $`W_{c2}11V`$. This estimate is in agreement with numerical simulations of the model (1) . Of course, since the size in the third direction is finite the eigenstates will be eventually localized. But their localization length $`l_c`$ will make a sharp jump from $`l_c1`$ at $`W>W_{c2}`$ to $`l_c\mathrm{exp}(g)1`$ at $`W<W_{c2}`$ that follows from the standard scaling theory in 2d . Here $`g`$ is the conductance of the quasi-two-dimensional layer of width $`L_z=M_R`$. As usual $`g=E_c/\mathrm{\Delta }_1`$ where $`E_c=D/L^2`$ is the Thouless energy, $`\mathrm{\Delta }_1V/(L^2L_z)`$ is the level spacing and the diffusion rate in the lattice model is $`DV(V/W)^2`$. As a result for $`W<W_{c2}`$ the TIP delocalization length jumps to exponentially large value $`l_c\mathrm{exp}(2\pi R\mathrm{\Delta }R(W_{c2}/W)^2)`$. In these estimates we assumed that $`l_1>\mathrm{\Delta }R>1`$ since if $`\mathrm{\Delta }Rl_1`$ the majority of states inside the ring are noninteracting and can be presented as the product of one-particle eigenstates. We also note that for $`W<W_{c2}`$ there is an energy interval around the band center with delocalized states where the TIP ring diffuses with the rate $`D_2V(W_{c2}/V)^2`$. When $`W`$ decreases the mobility edge approaches the bottom of the band as it happens in 3d Anderson model.
The above arguments presented for the case $`|U|V`$ indicate that it is possible to have a similar TIP delocalization at moderate value of $`UV`$ near the Fermi level. To investigate this case we rewrite the equation (1) in the basis of the noninteracting eigenstates that gives
$`(E_{m_1}+E_{m_2})\chi _{m_1,m_2}`$ $`+`$ $`U{\displaystyle \underset{m_1^{^{}},m_2^{^{}}}{}}Q_{m_1,m_2,m_1^{^{}},m_2^{^{}}}\chi _{m_1^{^{}},m_2^{^{}}}`$ (2)
$`=`$ $`E\chi _{m_1,m_2}.`$ (3)
Here $`\chi _{m_1,m_2}`$ are eigenfunctions of the TIP problem written in the basis of one-particle eigenstates $`\varphi _m`$ with eigenenergies $`E_m`$. The matrix $`UQ_{m_1,m_2,m_1^{^{}},m_2^{^{}}}`$ represents the two-body matrix elements of interaction $`U(𝐧_1n_2)`$ between noninteracting eigenstates $`|\varphi _{m_1}\varphi _{m_2}`$ and $`|\varphi _{m_1^{^{}}}\varphi _{m_2^{^{}}}`$. The Fermi sea is determined by the restriction of the summation in (2) to $`m_{1,2}^{(^{})}>0`$ with energies $`E_{m_{1,2}^{(^{})}}>E_F`$, where $`E_F`$ is the Fermi energy related to the filling factor $`\mu `$. We choose the case with half filling $`\mu =1/2`$ for which $`E_F0`$. In this way our model corresponds to the approximation of frozen Fermi sea successfully used for the Cooper problem . As it was done by Cooper we also introduce the high energy cut-off defined by the condition $`1m_1^{^{}}+m_2^{^{}}M`$. This rule determines an effective phonon frequency $`\omega _DM/L^2`$. We fix $`\alpha =L^2/M15`$ since $`\omega _D`$ should be independent of the system size $`L`$ . We checked that the results are not affected by a variation of $`\alpha `$ in few times. The first studies of the TIP model with frozen Fermi sea was done by Imry with the aim to take into account the effect of finite fermionic density and then was also analyzed in . Recently a similar model was investigated for the case of Hubbard attraction in 3d .
To study the eigenstate properties of our model we diagonalize numerically the Hamiltonian (2) and rewrite the eigenfunctions in the original lattice basis. In this way we determine the two-particle probability distribution $`F(𝐧_1,n_2)`$ from which we extract the one particle probability $`f(𝐧_1)=_{𝐧_2}F(𝐧_1,n_2)`$ and the probability of inter-particle distance $`f_d(𝐫)=_{𝐧_2}F(𝐫+n_2,n_2)`$ with $`𝐫=n_1n_2`$. The binding energy of an eigenstate in (2) is $`\mathrm{\Delta }E=E2E_FE`$ since $`E_F0`$. For the ground state with energy $`E_g`$ the coupling energy is $`\mathrm{\Delta }=2E_FE_g`$. The typical examples of probability distributions are shown in Fig. 1. They clearly show that the ground state in the presence of interaction remains localized and the particles stay on distance $`R`$ from each other. However, there are states with negative binding energy ($`\mathrm{\Delta }E<0`$) which are delocalized by interaction and for which the particles move around the ring in agreement with discussion of model (1) at $`|U|V`$. We stress that this delocalization of coupled states ($`\mathrm{\Delta }E<0`$) takes place in the well localized one-particle phase. However, at very strong disorder this delocalization disappears (see top right case in Fig. 1).
To analyze the delocalization of states with negative binding energy $`\mathrm{\Delta }E`$ in a more quantitative way we determine the inverse participating ratio (IPR) $`\xi `$ for one-particle probability $`1/\xi =_𝐧f^2(𝐧)`$, where brackets mark the averaging over 100 disorder realisations. In this way $`\xi `$ gives the number of lattice sites occupied by one particle in an eigenstate. The dependence of $`\xi `$ on $`\mathrm{\Delta }E`$ and $`W`$ is shown in Fig. 2 for different lattice sizes $`L`$ in the presence of interaction. This figure shows that near the ground state the interaction creates states which are even more localized than in the absence of interaction ($`\xi `$ is significantly smaller than at $`U=0`$, see insert Fig. 2).
In fact for $`\mathrm{\Delta }<\mathrm{\Delta }E<\mathrm{\Delta }E_c<0`$ the IPR value even slightly drops with the increase of $`L`$. However for the states with binding energy $`\mathrm{\Delta }E_c<\mathrm{\Delta }E<0`$ the situation becomes different and $`\xi `$ grows significantly with $`L`$ while the change of IPR at $`U=0`$ with $`L`$ is rather weak (see shaded band in Fig. 2). The critical value of the binding energy $`\mathrm{\Delta }E_c`$ can be defined as such an energy at which $`\xi `$ remains independent of $`L`$. In this way $`\mathrm{\Delta }E_c`$ determines the mobility edge for coupled states so that at given $`U`$ and $`W`$ the TIP eigenstates are localized for $`\mathrm{\Delta }<\mathrm{\Delta }E<\mathrm{\Delta }E_c`$ while for $`\mathrm{\Delta }E_c<\mathrm{\Delta }E<0`$ the states becomes delocalized (see an example in Fig. 1). In agreement with this picture $`\xi `$ varies up to 30 times when $`\mathrm{\Delta }E`$ changes from $`\mathrm{\Delta }`$ up to $`0`$. This variation grows with $`L`$ and the interaction radius $`R`$ since the system becomes more close to the effective 3d Anderson model as it was discussed above. The qualitative change of the structure of the eigenstates leads also to a change in the level spacing statistics $`P(s)`$ (Fig. 3). Near the ground state the statistics is close to the Poisson distribution $`P_P(s)=\mathrm{exp}(s)`$ typical for the localized Anderson phase while for $`\mathrm{\Delta }E_c<\mathrm{\Delta }E<0`$ it approaches to the Wigner surmise $`P_W(s)=\pi s\mathrm{exp}(\pi s^2/4)/2`$ corresponding to the delocalized phase .
The variation of the delocalization border $`\mathrm{\Delta }E_c`$ for TIP coupled states with disorder strength and interaction is shown in Fig. 4. While the coupling energy $`\mathrm{\Delta }`$ grows with $`U`$ and $`W`$, the mobility edge $`\mathrm{\Delta }E_c<0`$, on the contrary, disappears at strong $`W`$. According to the data of Fig. 4 all states with binding energy $`\mathrm{\Delta }E<0`$ become localized for $`W>W_{c2}9.5V`$ $`(U=2V)`$ and $`W>W_{c2}8V`$ $`(U=V)`$. This shows that at weaker interaction a weaker disorder is required to have delocalized coupled states. As it follows from Fig. 4, at small disorder $`W`$ the delocalization border $`\mathrm{\Delta }E_c`$ becomes closer and closer to the ground state. This means that at weak disorder the delocalization will take place for excited states with low energy. For $`WW_{c2}`$ and $`UV`$ the diffusion rate of delocalized TIP ring can be estimated as $`D_2V(W_{c2}/W)^2`$ . Further studies are required to determine the dependence of $`W_{c2}`$ on $`W`$ at $`|U|V`$.
In conclusion, our results show that long range attractive interaction between two particles in 2d leads to the appearance of delocalized diffusive states near the Fermi level inside the well localized noninteracting phase. It would be interesting to understand what will be the consequences of this delocalization for real many-body fermionic problem with attractive interaction. It is possible that obtained results will be also relevant for electrons with Coulomb repulsion. Indeed, in this case at very weak disorder each electron oscillates near an equilibrium position and the two-body interaction can be considered as an effective harmonic attraction .
We thank O.P.Sushkov for stimulating discussions, and the IDRIS in Orsay and the CICT in Toulouse for access to their supercomputers.
Appendix
FIG. 1bis. Color 2d density plots for the data of Fig.1 with the same ordering of figures. Blue corresponds to the minimum of the probabilty distribution and red to the maximum. The first four figures are drawn in logarithmic scale while two figures at the bottom are in linear scale. Blue/Red color corresponds to: $`f=1.3\times 10^{10}/f=0.2`$ (top left), $`f=1.1\times 10^9/f=0.26`$ (middle left), $`f=1.5\times 10^6/f=0.014`$ (bottom left), $`f=1.14\times 10^{11}/f=0.073`$ (top right), $`f_d=3.8\times 10^7/f_d=0.1`$ (middle right), $`f_d=1.5\times 10^5/f_d=0.0032`$ (bottom right). |
warning/0002/cond-mat0002213.html | ar5iv | text | # Origin of the Quasiparticle Peaks of Spectral Functions in High 𝑇_𝑐 Cuprates
## Abstract
Based on the SU(2) slave-boson approach to the t-J Hamiltonian, we examine the cause of the sharp peaks(’quasiparticle’ peaks) in the observed spectral functions in high $`T_c`$ cuprates. The computed results reveal that the spectral weight of the sharp peaks increases with hole doping rate in agreement with observation. It is shown that the observed sharp peaks are attributed to the enhancement of spinon pairing(spin singlet pair formation) by the presence of holon pair bosons in the superconducting state.
Recently we reported a study of phase diagram involving holon pair condensation for high $`T_c`$ cuprates based on an improved approach of the SU(2) slave-boson theory over a previous study of the U(1) slave-boson theory that we recently made. In this approach, both the spinon and holon degrees of freedom are introduced into the Heisenberg exchange term in the t-J Hamiltonian, by considering the possibility of on-site charge fluctuations which arise as a result of site to site electron(and thus holon) hopping for the quantum systems of hole doped high $`T_c`$ cuprates. Unlike the SU(2) theory, the phase fluctuation effects of order parameters are not taken into account in the U(1) mean field approach. Thus it is of great interest to study how the phase fluctuations affect the observed spectral functions by applying the SU(2) theory with the above considerations. Currently there exists a lack of understanding the microscopic cause of the sharp peaks(’quasiparticle’ peaks) in the ARPES (angle resolved photoemission spectroscopy). In the present study, using the improved approach of the SU(2) slave-boson theory we evaluate one particle spectral functions for the normal and superconducting states and focus on the cause of the sharp quasiparticle peaks which appear in the superconducting state. In addition we examine the role of phase fluctuations of the spinon pairing order parameters on the spectral functions based on the SU(2) theory.
In the slave-boson representation, the electron annihilation operator of spin $`\sigma `$, $`c_\sigma `$ can be written as a composite of spinon and holon operators. That is, $`c_\sigma =b^{}f_\sigma `$ in the U(1) representation and $`c_\alpha =\frac{1}{\sqrt{2}}h^{}\psi _\alpha `$ in the SU(2) theory with $`\alpha =1,2`$, where $`f_\sigma `$($`b`$) is the spinon(holon) annihilation operator in the U(1) theory, and $`\psi _1=\left(\begin{array}{c}f_1\\ f_2^{}\end{array}\right)`$ $`\left(\psi _2=\left(\begin{array}{c}f_2\\ f_1^{}\end{array}\right)\right)`$ and $`h=\left(\begin{array}{c}b_1\\ b_2\end{array}\right)`$ are the doublets of spinon and holon annihilation operators respectively in the SU(2) theory.
Introducing Hubbard Stratonovich transformations for direct, exchange and pairing channels and a subsequent saddle point approximation, the t-J Hamiltonian is decomposed into the spinon sector, $`H^f`$ and the holon sector, $`H^b`$,
$`H^f={\displaystyle \frac{J(1\delta )^2}{2}}{\displaystyle \underset{<i,j>}{}}[\mathrm{\Delta }_{ij}^f(f_{1j}f_{2i}f_{2j}f_{1i})+c.c.]`$ (2)
$`{\displaystyle \frac{J(1\delta )^2}{4}}{\displaystyle \underset{<i,j>}{}}[\chi _{ij}(f_{\sigma i}^{}f_{\sigma j})+c.c.],`$
$`H^b={\displaystyle \frac{t}{2}}{\displaystyle \underset{<i,j>}{}}\left[\chi _{ij}(b_{1i}^{}b_{1j}b_{2j}^{}b_{2i})\mathrm{\Delta }_{ij}^f(b_{1j}^{}b_{2i}+b_{1i}^{}b_{2j})\right]c.c.`$ (4)
$`{\displaystyle \underset{<i,j>,\alpha ,\beta }{}}{\displaystyle \frac{J}{2}}|\mathrm{\Delta }_{ij}^f|^2[\mathrm{\Delta }_{ij;\alpha \beta }^b(b_{\alpha i}b_{\beta j})+c.c.]{\displaystyle \underset{i,\alpha }{}}\mu _ib_{\alpha i}^{}b_{\beta j},`$
where $`\chi _{ij}=<f_{\sigma j}^{}f_{\sigma i}+\frac{2t}{J(1\delta )^2}(b_{1j}^{}b_{1i}b_{2i}^{}b_{2j})>`$ is hopping order parameter, $`\mathrm{\Delta }_{ij}^f=<f_{1j}f_{2i}f_{2j}f_{1i}>`$, spinon pairing order parameter, $`\mathrm{\Delta }_{ij;\alpha \beta }^b=<b_{i\alpha }b_{\beta j}>`$, holon pairing order parameter, and $`\mu _i`$, the effective chemical potential. With the uniform hopping order parameter, $`\chi _{ij}=\chi `$, the d-wave spinon pairing order parameter, $`\mathrm{\Delta }_{ij}^f=\pm \mathrm{\Delta }_f`$ with the sign $`+()`$ for the nearest neighbor link parallel to $`\widehat{x}`$ ($`\widehat{y}`$) and the s-wave holon pairing order parameter, $`\mathrm{\Delta }_{ij;\alpha \beta }^b=\mathrm{\Delta }_b(\delta _{\alpha ,1}\delta _{\beta ,1}\delta _{\alpha ,2}\delta _{\beta ,2})`$, the quasiparticle energy for spinon is given by
$`E_k^f`$ $`=`$ $`\sqrt{(ϵ_k^f)^2+(\mathrm{\Delta }_f^{^{}})^2},`$ (5)
where the spinon single particle energy is given by,
$`ϵ_k^f`$ $`=`$ $`{\displaystyle \frac{J(1\delta )^2}{2}}\chi (\mathrm{cos}k_x+\mathrm{cos}k_y),`$ (6)
and the spinon pairing gap,
$`\mathrm{\Delta }_f^{^{}}`$ $`=`$ $`J(1\delta )^2\mathrm{\Delta }_f(\mathrm{cos}k_x\mathrm{cos}k_y).`$ (7)
The single particle(electron) propagator of interest is given by a convolution integral of spinon and holon propagators in the momentum space,
$`G_{\alpha \beta }(𝐤,\omega )`$ $`=`$ $`i{\displaystyle \frac{d𝐤^{^{}}d\omega ^{^{}}}{(2\pi )^3}G_{\alpha \beta }^f(𝐤+𝐤^{^{}},\omega +\omega ^{^{}})G^b(𝐤^{^{}},\omega ^{^{}})}`$ (9)
in the U(1) theory, and
$`G_{\alpha \beta }(𝐤,\omega )`$ $`=`$ $`{\displaystyle \frac{i}{2}}{\displaystyle }{\displaystyle \frac{d𝐤^{^{}}d\omega ^{^{}}}{(2\pi )^3}}[{\displaystyle \underset{l,m}{}}G_{\alpha \beta lm}^f(𝐤+𝐤^{^{}},\omega +\omega ^{^{}})\times `$ (11)
$`G_{ml}^b(𝐤^{^{}},\omega ^{^{}})]\text{in the SU(2) theory}.`$
Here the spinon Green’s function is $`G_{\alpha \beta }^f(𝐤,\omega )=i𝑑t_𝐱e^{i\omega ti𝐤𝐱}<T[f_\alpha (𝐱,t)f_\beta ^{}(0,0)]>`$ and the holon Green’s function, $`G^b(𝐤,\omega )=i𝑑t_𝐱e^{i\omega ti𝐤𝐱}<T[b(𝐱,t)b^{}(0,0)]>`$. They are the mean field Green’s functions for the U(1) Hamiltonian. The mean field Green’s functions for the SU(2) Hamiltonian are $`G_{\alpha \beta lm}^f(𝐤,\omega )=i𝑑t_xe^{i\omega ti𝐤𝐱}<T[\psi _{\alpha l}(𝐱,t)\psi _{\beta m}^{}(0,0)]>`$ and $`G_{lm}^b(𝐤,\omega )=i𝑑t_xe^{i\omega ti𝐤𝐱}<T[b_l(𝐱,t)b_m^{}(0,0)]>`$ respectively. The symbol $`<`$ $`>`$ refers to the finite temperature ensemble average of an observable quantity $`O`$, $`<O>\frac{1}{Z}\mathrm{tr}(e^{\beta H}O)`$.
The one electron removal spectral function, $`I(𝐤,\omega )`$ is obtained from ,
$`I(𝐤,\omega )={\displaystyle \frac{1}{\pi }}\mathrm{ImG}(𝐤,\omega +i0^+)f(\omega ),`$ (12)
where $`f(x)`$ is the Fermi distribution function. In the present study, we choose the Heisenberg coupling constant, $`J=0.2`$ $`t`$ and the hopping strength $`t=0.44`$ $`eV`$. Using the SU(2) theory , the predicted values of optimal hole doping rate $`\delta `$, pseudogap temperature $`T^{}`$ and bose condensation temperature $`T_c`$ are $`\delta =0.13`$, $`T^{}=0.029t`$($`148K`$) and $`T_c=0.021t`$ ($`107.2K`$) respectively. To compute the spectral function above we first evaluate the convolution integral of the holon and spinon Green’s functions $`G(𝐤,\omega )`$, based on the effective Hamiltonians of the spinon and holon sectors respectively. All of the computed results are based on the square lattice of $`100\times 100`$ in momentum space which is found to be sufficient for numerical convergence.
Fig. 1 displays the momentum dependence of the computed spectral functions at optimal doping for the ranges of momentum from $`𝐤=(0,0)`$ to $`𝐤=(\pi ,\pi )`$ and $`(\pi ,0)`$ to $`(\pi ,0.4\pi )`$, by using the SU(2) theory. Variation of the predicted spectral peak positions with momentum is in qualitative agreement with the ARPES data. The predicted spectral functions at temperatures below $`T_c`$ are characterized by the presence of sharp peaks(quasiparticle peaks) with shoulders(humps). Although not shown here, the U(1) slave-boson theory also predicts similar structures with higher spectral peaks and lower humps. In Fig.1 (a) the spectral peak position or the gap is seen to shift from a high value of $`150meV`$ to a low value of $`30meV`$ for the range of momentum from $`𝐤=(0,0)`$ to $`(\pi ,0)`$ at a temperature of $`T=0.004t`$($`20.4K`$ with $`t=0.44eV`$). Encouragingly the predicted value of the gap $`30meV`$ at $`𝐤=(\pi ,0)`$ is close to the value of spinon pairing gap $`31meV`$ obtained from Eq.(7). This indicates that the (leading energy) gap is now identified as the spinon pairing gap. The formation of the spin singlet pairs(spinon pairs) is attributed to the opening of the pseudogap at the transition temperature(pseudogap temperature $`T^{}`$). Although not completely shown in the figure, we find that the predicted gap size undergoes a continuous change(increase) as temperature is decreased from $`T^{}`$ to a superconducting temperature $`T_c`$ and even to temperatures below $`T_c`$. This indicates that the observed leading edge gaps of the spectral functions in the superconducting state should have the same origin as the ones observed in the pseudogap phase. Thus the remaining problem is to explain how the sharply enhanced quasiparticle peaks occur in the superconducting state while such distinctively sharp peaks are not manifest above $`T_c`$, as is shown in Fig. 1(b). However, the hump feature with no peak is observed in the normal state of $`Bi2212`$ system, while there exists no such measured reports regarding other high $`T_c`$ cuprates such as $`YBCO`$. Indeed, it will be of great importance to verify whether this hump feature is a universal nature of the normal states of the high $`T_c`$ cuprates. We argue that the robustness of the sharp quasiparticle peaks is attributed to the enhancement of spinon pairing(spin singlet pair formation) owing to the influence of holon pair bosons present in the superconducting state. In the following we will provide an explanation.
In order to see the role of the holon pair bosons on the appearance of the sharp peaks in the superconducting state, we first choose momenta only at the bottom of holon band in the convolution integral. They are $`𝐤=(0,0)`$ for $`b_1`$ bosons, and $`𝐤=(\pi ,\pi )`$ for $`b_2`$ bosons, which will allow for the formation of the holon pair bosons of the zero center of mass momentum $`𝐪=(0,0)`$ at temperatures below $`T_c`$. With such allowance of only the zero center of mass momentum for the holon pairs, sharp quasiparticle peaks with no shoulders are predicted to occur, as indicated by a dashed line in Fig. 2(a). On the other hand, with its removal such sharp peaks tend to be suppressed, as is shown by the solid line in the figure. However, with the inclusion of all possible values of momenta including the zero center of mass momentum, i.e., $`𝐤=(0,0)`$ and $`𝐤=(\pi ,\pi )`$, both the sharp peaks and the broad shoulders simultaneously appear, as is shown in Fig. 2 (b). Thus we argue that the sharp spectral peaks with leading edge gaps are caused by the enhancement of spinon pairing(spin singlet pair formation) by the presence of the holon pair bosons in the superconducting state. Although not displayed here, we find that physics is unchanged even with the U(1) mean field treatment. However some differences are observed in peak positions and heights. For completion the differences will be discussed shortly.
In Fig. 3 the doping dependence of spectral functions is displayed for underdoped, optimally doped and overdoped cases at a low temperature below $`T_c`$, $`T=0.004t`$ ($`T=20.4K`$ with $`t=0.44`$ $`eV`$) with a choice of $`𝐤=(0.8\pi ,0)`$(that is, close to $`𝐤=(\pi ,0)`$) for a qualitative comparison with observation. We find that the predicted spectral weight of the sharp quasiparticle peaks in the underdoped region decreases with decreasing hole concentration, showing agreement with the ARPES. The peaks in the overdoped region, e.g., $`\delta =0.20`$ are found to be higher than the ones in the underdoped region, in agreement with observation, revealing that the spectral weight increases with hole concentration.
Fig. 4 illustrates the role of the phase fluctuations of the spinon pairing order parameters on the spectral functions at an underdoping rate $`\delta =0.08`$ by observing differences between the U(1) and SU(2) theories. It is reminded that the phase fluctuation effects of the order parameters are incorporated only in the SU(2) theory, but not in the U(1) mean field theory. The SU(2) theory predicted a lower peak with a higher shoulder(hump) compared to the U(1) theory. Thus phase fluctuations cause to enhance the shoulder by lowering the spectral weight of the quasiparticle peaks as a compensation. In addition we find that the predicted spectral position or the spin gap size is shifted to a larger binding energy compared to the U(1) result.
In the present study we examined how sharp quasiparticle peaks and broad shoulders in the one electron removal spectral functions are predicted based on both the SU(2) and U(1) slave-boson approaches to the t-J Hamiltonian. We found that the gap(or spectral peak position) undergoes a continuous change, manifesting a gradual increase with decreasing temperature from the pseudogap temperature $`T^{}`$ to temperatures below the superconducting temperature $`T_c`$. This indicates that the origin of the observed leading edge gaps with the sharp(quasiparticle) peaks in the superconducting state is the same as the ones in the pseudogap phase. We showed that this gap is caused by the formation of spin singlet pairs(spinon pairs) that exist in both the pseudogap and superconducting phases. However the appearance of the distinctively sharp quasiparticle peaks below $`T_c`$ is attributed to the enhanced probability of spinon pairing(spin singlet pair formation) by the presence of the abundant holon pair bosons in the superconducting state, in contrast to the highly suppressed peaks predicted for the normal state. It is noted that broad shoulders with no peaks are observed in the normal state of $`Bi2212`$ system, while there exist no such reported measurements for other high $`T_c`$ cuprates such as $`YBCO`$. Thus it awaits further scrutiny to see whether such disappearance of the quasiparticle peaks is a common ’rule’ for all the high $`T_c`$ cuprates. It was shown that the spectral weight of the quasiparticle peaks increases with increasing hole concentration. This trend is in complete agreement with the ARPES measurements. Finally we note from the comparison of the SU(2) and U(1) theories that the effects of phase fluctuations in the spinon paring order parameters result in the enhancement of shoulders by lowering the height of the sharp peaks(’quasiparticle’ peaks).
One(SHSS) of us acknowledges the generous supports of Korea Ministry of Education(BSRI-98 and 99) and the Center for Molecular Science at Korea Advanced Institute of Science and Technology. |
warning/0002/cond-mat0002156.html | ar5iv | text | # I General Discussion
## I General Discussion
The purpose of this article is to analyze the dependence of the energy of an elementary excitation on the strength of the confinement potential, which exists in a planar semiconductor heterostructure. Due to the fascinating technological progress in the field of man-made structures, it has become possible to fabricate e.g. quantum wells of a widely varying shape. It is an interesting theoretical task to discuss the excitation spectrum of such semiconductor structures as function of the tunable parameters, such as well width, well height, etc.. Concerning the excitations of interest, we concentrate on particle-phonon systems, the particles being electrons or holes. The simplest example is that of a single polaron, that is an electron, coupled to a certain branch of lattice vibrations. Another example is that of a polaronic exciton, that is an electron-hole pair, coupled to phonons. Whereas the latter one is important to characterize optical properties, the former one has direct implications for the transport behavior of the materials of interest.
We assume that the interface(s)-induced confinement can be mimicked by a simple potential $`V_n(z_n)`$, $`n`$ being the particle number, $`z_n`$ the corresponding coordinate (the growth direction of the heterostructure will always be assumed as $`z`$-direction). Explicit forms of $`V_n(z_n)`$ may be rectangular wells, parabolas etc.. In addition, we suppose translation invariance to hold within the $`xy`$-plane. We remark that effects as surface roughness would destroy this property and could lead to the appearanc of new phenomena (e.g. localized states).
In the following equation, we define the class of models under discussion:
$`H:`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{n=1}{\overset{N}{}}}𝐩_nm_{n}^{}{}_{}{}^{1}𝐩_n+U(𝐫_1,..,𝐫_N)+`$ (2)
$`{\displaystyle \underset{𝐤}{}}\mathrm{}\omega _𝐤a_𝐤^+a_𝐤+{\displaystyle \frac{1}{\sqrt{V}}}{\displaystyle \underset{n=1}{\overset{N}{}}}{\displaystyle \underset{𝐤}{}}(g_{𝐤,𝐧}e^{i\mathrm{𝐤𝐫}_𝐧}a_𝐤+h.c.)`$
$`=`$ $`:H_{el}+H_{ph}+H_{int}.`$ (3)
The nomenclature is self explaining. The quantity $`U(𝐫_\mathrm{𝟏},𝐫_\mathrm{𝟐})`$ is to contain the confinement potentials as well as the particle interaction:
$$U(𝐫_1,..,𝐫_N):=\underset{n=1}{\overset{N}{}}V_n(z_n)+\frac{1}{2}\underset{\begin{array}{c}n,n^{}=1,\hfill \\ nn^{}\hfill \end{array}}{\overset{N}{}}V_{n,n^{}}(𝐫_n,𝐫_n^{}),$$
(4)
where $`V_{n,n^{}}`$ has to be calculated as potential energy of particle $`n`$, exposed to the electrostatic potential of particle $`n^{}`$. Because of the boundary conditions, $`V_{n,n^{}}(𝐫_n,𝐫_n^{})`$ itself is not translation invariant (see e.g. Ref. ). The particle-phonon coupling is of Fröhlich type. The most prominent example to be used here is that of a coupling to (LO)-phonons.
The model has two relevant limiting cases, which should be reproduced by any theory. Let the maximum of the well widths be $`L`$ and the minimum $`L^{}`$. If $`L^{}`$ tends to infinity, the confinement is irrelevant and the energy spectrum of $`H`$ is that of a three-dimensional well-material excitation. If $`L`$ tends to zero, the (finite height) well is irrelevant, leaving us with the spectrum of a three-dimensional barrier-material excitation. The behaviour for intermediate values of the well widths can qualitatively be discussed as follows. Varying $`L,L^{}`$ from sufficiently large values to smaller ones, the binding energy should increase due to the higher Coulomb correlation (for instance, the reader should recall that the energy of the two-dimensional hydrogen ground state is four times larger than that of a three-dimensional one). When $`L,L^{}`$ become smaller and smaller, the ground-state wave function will more and more effectively tunnel into the barrier material — the energy approaches the barrier limit.
Thus, we might expect a maximum of the binding energy to appear at intermediate values of $`L,L^{}`$. It was a controversely discussed question whether or not this maximum appears at relevant (that is not too small) values of $`L`$. The answer to this question might be not the same for different systems.
## II Polarons
The physics of polarons, confined to quantum wells, passed a few stages, and it is not possible to present here even a brief list of references. In particular, it was found that different phonon modes contribute to the polaron binding energy — confined bulk 2phonons inside the well, interface phonon mode and half-space bulk phonon mode in the barrier. We cite only papers concerning polarons confined to a finite rectangular potential (one layer heterostructure) where contribution of all phonon modes were taken into account. Anyway, there are problems to be addressed while dealing with multilayered heterostructures. Namely, we have to answer the following questions:
1) How to deal with multilayered heterostructures? The total number of phonon modes becomes too large to make numerical calculations even with modern computers. Besides, a multilayered heterostructure can generate a confining potential of rather complicated form, not only the rectangular one.
2) How to deal with mass- and dielectric mismatches in different layers? The polaron effective mass $`m(z)`$, the electron-phonon coupling constant $`\alpha (z)`$ and the phonon dispersion law do depend on a layer, that is, on the electron position. To glue solutions in different layers seems to be a cumbersome job.
To tackle these problems we suggest specific approximations, which will briefly be indicated here.
* A multilayered $`GaAs/Al_xGa_{1x}As`$ heterostructure is considered as an effective medium. Its mean parameters are to be defined by averaging over different layers according to the way they enter the Hamiltonian.
* The bulk phonon mode only inhabits an effective medium with mean characteristics.
We specify the electronic part of the Hamiltonian:
$`H_{el}=H_{el,}+H_{el,}={\displaystyle \frac{\stackrel{}{p}_{}^{\mathrm{\hspace{0.17em}2}}}{2m}}+{\displaystyle \frac{p_z^{\mathrm{\hspace{0.17em}2}}}{2m}}+V(z),`$ (5)
The mean electron band mass $`m`$ is defined by the equation
$`H_{el,}\psi _1=E_1\psi _1,{\displaystyle \frac{1}{m}}={\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}𝑑z{\displaystyle \frac{|\psi _1(z)|^2}{m(z)}},`$ (6)
where $`\psi _1(z),E_1`$ are the ground state wave function and the energy for the electron motion in $`z`$ direction. As $`\psi _1`$ and $`E_1`$ depend on $`m`$, we actually have the system of two equations (6) to calculate the mean band mass $`m`$.
The free LO-phonon Hamiltonian reads as follows:
$`H_{ph}=\mathrm{}\omega _{\text{LO}}{\displaystyle \underset{\stackrel{}{k}}{}}a_\stackrel{}{k}^{}a_\stackrel{}{k},\omega _{\text{LO}}={\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}dz\omega (z)|\psi _1(z)|^2.`$ (7)
As $`m`$ is found already, we define here the mean phonon frequency $`\omega _{\text{LO}}`$. Note that in this paper we are not interested in processes of emission, absorption or scattering of phonons. Instead we concentrate on virtual phonons in a cloud around an electron. Subsequently, the properties of the effective phonons do depend on the position of the electron as it follows from Eq. (7).
In the same way we define the effective electron-phonon interaction Hamiltonian in the standard Fröhlich form with the mean Fröhlich coupling constant $`\alpha `$:
$`\sqrt{\alpha }={\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}𝑑z|\psi _1(z)|^2{\displaystyle \frac{\omega (z)}{\omega _{\text{LO}}}}\left(\alpha (z)\sqrt{{\displaystyle \frac{m\omega _{\text{LO}}}{m(z)\omega (z)}}}\right)^{1/2}.`$ (8)
Evidently, this model belongs to the class defined in Eq. (2) As examples we studied 1) a one-layer heterostructure described by a rectangular confining potential
$`V(z)=\{\begin{array}{cc}0,\hfill & |z|L/2\hfill \\ V_0,\hfill & |z|>L/2\hfill \end{array}.`$ (11)
(the $`z`$-dependence of the masses and dielectric parameters is completely analogous) and 2) a multilayered heterostructure corresponding to the Rosen-Morse potential
$`V(z)`$ $`=`$ $`V_0\mathrm{tanh}^2\left({\displaystyle \frac{z}{L_{RM}}}\right).`$ (12)
We use perturbation theory in powers of $`\alpha `$ for both potentials, but in the first case we perform the summation over all virtual states while in the case of the Rosen-Morse potential the Green function (see ) can be used. To compare results for the Rosen-Morse and the rectangular potentials, an effective width $`L`$ of the Rosen-Morse potential has to be found. We define it as the width of a rectangular potential of the same height $`V_0`$ with the same ground-state energy. The dependence $`L(L_{RM})`$ can then be calculated. The parametrization for experimental data concerning $`GaAs/Al_xGa_{1x}As`$ heterostructure is based on the results reported in Ref. with some modifications, which are discussed in our paper. Actually we use the dependence of the parameters on the $`Al`$ mole fraction $`x`$ which depends in turn on the coordinate $`z`$ via the relation $`V(z)=600(1.155x+0.37x^2)\mathrm{meV}`$. The confining potential $`V(z)`$ being given, we know the dependence $`x(z)`$ and, subsequently, the values of the parameters $`\alpha ,m,\omega `$ at each point of the heterostructure which are averaged then following Eqs. (6), (7) and (8).
The polaron energy and effective mass are calculated for $`x=0.3`$. Peaks are found for the effective mass at some potential widths, while the energy demonstrates rather a smooth behavior between the correct 3D-limits as is seen in Fig. 1. As to the Rosen-Morse potential, the results are presented in Fig. 2 together with those for the rectangular potential of the corresponding effective width. One can see an excellent coincidence of the results obtained within the different techniques; clearly, this fact increases their reliability. A comparison is also made with the results of the papers, and the details are discussed in our paper.
## III Excitons
Sampling the previous literature, most work has been done on rectangular quantum wells with confinement potentials of type (11). The electron-hole potential can be calculated as indicated above and was given e.g. in Ref. .
To treat eigenvalue problems as the present one, we use tractable decompositions of the Hamiltonian to generate lower bounds for the ground-state energy. The basic idea is as follows: Assume we study the Hamiltonian $`H=p_z^2/2m+V_1(z)+V_2(z)`$ to find its ground-state energy $`E`$. Then we use the decomposition
$`H_1=x{\displaystyle \frac{p_z^2}{2m}}+V_1(z),H_2=(1x){\displaystyle \frac{p_z^2}{2m}}+V_2(z),0x1.`$ (13)
If $`E_1(x),E_2(x)`$ are the corresponding ground-state energies of $`H_1,H_2`$, then a lower bound for $`E`$ is: $`E\mathrm{max}_x(E_1(x)+E_2(x))`$.
Upper bounds are produced by variational methods: The trial wave-function used in our calculations had the form:
$`\mathrm{\Psi }(\stackrel{}{r}_{},z_1,z_2)=\mathrm{\Phi }_1(z_1)\mathrm{\Phi }_2(z_2)e^{a\sqrt{r_{}^2+b(z_1z_2)^2}},`$ (14)
where $`\mathrm{\Phi }_i(z_i)`$ are the ground-state eigenfunctions of the free electron ($`i=1`$) or the hole ($`i=2`$) in the confining potentials of the type (11). Evidently, the variational parameters $`a,b`$ can be used to fit 3D and 2D limiting cases. If the masses can be assumed as constant over the heterostructure, these methods can profitably be combined with functional-integral techniques. Fig. 3 shows our result for $`Al_{0.3}Ga_{0.7}As/GaAs/Al_{0.3}Ga_{0.7}As`$. in comparison with experimental and previous theoretical results . Clearly, the maximum appears at a relevant width.
A second class of confinement potentials is of parabolic type, that is
$$V_i(z)=\frac{m_iR_{\mathrm{}}^2\lambda _i^2}{2\mathrm{}^2}z_i^2,$$
(15)
where $`\lambda _i`$ denotes the dimensionless confinement strength, $`R_{\mathrm{}}`$ is the Rydberg unit, which was extracted for reasons of convenience. To study the confinement-induced effects on the spectrum as accurately as possible, we disregarded any parameter mismatch. The quantity of interest is the diagonal element of the reduced density operator, namely
$$P_\beta (𝐂):=<𝐂|tr_{Ph}e^{\beta H}|𝐂>.$$
(16)
In this formula $`𝐂`$ is an abbreviation for an arbitrary (but fixed) set of the position coordinates of the particles involved. The right-hand side of Eq. (16) can be represented as a functional integral
$$P_\beta (𝐂)=Z_{Ph}\delta ^6Re^{S[𝐑]}.$$
(17)
In Eq. (17) $`Z_{Ph}`$ is the free-phonon partition function, and $`S`$ reads as follows:
$`S[𝐑]`$ $`:=`$ $`{\displaystyle _0^\beta }𝑑\tau \left({\displaystyle \underset{n=1}{\overset{2}{}}}{\displaystyle \frac{m_n}{2}}\dot{𝐑}_n^2(\tau )+U(𝐑_1(\tau ),𝐑_2(\tau ))\right)`$ (19)
$`{\displaystyle \underset{n,n^{}=1}{\overset{2}{}}}{\displaystyle \underset{𝐤}{}}{\displaystyle \frac{g_{𝐤,n}g_{𝐤,n^{}}}{V}}{\displaystyle \underset{0}{\overset{\beta }{}}}{\displaystyle \underset{0}{\overset{\beta }{}}}𝑑\tau 𝑑\tau ^{}G(\tau \tau ^{})e^{i𝐤[𝐑_n(\tau )𝐑_n^{}(\tau ^{})]}.`$
Within the functional integral (17), $`\delta ^6R\mathrm{}.`$ is to indicate integration over all real, $`6`$-dimensional paths $`𝐑(\tau )`$, which start and end at the point $`𝐂`$. The kernel function $`G(\tau \tau ^{})`$ is defined as
$$G(\tau ):=\frac{e^{\mathrm{}\omega (\beta |\tau |)}+e^{\mathrm{}\omega |\tau |}}{2[e^{\beta \mathrm{}\omega }1]}.$$
(20)
It is well known that functional integrals of type (17) with an action (19) cannot be evaluated in analytical form. Starting from the exact expression, we use variational procedures as in Feynman’s famous paper on polarons to find upper bounds on the ground-state energy. The necessary input is a trial action, which is accessible to a numerical treatment.
The trial companions of the exact action (19) were combinations of oscillator trial actions for the center-of-mass and the $`z`$-coordinate and three-dimensional (two-dimensional) Coulomb potentials for the three-dimensional (two-dimensional in-plane) relative coordinates. The corresponding results (see Ref. ) can be found in the following figures and are denoted as quasi three-dimensional (Q3D) and quasi two-dimensional (Q2D or Q2Dalt) ansatz. In Fig. 4 we neglect any phonon influence to demonstrate the smooth interpolation of the limiting values $`1R_{\mathrm{}}`$ and $`4R_{\mathrm{}}`$ of the binding energy (actually we plotted there the ground-state energy with the continuum edge being subtracted, that is, the quantity $`E_B`$). Fig. 5 shows results for the general case; we present data for the ground-state energy as well as the continuum edge, which is the reference for the binding energy and has to be calculated separately.
The results reported have been obtained in collaboration with M. Dzero and J. Wüsthoff; we gratefully thank both of them. We are indebted to J. T. Devreese, V. Gladilin, H. Leschke, V. M. Fomin, F. M. Peeters, and E. P. Pokatilov for useful discussions and remarks. The support of Deutsche Forschungsgemeinschaft and the Germany-JINR Heisenberg-Landau program is acknowledged. |
warning/0002/math0002229.html | ar5iv | text | # Locally trivial quantum vector bundles and associated vector bundles
## 1 Coverings of modules
First we give a definition of modules over an algebra being equivalent to the usual one. For a vector space $`E`$ we denote by $`End(E)`$ the algebra of linear endomorphisms of $`E`$.
###### Definition 1
Let $`E`$ be a vector space, let $`B`$ be an algebra and $`\kappa :BEnd(E)`$ a linear map.
$`(E,B,\kappa )`$ is called left module if $`\kappa `$ satisfies $`\kappa (ab)=\kappa (a)\kappa (b),;a,bB`$.
$`(E,B,\kappa )`$ is called right module if $`\kappa `$ satisfies $`\kappa (ab)=\kappa (b)\kappa (a)`$.
For a linear subspace $`QE`$, $`(Q,B,\kappa )`$ is a submodule if $`\kappa (B)(Q)Q`$.
###### Definition 2
Let $`(E,B,\kappa )`$ be a left (right) module. $`(E,B,\kappa )`$ is called faithful if $`ker\kappa =0`$.
In analogy to the case of algebras one can define coverings of modules. They will be needed in the definition of quantum vector bundles.
###### Definition 3
Let $`(E,B,\kappa )`$ be a left (right) module and let $`(Q_i)_{iI}`$ be a finite family of left (right) submodules of $`E`$. $`(Q_i)_{iI}`$ is called covering of $`E`$ if $`_iQ_i=0`$.
Obviously, as in the case of algebras, for a given covering $`(Q_i)_{iI}`$ of a module $`(E,B,\kappa )`$ one obtains a family of vector spaces $`E_i:=E/Q_i`$ with corresponding projections $`q_i:EE_i`$. Since $`Q_i`$ are submodules, there exist linear maps $`\kappa _i:BEnd(E_i)`$ defined by
$$\kappa _i(a)q_i=q_i\kappa (a)aB$$
such that $`(E_i,B,\kappa _i)`$ are left or right modules respectively.
Remark: Since $`\kappa _i`$ are homomorphisms for left modules and antihomomorphism for right modules, $`ker\kappa _i`$ are ideals in $`B`$.
###### Definition 4
A covering $`(Q_i)_{iI}`$ of a left (right) module $`(E,B,\kappa )`$ is called nontrivial if $`ker\kappa _i0iI`$.
###### Proposition 1
Let $`(E,B,\kappa )`$ be a faithful left (right) module and let $`(Q_i)_{iI}`$ be a covering of $`E`$.
Then $`(ker\kappa _i)_{iI}`$ is a covering of $`B`$.
The covering $`(Q_i)_{iI}`$ is nontrivial if the covering $`(ker\kappa _i)_{iI}`$ is nontrivial.
Proof: We have to prove that $`_iker\kappa _i=0`$. It is easy to verify that
$$ker\kappa _i=\{aB|\kappa (a)(E)Q_i).$$
Now it is clear that
$$\underset{i}{}ker\kappa _i=\{aB|\kappa (a)(E)Q_iiI\}$$
and since $`_iQ_i=0`$ and $`(E,B,\kappa )`$ is faithful it follows $`ker\kappa _i=0`$. $`\mathrm{}`$
Let $`B_i:=B/ker\kappa _i`$. It is obvious that the modules $`(E_i,B_i,\stackrel{~}{\kappa _i})`$, where $`\stackrel{~}{\kappa _i}`$ is defined by
$$\stackrel{~}{\kappa }_i\pi _i=\kappa _i,$$
are faithful. Let
$`q_{ij}:E`$ $``$ $`E/(Q_i+Q_j):=E_{ij}`$
$`q_j^i:E_i`$ $``$ $`E_{ij}`$
be the canonical projections. Assume that $`(E,B,\kappa )`$ is a faithful left (right) module and $`(Q_i)_{iI}`$ is a covering of $`E`$. One has the vector space
$$E_c:=\{(e_i)_{iI}\underset{iI}{}E_i|q_j^i(e_i)=q_i^j(e_j)\}$$
(1)
and an injective homomorphism
$$𝒦:EE_c$$
by $`𝒦(e)=(q_i(e))_{iI}`$.
###### Proposition 2
Let $`(E,B,\kappa )`$ be a faithful left (right) module and let $`(Q_i)_{iI}`$ be a covering of $`E`$. Let $`B_c`$ be the covering completion of $`B`$ with respect to $`(ker\kappa _i)_{iI}`$. Then there exists a linear map $`\kappa _c:B_cEnd(E_c)`$ such that $`(E_c,B_c,\kappa _c)`$ is a faithful left (right) module satisfying
$$𝒦(\kappa (a)(e))=\kappa _c(K(a))(𝒦(e)),aB,eE,$$
where $`K:BB_c`$ is the injektive Homomorphism defined by $`K(a)=(\pi _i(a))_{iI}`$. Proof: Since $`(Q_i)_{iI}`$ is a family of submodules, there exist linear maps $`\kappa _{ij}:BEnd(E_{ij})`$ defined by
$$\kappa _{ij}(a)q_{ij}=q_{ij}\kappa (a),aB$$
such that $`(E_{ij},B,\kappa _{ij})`$ is a left (right) module. Let $`B_{ij}:=B/(ker\kappa _i+ker\kappa _j)`$. Now one can define the linear map $`\stackrel{~}{\kappa }_{ij}:B_{ij}End(E_{ij})`$ by the formula
$$\stackrel{~}{\kappa }_{ij}\pi _{ij}=\kappa _{ij},$$
thus $`(E_{ij},B_{ij},\stackrel{~}{\kappa }_{ij})`$ is a left (right) module. It is easy to verify that $`\stackrel{~}{\kappa }_{ij}`$ satisfies
$$q_j^i\stackrel{~}{\kappa }_i(a)=\stackrel{~}{\kappa }_{ij}(\pi _j^i(a))q_j^i,aB_i.$$
Recall that $`B_c:=\{(a_i)_{iI}_{iI}B_i|\pi _j^i(a_i)=\pi _i^j(a_j)\}`$. We define $`\kappa _c:B_cEnd(E_c)`$ as follows:
$$\kappa _c((a_i)_{iI})((e_i)_{iI})=(\stackrel{~}{\kappa }_i(a_i)(e_i))_{iI}.$$
(2)
One has to prove that the image of $`\kappa _c(a_c)`$ lies in $`E_c`$ for all $`a_cB_c`$:
$`q_j^i(\stackrel{~}{\kappa }_i(a_i)(e_i))`$ $`=`$ $`\stackrel{~}{\kappa }_{ij}(\pi _j^i(a_i))(q_j^i(e_i))`$
$`=`$ $`\stackrel{~}{\kappa }_{ij}(\pi _i^j(a_j))(q_i^j(e_j))`$
$`=`$ $`q_i^j(\stackrel{~}{\kappa }_j(a_j)(e_j)).`$
The other properties of $`\kappa _c`$ follow from the definition. $`\mathrm{}`$
###### Definition 5
Let $`(E,B,\kappa )`$ be a faithful left (right) module and let $`(Q_i)_{iI}`$ be a covering of $`E`$. The covering $`(Q_i)_{iI}`$ is called complete, if the family of ideals $`(ker\kappa _i)_{iI}`$ is a complete covering of $`B`$ and the injective linear map $`𝒦:EE_c`$ is a left (right) module isomorphism.
## 2 Locally trivial quantum vector bundles and associated vector bundles
On the algebraic level, the notion of vector bundle is in the classical case related to the notion of section of a vector bundle. Let $`M`$ be a manifold, let $`C(M)`$ be the algebra of continous functions over $`M`$ and let $`V`$ be a vector space. The corresponding trivial vector bundle has the form $`M\times V`$. It is known, that the set of sections of $`M\times V`$ is the set of all $`V`$-valued functions denoted by $`C(M)V`$. This classical background leads us to the following definition of a locally trivial vector bundle.
###### Definition 6
A locally trivial quantum vector bundle (QVB) is a tupel
$$\{(E,B,\kappa ),V,(\zeta _i,J_i)_{iI}\}$$
(3)
where $`(E,B,\kappa )`$ is a faithful left module, $`(J_i)_{iI}`$ is a complete covering of $`B`$, $`V`$ is a vector space and $`\zeta _i:EB_iV`$ are surjective left module homomorphisms with the properties
$`(ker\zeta _i)_{iI}`$ complete covering of E (4)
$`\zeta _i(\kappa (a)(e))`$ $`=`$ $`\pi _i(a)\zeta _i(e)aBeE`$ (5)
$`ker\zeta _i+ker\zeta _j`$ $`=`$ $`ker((\pi _j^iid)\zeta _i)=ker((\pi _i^jid)\zeta _j).`$ (6)
Remark: In this definition we have used the left module structure $`(B_iV),B,\kappa _i)`$, which is defined by
$$\kappa _i(a)(b_iv)=\pi _i(a)b_ivaB,b_iB_i,vV.$$
In the following we want to denote such a vector bundle by $`E`$.
By definition, for a locally trivial QVB $`E`$ the family of submodules $`(ker\zeta _i)_{iI}`$ is a complete covering of $`E`$, i.e. $`E`$ is isomorphic to its covering completion. Note that there are isomorphisms $`\stackrel{~}{\zeta }_i:E/ker\zeta _iB_iV`$ defined by
$$\stackrel{~}{\zeta }_i(e+ker\zeta _i)=\zeta _i(e).$$
Because of (6) there exist also isomorphisms $`\zeta _{ij}^i:E/(ker\zeta _i+ker\zeta _j)B_{ij}V`$ defined by
$$\zeta _{ij}^i(e+ker\zeta _i+ker\zeta _j)=\zeta _i(e)+\zeta _i(ker\zeta _j),$$
thus there are left $`B_{ij}`$-module isomorphisms $`\varphi _{ij_E}`$ defined by $`\varphi _{ij_E}:=\zeta _{ij}^i\zeta _{ij}^{j}{}_{}{}^{1}`$ such that the covering completion of $`E`$ has the form
$$E=\{(e_i)_{iI}\underset{iI}{}B_iV|(\pi _j^iid)(e_i)=\varphi _{ij_E}(\pi _i^jid)(e_j)\}.$$
###### Proposition 3
Let $`N`$ be a locally trivial topological vector bundle over a compact topological space $`M`$.
Then the set of continuous sections $`\mathrm{\Gamma }(N)=\{s:MN\}`$ is a locally trivial QVB.
Let $`B=C(M)`$ be the algebra of continuous functions over a compact topological space $`M`$. Let $`(J_i)_{iI}`$ be a finite covering of $`B`$ coming from a finite covering $`(U_i)_{iI}`$ of $`M`$ by closed sets with nonempty open interior. Let $`E`$ be a locally trivial QVB over $`B`$ corresponding to this covering.
Then $`E`$ is the space of sections of a locally trivial vector bundle $`N`$.
Proof: We want to give here only the idea of the proof. To prove the first assertion one defines the module homomorphisms $`\zeta _i`$ in terms of the trivializations $`\psi _i:NU_i\times V`$ by
$$\zeta _i(s)=\psi _is$$
(identifying $`\psi _i(s(x))=(x,v)`$ with $`v`$) and shows the conditions claimed for $`\zeta _i`$.
To prove the second assertion one construct in term of the given locally trivial QVB a locally trivial vectorbundle in the following way. Let $`X=_i\{i\}\times U_i\times V`$. One obtains a locally trivial vector bundle $`N`$ over $`M`$ by factorizing $`X`$ with respect to the following relation:
$$(i,x,v)(j,x^{},v^{})\text{if}x=x^{}\text{and}v=\varphi _{ij_E}(1v^{})(x).$$
One proves that the module of sections $`\mathrm{\Gamma }(N)`$ is isomorphic to $`E`$. $`\mathrm{}`$.
Assume that there is given an LC differential algebra $`\mathrm{\Gamma }(B)`$ with complete covering $`(J_{i_\mathrm{\Gamma }})_{iI}`$ such that $`pr_0(J_{i_\mathrm{\Gamma }})=J_i`$, i.e. the factor algebras $`\mathrm{\Gamma }(B)/J_{i_\mathrm{\Gamma }}=\mathrm{\Gamma }(B_i)`$ are differential calculi over the factor algebras $`B_i=B/J_i`$. One can construct a locally trivial QVB $`((E_\mathrm{\Gamma },\mathrm{\Gamma }(B),\kappa _\mathrm{\Gamma }),V,(\zeta _{i_\mathrm{\Gamma }},J_{i_\mathrm{\Gamma }})_{iI})`$ from $`E`$ in the following way. One extends $`\varphi _{ij_E}`$ to $`\mathrm{\Gamma }(B_{ij})V`$ by
$$\varphi _{ij_E}(\gamma v)=\gamma \varphi _{ij_E}(1v),\gamma \mathrm{\Gamma }(B_{ij}),vV.$$
In terms of this extended module isomorphism one defines $`E_\mathrm{\Gamma }`$ by
$$E_\mathrm{\Gamma }:=\{(e_{i_\mathrm{\Gamma }})_{iI}\underset{iI}{}\mathrm{\Gamma }(B_i)V|(\pi _{j_\mathrm{\Gamma }}^iid)(e_{i_\mathrm{\Gamma }})=\varphi _{ij_E}(\pi _{i_\mathrm{\Gamma }}^jid)(e_{i_\mathrm{\Gamma }})\},$$
(7)
where the homomorphisms $`\pi _{j_\mathrm{\Gamma }}^i:\mathrm{\Gamma }(B_i)\mathrm{\Gamma }(B_{ij}):=\mathrm{\Gamma }(B)/(J_{i_\mathrm{\Gamma }}+J_{j_\mathrm{\Gamma }})`$ are the canonical projections.
By defining $`\kappa _\mathrm{\Gamma }:\mathrm{\Gamma }(B)End(E_\mathrm{\Gamma })`$ as
$$\kappa _\mathrm{\Gamma }(\gamma )((e_{i_\mathrm{\Gamma }})_{iI}):=(\pi _{i_\mathrm{\Gamma }}(\gamma )e_{i_\mathrm{\Gamma }})_{iI}$$
($`\pi _{i_\mathrm{\Gamma }}:\mathrm{\Gamma }(B)\mathrm{\Gamma }(B_i)=\mathrm{\Gamma }(B)/J_{i_\mathrm{\Gamma }}`$ is the canonical projection.) and $`\zeta _{i_\mathrm{\Gamma }}:E_\mathrm{\Gamma }\mathrm{\Gamma }(B_i)V`$ as the i-th projection one obtains a locally trivial QVB $`((E_\mathrm{\Gamma },\mathrm{\Gamma }(B),\kappa _\mathrm{\Gamma }),V,(\zeta _{i_\mathrm{\Gamma }},J_{i_\mathrm{\Gamma }})_{iI})`$.
Now one can consider connections on such locally trivial QVB. We add to the usual definition of a connection in a “vector bundle” as a covariant derivative () a condition of compatibility with the covering.
###### Definition 7
Let $`(E_\mathrm{\Gamma },\mathrm{\Gamma }(B),\kappa _\mathrm{\Gamma }),V,(\zeta _{i_\mathrm{\Gamma }},J_i)_{iI})`$ be the locally trivial QVB just defined. A connection on $`E_\mathrm{\Gamma }`$ is a linear map $`:E_\mathrm{\Gamma }E_\mathrm{\Gamma }`$ satisfying
$`(\gamma e)`$ $`=`$ $`(d\gamma )e+(1)^n\gamma (e),\gamma \mathrm{\Gamma }^n(B),eE_\mathrm{\Gamma }`$ (8)
$`(ker\zeta _{i_\mathrm{\Gamma }})`$ $``$ $`ker\zeta _{i_\mathrm{\Gamma }},iI.`$ (9)
It is easy to see that from the property (9) follows that there exist connections $`_i`$ on $`\mathrm{\Gamma }(B_i)V`$ and $`_{ij}^i`$ on $`\mathrm{\Gamma }(B_{ij})V`$ such that
$`_i\zeta _{i_\mathrm{\Gamma }}`$ $`=`$ $`\zeta _{i_\mathrm{\Gamma }}`$
$`_{ij}^i(\pi _{j_\mathrm{\Gamma }}^iid)`$ $`=`$ $`(\pi _{j_\mathrm{\Gamma }}^iid)_i.`$
###### Proposition 4
Connections on $`E_\mathrm{\Gamma }`$ are in one to one correspondence with families of connections $`_i:\mathrm{\Gamma }(B_i)V\mathrm{\Gamma }(B_i)V`$ satisfying
$$_{ij}^i=\varphi _{ij_E}_{ij}^j\varphi _{ji_E}.$$
(10)
Proof: Let $``$ be a connection on $`E_\mathrm{\Gamma }`$. There exists a family of connections $`_i`$ on $`\mathrm{\Gamma }(B_i)V`$ such that
$$((e_{i_\mathrm{\Gamma }})_{iI})=(_i(e_{i_\mathrm{\Gamma }}))_{iI}.$$
Because of the identities
$`(\pi _{j_\mathrm{\Gamma }}^iid)(e_{i_\mathrm{\Gamma }})`$ $`=`$ $`\varphi _{ij_E}(\pi _{i_\mathrm{\Gamma }}^jid)(e_{j_\mathrm{\Gamma }})`$
$`(\pi _{j_\mathrm{\Gamma }}^iid)_i(e_{i_\mathrm{\Gamma }})`$ $`=`$ $`\varphi _{ij_E}(\pi _{i_\mathrm{\Gamma }}^jid)_j(e_{j_\mathrm{\Gamma }})`$
one obtains
$`_{ij}^i(\pi _{j_\mathrm{\Gamma }}^iid)(e_{i_\mathrm{\Gamma }})`$ $`=`$ $`\varphi _{ij_E}_{ij}^j(\pi _{i_\mathrm{\Gamma }}^jid)(e_{j_\mathrm{\Gamma }})`$
$`=`$ $`\varphi _{ij_E}_{ij}^j\varphi _{ji_E}(\pi _{j_\mathrm{\Gamma }}^iid)(e_{i_\mathrm{\Gamma }}).`$
Conversely, if there is given a family of connections $`_i`$ satisfying property (10) the image of the linear map $``$ defined by
$$((e_{i_\mathrm{\Gamma }})_{iI})=(_i(e_{i_\mathrm{\Gamma }}))_{iI}$$
lies in $`E_\mathrm{\Gamma }`$ and has the properties of a connection. $`\mathrm{}`$
Remark: Let the family $`(e_i)_{iI}`$ be a linear basis in $`V`$. Let $``$ be a connection on $`\mathrm{\Gamma }(B)H`$. Then there is a family of one forms $`(A_i^j)`$ such that $`(1e_i)=_jA_i^je_j`$.
###### Definition 8
Let $``$ be a connection on $`E_\mathrm{\Gamma }`$. We call the linear map $`^2`$ the curvature of $``$.
Note that $`^2(\gamma e)=\gamma ^2(e)`$.
In the sequel we will be interested in QVB associated to a QPFB. We define these as follows (see also ):
###### Definition 9
Let $`𝒫`$ be a locally trivial QPFB, let $`F`$ be a vector space and let $`\rho :FHF`$ be a left $`H`$ coaction. The vector bundle $`E(𝒫,F)`$ associated to $`𝒫`$ and $`\rho `$ is defined as
$$E(𝒫,F):=\{e𝒫F|(\mathrm{\Delta }_𝒫id)(e)=(id\rho )(e)\}.$$
(11)
Remark: This is also called co-tensor product of $`𝒫`$ and $`F`$.
###### Proposition 5
The associated vector bundle $`E(𝒫,F)`$ is a locally trivial QVB.
Proof: By formula (10) of . $`E(𝒫,F)`$ has the form
$`E(𝒫,F)=\{(e_i)_{iI}{\displaystyle \underset{iI}{}}B_iHF|(\pi _j^iidid)(e_i)`$ $`=`$ $`(\varphi _{ij}id)(\pi _i^jidid)(e_j);`$
$`(id\mathrm{\Delta }id)(e_i)`$ $`=`$ $`(idid\rho )(e_i)\}.`$
There are isomorphisms
$$id\epsilon id:E_i:=\{e_iB_iHF|id\mathrm{\Delta }id(e_i)=idid\rho (e_i)\}B_iF,$$
where $`(id\epsilon id)^1=id\rho `$, such that all $`e_iE_i`$ are of the form
$$\underset{k}{}a_kf_{k_{(1)}}f_{k_{(0)}}.$$
One easily verifies that $`((\stackrel{~}{E},B,\stackrel{~}{\kappa }),F,(\stackrel{~}{\zeta _i},J_i)_{iI})`$ defined by
$$\stackrel{~}{E}(𝒫,F):=\{(\underset{k}{}a_{k_i}^if_{k_i}^i\}\underset{iI}{}B_iF|\underset{k}{}\pi _j^i(a_{k_i}^i)f_{k_i}^i=\underset{k}{}\pi _i^j(a_{k_i}^j)\tau _{ji}(f_{k_{j(1)}}^j)f_{k_{j(0)}}^j\}$$
$`\kappa (a)((\stackrel{~}{e}_i)_{iI})`$ $`=`$ $`(\pi _i(a)\stackrel{~}{e}_i)_{iI},aB,(\stackrel{~}{e}_i)_{iI}\stackrel{~}{E}`$
$`\zeta _i`$ i-th projection
is a locally trivial QVB (For the definition of the transition functions $`\tau _{ij}`$ see ). In terms of the isomorphisms $`id\epsilon id`$ one can define a module isomorphism $`ϵ:E(𝒫,F)\stackrel{~}{E}(𝒫,F)`$ by
$$ϵ((e_i)_{iI}):=(id\epsilon id(e_i))_{iI}.$$
This isomorphism exists due to the glueing properties. This is easy to see: An element $`(e_i)_{iI}E(𝒫,F)`$ is of the form $`(e_i)_{iI}=(_{k_i}a_{k_i}^if_{k_{i(1)}}^if_{k_{i(0)}}^i)_{iI}`$ satisfying
$$\underset{k_i}{}\pi _j^i(a_{k_i}^i)f_{k_{i(1)}}^if_{k_{i(0)}}^i=\underset{k_j}{}\underset{k_j}{}\pi _i^j(a_{k_j}^j)\tau _{ji}(f_{k_{j(2}}^j)f_{k_{j(1)}}^jf_{k_{j(0)}}^j.$$
(see formula (10) of )
Applying $`id\epsilon id`$ to both sides of this equation, one obtains for the element $`(id\epsilon id(e_i))_{iI}`$ the property
$$\underset{k_i}{}\pi _j^i(a_{k_i}^i)f_{k_i}^i=\underset{k_j}{}\pi _i^j(a_{k_j}^j)\tau _{ji}(f_{k_{j(1)}}^j)f_{k_{j(0)}}^j,$$
i.e. the image of $`ϵ`$ lies in $`\stackrel{~}{E}`$. The inverse of $`ϵ`$ is the map $`(\stackrel{~}{e}_i)_{iI}(id\rho (\stackrel{~}{e}_i))_{iI}`$, i.e $`ϵ`$ is an isomorphism. We obtain the locally trivial QVB $`((E(𝒫,F),B,\kappa ),F,(\zeta _i,J_i)_{iI})`$ with
$`\kappa `$ $`=`$ $`\stackrel{~}{\kappa }ϵ`$
$`\zeta _i`$ $`=`$ $`\stackrel{~}{\zeta }_iϵ.`$
$`\mathrm{}`$
As in the general case of locally trivial vector bundles one can construct a locally trivial vector bundle $`((E_\mathrm{\Gamma }(𝒫,F),\mathrm{\Gamma }_m(B),\kappa _\mathrm{\Gamma }),(\zeta _{i_\mathrm{\Gamma }},ker\pi _{i_{\mathrm{\Gamma }_m}})_{iI})`$ from the associated vector bundle $`E(𝒫,F)`$. The LC differential algebra $`\mathrm{\Gamma }_m(B)`$ is the maximal embeddable LC-differential algebra induced from the differential structure $`\mathrm{\Gamma }_c(𝒫)`$ on the locally trivial QPFB $`𝒫`$ (see and ). Let $`\varphi _{ij_E}:B_{ij}FB_{ij}F`$ be defined by
$$\varphi _{ij_E}(af):=a\tau _{ji}(f_{(1)})f_{(0)}.$$
By definition,
$$E_\mathrm{\Gamma }(𝒫,F)=((e_{i_\mathrm{\Gamma }})_{iI}\underset{iI}{}\mathrm{\Gamma }(B_i)F|(\pi _{j_{\mathrm{\Gamma }_m}}^iid)(e_{i_\mathrm{\Gamma }})=\varphi _{ij_E}(\pi _{i_{\mathrm{\Gamma }_m}}^jid)(e_{j_\mathrm{\Gamma }})\}.$$
(12)
###### Proposition 6
Let
$$horE(𝒫,F):=\{\gamma _Ehor\mathrm{\Gamma }_c(𝒫)F|(\mathrm{\Delta }_{𝒫_\mathrm{\Gamma }}id)(\gamma _E)=(id\rho )(\gamma _E)\}.$$
$`horE(𝒫,F)`$ and $`E_\mathrm{\Gamma }(𝒫,F)`$ are isomorphic as left $`\mathrm{\Gamma }_m(B)`$-modules.
Remark: $`horE(𝒫,F)`$ is in the classical situation the space of horizontal forms “of type $`\rho `$”.
Proof: By definition of $`hor\mathrm{\Gamma }_c(𝒫)`$,
$`horE(𝒫,F)`$ $`=`$ $`\{(\gamma _{i_E})_{iI}{\displaystyle \underset{iI}{}}\mathrm{\Gamma }(B_i)\widehat{}HF|`$
$`((\pi _j^iid)_\mathrm{\Gamma }id)(\gamma _{i_E})=(\varphi _{ij_\mathrm{\Gamma }}id)((\pi _i^jid)_\mathrm{\Gamma }id)(\gamma _{j_E});`$
$`(id\mathrm{\Delta }id)(\gamma _{i_E})=(idid\rho )(\gamma _{i_E})\}.`$
The last condition means that the i-th components have the form
$$\gamma _{i_E}=\underset{k_i}{}\gamma _{k_i}^if_{k_{i(1)}}^if_{k_{i(0)}}^i.$$
Now one defines again an isomorphism $`ϵ_\mathrm{\Gamma }:horE(𝒫,F)E_\mathrm{\Gamma }(𝒫,F)`$ by
$$ϵ_\mathrm{\Gamma }((\gamma _{i_E})_{iI})=(id\epsilon id(\gamma _{i_E}))_{iI}.$$
This isomorphism exists, if one can show that from the gluing conditions in $`horE(𝒫,F)`$ the gluing conditions in $`E_\mathrm{\Gamma }(𝒫,F)`$ follow. To this end apply $`P_{inv}id`$ (formula (66) of ) to
$$\underset{k_i}{}(\pi _j^iid)_\mathrm{\Gamma }(\gamma _{k_i}^if_{k_{i(1)}}^i)f_{k_{i(0)}}^i=\underset{k_j}{}\varphi _{ij_\mathrm{\Gamma }}(\pi _i^jid)_\mathrm{\Gamma }(\gamma _{k_j}^jf_{k_{j(1)}}^j)f_{k_{j(0)}}^j.$$
By the definition $`\pi _{j_{\mathrm{\Gamma }_g}}^i(\gamma )=(\pi _j^iid)_\mathrm{\Gamma }(\gamma \widehat{}1)`$ one obtains the gluing condition in $`E_\mathrm{\Gamma }(𝒫,F)`$, i.e.
$$\underset{k_i}{}\pi _{j_{\mathrm{\Gamma }_g}}^i(\gamma _{k_i}^i)f_{k_i}^i=\underset{k_j}{}\varphi _{ij_{E_\mathrm{\Gamma }}}(\pi _{i_{\mathrm{\Gamma }_g}}^j(\gamma _{k_j}^j)f_{k_j}^j),$$
which means that the image of $`ϵ_\mathrm{\Gamma }`$ lies in $`E_\mathrm{\Gamma }(𝒫,F)`$. Conversely, the inverse of $`ϵ_\mathrm{\Gamma }`$ is obviously defined by
$$ϵ_\mathrm{\Gamma }^1((e_{i_\mathrm{\Gamma }})_{iI})=(id\rho (e_{i_\mathrm{\Gamma }}))_{iI}.$$
$`\mathrm{}`$
Now we are interested in connections on such associated vector bundles. An important class of connections are the connections induced from left left covariant derivations on the locally trivial QPFB.
###### Proposition 7
Every left covariant derivation on the locally trivial QPFB $`𝒫`$ determines uniquely a connection on $`E_\mathrm{\Gamma }(𝒫,F)`$.
Proof: One defines a linear map $`:E_\mathrm{\Gamma }E_\mathrm{\Gamma }`$ by
$$:=ϵ_\mathrm{\Gamma }(D_lid)ϵ_\mathrm{\Gamma }^1,$$
which is easily seen to be a connection on $`E_\mathrm{\Gamma }(𝒫,F)`$. $`\mathrm{}`$
Remark: Because of the bijection between left and right covariant derivations also right covariant derivations on $`𝒫`$ determine connections on $`E_\mathrm{\Gamma }(𝒫,F)`$.
The curvature of such a connection is related to the curvature on the locally trivial QPFB $`𝒫`$ by
$$^2=ϵ_\mathrm{\Gamma }(D_l^2id)ϵ_\mathrm{\Gamma }^1.$$
The corresponding connections $`_i:\mathrm{\Gamma }(B_i)F\mathrm{\Gamma }(B_i)F`$ have the form
$$_i(\gamma f)=d\gamma f(1)^n\gamma A_{l_i}(f_1)f_0,\gamma \mathrm{\Gamma }^n(B_i),fF.$$
The curvatures of these connections are
$$_i^2(\gamma f)=\gamma F_{l_i}(f_1)f_0.$$ |
warning/0002/astro-ph0002092.html | ar5iv | text | # Fragmentation Instability of Molecular Clouds: Numerical Simulationsin press, ApJ 532, April 1, 2000
## 1 Introduction
The interstellar magnetic field plays an important role in the dynamics of molecular clouds and the collapse of dense cloud cores into protostars. Magnetic pressure and tension combine with thermal and turbulent kinetic pressure to resist gravitational collapse. The role of the magnetic field is often simply characterized by a critical mass $`M_{\mathrm{crit}}`$ which depends on the magnetic flux threading the cloud (Mestel & Spitzer, 1956; Mouschovias & Spitzer, 1976; Tomisaka, Ikeuchi, & Nakamura, 1988; McKee et al., 1993; Mestel, 1999). Clouds with $`M>M_{\mathrm{crit}}`$ are termed supercritical and collapse on a dynamical time-scale: the magnetic field can have a moderate effect on the morphology of collapse and can slow collapse in the cloud envelope (e.g. Black & Scott, 1982), but cannot significantly slow collapse in the core. Clouds with $`M<M_{\mathrm{crit}}`$ are termed subcritical, and evolve on a longer time-scale as magnetic support is lost due to ambipolar diffusion. The magnetic field is redistributed within the cloud so that the inner parts become supercritical. The cloud is then differentiated into a dynamically collapsing core with a magnetically supported envelope (Ciolek & Mouschovias, 1993; Fiedler & Mouschovias, 1993; Basu & Mouschovias, 1994; Ciolek & Mouschovias, 1994; Safier, McKee, & Stahler, 1997; Ciolek & Königl, 1998). Much progress has been made in following this type of evolution through 6 or more orders of magnitude of increase in central density, including the effects of rotation as well as detailed chemistry and grain physics. The outcomes of these calculations include detailed density, bulk velocity, ion-neutral drift velocity, magnetic field, and grain and ion abundance profiles in axisymmetric clouds, as well as an appreciation of the timescales, rate of magnetic flux loss, and role of magnetic braking in this mode of isolated star formation.
This theoretical picture can be observationally tested (see recent reviews by Evans, 1999; Myers, Evans, & Ohashi, 1999). So far, the results are ambiguous. Flow toward an isolated infrared source in the Bok globule B335 is well described by an inside-out collapse model (Zhou et al., 1990, 1993; Zhou, 1995). The density distribution and measured magnetic fieldstrength in the cloud B1 have been fit by a model with a subcritical envelope and a core which has evolved to a supercritical state by ambipolar drift (Crutcher et al., 1994). On the other hand, in some respects the existing models appear to be incomplete. Clumps and cores are not axisymmetric; Myers et al. (1991) surveyed 16 dense cores a few tenths of a parsec in size and pointed out that at least 6 of them are likely to be prolate. Ryden (1996) made a statistical argument, based on a larger sample, that clumps and globules are more likely prolate or triaxial than oblate. Ward-Thompson, Motte, & André (1999) showed that asphericity appears also at smaller scales. Collapse guided by a magnetic field could produce oblate clouds, but not prolate ones. It also appears unlikely that rotation accounts for the flattening; this has been shown quantitatively in the case of L1527 (Ohashi et al., 1997). Thus, the shapes are unexplained. Moreover, in some cases in which infall has been measured directly, it is more spatially extended, with faster velocities in the outer parts, than expected from the standard models of inside-out collapse or gravitational motion driven by ambipolar drift. This has been shown in the case of L1544 by Tafalla et al. (1998), and for 6 other starless cloud cores by Gregersen (1998). (Interestingly, a model of L1544 has been recently constructed by Ciolek & Basu (2000) to match the observations of Tafalla et al. (1998). The model requires a mass-to-flux ratio more nearly critical than previously published models by the same authors, and appears to display the same intermediate collapse discussed in this paper, although the authors do not call it out as such.) Finally, there are observations which relate to the timescale for collapse of molecular clouds into protostars. Lee & Myers (1999) find collapse timescales of $``$ 0.3-1.6 Myr from the ratio of the numbers of starless cores to cores with embedded young stellar objects. They state that this requires collapse 2-44 times faster than ambipolar drift models. These observations suggest that another ingredient may be required to explain the collapse of molecular clouds: the decay of turbulent support (Myers & Lazarian, 1998), or, as we explore in this paper, a magnetogravitational instability mediated by ambipolar drift.
The magnetic field likely also plays an important role in other aspects of collapse, which are still incompletely understood. The mechanism by which molecular clouds fragment, and the masses and morphology of those fragments, is clearly of importance to the stellar initial mass function and the origin of binary systems. The magnetic field can exert strong forces on many scales, affecting fragmentation (e.g. Boss, 1997, 1999). The field may also be a source of kinetic energy in the cloud, if the free energy of an ordered field can be released as turbulent motions. These issues have not been studied thoroughly in the detailed axisymmetric models of gravitational contraction with ambipolar drift referenced above because they would require calculations with no constraints on spatial symmetry.
In this paper, we address some of these issues by studying a simple problem: the evolution of small perturbations to an initially uniform, magnetically subcritical sheet of weakly ionized gas with a uniform magnetic field perpendicular to its plane. The perturbations evolve under the influence of magnetic tension, self gravity, thermal pressure, and ambipolar drift. Typically the sheet breaks up into a small number of fragments of elongated shape which are collapsing, losing magnetic flux through ambipolar drift, and interacting gravitationally with one another. The magnetic fields associated with these asymmetric clumps generate local vorticity (the net angular momentum of the sheet is identically zero). The characteristic vorticity structure is a vortex pair which flanks each clump and is associated with strong streaming motions along it. The clump masses are typically of order 1-10 $`M_{}`$, and scale with temperature like the Jeans mass, but are typically larger because of magnetic support. The main features of collapse in this geometry were predicted by the linear stability analysis of Zweibel (1998; hereafter Z98); there is also some overlap with the earlier stability analysis in 3D by Langer (1978). Here we follow the evolution into the nonlinear regime and follow the growth of density fluctuations from .01 to up to 10 times the mean surface density. Collapse occurs on an intermediate time-scale, slower than the dynamical or free-fall time-scale, but faster than the ambipolar diffusion time-scale. As noted above, this intermediate collapse rate may be observed in some clouds (e.g. Lee & Myers, 1999; Gregersen, 1998) . As no restrictions are placed on the morphology of the clumps, we can begin to explore the nature of flux loss and collapse in more complicated geometries than isolated axisymmetric clouds. One outcome suggested by linear theory which we have not resolved is whether stored magnetic energy is converted to turbulence, as there is no stored magnetic energy in the system initially.
The unperturbed initial geometry, governing equations, and linear theory are described in §2. Section 3 contains a description of the numerical method and main results: the collapse rate is discussed in §3.1, the relationship between magnetic field B and density $`\rho `$ or surface density $`\sigma `$ in §3.2, size of fragments in §3.3, velocity structure in §3.4, and distribution of energy in §3.5. The validity of the approximations is discussed in §4, and the summary and conclusions are in §5.
## 2 Governing Equations and Linear Theory
We simulate a flat slab of cold gas with the magnetic field initially perpendicular to the slab (the $`\widehat{z}`$ direction). Z98 discusses the linear theory for this model when the cloud temperature T equals 0. In this section we review the model in the more general case T$`>`$0. We recall the results of the linear theory and discuss the consequences of adding thermal pressure.
A flat slab-like model has observational and theoretical motivation: molecular clouds commonly have sheet or filament-like structure (although detailed, high-resolution information on the field orientation in such structures is not yet available for many objects). In a fairly quiescent environment, a roughly spherical molecular cloud with a large-scale, dynamically significant, ordered magnetic field will relax into a pancake or slab as matter drains down the field lines. Magnetic forces will allow comparatively little contraction perpendicular to the field direction, resulting in a slab with a predominantly perpendicular field. Such slabs could also be formed by shock waves propagating parallel to the local magnetic field.
We use a simple initial state with small ($`<`$1%) perturbations. The boundary conditions are periodic in the horizontal ($`\widehat{x}`$ and $`\widehat{y}`$) directions and the initial surface density $`\sigma _0`$ and magnetic field $`B_{z0}`$ are uniform. The unperturbed initial state was chosen for computational simplicity, as self-consistent finite disk and slab-like equilibrium states cannot generally be described by simple expressions in closed form. (e.g. Parker, 1974; Mouschovias, 1976; Baureis et al., 1989; Mestel & Ray, 1985). We note that our initial state is not technically one of static equilibrium, but rather a version of the commonly used Jeans swindle, in which the unperturbed gravitational potential is discarded (see e.g. discussion in Binney & Tremaine (1987)).
### 2.1 Governing Equations
We begin with the equations of ideal magnetohydrodynamics for two inviscid, non-resistive, interacting, magnetized fluids, one charged and the other neutral. Since sources and sinks are expected to dominate advection in the ion continuity equation, we treat directly only the continuity equation for neutrals, and parameterize the ion behavior through equation (2-10). The governing MHD equations are thus the equation of continuity for neutral particles,
$$\frac{\rho _n}{t}+\mathbf{}\left(\rho _n𝐯\right)=0,$$
(2-1a)
equations of motion for the two species,
$`\rho _n{\displaystyle \frac{D_n𝐯_n}{D_nt}}+\mathbf{}P_n+\rho _n\mathbf{}\mathrm{\Phi }_G`$ (2-1b)
$`+\rho _n\nu _{ni}(𝐯_n𝐯_i)`$ $`=`$ $`0,`$
$`\rho _i{\displaystyle \frac{D_i𝐯_i}{D_it}}+\mathbf{}P_i+\rho _i\mathbf{}\mathrm{\Phi }_G`$ (2-1c)
$`\rho _i\nu _{in}(𝐯_n𝐯_i)`$ $`=`$ $`{\displaystyle \frac{(\mathbf{}\times 𝐁)\times 𝐁}{4\pi }},`$
the induction equation
$$\frac{𝐁}{t}=\mathbf{}\times \left(𝐯_i\times 𝐁\right),$$
(2-1d)
and Poisson’s equation
$$^2\mathrm{\Phi }_G=4\pi G\rho .$$
(2-1e)
Subscripts $`i`$ and $`n`$ denote ions and neutral particles respectively. $`D_\alpha /D_\alpha t`$ is the convective derivative for species $`\alpha `$. The ion-neutral collision frequency is $`\nu _{in}=\rho _n\sigma v/(m_i+m_n)`$, and $`\nu _{ni}`$ is the neutral-ion collision frequency $`(\rho _i\nu _{in}=\rho _n\nu _{ni})`$. In the context of collision rates only, the symbol $`\sigma `$ represents the cross-section for elastic collisions; elsewhere it represents surface density. The gravitational potential is $`\mathrm{\Phi }_G`$, and $`𝐯`$, $`\rho `$, $`P`$, and $`𝐁`$ are the velocity, density, pressure, and magnetic field, respectively. We work on large scales and at low temporal frequencies for which the ions and electrons are coupled.
We assume that the ionization fraction in the cloud is low. For dense molecular gas which is ionized by cosmic rays and recombines on grains, $`n_i/n_nKn_n^{1/2}`$, where $`K1.1\times 10^5`$ (McKee et al., 1993) Departures from and generalizations of this ionization law are discussed below, see eq. \[2-10\]. Therefore, $`\rho _n\rho `$ and $`𝐯_n𝐯`$. If the neutral-ion collision time is much less than a dynamical time, the ambipolar drift velocity $`𝐯_D`$ can be written in the standard form (Shu, 1983):
$$𝐯_D𝐯_i𝐯_n=\frac{(\mathbf{}\times 𝐁)\times 𝐁}{4\pi \nu _{in}\rho _i}.$$
(2-2)
The flat geometry allows significant simplification of the equations by taking a vertical integral in the limit of infinitesimal vertical thickness. For example, the magnetic force simplifies to:
$`\underset{ϵ0}{lim}{\displaystyle \underset{ϵ}{\overset{+ϵ}{}}}𝑑z{\displaystyle \frac{(\mathbf{}\times 𝐁)\times 𝐁}{4\pi }}`$ $`=`$ (2-3)
$`\underset{ϵ0}{lim}{\displaystyle \frac{B_z}{4\pi }}\left[B_x\widehat{x}+B_y\widehat{y}\right]_ϵ^{+ϵ}`$ $`=`$ $`{\displaystyle \frac{B_z𝐁_h}{2\pi }}.`$ (2-4)
In the limit of an infinitesimally thin disc or slab, the vertical component of the magnetic field $`B_z`$ is continuous with respect to the plane of the slab (the $`\widehat{z}`$ direction), and the horizontal component $`𝐁_h`$ is antisymmetric (reverses sign) with respect to the plane of the slab.
In addition, we assume the vertical component of the velocity is negligible compared to the horizontal components $`(v_zv_x,v_y)`$. Lovelace & Zweibel (1997) found that thin disks are generally stable to warping modes, so we expect predominantly horizontal motion.
In the limit of zero gas density and thermal pressure outside the slab, the external magnetic field relaxes instantaneously to an equilibrium state, shown in Z98 to be a current-free or potential field state. We therefore assume that the magnetic field at $`z0`$ is a potential field. This allows us to calculate only the vertical part of the magnetic field in the disc, rather than all three components in a three-dimensional domain, a tremendous simplification. Limitations of this approximation, and corrections to it, are discussed in more detail in §4.1 and the Appendix.
The resulting system of 2-dimensional equations are as follows: The equation of continuity
$$\frac{\sigma }{t}+\mathbf{}_h(\sigma 𝐯_h)=0,$$
(2-5a)
of motion
$`\sigma {\displaystyle \frac{D𝐯_h}{Dt}}+\mathbf{}_hP+\sigma \mathbf{}_h\mathrm{\Phi }_G`$ $`=`$ $`{\displaystyle \frac{𝐁_hB_z}{2\pi }},`$ (2-5b)
$`v_z`$ $`=`$ $`0,`$ (2-5c)
the definition of the ambipolar drift velocity
$$𝐯_{Dh}=\frac{B_z𝐁_h}{2\pi \nu _{in}\sigma _i},$$
(2-5d)
the induction equation
$$\frac{B_z}{t}=\mathbf{}_h\left[(𝐯_h+𝐯_{Dh})B_z\right],$$
(2-5e)
the potential field equations
$`B_z`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Phi }_M}{z}},`$ (2-5f)
$`𝐁_h`$ $`=`$ $`_h\mathrm{\Phi }_M,`$ (2-5g)
and Poisson’s equation
$$\frac{\mathrm{\Phi }_G}{z}=2\pi G\sigma .$$
(2-5h)
The gravitational and magnetic potentials are $`\mathrm{\Phi }_G`$ and $`\mathrm{\Phi }_M`$ respectively, $`\sigma `$ is the surface density, and $`h`$ denotes horizontal components (e.g., $`𝐯_{Dh}`$ is the horizontal drift velocity). We assume an isothermal equation of state $`P=a^2\sigma `$.
Equations for an axisymmetric disk of small but finite half thickness $`Z`$ were derived by Ciolek & Mouschovias (1993). Their equations contain correction terms of order $`Z/R`$, where $`R`$ is the distance from the axis of symmetry; these terms include magnetic pressure, which provides a restoring force which we neglect, and corrections to the normal direction, which is tilted slightly from the vertical because $`Z`$ depends on $`R`$. These terms go smoothly to zero in the limit $`Z/R0`$. For canonical values of physical parameters in dense molecular clouds, we find $`Z/R<1/10`$. Strictly speaking, we should retain the magnetic pressure gradient, as Ciolek & Mouschovias (1993) do, because it is of the same order as, although generally less than, the thermal pressure gradient, and also provides a restoring force. However, magnetic tension clearly dominates magnetic pressure at the long wavelengths of greatest interest here, while as the wavenumber increases the rate of ambipolar drift increases as well, so that the magnetic pressure force at short wavelengths is less than it would be if the magnetic field were frozen in. Moreover, we know that there is a 3D version of the instability which is driven by magnetic pressure alone. We study the instability driven by tension, but the existence and nature of the instability should be the same whether it is driven by tension or pressure. For all these reasons, we think that the neglect of magnetic pressure is not a major source of error.
### 2.2 Nondimensionalization
All quantities in the problem are scaled by a self-consistent set of characteristic quantities. Given an initial vertical magnetic field $`B_{z0}`$, we choose as the characteristic surface density that which is marginally stable to collapse in the zero temperature limit (Nakano & Nakamura, 1978),
$$\sigma _{c0}\frac{B_{z0}}{2\pi G^{1/2}}.$$
(2-6)
In a 3-dimensional model of a cold magnetized molecular cloud, one logical choice would be to use the Alfvén velocity as a characteristic velocity. The 2-dimensional geometry precludes this approach, because the quantity $`B_z/\sqrt{2\pi \sigma }`$ which arises naturally has dimensions not of (a 2-D Alfvén) velocity but rather of length<sup>1/2</sup>/time. A length scale is thus required. A logical choice in this geometry is the scale height of the slab, $`H=a^2/2\pi \sigma G`$, but this is undesirable because the problem of most interest is a cold cloud, in which the isothermal sound speed $`a^20`$. Instead we choose a characteristic length scale $`L`$, which will be the horizontal domain size, or equivalently the largest spatial wavelength in the simulation (Of course, $`L`$ scales out of all final results when expressed in dimensional units). The characteristic velocity is the Alfvén speed for the critical surface density and magnetic field
$$v_{a0}B_{z0}\sqrt{\frac{L}{2\pi \sigma _{c0}}},$$
(2-7)
and the characteristic time is simply
$$t_{c0}\frac{L}{v_{a0}}.$$
(2-8)
The nondimensionalized variables are as follows:
$$\begin{array}{cccccc}\hfill \omega & & \frac{\sigma }{\sigma _{c0}},\hfill & \hfill 𝜷& & \frac{𝐁}{B_{z0}},\hfill \\ \hfill \tau & & \frac{t}{t_{c0}},\hfill & \hfill (\mu ,\nu )=𝝂& & \frac{𝐯}{v_{a0}},\hfill \\ \hfill _h& & \frac{_h}{L},\hfill & \hfill \xi ,\eta ,\zeta & & \frac{x,y,z}{L},\hfill \\ \hfill \varphi _G& & \frac{t_{c0}^2}{L^2}\mathrm{\Phi }_G,\hfill & \hfill \text{and}\varphi _M& & \frac{\mathrm{\Phi }_M}{B_{z0}L}.\hfill \end{array}$$
(2-9)
### 2.3 Parameterization of Ambipolar Drift
We assume that the product of the ion surface density $`\sigma _i`$ and the ion-neutral collision frequency $`\nu _{in}`$ is related to the surface density $`\sigma \sigma _n`$ according to the simple ansatz:
$$\frac{\delta (\sigma _i\nu _{in})}{\sigma _i\nu _{in}}=\alpha \frac{\delta \sigma }{\sigma }.$$
(2-10)
If the ionization fraction $`x`$ scales as the neutral density $`n_n^q`$, and the scale height $`H`$ scales as $`\sigma ^1`$, then $`\alpha =32q`$. Often $`q`$ is taken to be 0.5 (McKee et al., 1993), but a detailed treatment of grain dynamics in contracting cores (Ciolek & Mouschovias, 1994, 1995, 1998) shows that the parameter $`q`$ continuously decreases throughout contraction, and may range from $`0.6`$ to less than $`0.1`$ as the central density increases by about 6 orders of magnitude. We have tested the sensitivity of our results to the value of $`\alpha `$ by comparing models with $`\alpha `$ ranging from 0 to 3, and find that the results differ by less than 5% as the density perturbations grow from .01 to 1.5 times the mean density (This is consistent with linear perturbation theory, which predicts that the value of $`\alpha `$ enters only if the unperturbed initial state has an inclined magnetic field; Z98). In view of the insensitivity of the results to $`\alpha `$, as well as the fact that the volume density increases by less than 2 orders of magnitude in our calculations, we regard the ansatz equation (2-10) as adequate.
To compare with the linear theory, we use a nondimensional form of the drift frequency $`\mathrm{\Gamma }=t_{c0}kB_{z0}^2/2\pi \sigma _i\nu _{in}`$. (Z98 uses the dimensional $`\mathrm{\Gamma }=kB_{z0}^2/2\pi \sigma _i\nu _{in}`$.) In the simulation, it is convenient to compute the evolution of the drift independent of spatial wavenumber $`k`$. We use the quantity $`\mathrm{\Gamma }/kL`$, which is $`\mathrm{\Gamma }/2\pi `$ for the lowest spatial wavenumber $`k=2\pi /L`$.
The full set of governing equations in conservative form are as follows:
$`{\displaystyle \frac{\omega }{\tau }}`$ $`=`$ $`\mathbf{}_h(\omega 𝝂_h),`$ (2-11a)
$`{\displaystyle \frac{}{\tau }}\omega 𝝂_h`$ $`=`$ $`\mathbf{}_h(\omega 𝝂_h𝝂_h)a^2\mathbf{}_h\omega `$ (2-11b)
$`\omega \mathbf{}_h\varphi _G+\beta _z𝜷_h,`$
$`\omega `$ $`=`$ $`{\displaystyle \frac{\varphi _G}{\zeta }}|^+,`$ (2-11c)
$`𝜷_h`$ $`=`$ $`\mathbf{}_h\varphi _M,`$ (2-11d)
$`\beta _z`$ $`=`$ $`{\displaystyle \frac{\varphi _M}{\zeta }},`$ (2-11e)
$`{\displaystyle \frac{\beta _z}{\tau }}`$ $`=`$ $`\mathbf{}_h\left[\beta _z(𝝂+𝝂_D)\right],`$ (2-11f)
$`𝝂_D`$ $`=`$ $`\beta _z𝜷_h{\displaystyle \frac{\mathrm{\Gamma }}{kL}},`$ (2-11g)
$`{\displaystyle \frac{}{\tau }}\left({\displaystyle \frac{kL}{\mathrm{\Gamma }}}\right)`$ $`=`$ $`{\displaystyle \frac{kL}{\mathrm{\Gamma }}}{\displaystyle \frac{\alpha }{\omega }}{\displaystyle \frac{\omega }{\tau }}.`$ (2-11h)
Note that we interpret equation (2-10) as an Eulerian relation in equation (2-11h).
### 2.4 Physically Reasonable Parameter Regime
The input parameters for the model are the sound speed $`a`$, the strength of the initial magnetic field relative to the surface density $`B_{z0}/2\pi G^{1/2}\sigma _0=1/\omega _0`$, the drift parameter $`\mathrm{\Gamma }`$, and the constant $`\alpha `$ which determines the perturbation to the collision rate (see eq. \[2-10\]).
Typical magnetic fields in dense clouds are 30$`\mu `$G (Myers & Goodman, 1988b), and we choose the nondimensionalization length scale (see §2.2) to be 1pc, a typical size for a dense cloud or cloud core and its close neighborhood. A typical cold cloud temperature is 10K, and the average molecular weight is that of molecular hydrogen with 10% helium, $`m_n`$ = 3.9$`\times `$10<sup>-24</sup>g. This yields the following expressions for the Alfvén and sound speeds:
$`v_{a0}^2`$ $`=`$ $`B_{z0}^2{\displaystyle \frac{L}{2\pi \sigma _{c0}}}`$
$`=`$ $`B_{z0}LG^{1/2}`$
$`=`$ $`2.4\times 10^{10}{\displaystyle \frac{\mathrm{cm}^2}{\mathrm{s}^2}}\left({\displaystyle \frac{B_{z0}}{30\mu \mathrm{G}}}\right)\left({\displaystyle \frac{L}{\mathrm{pc}}}\right),`$
$`a^2`$ $`=`$ $`{\displaystyle \frac{k_BT}{m,}}`$ (2-12b)
$`{\displaystyle \frac{a^2}{v_{a0}^2}}`$ $``$ $`0.02\left({\displaystyle \frac{T}{10\mathrm{K}}}\right)\left({\displaystyle \frac{\mathrm{pc}}{L}}\right)\left({\displaystyle \frac{30\mu \mathrm{G}}{B_{z0}}}\right)`$ (2-12c)
$`\times \left({\displaystyle \frac{3.9\times 10^{24}\mathrm{g}}{m_n}}\right).`$
From now on, $`a^2`$ will be given in units of $`v_{a0}^2`$.
The initial surface density is chosen to be $``$ 0.02 g cm<sup>-2</sup>, or about 100 M pc<sup>-2</sup>. This corresponds not only to a typical column density for a dense cold cloud ($`N_H10^{22}`$ cm<sup>-2</sup>), but also to the surface density that would result if a spherical cloud of typical number density (10<sup>4</sup> cm<sup>-3</sup>) and typical size (several 10<sup>18</sup>cm) collapsed along a large-scale magnetic field into a thin pancake or disc. The following expression for the normalized surface density results:
$`\omega _0`$ $`=`$ $`{\displaystyle \frac{2\pi G^{1/2}\sigma _0}{B_{z0}}}`$
$``$ $`1.0\left({\displaystyle \frac{30\mu \mathrm{G}}{B_{z0}}}\right)\left({\displaystyle \frac{\sigma _0}{0.02\mathrm{g}\mathrm{cm}^2}}\right).`$
Thermal pressure raises the critical surface density for gravitational collapse, and for a cloud at T = 10K the value of $`\omega _0`$ for the fiducial parameters in equation (2.4) is 0.864 of that critical surface density $`\omega _{\mathrm{crit}}`$.
Thermal pressure imparts a finite scale height to the disk
$`H`$ $`=`$ $`{\displaystyle \frac{a^2}{2\pi \sigma _0G}}`$
$``$ $`4\times 10^{16}\mathrm{cm}\left({\displaystyle \frac{T}{10\mathrm{K}}}\right)\left({\displaystyle \frac{0.02\mathrm{g}\mathrm{cm}^2}{\sigma _0}}\right)`$
$`\times \left({\displaystyle \frac{3.9\times 10^{24}\mathrm{g}}{m_n}}\right).`$
Clearly, $`HL`$. An inclined magnetic field exerts a pinching force, compressing the disk and reducing $`H`$ further (Wardle & Königl, 1993).
Determination of the drift parameter $`\mathrm{\Gamma }`$ requires an expression for the neutral-ion collision frequency, $`\nu _{ni}2\times 10^9n_i`$ cm<sup>3</sup> s<sup>-1</sup> (Draine, Roberge, & Dalgarno, 1983) and the density of ions. We assume $`n_i=Kn_n^{1/2}`$, where $`K1.1\times 10^5n_n^{1/2}`$ cm<sup>-3/2</sup> (McKee et al., 1993). Using these relations
$`\mathrm{\Gamma }`$ $`=`$ $`{\displaystyle \frac{kt_{c0}B_{z0}^2}{2\pi \sigma _i\nu _{in}}}`$
$`=`$ $`{\displaystyle \frac{kt_{c0}B_{z0}^2}{2\pi \sigma _0\nu _{ni}}}`$
$``$ $`0.059\left({\displaystyle \frac{\mathrm{pc}}{L}}\right)^{\frac{1}{2}}\left({\displaystyle \frac{B_{z0}}{30\mu \mathrm{G}}}\right)^{\frac{3}{2}}\left({\displaystyle \frac{0.02\mathrm{g}\mathrm{cm}^2}{\sigma _0}}\right)`$
$`\times \left({\displaystyle \frac{5\times 10^4\mathrm{cm}^3}{n_n}}\right)^{\frac{1}{2}}\left({\displaystyle \frac{k}{2\pi /L}}\right).`$
The fiducial value of $`n_n`$ which appears here is consistent with the other parameters: $`n_n=\pi G\sigma _0^2/k_BT`$.
Although equation (2.4) shows that $`\mathrm{\Gamma }`$ does not depend on the scale height $`H`$, since the temperature of molecular clouds is quite well determined it is useful to rewrite $`\mathrm{\Gamma }`$ in a way which does depend on $`H`$ and suppresses the dependence on some of the other parameters. We have
$`\mathrm{\Gamma }`$ $`=`$ $`{\displaystyle \frac{4\pi ^{3/2}(m_nG)^{3/2}}{K\sigma v\omega _0^{3/2}}}\left({\displaystyle \frac{H}{L}}\right)^{1/2}`$
$`=`$ $`0.56\left({\displaystyle \frac{H}{L}}\right)^{1/2},`$
where we have used the standard values for all the constants. Equations (2.4) and (2.4) suggest $`\mathrm{\Gamma }0.1`$ for typical parameters.
### 2.5 Linear Theory
The collapse rate for a linear perturbation with positive thermal pressure can be easily calculated (see Z98 for the calculation at T=0K). The physical quantities $`\omega `$, $`\beta _z`$, $`𝝂`$ and governing equations (2-5a-2-5h) are linearized and reduced to one horizontal spatial dimension. Assuming a single Fourier mode
$`\omega `$ $``$ $`\omega _0+\omega e^{\gamma \tau +ik\xi k|\zeta |},`$ (2-17a)
$`\beta _z`$ $``$ $`1+\beta _ze^{\gamma \tau +ik\xi k|\zeta |},`$ (2-17b)
$`\nu `$ $``$ $`\nu e^{\gamma \tau +ik\xi k|\zeta |}`$ (2-17c)
leads to the dispersion relation
$$\gamma ^3+\mathrm{\Gamma }\gamma ^2+\left[\gamma _G^2\left(\frac{1}{\omega _0^2}1\right)+\gamma _T^2\right]\gamma \gamma _G^2\mathrm{\Gamma }+\gamma _T^2\mathrm{\Gamma }=0,$$
(2-18)
where $`\gamma _G=\sqrt{\omega _0kL}`$ is the nondimensionalized gravitational frequency, and $`\gamma _T=akL`$ is the nondimensionalized thermal frequency.
Equation (2-18) has two limits which provide some insight into what follows. In the limit $`\mathrm{\Gamma }=0`$ we recover the dispersion relation for waves driven by magnetic tension, thermal pressure, and self-gravity; the first two forces are stabilizing and the last is destabilizing. The system is stable for all wavenumbers if $`\omega _0<1`$, but if $`\omega _0>1`$, the system is unstable for wavenumbers $`k<k_c(\mathrm{\Gamma }=0)H^1(\omega _0^21)/\omega _0^2`$. (In a system of finite size $`L`$, $`k`$ is bounded from below by $`2\pi /L`$, leading to the absolute stability criterion $`\omega _0<\omega _{crit}`$ which we present below). The maximum growth rate, in dimensional form (recall that in equation (2-18), $`\gamma `$ is given in units of $`t_{c0}^1`$ ), is $`\gamma _{max}(\mathrm{\Gamma }=0)=\pi G\sigma _0(\omega _0^21)/(a\omega _0^2)`$, and occurs at a wavenumber $`k_m(\mathrm{\Gamma }=0)=\frac{1}{2}k_c(\mathrm{\Gamma }=0)`$ (in these dimensional expressions, $`a`$ is the dimensional sound speed).
In the limit of large $`\mathrm{\Gamma }`$, the magnetic field is uncoupled from the gas, and the dispersion relation reverts to that of an unmagnetized slab. The system is unstable for $`\gamma _G^2>\gamma _T^2`$, or $`k<k_c(\mathrm{\Gamma }\mathrm{})2\pi G\sigma _0/a^2=H^1`$. The maximum growth rate, which occurs at $`k_m(\mathrm{\Gamma }\mathrm{})=\frac{1}{2}k_c(\mathrm{\Gamma }\mathrm{})`$, is $`\gamma _{max}(\mathrm{\Gamma }\mathrm{})=\pi G\sigma _0/a`$. (Again, in these dimensional expressions, $`a`$ is dimensional).
The maximum growth rate, and the wavenumber at which it occurs, is always less for a magnetized but supercritical cloud than for an unmagnetized cloud, and, as expected, the supercritical case approaches the unmagnetized case as the magnetic fieldstrength decreases to zero.
In this paper we are interested in clouds which are magnetically subcritical, so that they would be stable in the limit $`\mathrm{\Gamma }=0`$, but would be unstable if the magnetic field were removed. That is, we are interested in clouds with a length much larger than the unmagnetized Jeans length. A small but nonzero $`\mathrm{\Gamma }`$ destabilizes a cloud to perturbations which are stabilized by magnetic fields in the absence of ambipolar drift, but would be unstable to the Jeans instability in the absence of magnetic fields.
For small sound speeds, the dispersion relation shows the same behavior seen for zero temperature in Z98: at low values of the drift parameter $`\mathrm{\Gamma }`$, the growth rate of the perturbation $`\gamma `$ is proportional to $`\mathrm{\Gamma }`$. At higher $`\mathrm{\Gamma }`$, however, $`\gamma \mathrm{\Gamma }^{1/3}`$. As surface densities $`\omega `$ and wavenumbers $`k`$ depart from the critical values for stability, more ambipolar drift is required for the system to show $`\gamma \mathrm{\Gamma }^{1/3}`$ behavior. To quantify these statements with an example, at $`\omega _0`$ = 1.1, $`a^2`$ = 0.02, $`\mathrm{\Gamma }`$ = 0, the critical $`k`$ below which ideal perturbations are unstable is $`k`$ = 9.5455. Very close to criticality, $`k`$ = 9.6, the $`\gamma \mathrm{\Gamma }^{1/3}`$ scaling holds for $`\mathrm{\Gamma }`$ as small as .001. At $`k`$ = 10, $`\gamma `$ increases with $`\mathrm{\Gamma }`$ faster than $`\mathrm{\Gamma }^{1/3}`$, but much slower than linearly, for $`.001<\mathrm{\Gamma }<.01`$, but approaches the $`\mathrm{\Gamma }^{1/3}`$ scaling for $`.01<\mathrm{\Gamma }<.1`$. As $`\mathrm{\Gamma }`$ is increased from .001 to .1, $`\gamma `$ increases from .378 to 1.75, which is $`\mathrm{\Gamma }^{.33}`$ scaling. At $`k`$ = 4$`\pi `$, this scaling has broken down noticeably: as $`\mathrm{\Gamma }`$ is increased from .001 to .1, $`\gamma `$ increases from .170 to 1.93, which is $`\mathrm{\Gamma }^{.53}`$ scaling. Most of the deviation occurs for small values of $`\mathrm{\Gamma }`$; for $`\mathrm{\Gamma }`$ between .01 and .1, $`\gamma \mathrm{\Gamma }^{.36}`$. Thus, for reasonable values of $`a^2`$ and $`\mathrm{\Gamma }`$, the $`\mathrm{\Gamma }^{1/3}`$ scaling law holds quite well even when $`k`$ is as much as 30% below the critical value.
The addition of thermal pressure increases the stability of the disk; more drift (larger $`\mathrm{\Gamma }`$) is required for collapse, and more is required to reach the transition from $`\gamma \mathrm{\Gamma }`$ to $`\gamma \mathrm{\Gamma }^{1/3}`$. An approximate value for the critical surface density for collapse, with positive thermal pressure, is obtained from solving the dispersion relation (see eq. \[2-18\]) for $`\mathrm{\Gamma }`$=0:
$`{\displaystyle \frac{1}{\omega _{\mathrm{crit}}}}`$ $`=`$ $`\pi a^2+\sqrt{1+\pi ^2a^4}`$
$``$ $`1\pi a^2,`$
$`\omega _{\mathrm{crit}}`$ $``$ $`1+\pi a^2,`$ (2-19b)
where the approximations hold for small sound speed. Solution of the dispersion relation also shows that for a given sound speed and drift parameter, there is a single mode with a maximal growth rate, and that the modes above a certain wavenumber are acoustically suppressed. Figure 1 shows the growth rate as a function of wavenumber for $`\mathrm{\Gamma }`$ = 0.1, and most unstable wavenumber, over a range of sound speeds. Figure 1 also shows that the stability boundary is very near the thermal Jeans stability boundary, while the fastest growing mode has a much longer wavelength than the Jeans wavelength.
As we will see later, the wavenumber at which the growth rate is maximized dominates the structure of clumps even into the nonlinear regime. Numerical solution of the dispersion relation equation (2-18) shows that $`k_m`$ is always less than $`k_m(\mathrm{\Gamma }\mathrm{})`$, the fastest growing wavenumber for the Jeans instability, but also scales with temperature as $`1/T`$. Therefore, we expect the fragment mass to be larger than the thermal Jeans mass, but to have the same $`T^2`$ temperature scaling. This is borne out by Figures 1 and 6.
## 3 Numerical Simulation
In order to follow the instability into the nonlinear regime we have carried out a numerical simulation. We use a Fourier collocation pseudo-spectral method (Canuto et al., 1988) to solve the governing equations. Values of physical quantities are stored at discrete points in physical space (known as collocation points), and spatial derivatives are evaluated in Fourier spectral space (hence the name “Fourier pseudo-spectral”). This particular method is well adapted to this problem for several reasons. Calculating a spatial derivative in Fourier space simply requires multiplication of each Fourier component by $`i`$ times its wavenumber. All terms involving the horizontal magnetic field $`𝐁_h`$ or gravitational potential $`\mathrm{\Phi }_G`$ are trivial to evaluate in Fourier space due to the simple form of the magnetic and gravitational potentials. Nonlinear terms, on the other hand, are trivial to evaluate in physical space by simple multiplication. Finally, growth of different Fourier modes can be monitored and controlled explicitly, simplifying comparison between the nonlinear numerical model and the single-wavenumber linear analytic results.
We use a Bulirsh-Stoer time-stepping routine with Richardson extrapolation (Press et al., 1986). The routine performs several modified midpoint method integrations at sub-intervals of the desired time-step. It then attempts to extrapolate to an infinite number of sub-intervals. The routine varies the number of explicit (calculated) sub-intervals based on the estimated error. In general, the full set of governing equations for this problem can be integrated with $``$ 5 explicit subintervals for each time-step of $`\delta \tau =0.1`$.
We tested convergence by running the code at increasing spatial resolution with the same initial conditions. Figure 2 shows the amplitude of a density perturbation, computed at different resolutions, as a function of time. The initial conditions had a single wavenumber density perturbation in each direction, forming a “checkerboard” pattern. (Collapse is described in detail below.) A 16<sup>2</sup> grid is sufficient to resolve the collapse of such a single wavenumber, from the linear regime (exponential growth) into the nonlinear regime. Use of a 32<sup>2</sup> or 64<sup>2</sup> grid changes the solution by less than 0.1% over most of the time period plotted. Finer grids are required to resolve collapse of smaller structures, and in runs which contain a spectrum of wavenumbers, a 64<sup>2</sup> grid was used.
We have also verified that the code reproduces the results of linear theory. If the initial condition corresponds to an eigenfunction of a growing mode calculated according to linear perturbation theory, with an amplitude of a few percent or less, then the disturbance initially grows at the exponential rate predicted by the linear theory. This is shown in Figure 3, which compares the growth rates measured from the code (discrete symbols) with the continuous curve obtained from solving the dispersion relation equation (2-18). The agreement is generally excellent.
The model is numerically stable until such time as power in the higher order wavenumbers grows to overwhelm power in the Fourier modes of interest. At that time, the simulation develops a “sawtooth” instability, with large variation between alternating collocation points. Power grows in these short-wavelength modes from numerical noise, whose magnitude is about 10<sup>-8</sup> compared to the power in the principal mode (measured in simulations whose initial conditions contained a single mode). Higher wavenumber modes are also driven by the nonlinearity of the problem, and this is the dominant physical source of power in those modes.
Growth of high-order modes can be controlled in several ways. Thermal pressure will stabilize high order modes, as was seen in Figure 1. In most cases, a physically reasonable finite cloud temperature of $``$ 10 K will stabilize the simulation long enough to follow the collapse well into the nonlinear regime. The problem of high-order mode stabilization is nearly independent of the drift parameter $`\mathrm{\Gamma }`$ because an increase (decrease) in ambipolar diffusion increases (decreases) the growth rate of all modes similarly. Thus the entire physically interesting part of parameter space is numerically accessible and numerically stable in this model.
### 3.1 Collapse Rate
Many runs were computed with initial conditions corresponding to the eigenfunction of the fastest-growing solution of the 3 modes present, at each wavenumber, in linear theory. This initial perturbation has a single initial wavenumber in each direction. Growth initially proceeds exponentially, with faster nonlinear collapse occurring as higher order spatial modes are driven. Nonlinear behavior is typically seen when the density enhancement associated with a perturbation has grown to 75%-100% of the mean density. We were able to follow the evolution of the system to peak surface densities about 10 times larger than the mean density (corresponding to a peak volume density about 100 times larger than the mean).
We made a detailed comparison with linear theory and with simple drift and collapse models by studying the early collapse of a perturbation with a single spatial wavenumber. Figure 4 shows how the growth of a fully nonlinear density perturbation depends on the drift parameter $`\mathrm{\Gamma }`$. The lower three curves describe collapse in subcritical clouds, which would be stabilized by the magnetic field in the absence of ambipolar drift. The initial surface density (normalized to the vertical magnetic field) is $``$ 86% of the critical surface density for collapse. The growth rate of perturbations in the nonlinear simulation (calculated from the time for the central density of the perturbation to grow from 1% to 100% of the mean density) is compared to the predicted growth rate $`\gamma `$ for linear perturbations, and to the ambipolar drift rate $`\mathrm{\Gamma }`$. Clearly, for the physically expected value of $`\mathrm{\Gamma }`$ ($``$ 0.05 - 0.10), §2.4), the collapse due to this instability is several times faster than simple collapse due to loss of magnetic support on a diffusive ambipolar drift time-scale. For example, when $`\mathrm{\Gamma }`$ = 0.1, the ambipolar drift rate is 0.1 and the unmagnetized collapse rate is 10. The rate of contraction found from the simulation is 0.4. At larger values of $`\mathrm{\Gamma }`$, the drift rate comes closer to the unmagnetized collapse rate. The collapse rate in the simulation approaches the unmagnetized collapse rate because the coupling between the magnetic field and the gas is weak. When the drift becomes very important ($`\mathrm{\Gamma }10`$), the drift timescale is very short, and the intermediate instability described in this paper is less significant. Even for the moderate sized density perturbations used to create Figure 4, ($`\delta \omega \omega `$), the growth rate is larger in the nonlinear simulation than in the linear problem, especially at relatively large values of $`\mathrm{\Gamma }`$, showing that the collapse is accelerated by the nonlinearity. It is important to note, however, that the growth rate shows the same dependence on ambipolar drift in both the linear theory and the nonlinear simulation: at low $`\mathrm{\Gamma }`$, $`\gamma \mathrm{\Gamma }`$, and at higher $`\mathrm{\Gamma }`$, $`\gamma \mathrm{\Gamma }^{1/3}`$.
The top curve in Figure 4 describes collapse in supercritical clouds, in which the magnetic field would be insufficient to prevent collapse even if it were frozen to the matter. The collapse rate depends only weakly on the strength of the ambipolar drift, as expected since the magnetic field is dynamically less important. When the ambipolar drift strength $`\mathrm{\Gamma }`$ becomes large, the drift timescale becomes comparable to the collapse timescale, and the subcritical and supercritical cases converge. Rapid collapse occurs in supercritical clouds due to the dynamical weakness of the field, and in subcritical clouds rapid diffusion removes magnetic support, quickly rendering them supercritical.
The evolution of self gravitating, subcritical disks with ambipolar drift was studied previously by Ciolek & Mouschovias (1994) and Basu & Mouschovias (1995). Ciolek & Mouschovias (1994) began with a subcritical ($`\omega _0`$ = .256), centrally condensed equilibrium state - the central surface density is 16 times the mean density. Thus, this model is more centrally condensed even initially than our models are when we terminate the simulation. The initial ambipolar drift time is 10 times the initial free fall time, which corresponds on Figure 4 to $`\mathrm{\Gamma }1`$. The evolutionary timescale in the subcritical, quasistatic phase is well estimated by the initial ambipolar drift time; after the cloud becomes supercritical its collapse rate approaches the freefall rate. Although we can extrapolate our results to this model only with caution, because the initial conditions are so different from ours, it does not surprise us that such a subcritical disk shows no evidence for the blending of dynamical and drift effects that we observe closer to criticality.
Basu & Mouschovias (1995) carried out a parameter study to determine the effects of the degree of criticality on the rate of evolution to a critical state. They also began with centrally condensed equilibrium models, forming a sequence in which the criticality parameter varied from 0.1 to 0.5. They found that the timescale for evolution to the critical state decreased by about a factor of 1.5 along this sequence, from somewhat longer than the estimated drift time to about 25% shorter (another, marginally critical model, collapsed at once). Although again a quantitative comparison of our models with theirs is difficult because of the different initial conditions, it is possible that the intermediate contraction rates which they see are a manifestation of the coupling between dynamical and ambipolar drift effects seen in our models. It may also be due to the increased central concentration of the initial equilibrium states along their sequence of models.
The nature and rate of collapse is observable in molecular clouds (Evans, 1999; Myers, Evans, & Ohashi, 1999), and an instability with an intermediate growth rate such as this one can help to explain observations that do not fit either of the classical scenarios - dynamical collapse or slow diffusive contraction. Our simulated cloud cores collapse with slower velocities and on larger physical scales than the dynamical inside-out collapse predicted when the magnetic field is unimportant, as in Shu (1977). Tafalla et al. (1998) and Gregersen (1998) have observed cores that appear to have such behavior; they find that the regions of inflow are too large to fit dynamical inside-out collapse models, but that the inflow velocities are too large for quasistatic diffusion models.
### 3.2 B-$`\rho `$ relation
The correlation between fieldstrength and density is an observable quantity which can provide insight into the manner in which the magnetic field evolves. It is useful to parameterize this relation as a power law
$$|𝐁|\rho ^\kappa .$$
(3-1)
Observations find $`\kappa `$ 0.5 (Troland & Heiles, 1986; Crutcher, 1999) over several orders of magnitude in density. In the case of a highly flattened cloud such as we simulate here, it is more convenient to define a magnetic field - surface density relation
$$|𝐁|\sigma ^\lambda .$$
(3-2)
If the slab scale height $`H`$ remains constant throughout the evolution, then $`\lambda =\kappa `$. If the scale height is determined by a balance between self-gravity and thermal pressure alone, then an isothermal slab obeys $`H1/\sigma `$, and $`\sigma \rho H`$ implies $`\sigma \rho ^{1/2}`$, or $`\kappa =0.5\lambda `$ (Crutcher, 1999; Spitzer, 1942).
If the magnetic field is frozen to the matter but not dynamically important, so contraction is isotropic, conservation of flux $`\mathrm{\Phi }_{mag}L^2|𝐁|`$ and mass $`ML^3\rho `$ requires $`|𝐁|\rho ^{2/3}`$ ($`\kappa =`$ 2/3). If the field is so strong that matter moves one dimensionally, parallel to the fieldlines, then $`\kappa 0`$. Calculations in which a cloud condenses to magnetohydrostatic equilibrium from a uniform initial state, with frozen in magnetic flux of a magnitude appropriate to the ISM, show anisotropic contraction, and the central values of $`B`$ and $`\rho `$ in the initial and final states are related by $`\kappa 0.5`$ (Mouschovias & Spitzer, 1976; Tomisaka, Ikeuchi, & Nakamura, 1988). If a cloud is already flattened and shrinks transversely, $`ML^2\sigma `$, and flux freezing implies $`|𝐁|\sigma ^1`$ ($`\lambda `$=1, $`\kappa `$=0.5). However, rather different input physics leads to a similar exponent: simulations of supersonic magnetized turbulence, without self-gravity, produce $`\kappa 0.4`$ if the field is not too strong (Padoan & Norlund, 1999).
Ambipolar drift generally reduces $`\kappa `$ below the value it would have if the field were frozen in. In the models of Fiedler & Mouschovias (1993), $`\kappa `$ averages about 0.2 during the so-called quasistatic phase. After the quasistatic phase ends, the mean value of $`\kappa `$ is 0.3 as $`\rho `$ increases by more than 5 orders of magnitude. In the highly flattened models of Ciolek & Mouschovias (1994), $`\kappa `$ increases smoothly as the cloud evolves from subcritical to supercritical, reaching peak values between 0.4 and 0.5 and being about 0.3 at criticality. In our simulations the magnetic field - surface density exponent $`\lambda `$ varies between 0.35 and 0.65. As shown in Figure 5, the exponent decreases as the central density of a clump increases. As collapse proceeds, not only does the density increase, but the magnetic field curvature also increases as fieldlines are dragged into the condensation. Both effects increase the ambipolar drift velocity $`𝐯_D=B_z𝐁_h/2\pi \sigma \nu _{ni}`$ and thus the rate of flux loss from the clump. Our models do not show the increase of $`\lambda `$ toward its frozen flux value as the central density increases seen in Ciolek & Mouschovias (1994), because we follow only the early stages of contraction, in which the velocity is well below the freefall value. The exponent $`\lambda `$ also decreases as the amount of ambipolar drift $`\mathrm{\Gamma }`$ increases, as would be expected, and as the cloud temperature increases. The latter effect results from the decreased efficiency of this instability in warm clouds. The collapse rate $`\gamma `$ decreases with increasing temperature as shown in §2.5, and the collapse time is longer relative to the ambipolar drift time, so more flux can leak from the clump as it collapses.
Comparison with observations depends on the assumed disc scale height $`H`$ and the central densities of observed cloud cores. Crutcher (1999) finds $`\kappa `$ = 0.47 for cloud cores with densities 10<sup>2.5</sup>cm<sup>-3</sup> $`n_H`$ 10<sup>7.5</sup>cm<sup>-3</sup>. His data are strongly weighted by observations which only measure upper limits for the magnetic fieldstrength, and omission of those data points results in an exponent of $`\kappa `$ = 0.3 over the same range of densities. We measure 0.3 $`\lambda `$ 0.7 in simulated cores. If the slab scale height $`H`$ is constant, then 0.3 $`\kappa `$ 0.7, but if the slab obeys $`H1/\sigma `$, then 0.2 $`\kappa `$ 0.35. Our simulations are thus consistent with the small number of available observations.
### 3.3 Clumps and Fragments
Real molecular clouds have density variations on many scales, or a spectrum of spatial wavenumbers. The linear theory for this instability (§2.5) predicts that a single mode will dominate collapse. Fragments of a single mass will form, with that mass depending on the temperature and degree of criticality, with little dependence on the strength of ambipolar drift, provided that it is small. We have simulated many ($``$ 100) clouds in which the initial density perturbations have a broad spatial Fourier spectrum and random phase. The real part of the initial Fourier spectrum of the density perturbation is a Gaussian centered on $`k`$ = 0 (with the omission of the $`k`$ = 0 mode itself). The FWHM is 0.33 $`k_{max}`$, where $`k_{max}`$ is the highest spatial mode in the simulation, so a range of low-order modes have similar initial strength, and there is significant initial power even in some acoustically damped (high $`k`$) modes. After some time, the clouds coalesce into a small number of fragments, each of which we define to be a region with $`\sigma >\sigma _0`$. The size of the fragments is well-predicted by the linear theory. For example, when $`a^2`$ = 0.02, the fastest linearly growing mode is $`k/2\pi 2.5`$, and in the simulation, the modes ($`k_x,k_y`$)$`/2\pi `$ = (1,3) and (3,1) have much larger amplitudes than other spatial modes.
Our simulations produce a wide variety of clumps and fragments, as expected for systems with random initial density fluctuations. Figure 6 shows the fragment masses at a fairly early stage of collapse (the density perturbation is $``$ 50% greater than the mean density). There is considerable scatter, but the linear theory is a good guide to the average fragment mass. The clump masses are similar to the masses of the supercritical cores formed in some previous simulations (Fiedler & Mouschovias, 1993; Basu & Mouschovias, 1994; Ciolek & Mouschovias, 1994; Ciolek & Königl, 1998).
At sufficiently large $`T`$, the spatial mode with the largest linear growth rate is the fundamental mode in the simulation domain ($`\lambda =L`$), and fragments larger than this cannot form in a periodic simulation. This explains the apparent flattening of the clump mass – temperature relation seen at the higher temperatures, as clump masses are bounded above by the mass of that “lowest mode”.
The question of whether there is a characteristic mass for molecular cloud substructure is an important one. Several observational studies have found a power law distribution of clump masses over several orders of magnitude, $`dN/dMM^p`$, where $`p1.51.7`$ (Blitz, 1993). Kramer et al. (1998) present evidence that the power law extends far below 1 $`M_{}`$. However, in the Taurus molecular cloud there is evidence of a minimum scale of a few tenths of a parsec, corresponding to several solar masses (Blitz & Williams, 1997). Goodman et al. (1998) have argued for an inner scale of 0.1 pc, which they identify with a transition to what they term “velocity coherence”. These inner scales are of the same order as the thermal Jeans length, and also close to the cutoff wavelength below which Alfvén waves are critically damped due to strong ambipolar drift (McKee et al., 1993). Our results suggest that there is another scale, which is somewhat larger, in magnetically subcritical clouds with weak ambipolar drift.
### 3.4 Velocity Structure
Unlike axisymmetric collapse models, these simulations allow the study of asymmetric collapse, clumps with complicated morphology, relative motion of clumps, and their internal velocity and vorticity fields. We find that collapse is often asymmetric and that significant vorticity is generated by the instability (although of course the net angular momentum remains zero in our simulations; its absolute value is more than an order of magnitude smaller than the estimated numerical errors).
Figure 7 shows four examples of collapsing fragments. The left panels show contours of constant surface density, arrows indicating the local direction and magnitude of velocity, and bold arrows indicating the center of mass motion of each clump. The velocity field shows infall towards clumps, but there is also a visual impression of “swirling” or rotational motion. This is borne out by the right panels of Figure 7, which show contours of constant vertical vorticity overlaid on contours of constant density. We show below that the magnetic field generates local vorticity.
The velocity field indicates that the collapse is in most cases asymmetric, with e.g. much greater infall velocities on one side of the clump than the other, and significant nonradial motion. Clump mergers are possible - one is quite likely taking place in the bottom panel of Figure 7 \- but the bulk motions are significantly slower than the infall velocities internal to clumps. The reverse is generally true in molecular clouds (e.g. Blitz, 1993), and this result in our simulations is a manifestation of the fact that the velocity field in the system is weak. We return to this point in §3.5. Sometimes the clumps move apart, but this is ambiguous in a periodic domain. The infall velocities within individual clumps are of order 0.25 $`a`$, and are more consistent with observations (Evans, 1999; Myers, Evans, & Ohashi, 1999).
Visual inspection of Figure 7 shows that the clumps are distinctly noncircular and quite elongated in shape. Accounting for the third dimension, our clumps should be considered triaxial or prolate. These shapes are consistent with the measured and inferred shapes reported by Myers et al. (1991), Ryden (1996), and Ward-Thompson, Motte, & André (1999).
Figure 8 shows a cross section of the surface density and velocity profiles across the short axis of one particular collapsing clump. The density is much more peaked than the velocity at this early stage of collapse; the FWHM of the infall speed is several times larger than the FWHM of the density peak. Both the velocity and density profiles are clearly, but not grossly, asymmetric. Although it is premature to compare these density and velocity profiles with observations of infall, it is encouraging that we see evidence for extended inward motions as have been reported (Tafalla et al., 1998; Gregersen, 1998; Evans, 1999; Myers, Evans, & Ohashi, 1999). A general feature of infall onto a line mass such as a filament or strongly prolate object is that the velocity decays more slowly with distance from the mass centroid than for infall onto a spherically symmetric, centrally concentrated object. This may be the main effect that produces the extended infall.
The vorticity generated in swirling motions is $`\times v`$ 5/$`t_{c0}`$ 8/Myr. Collapse proceeds on a timescale of several Myr, so the swirling and rotation of clumps is not insignificant, although it is dominated by infall and we have not found evidence for clumps torn apart by shear. Magnetic braking, which is excluded from our calculation by the potential field approximation, would reduce rotation. We estimate the magnetic braking rate in §4.1.
Figure 7 shows that the vorticity maxima are displaced from the density maxima. In order to understand this, we derive an evolution equation for the $`\widehat{z}`$ component of vorticity, $`\omega _z`$, by taking the curl of the equation of motion (2-5b)
$$\frac{\omega _z}{t}+\mathbf{}_h(\omega _z𝐯_h)=𝐁_h\times \mathbf{}_h\frac{B_z}{2\pi \sigma }.$$
(3-3)
According to equation (3-3), the generation of vorticity is second order in the amplitude of the fluctuation. There is generation of vorticity to first order only if there is a zero-order inclined field (Z98). We can understand the spatial pattern of vorticity as follows. Despite ambipolar drift, the contours of constant $`B_z`$ track the contours of constant $`\sigma `$ quite well. Therefore, the gradient of $`B_z/\sigma `$ is maximized toward the outer edge of a clump, not at its center. Note that if the magnetic flux were perfectly frozen to the matter, $`B_z/\sigma `$ would retain its initial constant value and there would be no vorticity production at all. However, real clouds probably have spatially varying $`B_z/\sigma `$, so in general, vorticity production does not require ambipolar drift. The maximum of $`𝐁_h`$, like the maximum gradient of $`B_z/\sigma `$, is displaced from the clump center. Equation (3-3) shows that therefore vorticity is generated off-center as well, and generally changes sign across the clump, so that clumps are associated with vortex pairs. Moreover, equation (3-3) shows that an axisymmetric clump does not generate vorticity. The dynamical pressure of the vortices accentuates the non-axisymmetric nature of clump contraction, and appears as streaming motions along the major axis of the clump. Although in principle equation (3-3) suggests that the vortical velocity is scaled by the Alfvén speed, in our numerical models the vortical velocities are rather small, somewhat less than the infall velocities. This is large enough to noticeably elongate the clumps, but not enough to tear them apart by shear.
### 3.5 Energy Redistribution
Z98 suggested that this magneto-gravitational instability might generate significant turbulent kinetic energy by releasing energy contained in the background magnetic field. Analysis of the total gravitational, magnetic, and kinetic energy (Fig. 9) in these simulations shows that the absolute values of all three forms of energy grow exponentially during collapse. The magnetic energy, which is the fluctuation energy integrated over the space outside the disc, dominates at all stages of collapse in these clouds, which would be stabilized by the magnetic field in the absence of ambipolar drift. The initial magnetic field is uniform in these simulations, and so there is no stored magnetic energy available for conversion to turbulent motions. The relatively low kinetic energy in the models is a signature of the importance of diffusive, as opposed to dynamical, effects. In ideal MHD turbulence with self gravity one would expect equipartition between the kinetic and potential energies (Zweibel & McKee, 1995). In this case the potential energy is the sum of the gravitational and magnetic energies, but Figure 9 shows that the kinetic energy is about an order of magnitude less than the equipartition value.
## 4 Discussion of Approximations
It is difficult and computationally expensive to simulate the evolution of magnetized, self gravitating molecular clouds in three dimensions with sufficient resolution to capture all the relevant physical processes. In this study we focussed on the two-dimensional instabilities of a sheetlike cloud surrounded by a conducting medium without inertia. This allowed us to approximate the external magnetic field as current free, and to work in only two spatial dimensions. In the next two subsections we discuss the accuracy of these approximations.
### 4.1 Potential Field Approximation
In the Appendix, we show how to extend each Fourier component of the magnetic and velocity perturbations above and below the sheet. Although the perturbations become highly nonlinear within the cloud, nonlinear effects outside the cloud are weak as long as the velocities inside the cloud are sub-Alfvénic with respect to the intercloud medium, which is expected for low ambient density. The compressive part of the cloud velocity field generates evanescent, fast magnetosonic waves which decay exponentially away from the disc, while the vortical part generates Alfvén waves which propagate away from the disc.
Both the magnetosonic and Alfvén waves slightly change the horizontal magnetic field perturbation in the disc, thereby changing the magnetic force from its value in the potential approximation. If we define an external Alfvén timescale $`\tau _{Ae}(kv_{Ae})^1`$ for wavenumber $`k`$ in the disc, and let the timescale for the perturbation in the cloud be $`\tau _c`$, then according to equation (A8) the correction to the force due to compressive motion is of order $`(\tau _{Ae}/\tau _c)^2`$, while the correction due to vortical motion is of order $`(\tau _{Ae}/\tau _c)𝒞_0/𝒟_0`$, where $`𝒞_0/𝒟_0`$ is the ratio of the amplitude of compressive to noncompressive motion. If the motions in the disc were Alfvénic, $`\tau _{Ae}/\tau _c`$ would be of order the ratio of the cloud density to external density $`(\rho _c/\rho _e)^{1/2}`$, but the perturbation frequency is sub-Alfvénic, so the ratio of timescales is even larger. Moreover, the perturbations are primarily compressive rather than vortical. Thus, the error in the force incurred by the potential approximation is likely to be small. Even if the density contrast were only 10<sup>2</sup>, and the motions in the disk were Alfvénic and purely vortical, the potential field approximation would still be accurate to 10%.
The Alfvén wave flux tends to suppress the instability, and removes vorticity from the cloud, at a rate that we can quantify. We define an energy damping time $`\gamma _d`$ as the ratio of outward propagating energy flux to the vertically integrated wave energy in the disc. From equation (A9),
$$\gamma _d=\frac{2\rho _ev_{Ae}}{\sigma }\frac{𝒞_0^2}{(𝒞_0^2+𝒟_0^2)}.$$
(4-1)
If we replace $`\sigma `$ by the critical surface density $`B_{0z}/2\pi G^{1/2}`$ and assume that the vertical fields inside and outside the cloud are the same then $`\gamma _d`$ is just the gravitational frequency for the intercloud medium, reduced by the ratio of vortical to total kinetic energy
$$\gamma _d=(4\pi G\rho _e)^{1/2}\frac{𝒞_0^2}{(𝒞_0^2+𝒟_0^2)}.$$
(4-2)
Since self gravity is presumably negligible in the low density intercloud medium, the energy loss rate is negligible as well. Loss of vorticity is measured by the magnetic braking rate $`\gamma _{mb}`$, which can be shown by a similar argument to be $`\gamma _{mb}=(4\pi G\rho _e)^{1/2}`$.
We thus see how the potential field limit is approached as the density contrast between the cloud and intercloud medium increases. At the late stages of clump formation the potential field becomes highly distorted and develops partially closed topology, but we can ignore this for the relatively mild density contrasts studied in this paper.
### 4.2 Approximations to the Gas Physics
The two dimensional approximation has a long and venerable history in galactic dynamics and accretion disc theory, as well as in studies of molecular clouds, and its errors for self gravitating systems are reasonably well understood. We expect the approximation to be reasonably good as long as the clump diameters exceed the disc thickness. The instability discussed in this paper has a 3D analog (Z98), but it must be treated by other means.
We assumed that the gas has an isothermal equation of state. This is a reasonable description of the kinetic pressure - density relationship, because of the high radiative efficiency of molecular gas. However, if the pressure were due to unresolved turbulence the medium would generally be less compressible; for example, Alfvén wave pressure follows density according to a 3/2 law (McKee & Zweibel, 1995). This would make the medium more stable by increasing the value of $`\gamma _T`$ (see eq. \[2-18\]), as would retention of magnetic pressure.
We took a uniform sheet at rest as an initial condition. This has the advantage of simplicity, but it means that there is no free energy stored in the background magnetic field. Thus, we have not tested the conjecture that the instability can convert magnetic free energy to turbulent energy, which was proposed in Z98.
We have treated ambipolar drift in the strong coupling approximation, and have implicitly assumed that $`v_D<`$ 20 km s<sup>-1</sup> (otherwise the rate coefficient would change). This is reasonable as long as the ion-neutral collision time, which is of order $`5\times 10^9n_n^1`$ s, is less than other timescales in the problem.
We have parameterized the relationship between the collision rate and the surface density by an exponent $`\alpha `$. In order to do better we would need a three dimensional model of the sheet and might need to follow the ionization as well. The results of linear theory are rather insensitive to the value of $`\alpha `$, which suggests that it need not be calculated very accurately in the present models.
## 5 Summary
The study of axisymmetric contraction of weakly ionized, self gravitating, magnetized clouds has proceeded quite far (Basu & Mouschovias, 1994; Ciolek & Mouschovias, 1993, 1994; Safier, McKee, & Stahler, 1997; Ciolek & Königl, 1998). In this paper, our emphasis has been on the initial breakup of a cold magnetized cloud gas into fragments and the early stages of their magnetic flux loss and contraction. We include self gravity, magnetic tension, and ambipolar drift, but we do not include detailed chemistry of grain physics, choosing instead a simple parameterization of the ionization. We follow the evolution for a shorter time than the isolated cloud collapse studies, but we impose no symmetry constraints or initial density or velocity structure.
We study highly flattened clouds, with an initially perpendicular magnetic field, which are slightly subcritical. The linear theory of collapse in such geometry (Z98, T=0; this work, T$`>`$0) predicts collapse on an intermediate timescale, faster than the diffusive timescale set by ambipolar drift, but slower than the dynamical timescale of free-falling inside-out collapse. The linear theory also predicts the existence of a single spatial wavenumber with maximal growth rate, with sufficiently short wavelengths stabilized by thermal pressure. This naturally suppresses power at short wavelengths, which is important for the success of the spectral method we employ in the simulations.
We simulate collapse in clouds with random initial density perturbations which grow from $`<`$0.1% of the mean density to 5-10 times the mean density. We confirm the intermediate collapse rate predicted by linear theory (§3.1), although the nonlinear collapse rate is faster than the linear rate. These intermediate rates are consistent with some recent observations of infall in molecular clouds (Evans, 1999; Myers, Evans, & Ohashi, 1999).
We show that clouds fragment into clumps with size corresponding to the wavelength of the spatial mode of maximal linear growth rate (§3.3), generally 1-10 $`M_{}`$. Collapse is asymmetric and complex (§3.4), and generally forms prolate clumps, for which there is observational evidence (Myers et al., 1991; Ryden, 1996; Ward-Thompson, Motte, & André, 1999). Sometimes the clumps are in mutual orbit, although the typical clump separation, a few tenths of a parsec, is too large to be relevant to the formation of binary stars. The magnetic field drives the growth of local vorticity, typically in the form of vortex pairs which straddle the clumps and are associated with streaming motions along them.
Considerable magnetic flux is lost from the collapsing clumps, consistent with the observationally determined $`|𝐁|\rho ^{0.5}`$ (Troland & Heiles, 1986; Crutcher, 1999) (see also §3.2). This flux loss is consistent with other calculations of cloud evolution (Fiedler & Mouschovias, 1993; Ciolek & Mouschovias, 1994) (although ambipolar drift is not necessary to bring about this relationship, either for isolated (Mouschovias, 1976) or turbulent, highly structured (Padoan & Norlund, 1999) clouds). As magnetic flux is lost and the surface density increases in the central regions of a contracting core, further fragmentation might ensue.
One prediction of the linear theory, namely that the instability could convert magnetic free energy to turbulence, has not been borne out by the simulations. This may be due to the fact that the initial magnetic field is completely uniform and therefore carries no free energy. Although significant magnetic curvature develops late in the runs, the cloud has already become quite dynamical. This prediction awaits future tests with a more stressed initial state. The relative motions of the clumps shown in Figure 7 are about an order of magnitude less than the relative motions of clumps separated by a few tenths of a parsec in real clouds (Goodman et al., 1998), although the infall velocities in the simulation are comparable to measured velocities (Myers, Evans, & Ohashi, 1999).
An interesting area of future work would be to extend this study to true three-dimensional clouds. The linear theory (Z98) indicates that this instability exists in three as well as two dimensions, and that the growth rate is still intermediate to slow diffusive contraction and fast dynamical collapse. Construction of a nonlinear model in three dimensions would be more difficult than the two dimensional models developed here, but could prove interesting. In this vein, we find the recent successful fit of observations of L1544 with a nearly critical model (Ciolek & Basu, 2000) encouraging.
We are happy to acknowledge support by NSF grant AST 9800616, a 3-year NSF Graduate Research Fellowship to R.I., and NASA grant NAG 5-4063 to the University of Colorado, as well as discussions with Neal Evans and comments by an anonymous referee.
## Appendix A Appendix
In this Appendix we drop the potential field approximation and calculate the response of the intercloud medium to motions within the cloud. This allows us to estimate the errors incurred by assuming a potential field.
We carry out the estimate using linearized, ideal MHD theory. This is more accurate in the intercloud medium than it would be in the disc, because the intercloud or external Alfvén speed $`v_{Ae}`$ is relatively high while ion-neutral friction is weak. We assume the equilibrium intercloud field is uniform and vertical ($`𝐁_\mathrm{𝟎}=\widehat{z}B_0`$). We choose to work in the half space $`z>0`$ (the results are similar in the other half space). In linear theory, the motions are purely horizontal, and we can derive the following pair of decoupled equations for the divergence $`𝒟`$ and vertical component $`𝒞`$ of the curl of the velocity
$`\left({\displaystyle \frac{^2}{t^2}}v_{Ae}^2^2\right)𝒟`$ $`=`$ $`0,`$ (A1a)
$`\left({\displaystyle \frac{^2}{t^2}}v_{Ae}^2{\displaystyle \frac{^2}{z^2}}\right)𝒞`$ $`=`$ $`0.`$ (A1b)
Equations (A1a) and (A1b) represent fast magnetosonic waves and Alfvén waves, respectively. In general, both types of waves are generated by the motions in the cloud.
In order to make progress, we assume plane wave horizontal behavior and exponential behavior in time, so that all perturbations depend on ($`x`$, $`y`$, $`t`$) as $`\mathrm{exp}(\gamma t+ik_xx+ik_yy)`$, where $`\gamma `$ may be complex: $`\gamma =i\omega +\nu `$ with both $`\omega `$, $`\nu >0`$. Then
$$𝒞=i(k_xv_yk_yv_x);𝒟=i(k_xv_x+k_yv_y)$$
(A2)
can be calculated at $`z=0`$ in terms of the motions on the disc. The vertical extensions of these quantities can be found from equations (A1a) and (A1b), choosing outward going or exponentially decaying wave solutions. For the Alfvénic part,
$$𝒞=𝒞_0e^{ik_Az};k_A\frac{\gamma }{v_{Ae}},$$
(A3)
where $`𝒞_0`$ is the value of $`𝒞`$ at $`z=0`$. For the magnetosonic part,
$$𝒟=𝒟_0e^{k_Mz};k_Mk_{}\left(1+\frac{\gamma ^2}{k_{}^2v_{Ae}^2}\right)^{1/2},$$
(A4)
where $`k_{}^2k_x^2+k_y^2`$. In equation (A3) we have written the vertical dependence as a propagating wave, and in equation (A4) as an evanescent wave. Although both $`k_A`$ and $`k_M`$ are complex, because $`\gamma `$ is complex, our notation reflects the salient aspects of their behavior. The magnetosonic wave is almost purely evanescent because the wave frequency is much less than the disc Alfvén frequency, which in turn is much less than the intercloud Alfvén frequency. The Alfvén wave has a substantial propagating component and decays in space as long as the disturbance is growing in time, which is purely a result of causality.
We now calculate the perturbed magnetic field components $`\delta 𝐁`$ at the disc. According to the linearized induction equation,
$`{\displaystyle \frac{\delta 𝐁_{}}{t}}`$ $`=`$ $`B_0{\displaystyle \frac{𝐯_{}}{z}},`$ (A5a)
$`{\displaystyle \frac{\delta B_z}{t}}`$ $`=`$ $`B_0𝒟.`$ (A5b)
Inverting the definitions of $`𝒞`$ and $`𝒟`$ for the velocity components gives
$$v_x=\frac{i}{k_{}^2}(k_y𝒞k_x𝒟);v_y=\frac{i}{k_{}^2}(k_x𝒞+k_y𝒟).$$
(A6)
Using equations (A3), (A4), and (A6) in equations (A5a) and (A5b) gives the field components at $`z=0`$
$`\delta B_x`$ $`=`$ $`{\displaystyle \frac{B_0}{k_{}^2\gamma }}(k_Ak_y𝒞_0+ik_Mk_x𝒟_0),`$ (A7a)
$`\delta B_y`$ $`=`$ $`{\displaystyle \frac{B_0}{k_{}^2\gamma }}(k_Ak_x𝒞_0+ik_Mk_y𝒟_0),`$ (A7b)
$`\delta B_z`$ $`=`$ $`{\displaystyle \frac{B_0}{\gamma }}𝒟_0.`$ (A7c)
We can use equations (A7a-A7c) to compare the MHD solution with the potential field limit. The ratios of the perturbed horizontal to vertical field components at $`z=0`$ can be written as
$$\frac{\delta 𝐁_{}}{\delta B_z}=i\widehat{k}_{}\left(1+\frac{\gamma ^2}{k_{}^2v_{Ae}^2}\right)^{1/2}+(\widehat{z}\times \widehat{k}_{})\frac{\gamma }{k_{}v_{Ae}}\frac{𝒞_0}{𝒟_0}.$$
(A8)
Equation (A8), together with equation (A3), shows that in the limit $`v_{Ae}\mathrm{}`$ the potential field solution is exact.
The Alfvénic part of the disturbance, as a propagating wave, removes both energy and angular momentum from the cloud. The energy flux $`_W`$ (accounting for both kinetic and electromagnetic energy, and for waves propagating in both directions away from the disc) is
$$_W=2\rho _e\frac{𝒞_0^2}{k_{}^2}v_{Ae}.$$
(A9)
In §4.1 we use equation (4-1) to derive the rate at which the perturbation in the disc is damped by outgoing waves. |
warning/0002/hep-ph0002132.html | ar5iv | text | # Soft-Gluon Resummation and PDF Theory Uncertainties
## 1 FACTORIZATION & THE NLO MODEL
A generic inclusive cross section for the process $`A+BF+X`$ with observed final-state system $`F`$, of total mass $`Q`$, can be expressed as
$`Q^4{\displaystyle \frac{d\sigma _{ABFX}}{dQ^2}}`$ $`=`$ $`\varphi _{a/A}(x_a,\mu ^2)\varphi _{b/B}(x_b,\mu ^2)`$ (1)
$`\widehat{\sigma }_{abFX}(z,Q,\mu ),`$
with $`z=Q^2/x_ax_bS`$. The $`\widehat{\sigma }_{ab}`$ are partonic hard-scattering functions, $`\widehat{\sigma }=\sigma _{\mathrm{Born}}+(\alpha _s(\mu ^2)/\pi )\widehat{\sigma }^{(1)}+\mathrm{}.`$ They are known to NLO for most processes in the standard model and its popular extensions. Corrections begin with higher, uncalculated orders in the hard scattering, which respect the form of Eq. (1). The discussion is simplified in terms of moments with respect to $`\tau =Q^2/S`$,
$`\stackrel{~}{\sigma }_{ABFX}`$ $`=`$ $`{\displaystyle _0^1}𝑑\tau \tau ^{N1}Q^4𝑑\sigma _{ABFX}/𝑑Q^2`$ (2)
$`={\displaystyle \underset{a,b}{}}\stackrel{~}{\varphi }_{a/A}(N,\mu ^2)\stackrel{~}{\sigma }_{abFX}(N,Q,\mu )\stackrel{~}{\varphi }_{b/B}(N,\mu ^2),`$
where the moments of the $`\varphi `$’s and $`\widehat{\sigma }_{abFX}`$ are defined similarly.
Eqs. (1) and (2) are starting-points for both the determination and the application of parton distribution functions (PDFs), $`\varphi _{i/H}`$, using 1-loop $`\widehat{\sigma }`$’s We may think of this collective enterprise as an “NLO model” for the PDFs, and for hadronic hard scattering in general. For precision applications we ask how well we really know the PDFs . Partly this is a question of how well data constrain them, and partly it is a question of how well we could know them, given finite-order calculations in Eqs. (1) and (2). We will not attempt here to assign error estimates to theory. We hope, however, to give a sense of how to distinguish ambiguity from uncertainty, and how our partial knowledge of higher orders can reduce the latter.
## 2 UNCERTAINTIES, SCHEMES & SCALES
It is not obvious how to quantify a “theoretical uncertainty”, since the idea seems to require us to estimate corrections that we haven’t yet calculated. We do not think an unequivocal definition is possible, but we can try at least to clarify the concept, by considering a hypothetical set of nucleon PDFs determined from DIS data alone . To make such a determination, we would invoke isospin symmetry to reduce the set of PDF’s to those of the proton, $`\varphi _{a/P}`$, and then measure a set of singlet and nonsinglet structure functions, which we denote $`F^{(i)}`$. Each factorized structure function may be written in moment space as
$$\stackrel{~}{F}^{(i)}(N,Q)=\underset{a}{}\stackrel{~}{C}_a^{(i)}(N,Q,\mu )\stackrel{~}{\varphi }_{a/P}(N,\mu ^2),$$
(3)
in terms of which we may solve for the parton distributions by inverting the matrix $`\stackrel{~}{C}`$,
$$\stackrel{~}{\varphi }_{a/P}(N,\mu ^2)=\underset{i}{}\stackrel{~}{C}^1{}_{a}{}^{(i)}(N,Q,\mu )\stackrel{~}{F}^{(i)}(N,Q).$$
(4)
With “perfect” $`\stackrel{~}{F}`$’s at fixed $`Q`$, and with a specific approximation for the coefficient functions, we could solve for the moment-space distributions numerically, without the need of a parameterization. In a world of perfect data, but of incompletely known coefficient functions, uncertainties in the parton distributions would be entirely due to the “theoretical” uncertainties of the $`C`$’s:
$$\delta \stackrel{~}{\varphi }_{a/P}(N,\mu )=\underset{i}{}\delta \stackrel{~}{C}^1{}_{a}{}^{(i)}(N,Q,\mu )\stackrel{~}{F}^{(i)}(N,Q).$$
(5)
Our question now becomes, how well do we know the $`C`$’s? In fact this is a subtle question, because the coefficient functions depend on choices of scheme and scale.
Factorization schemes are procedures for defining coefficient functions perturbatively. For example, choosing for $`F_2`$ the LO (quark) coefficient function in Eq. (4) defines a DIS scheme (with $`\stackrel{~}{C}`$ independent of $`\mu `$, which is then to be taken as $`Q`$ in $`\stackrel{~}{\varphi }`$). Computing the $`C`$’s from partonic cross sections by minimal subtraction to NLO defines an NLO $`\overline{\mathrm{MS}}`$ scheme, and so on. Once the choices of $`C`$’s and $`\mu `$ are made, the PDF’s are defined uniquely.
Evolution in an $`\overline{\mathrm{MS}}`$ or related scheme, enters through
$`\mu {\displaystyle \frac{d}{d\mu }}\stackrel{~}{\varphi }_{a/H}(N,\mu ^2)`$ $`=`$ $`\mathrm{\Gamma }_{ab}(N,\alpha _s(\mu ^2))\stackrel{~}{\varphi }_{b/H}(N,\mu ^2)`$
$`\mu {\displaystyle \frac{d}{d\mu }}\stackrel{~}{C}_c^{(i)}(N,Q,\mu )`$ $`=`$ $`\stackrel{~}{C}_d^{(i)}(N,Q,\mu )\mathrm{\Gamma }_{dc}(N,\alpha _s(\mu ^2)).`$ (6)
In principle, by Eq. (6), the scale-dependence of the $`C_a^{(i)}`$ exactly cancels that of the PDFs in Eq. (3) and, by extension, in Eq. (1). This cancelation, however, requires that each $`C`$ and the anomalous dimensions $`\mathrm{\Gamma }`$ be known to all orders in perturbation theory.
To eliminate $`\mu `$-dependence up to order $`\alpha _s^{n+1}`$, we need $`\widehat{\sigma }`$ to order $`\alpha _s^n`$ and the $`\mathrm{\Gamma }_{ab}`$ to $`\alpha _s^{n+1}`$. One-loop (NLO) QCD corrections to hard scattering require two-loop splitting functions, which are known. The complete form of the NNLO splitting functions, is still somewhere over the horizon . Even when these are known, it will take some time before more than a few hadronic hard scattering functions are known at NNLO.
We can clarify the role of higher orders by relating structure functions at two scales,$`Q_0`$ and $`Q`$. Once we have measured $`F(N,Q_0)`$, we may predict $`F(N,Q)`$ in terms of the relevant anomalous dimensions and coefficient functions by
$`F(N,Q)`$ $`=`$ $`F(N,Q_0)\mathrm{e}^{_{Q_0}^Q\frac{d\mu ^{}}{\mu ^{}}\mathrm{\Gamma }(N,\alpha _s(\mu ^{}{}_{}{}^{2}))}`$ (7)
$`\times \left[{\displaystyle \frac{\stackrel{~}{C}(N,Q,Q)}{\stackrel{~}{C}(N,Q_0,Q_0)}}\right].`$
This prediction, formally independent of PDFs and independent of the factorization scale, has corrections from the next, still uncalculated order in the anomalous dimension and in the ratio of coefficient functions. The asymptotic freedom of QCD gives a special role to LO: only the one-loop contribution to $`\mathrm{\Gamma }`$ diverges with $`Q`$ in the exponent, and contributes to the leading, logarithmic scale breaking. NLO corrections already decrease as the inverse of the logarithm of $`Q`$, NNLO as two powers of the log. Thus, the theory is self-regulating towards high energy, where dependence on uncalculated pieces in the coefficients and anomalous dimensions becomes less and less important.
The general successes of the NLO model strongly suggest that relations like (7) are well-satisfied for a wide range of observables and values of $`N`$ (or $`x`$) in DIS and other processes. This does not mean, however, that we have no knowledge of, or use for, information from higher orders. In particular, near $`x=1`$ PDFs are rather poorly known . At the same time, the ratio of $`C`$’s depends on $`N`$, and if $`\alpha _s\mathrm{ln}N`$ is large, it becomes important to control higher-order dependence on $`\mathrm{ln}N`$. This is a task usually referred to as resummation, to which we now turn.
## 3 RESUMMATION
Let us continue our discussion of DIS, describing what is known about the $`N`$-dependence of the coefficient functions $`C`$, as a step toward understanding the role of higher orders. Specializing again for simplicity to nonsinglet or valence, the resummed coefficient function may be written as
$$\stackrel{~}{C}^{\mathrm{res}}(N,Q,\mu )=\stackrel{~}{C}_{sub}^{\mathrm{NLO}}(N,Q,\mu )+C_\delta ^{\mathrm{DIS}}\mathrm{e}^{E_{\mathrm{DIS}}(N,Q,\mu )},$$
(8)
where “sub” implies a subtraction on $`\stackrel{~}{C}^{\mathrm{NLO}}`$ to keep $`\stackrel{~}{C}^{\mathrm{res}}`$ exact at order $`\alpha _s`$, and where $`C_\delta ^{\mathrm{DIS}}`$ corresponds to the NLO $`N`$-independent (“hard virtual”) terms. The exponent resums logarithms of $`N`$:
$`E_{\mathrm{DIS}}(N,Q,\mu )`$ $`=`$
$`{\displaystyle _{Q^2/\overline{N}}^{\mu ^2}}{\displaystyle \frac{d\mu ^{}^2}{\mu ^{}^2}}[A(\alpha _s(\mu ^{}{}_{}{}^{2}))\mathrm{ln}(\overline{N}\mu ^{}{}_{}{}^{2}/Q^2)+B(\alpha _s(\mu ^{}{}_{}{}^{2}))],`$
with $`\overline{N}N\mathrm{e}^{\gamma _E}`$, and with
$`A(\alpha _s)`$ $`=`$
$`{\displaystyle \frac{\alpha _s}{\pi }}C_F\left[1+{\displaystyle \frac{\alpha _s}{2\pi }}\left(C_A\left({\displaystyle \frac{67}{18}}{\displaystyle \frac{\pi ^2}{6}}\right){\displaystyle \frac{10}{9}}T_F\right)\right]`$
$`B(\alpha _s)`$ $`=`$ $`{\displaystyle \frac{3}{2}}C_F{\displaystyle \frac{\alpha _s}{2\pi }}.`$ (10)
Eq. (3) is accurate to leading (LL) and next-to-leading logarithms (NLL) in $`N`$ in the exponent: $`\alpha _s^m\mathrm{ln}^{m+1}N`$ and $`\alpha _s^m\mathrm{ln}^mN`$, respectively. The $`N`$ dependence of the ratio $`\stackrel{~}{C}_2^{\mathrm{res}}(N,Q,Q)/\stackrel{~}{C}_2^{\mathrm{NLO}}(N,Q,Q)`$ is shown in Fig. 1, with $`Q^2`$ = 1, 5, 10, 100 GeV<sup>2</sup>. At $`N=1`$ the ratio is unity. It is less than unity for moderate $`N`$, but then begins to rise, with a slope that increases strongly for small $`Q`$. At low $`Q^2`$ and large $`N`$, higher orders can be quite important. What does this mean for PDFs? We can certainly refit PDFs with resummed coefficient functions, and we see that the high moments of such PDFs are likely to be quite different from those from NLO fits.
To get a sense of how such an NLL/NLO-$`\overline{\mathrm{MS}}`$ scheme might differ from a classic NLO-$`\overline{\mathrm{MS}}`$ scheme, we resort to a model set of resummed distributions, determined as follows. We define valence PDFs in the resummed scheme by demanding that their contributions to $`F_2`$ match those of the corresponding NLO valence PDFs at a fixed $`Q=Q_0`$, which is ensured by
$$\stackrel{~}{\varphi }^{\mathrm{res}}(N,Q_0^2)=\stackrel{~}{\varphi }^{\mathrm{NLO}}(N,Q_0^2)\frac{\stackrel{~}{C}_2^{\mathrm{NLO}}(N,Q_0,Q_0)}{\stackrel{~}{C}_2^{\mathrm{res}}(N,Q_0,Q_0)}.$$
(11)
Using the resummed parton densities from Eq. (11), we can generate the ratios $`F_2^{\mathrm{res}}(x,Q)/F_2^{\mathrm{NLO}}(x,Q)`$.
The result of this test, picking $`Q_0^2=100`$ GeV<sup>2</sup> is shown in Fig. 2, for the valence $`F_2(x,Q)`$ of the proton, with $`x`$ = 0.55, 0.65, 0.75 and 0.85. The NLO distributions were those of , and the inversion of moments was performed as in . The effect of resummation is moderate for most $`Q`$. At small values of $`Q`$, and large $`x`$, the resummed structure function shows a rather sharp upturn. One also finds a gentle decrease toward very large $`Q`$ . We could interpret this difference as the uncertainty in the purely NLO valence PDFs implied by resummation.
From this simplified example, we can already see that the use of resummed coefficient functions is not likely to make drastic differences in global fits to PDFs based on DIS data, at least so long as the region of small $`Q^2`$, of 10 GeV<sup>2</sup> or below, is avoided at very large $`x`$. At the same time, it is clear that a resummed fit will make some difference at larger $`x`$, where PDFs are not so well known. We stress that a full global fit will be necessary for complete confidence.
## 4 RESUMMED HADRONIC SCATTERING
Processes other than DIS play an important role in global fits, and in any case are of paramount phenomenological interest. Potential sources of large corrections can be identified quite readily in Eq. (2). At higher orders, factors such as $`\alpha _s\mathrm{ln}^2N`$, can be as large as unity over the physically relevant range of $`z`$ in some processes. In this case, they, and their scale dependence can be competitive with NLO contributions. Since they make up well-defined parts of the correction at each higher order, however, it is possible to resum them. To better determine PDFs in regions of phase space where such corrections are important, we may incorporate resummation in the hard-scattering functions that determine PDFs.
The Drell-Yan cross section is the benchmark for the resummation of logs of $`1z`$, or equivalently, logarithms of the moment variable $`N`$ ,
$`\widehat{\sigma }_{q\overline{q}}^{\mathrm{DY}}(N,Q,\mu )`$ $`=`$ $`\sigma _{\mathrm{Born}}(Q)C_\delta ^{\mathrm{DY}}\mathrm{e}^{E_{\mathrm{DY}}(N,Q,\mu )}`$ (12)
$`+𝒪(1/N).`$
The exponent is given in the $`\overline{\mathrm{MS}}`$ scheme by
$`E_{\mathrm{DY}}(N,Q,\mu )`$ $`=`$ $`2{\displaystyle _{Q^2/\overline{N}^2}^{\mu ^2}}{\displaystyle \frac{d\mu ^{}^2}{\mu ^{}^2}}A(\alpha _s(\mu ^{}{}_{}{}^{2}))\mathrm{ln}\overline{N}`$ (13)
$`+2{\displaystyle _{Q^2/\overline{N}^2}^{Q^2}}{\displaystyle \frac{d\mu ^{}^2}{\mu ^{}^2}}A(\alpha _s(\mu ^{}{}_{}{}^{2}))\mathrm{ln}\left({\displaystyle \frac{\mu ^{}}{Q}}\right),`$
with $`A`$ as in Eq. (10), and where we have exhibited the dependence on the factorization scale, setting the renormalization scale to $`Q`$. Just as in Eq. (3) for DIS, Eq. (13) resums all leading and next-to-leading logarithms of $`N`$.
It has been noted in several phenomenological applications that threshold resummation, and even fixed-order expansions based upon it, significantly reduce sensitivity to the factorization scale . To see why, we rewrite the moments of the Drell-Yan cross section in resummed form as
$`\sigma _{AB}^{\mathrm{DY}}(N,Q)`$ (14)
$`={\displaystyle \underset{q}{}}\varphi _{q/A}(N,\mu )\widehat{\sigma }_{q\overline{q}}^{\mathrm{DY}}(N,Q,\mu )\varphi _{\overline{q}/B}(N,\mu )`$
$`={\displaystyle \underset{q}{}}\varphi _{q/A}(N,\mu )e^{E_{\mathrm{DY}}(N,Q,\mu )/2}\sigma _{\mathrm{Born}}(Q)C_\delta ^{\mathrm{DY}}`$
$`\times \varphi _{\overline{q}/B}(N,\mu )e^{E_{\mathrm{DY}}(N,Q,\mu )/2}+𝒪(1/N).`$
The exponentials compensate for the $`\mathrm{ln}N`$ part of the evolution of the parton distributions, and the $`\mu `$-dependence of the resummed expression is suppressed by a power of the moment variable,
$$\mu \frac{d}{d\mu }\left[\varphi _{q/A}(N,\mu )e^{E_{\mathrm{DY}}(N,Q,\mu )/2}\right]=𝒪(1/N).$$
(15)
This surprising relation holds because the function $`A(\alpha _s)`$ in Eq. (10) equals the residue of the $`1/(1x)`$ term in the splitting function $`P_{qq}`$. Thus, the remaining $`N`$-dependence in a resummed cross section still begins at order $`\alpha _s^2`$, but the part associated with the $`1/(1x)`$ term in the splitting functions has been canceled to all orders. Of course, the importance of the remaining sensitivity to $`\mu `$ depends on the kinematics and the process. In addition, although resummed cross sections can be made independent of $`\mu `$ for all $`\mathrm{ln}N`$, they are still uncertain at next-to-next-to leading logarithm in $`N`$, simply because we do not know the function $`A`$ at three loops. Notice that none of these results depends on using PDFs from a resummed scheme, because $`\overline{\mathrm{MS}}`$ PDFs, whether resummed or NLO, evolve the same way. The remaining, uncanceled dependence on the scales leaves room for an educated use of scale-setting arguments . The connection between resummation and the elimination of scale dependence has also been emphasized in .
Scale dependence aside, can we in good conscience combine resummed hard scattering functions in Eq. (1) with PDFs from an NLO scheme? This wouldn’t make much sense if resummation significantly changed the coefficient functions with which the PDFs were originally fit. As Fig. 2 shows, however, this is unlikely to be the case for DIS at moderate $`x`$. Thus, it makes sense to apply threshold resummation with NLO PDFs to processes and regions of phase space where there is reason to believe that logs are more important at higher orders than for the input data to the NLO fits.
At the same time, a set of fits that includes threshold resummation in their hard-scattering functions can be made , and their comparison to strict NLO fits would be quite interesting. Indeed, such a comparison would be a new measure of the influence of higher orders. A particularly interesting example might be to compare resummed and NLO fits using high-$`p_T`$ jet data .
## 5 POWER-SUPPRESSED CORRECTIONS
In addition to higher orders in $`\alpha _s(\mu ^2)`$, Eq. (1) has corrections that fall off as powers of the hard-scattering scale $`Q`$. In contrast to higher orders, these corrections require a generalization of the form of the factorized cross section. Often power corrections are parameterized as $`h(x)/[(1x)Q^2]`$ in inclusive DIS, where they begin at twist four. In DIS, this higher twist term influences PDFs when included in joint fits with the NLO and NNLO models, and vice-versa . As in the case with higher orders, such “power-improved” fits should be treated as new schemes.
## 6 CONCLUSIONS
The success of NLO fits to DIS and the studies of resummation above suggest that over most of the range of $`x`$, theoretical uncertainties of the NLO model are not severe. At the same time, to fit large $`x`$ with more confidence than is now possible may require including the resummed coefficient functions.
Resummation is especially desirable for global fits that employ a variety of processes, such as DIS and high-$`p_T`$ jet production, which differ in available phase space near partonic threshold. In a strictly NLO approach, uncalculated large corrections are automatically incorporated in the PDFs themselves. As a result, the NLO model cannot be expected to fit simultaneously the large-$`x`$ regions of processes with differing logs of $`1x`$ in their hard-scattering functions, unless these higher-order corrections are taken into account.
The results illustrated in the figures suggest that these considerations may be important in DIS with $`Q^2`$ below a few GeV<sup>2</sup> and at large $`x`$, where they may have substantial effects on estimates of higher twist in DIS. In hadronic scattering, large-$`N`$ ($`x1`$) resummation, which automatically reduces scale dependence, may play an even more important role than in DIS.
### Acknowledgments
We thank Andreas Vogt and Stephane Keller for useful discussions. |
warning/0002/cond-mat0002196.html | ar5iv | text | # Abrupt Change of Josephson Plasma Frequency at the Phase Boundary of the Bragg Glass in Bi2Sr2CaCu2O8+δ
\[
## Abstract
We report the first detailed and quantitative study of the Josephson coupling energy in the vortex liquid, Bragg glass and vortex glass phases of Bi<sub>2</sub>Sr<sub>2</sub>CaCu<sub>2</sub>O<sub>8+δ</sub> by the Josephson plasma resonance. The measurements revealed distinct features in the $`T`$\- and $`H`$-dependencies of the plasma frequency $`\omega _{pl}`$ for each of these three vortex phases. When going across either the Bragg-to-vortex glass or the Bragg-to-liquid transition line, $`\omega _{pl}`$ shows a dramatic change. We provide a quantitative discussion on the properties of these phase transitions, including the first order nature of the Bragg-to-vortex glass transition.
\]
The vortex matter in high-$`T_c`$ superconductors exhibits a fascinatingly rich phase diagram with a variety of phase transitions. There, thermal fluctuation and disorder alter dramatically the vortex phase diagram which has been observed in the conventional superconductors. At high temperature, the strong thermal fluctuation melts a vortex lattice into a vortex liquid well below the upper critical field. On the other hand, at low temperature or low field where the vortex liquid freezes into a solid phase, disorder plays an important role. The disorder is known to destroy the long-range order of the Abrikosov lattice . Recent investigations have revealed that the vortex solid phase is comprised of two distinct phases; a highly disordered phase at high field and a rather ordered phase at low field . The former phase is the vortex glass or entangled solid phase which is characterized by divergent barriers for vortex motion . The latter phase is the Bragg glass or quasilattice phase in which no dislocation exists and quasi-long-range translational order is preserved . In very clean single crystals, thermodynamical measurements have revealed that the Bragg glass undergoes a first order transition (FOT) to the vortex liquid . The transition from the Bragg glass to the vortex glass, on the other hand, is characterized by the second magnetization peak at which the critical current shows a sharp increase . It was proposed that the crossover from the FOT to the second peak regime is governed by a critical point $`T_{cp}`$ in the phase diagram, which in Bi<sub>2</sub>Sr<sub>2</sub>CaCu<sub>2</sub>O<sub>8+δ</sub> is located near 40 K. While the nature of the vortex liquid has been extensively studied, the properties of the Bragg glass and the nature of the thermally induced FOT from the Bragg glass to the vortex liquid are still not quite understood. Moreover, the phase transition from the Bragg glass to the vortex glass at lower temperatures has been a longstanding issue, though this transition is proposed to be disorder driven, caused by competition between the elastic and pinning energies . A major obstacle has been that most of the previous experiments had to been performed under a strongly nonequilibrium condition because most part of the Bragg and vortex glasses are located deep inside the irreversibility line $`T_{irr}`$.
The most direct way to clarify the nature of these phases and the phase transitions among them is to measure the interlayer phase coherence for each vortex phases, because the CuO<sub>2</sub> layers are connected by the Josephson effect. One of the most powerful probes for the interlayer phase coherence is the Josephson plasma resonance (JPR) which provides a direct measurement of the Josephson plasma frequency $`\omega _{pl}`$ related to the maximum Josephson critical current $`j_J=\epsilon _0\mathrm{\Phi }_0\omega _{pl}^2/8\pi ^2cd`$ and the Josephson coupling energy $`U_J=\mathrm{\Phi }_0j_J/2\pi c`$, where $`\epsilon _0`$ and $`d`$ are the dielectric constant and interlayer spacing, respectively. . Especially in Bi<sub>2</sub>Sr<sub>2</sub>CaCu<sub>2</sub>O<sub>8+δ</sub> with large anisotropy, a very precise determination of $`\omega _{pl}`$ is possible because $`\omega _{pl}`$ falls within the microwave window.
All of the JPR measurements of Bi<sub>2</sub>Sr<sub>2</sub>CaCu<sub>2</sub>O<sub>8+δ</sub> up to now have been carried out in the cavity resonator by reducing $`\omega _{pl}`$ by $`H`$ . Unfortunately, sweeping $`H`$ below $`T_{irr}`$ drives the vortex system into a strongly nonequilibrium state due to the Bean critical current induced by the field gradient inside the crystal, as was demonstrated in Refs. and . Therefore, in order to investigate the Bragg and vortex glass phases, it is crucial to measure the JPR as a function of the microwave frequency $`\omega _{pl}`$ while holding $`H`$ at a constant value. In this Letter, we report the first detailed and quantitative study of the Josephson coupling energy in the Bragg glass, the vortex glass and the vortex liquid phases and the phase transitions among them by the JPR which has been preformed by sweeping $`\omega `$ continuously. The measurements revealed distinct features in the $`T`$\- and $`H`$-dependencies of $`\omega _{pl}`$ for each of the three different vortex phases. When going across either the Bragg-to-vortex glass
or the Bragg-to-liquid transition line, $`\omega _{pl}`$ shows a dramatic change. We provide a quantitative discussion on the nature of these phase transitions in the light of these results.
All experiments were performed on a slightly underdoped Bi<sub>2</sub>Sr<sub>2</sub>CaCu<sub>2</sub>O<sub>8+δ</sub> single crystals ($`T_c`$=82.5 K) with dimensions $`1.2\times 0.5\times 0.03`$mm<sup>3</sup> grown by the traveling floating zone method. The inset of Fig.1 shows a typical magnetization step measured by SQUID magnetometer which can be attributed to the FOT of the vortex lattice. This FOT terminates at $`40`$ K and the step is followed by the second magnetization peak located at $``$230 Oe. Figure 1 shows the phase diagram obtained by the magnetization measurements. The JPR is measured by sweeping $`\omega `$ continuously from 20 GHz to 150 GHz . The sample was placed at the center of the broad wall of the waveguide in the traveling wave TE<sub>01</sub> mode. We used a bolometric technique to detect very small microwave absorption by the sample and employed a leveling loop technique to ckeep the microwave power constant when sweeping frequency. For this crystal $`\omega _{pl}`$=125 GHz at $`T`$=0, corresponding to the anisotropy parameter $`\gamma =\lambda _c/\lambda _{ab}550`$, where $`\lambda _{ab}`$ and $`\lambda _c`$ are the in-plane and out-of-plane penetration lengths, respectively. Here we used $`\lambda _{ab}200`$ nm and $`\lambda _c=c/\omega _{pl}\sqrt{\epsilon _0}`$110 $`\mu `$m. We determined $`\epsilon _0=11.5\pm 1`$ from the dispersion of the transverse plasma mode. All
JPR measurements were performed in H$`c`$ under the field cooling condition (FCC) where the field is very uniform. In this condition, the system is in equilibrium or at worst is trapped in a metastable state which we expect should be much closer to equilibrium compared to the state obtained in the field sweeping condition (FSC). In fact, while the resonance frequency below $`T_{irr}`$ did not change at all with time for more than 48 hours in the FCC, it increases gradually with time in the FSC. We also confirmed that the resonance curves are exactly the same in different cooling cycles.
Figures 2(a) and (b) depict the resonant absorption as a function of $`\omega `$ when crossing the second peak field $`H_{sp}`$ and the FOT field $`H_m`$, respectively. When $`\omega `$ coincides with $`\omega _{pl}`$, the resonant absorption of the microwave occurs. These are the JPR measured in the Bragg and vortex glass phases under the FCC for the first time. In the magnetic field, $`\omega _{pl}`$ can be written as
$$\omega _{pl}^2(B,T)=\omega _{pl}^2(0,T)\mathrm{cos}\varphi _{n,n+1}.$$
(1)
Here $`\mathrm{cos}\varphi _{n,n+1}`$ represents the thermal and disorder average of the cosine of the gauge invariant phase difference between layer $`n`$ and $`n+1`$. If the vortex forms a straight line along the $`c`$-axis, $`\mathrm{cos}\varphi _{n,n+1}`$ is unity. The reduction of $`\mathrm{cos}\varphi _{n,n+1}`$ from unity is caused by the Josephson strings that are created by the deviation from the straight alignment of the pancake vortices along the $`c`$-axis. Thus $`\omega _{pl}`$ gives a direct information on the vortex alignment and therefore the phase transition of the vortex matter. After gradual decrease with $`H`$ at lower $`H`$, $`\omega _{pl}`$ shows a sharp decrease in the field range between 215 Oe and 220 Oe at 6.5 K and between 140 0e and
160 Oe at 50 K. At 217.5 Oe in Fig.2(a) and at 150 Oe in Fig.2(b), the resonance lines become broader, indicating a very rapid change of $`\omega _{pl}`$ with $`H`$. At higher $`H`$, $`\omega _{pl}`$ again decreases gradually. In Fig.1, we plot the fields at which $`\omega _{pl}`$ shows an abrupt change. These fields coincide well with the second peak and FOT fields determined by magnetization measurements.
We first discuss the resonance when going across the FOT. The inset of Fig.3 depicts the $`H`$-dependence of $`\mathrm{cos}\varphi _{n,n+1}`$ obtained from $`\omega _{pl}^2(B,T)/\omega _{pl}^2(0,T)`$. Although similar results have been reported , quantitative analysis was very difficult because the JPR measurements in the Bragg glass had been done under the strongly nonequilibrium condition, as we have already mentioned. Figure 3 depicts $`\mathrm{cos}\varphi _{n,n+1}`$ as a function of $`H`$ normalized by $`H_m`$. Interestingly, $`\mathrm{cos}\varphi _{n,n+1}`$ exhibits very similar $`H/H_m`$-dependence at all temperatures. Obviously, the $`H`$-dependence of $`\mathrm{cos}\varphi _{n,n+1}`$ above FOT is very different from that below FOT; the curvature changes from negative to positive. We found that $`\mathrm{cos}\varphi _{n,n+1}`$ in the Bragg glass phase can be fitted as,
$$\mathrm{cos}\varphi _{n,n+1}=1A_1\frac{H}{H_m}A_2\left(\frac{H}{H_m}\right)^2,$$
(2)
with $`A_1`$=0.16 and $`A_2`$=0.19 above 40 K as shown in the dashed line in Fig.3. On the other hand, according to high temperature expansion theory , $`\mathrm{cos}\varphi _{n,n+1}`$ in the liquid phase above FOT can be written as,
$$\mathrm{cos}\varphi _{n,n+1}=\frac{U_J\mathrm{\Phi }_0}{2k_BTH},$$
(3)
when the Josephson energy is negligible compared with the energy of thermal fluctuations, i.e.$`U_Jk_BTH/\mathrm{\Phi }_0`$. It has been shown experimentally that $`\mathrm{cos}\varphi _{n,n+1}`$ is inversely proportional to $`H`$ in the liquid phase . The present results provide a further rigorous test to Eq.3, because we now have no ambiguous fitting parameter and also have the data of the very detailed $`H`$-dependence of $`\mathrm{cos}\varphi _{n,n+1}`$ obtained by sweeping $`\omega `$. The solid line in Fig.3 shows the result of the calculation. In the calculation, we used $`\epsilon _0`$=12.0. The fit to the data is excellent in the whole $`H`$-range at $`H>1.2H_m`$, indicating that the vortex liquid is decoupled on the scale of the interlayer distance. Small deviation from Eq.3 is observed at $`H1.2H_m`$. This suggests that the vortex-vortex correlation effect in the $`ab`$-plane which gives rise to the deviation from $`1/H`$-dependence of $`\mathrm{cos}\varphi _{n,n+1}`$ in the liquid phase is important just above the FOT .
At $`H_m`$, $`\mathrm{cos}\varphi _{n,n+1}`$ is reduced to $``$0.7 at all temperatures, showing an occurrence of large vortex wandering in the Bragg glass. Near $`T_c`$, we note that $`\omega _{pl}`$ at $`H`$=0 is already suppressed by the phase fluctuations. If this effect is taken into account, it is expected that $`\mathrm{cos}\varphi _{n,n+1}`$ at $`H_m`$ slowly increases with $`T`$, indicating that the melting becomes more linelike at higher $`T`$. The values of $`\mathrm{cos}\varphi _{n,n+1}`$ at $`H_m`$ are close to the recent results of computer simulations for systems with small anisotropies.
The internal energy $`U`$ experiences a jump $`\mathrm{\Delta }U`$ at the FOT. This latent heat $`\mathrm{\Delta }U`$ can be represented as a sum of the jumps in the in-plane energy, in the electromagnetic coupling energy, and in the Josephson energy $`\mathrm{\Delta }U_J`$ . To understand the nature of the FOT in detail, it is important to establish the relative jump in Josephson energy $`\mathrm{\Delta }U_J/\mathrm{\Delta }U`$. At 60 K, $`\mathrm{\Delta }U_J/T`$ can be estimated to be $`0.21k_B`$ from $`\mathrm{cos}\varphi _{n,n+1}`$ which drops approximately from 0.70 to 0.45. On the other hand, $`\mathrm{\Delta }U/T`$ at 60 K obtained from the magnetization step $`\mathrm{\Delta }M`$ using the Clausius-Clapeyron relation,
$$\mathrm{\Delta }U/T_m=\mathrm{\Delta }S=d\mathrm{\Phi }_0\frac{\mathrm{\Delta }M}{B_m}\frac{dB_m}{dT},$$
(4)
is $`1.34k_B`$. Here $`\mathrm{\Delta }S`$ is the entropy jump at the FOT point ($`T_m`$, $`B_m`$). Thus we find that $`\mathrm{\Delta }U_J`$ constitutes approximately 16% of the latent heat, showing that $`\mathrm{\Delta }U_J`$ occupies a substantial part of the latent heat at the FOT even in Bi<sub>2</sub>Sr<sub>2</sub>CaCu<sub>2</sub>O<sub>8+δ</sub> with very large anisotropy.
We now move on to the subject of the Josephson coupling at low temperatures when going across the transition from the Bragg glass to the vortex glass. Figure 4 shows the $`H`$-dependence of $`\mathrm{cos}\varphi _{n,n+1}`$ below 35 K. Below 100 Oe, $`\mathrm{cos}\varphi _{n,n+1}`$ shows a hump structure which may be related with the lower critical field. Above 100 Oe, the $`H`$-dependence of $`\mathrm{cos}\varphi _{n,n+1}`$ is very similar to that at high $`T`$ when crossing the FOT. At all temperatures, $`\mathrm{cos}\varphi _{n,n+1}`$ shows an abrupt change at the second peak field $`H_{sp}`$220 Oe. In similarity to the high temperature behavior, the $`H`$-dependence of $`\mathrm{cos}\varphi _{n,n+1}`$ below and above $`H_{sp}`$ are very different, showing clearly that $`H_{sp}`$ separates two distinct vortex phases. In a very narrow field interval less than 5 Oe at $`H_{sp}`$, $`\mathrm{cos}\varphi _{n,n+1}`$ drops from approximately 0.7 to 0.5 (see also Fig.1(a)), corresponding to a nearly 20%-reduction of $`U_J`$. This strong reduction of $`\mathrm{cos}\varphi _{n,n+1}`$ provides a direct evidence of the decoupling nature of the Bragg-to-vortex glass transition . At $`H_{sp}`$, $`\mathrm{cos}\varphi _{n,n+1}`$ is reduced to $``$0.7 from the zero field value similar to that below FOT. Although we do not show here, $`\mathrm{cos}\varphi _{n,n+1}`$ in the vortex glass phase deviates from the $`1/H`$ dependence in the whole $`B`$-regime above $`H_{sp}`$, which is to be contrasted to the behavior in the vortex liquid phase.
We finally discuss the phase transition from the Bragg glass to the vortex glass inferred from the JPR. The first question is the order of the transition. The abrupt change of $`\mathrm{cos}\varphi _{n,n+1}`$ shown in Figs.3 and 4 provides a direct evidence of the abrupt changes of the $`c`$-axis correlation length of the pancakes and of $`U_J`$ which composes a substantial part in the free energy. In Fig.3 we plot the change of $`\mathrm{cos}\varphi _{n,n+1}`$ at $`H_{sp}`$ ($`T`$=30 K), for the comparison with the change of the same quantity at the FOT. The change of $`\mathrm{cos}\varphi _{n,n+1}`$ at $`H_{sp}`$ is comparable or even sharper than that at the FOT. This fact strongly indicates the first order nature of the phase transition from the Bragg glass to the vortex glass. The second issue is the critical point $`T_{cp}`$ of the FOT which has been proposed to terminate at $``$40 K . This proposal was made from the observation that $`\mathrm{\Delta }S`$ becomes extremely small which can be seen from the $`T`$-independence of FOT line below $`T_{cp}`$. However, the vanishing of $`\mathrm{\Delta }S`$ does not immediately imply the termination of the FOT, which suggests that the issue of the termination is nontrivial. As seen in Figs.3 and 4, there is no discernible difference in the $`H`$-dependence of $`\mathrm{cos}\varphi _{n,n+1}`$ as we go through the Bragg-to-liquid transition regime, into the Bragg-to-vortex glass transition regime, except for a gradual decrease of the change of $`\mathrm{cos}\varphi _{n,n+1}`$ at the transition. These results imply that the FOT does not terminate at $``$40 K, but that there is no critical point or the FOT persists at least below 6.4 K. We note that a similar conclusion has been reached very recently using the magneto-optical imaging technique .
In summary, we have performed the JPR experiments in the Bragg glass, vortex glass, and vortex liquid phases in the FCC. We found an abrupt change in the Josephson coupling energy when going across either the FOT line or the second magnetization peak line. We showed that this change occupies a substantial part of the latent heat at the FOT. The results suggest that the Bragg-to-vortex glass transition is first order in nature and that the critical point of the FOT does not terminate at $``$40 K.
We thank B. Horovitz, X. Hu, Y. Kato, P.H. Kes, T. Onogi, A. Sudbø, and A. Tanaka for discussions. We are indebted to L.N. Bulaevskii for several valuable comments. We also thank A.E.Koshelev for the critical reading of the manuscript. |
warning/0002/astro-ph0002041.html | ar5iv | text | # Cosmic Ray Rejection by Linear Filtering of Single Images1footnote 11footnote 1 Accepted for publication in the May 2000 issue of the Publications of the Astronomical Society of the Pacific.
## 1 Introduction
Images from most current-day astronomical instruments have tractable noise properties. An exemplary case is optical images from CCD detectors, whose uncertainties are generally dominated by the Poisson statistics of the detected photons, with (usually smaller) contributions from detector read noise, dark current, and other comparatively minor nuisances. Most of these noise sources are well approximated by Gaussian distributions, and their sum is therefore also well approximated by a Gaussian.
Cosmic rays impinging on a detector can yield large signals over single pixels or small groups of pixels, thereby introducing a distinctly non-Gaussian tail to the noise distribution. The most common approach to removing cosmic rays from astronomical images is to take multiple exposures and combine them with some sort of outlier rejection. Real astronomical objects should (usually) be present on multiple frames, while cosmic ray hits will not generally repeat. Such methods have been presented in the literature by (e.g.) Shaw & Horne (1992) and Windhorst, Franklin, & Neuschaefer (1994), and are widely implemented in astronomical image processing packages.
However, there are times when multiple images are not available, or when the sources of interest may be moving or varying on timescales short compared to the interval between exposures. In these cases, a cosmic ray rejection method capable of operating on single exposures is necessary. Cosmic ray rejection in single frames can also be useful even when multiple exposures are to be stacked, since stacking often requires spatial interpolation of the input images, and any cosmic rays not previously identified can be spread over many pixels by spatially extended interpolation kernels. Additionally, if a stack of images has widely different point spread function (PSF) widths, rejection algorithms used while stacking tend either to be overly lenient, potentially admitting cosmic rays; or overly strict, discarding valid data from images with very good or very bad seeing. Examples of both these behaviors are offered by sigma clipping algorithms, where the contribution of a particular exposure to a stack is discarded if it differs from the mean (or median) intensity at that location by more than $`k\sigma `$, where $`k`$ is a constant (generally with $`2k5`$) and $`\sigma `$ measures the intensity uncertainty at that location. If $`\sigma `$ is measured directly from the list of exposure intensities at a fixed sky position, a lenient rejection results, while if $`\sigma `$ is taken from the known Poisson statistics of electrons in single exposures, a strict rejection results.
To identify cosmic rays in single exposures, rejection algorithms rely on the sharpness of cosmic rays relative to true astronomical objects. That is, any legitimate object in our astronomical image is blurred by the PSF, but there is no such requirement on cosmic ray hits. Provided the image is well-sampled (in practice, $`2`$ pixels across the PSF full width at half maximum), cosmic ray hits can be identified as those features with spatial variations too rapid for consistency with the PSF. Murtagh (1992) and Salzberg et al (1995) have explored trainable classifier approaches to single-image cosmic ray rejection. Their methods have the advantage of applicability to substantially undersampled data (from the WF/PC-I instrument on the Hubble Space Telescope). On the other hand, these methods ultimately rely on a training set, which may be subjectively defined.
The present paper explores a method suggested by Fischer and Kochanski (1994), who remark that the optimal filter for detecting \[single-pixel\] cosmic ray hits is the point spread function minus a delta function. This can be regarded as a difference between the matched filter for detecting point sources (i.e. the PSF) and that for detecting single pixels (i.e. a delta function). There is one free parameter in such a filter, which is the amplitude ratio of the two functions. We develop this filtering method in detail by considering the cosmic ray rejection rates and false alarm rates. Much of our analysis is devoted to choosing the delta function amplitude appropriately. With a careful choice of this parameter, it is possible to ensure that the false alarm rate nowhere exceeds its value in blank sky regions.
In section 2, we derive the noise properties of our filtered image, and explain how to tune the filter to avoid excessive rejection of valid data. In section 3, we discuss practical issues that arise when implementing our algorithm. Section 4 presents simulations used to verify the algorithm’s performance. Finally, in section 5 we summarize our work, describe our usual application for our algorithm, and comment on a desirable future direction for cosmic ray rejection algorithms.
## 2 Mathematical formalism
Suppose we have an image $`I`$ with the following properties: First, it has some background level $`B_I`$ and noise $`\sigma _I`$, and the sky noise is uncorrelated between any pair of pixels. Second, it is linear in the input signal with a gain $`g`$ photons per count, so that a pixel containing object flux $`S`$ counts will have a noise contribution of $`\sqrt{S/g}`$ counts from Poisson noise in the object signal. Third, it has a point spread function that can be well approximated by a Gaussian of characteristic width $`\xi `$ (i.e., the stellar profiles have a functional form $`\mathrm{exp}[(\mathrm{\Delta }x^2+\mathrm{\Delta }y^2)/(2\xi ^2)]`$), and is well sampled (i.e., $`\xi 1`$ pixel). This third property is an analytical convenience that is reasonably near truth for seeing-limited optical images from ground-based telescopes. Other PSF models would complicate the mathematical analysis that follows, but would not greatly change either its flavor or its quantitative results.
Now consider convolving this image with a spatial filter $`F`$ consisting of a unit-normalized point spread function $`A=\mathrm{exp}[(\mathrm{\Delta }x^2+\mathrm{\Delta }y^2)/(2\xi ^2)]/(2\pi \xi ^2)`$ minus a scaled delta function: $`F=A\alpha \delta (\mathrm{\Delta }x)\delta (\mathrm{\Delta }y)`$. Call the convolved image $`J`$, so that $`J=IF=IA\alpha I`$. (Here and throughout the paper, “$``$” is the convolution operator.)
If we regard the convolution kernel as a matched filter, it is clear that a broader kernel (likely using a functional form besides the Gaussian) would be more effective at separating cosmic rays from faint galaxies or other extended sources. However, almost all astronomical images contain some legitimate pointlike sources, which should not be rejected. Using a template more extended than a point source would risk rejecting stars, and such templates are therefore not explored further.
### 2.1 Noise properties of the filtered image
We calculate the noise in the convolved image in two steps, first determining the noise in $`IA`$ and then modifying the result to account for the second term in filter $`F`$. Treating the noise in each pixel as an independent random variable with variance $`\sigma _I^2`$, the variance in the convolved image is simply a weighted sum $`\sigma _J^2=_{\mathrm{}}w_{\mathrm{}}^2\sigma _{I,\mathrm{}}^2`$, where the sum runs over pixels and $`w_{\mathrm{}}`$ is simply $`F`$ evaluated at the location $`(\mathrm{\Delta }x_{\mathrm{}},\mathrm{\Delta }y_{\mathrm{}})`$ of pixel $`\mathrm{}`$. Now, in regions of blank sky, $`\sigma _{I,\mathrm{}}\sigma _I`$ is constant, so $`\sigma _{IF}=\sigma _I^2_{\mathrm{}}w_{\mathrm{}}^2`$.
We can calculate the noise level in $`IA`$ by defining weights $`v_{\mathrm{}}`$ as $`A`$ evaluated at $`(\mathrm{\Delta }x_{\mathrm{}},\mathrm{\Delta }y_{\mathrm{}})`$, and noting that
$$\frac{\sigma _{IA}^2}{\sigma _I^2}=\underset{\mathrm{}}{}v_{\mathrm{}}^2_0^{\mathrm{}}\frac{2\pi rdr}{(2\pi \xi ^2)^2}\left[\mathrm{exp}\left(\frac{r^2}{2\xi ^2}\right)\right]^2=\frac{1}{4\pi \xi ^2},$$
(1)
so that $`\sigma _{IA}^2=\sigma _I^2/(4\pi \xi ^2)`$. The continuous approximation to the discrete sum made here should be reasonably accurate for well-sampled data.
Modifying this for the central pixel, which has weight $`w=1/(2\pi \xi ^2)\alpha `$ rather than $`v=1/(2\pi \xi ^2)`$ as used above, we find
$$\frac{\sigma _J^2}{\sigma _I^2}=\underset{\mathrm{}}{}w_{\mathrm{}}^2\frac{1}{4\pi \xi ^2}\left(\frac{1}{2\pi \xi ^2}\right)^2+\left(\frac{1}{2\pi \xi ^2}\alpha \right)^2=\frac{1}{4\pi \xi ^2}\frac{\alpha }{\pi \xi ^2}+\alpha ^2.$$
(2)
That accomplished, we can determine the significance level that a cosmic ray with amplitude $`n\sigma _I`$ will have in image $`J`$. A single pixel cosmic ray with $`C_I`$ counts will result in a pixel with expectation value $`C_J=\left[1/(2\pi \xi ^2)\alpha \right]C_I`$ below the sky level of $`J`$ (which is $`B_J=(1\alpha )B_I`$). If $`C_I=n\sigma _I`$, then the final significance level is
$$\frac{C_J}{\sigma _J}=\frac{C_I}{\sigma _I}\times \frac{1/(2\pi \xi ^2)\alpha }{\left(1/(4\pi \xi ^2)\alpha /(\pi \xi ^2)+\alpha ^2\right)^{1/2}}=\frac{C_I}{\sigma _I}\times \left[1+\frac{\pi \xi ^21}{(2\pi \xi ^2\alpha 1)^2}\right]^{1/2}.$$
(3)
In general, this is a lower significance level than in the original image. In the limit of very well sampled data ($`\xi \mathrm{}`$) this reduces to a significance level of $`C_I/\sigma _I`$, recovering the input as one might expect. The gradual approach to this limit simply reflects the dependence of cosmic ray identification on the sampling of an image. Figure 1 shows contours of $`C_J\sigma _I/(C_I\sigma _J)`$.
When we reject cosmic rays, we need to be careful not to reject the cores of legitimate point sources. In order to avoid doing so, we calculate the noise level at a location near a point source, making the same continuous approximation to discrete sums used in deriving equation 1. The complication arising in this procedure is that the presence of a source changes the noise properties of the image. A pixel containing object flux $`I_{\mathrm{}}`$ has noise level given by $`\sigma _{I,\mathrm{}}^2=\sigma _I^2+I_{\mathrm{}}/g`$.
We consider below the noise at a location $`q`$ pixels from the location of a star with peak counts $`S_0`$, and define function $`S_I(q)=S_0\mathrm{exp}\left[q^2/(2\xi ^2)\right]`$. In the convolved image, this pixel has expected flux
$$S_J(q)=S_0\left(\frac{1}{2}\mathrm{exp}\left[\frac{q^2}{4\xi ^2}\right]\alpha \mathrm{exp}\left[\frac{q^2}{2\xi ^2}\right]\right).$$
(4)
Since variances add linearly, modifying our earlier analysis for the additional noise term is relatively straightforward. For image $`IA`$, we find
$`\sigma _{AI}^2(q)={\displaystyle \underset{\mathrm{}}{}}v_{\mathrm{}}^2\sigma _{I,\mathrm{}}^2`$ (5)
$`=`$ $`{\displaystyle \underset{\mathrm{\Delta }x}{}}{\displaystyle \underset{\mathrm{\Delta }y}{}}\left({\displaystyle \frac{1}{2\pi \xi ^2}}\mathrm{exp}\left[{\displaystyle \frac{(\mathrm{\Delta }x^2+\mathrm{\Delta }y^2)}{2\xi ^2}}\right]\right)^2\left(\sigma _I^2+{\displaystyle \frac{S_0}{g}}\mathrm{exp}\left[{\displaystyle \frac{([\mathrm{\Delta }xq]^2+\mathrm{\Delta }y^2)}{2\xi ^2}}\right]\right)`$ (6)
where we have assumed (without loss of generality) that the offset to the star is along the $`x`$-axis. Defining $`\mathrm{\Delta }x^{}=\mathrm{\Delta }xq/3`$, and substituting our result from equation 1, this becomes
$`\sigma _{AI}^2(q)`$ $`=`$ $`{\displaystyle \frac{\sigma _I^2}{4\pi \xi ^2}}+{\displaystyle \frac{S_0}{(2\pi \xi ^2)^2g}}\mathrm{exp}\left[{\displaystyle \frac{q^2}{3\xi ^2}}\right]{\displaystyle \underset{\mathrm{\Delta }x}{}}{\displaystyle \underset{\mathrm{\Delta }y}{}}\mathrm{exp}\left[{\displaystyle \frac{3(\mathrm{\Delta }x^2+\mathrm{\Delta }y^2)}{2\xi ^2}}\right]`$ (7)
$``$ $`{\displaystyle \frac{\sigma _I^2}{4\pi \xi ^2}}+{\displaystyle \frac{S_0}{(2\pi \xi ^2)^2g}}\mathrm{exp}\left[{\displaystyle \frac{q^2}{3\xi ^2}}\right]{\displaystyle _0^{\mathrm{}}}2\pi r\mathrm{exp}\left[{\displaystyle \frac{3r^2}{2\xi ^2}}\right]𝑑r`$ (8)
$`=`$ $`{\displaystyle \frac{\sigma _I^2}{4\pi \xi ^2}}+\mathrm{exp}\left[{\displaystyle \frac{q^2}{3\xi ^2}}\right]{\displaystyle \frac{S_0/g}{6\pi \xi ^2}}.`$ (9)
Again modifying the result to account for the delta function in the convolution kernel, we obtain for the noise at distance $`q`$ from the point source
$`\sigma _J^2(q)=\sigma _{AI}^2+\left(\sigma _I^2+S_I(q)/g\right)\left(\left[{\displaystyle \frac{1}{2\pi \xi ^2}}\alpha \right]^2\left[{\displaystyle \frac{1}{2\pi \xi ^2}}\right]^2\right)`$ (10)
$`=`$ $`\sigma _I^2\left({\displaystyle \frac{1}{4\pi \xi ^2}}{\displaystyle \frac{\alpha }{\pi \xi ^2}}+\alpha ^2\right)+{\displaystyle \frac{S_0}{g}}\mathrm{exp}\left[{\displaystyle \frac{q^2}{2\xi ^2}}\right]\left({\displaystyle \frac{\mathrm{exp}\left[+q^2/(6\xi ^2)\right]}{6\pi \xi ^2}}{\displaystyle \frac{\alpha }{\pi \xi ^2}}+\alpha ^2\right).`$ (11)
In the limit $`q\mathrm{}`$, this expression reproduces our blank sky result (equation 2), while at the peak of the star, it simplifies somewhat as the exponential terms go to unity.
These results can easily be generalized to a superposition of point sources; the second terms on the right hand sides of equations 4 and 11 would simply be replaced by a sum of such terms, each with its own value of the intensity parameter $`S_0`$ and distance parameter $`q`$.
### 2.2 Keeping the valid peaks
To identify cosmic rays in our image, we plan to threshold the convolved image $`J`$, flagging all pixels with excessively negative values in $`J`$. There are two probabilities of interest here, namely the probability that we will correctly flag a cosmic ray with intensity $`C_I`$ (the detection rate), and the probability that we will incorrectly flag a pixel without cosmic ray flux (the false alarm rate). We have chosen to concentrate our efforts on controlling the false alarm rate, and to accept the resulting detection rate. From a hypothesis testing perspective (e.g., Kendall & Stuart 1967, chapter 22), this approach corresponds to making the null hypothesis that a given pixel is uncontaminated by cosmic ray flux. The false alarm rate is then the probability of a type I error. Missed cosmic ray events are type II errors, and their probability can be calculated as a function of cosmic ray intensity. The tradeoffs between these two errors for a variety of cosmic ray rejection algorithms are reviewed by Murtagh & Adorf (1991).
We can never set the probability of rejecting valid pixels to be precisely zero so long as we have noise in our image and we reject any pixels at all. Instead we note that there is some finite probability $`p_{sky}`$ of rejecting an arbitrary sky pixel, and demand that the probability $`p_{obj}`$ of rejecting a pixel containing positive object flux not exceed $`p_{sky}`$.
Consider a threshold level in image $`J`$, $`t=k\sigma _J`$. The expected count level in $`J`$ is given by equation 4, and the noise level there is given by equation 11. We demand that
$$S_J(q)k\sigma _J(q)k\sigma _J(\mathrm{})$$
(12)
in order to ensure that $`p_{obj}p_{sky}`$. By using our previous expressions for $`S_J(q)`$ and $`\sigma _J(q)`$, we convert this into a constraint on $`\alpha `$. An immediate (though weak) constraint is that $`0\alpha <1/2`$, since $`\sigma _J(0)>\sigma _J(\mathrm{})`$, and the expected count rate must be positive to compensate for the increased noise at the star’s location.
In the remainder of section 2, we derive conditions guaranteeing that inequality 12 will hold for all values of $`S_0`$ and $`q`$. Readers who are not interested in the mathematical details may wish to skim section 2.3, which explains how to choose the parameter $`\alpha `$, and then move on to section 3, where we discuss implementation of the cosmic ray rejection algorithm.
By construction, condition 12 is fulfilled as an equality for $`S_0=0`$ and for any value of $`q`$. To ensure that 12 holds for all $`S_0>0`$, it is sufficient to show that
$$\frac{d}{dS_0}\left[S_J(q)k\sigma _J(q)\right]0$$
(13)
for all $`S_0>0`$ and for arbitrary $`q`$. Multiplying relation 13 by $`\mathrm{exp}[+q^2/(4\xi ^2)]`$, substituting previous results for $`S_J(q)`$ and for $`\sigma _J(q)`$, and using $`d\sigma _J(q)^2/dS_0=2\sigma _J(q)\times d\sigma _J(q)/dS_0`$, we obtain
$$=(\frac{1}{2}\alpha \mathrm{exp}\left[\frac{q^2}{4\xi ^2}\right])\frac{k}{2g}\mathrm{exp}\left[\frac{q^2}{4\xi ^2}\right]\{\frac{\mathrm{exp}\left[+q^2/6\xi ^2\right]}{6\pi \xi ^2}\frac{\alpha }{\pi \xi ^2}+\alpha ^2\}/$$
$$\left\{\sigma _I^2\left[\frac{1}{4\pi \xi ^2}\frac{\alpha }{\pi \xi ^2}+\alpha ^2\right]+\frac{S_0}{g}\mathrm{exp}\left[\frac{q^2}{2\xi ^2}\right]\left[\frac{\mathrm{exp}\left[+q^2/6\xi ^2\right]}{6\pi \xi ^2}\frac{\alpha }{\pi \xi ^2}+\alpha ^2\right]\right\}^{1/2}0$$
(14)
for all $`S_0>0`$.
We now assert that for well-behaved images, it is possible to choose $`\alpha `$ so that relation 14 is fulfilled as an equality for $`S_0=0`$ and $`q=0`$. We will justify this assertion in section 2.3 below.
Taking as a hypothesis for now that $`=0`$ for $`S_0=0`$ and $`q=0`$, we first examine the case $`S_0=0`$, $`q>0`$. Substituting $`S_0=0`$ in equation 14 and rearranging,
$$1/2\mathrm{exp}\left[\frac{q^2}{4\xi ^2}\right]\left\{\alpha +\frac{k}{2g\sigma _I}\left(\alpha ^2\frac{\alpha }{\pi \xi ^2}\right)/\sqrt{\frac{1}{4\pi \xi ^2}\frac{\alpha }{\pi \xi ^2}+\alpha ^2}\right\}$$
$$+\mathrm{exp}\left[\frac{q^2}{12\xi ^2}\right]k/\left\{12\pi g\sigma _I\xi ^2\sqrt{\frac{1}{4\pi \xi ^2}\frac{\alpha }{\pi \xi ^2}+\alpha ^2}\right\}.$$
(15)
We see that the right hand side contains two exponentially decreasing terms. Provided both are positive, the right hand side of equation 15 will clearly be a decreasing function of $`q`$ for all $`q0`$. The second term is positive since all of its factors are positive by definition. Now, by hypothesis, relation 15 is an equality for $`q=0`$, so that the first term in 15 will also be positive provided that the second is $`<1/2`$ for $`q=0`$. This yields a quadratic constraint on $`\alpha `$:
$$\alpha ^2\frac{\alpha }{\pi \xi ^2}+\frac{1}{4\pi \xi ^2}>\left[\frac{k}{6\pi g\sigma _I\xi ^2}\right]^2$$
(16)
This now becomes our sufficient condition for inequality 15 to be fulfilled for all $`q`$.
Turning our attention to $`S_0>0`$, we observe by inspecting 14 that $``$ is an increasing function of $`S_0`$ for any fixed $`q`$ provided only that
$$\alpha ^2\frac{\alpha }{\pi \xi ^2}+\frac{1}{6\pi \xi ^2}>0.$$
(17)
Thus, if we can find a value of $`\alpha `$ that simultaneously fulfills relation 14 as an equality, and fulfills inequalities 16 and 17, we have a convolution kernel that will allow rejection of cosmic rays without any excess risk of rejecting a valid pixel just because it contains flux from an object. In the next section, we determine the parameter space over which this is possible.
### 2.3 The choice of $`\alpha `$
We now turn to deriving the value of $`\alpha `$ that fulfills our earlier assertion, satisfying 14 as an equality for $`S_0=0`$ and $`q=0`$. This is easier if we first define the auxiliary parameter $`\beta =1/2\alpha `$. The intuitive significance of $`\beta `$ is that the expected counts in the filtered image $`J`$ are $`S_J(0)=\beta S_0`$ at the location of a star having $`S_0`$ counts in original image $`I`$. Substituting $`1/2\beta `$ for $`\alpha `$, setting $`S_0`$ and $`q`$ to zero, and requiring exact equality, expression 14 becomes
$$\beta =\frac{k}{2g\sigma _I}\frac{\left(1\frac{1}{\pi \xi ^2}\right)\left(\frac{1}{4}\beta \right)+\beta ^2\frac{1}{12\pi \xi ^2}}{\sqrt{\left(1\frac{1}{\pi \xi ^2}\right)\left(\frac{1}{4}\beta \right)+\beta ^2}}.$$
(18)
This equation can be rearranged into a quartic in $`\beta `$ (or equivalently $`\alpha `$). Rather than doing so, we note that the present version can be solved iteratively for $`\beta `$ by calculating the right hand side a few times, inserting $`\beta =0.25`$ (or indeed any number in $`[0,0.25]`$) the first time and using the previous result at each successive iteration. $`\beta `$ is effectively a function of two parameters, $`k/(g\sigma _I)`$ and $`\xi `$. Now, for most imaging CCD data, we expect reasonable choices of these parameters to be $`3k5`$, $`1g10`$, $`\sigma _I5/g`$, and $`\xi 1`$ pixel. (Our estimate for $`\sigma _I`$ is based on the assumption that the read noise is $`5`$ electrons and the observer will typically ensure that sky noise is greater than read noise.) This leads to $`k/(g\sigma _I)1`$ under typical circumstances. In this limit, we expect
$$\beta \beta _{\text{lim}}=\frac{1}{2}\left(\frac{k}{g\sigma _I}\right)\frac{1/41/(3\pi \xi ^2)}{1/41/(4\pi \xi ^2)}.$$
(19)
One can further show that $`\beta >\beta _{\text{lim}}`$ under then nearly generic conditions that $`0\beta <1/4`$ and $`\beta <1/21/(2\pi \xi ^2)`$. We therefore have a choice when implementing the algorithm between assuming $`\beta =\beta _{\text{lim}}`$ or solving equation 18 iteratively.
We have not derived analytically the range of parameter space over which $`\beta `$ can be found iteratively, but empirically, the iterative solution will converge to a sensible result (fulfilling conditions 14, 16, and 17) provided that $`0<k/(g\sigma _I)<2`$ and that $`\xi >2/\sqrt{3\pi }`$. These are our final set of sufficient conditions for this algorithm to work as desired. They are not as rigorously derived as conditions 14, 16, and 17, but do provide a quick check on when the method is likely to be applicable. Figure 2 shows contours of $`\beta `$ as a function of $`\xi `$ and $`k/(g\sigma _I)`$.
Given our formula for $`\beta `$, it is now also possible to determine $`C_J\sigma _I/(C_I\sigma _J)`$ (the multiplicative reduction in significance level of a cosmic ray after convolution) as a function of $`\xi `$ and $`k/(g\sigma _I)`$. Figure 3 shows contours of this efficiency factor.
## 3 Practical Implementation
We have implemented this algorithm and applied it to several data sets during the past year. In doing so, we introduced several enhancements of the basic algorithm that allow the method to run gracefully on our real data. Two particular artifacts were addressed by these enhancements. First, some bad pixels in some CCD cameras give data values far below the sky level $`B_I`$. If left alone, such pixels will cause many of their neighbors to be flagged as cosmic rays, since the wings of the convolution kernel will spread a strongly negative pixel in image $`I`$ over many neighboring pixels in image $`J`$. Second, cosmic ray hits are often multiple-pixel events. In this case, a cosmic ray pixel may shield less strongly contaminated neighboring cosmic ray pixels from identification. An additional complication is that the background may not be spatially uniform, which hinders measuring the noise level in an image and defining a sensible threshold level for cosmic ray rejection.
The problems of low pixel values and spatially variable background levels can both be handled with preprocessing steps applied prior to the spatial convolution. Nonuniform sky level can be removed by generating and subtracting a smoothed map of the background intensity. I have chosen to use a large spatial median filter for this background generation, but any method working on larger spatial scales than the largest object in the frame would work. The main caveat is that a spatially constant rejection threshold should not be used if the background varies enough to introduce substantial spatial variations in the local Poisson sky noise. Low pixels can be flagged and replaced before the spatial convolution, using a simple threshold operation. This is of course best done after any variable background is subtracted, since the sky level should be uniform for the thresholding operation to be well behaved. Any previously known bad pixels can also be replaced with the background level or an interpolation from their good neighbors at this stage. This approach to background estimation has yielded good results for our data, in which large scale intensity variations are weak ($`<10\%`$ of $`B_I`$) and due primarily to flatfielding errors. For cases where the intensity level contains structure on a wide range of spatial scales, multiscale transform methods (Starck, Murtagh, & Bijaoui 1995) can provide a natural treatment of the background level.
We estimate the variance in the input image empirically, using the iteratively clipped sample variance of the background-subtracted image to determine $`\sigma _I`$, which is taken to be spatially uniform. Spatial variations in the noise level could become a problem for some images. Such variations can be handled with a minor modification to the algorithm, by making the rejection threshold in the filtered image depend on the locally measured noise, provided only that the noise level variations occur on spatial scales large compared to the convolution kernel. Multiscale methods can again be used for accurate noise estimation in the presence of spatially variable backgrounds or extended objects (Starck & Murtagh 1998).
Some multiple-pixel cosmic ray hits will be well handled by a single convolution and flagging step, provided that they remain smaller than the PSF and that they are sufficiently strong (with intensities substantially exceeding the threshold for single pixel events). However, multiple-pixel events usually contain pixels with a range of intensities. When two contaminated pixels of very different intensity lie side by side, the stronger pixel will be flagged but the weaker one will be “shielded” from detection by its prominent neighbor. To identify such “shielded” cosmic ray pixels, the convolution and flagging algorithm can be run iteratively. After every flagging iteration, the newly identified cosmic ray pixels are replaced with the sky level, thereby exposing their less prominent neighbors to scrutiny. This iterative approach is highly successful at flagging all parts of a multi-pixel cosmic ray hit lying above the requested detection threshold.
When setting threshold levels for rejection, we have chosen to use empirical measures of the sky variance in both input and convolved image for convenience. Comparing the results of these empirical measures to the predicted relation given by equation 2 gave agreement at the $`5`$$`7\%`$ level for a test case with $`3`$ pixel FWHM seeing (i.e., $`\xi =1.27`$), with the measured variance of the convolved image slightly exceeding the prediction. Disagreements at this level could be due to several expected effects, e.g., the influence of real objects on the pixel histogram (which increases after smoothing), or the continuous approximation to discrete sums made in deriving equation 2.
This iterative cosmic ray rejection is of course computationally expensive when compared to basic image reduction steps like bias subtraction and flatfielding. Presently, eight iterations of cosmic ray rejection for a 2048$`\times `$4096 pixel image requires of order 10 minutes to run on a 295 MHz Sun Ultra-30 with 248 megabytes of main memory. This speed could be improved by implementing the algorithm entirely in a compiled programming language (the present implementation being an interpreted IRAF script). It is nevertheless fast enough that I have routinely applied the algorithm to large data sets (tens of $`2018\times 4096\times 8`$ pixel images from the Kitt Peak National Observatory CCD Mosaic camera). In principle, the computational requirements should scale as $`n\mathrm{log}(n)`$ for $`n`$ pixels in the large-$`n`$ limit, since the convolution can be implemented using fast Fourier transforms, while the remaining steps should all be linear in the number of pixels. Timing tests on a 400 MHz Intel Pentium-II computer with 128 megabytes of memory yielded a scaling of approximately $`n^{1.4}`$ for images with $`\mathrm{log}_2(n)=22\pm 1`$ (i.e. roughly 2k by 2k pixels). The difference between this scaling and the $`n\mathrm{log}(n)`$ scaling suggested from first principles is perhaps due to the variety of different computational demands (memory, i/o, cpu speed) which can limit the performance of the algorithm for different image sizes.
## 4 Simulations
To verify the analytical results of section 2 and study the effectiveness of the iterated algorithm on multiple pixel events, we carried out three types of artificial data simulations. The first type tested the algorithm’s detection rate for single pixel events; the second tested the false alarm rate at the locations of point sources; and the third tested rejection of multi-pixel events. For both tests of detection rates, the empirically measured detection threshold was taken as the intensity of added cosmic rays for which 50% of the affected pixels were correctly flagged. This is the appropriate cutoff because a simulated cosmic ray of precisely threshold intensity will be boosted above the cutoff by Poisson noise half the time, and will fall below threshold the other half.
The cosmic ray detection rate test added single pixel cosmic ray hits to a noise field and counted the number of hits correctly flagged by the algorithm as a function of CR intensity, rejection threshold, and PSF width. This allows a check of the analytic results plotted in figure 3. The agreement is good, with empirically measured detection thresholds falling between 100% and 110% of the theoretical expectations throughout the tested parameter space. In particular, the measured detection threshold is within 3% of the predicted value for $`k/(g\sigma _I)0.3`$, which is the regime of greatest interest for broadband astronomical imaging.
An interesting variant on the detection test is to run it on a field with stars or other astronomical objects. Our tests showed an appreciable degradation ($`7\%`$) of the CR detection threshold averaged over the image in the presence of a reasonably dense star field. This is expected, since the detection efficiency decreases in the wings of a point source (see section 2). However, there is no good way to characterize this effect for all possible images. If a precise measurement of the detection threshold in some particular image is needed, it can be obtained through simulations by adding “cosmic ray hits” to that exact image and studying their recovery rates.
The false alarm rate test examined the probability of rejecting a given pixel in a pure Poisson noise field with and without point sources. Only the central pixels of the point source locations were considered, since the wings of stellar profiles are less likely than their cores to be incorrectly rejected. This test confirmed that the probability of rejecting the central pixel of a star does not exceed the probability of rejecting an arbitrary sky pixel under our algorithm.
Finally, the multiple pixel event tests placed artificial bad columns onto noise fields and measured the fraction of rejected pixels. Bad columns were chosen as a suitably conservative limiting case of multiple-pixel cosmic ray events, since such events usually have a linear morphology. The realized rejection threshold was determined as a function of stellar FWHM, number of adjacent bad columns, and requested sigma clipping level. These simulations were run with a sky noise of $`41`$ ADU and gain of $`3`$, so that they are restricted to low values of $`k/(g\sigma _I)`$. The general result is (unsurprisingly) that features comparable in size to the PSF cannot be rejected, while features much smaller than the PSF on only one axis can be rejected with relatively modest increases in the intensity threshold for rejection. Results of the multi-pixel event simulations are summarized in table 1.
## 5 Summary and Discussion
We have presented a cosmic ray rejection algorithm based on a convolution of the input image. The advantages of the method spring from the linear nature of the spatial filter, which allows us to determine the noise properties of the filtered image and so to calculate and control the probability of rejecting the central pixel (or indeed any pixel) of a point source. This safety mechanism ensures that cosmic ray rejection can be applied throughout the image, without special treatment for the locations of sources. The sensitivity to cosmic rays is of course reduced at the locations of objects, because of the added Poisson noise contributed by object photons and the resulting need to maintain a positive expectation value in the filtered image.
We usually apply our method conservatively, considering pixels innocent until proven guilty beyond any reasonable doubt. This means that given some uncertainty in the measured point spread function, we use a convolution kernel that is slightly narrower than our best estimate of the PSF (generally by about 10%). This choice depends on the relative importance of keeping legitimate sources and rejecting spurious ones for the scientific problem at hand.
Our original goal in developing this algorithm was to flag and replace cosmic ray hits in individual exposures that are later aligned and stacked. The alignment procedure requires interpolating the original images, and we use sinc interpolation to preserve the spatial resolution and noise properties of the input image. However, sinc interpolation assumes well sampled data and responds badly to cosmic rays, spreading their effects over many more pixels than were originally affected and motivating us to replace them at an early stage. We nevertheless have a second chance to reject cosmic rays by looking for consistency among our different exposures when we stack them, and this second chance helps motivate our generally conservative approach to cosmic ray flagging.
By applying the algorithm developed here followed by sigma rejection during image stacking, we exploit two distinct properties of cosmic rays: They are sharper than the point spread function, and they do not repeat from exposure to exposure. However, we are using these two tests in sequence. An algorithm exploiting both pieces of information simultaneously could potentially yield more sensitive cosmic ray rejection. For general data sets, such an algorithm would have to handle stacks of unregistered images with different PSFs, making its development difficult but potentially rewarding. An interesting effort in this regard is Freudling’s (1995) algorithm, which identifies cosmic rays in the course of deconvolving and coadding images with Hook & Lucy’s (1992) method.
This work was supported by a Kitt Peak Postdoctoral Fellowship and by an STScI Institute Fellowship. Kitt Peak National Observatory is part of the National Optical Astronomy Observatories, operated by the Association of Universities for Research in Astronomy (AURA) under cooperative agreement with the National Science Foundation. The Space Telescope Science Institute (STScI) is operated by AURA under NASA contract NAS 5-26555. I thank an anonymous referee for their remarks. |
warning/0002/physics0002007.html | ar5iv | text | # Teaching the EPR–Paradox at High School ?
## 0.1 Introduction
The first question to be answered is: why should quantum theory be taught at school at all? For choosing this topic there are the following three reasons:
1. Quantum theory is the fundamental theory of modern physics. It plays a significant role in nearly all modern developments of physics. Many recent experiments and research in nanostructures with large applicability in technology rely on quantum effects.
2. Quantum theory has important philosophical aspects. Many people are highly interested in interpretation and understanding quantum theory as shows up in the many popular books about this subject.
3. Pupils at the age from 16+ on are searching for their place in the world. They are trying to understand the world and are open for philosophical hints that help them in building their own world view.
But often this highly fascinating subject is avoided at school because of mathematical and conceptual difficulties. I therefore want to show a possible way to introduce quantum theory in a manner suitable for interested pupils. In this article I concentrate on the mathematical part because here lie some difficulties. For an introduction into the philosophical aspects I have developed a dialogue between philosophers from different times - classical antiquity (Parmenides), the Enlightenment (Kant) and from our century, published elsewhere, (\[Pos98\]). The goal of the following is to clarify the main difficulties in teaching the physical basis of quantum theory and how to keep them minor.
Speaking qualitatively about quantum theory in an adequate way is nearly impossible in itself since all our concepts and terms have been developed along everyday experience. Hence our language is well suited to communicate about concrete physical objects with well determined properties or about psychological issues. In my opinion this last property should be used in dealing with the interfering and superposing objects of quantum theory that may have more similarities or associations with pyschological feelings than with concrete balls or waves occurring in classical physics. In every attempt to talk about quantum theory one has to be aware of this principal difficulty already recognized by Bohr, Heisenberg, Pauli and others. In the complementary worlds of the quantum regime and the classical regime we only are at home in the classical regime. The other regime remains accessible only through sophisticated experiments - even if an experimentalist would call them easy and simple…. One way out might be to talk in images - but soon one arrives at poor analogies. Hence, in order to reach more than only a superficial knowledge at least some hints to the mathematical background of quantum theory must be given. On the other hand at least in the beginning some themes should be avoided to facilitate the pupils the understanding of the peculiarities of quantum theory.
## 0.2 What can be done without too many technicalities?
It is nearly impossible to understand quantum theory without considering its mathematical structure. Nevertheless at school the mathematical apparatus of quantum mechanics has to be abandoned for the major part. The main ideas, however, can be presented quite easily with help of the typical quantum phenomenon “spin” having no classical analogues. Experiments with polarized photons may help in conveying the essentials. In the following I describe the reasons for taking spin as the first subject in treating quantum theory in more detail. Furthermore I explain with the example of EPR–gedanken experiment how to proceed.
Treating the phenomenon “spin” right in the beginning has several advantages:
* Spin lies at the heart of quantum theory. Its properties are used to explain the different statistics, the fine structure of spectra, the splitting of spectra in external (electrical or magnetic) fields. Already a spin system consisting of two particles, i.e. living in a four–dimensional Hilbert space can no longer be described classically. This proof is similar to the proof of Bell inequality, (\[Bau\]).
* The procedure to describe spin mainly by its structure is typical of quantum theory. Furthermore, the mathematics of spin is quite simple, using mainly the well–known Pauli– matrices:
$$\sigma _x=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right),\sigma _y=\left(\begin{array}{cc}0& i\\ i& 0\end{array}\right),\sigma _z=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right).$$
With the help of these quite simple looking matrices, acting on two–dimensional Hilbert space, the most essential mathematical structures of quantum theory can be explained and interpreted, see table 1. Some details are explained in the next section.
* The meaning of Heisenberg uncertainty relation can be explained as principal non–existence of fixed values for properties hence of their non–determination, \[Pos99\]. Therefore the danger that the uncertainty is perceived as measurement mistakes can be drastically diminished. Some helpful constructions even can be visualized on the blackboard.
* Spin is a phenomenon of special importance in modern experiments reaching from Nuclear Magnetic Resonance used in medical applications to realizations of the Einstein–Podolsky–Rosen gedankenexperiment. Its treatment opens the way to a discussion of philosophical aspects of quantum theory which quickly reaches the main points: the question of reality and objectivity in nature treated on a mathematical and physical foundation.
### 0.2.1 An Example
As an example I show how the arguments of Einstein, Podolsky and Rosen in their famous paper \[EPR35\] can be used to show the power of the mathematical formalism and - even more important - how the mathematical constructions can be interpreted in this framework. This allows a bridge to be built from the mathematical structures over the physical phenomena and connecting to a philosophical discussion. Instead of arguing very sophisticated within the mathematical formalism the main goal should be to uncover the main aspects of quantum theory and in this way to build a solid fundament from which the mathematics can be developed further, (see also table 1).
The argumens of EPR can be developed the following way:
#### Step 1: The mathematical tools
In 1935, the year of the EPR-paper the mathematical framework has just been settled implying the following main points:
* The state of a quantum object is given by a state vector $`\psi `$ containing all the available information, i.e. a complete description of the physical properties of the quantum object.
* Each physical quantity is given together with all the possible results of a measurement of that quantity and corresponding eigenstates, i.e. all the states a quantum object can attain after a measurement. A mathematical realization of this concept is given for instance by matrices.
* Arbitrary states can be expressed with aid of the eigenstates of such a matrix resp. physical quantity.
In a deviation from the original argument of EPR I would advise taking the spin realized with the above–mentioned Pauli–matrices as a concrete example. The students can compute the eigenvalus and eigenstates easily from the matrices. The possible measurement results (eigenvalues) are $`+1`$ and $`1`$ together with the corresponding eigenstates. The first expriment showing this property directly has been the Stern–Gerlach–experiment. Furthermore the Pauli–matrices fulfill the condition crucial for the next step of argument of EPR: they do not have any eigenstates in common. Hence there always are several possibilities to represent the spin state of a quantum object, namely with respect to the respective eigenstates of the different matrices correponding to the spin di rections. The representation of an arbitrary spin state $`\psi (s)`$ with respect to the eigenstates of $`\sigma _x`$ would be
$$\psi (s)=c_1\left(\begin{array}{c}\frac{1}{\sqrt{2}}\\ \frac{1}{\sqrt{2}}\end{array}\right)+c_2\left(\begin{array}{c}\frac{1}{\sqrt{2}}\\ \frac{1}{\sqrt{2}}\end{array}\right)$$
and with respect to the eigenstates of $`\sigma _z`$ the same state $`\psi `$ would look like:
$$\psi (s)=k_1\left(\begin{array}{c}1\\ 0\end{array}\right)+k_2\left(\begin{array}{c}0\\ 1\end{array}\right)$$
with different coefficients $`(c_i)(k_i)`$. (For a concrete example look at the table.) This fact perhaps does not matter too much since we always can change coordinates (here it would be the rotation of a coordinate system by an angle of 45 degree). But here it means that the spin state of a given system possesses two different representations belonging to the “same piece of reality”(EPR). The most interesting thing happens in the next step!
#### Step 2: The experimental setup
Two quantum objects, e.g. photons, are brought into interaction or produced in a single process and hence become entangled i.e. they share a common “history”. After that they are separated from each other without any further manipulation, let us say one is brought to the moon, the second stays on earth.
Because of their common “history” they are described by one common state $`\psi `$ which is not just the addition of the states of the single photons. This consideration is central for the whole argument of EPR. The development of the entangled state of both photons into eigenstates with respect to eigenstates of $`\sigma _x`$ is given by:
$$\psi (s_1,s_2)=\psi _1(s_1)\left(\begin{array}{c}1\\ 1\end{array}\right)+\psi _2(s_1)\left(\begin{array}{c}1\\ 1\end{array}\right)$$
and with respect to the eigenstates of $`\sigma _z`$:
$$\psi (s_1,s_2)=\varphi _1(s_1)\left(\begin{array}{c}1\\ 0\end{array}\right)+\varphi _2(s_1)\left(\begin{array}{c}0\\ 1\end{array}\right)$$
The only difference to the representations above is that the coefficients now depend on $`s_1`$ The meaning of these two representations is that photon 1 is described differently depending on the description chosen for photon 2, namely $`\psi _i(s_1)`$ resp. $`\varphi _i(s_1)`$. This is called the entanglement of the two photons. Therefore I would prefer to call the whole system consisting out of these two photons rather a “diphoton” in order to emphasize that they build one whole (also see step 5 below).
#### Step 3: Classical Assumptions
Assuming a fixed objective reality and demanding that physics has to give a complete description of reality Einstein arrives at a contradiction to the predictions of quantum theory. More precisely, Einstein assumes:
1. Separability
Classical Physics only knows action between objects in direct contact with each other. With “object” in this sense I also denote e.g. fields. Hence if two objects are separated in space, including intermediating fields, all future manipulations on them are absolutely independent from each other. We could summarize this in the sentence: Spatially separated objects also are physically separated. This is an implicit assumption of EPR that is not spoken out directly, but is underlying the whole argument as can be seen in the last paragraph of the famous EPR-paper \[EPR35\].
“Separability” hence means that the respective descriptions of two spatially separated photons should be totally independent from each other.
2. Physical Reality
Einstein defines a pragmatic criterion for reality: Every well determined physical quantity has to have a representation in the theory. The point herein lies in the question: Which properties are well determined? Einstein regarded every physical quantity that can be measured as well determined. But quantum theory deviates in so far from classical physics as not all (in principle) measurable quantites have well determined properties at the same instant. They only possess them as a potentiality.
From this view point the different descriptions from above (step 2) should not occur in a “good” physical theory.
#### Step 4: Quantum Theoretical Outcome
We can get information about the photons only after a measurement. What can possibly happen then? There are several possibilities (as an example):
1. The spin of photon 2 is measured in $`x`$-direction. At the same instant the spin state of photon 1 is $`\psi _1(s_1)`$ or $`\psi _2(s_1)`$ according to the result of the measurement at photon 2.
2. The spin of photon 2 is measured in $`z`$-direction. At the same instant the spin state of photon 1 is $`\varphi _1(s_1)`$ or $`\varphi _2(s_1)`$ according to the result of the measurement performed on photon 2.
That means that photon 1 immediately “knows” the kind of measurement done on photon 2 far away as well as its result. Einstein calls this a “spooky action at a distance”, which may not occur in classical physics.
#### Step 5: Interpretation
The behaviour of both entangled photons is strongly connected to each other, they behave in spite of their spatial separation as one single quantum object. Therefore I propose to call these both a “diphoton” which suggests more clearly that there is only one common state of the whole system, and not an addition of states of separated photons. Furthermore, the outcome of measurements demonstrates that we may not assume that photon 1 or photon 2 had fixed values for their spin directions before measurement. For this purpose one could use the comfortable Dirac–notation for spin states e.g.: $`\psi (s_1,s_2)=|1,0|0,1`$ for an entangled spin state instead of the above used vector–notation. The Dirac–notation has the advantage of showing only the relative directions of spins of both photons, which is the only property that is fixed and well determined (in absence of manipulations). The directions themselves are not determined, they only show up after a measurement. Fixed values of properties do not exist in general, they only emerge in measurements. Once this essential point is grasped the way is open for applications.
This access consequently avoids possible pitfalls which in general erschweren understanding quantum theory.
## 0.3 Which Themes to Avoid in a First Approach?
From historical reasons, having their roots in the development of quantum theory, most ways of teaching the concepts of quantum physics refer to classical models. This “procedere” causes principal difficulties in understanding. Therefore every reference to classical concepts should be avoided as far as possible. The most important points to avoid are:
* Speaking about position and momentum, i.e. about trajectories
If the concepts of position and momentum — well–known from everyday experience — are used at the very beginning of a course in quantum theory there exists the danger of transferring classical thinking to quantum theory, although everybody would say: clearly, in quantum theory there are no trajectories. Conceptual difficulties arising from use of the terms “ position” and “velocity” can be avoided in the simplest manner if these fundamental classical concepts do not play any role in the beginning of a course in quantum theory. Then any association of classical ideas might disappear and students might recognize the philosophical significance of quantum theory far more easily. The most prominent example is the famous Heisenberg uncertainty relation for position and momentum which easily is misunderstood in a sense that the uncertainty simply relies on disturbance by measurement in the usual sense. Instead the uncertainty relations are kind of measure for distinguishing classical behaviour from quantum behaviour in that they determine whether two physical quantities can attain fixed values at the same instant. If two physical quantities can attain fixed values at the same instant the quantum object in question behaves “classically”, if not it displays quantum behaviour as e.g. spin. Hence the role and the implications of the non–existence of fixed values for some properties at the same instant - as expressed in uncertainty relations - might not be fully appreciated in their revolutionary potential if one concentrates on “position” and “momentum”. Besides undesired analogies to Newtonian mechanics the corresponding operators for position and momentum and their eigenstates are mathematically far more difficult to handle than the $`2\times 2`$–spin–matrices.
* Speaking about particle–wave–dualism
Waves and particles both are classical concepts, complementary to each other. One could illustrate their relation by looking at the same object from different sights. E.g. a cylinder standing upright appears completely different from the above (a circle) compared to a look from the side (a rectangle). But this observation does not meet the essential point in quantum theory.
A first step to avoid analogy to classical phenomena would be to use the term “quantum object” instead of wave or particle. Only after the quantum mechanical concepts are fixed there might be a careful use of those “classical” terms be allowed where unevitable. Perhaps the importance of using suitable terms may become clear with the example of the double–slit–experiment. If it is replaced by the so called Taylor–experiment in which photons display at the same time wave properties - they show interference - as well as particle properties - they arrive at distinct points on the film, the necessity of changing concepts gets far more obvious.
* Speaking about spin as sort of spinning around
One should not give an image of spin. Especially one may not think in terms of an electron spinning around. The quantum mechanical spin is simply structure manifesting itself and its behaviour through experiments, especially in the Stern–Gerlach–experiment which may serve as an introductory experiment. As shown above the abstract structure of spin can de introduced to a certain extent, depending on the mathematical capabilities of students.
Those three points are mentioned here because their avoidance breaks with the tradition of teaching and speaking about quantum theory. The preceding sections showed an alternative.
## 0.4 Conclusion and Perspectives
The recent EPR–experiments are the starting point for all the current developments concerning the fundamentals of quantum theory as well as technological utopies in the area of quantum computing and teleportation. In addition it widely opens the door to philosophical discussions. In so far the EPR–gedanken experiment lies at the heart of quantum theory and its interpretation.
The entrance to quantum mechanics with help of the phenomenon spin quickly gives gifted or interested pupils a possibility of discussing the properties of quantum objects, the mathematical structures and the interpretation of quantum mechanics on a technically very modest level but nonetheless quite precise. As the spin inevitably is a purely quantum mechanical phenomenon this opens via the EPR–gedanken–experiment a short way into the crucial points of understanding concepts as well as philosophical implications of quantum theory and hence gives the possibility for people to revisit their view of nature, their Weltbild. I regard this an important contribution to general education. |
warning/0002/astro-ph0002033.html | ar5iv | text | # Multiwavelength Observations of the Second Largest Known FR II Radio Galaxy, NVSS 2146+82
## 1 Introduction
The “giant” radio galaxies (GRGs), which we define as double radio sources whose overall projected linear extents exceed 2$`h_{50}^1`$ Mpc, are interesting as extreme examples of radio source development and evolution. Members of this class, which comprise only a few percent of all powerful extragalactic radio sources, have been documented for almost 25 years (e.g., Willis, Strom, & Wilson (1974)). They have been used to constrain the spectral aging and evolution of radio sources and as tests for the evolution of conditions in intergalactic environments on Mpc scales (Strom & Willis (1980); Subrahmanyan & Saripalli (1993); Cotter, Rawlings, & Saunders (1996)). Their 1.4 GHz radio powers are generally in the regime $`10^{24.5}<P_{1.4}<10^{26}`$ $`h_{50}^2`$ W Hz<sup>-1</sup>, just above the transition between Fanaroff-Riley Types I (plumed) and II (lobed) radio structures (Fanaroff & Riley (1974)). It is unclear whether the giant sources are examples of unusually long-lived (and directionally stable) nuclear activity in radio-loud systems, or of the development of sources in unusually low-density environments.
Because of their large angular sizes, nearby giant radio galaxies can in principle be studied in great detail, but their largest-scale structures may be over-resolved and undersampled by interferometers. They have traditionally been discovered through sky surveys with compact interferometers or single dishes at relatively low frequencies, where angular resolution is modest but large fields of view and diffuse steep-spectrum structures can be imaged more easily. The source NVSS 2146+82 was noted as a candidate giant radio galaxy when it appeared in the first 4 by 4 field surveyed by the NRAO VLA Sky Survey (NVSS: Condon et al. (1998)), a northern-hemisphere survey at 1.4 GHz using the VLA D configuration at 45″ (FWHM) resolution.
Figure 1 shows contours of the NVSS image at 45″ resolution. There are two symmetric, extended lobes (D and E) on either side of an unresolved component C, plus an unusually large number of other radio sources within 10′ of C. Two of these (A and B) are also symmetrically located around C.
Comparison with the Digital Sky Survey (DSS) showed that source C coincides with an $``$18<sup>th</sup> mag elliptical galaxy to within the uncertainties in the NVSS and DSS positions. If the elliptical galaxy is the host of an unusually large radio source (C+D+E), then the apparent magnitude suggests that the whole structure may be similar in linear scale to 3C 236. The DSS also shows a nearby image that might be another galactic nucleus, and a faint extended feature suggesting a possible “tail” or interaction.
We have undertaken several observational studies of the radio and optical objects in the field to determine their nature and to clarify the relationships between the optical and radio sources. These studies include:
1. High resolution radio imaging at 4.9 and 8.4 GHz to locate any compact flat-spectrum radio components in the field, and thus to identify any AGN that could be responsible for some or all of the other radio emission,
2. A search for fainter diffuse radio emission between the D and E components that might link them together or to other sources in the field and thus clarify their physical relationship,
3. Higher-resolution radio imaging of the other radio sources in the field to explore whether they might be physically related to the diffuse components, or to each other by gravitational lensing,
4. Optical spectroscopy of both optical “nuclei” and other galaxies in the field,
5. UBVRI optical photometry of the field, and
6. X-ray imaging using ROSAT HRI observations to search for any hot X-ray emitting gas which might be associated with an overdensity of galaxies or non-thermal X-ray emission from an AGN.
Throughout this paper, we assume a Hubble constant $`H_0=50h_{50}`$ km s<sup>-1</sup> Mpc<sup>-1</sup> . At a redshift of $`z=0.145`$, the angular diameter distance to the radio galaxy is 708.4$`h_{50}^1`$ Mpc, the luminosity distance is 928.7$`h_{50}^1`$ Mpc, and $`1\mathrm{}`$ corresponds to 206$`h_{50}^1`$ kpc.
## 2 Radio Observations
Table 1 gives a journal of our VLA observations. The observations in the A configuration were designed to locate any compact radio components in the field. Those in the B, C, and D configurations were intended to image the largest scale emission in enough detail to reveal any relationships and connections between the extended components, as well as to determine their spectral and Faraday rotation/depolarization properties. The BnC configuration data were designed as a sensitive search for connections, such as jets, between the central radio source and the extended features.
The flux density calibration was based on 3C 48 and 3C 286. The on-axis instrumental polarization corrections were determined from observations of the unresolved synthesis phase calibrator 2005+778, and the absolute polarization position angle scale from observations of 3C 286. Multiple observations of 3C 286 and other polarized sources were used to detect problems with ionospheric Faraday rotation, but none was noted in any of the sessions. The data were calibrated using the source 2005+778 as an intermediate phase reference, then self-calibrated using AIPS software developed by W. D. Cotton for the NVSS survey.
Due to the large size of this source, 1.4 and 1.6 GHz observations used three pointings; one on the central source C, and one near the center of each putative lobe. The B, C, and D VLA configuration observations were made at 1.365 and 1.636 GHz to measure rotation measure and spectral index. The data from these frequencies were calibrated and imaged separately. Data taken in the BnC configuration were in two adjacent 50 MHz bands centered on 1.4 GHz. Since the source extent is comparable to that of the antenna pattern and the bandwidth used was relatively large, the deconvolution (CLEAN) and self calibration applied corrections for the frequency dependence of the antenna pattern. Data from each of the three pointings were imaged independently and combined into a single image by interpolating the images onto a common grid, averaging weightings by the square of the antenna power pattern, and correcting for the effects of the antenna pattern. The 0.3 GHz observations were of limited use owing to interference.
### 2.1 Radio Results
The most sensitive image of NVSS 2146+82 is derived from our BnC configuration data at 1.4 GHz which has a resolution of 13″ (FWHM). Figure 2 shows logarithmic contours of the total intensity in the region around the source in this image; the rms noise is 20 $`\mu `$Jy per CLEAN beam area. A gray scale representation of the same image showing the filamentary structure of the lobes is given in Figure 3. Figure 4 shows the inner region of this image contoured to lower levels using an initially linear contour interval.
### 2.2 Association of Features
The structures of the extended features D and E shown in Figures 2 and 3 are entirely consistent with their being associated with each other as the two lobes of a large FR II double source of overall angular size 19$`\stackrel{}{\mathrm{.}}`$5. Both features are brightest in the regions furthest from C, contain bright (but resolved) substructure near their outer edges resembling the hot spots of FR II sources, and have their steepest brightness gradients on their outer edges. The overall length of the two lobes is the same to within 5%. Although features A and B in Figure 1 appear symmetric around feature C, the higher resolution VLA images (Figures 2 and 3) reveal them to be background sources, unrelated to NVSS 2146+82.
The northern feature (D) contains a region of enhanced emission (hot spot) at its northern extremity with about 65 mJy in an area 30″ by 18″ and an L-shaped extension to the West. The southern feature (E) has 75 mJy in a region of enhanced emission 50″ by 30″ (a “warm spot”) recessed by 10% of the distance from the core and sharp brightness gradients around its southern and western boundaries. Both regions of enhanced emission show evidence of finer, but resolved, structure in our data taken in the B configuration (see contour plots in Figure 5). Figure 3 clearly shows that the internal brightness distributions of both lobes are non-uniform, and suggest the presence of filamentary structures, again a common characteristic of FR II radio lobes at this relative resolution.
Most importantly, Figures 2, 3, and 4 also show that these lobes are linked to the central compact feature C by elongated features that are plausibly the brightest segments of a weak jet-counterjet system. These features are labeled in Figure 4.
We interpret the following features as belonging to the jet in the south lobe.
J1. This feature is clearly part of a jet that points towards the south lobe but not directly at the peak of feature E.
J2. This feature (1$`\stackrel{}{\mathrm{.}}`$5 from C) and feature K (1$`\stackrel{}{\mathrm{.}}`$4 to the north of C) are roughly symmetric in distance from C and in intensity but are not quite collinear with C. On both sides of the source the jet becomes harder to trace further into the lobe. J2 appears to be south of the C–J1 direction, suggesting a southward bend, however.
J3. This feature is plausibly a knot in the continuation of the jet into the south lobe. The lobe brightens beyond J3 and contains a diffuse ridge that is a plausible continuation of the (possibly decollimated) jet in the direction of the “warm spot” E. The north lobe also brightens at about the same distance from C although there is no feature corresponding to J3 in the north.
Table 2 gives flux density estimates for the main features of the source. We estimate that the jet and counterjet together comprise about 1% of the total flux density of the extended lobes, a typical jet “prominence” for radio galaxies slightly above the FR I–II transition.
The higher-resolution radio images provide no evidence that sources A, B, or F in Figure 1 are physically related to each other, or to C, D and E. Although none can be optically identified, we consider it likely that these are three (or more) unrelated background sources. The symmetrical alignment of A and B around C is apparently coincidental, and there is no evidence for any radio “bridge” between these sources and component C.
### 2.3 Polarimetry
The polarization structure derived from the sensitive BnC configuration observations is shown in Figure 6. The 1.4 and 1.6 GHz data are sufficiently separated in frequency to enable us to measure Faraday rotation but still maintain comparable surface brightness sensitivity. The derived rotation measure images of the two lobes are shown in Figure 7. The rotation measure distribution over the north lobe is featureless but several filamentary rotation measure structures can be seen over the southern lobe. The average rotation measure is about the same in the two lobes, -9 rad m<sup>-2</sup> in the north and -8 rad m<sup>-2</sup> in the south. The Faraday rotation measure in the south lobe has a somewhat larger root mean square variation, 8 rad m<sup>-2</sup> compared to 5 rad m<sup>-2</sup> in the north.
### 2.4 Spectral Index Distribution
Figure 8 shows the 0.35 to 1.4 GHz spectral index distribution inferred from comparing the WENSS (Rengelink et al. (1997)) image with our BnC configuration image convolved to the same resolution. The northern and southern warm spots have spectral indices<sup>1</sup><sup>1</sup>1Spectral index, $`\alpha `$, as used here is given by $`S=S_0\nu ^\alpha `$. $`\alpha _{0.35}^{1.4}`$ of -0.6 and -0.55, not unusual for the hot spots of FR II sources in this frequency regime. The background sources also exhibit spectral indices that are quite typical of extragalactic sources (A, -0.68; B, -1.0; F, -0.7). Near the centers of the north and south lobes of NVSS 2146+82, however, this comparison shows regions of unusually “flat” spectral index ($`\alpha _{0.35}^{1.4}0.3\pm 0.02`$ in the north lobe, $`\alpha _{0.35}^{1.4}0.4\pm 0.03`$ in the south lobe).
The spectral index variations across the lobes can also be studied from our 1.36 and 1.63 GHz data. Due to the low surface brightness the data were tapered to 55$`\mathrm{}`$ resolution before imaging for this comparison. To eliminate any complication from the mosaicing technique, only data derived from the pointing on a given lobe were used to determine the spectral index variations for that lobe. Thus, the data from two pointings were imaged independently at 1.36 and 1.63 GHz, corrected for the antenna power pattern, and spectral index images were derived independently for the two lobes. These results are shown in Figure 9. The close spacing of the frequencies makes determining the spectral index more difficult; but this is compensated to some extent by the nearly identical imaging properties at the two frequencies, which reduce systematic errors. These data sets are fully independent of those used for the spectral index image in Figure 8, but also reveal symmetric regions of unusually flat spectral index, $`\alpha _{1.4}^{1.6}`$ -0.3$`\pm 0.08`$, in both lobes.
We conclude that four independent data sets show evidence for regions with $`\alpha _{0.35}^{1.4}`$ -0.3 in regions of relatively high signal to noise ratio. These regions are not artifacts of “lumpiness” in the zero levels of the images.
### 2.5 Source Alignment
NVSS 2146+82 is not aligned along a single axis. The two warm regions (E and D) and the core (C) are not collinear. The jet in the south appears to have several bends; one near the end of J1 (see Figure 4) where it bends toward J2, a change in position angle from -150 to -170. Beyond J3, the ridge line of the lobe is fairly well defined and is again at position angle -150, consistent with a second bend (apparently $``$ 20) in the neighborhood of J3. The jet is not so prominent in the north but feature K, which may be the brightest part of a counterjet, is elongated along position angle of -169.
The general “C” shape of the source suggests that the overall misalignment is due to environmental effects that have bent the jets, rather than to a changing initial jet direction which is likely to produce overall “S” symmetry.
We consider it beyond doubt that C, D and E comprise a single large FR II radio source with weak radio jets, whose parent object is the galaxy identified with C.
## 3 Optical Observations of NVSS 2146+82 and its Environs
Optical photometric and spectroscopic observations were obtained to identify the host galaxy of the radio emission and to measure its redshift. We began the search for the optical counterpart to the radio source using the Digitized Sky Survey (Lasker et al. (1990); hereafter DSS). The radio core is aligned with an elliptical galaxy on the DSS image to within the astrometric accuracy of the radio and optical positions from the NVSS and DSS. There is also a second, equally bright object a few arcseconds east of the galaxy at the radio core position. Finally, in the DSS image, there appears to be S–shaped diffuse emission that passes through both bright “nuclei”. Therefore, our initial assumption was that the host galaxy of NVSS 2146+82 was possibly a disturbed, double nucleus galaxy. In the following sections, we summarize the optical imaging of the field surrounding the candidate host galaxy and the spectroscopic observations of this host galaxy and its candidate galactic companions.
### 3.1 Photometric Observations
U, B, V, R, and I CCD observations were obtained at the 1.52-m telescope at Palomar Observatory on the nights of 7-9 January 1997. In addition, U, B, V, and I CCD observations were made at Kitt Peak National Observatory on 4 April 1997. The Palomar 1.52-m observations were made with a 2048 $`\times `$ 2048 CCD with a pixel scale of 0$`\stackrel{}{\mathrm{.}}`$37 per pixel, resulting in a 12$`\stackrel{}{\mathrm{.}}`$63 field of view. Though photometric, the seeing was poor ($`25\mathrm{}`$ on 7 January, $`1.52.5\mathrm{}`$ on 8,9 January) during the Palomar run, so higher resolution ($`1.21.4\mathrm{}`$ seeing) images were obtained with the KPNO 4-m telescope in April. The KPNO observations were made with the prime focus T2KB CCD with a pixel scale of $`0\stackrel{}{\mathrm{.}}47`$ per pixel, resulting in a 16$`^{^{}}`$ field of view. Because the KPNO data were not taken in photometric conditions, the Palomar data remained useful for calibration. Data from both observing runs were reduced using the standard IRAF CCDRED reduction tasks.
After the initial reduction, aperture photometry was performed on the host galaxy of NVSS 2146+82 using the IRAF package APPHOT. Unfortunately, due to the poor seeing on the first night of the Palomar run and the proximity of the foreground star (see §3.3) to the AGN host, it was impossible to photometer NVSS 2146+82 without significant flux contamination from the foreground star. Therefore, we used the DAOPHOT II package (Stetson (1987)) to PSF fit and subtract stars from the Palomar NVSS 2146+82 images.
After the foreground star was subtracted, photometry of the galaxy was performed identically to the photometry of several Landolt (1992) standard stars. Approximately 20 stars were selected from each frame containing the AGN host galaxy. A circular aperture 2.5 times the average FWHM of these stars was used to measure the flux of the host galaxy. This aperture was chosen to be consistent with the standard star photometry and because it completely enclosed the host without including contaminating flux from other nearby objects.
Once instrumental magnitudes for the galaxy were determined, they were transformed to the standard system using transformation equations incorporating an airmass and color term that were determined for the Landolt standard stars. The results of our U,B,V,R,& I photometry of the host galaxy are listed in Table 3.
### 3.2 Spectroscopic Observations
Optical spectra of NVSS 2146+82 were obtained at Kitt Peak National Observatory on 9 December 1996. The spectroscopic observations were made with the RC Spectrograph on the KPNO Mayall 4-meter telescope. The detector in use was the T2KB CCD in a 700 $`\times `$ 2048 pixel format. All exposures were made with a $`1\mathrm{}`$ slit width and a 527 lines/mm grating. The spectral resolution, measured using unresolved night sky lines, is $``$3.4 Å. The data were reduced using the standard IRAF reduction tasks. The extracted spectra were wavelength calibrated using a solution determined from the spectrum of a HeNeAr comparison source. Finally, spectrophotometric calibration was applied using a flux scale extrapolated from several standard star spectra.
Spectra of candidate galactic companions to NVSS 2146+82 (see §3.5 below) were obtained with the HYDRA multi-fiber positioner and the Bench Spectrograph as part of the WIYN<sup>6</sup><sup>6</sup>6The WIYN Observatory is a joint facility of the University of Wisconsin-Madison, Indiana University, Yale University, and the National Optical Astronomy Observatories. Queue Experiment over the period of 14-22 September 1998. The T2KC CCD was used as the spectrograph detector in its spatially binned 1024 $`\times `$ 2048 pixel mode. All exposures were made with the red fibers, the Simmons camera, and a 400 lines/mm grating. The spectral resolution in this configuration is $``$4.5 Å.
We calculated an astrometric solution for the KPNO 4-m frame of the NVSS 2146+82 field using positions for stars in the frame taken from the USNO A1.0 catalog (Monet et al. (1996)). Using this solution, we derived positions with the accuracy required by the HYDRA positioner for our target galaxies. Due to fiber placement restrictions and the density of our target galaxies on the sky, we were only able to place 46 fibers on targets. The remaining 50 fibers were randomly placed on blank sky, and they were used during the reduction process for night sky subtraction.
The nine 30 minute program exposures were reduced using the IRAF DOHYDRA script. The weather conditions during the last two nights were poor, and the spectra from these nights were not usable. Therefore the final spectra were obtained by co-adding only the data from nights one and two, a total of two hours of integration.
### 3.3 Redshifts and Line Luminosities
In Figure 10, we present a contour plot of the V band surface brightness from the central $`40\mathrm{}\times 40\mathrm{}`$ region of the KPNO 4-m image after smoothing with a 3 pixel by 3 pixel boxcar kernel. Although we find that the elliptical galaxy at the radio core position ($`\alpha =21^\mathrm{h}45^\mathrm{m}30^\mathrm{s}`$, $`\delta =+81\mathrm{°}54\mathrm{}55\mathrm{}`$ J2000.0) has a narrow line AGN emission spectrum with a redshift of $`z=0.145`$, we find that the object just to the east, which was assumed to be potentially a second nucleus, has a zero-redshift stellar spectrum, indicating it is a foreground star. Figure 11 shows two plots of the wavelength and flux calibrated spectrum of the host galaxy of NVSS 2146+82.
An unusual feature of the spectrum (Figure 11) of the AGN is that all of the emission lines appear to be double peaked. The second panel in Figure 11 shows an expanded view of the $`[\mathrm{OIII}]`$ doublet clearly showing the double peaked profile of the emission lines. Each emission line was easily fit with a blend of two gaussians, indicating that AGN line emission is coming from two sources with a velocity separation of $``$450 km s<sup>-1</sup>.
Since the AGN emission line spectrum gives two different velocities, we have decided to take the velocity of the stellar component of the galaxy as the systemic velocity of the galaxy. The stellar absorption line redshift, calculated by cross-correlating the host galaxy spectrum with the spectrum of the star immediately to the east, is 0.1450$`\pm `$0.0002.
Table 4 lists the properties of the observed emission features in the spectrum of NVSS 2146+82. The redshifts of the AGN emission line components were calculated by identifying features and taking the average redshift of all of the identified features. In this way, the two AGN emission line components have been measured to be at velocities of 40070$`\pm `$50 km s<sup>-1</sup> and 40520$`\pm `$50 km s<sup>-1</sup>, which corresponds to redshifts of 0.1440$`\pm `$0.0002 and 0.1456$`\pm `$0.0002 respectively. This indicates that the gas which is giving rise to the bluer component of the AGN emission line spectrum is moving relative to the stars in the AGN host galaxy at $`280`$ km s<sup>-1</sup> and the gas emitting the redder lines is moving at 170 km s<sup>-1</sup> with respect to the stars.
Each emission feature identified in Table 4 was fit with a blend of two gaussian components (except for the two weak lines $`[`$Ne III$`]`$ $`\lambda 3967`$ and $`[`$O III$`]`$ $`\lambda 4363`$, where a single gaussian was used) to determine the line flux. The fluxes listed in Table 4 were measured after the spectrum of NVSS 2146+82 was flux calibrated using the average of four measurements of the calibrator Feige 34. The flux of the calibrator varied significantly among our four separate exposures, and we therefore estimate our spectrophotometry is only accurate to about 20%. In addition to calibration error, there is an additional error in the profile fitting, and therefore the errors listed for the fluxes include both calibration and measurement error.
We derived an extinction of $`A_V=0.9\pm 0.9`$ (the large error is due mostly to the calibration error in the fluxes) using the standard Balmer line ratios for Case B recombination (Osterbrock (1989)) and the extinction law of Cardelli et al. (1989). The Galactic extinction at the position of NVSS 2146+82 is given as $`A_V=0.5`$ on the reddening maps of Schlegel et al. (1998). This value is consistent with our Balmer line derived value, but possibly indicates that there may be some dust in the host galaxy itself. We decided to correct the measured line fluxes for reddening using the mean value we derived of $`A_V=0.9`$. The errors listed in Table 4 for the fluxes do not include the error in the extinction determination.
### 3.4 Optical Properties of the Host Galaxy
Sandage (1972) found that the optical luminosity function of radio galaxy hosts was similar to that of first ranked cluster members, and he noted that their optical morphology was similar to bright E galaxies. Although it was therefore generally believed that the hosts of all radio galaxies were gE types, subsequent large surveys of radio galaxies showed a good deal of evidence for peculiar morphologies (e.g., Heckman et al. (1986)). We find that the host galaxy of NVSS 2146+82 is likely typical, i.e. it is a gE galaxy, but with evidence of some peculiar morphological features.
The broadband colors of NVSS 2146+82 are typical of bright FR II host galaxies. The absolute magnitude we derive for the host is $`M_V=22.9`$ at $`z=0.145`$ if we adopt a K correction of 0.46 magnitudes in the V passband (Kinney et al. (1996)). This magnitude is consistent with the host being a gE galaxy, and also is very similar to the mean V magnitude for 50 low redshift FR IIs of -22.6 (Zirbel (1996)).
Similar to other FR II host galaxies, we find the optical morphology of the host elliptical of NVSS 2146+82 to be disturbed. In Figure 10, the four distinct objects besides the host galaxy and foreground star have been identified as having non-stellar morphologies with the Faint Object Classification and Analysis System (FOCAS, Valdes (1982)). If these four galaxies share the same redshift as the gE host of NVSS 2146+82, they all lie 50–100 kpc away from its nucleus, a distance that implies that they may be dynamically interacting with it. Figure 10 also shows what appears to be a bridge of diffuse optical light that almost connects NVSS 2146+82 to the galaxy to the southwest. This bridge may indicate that this smaller galaxy has recently passed close enough to NVSS 2146+82 to interact with it gravitationally. There is also a fifth object $`5\mathrm{}`$ to the southeast of the center of NVSS 2146+82, which could be in the process of merging with the gE galaxy. However, due to the faintness of this object and its proximity to the nucleus of 2146+82, we are unable to classify this object definitively as a galaxy with the FOCAS software. Although we cannot conclude based on this image that NVSS 2146+82 is undergoing a merger, its outer isophotes do show evidence that it has been disturbed.
Correlations between the radio power and optical emission line luminosities in radio galaxies have been established in several studies (e.g., Rawlings & Saunders (1991); Zirbel & Baum (1995); Tadhunter et al. (1998)). These radio/optical correlations are assumed to arise primarily due to the fact that both the radio jet and the ionization source originate in the central engine. The radio core power at 5 GHz ($`\mathrm{log}P[W/Hz]=23.85`$) and the H$`\alpha +[`$N II$`]`$ luminosity ($`\mathrm{log}L[W]=35.2`$) for NVSS 2146+82 lie well within the dispersion in the correlation in these quantities found for low redshift FR IIs (Zirbel & Baum (1995)). This apparently indicates that the physical conditions that cause this radio/optical correlation to arise may be similar in this GRG and in “normal” FR IIs.
The shape of the emission line profiles in NVSS 2146+82 are not unique; emission line profiles and narrow band imaging of Seyfert galaxies and radio galaxies have shown evidence for interaction between the radio synchrotron emitting plasma and the optically emitting ionized gas (see e.g., Whittle (1989)). Although the majority of objects that show kinematic evidence for interactions between the radio jets and ionized gas clouds tend to have more compact radio structures, the double peaked line profiles seen in NVSS 2146+82 appear similar to those seen in radio galaxies with jet/cloud interactions. A recent model (Taylor, Dyson, & Axon (1992)) for interactions between the nuclear radio emission and NLR gas in Seyferts produces $`[`$O III$`]`$ profiles for objects near the plane of the sky that are very similar to the double peaked profiles seen in NVSS 2146+82. The model of Taylor et al. (1992) produces double peaks in the line profiles of objects oriented close to the plane of the sky because the emission lines are postulated to arise from gas that is being accelerated as a bowshock expands into the ionized medium surrounding the nucleus. They model the bowshock as a series of annuli, and each annulus contributes most of its luminosity at the two extreme radial velocities found along the line of sight. Although the specifics of the model of Taylor et al. (1992), such as the discrete plasmon emission from the radio nucleus, may not necessarily apply in the case of NVSS 2146+82, it suggests that the narrow line profiles observed for this FR II (which is assumed to be very near the plane of the sky) can be produced plausibly in a model where the ionized gas is in a cylindrical geometry around the radio jet.
Double peaked broad lines have been observed in radio galaxies (e.g., Pictor A \[Halpern & Eracleous (1994)\]), however the model that is typically invoked to explain the broad line profiles requires the radio galaxy to be oriented close to the line of sight. Since NVSS 2146+82 does not show a broad line component and is unlikely to be oriented close to the line of sight, the accretion disk model relied on to fit double peaked broad lines in AGN is probably unrelated to the emission line profiles observed in NVSS 2146+82.
Although a jet/cloud interaction appears to be the most reasonable explanation for the double peaked narrow emission lines observed in the spectrum of NVSS 2146+82, it is also plausible that a gravitational interaction between the FR II host galaxy and its nearest companions may be the source of the $``$450 km/sec separation between the blue and red emission line peaks. Higher spatial resolution long slit spectroscopy is necessary to determine which cause is more likely.
### 3.5 Environment
Deep CCD imaging of the region surrounding the host galaxy of NVSS 2146+82 has revealed a large number of nearby galaxies. These galaxies are near the limiting magnitude of the POSS/DSS images, so NVSS 2146+82 appears to lie in a sparsely populated region of the sky in the DSS. However, photometry from the deeper Palomar 1.52-m images gives $`22M_V19.5`$ for these nearby galaxies if they also lie at $`z=0.145`$, indicating a possible association with NVSS 2146+82. In Figure 12, we present a region of the KPNO 4-m image of NVSS 2146+82 that is 0.5 Mpc on a side and that has all identified galaxies with $`m_v21.3`$ (corresponding to $`M_V19`$ at $`z=0.145`$) circled. These images do not go deep enough to allow accurate identification and photometry of all galaxies to $`M_V=19`$, so this sample is not complete. However, even though the sample shown in Figure 12 is probably only complete to $`M_V20.5`$, we have identified 34 candidate galaxies surrounding NVSS 2146+82.
Although there are no previous identifications of the cluster around NVSS 2146+82 (at $`b=21\stackrel{}{\mathrm{.}}5`$, it is too close to the Galactic Plane to have been included in the Abell catalog), there is a Zwicky cluster to the north, with NVSS 2146+82 lying only $`5\mathrm{}`$ south of the southern border of the Zwicky cluster. The Zwicky cluster 2147.0+8155 (B1950.0 coordinates) is a compact group with 56 members classified as “extremely distant” or $`z>0.22`$ (Zwicky et al. (1961)). While this gives a redshift for the Zwicky cluster larger than that of NVSS 2146+82, it is close enough to $`z=0.145`$ ($`<400`$ Mpc more distant) that we may be seeing NVSS 2146+82 in projection against a background rich cluster.
In September of 1998 WIYN/HYDRA spectra were obtained of 46 candidate galactic companions of NVSS 2146+82 to determine their redshifts. The sample of 46 was selected in the following way: (1) We selected all objects morphologically classified as galaxies in the KPNO 4-m frame by FOCAS with aperture magnitudes $`<21`$, resulting in an initial sample of 205 galaxies. (2) We divided this group into two subdivisions: the first being all galaxies within 0.5 Mpc of 2146+82 in projected radius, and the second being all those outside of the 0.5 Mpc radius. However, due to exposure time limitations, the available sample taken from the 34 galaxies identified in Figure 12 within 0.5 Mpc of the host was reduced to the 17 brightest galaxies. Fiber placement restrictions allowed us to observe only 11 of these 17 galaxies. Objects from the sample outside of the 0.5 Mpc radius from NVSS 2146+82 were assigned to 35 of the remaining fibers, leaving about 45 fibers on blank sky to allow accurate sky subtraction. Unfortunately, as mentioned in §3.2 above, the weather conditions during some of the queue observing were poor, and this limited the success of the program. There was enough signal-to-noise to identify features in the spectra of only 24 of the 46 objects successfully. We found that 7 of the 24 objects with good spectra were actually misidentified stars.
Nonetheless, from the remaining 17 spectra of galaxies in the field surrounding NVSS 2146+82, we were successful in identifying what we believe to be a true cluster that contains the radio source host galaxy. Figure 13 presents an image with the 17 galaxies with measured redshifts marked. The positions, redshifts, and magnitudes for these objects are listed in Table 5. A quality factor is assigned for each redshift using the 0 (unreliable) to 6 (highly reliable) scale of Munn et al. (1997). The quality is determined using: $`q=\mathrm{min}[6,\mathrm{min}(1,N_{def}),+2N_{def}+N_{prob}]`$, where $`N_{def}`$ is the number of spectral features that are accurately identified (less than 5% chance of being incorrect) and $`N_{prob}`$ is the number of spectral features that are probably correct (about a 50% chance of being correct). If $`q>3`$ is adopted as the requirement for a reliable redshift, 5 of the 17 galaxies have unreliable redshifts. The histogram plotted in Figure 14 is a redshift distribution for the 17 galaxies, and it shows that 50% (6) of the reliable redshifts fall in the range of $`z=0.1350.148`$, with 5 of those having redshifts of $`z=0.1440.148`$.
Extrapolating the redshift distribution for the sample of galaxies identified around NVSS 2146+82 from the redshift distribution of the 17 reliable galaxy spectra suggests that the 2146+82 cluster may be Abell richness class 0 or 1. Of course, the statistics are very uncertain. Of the 11 galaxies within a projected distance of 0.5 Mpc of NVSS 2146+82 that were in the WIYN/HYDRA sample, redshifts were measured for three of them. Two of these have $`z=0.1440.145`$, while the third has $`z=0.135`$. We identified features in 21 of the remaining 35 spectra that were measured for objects outside of the projected 0.5 Mpc radius. We found that 7 were misclassified stars, and 3 of the 14 galaxies with reliable redshifts had $`0.144<z<0.148`$. Abell’s (1958) richness criterion was based on the number of cluster galaxies within the range $`m_3`$ to $`m_3+2`$ ($`m_3`$ is the magnitude of the third brightest cluster member). For the NVSS 2146+82 cluster, $`m_3`$ should be $`<18.3`$, since the third brightest galaxy of the 7 (which includes NVSS 2146+82) we have found at $`z=0.145`$ has $`m=18.3`$. Of the 205 galaxies originally found in the KPNO 4-m field containing NVSS 2146+82, 123 of these fall within the $`m_3`$ to $`m_3+2`$ range used for estimating the Abell richness. If we apply the percentages above to this sample of 123 galaxies, then $`37\pm 13`$ might be at the same redshift as NVSS 2146+82. To this point, we have been considering the cluster richness inside of 0.5 Mpc, for comparison with the $`N_{0.5}^{19}`$ richnesses of Allington-Smith et al. (1993) and Zirbel (1997), and also within an area $``$3.8 Mpc on a side, which is the size of the KPNO 4-m field at $`z=0.145`$. However, we must note that the original richness criterion for Abell class 1 clusters was that 50 or more galaxies were contained in a radius of 3 Mpc for $`H_0=50`$ km sec<sup>-1</sup> Mpc<sup>-1</sup> (Abell (1958)). A circle of radius 3$`h_{50}^1`$ Mpc at $`z=0.145`$ subtends 507 square arcminutes on the sky, nearly twice the amount of area covered in our image. If the calculated optical richness from the 4-m image galaxy sample is taken as a lower limit to the number of galaxies within an Abell radius, the richness class of the group surrounding NVSS 2146+82 appears to be at least Abell class 0.
## 4 X–ray Observations and Constraints
Richness class 0 clusters of galaxies typically have luminosities with L$`{}_{x}{}^{}10^{4345}`$ ergs s<sup>-1</sup> (Ebeling et al. (1998)), while X-ray AGN range from L$`{}_{x}{}^{}10^{4044}`$ ergs s<sup>-1</sup> (Green, Anderson & Ward (1992)), so a cluster or bright AGN will easily be seen with a medium length exposure with ROSAT. NVSS 2146+82 was observed with the ROSAT High Resolution Imager (HRI) between 1998 February 24 and 1998 March 13 for a duration of 30.3 ksec to search for any hot gas that might be associated with the apparent overdensity of galaxies or for an X-ray luminous AGN.
The data were analyzed with the IRAF Post-Reduction Off-line Software (PROS). The HRI data were filtered for periods of high background and corrected for non-X-ray background, vignetting, and exposure using the computer programs developed by Snowden (Plucinsky et al. (1993); Snowden (1998)). After filtering, the live exposure was 29.8 ksec. The resulting X-ray image was convolved with a gaussian beam with $`\sigma =2\mathrm{}`$ to recover diffuse X-ray emission. The contours of the image are shown superposed on the DSS image in Figure 15.
A few sources were visible near the edge of the field, but there seem to be no significant sources of X-ray emission associated with any optical or radio sources within the $`20\mathrm{}`$ extent of NVSS 2146+82 (Figure 15). We derived upper limits on both the AGN or cluster emission by extracting the X-ray counts from the corrected X-ray image using circular regions centered on the host galaxy of $`20\mathrm{}`$ and $`2\stackrel{}{\mathrm{.}}25`$, respectively. The region sizes were chosen simply because $`20\mathrm{}`$ represents the size of a typical HRI point source and $`2\stackrel{}{\mathrm{.}}25`$ is roughly 1-2 times the typical size of a cluster core at the distance of the radio galaxy. The X-ray background was determined by extracting the X-ray counts from an annulus of $`2.255\mathrm{}`$ centered on the nucleus of the radio host and removing 3 point sources using $`20\mathrm{}`$ circular regions. We used PIMMS (Mukai (1993)) to convert the HRI count rate into an unabsorbed flux in the 0.1-2.0 keV band, assuming an emission model and a Galactic photoelectric absorption column of $`1.058\times 10^{21}`$ cm<sup>2</sup> (Stark et al. (1992)). For the AGN, we assumed a power law with a photon index, $`\mathrm{\Gamma }`$, of 2.0 and derived an upper limit at the $`90\%`$ confidence level of $`3.52\times 10^{14}`$ ergs cm<sup>-2</sup> s<sup>-1</sup>, or $`3.63\times 10^{42}h_{50}^2`$ ergs s<sup>-1</sup> at the distance of the radio galaxy. Similarly for the cluster, we assumed a Raymond-Smith thermal emission spectrum characterized by $`kT=1`$ keV which yielded an upper limit of $`1.33\times 10^{13}`$ ergs cm<sup>-2</sup> s<sup>-1</sup>, or $`1.37\times 10^{43}h_{50}^2`$ ergs s<sup>-1</sup>.
Unfortunately, our limit on the X-ray emission from the radio galaxy is not very stringent. Fabbiano et al. (1984) studied the X-ray properties of several 3CR radio galaxies with the Einstein Observatory. They found that the FR II’s radio and X-ray luminosities are strongly correlated. Thus with a radio flux of 6.8 mJy at 5 GHz, NVSS 2146+82 should have a nuclear X-ray flux of a few times $`10^{42}`$ ergs s<sup>-1</sup>. This flux is comparable to our upper limit. Taking into account the intrinsic scatter in the radio/X-ray correlation, our non-detection of the AGN is quite reasonable.
Our upper limit on the X-ray emission from hot cluster gas provides a much stronger constraint. Most Abell richness class 0 clusters have X-ray luminosities of $`10^{4345}`$ erg/s (Ebeling et al. (1998)). Therefore any cluster of galaxies associated with the radio galaxy must be either intrinsically weak in X-rays or must be poorer than our optical estimate. Wan & Daly (1996) studied the X-ray emission of low-redshift FR II galaxies and found that poor clusters that contain FR II sources are underluminous in X-rays compared to similar clusters that do not contain FR IIs. The median X-ray luminosity for low-$`z`$ clusters with FR IIs was found to be $`1.3\times 10^{42}h_{50}^2`$ ergs s<sup>-1</sup> while it is $`1.33\times 10^{43}h_{50}^2`$ ergs s<sup>-1</sup> for a sample of low-$`z`$ clusters without FR IIs (Wan & Daly 1996). Assuming that the group surrounding NVSS 2146+82 is similar to that of other clusters found around low-$`z`$ FR IIs and is underluminous in X-rays, the optical richness estimate is probably correct.
## 5 Discussion
### 5.1 Physical Properties of the Radio Source
#### 5.1.1 Size and Luminosity
Our observations of NVSS 2146+82 clearly show that it is an unusually large FR II radio galaxy. Its angular distance from the north lobe to the south lobe gives an unusually large extent of $`\theta =19\stackrel{}{\mathrm{.}}5`$. For our assumed cosmology and our measured redshift of $`z=0.145`$, the linear extent of the radio structure is 4$`h_{50}^1`$ Mpc, placing it in the Giant Radio Galaxy (GRG) class, which we define as sources larger than 2$`h_{50}^1`$ Mpc. NVSS 2146+82 is therefore the second largest FR II known, surpassed only by 3C236 which is $``$6$`h_{50}^1`$ Mpc in extent. FR II galaxies of this size are extremely rare; a literature search by Nilsson et al. (1993) of 540 FR IIs contains only 27 objects with sizes greater than 1$`h_{50}^1`$ Mpc. Of this sample of 27 large FR IIs, only 5 are larger than 2$`h_{50}^1`$ Mpc. For comparison, the other known giant radio sources are shown in Table 6. The log radio luminosity of NVSS 2146+82 at 1.4 GHz is 25.69, in the middle of the range for giant radio sources.
It remains unclear if there are fundamental differences between GRGs and “normal” radio galaxies. The relative paucity of known GRGs may be in part due to observational selection effects in past radio surveys. An alternative reason for the rarity of giant radio galaxies may be that the physical conditions necessary for the creation of a GRG are uncommon in the universe. Although the similarity between NVSS 2146+82 and other FR IIs suggests that it is a typical FR II radio galaxy at the extreme end of the size distribution, a study of a complete sample of radio galaxies that includes GRGs will have to be made to determine if GRGs are part of a continuous distribution in size of normal radio galaxies or if there are fundamental differences between GRGs and smaller FR IIs.
#### 5.1.2 Equipartition calculations
If the usual equipartition assumptions are made, then it is possible to estimate the magnetic field strength and pressure in the lobes. Assuming that the observed spectral index is maintained from 10 MHz to 100 GHz, that there are equal energies in the radiating electrons and other particles, and that the filling factor is unity, the derived magnetic field is B$`{}_{min}{}^{}5\times 10^6h_{50}^{2/7}`$ Gauss and p$`{}_{min}{}^{}3.5\times 10^3h_{50}^{4/7}`$ cm<sup>-3</sup>K for the hot spots. At the midpoint of the lobes these values are B$`{}_{min}{}^{}8\times 10^7h_{50}^{2/7}`$ Gauss and p$`{}_{min}{}^{}2.3\times 10^2h_{50}^{4/7}`$ cm<sup>-3</sup>K. At this redshift, the 3 K microwave background has an equivalent magnetic field of 4.2$`\times 10^6`$ Gauss so the energy loss in the lobes should be dominated by inverse Compton scattering of this background, and the time for the electrons radiating at 1400 MHz to lose half of their energy will be $`10^8h_{50}^{3/7}`$ years.
#### 5.1.3 Magnetic Field and Faraday Rotation
The mean Faraday rotation of $``$ $`9`$ rad m<sup>-2</sup> shown in Figure 7 is consistent with the results of Simard-Normandin, Kronberg, & Button (1981) for other extragalactic sources seen through this region of the Galaxy ($`l=116.^{}7,b=21.^{}5`$). It is therefore likely that the rotation measure screen seen in Figure 7 is primarily the foreground screen of our Galaxy. The low apparent rotation measure and the smoothness of the polarization structure shown in Figure 6 suggests that the magnetic field in this source is well ordered. The field configuration is entirely typical of older extended FR II sources, with the E vectors lying approximately perpendicular to the ridge line of the radio emission in most features.
We note that the greater variance and evidence for organized structure in the Faraday rotation of the southern lobe is the opposite of what would be expected if the jet sidedness were due to Doppler favoritism and the Faraday rotating medium were local to the source. We think it more likely that the Faraday rotation structure arises along the line of sight in our Galaxy.
#### 5.1.4 Spectral Index Variations
The spectral index variations shown in Figure 9 indicate that there are regions $`2\stackrel{}{\mathrm{.}}4`$ back towards C from the brightest region in each lobe that have unusually flat spectra ($`\alpha _{0.35}^{1.4}`$ -0.3), flatter even than the hot spots. The only extended synchrotron sources known with spectra this flat are a few Galactic supernova remnants (Berkhuijsen (1986)).
The spectral index structure in NVSS 2146+82 is unlike the systematic steepening of the spectrum away from the hot spots that is usually interpreted as an effect of spectral aging in extended lobes. In such interpretations, electrons are presumed to be injected into a high field region in or around the hot spots, and their energy spectrum steepens with distance as they diffuse into lower field regions of the extended lobes. Clearly no such interpretation can be made here.
These flatter spectrum regions occur in the transition zone from the featureless parts of the lobes (closer to the core) to the parts near the regions of enhanced emission that contain significant filamentary structure. The anomalous regions are near the midline of the lobes; the southern region is centered on the path of the jet and the northern region is at one end of a prominent filament (the path of the jet is uncertain). The relative symmetry of the flatter spectrum regions of the lobes suggests that they might be produced by an intrinsic property of the source, such as a variable spectral index in the injection spectrum of the relativistic electrons from the jet, rather than local environmental effects.
If the magnetic field has values near those estimated by the equipartition calculations given above, then the energy loss of the radiating electrons is dominated by inverse Compton scattering against the Cosmic Microwave Background. In the low density, low magnetic fields in these lobes, the aging effects will be slow and the history of a variable electron spectrum could be maintained along the length of the lobe.
#### 5.1.5 Size Scales of Symmetry in the Radio Source
There are three size scales on which symmetries appear or change in the radio structure: The first is $`1\stackrel{}{\mathrm{.}}5=300h_{50}^1`$ kpc. The jets appear to become symmetric on this scale but are asymmetric on smaller scales. If the J2 and K components (Figure 4) are symmetric features in the jet and counterjet, any Doppler boosting from relativistic motion must have disappeared by this point in the jet. The second scale is $`3\stackrel{}{\mathrm{.}}2=640h_{50}^1`$ kpc. On this scale, there is a dramatic brightening of both lobes. The third scale is $`6\stackrel{}{\mathrm{.}}5=1300h_{50}^1`$ kpc. At this distance, the lobes become even brighter and strong filamentary structure appears. This is the distance at which regions of spectral anomaly appear in the extended emission.
The largest scale symmetries thus suggest a symmetric overall environment, apart from the slight non-collinearity (C-symmetry) of the structure. The small scale brightness asymmetries of the jet and counterjet might be attributed to Doppler boosting and dimming by relativistic motion which effectively disappears by $`300h_{50}^1`$ kpc, i.e. on a scale more typical of a “non-giant” FR II source. We reiterate however that the small asymmetry in rotation measure dispersion (variance) between the lobes is opposite in sign to that expected on this interpretation. This asymmetry seems more likely to reflect an intrinsic asymmetry (or gradient) in the foreground magnetoionic medium.
### 5.2 The Optical Environment
One possibility for the origin of GRGs is that they are otherwise normal FR II sources that reside in extremely low density gaseous environments. The environments in which radio galaxies reside have been studied in depth (e.g. Longair & Seldner (1979); Heckman et al. (1986); Prestage & Peacock (1988); Hill & Lilly (1991); Allington-Smith et al. (1993); Zirbel (1997)) because the gas density and pressure in the host galaxy’s ISM, any intracluster medium, and the IGM are at least partly responsible for determining the resulting radio morphology.
An intriguing result of recent studies (Hill & Lilly (1991); Allington-Smith et al. (1993), Zirbel (1997)) is that FR II galaxies are found in a range of cluster richnesses at moderate redshifts, but they are only found in poor to very poor groups at low redshift. The “richness” of the cluster associated with a radio galaxy can be estimated in a statistical sense in the absence of redshift data on nearby galaxies. Allington-Smith et al. (1993) define the richness parameter $`N_{0.5}^{19}`$ as the number of galaxies within a projected radius of 500 kpc and with $`M_V19.0`$ assuming the same redshift as the AGN. The number counts are corrected for contamination by foreground and background galaxies by subtracting number counts from a field offset from the radio galaxy. Zirbel (1997) gives a conversion of $`N_{0.5}^{19}`$ to Abell class as $`N_{Abell}=2.7(N_{0.5}^{19})^{0.9}`$. With this conversion, the thresholds for Abell Classes 0 and 1 are $`N_{0.5}^{19}=15`$ and 26 respectively. Using this richness estimation technique, Zirbel (1997) found that of a sample of 29 low redshift ($`z<0.2`$) FR IIs: (1) 41% of the sample of low $`z`$ FR IIs reside in very poor groups ($`N_{0.5}^{19}<3.5`$), and (2) more importantly, no low redshift FR II was found in a rich group with $`N_{0.5}^{19}>20`$. Based on the results given in §3.5, NVSS 2146+82 appears to reside in a group with an anomalously high galaxy richness compared to other low redshift FR IIs. Although the galaxy counts from the field surrounding NVSS 2146+82 were not calculated identically to those of Zirbel (1997), the value of $`N_{0.5}^{19}`$ is likely $`>2530`$ for NVSS 2146+82.
The upper limit on the cluster X-ray emission is consistent with the NVSS 2146+82 group being at the low end of the X-ray luminosity distribution for poor clusters. Wan & Daly (1996) found that in a comparison of low redshift clusters with and without FR II sources, clusters that contained FR IIs were underluminous in X-rays compared to clusters without FR IIs. Although the cluster surrounding NVSS 2146+82 may be Abell Class 0, its lack of associated X-ray gas suggests that the pressure in the surrounding medium is low enough for a giant radio source to form with little disruption of the FR II jet.
Curiously, several other GRGs listed in Table 6 also appear to lie in regions with overdensities of nearby galaxies. The GRG 0503-286 appears to lie in a group of 30 or so galaxies (Saripalli et al. (1986)). These companions are concentrated to the northeast of the host galaxy of 0503-286, and may have caused the asymmetric appearance of the northern lobe of the radio structure. Overdensities of nearby galaxies are also reported for 1358+305 (Parma et al. (1996)) and 8C 0821+695 (Lacy et al. (1993)); however, in both cases there is no spectroscopic confirmation of the redshifts of the candidate cluster galaxies. In a recent study of the optical and X-ray environments of radio galaxies, Miller et al. (1999) find that for a sample of FR I sources, all have extended X-ray emission and overdensities of optical galaxies. However of their sample of seven FR II sources, none have overdensities of optical galaxies or extended X-ray emission except for the GRG DA240, which has no extended X-ray emission but does have a marginally significant excess of optical companions. Perhaps for at least some of the GRG population, the presence of the host galaxy in an optically rich group with little associated X-ray gas is related to the formation or evolution of the radio source?
## 6 Summary and Conclusions
We have presented multi-wavelength observations of the unusually large FR II radio galaxy NVSS 2146+82. The overall size of the radio source is $`4h_{50}^1`$ Mpc, making it the second largest known FR II source. We have found the host galaxy to be similar in both luminosity and morphology to a sample of other low redshift FR II galaxies. Emission line profiles seen in the spectrum of the host galaxy are double peaked, which may indicate that the ionized gas may be being accelerated by the bipolar radio jet.
We have also found evidence for an anomalously rich group of galaxies at the same redshift as NVSS 2146+82 that has little associated X-ray emitting gas. Though unusual in having a rich environment, this source is similar to other low redshift FR IIs in clusters; the NVSS 2146+82 group is underluminous in X-rays compared to clusters of similar richness that contain no FR II. The large radio size, lack of significant Faraday rotation and non detection of X-rays all suggest that in spite of the richness of the cluster in which this galaxy resides, it has a low gas density.
There is some morphological evidence that the host galaxy of NVSS 2146+82 may be undergoing tidal interaction with one or more of its nearest companions. Also, an interaction may be responsible for the double-peaked emission line profiles, however the spatial resolution of the spectrum of the nucleus is not high enough to distinguish between a merger origin or radio jet/cloud interaction origin for the peculiar profiles.
Apart from the radio spectral index anomaly, the radio properties of this source are like a normal FR II source scaled up by a factor of ten, preserving the standard overall morphology and polarization structure. In the outer regions of the source the magnetic field is likely to be so weak that inverse Compton losses to the Cosmic Microwave Background dominate synchrotron losses.
We are grateful to Mark Whittle for many helpful conversations. We are grateful to Matt Bershady, Randy Phelps, and Mike Siegel for either sharing observing time or taking observations in support of this research. CP acknowledges the support of a Grant-in-aid of Research from Sigma Xi, the Scientific Research Society. This research has made use of the NASA/IPAC Extragalactic Database (NED) which is operated by the Jet Propulsion Laboratory, California Institute of Technology, under contract with the National Aeronautics and Space Administration. We acknowledge the use of NASA’s SkyView facility (http://skyview.gsfc.nasa.gov) located at NASA Goddard Space Flight Center. |
warning/0002/math-ph0002045.html | ar5iv | text | # The Standard Model – the Commutative Case: Spinors, Dirac Operator and de Rham Algebra
## 1. The theorems by Gel’fand and Serre-Swan
One of the corner stones of the beginning of noncommutative geometry was I. M. Gel’fand’s theorem published in 1940. He established an equivalence principle between some topological objects and algebraic-axiomatic structures that can be expressed in the following way (cf. ):
###### Theorem 1.1.
(I. M. Gel’fand)
Let $`A`$ be a commutative C\*-algebra and $`X`$ the set of its characters. The topology on $`X`$ should be that one induced by the weak\* topology on the dual space $`A^{}`$.
Then $`X`$ is a locally compact Hausdorff space, and $`X`$ is compact iff $`A`$ is unital.
The C\*-algebra $`A`$ is $``$-isomorphic to the commutative C\*-algebra $`C_0(X)`$ of all continuous functions on $`X`$ vanishing at infinity.
In a more contemporary language this bijection can be expressed as a categorical equivalence. We have to add a set of suitable morphisms to the sets of objects ’commutative C\*-algebras’ and ’locally compact Hausdorff spaces’. They are called proper morphisms: for C\*-algebras we have to take $``$-homomorphisms that map approximate identities to approximate identities, and for locally compact Hausdorff spaces we have to select those continuous maps for which the pre-image of a compact set is always compact. Then we can summarize the categorical equivalence:
| commutative C\*-algebras $`C(X)`$ | $``$ | locally compact Hausdorff spaces $`X`$ |
| --- | --- | --- |
| proper $``$-homomorphisms | | proper continuous homomorphisms |
The noncommutative viewpoint enters the picture removing the commutativity condition on the multiplication in C\*-algebras. Algebraically the left side is still a proper category, and many theorems for commutative C\*-algebras can be generalized to the noncommutative situation. (But, there are also pure noncommutative structures like those described by Tomita-Takesaki theory.) However, the right side possesses no obvious candidate for a counterpart of the left side generalization to preserve the categorical equivalence. One reason is that the notion of a point that is crucial for any geometry becomes a vacuous notion under such an extension of the theory. Consequently, what we are left with is the algebraic noncommutative picture on the left side.
Looking for further topological and geometrical structures that can be categorically replaced by appropriate algebraic structures J.-P. Serre (1957/58) and R. G. Swan (1962) independently established a categorical equivalence between projective finitely generated $`C(X)`$-modules and locally trivial vector bundles over $`X`$ for compact Hausdorff spaces $`X`$. To describe it in greater detail some preparation is necessary.
To introduce both the notions, first, define a (left) unital $`A`$-module $``$ over a unital algebra to be projective finitely generated if it is a direct summand (in an $`A`$-module sense) of a free $`A`$-module $`A^n`$ for $`n`$, where $`A^n`$ consists of all $`n`$-tuples of elements of $`A`$ equipped with coordinate-wise addition and an action of $`A`$ on $`A^n`$ given as (left) multiplication of any $`n`$-tuple entry by fixed elements of $`A`$. The set of projective finitely generated $`A`$-modules can be equipped with the structure of direct sums $``$ of $`A`$-modules. To introduce the structure of a module tensor product we have to consider them as $`A`$-bimodules defining another (right) action of $`A`$ on $`A^n`$ as a (right) multiplication of any $`n`$-tuple entry by fixed elements of $`A`$. The module tensor product $`_1_A_2`$ is the algebraic tensor product of the linear spaces $`_1`$, $`_2`$ factored by the module ideal generated as the linear hull
$$\mathrm{Lin}\{h_1ah_2h_1ah_2:h_1_1,h_2_2,aA\}.$$
We have associative and distributive laws for the addition and the tensor products and a commutative law for addition. The neutral element of addition is the $`A`$-module consisting only of the zero element, and the neutral element of the module tensor product is the bimodule $`A^1=A`$.
As the set of homomorphisms we consider all $`A`$-(bi-)module homomorphisms of projective finitely generated $`A`$-(bi-)-modules.
The second structure involved in the stressed for categorical equivalence consists of locally trivial vector bundles over compact Hausdorff spaces $`X`$.
###### Definition 1.2.
Given a topological space $`E`$, a compact Hausdorff space $`X`$ and a continuous mapping $`p:EX`$. Then $`E`$ is a locally trivial vector bundle $`(E,p,X)`$ over $`X`$ if for every $`xX`$ there exists a finite-dimensional vector space $`E_x`$ (equipped with the Euclidean topology) and a neighborhood $`U_xX`$ such that a homeomorphism $`\varphi :U_x\times E_xp^1(U_x)`$ exists and $`p\varphi (x,e)=x`$ for any $`xX`$. In case $`U_xX`$ the vector bundle is (globally) trivial.
A map $`\varphi :(E,p,X)(F,q,X)`$ is a vector bundle isomorphism in case $`\varphi `$ is bijective, $`\varphi `$ and $`\varphi ^1`$ are continuous and $`\varphi |_{E_x}:E_xF_x`$ is linear for any $`xX`$. The map $`\varphi `$ is a vector bundle homomorphism if $`\varphi `$ is continuous and $`\varphi |_{E_x}:E_xF_x`$ is a linear embedding as a subspace.
We call $`X`$ the base space, $`E`$ the total space, $`E_x=p^1(\{x\})`$ the fibre over $`x`$ and $`p`$ the projection map.
Note, that the compactness of $`X`$ implies $`sup(\mathrm{dim}(E_x))<\mathrm{}`$. As one of the alternative descriptions of vector bundles in geometry we can describe them in local terms : a vector bundle $`(E,p,X)`$ is given by an atlas $`\{U_\alpha \}X`$ of (open) charts and of coordinate homeomorphisms $`\{f_\alpha :U_\alpha \times E_xp^1(U_\alpha )\}`$ ($`xU_\alpha `$) such that the transition functions
$$f_{\alpha \beta }:=f_\beta ^1f_\alpha :(U_\alpha U_\beta )\times E_x(U_\alpha U_\beta )\times E_x$$
are described by $`f_{\alpha \beta }(x,e)=(x,\overline{f_{\alpha \beta }}(x)e)`$ with continuous functions $`\overline{f_{\alpha \beta }}:U_\alpha U_\beta GL(n,)`$ fulfilling the law
$$\overline{f_{\alpha \alpha }}=\mathrm{id}_{U_\alpha },\overline{f_{\alpha \gamma }}\overline{f_{\gamma \beta }}\overline{f_{\beta \alpha }}=\mathrm{id}_{U_\alpha U_\beta U_\gamma }.$$
We can show that the condition $`\overline{f_{\alpha \beta }}:U_\alpha U_\beta GL(n,)`$ can be always reduced to $`\overline{f_{\alpha \beta }}:U_\alpha U_\beta U(n)`$ (or, for real vector spaces, $`\overline{f_{\alpha \beta }}:U_\alpha U_\beta O(n)`$) changing the coordinate functions in a suitable way, cf. . The group $`U(n)`$ (or $`O(n)`$) is said to be the structural group of the vector bundle.
For further use we introduce the notion of an orientation on vector bundles over orientable compact manifolds.
###### Definition 1.3.
Let $`M`$ be an orientable compact manifold. The vector bundle $`(E,p,M)`$ is orientable if there exists an atlas $`\{U_\alpha \}`$ describing $`E`$ with transition functions $`\{\overline{f_{\alpha \beta }}\}`$ taking values in $`GL^+(n,)`$. The corresponding atlas is said to be an orientation of the vector bundle $`(E,p,M)`$.
For a fixed compact Hausdorff space $`X`$ the set of vector bundles with base space $`X`$ can be equipped with some algebraic structure. The Whitney sum of two vector bundles $`(E,p,X)`$ and $`(F,q,X)`$ is the vector bundle $`(EF,pq,X)`$, where
$`EF`$ $`:=`$ $`\{(e,f)E\times F:p(e)=q(f)X\},`$
$`(pq)(e,f)`$ $`:=`$ $`p(e)=q(f)X.`$
Local triviality is preserved under Whitney sums. The fibres are the vector spaces $`E_xF_x`$. The tensor product of two vector bundles $`(E,p,X)`$ and $`(F,q,X)`$ is the vector bundle $`(EF,pq,X)`$ with the fibres $`E_xF_x`$ for $`xX`$ and the transition functions $`\overline{f_{\alpha \beta }}(x):=\overline{f_{\alpha \beta ,E}}(x)\overline{f_{\alpha \beta ,F}}(x)`$ coming from a common atlas $`\{U_\alpha \}X`$ of the vector bundles $`(E,p,X)`$ and $`(F,p,X)`$. We observe that for trivial vector bundles $`X\times ^n=:\overline{n}`$ the two operations are related by the isomorphism $`E\overline{n}=_{i=1}^nE_{(i)}`$, where $`n`$ is arbitrary. Concerning the algebraic properties of the two operations both they are associative and fulfil the obvious distributivity laws, and the Whitney addition is commutative in the sense of an appropriate isomorphism of vector bundles. The neutral elements are $`\overline{0}`$ and $`\overline{1}`$, respectively.
One of the central observations is Swan’s theorem:
###### Theorem 1.4.
(R. G. Swan, 1962)
Let $`(E,p,X)`$ be a locally trivial vector bundle over a compact Hausdorff base space $`X`$. There exists a locally trivial vector bundle $`(F,q,X)`$ over $`X`$ such that $`(EF,pq,X)`$ is trivial (with finite-dimensional fibre).
The proof is elaborated, and we refer to R. G. Swan’s paper or to for different versions of proofs.
###### Definition 1.5.
A section in a vector bundle $`(E,p,X)`$ is a continuous map $`s:XE`$ such that $`(ps)(x)=x`$ for every $`xX`$. The set of sections of $`(E,p,X)`$ is denoted by $`\mathrm{\Gamma }(E)`$.
###### Proposition 1.6.
Let $`X`$ be a compact Hausdorff space. Every locally trivial vector bundle admits non-trivial sections. For every vector bundle $`(E,p,X)`$ the set $`\mathrm{\Gamma }(E)`$ has the algebraic structure of a $`C(X)`$-module.
Any isomorphism of vector bundles induces an isomorphism of the corresponding modules of sections. Whitney sums of vector bundles correspond to direct $`C(X)`$-module sums of the related modules of sections, tensor products of vector bundles correspond to bimodule tensor products.
For compact $`X`$ the $`C(X)`$-module $`\mathrm{\Gamma }(E)`$ is projective and finitely generated, in particular, $`\mathrm{\Gamma }(X\times ^n)C(X)^n`$ for every $`n`$.
###### Proof.
The existence of continuous sections can be proved applying Uryson’s Lemma to constant sections in the (trivial) part of the vector bundle over one chart $`U`$, getting continuous sections of the whole vector bundle supported in one chart $`U`$ over which the vector bundle is trivial.
Any bundle homomorphism $`\varphi :(E,p,X)(F,q,X)`$ maps sections in $`E`$ to sections in $`F`$. If $`\varphi `$ is a bundle isomorphism, then $`\varphi _{}:\mathrm{\Gamma }(E)\mathrm{\Gamma }(F)`$ is a $`C(X)`$-module isomorphism.
We observe that $`\mathrm{\Gamma }(X\times ^n)C(X)^n`$. These $`C(X)`$-modules are free and finitely generated. Since $`C(X)^n\mathrm{\Gamma }(EF)\mathrm{\Gamma }(E)\mathrm{\Gamma }(F)`$ for a given vector bundle $`(E,p,X)`$, some vector bundle $`(F,p,X)`$ and $`n<\mathrm{}`$ by Swan’s theorem, $`\mathrm{\Gamma }(E)`$ is projective and finitely generated. ∎
###### Theorem 1.7.
(J.-P. Serre, 1957/58, R. G. Swan, 1962)
Let $`X`$ be a compact Hausdorff space and $``$ be a finitely generated projective $`C(X)`$-module.
If $`𝒢C(X)^n`$ for some $`n<\mathrm{}`$, then let $`P`$ be the projection of $`C(X)^n`$ onto $``$ along $`𝒢`$. Interpreting $`P`$ as an element of $`M_n(C(X))C(X,M_n())`$ define
$$\mathrm{\Xi }():=\{(x,e)X\times ^n:e\mathrm{ran}(P)\}.$$
Then $`\mathrm{\Xi }()`$ is a locally trivial vector bundle over $`X`$, $`\mathrm{\Gamma }(\mathrm{\Xi }())`$. Moreover, if $`=\mathrm{\Gamma }(E)`$ for some vector bundle $`E`$, then $`\mathrm{\Xi }(\mathrm{\Gamma }(E))E`$.
###### Proof.
$`\mathrm{\Xi }(\mathrm{\Gamma }(E))E`$: Assume $`EFX\times ^n`$ by Swan’s theorem. Let $`\pi _x:^nE_x`$ be the fibrewise projection, $`xX`$. Define $`\pi :X\times ^nE`$ by $`\pi (x,e)=(x,\pi _x(e))`$ for $`xX`$, $`e^n`$. Then $`\pi `$ is a correctly defined surjective bundle homomorphism.
Let $`P=\pi _{}:\mathrm{\Gamma }(X\times ^n)\mathrm{\Gamma }(E)\mathrm{\Gamma }(F)`$ be the induced $`C(X)`$-module map, a projection onto $`\mathrm{\Gamma }(E)`$. Note, that $`P(x)=\pi _x`$ for every $`xX`$. Therefore, $`E=\mathrm{\Xi }(\mathrm{\Gamma }(E))`$ by construction.
$`\mathrm{\Gamma }(\mathrm{\Xi }())=`$: Note, that $`(\mathrm{\Xi }())_x=\{x\}\times \{e^n:e\mathrm{ran}(P(x))\}`$ are the fibres of $`\mathrm{\Xi }()`$. The family of projections $`\{P(x)\}`$ is continuous, and $`\mathrm{\Xi }()`$ becomes a locally trivial vector bundle. Thus, $`\mathrm{\Gamma }(\mathrm{\Xi }())=\{fC(X,^n):f\mathrm{ran}(P)=\}`$. ∎
Formulating the result in a categorical language we obtain a categorical equivalence between an algebraic and a geometric category if suitable sets of $`C(X)`$-module and bundle homomorphisms are chosen:
| projective, finitely generated $`C(X)`$-modules | $``$ | locally trivial vector bundles $`(E,p,X)`$ |
| --- | --- | --- |
| proper $`C(X)`$-module maps | | proper bundle homomorphisms |
We would like to point out that this categorical equivalence can be extended to the situation of infinite-dimensional fibres, however we will lose local triviality of the Banach bundles if we try to preserve a suitable category of $`C(X)`$-modules like Banach or Hilbert $`C(X)`$-modules on the left side. Moreover, most locally trivial bundles over compact Hausdorff spaces $`X`$ with fibre $`l_2`$ turn out to be automatically globally trivial.
Now, we specify the compact Hausdorff space $`X`$ to be a compact smooth manifold $`M`$. The observation to be made is that every locally trivial vector bundle over $`M`$ with continuous transition functions in some atlas is in fact equipped with an atlas containing smooth transition functions, i.e. there is no reason to distinguish between ’continuous’ and ’smooth’ vector bundles over smooth compact manifolds $`M`$, cf. for a proof.
###### Lemma 1.8.
For every vector bundle $`(E,p,M)`$ there exists an atlas on $`M`$ such that $`E`$ is trivial over every chart $`U_\alpha `$ and the transition functions $`\overline{f_{\alpha \beta }}:U_\alpha U_\beta GL(n,)`$ are smooth functions.
The Fréchet algebra $`C^{\mathrm{}}(M)`$ and the C\*-algebra $`C(M)`$ have the same set of characters: every character on $`C^{\mathrm{}}(M)`$ is automatically continuous and a measure on $`M`$ and hence, a character of $`C(M)`$. Consequently, $`C^{\mathrm{}}(M)^n\mathrm{\Gamma }^{\mathrm{}}(M\times ^n)`$, and the categorical equivalence between projective $`C^{\mathrm{}}(M)`$-modules and vector bundles over $`M`$ is a reduction of Serre-Swan’s categorical equivalence. For the Whitney sum and the bundle tensor product we get the following corresponding module operations on the $`C^{\mathrm{}}(M)`$-modules of smooth sections:
$`\mathrm{\Gamma }^{\mathrm{}}(EF)`$ $`=`$ $`\mathrm{\Gamma }^{\mathrm{}}(E)_{C^{\mathrm{}}(M)}\mathrm{\Gamma }^{\mathrm{}}(F),`$
$`\mathrm{\Gamma }^{\mathrm{}}(EF)`$ $`=`$ $`\mathrm{\Gamma }^{\mathrm{}}(E)_{C^{\mathrm{}}(M)}\mathrm{\Gamma }^{\mathrm{}}(F).`$
## 2. Hermitean structures and frames for sets of sections
As an essential tool we need the existence and the properties of a continuous field of scalar products on the fibres of vector bundles. This structure is not needed to prove the Serre-Swan’ theorem, it arises additionally.
###### Definition 2.1.
Let $`X`$ be a compact Hausdorff space and $`(E,p,X)`$ be a vector bundle with base space $`X`$. A $`C(X)`$-valued inner product on $`(E,p,X)`$ is a bilinear mapping $`.,.:\mathrm{\Gamma }(E)\times \mathrm{\Gamma }(E)C(X)`$ that is continuous in both the arguments, acts fibrewise (i.e. is $`C(X)`$-linear in the first argument) and its restriction to any fibre $`E_x`$ generates a scalar product on it. (Some authors refer to this structure as to a Hermitean structure on the vector bundle.)
###### Theorem 2.2.
Let $`X`$ be a compact Hausdorff space and $`(E,p,X)`$ be a vector bundle with base space $`X`$. Then $`(E,p,X)`$ admits $`C(X)`$-valued inner products $`.,.`$ on the $`C(X)`$-module $`\mathrm{\Gamma }(E)`$ such that $`\mathrm{\Gamma }(E)`$ is complete with respect to the resulting norm $`.:=.,.^{1/2}`$.
Any two $`C(X)`$-valued inner products $`.,._1`$, $`.,._2`$ are related by a positive invertible $`C(X)`$-linear operator $`S`$ on $`\mathrm{\Gamma }(E)`$ via the formula $`.,._1S(.),._2`$.
If $`X`$ is a smooth manifold then $`.,.`$ can be chosen in such a way that its restriction to $`\mathrm{\Gamma }^{\mathrm{}}(E)\times \mathrm{\Gamma }^{\mathrm{}}(E)`$ takes values in $`C^{\mathrm{}}(X)`$.
###### Proof.
Because of the categorical equivalence by J.-P. Serre and R. G. Swan it is sufficient to indicate the existence and the properties of $`C(X)`$-valued inner products on finitely generated projective $`C(X)`$\- or $`C^{\mathrm{}}(X)`$-modules. For $`C(X)^n`$ the $`C(X)`$-valued inner product is defined as $`(f_1,\mathrm{},f_n),(g_1,\mathrm{},g_n)=_{i=1}^nf_i\overline{g_i}`$. For direct summands $`P(C(X)^n)`$ of $`C(X)^n`$ we reduce this $`C(X)`$-valued inner product to elements of them.
The relation between two $`C(X)`$-valued inner products follows from an analogue of Riesz’ representation theorem for $`C(X)`$-linear bounded module maps from $`C(X)^n`$ into $`C(X)`$. (Attention: This may fail for more general $`C(X)`$-modules with $`C(X)`$-valued inner products.)
If $`X`$ is a smooth manifold, then the $`C(X)`$-valued inner product defined above maps elements with smooth entries to smooth functions on $`X`$. A perturbation of the $`C(X)`$-valued inner product by a positive invertible operator $`S`$ that preserves the range $`C^{\mathrm{}}(X)`$ of it or the restriction to a direct summand of $`C^{\mathrm{}}(X)^n`$ do not change this fact. ∎
We would like to remark that for more general $``$-algebras $`𝒜C^{\mathrm{}}(X)`$ that are closed under holomorphic calculus and contain the identity the property of $`𝒜`$-valued inner products on the correspondingly reduced set of sections $`\mathrm{\Gamma }^𝒜(E)\mathrm{\Gamma }^{\mathrm{}}(E)`$ to possess an analogue of the Riesz’ property has to be axiomatically supposed, in general.
Now, we indicate the existence of finite sets of generators of $`\mathrm{\Gamma }^{\mathrm{}}(E)`$ as a $`C^{\mathrm{}}(M)`$-module for vector bundles $`(E,p,M)`$ over smooth manifolds $`M`$. Consider the free $`C^{\mathrm{}}(M)`$-module $`\mathrm{\Gamma }^{\mathrm{}}(M\times ^n)=C^{\mathrm{}}(M)^n`$ for $`n`$ and a $`C^{\mathrm{}}(M)`$-valued inner product $`.,._0`$ on it. Then there exists an orthonormal with respect to $`.,._0`$ basis consisting of $`n`$ elements of this module. Indeed, on free $`C(M)`$-modules $`C(M)^n`$ every $`C(M)`$-valued inner product is related to the canonical $`C(M)`$-valued inner product by a bounded invertible positive module operator $`S`$ that fulfills the identity $`.,._{can.}S(.),.`$. The restriction of $`.,._{can.}`$ to $`C^{\mathrm{}}(M)^n`$ is $`C^{\mathrm{}}(M)`$-valued, and $`.,._0`$ can be extended to $`C(M)^n`$. So the linking operator $`S`$ exists on $`C(M)^n`$, and its restriction to $`C^{\mathrm{}}(M)^n`$ maps smooth elements to smooth elements. However, the canonical $`C^{\mathrm{}}(M)`$-valued inner product on $`C^{\mathrm{}}(M)^n`$ admits an orthonormal basis consisting of smooth elements:
$$\{e_1,\mathrm{},e_n:e_i=(0,\mathrm{},0,1_{(i)},0,\mathrm{},0)\}.$$
Consequently, $`\{S^{1/2}(e_i):i=1,\mathrm{},n\}`$ is an orthonormal basis of $`C^{\mathrm{}}(M)^n`$ with respect to the given $`C^{\mathrm{}}(M)`$-valued inner product $`.,._0`$.
Let $``$ be a projective finitely generated $`C^{\mathrm{}}(M)`$-module, i.e. $`=C^{\mathrm{}}(M)^n`$ for a finite integer $`n`$. Denote by $`P`$ the $`C^{\mathrm{}}(M)`$-linear projection onto $``$ along $``$. Then the set $`\{P(e_i):i=1,\mathrm{},n\}`$ of elements of $``$ has the remarkable property that
$$\xi =\underset{i=1}{\overset{n}{}}\xi ,P(e_i)_0P(e_i)$$
for every $`\xi `$. The engeneering literature on wavelets calls such sets of generators of Hilbert spaces (normalized tight) frames, whereas the literature on conditional expectations calls them quasi-bases or (module) bases. The notion ’basis’ is, however, misleading since the elements of the generator sequence $`\{P(e_i):i=1,\mathrm{},n\}`$ may allow a non-trivial $`C^{\mathrm{}}(M)`$-linear decomposition of the zero element of $``$. To see that let $``$ be simply the subset of all elements admitting only allover equal entries in their $`n`$-tuple representation. For more details we refer the reader to . To summarize the arguments we formulate
###### Theorem 2.3.
Let $`M`$ be a smooth compact manifold and $`(E,p,M)`$ be a vector bundle with base space $`M`$. Let $`.,.`$ be a Hermitean structure on it. Then the projective finitely generated $`C^{\mathrm{}}(M)`$-module $`\mathrm{\Gamma }^{\mathrm{}}(E)`$ possesses a finite subset $`\{\eta _i:i\}`$ such that $`\mathrm{\Gamma }^{\mathrm{}}(E)`$ is generated as a $`C^{\mathrm{}}(M)`$-module by this set and the equality
$$\xi =\underset{i=1}{\overset{n}{}}\xi ,\eta _i\eta _i$$
is satisfied for every $`\xi \mathrm{\Gamma }^{\mathrm{}}(E)`$.
## 3. Clifford and spinor bundles, spin manifolds
Let $`(M,g)`$ be a smooth Riemannian manifold, where the Riemannian metric $`g_x`$ induces a scalar product on $`T_xM`$ for any $`xM`$. Note, that the tangent space $`T_xM`$ and the cotangent space $`T_x^{}M`$ are isomorphic via the scalar product on $`T_xM`$ for any $`xM`$. If $`(T_xM,g_x)`$ denotes the Hilbert tangent space then let $`(T_x^{}M,g_x^1)`$ denote the resulting Hilbert cotangent space.
Let $`Cl(T_xM,g_x)`$ be the real Clifford algebra of the tangent space $`T_xM`$ with respect to the scalar product induced by the Riemannian metric $`g_x`$, $`xX`$ arbitrarily fixed. This algebra is defined to be a quotient of the tensor algebra $`𝒯(T_xM)`$ generated by the linear space $`T_xM`$, i.e. of
$$𝒯(T_xM)=T_xM(T_xMT_xM)\mathrm{}(T_xM\mathrm{}T_xM)\mathrm{}.$$
More precisely,
$$Cl(T_xM,g_x):=𝒯(T_xM)/\mathrm{Ideal}(eeg_x(e,e):eT_xM).$$
The real Clifford algebra $`Cl(T_xM,g_x)`$ possesses a $`_2`$-grading induced by the map $`\chi _x:(x,e)T_xM(x,e)T_xM`$, i.e. by the linear operator $`\chi `$ on $`Cl(T_xM,g_x)`$ with the property $`\chi ^2=\mathrm{id}`$, with eigen-values $`\{1,1\}`$ and isomorphic eigenspaces $`Cl^{even}(T_xM,g_x)`$, $`Cl^{odd}(T_xM,g_x)`$ summing up to the algebra itself. The words ’even’ and ’odd’ refer to the highest degree of the element under consideration and its property to be an even or odd number. If $`n=2m+1`$ then $`\chi `$ is realized as a multiplication by a central element Extending the isomorphism between tangent space and cotangent space via the scalar product $`g_x`$ on the first space we obtain a canonical algebraic isomorphism of the Clifford algebras $`Cl(T_xM,g_x)`$ and $`Cl(T_x^{}M,g_x)`$ for any fixed $`xX`$.
The Clifford algebra bundle $`l(M)`$ is defined fibrewise using the atlas on $`M`$ induced by the tangent bundle atlas of $`TM`$ (or the cotangent bundle atlas of $`T^{}M`$):
$$l_x(M):=Cl(T_xM,g_x)_{}\stackrel{\tau }{}\{\begin{array}{ccc}M_{2^m}()\hfill & :& n=2m\hfill \\ M_{2^m}()M_{2^m}()\hfill & :& n=2m+1\hfill \end{array}.$$
Note that the isomorphism $`\tau `$ is quite complicated, and in case $`n=2m+1`$ it maps both the even and the odd part of the Clifford algebra to both the blocks of the matrix sum at the right (see \[12, p. 15\] for details). The Clifford bundle possesses a $`_2`$-grading induced from that one on its fibres. The C\*-algebra structure of $`\mathrm{\Gamma }(l(M))`$ comes from the algebra structure of the Clifford algebra fibres and from the involution induced by $`_{}`$ from $``$. The C\*-norm exists and is uniquely defined since the multiplication and the involution are given and every fibre is finite-dimensional.
Some authors (cf. ) prefer to restrict the Clifford algebra bundle to the even part in case the dimension of the manifold is $`n=2m+1`$. The loss of that alternative definition is the $`_2`$-grading. The advantage of that approach is the structure of $`l(M)`$ as a continuous field of simple C\*-algebras allowing the attempt to interpret this bundle as a homomorphism bundle derived from some other vector bundle with base space $`M`$. We prefer to postpone this reduction until the spinor bundle has to be built up.
Consider either the Clifford algebra bundle $`l(M)`$ over $`M`$ for $`n=2m`$ or the first matrix block part $`l(M)^{}`$ of the Clifford algebra bundle $`l(M)`$ over $`M`$ for $`n=2m+1`$ (in its matrix representation) locally: for every $`xX`$ we find vector spaces $`S_x`$ such that the (first part of the) Clifford algebra bundle is locally isomorphic to the trivial homomorphism bundle $`\mathrm{Hom}(S_x)`$ of the trivial bundle $`(U\times S_x,p_U,U)`$. The C\*-algebra structure on $`\mathrm{\Gamma }(l_x(M))`$ (resp., $`\mathrm{\Gamma }(l_x(M)^{})`$) induces a unique scalar product on $`S_x`$ compatible with it. The dimension of the linear spaces $`S_x`$ is constant and equals $`\mathrm{dim}(S_x)=2^m`$ for any manifold dimensions $`n`$, $`m:=[n/2]`$.
Whether we can glue these trivial pieces together to obtain a vector bundle $`S`$ over the compact Riemannian manifold $`M`$ carrying an irreducible left action of the Clifford bundle (resp., the first part of it) that acts locally in the manner described, or not? Unfortunately, not always. If $`n=2m`$ the Clifford bundle $`l(M)`$ serves as a homomorphism bundle for some other vector bundle with the same compact base space $`M`$ if and only if the Dixmier-Douady class $`\delta (l(M))H^3(M,)`$ equals zero, where $`\delta (l(M))`$ also equals the third integral Stiefel-Whitney class $`w_3(TM)H^3(M,)`$. If $`n=2m+1`$ the first part $`l(M)^{}`$ of the Clifford bundle is a homomorphism bundle of some other vector bundle if and only if the second part of it does so, if and only if the Dixmier-Douady class $`\delta (l(M)^{})H^3(M,)`$ equals zero, where $`\delta (l(M)^{})`$ also equals the third integral Stiefel-Whitney class $`w_3(TM)H^3(M,)`$. The first fact was observed by J. Dixmier in \[8, Th. 10.9.3\] and again investigated in connection with spinor bundles by R. J. Plymen .
To formulate the definition of a spinor bundle on a given compact smooth Riemannian manifold $`M`$ or, equivalently, the definition of the property of $`M`$ to be a $`\mathrm{Spin}^{}`$-manifold we have to introduce the notion of Morita equivalence of certain unital $``$-algebras. We will do that only for the two $``$-algebras of interest, for more general cases we refer to . Let us fix the unital $``$-algebra
$$B=\{\begin{array}{ccc}C^{\mathrm{}}(M,l(M))=\mathrm{\Gamma }^{\mathrm{}}(l(M))\hfill & :& n=2m\hfill \\ C^{\mathrm{}}(M,l(M)^{})=\mathrm{\Gamma }^{\mathrm{}}(l(M)^{})\hfill & :& n=2m+1\hfill \end{array}.$$
###### Definition 3.1.
Let $`M`$ be a compact smooth Riemannian manifold. Consider the unital $``$-algebras $`A=C^{\mathrm{}}(M)`$ and $`B`$. They are Morita-equivalent as algebras if there exists a $`B`$-$`A`$ bimodule $``$ and an $`A`$-$`B`$ bimodule $``$ such that $`_AB`$ and $`_BA`$ as $`B`$\- and $`A`$-bimodules, respectively.
In our case $``$ can be chosen to be a projective and finitely generated (left) module over the unital $``$-algebra $`A=C^{\mathrm{}}(M)`$ denoted by $`\stackrel{~}{𝒮}`$. As a projective finitely generated $`C^{\mathrm{}}(M)`$-module $`\stackrel{~}{𝒮}`$ admits a $`C^{\mathrm{}}(M)`$-valued inner product $`.,._{C^{\mathrm{}}(M)}`$. Then $`B`$ can be realized as the $``$-algebra of bounded module operators over $`\stackrel{~}{𝒮}`$ generated as $`\mathrm{Lin}\{\xi ,\eta _{C^{\mathrm{}}(M,l(M))}:\xi ,\eta \stackrel{~}{𝒮}\}`$, where
$$\xi ,\eta _{C^{\mathrm{}}(M,l(M))}(\nu ):=\nu ,\xi _{C^{\mathrm{}}(M)}\eta \mathrm{for}\nu \stackrel{~}{𝒮}.$$
So the counterpart $``$ of $``$ can be described as the set $`\{\overline{\xi }:\xi ,\overline{\xi }a:=\overline{(a^{}\xi )},aA\}`$. Obviously, the right action of $`B`$ on $``$ is simultaneously swept to a left $`B`$-action on $``$. The $`C^{\mathrm{}}`$-module $``$ together with the $`C^{\mathrm{}}(M)`$-valued inner product $`.,._{C^{\mathrm{}}(M)}`$ is said to be a $`B`$-$`A`$ imprimitivity bimodule.
###### Definition 3.2.
(R. J. Plymen, 1982)
Let $`(M,g)`$ be a compact smooth Riemannian manifold, let $`A=C^{\mathrm{}}(M)`$ and $`B`$ as defined above in dependency on the dimension of $`M`$. Both $`A`$ and $`B`$ are unital $``$-algebras of smooth mappings.
We say that the tangent bundle $`TM`$ of $`M`$ admits a $`\mathrm{Spin}^{}`$-structure if $`TM`$ is orientable as a vector bundle and the Dixmier-Douady class $`\delta (l(M))`$ equals zero for $`n=2m`$ or, respectively, $`\delta (l(M)^{})=0`$ for $`n=2m+1`$.
If this condition is fulfilled then the $`\mathrm{Spin}^{}`$-structure on $`TM`$ is a pair $`(ϵ,\stackrel{~}{𝒮})`$ consisting of an orientation $`ϵ`$ of $`TM`$ and a $`B`$-$`A`$ imprimitivity bimodule $`\stackrel{~}{𝒮}`$.
The compact smooth Riemannian manifold $`M`$ is a $`\mathrm{Spin}^{}`$-manifold if the tangent bundle $`TM`$ of $`M`$ admits a $`\mathrm{Spin}^{}`$-structure.
Questions like existence or uniqueness of $`\mathrm{Spin}^{}`$-structures are complicated and depend on several properties of the manifold $`M`$. For an accessible and detailed geometrical account see .
By the Serre-Swan’ theorem the $`B`$-$`A`$ imprimitivity bimodule $`\stackrel{~}{𝒮}`$ can be realized as the $`C^{\mathrm{}}(M)`$-module of smooth sections $`\mathrm{\Gamma }^{\mathrm{}}(S)`$ of a uniquely determined vector bundle $`(S,p_S,M)`$ with base space $`M`$. The vector bundle $`(S,p_S,M)`$ is called the spinor bundle.
If $`n=2m`$ then the spinor bundle admits a non-trivial $`_2`$-grading arising from the grading of the Clifford bundle $`l(M)`$: $`S=S^+S^{}`$, where $`\mathrm{dim}(S_x^+)=\mathrm{dim}(S_x^{})=2^{m1}`$. Furthermore, the set of smooth sections of $`S`$ always admits a $`C^{\mathrm{}}(M)`$-valued inner product $`.,._{C^{\mathrm{}}(M)}`$. The smooth sections of the spinor bundles are called spinors, or chiral vector fields in physics.
###### Definition 3.3.
Let $`H`$ be the Hilbert space
$$H:=\overline{\{\xi \mathrm{\Gamma }^{\mathrm{}}(S):_M\xi ,\xi _{C^{\mathrm{}}(M)}𝑑g<+\mathrm{}\}},$$
Sometimes $`H`$ is referred to as the spinor Hilbert space. The Hilbert space $`H`$ consists of all square-integrable sections of the spinor bundle $`S`$, i.e. $`H=L_2(M,S)`$.
To obtain more well-behaved structure we have to assume additionally that the manifold $`M`$ is compact. The spinor Hilbert space $`H`$ inherits the non-trivial $`_2`$-grading arising from the grading of the spinor bundle in case $`n=2m`$: $`H=H^+H^{}`$, where $`H^\pm :=L_2(M,S^\pm )`$.
The sections of the Clifford bundle $`l(M)`$ act naturally on $`H`$. To look for details recall that $`\mathrm{\Gamma }(l(M))=C(M)+\mathrm{\Gamma }(T^{}M)+\mathrm{}`$. Then the elements of $`C(M)`$, i.e. of the zeroth component of $`\mathrm{\Gamma }(l(M))`$, act as multiplication operators on $`\mathrm{\Gamma }(S)`$ and, hence, on $`H`$ by continuity. Identifying $`\mathrm{\Gamma }^{\mathrm{}}(T^{}M)`$ by the $`C^{\mathrm{}}(M)`$-module $`A^1(M)`$ of $`1`$-forms on $`M`$, the images of $`1`$-forms under $`\gamma `$ fulfill the rule
$$\gamma (\alpha )\gamma (\beta )+\gamma (\beta )\gamma (\alpha )=2g^{ij}\alpha _i\beta _j\mathrm{for}\alpha ,\beta A^1(M).$$
Consequently, $`\gamma (dx^k)^2>0`$ and non-trivial $`1`$-forms are faithfully represented. The representation $`\gamma :\mathrm{\Gamma }(l(M))B(H)`$ is called the spin representation. We will use it again in the last part of the present survey.
## 4. Spin connection and Dirac operator
Let $`(M,g)`$ be a compact smooth Riemannian $`\mathrm{Spin}^{}`$-manifold, where the Riemannian metric $`g_x`$ induces a scalar product in the cotangent spaces $`T_x^{}M`$ for any $`xM`$. The Riemannian metric $`g`$ on $`M`$ gives rise to a unique Levi-Civita connection $`^g`$ (or Riemannian connection). It is defined on (contra-/covariant) $`C^{\mathrm{}}`$-tensor fields over $`M`$ of arbitrary order, $`^g`$ is symmetric, $`^g(g^{ij})=0`$ and the torsion of $`^g`$ vanishes.
In particular, $`^g:A^1(M)A^1(M)_AA^1(M)`$ obeying a Leibniz rule:
$$^g(\omega a)=^g(\omega )a+\omega da$$
for $`aC^{\mathrm{}}(M)`$ and arbitrary tensor fields $`\omega `$ on $`M`$. Lifting this Levi-Civita connection to the spinor bundle $`S`$ (where $`\mathrm{\Gamma }(S)`$ is equipped with the $`C(M)`$-valued inner product arising from $`\stackrel{~}{𝒮}`$) we obtain another Levi-Civita connection there, the spin connection.
###### Definition 4.1.
The spin connection is an operator $`^S:\mathrm{\Gamma }^{\mathrm{}}(S)\mathrm{\Gamma }^{\mathrm{}}(S)_AA^1(M)`$ that is linear and satisfies the two Leibniz rules
$$^S(\psi a)=^S(\psi )a+\psi da,$$
$$^S(\gamma (\omega )\psi )=\gamma (^g(\omega ))\psi +\gamma (\omega )^S(\psi )$$
for $`aA=C^{\mathrm{}}(M)`$, $`\omega A^1(M)=\mathrm{\Gamma }^{\mathrm{}}(T^{}M)`$, $`\psi \mathrm{\Gamma }^{\mathrm{}}(S)`$.
The spin connection on the spinor bundle $`(S,p_S,M)`$ gives rise to the Dirac operator acting on the spinor Hilbert space $`H`$.
###### Definition 4.2.
Let $`m:\mathrm{\Gamma }^{\mathrm{}}(S)_AA^1(M)\mathrm{\Gamma }^{\mathrm{}}(S)`$ be the mapping defined by the rule $`m(\psi \omega )=\gamma (\omega )(\psi )`$ for $`\omega A^1(M)=\mathrm{\Gamma }^{\mathrm{}}(T^{}M)\mathrm{\Gamma }^{\mathrm{}}(l(M))`$, $`\psi \mathrm{\Gamma }^{\mathrm{}}(S)`$. The Dirac operator on $`S`$ is the mapping $`D/:=m^S`$ that acts on the domain $`\mathrm{\Gamma }^{\mathrm{}}(S)H`$ of the spinor Hilbert space $`H`$ as an unbounded operator.
The Dirac operator $`D/`$ has a number of remarkable properties. We list them without proof. For references see :
* If $`n=2m`$ then $`D/:\mathrm{\Gamma }^{\mathrm{}}(S^\pm )\mathrm{\Gamma }^{\mathrm{}}(S^{})`$. Moreover, with respect to this decomposition of $`\mathrm{\Gamma }^{\mathrm{}}(S)`$ the Dirac operator can be represented as
$$D/=\left(\begin{array}{cc}0& D/^+\\ D/^{}& 0\end{array}\right),D/^+(h^+),h^{}=h^+,D/^{}(h^{})$$
for $`h^\pm \mathrm{\Gamma }^{\mathrm{}}(S^\pm )`$.
* If $`n=2m`$ and $`\chi :(h^+,h^{})H^+H^{}(h^+,h^{})H^+H^{}`$ is the grading operator on the spinor Hilbert space $`H`$ then $`\chi D/+D/\chi =0`$.
* $`D/`$ is symmetric and extends to an unbounded self-adjoint operator on $`H`$. (Same denotation.)
* $`[D/,a]`$ is compact and $`[[D/,a],b]=0`$ for every $`a,bC^{\mathrm{}}(M)`$.
* $`D/`$ is a Fredholm operator, i.e. $`\mathrm{ker}(D/)`$ is finite-dimensional.
* The operator $`D/^1`$ defined on the orthogonal complement of $`\mathrm{ker}(D/)`$ is compact. The eigenvalues $`\{\lambda _k\}`$ of $`D/^1`$ counted with multiplicity fulfil the relation $`\lambda _kCk^{1/n}`$ for some constant $`C`$ and $`n=\mathrm{dim}(M)`$.
* The spectrum of $`D/`$ is discrete and consists of eigenvalues of finite multiplicity.
* $`D/`$ is an elliptic first order differential operator.
* The algebra $`A=C^{\mathrm{}}(M)`$ is represented on the spinor Hilbert space $`H`$ by multiplication operators (via $`\gamma `$). We obtain
$$[D/,a]=D/(a\psi )aD/(\psi )=\gamma (da)\psi $$
for $`aA=C^{\mathrm{}}(M)`$, $`\psi H`$. In particular, since $`a`$ is smooth and $`M`$ is compact, the operator $`[D/,a]`$ is bounded with the sup-norm $`\gamma (da)_{\mathrm{}}`$ of the multiplication operator by $`\gamma (da)`$.
* For the geodesic distance of two points $`p,qM`$ we have
$$d(p,q)=sup\{|\widehat{p}(a)\widehat{q}(a)|:aC^{\mathrm{}}(M),[D/,a]1\},$$
where $`\widehat{p}`$ is the character on $`C^{\mathrm{}}(M)`$ induced by evaluation in $`pM`$ and $`\gamma (da)_{\mathrm{}}=a_{Lip}=[D/,a]`$.
* The Lichnérowicz formula is valid:
$$D/^2=\mathrm{\Delta }^S+\frac{1}{4}R,$$
where $`R`$ is the scalar curvature of the metric and $`\mathrm{\Delta }^S`$ is the Laplacian operator lifted to the spinor bundle that can be described in local coordinates by $`\mathrm{\Delta }^S=g^{ij}(_i^S_j^S\mathrm{\Gamma }_{ij}^k_k^S)`$ with $`\mathrm{\Gamma }_{ij}^k`$ the Christoffel symbols of the connection.
* For any $`fC^{\mathrm{}}(M)`$ one has the formula
$$_Mfdg=\frac{(1)^nn\mathrm{\Gamma }(n/2)}{2^{[n/2]+1n}\pi ^{n/2}}\mathrm{Tr}_\omega (f|D/|^n),$$
where $`\mathrm{Tr}_\omega `$ denotes the Dixmier trace.
## 5. The universal differential algebra $`\mathrm{\Omega }C^{\mathrm{}}(M)`$ and Connes’ differential algebra $`\mathrm{\Omega }_{D/}C^{\mathrm{}}(M)`$
As a good source for the commutative approach to differential algebras we can refer to the monograph of G. Landi . Complementary information can be found in .
###### Definition 5.1.
Let $`M`$ be a compact smooth manifold. Identify a suitable completion of the algebraic tensor product $`C^{\mathrm{}}(M)\mathrm{}C^{\mathrm{}}(M)`$ with $`C^{\mathrm{}}(M\times \mathrm{}\times M)`$, the same number of $`/\times `$ operations supposed.
The universal differential algebra $`\mathrm{\Omega }C^{\mathrm{}}(M)=_p\mathrm{\Omega }^pC^{\mathrm{}}(M)`$ is defined by the linear spaces:
$`\mathrm{\Omega }^0C^{\mathrm{}}(M)`$ $`:=`$ $`C^{\mathrm{}}(M)`$
$`\mathrm{\Omega }^pC^{\mathrm{}}(M)`$ $`:=`$ $`\{f\overline{_1^{p+1}C^{\mathrm{}}(M)}:`$
$`f(x_1,\mathrm{},x_{k1},x,x,x_{k+2},\mathrm{},x_{p+1})=0,k\}`$
The exterior differential $`\delta :\mathrm{\Omega }^p\mathrm{\Omega }^{p+1}`$ is defined by
$`(\delta f)(x_1,x_2)`$ $`:=`$ $`f(x_2)f(x_1)`$
$`(\delta f)(x_1,\mathrm{},x_{p+1})`$ $`:=`$ $`{\displaystyle \underset{k=1}{\overset{p+1}{}}}(1)^{k1}f(x_1,\mathrm{},x_{k1},x_{k+1},\mathrm{},x_{p+1})`$
The $`C^{\mathrm{}}(M)`$-bimodule structure on $`\mathrm{\Omega }C^{\mathrm{}}(M):=_p\mathrm{\Omega }^pC^{\mathrm{}}(M)`$ is given by:
$`(gf)(x_1,\mathrm{},x_{p+1})`$ $`:=`$ $`g(x_1)f(x_1,\mathrm{},x_{p+1})`$
$`(fg)(x_1,\mathrm{},x_{p+1})`$ $`:=`$ $`f(x_1,\mathrm{},x_{p+1})g(x_{p+1})`$
It extends to a general multiplication by the formula
$$(fh)(x_1,\mathrm{},x_{(p+q)+1}):=f(x_1,\mathrm{},x_{p+1})h(x_{p+1},\mathrm{},x_{(p+q)+1})$$
for $`f\mathrm{\Omega }^pC^{\mathrm{}}(M)`$, $`h\mathrm{\Omega }^qC^{\mathrm{}}(M)`$.
Key properties of the exterior differential are linearity, the Leibniz rule and the vanishing of its square:
$`\delta (ab)`$ $`=`$ $`(\delta a)b+(1)^pa(\delta b),\delta ^2=0,`$
$`\delta (\alpha a+\beta b)`$ $`=`$ $`\alpha (\delta a)+\beta (\delta b)`$
for $`a\mathrm{\Omega }^pC^{\mathrm{}}(M)`$, $`b\mathrm{\Omega }C^{\mathrm{}}(M)`$, $`\alpha ,\beta `$. These three properties give rise to another representation of the differential algebra as a linear hull of standard elements as it is used in the noncommutative case:
$$\mathrm{\Omega }^pC^{\mathrm{}}(M)=\mathrm{Lin}\{a_0\delta a_1\mathrm{}\delta a_p:a_iC^{\mathrm{}}(M)\},$$
$$\delta (a_0\delta a_1\mathrm{}\delta a_p)=\delta a_0\delta a_1\mathrm{}\delta a_p.$$
We take the parity of the degree $`p`$ as a grading for the differential algebra $`\mathrm{\Omega }C^{\mathrm{}}(M)=_p\mathrm{\Omega }^pC^{\mathrm{}}(M)`$.
To go further and to construct Connes’ differential algebra we need another property of our compact smooth manifold $`M`$ – it has to be Riemannian and $`\mathrm{Spin}^{}`$. Then we have a spectral triple $`(C^{\mathrm{}}(M),H=L_2(M,S),D/)`$ by construction, and we consider an algebraic representation of $`\mathrm{\Omega }C^{\mathrm{}}(M)`$ on $`B(H)`$:
$$\pi :\mathrm{\Omega }C^{\mathrm{}}B(H),\pi (a_0\delta a_1\mathrm{}\delta a_p):=a_0[D/,a_1]\mathrm{}[D/,a_p]$$
where $`a_iC^{\mathrm{}}(M)`$, and $`C^{\mathrm{}}(M)`$ acts on $`\stackrel{~}{𝒮}H`$ by the usual module action. (If one introduces an involution on the differential algebra then $`\pi `$ becomes a $``$-representation, however we do not need this additional structure for our purposes.) If we want $`\pi `$ to be a representation commuting with the action of the differential in some way we run into difficulty since $`\pi (\omega )=0`$ does not imply $`\pi (\delta \omega )=0`$, in general. Fortunately, there exists a differential ideal of $`\mathrm{\Omega }C^{\mathrm{}}(M)`$, the ’junk ideal’ $`J`$.
###### Lemma 5.2.
Let $`J_0:=_pJ_0^p`$ be the graded two-sided ideal of $`\mathrm{\Omega }C^{\mathrm{}}(M)`$ given by
$$J_0^p:=\{\omega \mathrm{\Omega }^pC^{\mathrm{}}(M):\pi (\omega )=0\}.$$
Then $`J:=J_0+\delta J_0`$ is a graded differential two-sided ideal of $`\mathrm{\Omega }C^{\mathrm{}}(M)`$.
###### Definition 5.3.
(A. Connes)
The graded differential algebra of Connes’ forms over the algebra $`C^{\mathrm{}}(M)`$ is defined by
$$\mathrm{\Omega }_{D/}C^{\mathrm{}}(M):=\mathrm{\Omega }C^{\mathrm{}}(M)/J\pi (\mathrm{\Omega }C^{\mathrm{}}(M))/\pi (\delta J_0).$$
The space of Connes’ $`p`$-forms is $`\mathrm{\Omega }_{D/}^pC^{\mathrm{}}(M)=\mathrm{\Omega }^pC^{\mathrm{}}(M)/J^p`$. On $`\mathrm{\Omega }_{D/}C^{\mathrm{}}(M)`$ there exists a differential induced by $`\delta `$ with the usual properties:
$$d:\mathrm{\Omega }_{D/}^pC^{\mathrm{}}(M)\mathrm{\Omega }_{D/}^{p+1}C^{\mathrm{}}(M),d([\omega ]):=[\delta \omega ]\pi ([\delta \omega ]).$$
## 6. The exterior algebra bundle $`\mathrm{\Lambda }(M)`$ and the de Rham complex
Let $`M`$ be a compact smooth manifold equipped with an atlas inherited from the cotangent bundle $`T^{}M`$. Denote by $`\mathrm{\Lambda }(T_x^{}M)`$ the real exterior algebra of the cotangent space $`T_x^{}M`$, $`xX`$. Recall that
$$\mathrm{\Lambda }(T_x^{}M):=𝒯(T_x^{}M)/\mathrm{Ideal}(ee:eT_x^{}M).$$
The real exterior algebra $`\mathrm{\Lambda }(T_x^{}M)`$ possesses a $`_2`$-grading, i.e. a linear operator $`\chi `$ on it with $`\chi ^2=\mathrm{id}`$, eigenvalues $`\{1,1\}`$ and isomorphic eigen-spaces $`\mathrm{\Lambda }^+(T_x^{}M)`$, $`\mathrm{\Lambda }^{}(T_x^{}M)`$ summing up to the algebra itself. The signs $`\pm `$ stand for the parity of the degree $`p`$ of the exterior form. An exterior $`p`$-form on $`M`$ is locally given by
$$\omega =\underset{i_1,\mathrm{},i_p}{}a_{i_1,\mathrm{},i_p}dx^{i_1}\mathrm{}dx^{i_p}$$
with smooth functions $`a_{i_1,\mathrm{},i_p}(x)`$ defined on a chart $`U`$.
###### Definition 6.1.
The exterior algebra bundle $`\mathrm{\Lambda }(M)`$ is fibrewise defined using the atlas on $`M`$ induced by the cotangent bundle atlas of $`T^{}M`$:
$$\mathrm{\Lambda }_x(M):=\mathrm{\Lambda }(T_x^{}M)_{},xX.$$
Consequently, $`\mathrm{\Lambda }^0(M)`$ is a trivial line bundle over $`M`$ and $`\mathrm{\Lambda }^1(M)=T^{}M`$ and $`\mathrm{\Lambda }^k(M)=0`$ for $`k>n=\mathrm{dim}(M)`$. The set $`\mathrm{\Lambda }^p(M)`$ is said to be the set of all $`p`$-forms, and $`\mathrm{\Lambda }(M):=_p\mathrm{\Lambda }^p(M)`$ is a linear space by definition. The multiplication is fibrewise defined by the $``$-multiplication of $`\mathrm{\Lambda }_x(M)`$, i.e. $`\omega _1\omega _2=(1)^{pq}\omega _2\omega _1`$ for $`\omega _1\mathrm{\Lambda }^p`$, $`\omega _2\mathrm{\Lambda }^q`$.
The exterior differential $`d:\mathrm{\Gamma }^{\mathrm{}}(\mathrm{\Lambda }^p(M))\mathrm{\Gamma }^{\mathrm{}}(\mathrm{\Lambda }^{p+1}(M))`$ induced by the local differential $`d_x`$ on $`\mathrm{\Lambda }(T_x^{}M)`$ is linear and obeys the rules
$`d(\omega _1\omega _2)`$ $`=`$ $`d\omega _1\omega _2+(1)^p\omega _1d\omega _2,`$
$`d(d\omega )`$ $``$ $`0`$
for $`\omega _1\mathrm{\Lambda }^p`$, $`\omega _2\mathrm{\Lambda }^q`$. Moreover, in local coordinates we have
$`df`$ $`=`$ $`{\displaystyle \underset{i}{}}{\displaystyle \frac{f}{x_i}}dx^i\mathrm{for}fC^{\mathrm{}}(M),`$
$`d\omega `$ $`=`$ $`{\displaystyle \underset{i_1,\mathrm{},i_p}{}}da_{i_1,\mathrm{},i_p}dx^{i_1}\mathrm{}dx^{i_p},\omega \mathrm{\Lambda }^p.`$
As a result we obtain a complex, the de Rham’ complex
$$0\mathrm{\Gamma }^{\mathrm{}}(\mathrm{\Lambda }^0(M))\stackrel{d}{}\mathrm{\Gamma }^{\mathrm{}}(\mathrm{\Lambda }^1(M))\stackrel{d}{}\mathrm{}\stackrel{d}{}\mathrm{\Gamma }^{\mathrm{}}(\mathrm{\Lambda }^n(M))0$$
that gives rise to cohomology groups that are isomorphic to the cohomology groups $`H^{}(M,)`$. The name of the complex comes from the application of de Rham’s theorem to this particular situation, cf..
## 7. $`\mathrm{\Omega }_{D/}C^{\mathrm{}}(M)`$ versus $`\mathrm{\Lambda }(M)`$
The final goal of the present notes is another theorem relating a structure formally depending on the Riemannian metric on the compact smooth manifold $`M`$ to another structure that does not depend even on its existence or absence. This observation by A. Connes was the starting point for his subsequent noncommutative generalizations, cf. .
###### Theorem 7.1.
(A. Connes)
Comparing the components of Connes’ differential algebra $`\mathrm{\Omega }_{D/}C^{\mathrm{}}(M)`$ and the smooth sections of components of the exterior algebra bundle $`\mathrm{\Lambda }(M)`$ we obtain an isomorphism $`\mathrm{\Omega }_{D/}^pC^{\mathrm{}}(M)\mathrm{\Gamma }^{\mathrm{}}(\mathrm{\Lambda }^p(M))`$ for every $`p0`$. Moreover, it extends to the commutative diagrams
$$\begin{array}{ccc}\mathrm{\Omega }_{D/}^pC^{\mathrm{}}(M)& \stackrel{d}{}& \mathrm{\Omega }_{D/}^{p+1}C^{\mathrm{}}(M)\\ & & \\ \mathrm{\Gamma }^{\mathrm{}}(\mathrm{\Lambda }^p(M))& \stackrel{d}{}& \mathrm{\Gamma }^{\mathrm{}}(\mathrm{\Lambda }^{p+1}(M))\end{array}$$
showing an equivalence of differential algebras.
Harald Upmeier kindly communicated a new approach to a proof of this theorem. We present here the variant that arose after some discussions and that preserves his basic ideas.
###### Proof.
For every $`p0`$ consider the subbundle $`l(M)^{pev}`$ that consists of the intersection of the subbundle of all elements of $`l(M)`$ of degree at most $`p`$ with either the subbundle $`l(M)^{even}`$ or $`l(M)^{odd}`$ in accordance with the parity of $`p`$. In the same manner we define $`\mathrm{\Omega }^{pev}C^{\mathrm{}}(M)=_{k=pev}\mathrm{\Omega }^kC^{\mathrm{}}(M)`$, where $`k`$ runs over all indices between $`0`$ and $`p`$ differing from $`p`$ by zero or an even number.
Claim 1: $`\pi (\mathrm{\Omega }^{pev}C^{\mathrm{}}(M))\gamma (\mathrm{\Gamma }^{\mathrm{}}(l^{pev}(M)))`$ for every $`p`$.
We prove the claim by induction. For $`p=0`$ a comparison of the definitions shows that $`\pi (f)=\gamma (f)=f\mathrm{id}_H`$ for every $`fC^{\mathrm{}}(M)`$. In case $`p=1`$ we obtain $`\pi (df)=[D/,\pi (f)]=[D/,\gamma (f)]=\gamma (df)`$ for every $`fC^{\mathrm{}}(M)`$.
To show the general argument recall that the complexified Clifford algebra of a real vector space $`V`$ and the complexified exterior algebra of $`V`$ are related by the isomorphisms $`Cl_{}^{pev}(V)/Cl_{}^{(p2)ev}(V)\mathrm{\Lambda }_{}^p(V)`$ for every $`p`$. So there exist induced symbol maps between the components of the Clifford bundle and the exterior algebra bundle over $`M`$,
$$\sigma ^p:\mathrm{\Gamma }^{\mathrm{}}(l^{pev}(M))\mathrm{\Gamma }^{\mathrm{}}(\mathrm{\Lambda }^p(M)),$$
$`p`$, with kernels $`\mathrm{ker}(\sigma ^p)=\mathrm{\Gamma }^{\mathrm{}}(l^{(p2)ev}(M))`$. Consequently, every smooth section of the $`p`$-th component $`l^p(M)`$ of the Clifford bundle can be represented as a finite linear combination of elementary elements of the form $`\{f_0df_1\mathrm{}df_p:f_iC^{\mathrm{}}(M)\}`$, where the central dot denotes the Clifford multiplication. To see this apply Swan’s theorem and take the projection $`P`$ of the canonical orthonormal basis $`\{e_i\}`$ of the trivial bundle that houses $`\mathrm{\Lambda }^p(M)`$ as a direct summand. By Theorem 2.3 the set $`\{P(e_i)\}`$ generates $`\mathrm{\Gamma }^{\mathrm{}}(\mathrm{\Lambda }^p(M))`$ as a $`C^{\mathrm{}}(M)`$-module. There exists a finite atlas of $`M`$ and a partition of unity $`\{u_\alpha \}`$ corresponding to it such that every component $`u_\alpha P(e_i)`$ of a certain generator $`P(e_i)`$ can be written as a finite sum of elements of the set $`\{f_{0,\alpha }df_{1,\alpha }\mathrm{}df_{p,\alpha }:f_{i,\alpha }C^{\mathrm{}}(M),\mathrm{supp}(f_{i,\alpha })\mathrm{supp}(u_\alpha )\}`$. Since all sums are finite we get the desired decomposition property for smooth sections of $`\mathrm{\Lambda }^p(M)`$. Pulling this system of generators back via $`\sigma ^p`$ we obtain it for smooth sections of $`l^p(M)`$, too.
To show the inclusion $`\pi (\mathrm{\Omega }^{pev}C^{\mathrm{}}(M))\gamma (\mathrm{\Gamma }^{\mathrm{}}(l^{pev}(M)))`$ we have only to check the canonical elements $`f_0df_1df_2\mathrm{}df_p\mathrm{\Omega }^pC^{\mathrm{}}(M)`$ since the inclusion is supposed to be already established for lower degrees by induction. We have
$`\pi (f_0df_1df_2\mathrm{}df_p)`$ $`=`$ $`\pi (f_0)[D/,\pi (f_1)]\mathrm{}[D/,\pi (f_p)]`$
$`=`$ $`\gamma (f_0)\gamma (df_1)\mathrm{}\gamma (f_p)`$
$`=`$ $`\gamma (f_0df_1\mathrm{}df_p).`$
Conversely, we have to show that $`\pi (\mathrm{\Omega }^{pev}C^{\mathrm{}}(M))\gamma (\mathrm{\Gamma }^{\mathrm{}}(l^{pev}(M)))`$ for every $`p`$.
By the results of our considerations on the symbol maps and by induction we have to verify the inclusion only for finite sums $`c=_{fin.,l}f_{0,l}df_{1,l}\mathrm{}df_{p,l}l^p(M)`$. We get
$$\gamma (c)=\underset{fin.,l}{}\gamma (f_{0,l})[D/,f_{1,l}]\mathrm{}[D/,f_{p,l}]=\pi \left(\underset{fin.,l}{}f_{0,l}df_{1,l}\mathrm{}df_{p,l}\right)\pi (\mathrm{\Omega }^pC^{\mathrm{}}(M)).$$
This establishes the statement of the first claim.
Claim 2: $`\pi (d\mathrm{ker}(\pi ^{p1}))\gamma (\mathrm{ker}(\sigma ^p))`$ for any $`p`$ with $`p2`$.
Suppose, $`\omega =_{fin.,l}f_{0,l}df_{1,l}\mathrm{}df_{p1,l}\mathrm{ker}(\pi ^{p1})\mathrm{\Omega }_{D/}^{p1}`$. By the first step
$$\gamma \left(\underset{fin.,l}{}f_{0,l}df_{1,l}\mathrm{}df_{p1,l}\right)=\pi (\omega )=0$$
and $`_{fin.,l}f_{0,l}df_{1,l}\mathrm{}df_{p1,l}=0`$ since $`\gamma `$ is injective on such elementary elements. Consider $`d\omega =_{fin.,l}df_{0,l}df_{1,l}\mathrm{}df_{p1,l}l^p(M)`$ :
$`\pi (d\omega )`$ $`=`$ $`\gamma \left({\displaystyle \underset{fin.,l}{}}df_{0,l}df_{1,l}\mathrm{}df_{p1,l}\right)`$
$`\sigma ^p\left({\displaystyle \underset{fin.,l}{}}df_{0,l}df_{1,l}\mathrm{}df_{p1,l}\right)`$ $`=`$ $`{\displaystyle \underset{fin.,l}{}}df_{0,l}df_{1,l}\mathrm{}df_{p1,l}`$
$`=`$ $`d\left({\displaystyle \underset{fin.,l}{}}f_{0,l}df_{1,l}\mathrm{}df_{p1,l}\right)`$
$`=`$ $`d\sigma ^{p1}\left({\displaystyle \underset{fin.,l}{}}f_{0,l}df_{1,l}\mathrm{}df_{p1,l}\right)`$
$`=`$ $`0.`$
Consequently, $`\pi (d\mathrm{ker}(\pi ^{p1}))\gamma (\mathrm{ker}(\sigma ^p))`$ for every $`p`$ with $`p2`$.
To show the reverse inclusion, let $`c\mathrm{\Gamma }^{\mathrm{}}(l^{(p2)}(M))`$. If $`\{u_\alpha \}`$ is a partition of unity corresponding to the selected atlas then we can assume $`\mathrm{supp}(c)U`$ since $`\gamma (c)=\gamma (_\alpha u_\alpha c)=_\alpha \gamma (u_\alpha c)`$. Furthermore, by the discussions in the first part of this proof
$$c=\underset{fin.,l}{}f_{0,l}df_{1,l}\mathrm{}df_{p2,l}$$
for some functions $`f_iC^{\mathrm{}}(M)`$.
Let $`hC^{\mathrm{}}(M)`$ with $`h(y)\lambda >0`$ for any $`yM`$ and $`d_xh,d_xh_{g^1}(x)\mu >0`$ for every $`xU`$. This forces $`h,h^1C^{\mathrm{}}(M)`$ and
$$\stackrel{~}{f}_{0,l}(x):=\frac{f_{0,l}(x)}{2d_xh,d_xh^1_{g^1}}=\frac{h(x)^2}{2d_xh,d_xh_{g^1}}f_{0,l}(x)C^{\mathrm{}}(M)$$
for every $`l`$. Furthermore,
(1)
$$\stackrel{~}{f}_{0,l}(dhdh^1+dh^1dh)=2\stackrel{~}{f}_{0,l}dh,dh^1_{g^1}=f_{0,l}$$
for every $`l`$, and by the first step we obtain
$$\pi (hdh^1+h^1dh)=\gamma (hdh^1+h^1dh)=\gamma (d(hh^1))=0.$$
Therefore, for any $`l`$
$$\omega _l:=(hdh^1+h^1dh)df_{1,l}df_{2,l}\mathrm{}df_{p2,l}\mathrm{ker}(\pi ^{p1}),$$
$$\stackrel{~}{f}_{0,l}d\omega _l(\mathrm{ker}(\pi ^p)+d\mathrm{ker}(\pi ^{p1})),$$
since the latter is an ideal in $`\mathrm{\Omega }^p`$. Finally, by (1) and the first step:
$`\gamma (c)`$ $`=`$ $`\gamma \left({\displaystyle \underset{fin.,l}{}}\stackrel{~}{f}_{0,l}(dhdh^1+dh^1dh)df_{1,l}df_{2,l}\mathrm{}df_{p2,l}\right)`$
$`=`$ $`\pi \left({\displaystyle \underset{fin.,l}{}}\stackrel{~}{f}_{0,l}(dhdh^1+dh^1dh)df_{1,l}df_{2,l}\mathrm{}df_{p2,l}\right)`$
$`=`$ $`\pi \left({\displaystyle \underset{fin.,l}{}}\stackrel{~}{f}_{0,l}d\omega _l\right)`$
$``$ $`\pi (\mathrm{ker}(\pi ^p)+d\mathrm{ker}(\pi ^{p1}))=\pi (d\mathrm{ker}(\pi ^{p1})).`$
We arrive at $`\pi (d\mathrm{ker}(\pi ^{p1}))\gamma (\mathrm{ker}(\sigma ^p))`$ for every $`p`$ with $`p2`$, and claim 2 is proved.
As the final step we list the following chain of identifications and isomorphisms:
$`\mathrm{\Omega }_{D/}^pC^{\mathrm{}}(M)`$ $`=`$ $`\pi (\mathrm{\Omega }^pC^{\mathrm{}}(M))/\pi (d\mathrm{ker}(\pi ^{p1}))`$
$``$ $`\mathrm{\Gamma }^{\mathrm{}}(l^{pev}(M)/\mathrm{ker}(\sigma ^p)`$
$``$ $`\mathrm{im}(\sigma ^p)`$
$`=`$ $`\mathrm{\Gamma }^{\mathrm{}}(\mathrm{\Lambda }^p(M)).`$
###### Corollary 7.2.
$`\mathrm{\Omega }_{D/}^pC^{\mathrm{}}(M)=0`$ for every $`p>\mathrm{dim}(M)=n`$.
Finishing we point out that there are new ideas and results in noncommutative geometry that are closely related to Theorem 7.1 and invent quantum de Rham cohomology on Poisson manifolds in a sense different from A. Connes’ work. Pioneering results have been obtained by M. Kontsevich , Huai-Dong Cao and Jian Zhou , among others. We refer to these sources for details.
Acknowledgements: I am very grateful to Harald Upmeier for the encouragements and the communication of his new proof of the last theorem of the present paper while the talk and these notes were under preparation. I would like to thank Peter M. Alberti and Rainer Matthes for their repeated valuable discussions and the exchange of literature that took place in Leipzig during the last two months. I am indebted to Christian Bär and Klaus Fredenhagen for comments on incorrect expression and a wrong assumption in the first version of the new proof of the main theorem that lead to the presented here improved article. |
warning/0002/cond-mat0002397.html | ar5iv | text | # On the Influence of Gravity on the Thermal Conductivity
## 1 Introduction
The simplest heat flow problem consists of a gas enclosed between two infinite, parallel plates kept at different temperatures. In the continuum limit, the Fourier law establishes a linear relation between the heat flux and the thermal gradient, i.e., $`𝐪=\kappa T`$, where $`\kappa `$ is the thermal conductivity coefficient. The continuum description applies when $`\lambda /L1`$ and $`\lambda /\mathrm{}1`$, where $`\lambda `$ is the mean free path, $`L`$ is the separation between the plates, and $`\mathrm{}T|T|^1`$ is the characteristic length over which the temperature changes. Nevertheless, exact results from the Boltzmann equation for Maxwell molecules and from the BGK model for general interactions show that the Fourier law still holds for large gradients ($`\lambda \mathrm{}`$), provided that the Knudsen number $`\text{Kn}=\lambda /L`$ is small.
In this paper we are interested in studying the influence on the heat flux of an external field $`g`$ (e.g., gravity) normal to the plates. We will also assume that the flow velocity vanishes. This absence of convection is possible if the Rayleigh number is smaller than its critical value ($`\text{Ra}<1700`$) . This precludes the existence of the Rayleigh-Bénard instability, that has been studied for dilute gases by other authors . In addition to $`\lambda `$, $`L`$, and $`\mathrm{}`$, the presence of gravity introduces a fourth characteristic length, namely $`hk_BT/mg`$, which represents the vertical distance over which the field produces a significant effect. In ordinary laboratory conditions, $`h`$ is several orders of magnitude larger than $`\lambda `$ and $`\mathrm{}`$, so that the constitutive equations are not affected by the action of gravity. However, discrepancies with respect to the Navier-Stokes predictions can be expected if $`h`$ is not extremely large. According to a recent perturbation solution of the Boltzmann equation through second order in $`g`$ , one can estimate that the heat flux deviates from the Fourier law as much as $`10\%`$ if $`\lambda \mathrm{}0.01h`$. The aim of this paper is to go beyond the second order in $`g`$ by using the BGK model of the Boltzmann equation. Specifically, we will obtain the hydrodynamic profiles as well as the pressure tensor and the heat flux through sixth order in gravity.
## 2 Description of the problem
Let us consider a dilute gas described by the BGK kinetic equation :
$$\frac{}{t}f+𝐯f+\frac{𝐅}{m}\frac{}{𝐯}f=\nu (ff_L).$$
(1)
Here, $`f(𝐫,𝐯,t)`$ is the one-particle distribution function, $`𝐅`$ is an external force, $`m`$ is the mass of a particle, $`\nu `$ is a collision frequency, and $`f_L`$ is the local equilibrium distribution function, that is characterized by the local density $`n`$, the local velocity $`𝐮`$, and the local temperature $`T`$, defined as
$$\{n,n𝐮,3nk_BT\}=𝑑𝐯\{1,𝐯,m(𝐯𝐮)^2\}f,$$
(2)
where $`k_B`$ is the Boltzmann constant. The collision frequency is proportional to the density and its dependence on the temperature models the influence of the interaction potential. For instance, for Maxwell molecules $`\nu n`$, while $`\nu nT^{1/2}`$ for hard spheres.
The problem we want to investigate is that of a gas enclosed in a slab between two plates at different temperatures. We assume the existence of a stationary state with spatial variation along the normal direction $`z`$ only and a constant external field $`𝐅=mg\widehat{𝐳}`$ along that direction. The constant $`g`$ can be interpreted as the gravitational acceleration. In addition, we assume that there is no convection, i.e., $`𝐮=0`$. In order to ease the notation, it is convenient to introduce dimensionless quantities. To that end, we choose an arbitrary point $`z_0`$ in the bulk as the origin and take the quantities at that point (denoted by a subscript 0) as reference values. Therefore, we define $`T^{}T/T_0`$, $`p^{}p/p_0`$, $`𝐯^{}𝐯/v_0`$, $`f^{}(k_BT_0/p_0)v_0^3f`$, $`g^{}g/v_0\nu _0`$, where $`p=nk_BT`$ is the hydrostatic pressure and $`v_0(k_BT_0/m)^{1/2}`$ is a thermal velocity. In these units, $`g^{}=\lambda _0/h_0`$, where we define the mean free path as $`\lambda _0=v_0/\nu _0`$. In the case of the spatial variable $`z`$, it is convenient to rescale it in a nonlinear way that takes into account the local dependence of the collision frequency. Consequently, we define
$$s=v_0^1_{z_0}^z𝑑z^{}\nu (z^{}).$$
(3)
Thus, the stationary BGK equation reads
$$\left(1+v_z^{}_sg^{}\frac{T^{}}{p^{}}D_v\right)f^{}=f_L^{},$$
(4)
where $`_s/s`$ and $`D_v/v_z^{}`$. Furthermore, for the sake of concreteness, we have restricted ourselves to the case of Maxwell molecules (i.e., $`\nu p/T`$). In the geometry of the problem, the relevant velocity moments are defined as
$$M_{\alpha \beta }=𝑑𝐯^{}v_{}^{}{}_{}{}^{2\alpha }v_{z}^{}{}_{}{}^{\beta }f^{}.$$
(5)
In particular, the corresponding moments at local equilibrium are
$$M_{\alpha \beta }^L=\frac{(2\alpha +\beta +1)!!}{\beta +1}p^{}T_{}^{}{}_{}{}^{\alpha 1+\beta /2}$$
(6)
for $`\beta `$ even, being zero otherwise. Conservation of momentum and energy implies that $`_sM_{02}=g^{}`$ and $`_sM_{11}=0`$.
In the absence of gravitation ($`g=0`$), Eq. (4) has an exact solution characterized by a constant pressure, $`p^{}=1`$, and a “linear” temperature profile, $`T^{}=1+ϵs`$, that applies to arbitrary values of the reduced thermal gradient $`ϵ=\lambda _0/\mathrm{}_0`$. The velocity moments are polynomials in both $`s`$ and $`ϵ`$. Their explicit expression for $`\beta +2(\alpha 1)0`$ is
$$M_{\alpha \beta }=(1)^\beta \underset{\stackrel{r=0}{(r+\beta )\text{even}}}{\overset{\beta +2(\alpha 1)}{}}\frac{(2\alpha +\beta +r+1)!!(\alpha 1+\frac{\beta +r}{2})!}{(\alpha 1+\frac{\beta r}{2})!(\beta +r+1)}ϵ^r(1+ϵs)^{\alpha 1+(\beta r)/2}.$$
(7)
In particular, $`M_{11}=5ϵ`$, which means that the Fourier law holds even for large thermal gradients.
The motivation of this paper is to analyze the influence of gravitation on the profiles and transport properties of the above steady Fourier flow. However, the presence of the operator $`D_v`$ in Eq. (4) complicates the problem significantly. A convenient strategy is to take the pure steady Fourier flow corresponding to a value of $`ϵ`$ equal to the actual thermal gradient at the point $`s=0`$ as a reference state. Consequently, we will carry out a perturbation expansion in powers of $`g`$:
$$f^{}=f^{(0)}+f^{(1)}g^{}+f^{(2)}g_{}^{}{}_{}{}^{2}+\mathrm{},$$
(8)
$$M_{\alpha \beta }=M_{\alpha \beta }^{(0)}+M_{\alpha \beta }^{(1)}g^{}+M_{\alpha \beta }^{(2)}g_{}^{}{}_{}{}^{2}+\mathrm{},$$
(9)
$$p^{}=p^{(0)}+p^{(1)}g^{}+p^{(2)}g_{}^{}{}_{}{}^{2}+\mathrm{},$$
(10)
$$T^{}=T^{(0)}+T^{(1)}g^{}+T^{(2)}g_{}^{}{}_{}{}^{2}+\mathrm{},$$
(11)
where $`M_{\alpha \beta }^{(0)}`$ is given by Eq. (7), $`p^{(0)}=1`$, and $`T^{(0)}=1+ϵs`$. By definition, $`p^{(k)}(0)=T^{(k)}(0)=_sT^{(k)}|_{s=0}=0`$ for $`k1`$. It must be emphasized that the terms of order $`g_{}^{}{}_{}{}^{k}`$ are nonlinear functions of $`ϵ`$ since no restriction to the order on $`ϵ`$ exists.
## 3 Perturbation expansion
In this section we obtain the hydrodynamic profiles $`p^{(k)}`$ and $`T^{(k)}`$, the momentum flux $`M_{02}^{(k)}`$, and the heat flux $`M_{11}^{(k)}`$ through order $`k=6`$. Insertion of Eq. (8) into Eq. (4) yields
$$f^{(k)}=\underset{j=0}{\overset{\mathrm{}}{}}(v_z^{}_s)^j\left[f_L^{(k)}+D_v\underset{i=0}{\overset{k1}{}}\left(\frac{T^{}}{p^{}}\right)^{(i)}f^{(ki1)}\right].$$
(12)
This is a formal solution, since $`f_L^{(k)}`$ is a functional of $`f^{(k)}`$ through its dependence on the pressure and temperature. Taking moments in Eq. (12), one has
$`\mathrm{\Delta }M_{\alpha \beta }^{(k)}`$ $``$ $`M_{\alpha \beta }^{(k)}M_{\alpha \beta }^{L(k)}={\displaystyle \underset{j=1}{\overset{\mathrm{}}{}}}(_s)^jM_{\alpha ,\beta +j}^{L(k)}{\displaystyle \underset{j=0}{\overset{\mathrm{}}{}}}(_s)^j{\displaystyle \underset{i=0}{\overset{k1}{}}}\left({\displaystyle \frac{T^{}}{p^{}}}\right)^{(i)}`$ (13)
$`\times \left(2\alpha M_{\alpha 1,\beta +j+1}^{(ki1)}+(\beta +j)M_{\alpha ,\beta +j1}^{(ki1)}\right).`$
The fact that $`f`$ and $`f_L`$ have the same hydrodynamic moments leads to the consistency conditions
$$\mathrm{\Delta }M_{00}^{(k)}=\mathrm{\Delta }M_{01}^{(k)}=\mathrm{\Delta }M_{10}^{(k)}=0,$$
(14)
for any $`k`$. In order to convert Eq. (13) into an explicit equation that can be
solved recursively, we need to know the spatial dependence of $`p^{(k)}`$ and $`T^{(k)}`$. It turns out that $`p^{(k)}`$ is a polynomial of degree $`k2`$ in $`s`$ (except $`p^{(1)}`$, that is linear in $`s`$), while $`T^{(k)}`$ is a polynomial of degree $`k+1`$, the coefficients being nonlinear functions of the reduced thermal gradient $`ϵ`$. Thus, at a given order $`k`$, insertion of these polynomials into Eq. (13) and application of the consistency requirements (14) allow one to determine the unknown coefficients and the problem can be recursively solved. Notice that, seen as functions of the actual space variable $`z`$, $`p^{(k)}`$ and $`T^{(k)}`$ are much more complicated than just polynomials. In fact, in the absence of gravitational force, the relationship between $`s`$ and $`z`$ is nonlinear: $`z=\lambda _0(s+\frac{1}{2}ϵs^2)`$. In general, such a relationship can be obtained inverting Eq. (3) as
$$z=z_0+\lambda _0_0^s𝑑s^{}\frac{T^{}(s^{})}{p^{}(s^{})}.$$
(15)
The above scheme is straightforward but tedious to carry out. Since all the manipulations are algebraic, they render themselves to the use of symbolic programming languages. In this paper, we have evaluated the perturbation expansion through sixth order in the field. Since the expressions of the hydrodynamic profiles become progressively longer, here we only give the explicit results through order $`k=4`$:
$$p^{}=1sg^{}\frac{276}{5}ϵ^2sg_{}^{}{}_{}{}^{3}\frac{1}{5}s\left[\frac{12}{5}ϵ\left(\mathrm{112\hspace{0.17em}973}ϵ^2+30\right)+588ϵ^2s\right]g_{}^{}{}_{}{}^{4}+𝒪(g_{}^{}{}_{}{}^{5}),$$
(16)
$`T^{}`$ $`=`$ $`1+ϵs+{\displaystyle \frac{1}{2}}ϵs^2g^{}ϵs^2({\displaystyle \frac{66}{5}}ϵ{\displaystyle \frac{1}{3}}s)g_{}^{}{}_{}{}^{2}ϵs^2[{\displaystyle \frac{16}{25}}(6624ϵ^2+5)`$ (17)
$`+{\displaystyle \frac{346}{15}}ϵs{\displaystyle \frac{1}{4}}s^2]g^{}^3{\displaystyle \frac{1}{5}}ϵs^2[{\displaystyle \frac{12}{25}}ϵ(\mathrm{50\hspace{0.17em}765\hspace{0.17em}962}ϵ^2+\mathrm{31\hspace{0.17em}445})`$
$`+{\displaystyle \frac{2}{15}}(\mathrm{399\hspace{0.17em}621}ϵ^2+200)s+{\displaystyle \frac{971}{6}}ϵs^2s^3]g^{}^4+𝒪(g_{}^{}{}_{}{}^{5}).`$
Once the hydrodynamic profiles are known, Eq. (13) can be used to obtain all the velocity moments at a given order. The most relevant moment is $`M_{11}`$, which is related to the heat flux, $`q_z=(p_0v_0/2)M_{11}`$. Another important quantity is $`M_{02}=P_{zz}/p_0`$, which measures the anisotropy of the pressure tensor $`𝖯`$. As said above, in the absence of gravitation the Fourier law applies exactly, i.e. $`q_z^{(0)}=\kappa T/z=(5p_0v_0/2)ϵ`$, and the pressure tensor is isotropic, i.e. $`P_{zz}^{(0)}=p_0`$. In order to characterize the deviations from these Navier-Stokes predictions due to gravity, we introduce the following reduced quantities:
$$\mathrm{\Lambda }(ϵ,g^{})=\frac{q_z}{q_z^{(0)}}=1+\underset{k=1}{\overset{\mathrm{}}{}}\mathrm{\Lambda }^{(k)}(ϵ)g_{}^{}{}_{}{}^{k},$$
(18)
$$\gamma (ϵ,g^{})=\frac{P_{zz}|_{z=z_0}}{p_0}=1+\underset{k=1}{\overset{\mathrm{}}{}}\gamma ^{(k)}(ϵ)g_{}^{}{}_{}{}^{k}.$$
(19)
The results show that $`\mathrm{\Lambda }^{(k)}`$ and $`\gamma ^{(k)}`$ are polynomials in $`ϵ`$ of degree $`k`$ and a defined parity,
$$\mathrm{\Lambda }^{(k)}(ϵ)=\underset{\mathrm{}=0}{\overset{k}{}}\mathrm{\Lambda }_{\mathrm{}}^{(k)}ϵ^{\mathrm{}},\gamma ^{(k)}(ϵ)=\underset{\mathrm{}=1}{\overset{k}{}}\gamma _{\mathrm{}}^{(k)}ϵ^{\mathrm{}}.$$
(20)
To second order in $`g^{}`$ we get $`\mathrm{\Lambda }_1^{(1)}=\frac{58}{5}`$, $`\mathrm{\Lambda }_0^{(2)}=\frac{16}{5}`$, $`\mathrm{\Lambda }_2^{(2)}=\frac{\mathrm{47\hspace{0.17em}968}}{25}`$, $`\gamma _1^{(1)}=0`$, and $`\gamma _2^{(2)}=\frac{84}{5}`$. The coefficients $`\mathrm{\Lambda }_{\mathrm{}}^{(k)}`$ and $`\gamma _{\mathrm{}}^{(k)}`$ for $`3k6`$ are given in Tables 1 and 2, respectively.
## 4 Discussion
The numerical coefficients appearing in Eqs. (16) and (17), as well as in Tables 1 and 2 clearly indicate that the expansion in powers of $`g^{}`$, Eq. (8), is only asymptotic. This does not pose a serious problem, at least from a practical point of view, except for values of $`g^{}`$ (say, $`g^{}>10^2`$) that correspond to gravitational fields unrealistically large. As a consequence, only the first few terms in the expansion are useful for small values of $`g^{}`$. In the subsequent analysis, terms of order $`g_{}^{}{}_{}{}^{3}`$ and higher will not be considered. This also allows us to make a closer comparison with results derived from the Boltzmann equation for Maxwell molecules . Notwithstanding, the knowledge of the remaining terms might be useful to attempt to resum the infinite series and get the transport properties for arbitrary values of $`g^{}`$.
The solution to the Boltzmann equation through order $`g_{}^{}{}_{}{}^{2}`$ exhibits the same structure as in Eqs. (16)–(19), so that only the numerical coefficients differ. As a matter of fact, the coefficient $`\frac{66}{5}`$ appearing in Eq. (17) is replaced by $`\frac{468}{45}`$; in addition, the Boltzmann solution yields $`\mathrm{\Lambda }_1^{(1)}=\frac{46}{5}`$, $`\mathrm{\Lambda }_0^{(2)}=\frac{12}{5}`$, $`\mathrm{\Lambda }_2^{(2)}=503.7`$, $`\gamma _1^{(1)}=0`$ and $`\gamma _2^{(2)}=\frac{128}{45}`$. The differences indicate that the influence of gravity is stronger in the BGK description than in the Boltzmann one, especially in the case of the pressure anisotropy. In order to carry out a more detailed comparison, it is convenient to construct Padé approximants for the generalized thermal conductivity coefficient $`\mathrm{\Lambda }`$. More specifically, we consider the approximants
$$\mathrm{\Lambda }_{[1,1]}(ϵ,g^{})=\frac{\mathrm{\Lambda }^{(1)}(ϵ)+\left[\mathrm{\Lambda }_{}^{(1)}{}_{}{}^{2}(ϵ)\mathrm{\Lambda }^{(2)}(ϵ)\right]g^{}}{\mathrm{\Lambda }^{(1)}(ϵ)\mathrm{\Lambda }^{(2)}(ϵ)g^{}},$$
(21)
$$\mathrm{\Lambda }_{[0,2]}(ϵ,g^{})=\left\{1\mathrm{\Lambda }^{(1)}(ϵ)g^{}+\left[\mathrm{\Lambda }_{}^{(1)}{}_{}{}^{2}(ϵ)\mathrm{\Lambda }^{(2)}(ϵ)\right]g_{}^{}{}_{}{}^{2}\right\}^1.$$
(22)
In the case of the BGK equation, these two approximants differ less than $`2\%`$ if $`g^{}<10^2`$ and $`|ϵ|g^{}<3\times 10^3`$. The differences are smaller in the case of the Boltzmann equation. In this range of values for $`ϵ`$ and $`g^{}`$, a reliable approximation is $`\mathrm{\Lambda }\frac{1}{2}\left(\mathrm{\Lambda }_{[1,1]}+\mathrm{\Lambda }_{[0,2]}\right)`$. Figure 1 shows this approximation for $`\mathrm{\Lambda }`$ in the interval $`3ϵ3`$ at $`g^{}=10^3`$, as given by the BGK and Boltzmann equations. We observe that the heat flux increases with respect to its Navier-Stokes value when one heats from above ($`ϵ>0`$), while the opposite happens when one heats from below ($`ϵ<0`$). This effect is not symmetric, since it is more significant if $`ϵ>0`$ than if $`ϵ<0`$. As said above, Fig. 1 also shows that the influence of gravity is less important in the Boltzmann description.
In summary, we have solved the BGK model for a gas simultaneously subjected to a thermal gradient and a parallel gravity field of magnitude $`g`$, in the absence of convection. The solution has been obtained by a perturbation expansion in powers of gravity, the reference state being the non-equilibrium pure Fourier flow with arbitrarily large thermal gradients. We have explicitly obtained the solution through sixth order in the field. The results clearly indicate that the expansion is not convergent, although it seems to be at least asymptotic. The work reported in this paper extends previous results derived from the Boltzmann equation to second order in the field . The similarity in the structure of the coefficients appearing in both descriptions suggests that the expansion obtained from the Boltzmann equation is also asymptotic. Nevertheless, given that at practical level, the values of $`g`$ are small, the usefulness of the expansion is restricted to the first few terms.
The main results concerning the transport of momentum and energy are that the external field induces (i) anisotropy in the pressure tensor, $`(P_{zz}p)/p\frac{84}{5}ϵ^2g_{}^{}{}_{}{}^{2}`$, and (ii) deviations from the Fourier law, $`q_z/q_z^{(0)}1\frac{58}{5}ϵg^{}`$. While the first effect is of second order, the correction to the heat flux is of first order, so that it depends on the sign of the thermal gradient. As a consequence, the heat transport is inhibited when the gas is heated from below ($`ϵ<0`$), while the opposite happens when the gas is heated from above ($`ϵ>0`$).
This work has been done under the auspices of the Agencia Española de Cooperación Internacional (Programa de Cooperación Interuniversitaria Hispano-Marroquí). V.G. and A.S. acknowledge partial support from the DGES (Spain) through Grant No. PB97-1501 and from the Junta de Extremadura (Fondo Social Europeo) through Grant No. PRI97C1041. The authors are grateful to Prof. Y. Sone for valuable discussions about the subject of this paper. |
warning/0002/nucl-ex0002003.html | ar5iv | text | # Phenomenology of the deuteron electromagnetic form factors
## Abstract
A rigorous extraction of the deuteron charge form factors from tensor polarization data in elastic electron-deuteron scattering, at given values of the 4-momentum transfer, is presented. Then the world data for elastic electron-deuteron scattering is used to parameterize, in three different ways, the three electromagnetic form factors of the deuteron in the 4-momentum transfer range 0-7 fm<sup>-1</sup>. This procedure is made possible with the advent of recent polarization measurements. The parameterizations allow a phenomenological characterization of the deuteron electromagnetic structure. They can be used to remove ambiguities in the form factors extraction from future polarization data.
1 Introduction
The deuteron, as the only two-nucleon bound state, has been the subject of many theoretical and experimental investigations. Since it has spin 1, its electromagnetic structure is described by three form factors, charge monopole $`G_C`$, charge quadrupole $`G_Q`$ and magnetic dipole $`G_M`$, assuming P- and T-invariance. Measurements of elastic electron deuteron scattering observables provide quadratic combinations of these form factors. Since most of the data available come from differential cross section measurements, it has been customary, both in the data presentation and in the comparison with theoretical models, to use the two structure functions $`A`$ and $`B`$ defined hereafter, extracted from the cross section data by a Rosenbluth separation . With the advent of tensor polarimeters and tensor polarized internal targets, polarization observables have been measured as well, which allow the separation of the two charge form factors.
The purpose of this work is twofold. First, in Sect. 2, the calculation of $`G_C`$ and $`G_Q`$, at given values of the 4-momentum transfer $`Q`$, from polarization data together with (interpolated) $`A`$ and $`B`$ data is reexamined and updated with respect to previous work.
Then, in Sect. 3, parameterizations of the three deuteron form factors, in the 4-momentum transfer range $`Q=07`$ fm<sup>-1</sup>, are provided. Above 7 fm<sup>-1</sup>, only small angle cross section data are available, preventing the separate determination of the three form factors. We have determined the three deuteron electromagnetic form factors by fitting directly the measured differential cross section and polarization observables. This procedure eliminates the need for an intermediate determination of $`A`$ and $`B`$, and results in a more realistic evaluation of errors for the form factors.
One parameterization is used for a determination of the node of the charge form factor $`G_C`$, while the application of the work of Ref. allows the determination of reduced form factors in a helicity basis. The accuracy in the determination of these form factors is limited by the assumption of a one-photon exchange mechanism in the first order Born approximation at low $`Q`$, and by the accuracy of the data at intermediate to high momentum transfers. A third parameterization was recently applied for a precise determination of the rms–charge radius of the deuteron . At low $`Q`$, Coulomb distortion was taken into account to extract precise values of $`G_C`$. Applying this correction resolved an old discrepancy between the deuteron radius determined via $`(e,e^{})`$ and N–N scattering. In the intermediate to high $`Q`$-range, other corrections such as the double scattering contribution to two photon exchange should be considered, but they are at present neither accurately calculated nor experimentally determined.
2 Observables and form factors
2.1 e-d observables
Assuming single photon exchange, the electron-deuteron unpolarized elastic differential cross section can be written as
$$\frac{d\sigma }{d\mathrm{\Omega }}=\sigma _{NS}\left[G_C^2(Q^2)+\frac{8}{9}\eta ^2G_Q^2(Q^2)+\frac{2}{3}\eta \epsilon ^1(Q^2,\theta _e)G_M^2(Q^2)\right]\sigma _{NS}S,$$
(1)
where $`\sigma _{NS}`$ is the Mott differential cross section multiplied by the deuteron recoil factor, $`\theta _e`$ the electron scattering angle, $`\eta =Q^2/4M_d^2`$, $`M_d`$ the deuteron mass; $`\epsilon =[1+2(1+\eta )\mathrm{tan}^2(\theta _e/2)]^1`$ is related to the virtual photon polarization. The quantity $`SA+B\mathrm{tan}^2(\theta _e/2)`$ defines the usual $`A`$ and $`B`$ elastic structure functions.
The tensor polarization observables $`t_{2q}`$, or equivalently the analyzing powers $`T_{2q}`$, have been measured as well. Their expression as a function of the three form factors, still in the one-photon exchange approximation, is given by:
$`\sqrt{2}St_{20}`$ $`=`$ $`{\displaystyle \frac{8}{3}}\eta G_CG_Q+{\displaystyle \frac{8}{9}}\eta ^2G_Q^2+{\displaystyle \frac{1}{3}}\eta \epsilon ^1G_M^2`$ (2)
$`\sqrt{3}St_{21}`$ $`=`$ $`2\eta \left(\eta +\eta ^2\mathrm{sin}^2{\displaystyle \frac{\theta _e}{2}}\right)^{1/2}G_MG_Q\mathrm{sec}{\displaystyle \frac{\theta _e}{2}}`$ (3)
$`2\sqrt{3}St_{22}`$ $`=`$ $`\eta G_M^2.`$ (4)
2.2 Calculation of $`G_C`$ and $`G_Q`$
The charge form factors are here extracted from $`t_{20}(Q,\theta _e)`$ data, together with $`A(Q)`$ and $`B(Q)`$ (interpolated) data. The analyses presented in need to be updated, because of new $`t_{20}`$ and $`A`$ data. In particular, the parameterization of $`A`$ used in gave a very small weight to the then only existing high $`Q`$ data and is lower than the new data around 4.5 fm<sup>-1</sup>. Furthermore, we present here a more compact solution and a more rigorous treatment of errors.
For our purpose, it is useful to define new quantities $`A_0AB/2(1+\eta )`$ and $`\stackrel{~}{t}_{20}`$ , derived respectively from $`A`$ and $`t_{20}`$ by eliminating the magnetic contribution:
$$\stackrel{~}{t}_{20}\frac{\frac{8}{3}\eta G_CG_Q+\frac{8}{9}\eta ^2G_Q^2}{\sqrt{2}(G_C^2+\frac{8}{9}\eta ^2G_Q^2)}=\frac{St_{20}+B/4\sqrt{2}\epsilon (1+\eta )}{A_0}$$
(5)
Using the reduced form factors $`g_C=G_C/\sqrt{A_0}`$ and $`g_Q=2\eta G_Q/3\sqrt{A_0}`$, (1,2,5) lead to:
$`g_C^2+2g_Q^2`$ $`=`$ $`1`$ (6)
$`2g_Cg_Q+g_Q^2`$ $`=`$ $`p\stackrel{~}{t}_{20}/\sqrt{2}`$ (7)
where $`p`$ (or conventionnally $`p_{ZZ}`$) is the tensor polarization in Cartesian notation (also called alignment). There are four solutions to these equations given by
$$(g_Q^\pm )^2=\frac{2+p\pm \sqrt{\mathrm{\Delta }}}{9}$$
(8)
with $`\mathrm{\Delta }=8(1p)(\frac{1}{2}+p)`$ and $`g_C^\pm `$ from (7). The physical solution is easily selected at small $`Q`$ from the static moments ($`g_C(0)=1,g_Q(0)=0`$). It corresponds to the choice of a minus sign in (8) and of $`g_Q>0`$. Since $`\stackrel{~}{t}_{20}`$ and $`t_{21}`$, both proportional to $`G_Q`$, do not cross zero at a same value of $`Q`$ , $`g_Q`$ has to remain positive over the whole range considered in this work. The two remaining solutions ($`g_Q^+,g_C^+`$) and ($`g_Q^{},g_C^{}`$) cross each other at values $`Q_{min}`$ and $`Q_{max}`$ where $`\stackrel{~}{t}_{20}`$ reaches its extrema $`\sqrt{2}`$ and $`+\sqrt{2}/2`$ ($`\mathrm{\Delta }=0`$). The physical solution must switch from “$``$” to “$`+`$” at $`Q=Q_{min}`$ and then back to “$``$” at $`Q_{max}`$ in order to ensure a continuity of the form factor derivatives. For polarization data close to these extrema, $`Q`$ may be below or above the a priori unknown $`Q_{min}`$ (or $`Q_{max}`$), and the choice of solution is ambiguous. $`Q_{min}`$, from our three global fits to the $`ed`$ data (see Sect. 3), is determined to be close to 3.3 fm<sup>-1</sup>. On the other hand, there are not enough polarization data to constrain the value of $`Q_{max}`$, so that the above mentioned ambiguity remains around $`Q68`$ fm<sup>-1</sup>. This is the case for the two points at highest $`Q`$ in .
An additional complication arises for five polarization data points in Refs. which lay partially outside the physical region $`\sqrt{2}\stackrel{~}{t}_{20}1/\sqrt{2}`$. This situation is quite probable for points with finite errors close to a physical limit . For the sake of extracting $`G_C`$ and $`G_Q`$, the interval of 68.3% confidence level $`[\stackrel{~}{t}_{20}\mathrm{\Delta }\stackrel{~}{t}_{20},\stackrel{~}{t}_{20}+\mathrm{\Delta }\stackrel{~}{t}_{20}]`$, and eventually the most probable value $`\stackrel{~}{t}_{20}`$, are then modified according to the method presented in . The resulting confidence interval is entirely within the physical region ($`\mathrm{\Delta }0`$). In this particular case, the modified values of $`p`$ are used in (7,8) instead of the measured ones. As a result of this procedure, the errors on the form factors may be asymmetric.
The calculated values of $`G_C`$ and $`G_Q`$, corresponding to all measurements of $`t_{20}`$, are presented in Table 1 and Fig. 1. The later also shows results of parameterizations to be discussed in Sect. 3. Uncertainties come from the quoted errors in $`t_{20}`$, combined quadratically with errors on $`A`$ and $`B`$ reflecting the spread of the data (for example, at 5 fm<sup>-1</sup>, 8.5 and 17 % respectively). For the two points of highest $`Q`$, the two solutions of (7,8) are given. The first one is preferred, based on theoretical guidance and on the parameterizations discussed below. Only parameterization I (Sect. 3.1) favors the second solution for the point at $`Q=6.64`$ fm<sup>-1</sup>. Note that $`\stackrel{~}{t}_{20}`$ need not necessarily reach its maximum allowed value, in which case the first (“$`+`$”) solution would prevail from $`Q=Q_{min}`$ up to the undetermined node of $`G_Q`$, or to the second minimum of $`\stackrel{~}{t}_{20}`$, whichever occurs first.
3 Parameterization of the form factors
The three paramaterizations described below are determined through a $`\chi ^2`$ minimization involving 269 cross section data points and 39 polarization data points . In most polarization data, and in some cross section data, the systematic uncertainties are dominant and may vary from point to point in a given experiment. The error considered in the $`\chi ^2`$ minimization is then the quadratic sum of the statistical and systematic uncertainties. The uncertainties on the parameters are given by the error matrix. For data where an overall normalization uncertainty may apply, the resulting systematic uncertainty of the fitted parameters have been evaluated by changing each individual data set by the quoted error and re-fitting the complete data set. This last procedure was carried on only with parameterization III (Sect. 3.3).
The $`\chi ^2`$ per degree of freedom ($`\chi ^2/N_{d.f.}`$) all exceed the value of 1, because of systematic differences between some data sets, at the limit or beyond the quoted systematic uncertainties. Among the most recent experiments, this is the case for the $`A`$ measurements of Refs. , and in a lesser extent for the $`t_{20}`$ measurements of Refs. . The fits then give an average representation of the data, though biased toward experiments with a larger number of data points.
3.1 Parameterization I
In the first parameterization (I), each form factor is given by:
$$G_X(Q^2)=G_X(0)\left[1\left(\frac{Q}{Q_X^0}\right)^2\right]\left[1+\underset{i=1}{\overset{5}{}}a_{X}^{}{}_{i}{}^{}Q^{2i}\right]^1,$$
(9)
with $`X=C,Q`$ or $`M`$. This expression has the advantage of displaying explicitly the first node $`Q_X^0`$ of each form factor. The normalizing factors $`G_X(0)`$ are fixed by the deuteron static moments. With 18 free parameters, a fit is obtained with $`\chi ^2/N_{d.f.}=1.5`$.
3.2 Parameterization II
Another parameterization (II) has been proposed by Kobushkin and Syamtomov . Each form factor is proportional to the square of a dipole nucleon form factor $`G_D`$ and to a linear combination of reduced helicity transition amplitudes $`g_0,g_1,g_2`$:
$$\left(\begin{array}{c}G_C\\ G_Q\\ G_M\end{array}\right)=G_D^2\left(\frac{Q^2}{4}\right)(\eta )\left(\begin{array}{c}g_0\\ g_1\\ g_2\end{array}\right).$$
(10)
Each of these amplitudes is parameterized as a sum of four Lorentzian factors:
$$g_k=Q^k\underset{i=1}{\overset{4}{}}\frac{a_{ki}}{\alpha _{ki}^2+Q^2}.$$
(11)
For each $`k`$, the $`\alpha _{ki}^2`$ follow an arithmetical suite defined by 2 independent parameters. In addition, an asymptotic behavior dictated by quark counting rules and helicity rules valid in perturbative quantum chromodynamics (pQCD), together with the normalization conditions at $`Q=0`$, imply 6 relations between the parameters $`a_{ki}`$ and $`\alpha _{ki}`$ . As a result, each amplitude is described by 4 independent parameters. New parameters are obtained here, due on one hand to a newer data base, and on the other hand to the fitting of the differential cross sections instead of $`A`$ and $`B`$. With 12 free parameters, a fit to the data set is obtained with $`\chi ^2/N_{d.f.}=1.8`$, whereas the original values of the parameters in Ref. yield $`\chi ^2/N_{d.f.}=7.5`$. This parameterization, in contrast with the two other ones presented in this paper, can be extrapolated well above 7 fm<sup>-1</sup>, albeit with some theoretical prejudice. We confirm the observation of Refs. that the double helicity flip transition amplitude $`g_2`$ has a magnitude comparable to the zero helicity flip amplitude $`g_0`$ in the $`Q`$-range considered here, which means that these amplitudes are not in the asymptotic regime expected from pQCD.
3.3 Parameterization III
The third parameterization (III) employs a Sum-of-Gaussians (SOG) . The form factors are written as
$$G_X(Q)=G_X(0)e^{\frac{1}{4}Q^2\gamma ^2}\underset{i=1}{\overset{25}{}}\frac{A_i}{1+2R_i^2/\gamma ^2}\left(\mathrm{cos}(QR_i)+\frac{2R_i^2}{\gamma ^2}\frac{\mathrm{sin}(QR_i)}{QR_i}\right)$$
(12)
Although our interest here lies in its $`Q`$-space version, the parameterization is better described in configuration space where it corresponds to a density $`\rho (R)`$ written as a sum of Gaussians placed at arbitrary radii $`R_i`$, with amplitudes $`A_i`$ fitted to the data, and a fixed width $`\gamma `$. The distance $`R`$ refers to the distance of the nucleons to the deuteron center of mass. The parameterization represents a totally general basis and the following applied restrictions are justified on physics grounds. First, one does not expect structures smaller than the size of the nucleon, which determines the width $`\gamma `$ to be the size of the proton ($`\gamma \sqrt{3/2}=0.8`$ fm). Second, the spacing between Gaussians is chosen slightly smaller than this width: 0.4 fm or 0.5 fm. Third, the Gaussians are placed at radii $`R_iR_{max}=10`$ fm, which is justified given the fact that one can easily specify the radius at which the tails of densities give no significant $`(<10^3)`$ contribution to $`G_X(Q)`$. In addition, outside the range of the NN–force, the deuteron wave functions have an analytic form which is well known and depends only on the deuteron binding energy. Thus, for radii $`R_i4`$ fm, one can impose this shape and fix the ratio of the amplitudes $`A_i`$. Each form factor is then determined with 11 free parameters: 10 Gaussian amplitudes $`A_1`$ to $`A_{10}`$, corresponding to $`R_i<4`$ fm, and one overall amplitude for the shape-given tail at $`R4`$ fm. With a total of 33 independent parameters, a $`\chi ^2/N_{d.f.}`$ of 1.5 is obtained in the fit.
3.4 Results and discussion
The resulting form factors from the three parameterizations are shown in Fig. 1. As functions of two variables ($`Q`$ and $`\theta _e`$), the fitted quantities cannot be easily represented together with the parameterizations. In order to illustrate the quality of the fits, we present plots of relative differences of $`A`$ and $`B`$, and of $`\stackrel{~}{t}_{20}(Q)`$ in Fig. 2. $`t_{21}`$ and $`t_{22}`$ are equally well fitted, which constitutes, within experimental uncertainties, an indication of the coherence of equations (1,2,3,4), and therefore of the consistency of the one-photon exchange approximation.
From the average and dispersion between the three parameterizations, combined with the fit uncertainty on $`Q_C^0`$, the node of the charge form factor is determined to be located at $`4.21\pm 0.08`$ fm<sup>-1</sup>, a value governed by the $`t_{20}`$ results of Refs. . Assuming as we do here implicitly that these two data sets have the same weight, the location of this node is not quite consistent with a relation between the two- and three-nucleon isoscalar charge form factors, established with various $`NN`$ potentials . The secondary maximum of $`|G_C|`$ is very flat, so that its location ($`5.3\pm .5`$ fm<sup>-1</sup>) is not determined very precisely. Its magnitude ($`.0038\pm .0003`$) is clearly inconsistent with the corresponding one of the three-nucleon isoscalar charge form factor, still within the same model calculations . The $`t_{21}`$ results of Ref. , though of limited accuracy, help confirm a node of the magnetic form factor at $`7.2\pm 0.3`$ fm<sup>-1</sup>. As for the first node of $`G_Q`$, according to most theoretical models, it should appear at a higher value of $`Q`$, above the range where our parameterization method applies. The value $`Q_Q^0=7.7\pm 0.6`$ fm<sup>-1</sup> given by parameterization I is probably the smallest possible value allowed by the present data. It is due to this parameterization following the downward trend of the $`t_{20}`$ data point at the highest $`Q`$ (see Fig. 2). This trend however is not statistically significant. Parameterization II, when extrapolated, suggests a much higher value of $`Q`$ for the node of $`G_Q`$. Finally, from
$$r^26\frac{dG_C}{dQ^2}|_{Q^2=0}=6\left[a_{C}^{}{}_{1}{}^{}+(Q_C^0)^2\right],$$
(13)
we calculate the root mean square charge radius of the deuteron $`r=2.094\pm 0.003`$ (stat.) $`\pm 0.009`$ (syst.) fm. The statistical uncertainty is given by the error matrix from parameterization I, while the systematic uncertainty is evaluated with parameterization III (see above remark about normalization uncertainties on individual data sets). This radius is 1.7% smaller than the value $`r=2.130`$ fm reported in , consistent with expectations in the absence of corrections due to Coulomb distortion.
4 Conclusion
The extraction of the charge form factors $`G_C`$ and $`G_Q`$ from experiment, at given values of $`Q`$, has been reexamined. The solutions were expressed in the most compact and physical way, while a new treatment of errors was applied to polarization data at or beyond physical limits. The existing electron-deuteron elastic scattering data were used for direct parameterizations of the three deuteron electromagnetic form factors, up to $`Q=7`$ fm<sup>-1</sup>. The numerical results may be requested from the authorsContacts: jball@cea.fr (parameterizations I and II), jourdan@ubaclu.unibas.ch (III). and will be updated as new data become available in the future. The inferred value of $`Q_{min}3.3`$ fm<sup>-1</sup> corresponding to the minimum of $`\stackrel{~}{t}_{20}`$ could be used, or recalculated with such global fits, for future experiments in this $`Q`$-range , in order to resolve the discussed ambiguities in the form factors calculation. These future experiments should help confirm, or adjust, the exact value of the node of the charge form factor: this location is sensitive to the strength of the $`NN`$ repulsive core, to the size of the isoscalar meson exchange contributions and to relativistic corrections. The observation of the node of the magnetic form factor should be confirmed in a more precise experiment . Together with the determination of the secondary maximum of $`|G_C|`$ , this would complete the full characterization of the deuteron electromagnetic structure up to $`Q7`$ fm<sup>-1</sup>.
###### Acknowledgements.
The authors gratefully acknowledge discussions with A. Kobushkin and I. Sick. This work was supported by the French Centre National de la Recherche Scientifique and Commissariat à l’Energie Atomique, the Swiss National Science Foundation, the U.S. Department of Energy and National Science Foundation, and the K.C. Wong Foundation. |
warning/0002/hep-th0002200.html | ar5iv | text | # 1. Introduction
## 1. Introduction
In this paper we continue our studies of multidimensional gravitational models based on $`D`$-dimensional Einstein equations with fields of antisymmetric forms of arbitrary rank (see \[1–3\] and references therein) as some low-energy limit of a future unified model (M-, F- or other type). Our main interest here will be in the stability properties of multidimensional black-hole (BH) and non-BH solutions with nonzero fields of forms, associated with charged $`p`$-branes. There exist a large number of such solutions in arbitrary dimensions — see e.g. \[2, 4–9\] and references therein. They are important in connection with studies of processes at early stages of the Universe, counts of micro-states in BH thermodynamics and now especially due to new developments in M-theory related to the AdS/CFT correspondence . For recent reviews of this rapidly developing field see, e.g., .
BH stability studies have a long history, of which we will only mention (more or less arbitrarily) some milestones, concerning spherically symmetric backgrounds. Regge and Wheeler considered the stability of the Schwarzschild space-time and developed the formalism of spherical harmonics for metric perturbations. Vishveshwara finally proved the linear stability of Schwarzschild BHs; Moncrief did the same for Reissner-Nordström ones. BHs with a conformally coupled scalar field were shown to be unstable under spherically symmetric perturbations , as well as minimally coupled scalar field configurations in general relativity possessing naked singularities . The monopole degree of freedom is present there due to the scalar field; it was argued that monopole perturbations were most likely to be unstable due to the absence of centrifugal terms in the effective potentials; catastrophic instabilities were indeed found and it was unnecessary to study other multipoles. On the other hand, coloured BHs, containing non-Abelian gauge fields, were shown to be, in general, unstable due to their sphaleronic degrees of freedom — see and references therein. A recent overview of 4-dimensional perturbation studies may be found in Ref. .
For BHs in multidimensional theories of gravity the situation is more complex since, on the one hand, there emerge new effective scalar fields (extra-dimension scale factors, sometimes called moduli fields) in the external space-time, and, on the other, instabilities may be caused by waves in extra dimensions. Instabilities of the latter kind were indeed found by Gregory and Laflamme for a limited class of neutral and charged black strings and branes, having a constant internal space scale factor. Furthermore, it was argued that compactification on a sufficiently small length scale should prevent the onset of instability, and, moreover, that extremal black branes are stable . It was concluded that only very light BHs, whose horizon radii have the same order of magnitude as their extra dimensions, manifest this form of instability.
It is therefore of interest to inquire whether or not there are other forms of instability, maybe “more dangerous”, on more general backgrounds, containing nontrivial internal space structures and/or several dilatonic scalars and brane charges. As was previously the case with backgrounds containing effective scalar fields, it is natural to consider first the simplest, monopole perturbations.
Earlier we analyzed the stability of static, spherically symmetric solutions to the Einstein-Maxwell-scalar equations with a dilatonic type coupling between scalar and electromagnetic fields in $`D`$-dimensional gravity . It was proved there that only BH configurations were stable under linear spherically symmetric perturbations, while non-BH solutions turned out to be catastrophically unstable. A similar result was obtained for dilatonic BHs with the inclusion of the Gauss-Bonnet curvature term due to one-loop quantum corrections . We will now show that in the simplest case of a single charged black brane the solution is stable under linear spherically symmetric perturbations, whereas single-brane solutions with naked singularities are unstable. So the results of are generalized.
We also present a tentative consideration of multi-brane BHs and conclude that in cases when the perturbation equations decouple, the stability conclusion is also valid. Two classes of such systems are indicated, both characterized by certain relations among brane charges, such that, in terms of Sec. 3, the constituent vectors $`\stackrel{}{Y}_s`$ form a single block of a block-orthogonal system (BOS) — single-block BHs for short. Namely, the stability is proved for arbitrary two-brane single-block BHs and multi-brane single-block BHs with mutually orthogonal vectors $`\stackrel{}{Y}_s`$ (see the details in Sec. 6.2). For many single-block configurations which do not belong to these classes, the stability can be proved as well, but their properties require individual studies; see an example in the Appendix, Eqs. (A.4)–(A.7). There are, however, numerous multi-brane BHs for which decoupling is impossible and one may expect that some of them show a new type of instability connected with mode interaction; a study of these systems is in progress.
The paper is organized as follows. Sec. 2 describes the general features of the field model to be considered. Sec. 3 presents some known static solutions, including BHs, on the basis of the target space $`𝕍`$ connected with dimensional reduction. In Sec. 4 a truncated target space $`\overline{𝕍}`$, more appropriate for treating the perturbations, is introduced, and wave equations for perturbations are derived. In Sec. 5 the stability properties of single-brane configurations are deduced, while in Sec. 6 the stability of some multi-brane BHs under spherically symmetric perturbations is established. The Appendix gives some examples from 11-dimensional supergravity.
The word “stable” throughout the paper means “stable under linear spherically symmetric perturbations”.
## 2. The model
Our starting point is, as in Refs. \[1–8\], the model action for $`D`$-dimensional gravity with several scalar dilatonic fields $`\phi ^a`$ and antisymmetric $`n_s`$-forms $`F_s`$:
$$S=\underset{𝕄}{}d^Dz\sqrt{|g|}\left\{[g]\delta _{ab}g^{MN}_M\phi ^a_N\phi ^b\underset{s𝒮}{}\frac{\eta _s}{n_s!}\mathrm{e}^{2\lambda _{sa}\phi ^a}F_s^2\right\},$$
(1)
in a pseudo-Riemannian manifold $`𝕄=_u\times 𝕄{}_{0}{}^{}\times \mathrm{}\times 𝕄_n`$, with factor space dimensions $`d_i`$, $`i=0,\mathrm{},n`$; $``$ is the scalar curvature. We will assume $`𝕄`$ to be spherically symmetric , so that the metric is
$`ds_D^2=g_{MN}dz^Mdz^N`$ $`=`$ $`\mathrm{e}^{2\alpha ^0}du^2+{\displaystyle \underset{i=0}{\overset{n}{}}}\mathrm{e}^{2\beta ^i}ds_i^2`$ (2)
$`=`$ $`\mathrm{e}^{2\alpha ^0}du^2+\mathrm{e}^{2\beta ^0}d\mathrm{\Omega }^2\mathrm{e}^{2\beta ^1}dt^2+{\displaystyle \underset{i=2}{\overset{n}{}}}\mathrm{e}^{2\beta ^i}ds_i^2.`$
Here $`u`$ is a radial coordinate ranging in $`_u`$; $`ds_0^2=d\mathrm{\Omega }^2`$ is the metric on a unit $`d_0`$-dimensional sphere $`𝕄{}_{0}{}^{}=S^{d_0}`$; $`t𝕄{}_{1}{}^{}_t`$ is time; the metrics $`g^i=ds_i^2`$ of the “extra” factor spaces ($`i2`$) are assumed to be Ricci-flat and can have arbitrary signatures $`\epsilon _i=signg^i`$; $`|g|=|detg_{MN}|`$ and similarly for subspaces; $`F_s^2=F_{s,M_1\mathrm{}M_{n_s}}F_s^{M_1\mathrm{}M_{n_s}}`$; $`\lambda _{sa}`$ are coupling constants; $`\eta _s=\pm 1`$ (to be specified later); $`s𝒮`$, $`a𝒜`$, where $`𝒮`$ and $`𝒜`$ are some finite sets. The “scale factors” $`\mathrm{e}^{\beta ^i}`$ and the scalars $`\phi ^a`$ are assumed to depend on $`u`$ and $`t`$ only.
The $`F`$-forms should be also compatible with spherical symmetry. A given $`F`$-form may have several essentially (non-permutatively) different components; such a situation is sometimes called “composite $`p`$-branes”<sup>3</sup><sup>3</sup>3There is an exception: two components, having only one noncoinciding index, cannot coexist since in this case there emerge nonzero off-block-diagonal components of the energy-momentum tensor (EMT) $`T_M^N`$, while the Einstein tensor in the l.h.s. of the Einstein equations is block-diagonal. See more details in Ref. .. For convenience, we will nevertheless treat essentially different components of the same $`F`$-form as individual (“elementary”) $`F`$-forms. A reformulation to the composite ansatz, if needed, is straightforward.
Each $`n_s`$-form $`F=dA_{[M_1}A_{M_2\mathrm{}M_{n_s}]}dz^{M_1}\mathrm{}dz^{M_{n_s}}`$ is then associated with a certain subset $`I=\{i_1,\mathrm{},i_k\}`$ ($`i_1<\mathrm{}<i_k`$) of the set of numbers labelling the factor spaces: $`\{i\}=I_0=\{0,\mathrm{},n\}`$. The forms $`F_s`$ are naturally classified as electric ($`F_{\mathrm{e}I}`$) and magnetic ($`F_{\mathrm{m}I}`$) ones. By definition, the potential $`A_I`$ of an electric form $`F_{\mathrm{e}I}`$ carries the coordinate indices of the subspaces $`𝕄{}_{i}{}^{},iI`$ and is $`u`$-dependent (since only a radial component of the field may be nonzero). A magnetic form $`F_{\mathrm{m}I}`$ is built as a form dual to a possible electric one, and its nonzero components carry coordinate indices of the subspaces $`𝕄{}_{i}{}^{},i\overline{I}\stackrel{\mathrm{def}}{=}I_0I`$, One can write:
$$n_{\mathrm{e}I}=rankF_{\mathrm{e}I}=d(I)+1,n_{\mathrm{m}I}=rankF_{\mathrm{m}I}=DrankF_{\mathrm{e}I}=d(\overline{I})$$
(3)
where $`d(I)=_{iI}d_i`$ are the dimensions of the subspaces $`𝕄{}_{I}{}^{}=𝕄{}_{i_1}{}^{}\times \mathrm{}\times 𝕄_{i_k}`$. The index $`s`$ will be used to jointly describe the two types of forms, so that
$$𝒮=\{s\}=\{\mathrm{e}I_s\}\{\mathrm{m}I_s\}.$$
(4)
We will make some more assumptions to assure that all $`F`$-forms behave like genuine electric or magnetic fields in the physical subspace $`𝕄{}_{\mathrm{phys}}{}^{}=_u\times _t\times 𝕄_0`$, namely:
$`(𝐢)`$ $`1I_s,s\text{(the subspaces }M_{I_s}\text{ contain the time axis }_t\text{)};`$ (5)
$`(\mathrm{𝐢𝐢})`$ $`0I_s,s\text{(the branes only “live” in extra dimensions)};`$ (6)
$`(\mathrm{𝐢𝐢𝐢})`$ $`T_t^t(F_s)>0,s\text{(the energy density is positive).}`$ (7)
By (i), the so-called quasiscalar forms (forms with $`1I_s`$, behaving as effective scalar or pseudoscalar fields in $`𝕄_{\mathrm{phys}}`$) are excluded. The reason for adopting (i) is that our interest here is mostly in BHs which do not admit nonzero quasiscalar forms (the no-hair theorem for brane systems ).
Assumption (iii) holds if all extra dimensions are spacelike ($`\epsilon _i=1`$, $`i2`$) and in (1) all $`\eta _s=1`$. In more general models, with arbitrary $`\epsilon _i`$, (iii) holds if
$`\eta _{\mathrm{e}I_s}=\epsilon (I_s),\eta _{\mathrm{m}I_s}=\epsilon (\overline{I}_s),\epsilon (I)\stackrel{\mathrm{def}}{=}{\displaystyle \underset{iI}{}}\epsilon _i.`$ (8)
We will consider static configurations and their small (linear) time-dependent perturbations. It turns out, however, that under the above assumptions the Maxwell-like field equations for $`F_s`$ may be integrated in a general form for their arbitrary dependence on $`u`$ and $`t`$. Indeed, for an electric $`m`$-form $`F_s`$ ($`s=\mathrm{e}I`$, $`m=d(I_s)+1`$) the field equations due to (1)
$$\left(\genfrac{}{}{0pt}{}{_u}{_t}\right)\left(F_s^{utM_3\mathrm{}M_m}\sqrt{|g|}\mathrm{e}^{2\lambda _{sa}\phi ^a}\right)=0$$
(9)
are easily integrated to give
$$F_s^{utM_3\mathrm{}M_m}=Q_s\mathrm{e}^{\alpha ^0\sigma _02\lambda _{sa}\phi ^a}\epsilon ^{M_3\mathrm{}M_{d(I)}}/\sqrt{|g_I|}\frac{1}{m!}F_s^2=\epsilon (I)Q_s^2\mathrm{e}^{2\sigma (\overline{I})2\lambda _{sa}\phi ^a}.$$
(10)
where $`\epsilon ^{\mathrm{}}`$ and $`\epsilon _{\mathrm{}}`$ are Levi-Civita symbols, $`|g_I|=_{iI}|g^i|`$, and $`Q_s=\mathrm{const}`$ are charges. In a similar way, for a magnetic $`m`$-form $`F_s`$ ($`s=\mathrm{m}I`$, $`m=d(\overline{I}_s)`$), the field equations and the Bianchi identities $`dF_s=0`$ lead to
$$F_{s,M_1\mathrm{}M_{d(\overline{I})}}=Q_s\epsilon _{M_1\mathrm{}M_{d(\overline{I})}}\sqrt{|g_{\overline{I}}|}\frac{1}{m!}F_s^2=\epsilon (\overline{I})Q_s^2\mathrm{e}^{2\sigma (\overline{I})+2\lambda _{sa}\phi ^a}.$$
(11)
We use the notations
$$\sigma _i=\underset{j=i}{\overset{n}{}}d_j\beta ^j(u,t),\sigma (I)=\underset{iI}{}d_i\beta ^i(u,t).$$
(12)
Evidently, the expressions (10) and (11) differ only in the signs before $`\lambda _{sa}`$ and the signature-dependent prefactors $`\epsilon `$. Due to (8), their energy-momentum tensors (EMTs) coincide up to the replacement $`\lambda _{sa}\lambda _{sa}`$, and their further treatment is quite identical. In what follows we therefore mostly speak of electric forms, but the results are easily reformulated for any sets of electric and magnetic forms. We also assume that all $`Q_s0`$.
## 3. Static systems
### 3.1. The target space $`𝕍`$
Under the above assumptions, the system is well described using the so-called $`\sigma `$ model representation ), to be briefly outlined here as applied to static, spherically symmetric systems. This formulation can be derived by reducing the action (1) to the $`(d_0+1)`$-dimensional space $`_u\times 𝕄_0`$.
As in and many later papers, we choose the harmonic $`u`$ coordinate ($`^M_Mu=0`$), such that
$$\alpha ^0(u)=\sigma _0(u)\underset{i=0}{\overset{n}{}}d_i\beta ^i.$$
(13)
Due to (6), the combination $`\left(\genfrac{}{}{0pt}{}{1}{1}\right)+\left(\genfrac{}{}{0pt}{}{2}{2}\right)`$ of the Einstein equations has a Liouville form and is integrated giving
$`\mathrm{e}^{\beta ^0\alpha ^0}=(d_01)s(k,u),s(k,u)\stackrel{\mathrm{def}}{=}\{\begin{array}{cc}k^1\mathrm{sinh}ku,\hfill & h>0,\hfill \\ u,\hfill & h=0,\hfill \\ k^1\mathrm{sin}ku,\hfill & h<0.\hfill \end{array}`$ (17)
where $`k`$ is an integration constant. With (17) the $`D`$-dimensional line element may be written in the form ($`\overline{d}\stackrel{\mathrm{def}}{=}d_01`$)
$$ds_D^2=\frac{\mathrm{e}^{2\sigma _1/\overline{d}}}{[\overline{d}s(k,u)]^{2/\overline{d}}}\left\{\frac{du^2}{[\overline{d}s(k,u)]^2}+d\mathrm{\Omega }^2\right\}\mathrm{e}^{2\beta ^1}dt^2+\underset{i=2}{\overset{n}{}}\mathrm{e}^{2\beta ^i}ds_i^2.$$
(18)
The $`u`$ coordinate is defined for $`0<u<u_{\mathrm{max}}`$ where $`u=0`$ corresponds to spatial infinity while $`u_{\mathrm{max}}`$ may be finite or infinite depending on the form of a particular solution.
The remaining set of unknowns $`\beta ^i(u),\phi ^a(u)`$ ($`i=1,\mathrm{},n,a𝒜`$) can be treated as a real-valued vector function $`x^A(u)`$ (so that $`\{A\}=\{1,\mathrm{},n\}𝒜`$) in an $`(n+|𝒜|)`$-dimensional vector space $`𝕍`$ (target space). The field equations for $`x^A`$ can be derived from the Toda-like Lagrangian
$`L=G_{AB}x_u^Ax_u^B+{\displaystyle \underset{s}{}}Q_s^2\mathrm{e}^{2y_s(u)}{\displaystyle \underset{i=1}{\overset{n}{}}}d_i(\beta _u^i)^2+{\displaystyle \frac{\sigma _{1,u}^2}{d_01}}+\delta _{ab}\phi _u^a\phi _u^b+{\displaystyle \underset{s}{}}Q_s^2\mathrm{e}^{2y_s(u)}`$ (19)
(the subscript $`u`$ means $`d/du`$), with the “energy” constraint
$$E=G_{AB}x_u^Ax_u^B\underset{s}{}Q_s^2\mathrm{e}^{2y_s}=\frac{d_0}{d_01}k^2signk.$$
(20)
The nondegenerate symmetric matrix
$$(G_{AB})=\left(\begin{array}{cc}d_id_j/\overline{d}+d_i\delta _{ij}& 0\\ 0& \delta _{ab}\end{array}\right)$$
(21)
defines a positive-definite metric in $`𝕍`$; the functions $`y_s(u)`$ are scalar products:
$$y_s=\sigma (I_s)\lambda _{sa}\phi ^aY_{s,A}x^A,(Y_{s,A})=(d_i\delta _{iI_s},\lambda _{sa}),$$
(22)
where $`\delta _{iI}=1`$ if $`iI`$ and $`\delta _{iI}=0`$ otherwise. The contravariant components and scalar products of the vectors $`\stackrel{}{Y}_s`$ are found using the matrix $`G^{AB}`$ inverse to $`G_{AB}`$:
$`(G^{AB})=\left(\begin{array}{cc}\delta ^{ij}/d_i1/(D2)& 0\\ 0& \delta ^{ab}\end{array}\right),(Y_s{}_{}{}^{A})=(\delta _{iI_s}{\displaystyle \frac{d(I_s)}{D2}},\lambda _{sa});`$ (23)
$`Y_{s,A}Y_s^{}{}_{}{}^{A}\stackrel{}{Y}{}_{s}{}^{}\stackrel{}{Y}{}_{s^{}}{}^{}=d(I_sI_s^{}){\displaystyle \frac{d(I_s)d(I_s^{})}{D2}}+\lambda _{as}\lambda _{as^{}}.`$ (24)
The equations of motion in terms of $`\stackrel{}{Y}_s`$ read
$$\ddot{x}{}_{}{}^{A}=\underset{s}{}Q_s^2Y_s{}_{}{}^{A}\mathrm{e}_{}^{2y_s(u)}.$$
(25)
### 3.2. Exact solutions: orthogonal systems (OS)
The integrability of the Toda-like system (19) depends on the set of vectors $`\stackrel{}{Y}_s`$. In many cases general or special solutions to Eqs. (25) are known. Here we will mention the simplest case of integrability: a general solution is available if all $`\stackrel{}{Y}_s`$ are mutually orthogonal in $`𝕍`$ , that is,
$$\stackrel{}{Y}{}_{s}{}^{}\stackrel{}{Y}{}_{s^{}}{}^{}=\delta _{ss^{}}Y_s^2,Y_s^2=d(I)[1d(I)/(D2)]+\lambda _s^2>0$$
(26)
where $`\lambda _s^2=_a\lambda _{sa}^2`$. Then the functions $`y_s(u)`$ obey the decoupled Liouville equations $`y_{s,uu}=Q_s^2Y_s^2\mathrm{e}^{2y_s}`$, whence
$$\mathrm{e}^{2y_s(u)}=Q_s^2Y_s^2s^2(h_s,u+u_s)$$
(27)
where $`h_s`$ and $`u_s`$ are integration constants and the function $`s(.,.)`$ has been defined in (17). For the sought functions $`x^A(u)`$ and the “conserved energy” $`E`$ we then obtain:
$`x^A(u)`$ $`=`$ $`{\displaystyle \underset{s}{}}{\displaystyle \frac{Y_s^A}{Y_s^2}}y_s(u)+c^Au+\underset{¯}{c}^A,`$ (28)
$`E`$ $`=`$ $`{\displaystyle \underset{s}{}}{\displaystyle \frac{h_s^2signh_s}{Y_s^2}}+\stackrel{}{c}{}_{}{}^{2}={\displaystyle \frac{d_0}{d_01}}k^2signk,`$ (29)
where the vectors of integration constants $`\stackrel{}{c}`$ and $`\underset{¯}{\overset{}{c}}`$ are orthogonal to all $`\stackrel{}{Y}_s`$: $`c^AY_{s,A}=\underset{¯}{c}^AY_{s,A}=0`$, or
$$c^id_i\delta _{iI_s}c^a\lambda _{sa}=0,\underset{¯}{c}^id_i\delta _{iI_s}\underset{¯}{c}^a\lambda _{sa}=0.$$
(30)
### 3.3. Exact solutions: block-orthogonal systems (BOS)
The above OS solutions are general for input parameters ($`D`$, $`d_i`$, $`\stackrel{}{Y}_s`$) satisfying Eq. (26): there is an independent charge attached to each (elementary) $`F`$-form. One can, however, obtain special solutions for more general sets of input parameters, under less restrictive conditions than (26). Namely, assuming that some of the functions $`y_s(u)`$ (22) coincide, one obtains the so-called BOS solutions , where the number of independent charges coincides with the number of different functions $`y_s(u)`$.
Indeed, suppose that the set $`𝒮`$ splits into several non-intersecting non-empty subsets,
$$𝒮=\underset{\omega }{}𝒮_\omega ,|𝒮_\omega |=m(\omega ),$$
(31)
such that the vectors $`\stackrel{}{Y}_{\mu (\omega )}`$ ($`\mu (\omega )𝒮_\omega `$) form mutually orthogonal subspaces $`𝕍_\omega 𝕍`$:
$$\stackrel{}{Y}{}_{\mu (\omega )}{}^{}\stackrel{}{Y}{}_{\nu (\omega ^{})}{}^{}=0,\omega \omega ^{}.$$
(32)
Then the corresponding result from can be formulated as follows:
Proposition 1. Let, for each fixed $`\omega `$, all $`\stackrel{}{Y}{}_{\nu }{}^{}𝕍_\omega `$ be linearly independent, and let there be a vector $`\stackrel{}{Y}{}_{\omega }{}^{}=_{\mu 𝒮_\omega }a_\mu \stackrel{}{Y}_\mu `$ with $`a_\mu >0`$ such that
$$\stackrel{}{Y}{}_{\mu }{}^{}\stackrel{}{Y}{}_{\omega }{}^{}=Y_\omega ^2\stackrel{\mathrm{def}}{=}\stackrel{}{Y}{}_{\omega }{}^{2},\mu 𝒮_\omega .$$
(33)
Then one has the following solution to the equations of motion (25), (20):
$`x^A`$ $`=`$ $`{\displaystyle \underset{\omega }{}}{\displaystyle \frac{Y_\omega ^A}{Y_\omega ^2}}y_\omega (u)+c^Au+\underset{¯}{c}^A,`$ (34)
$`\mathrm{e}^{2y_\omega }`$ $`=`$ $`\widehat{q}_\omega Y_\omega ^2s^2(h_\omega ,u+u_\omega ),\widehat{q}_\omega \stackrel{\mathrm{def}}{=}{\displaystyle \underset{\mu 𝒮_\omega }{}}Q_\mu ^2,`$ (35)
$`E`$ $`=`$ $`{\displaystyle \underset{\omega }{}}{\displaystyle \frac{h_\omega ^2signh_\omega }{Y_\omega ^2}}+\stackrel{}{c}{}_{}{}^{2}={\displaystyle \frac{d_0}{d_01}}k^2signk`$ (36)
where $`h_\omega `$, $`u_\omega `$, $`c^A`$ and $`\underset{¯}{c}^A`$ are integration constants; $`c^A`$ and $`\underset{¯}{c}^A`$ are constrained by the orthogonality relations (30) (so that the vectors $`\stackrel{}{c}`$ and $`\underset{¯}{\overset{}{c}}`$ are orthogonal to each individual vector $`\stackrel{}{Y}{}_{s}{}^{}𝕍`$).
Eqs. (33) form a set of linear algebraic equations with respect to the “charge factors” $`a_\nu =Q_\nu ^2/\widehat{q}_\omega `$, satisfying the condition $`_{\mu 𝒮_\omega }a_\mu =1`$. A solution to (33) for given $`\stackrel{}{Y}_\mu `$ can contain some $`a_\mu <0`$; according to , this would mean that such a $`p`$-brane is “quasiscalar”, violating the assumption (5). Solutions with such branes are possible but are rejected here since they do not lead to black holes. Furthermore, if a solution to (33) gives $`a_\mu =0`$ for some $`\mu 𝒮_\omega `$, this means that the block cannot contain such a $`p`$-brane, and then the consideration may be repeated without it.<sup>4</sup><sup>4</sup>4Geometrically, the vector $`\stackrel{}{Y}_\omega `$ solving Eqs. (33) is the altitude of the pyramid formed by the vectors $`\stackrel{}{Y}_\mu `$, $`\mu 𝒮_\omega `$ with a common origin. The condition $`a_\mu >0`$ means that this altitude is located inside the pyramid, while $`a_\mu =0`$ means that the altitude belongs to one of its faces.
The function $`y_\omega (u)`$ is equal to $`y_{\mu (\omega )}(u)=Y_{\mu (\omega ),A}x^A`$, which is, due to (33), the same for all $`\mu 𝒮_\omega `$. The BOS solution generalizes the OS one, (27), (28): the latter is restored when each block contains a single $`F`$-form.
Both kinds of solutions are asymptotically flat, and it is natural to normalize the functions $`y_s(u)`$ and $`y_\omega (u)`$ by the condition $`y_s(0)=0`$ or $`y_\omega (0)=0`$, so that the constants $`u_s`$ and $`u_\omega `$ are directly related to the charges.
Other solutions to the equations of motion are known, connected with Toda chains and Lie algebras .
### 3.4. Black-hole solutions
Black holes (BHs) are distinguished among other spherically symmetric solutions by the existence of horizons instead of singularities in the physical $`d_0+2`$-dimension space $`𝕄_{\mathrm{phys}}`$; the extra dimensions and scalar fields are also required to be well-behaved on the horizon to provide regularity of the $`D`$-dimensional metric. Thus BHs are described by the above solutions under certain constraints upon the integration constants. Namely, for the BOS solution (33)–(36), requiring that all the scale factors $`\mathrm{e}^{\beta ^i}`$ (except $`\mathrm{e}^{\beta ^1}=\sqrt{|g_{tt}|}`$ which should tend to zero) and scalars $`\phi ^a`$ tend to finite limits as $`uu_{\mathrm{max}}`$, we get :
$`h_\omega =k>0,\omega ;c^A=k{\displaystyle \underset{s}{}}Y_\omega ^2Y_\omega {}_{}{}^{A}k\delta _1^A`$ (37)
where $`A=1`$ corresponds to $`i=1`$ (time). The constraint (29) then holds automatically. The value $`u=u_{\mathrm{max}}=\mathrm{}`$ corresponds to the horizon. The same condition for the OS solution (27)–(30) is obtained by replacing $`\omega s`$.
Under the asymptotic conditions $`\phi ^a0`$, $`\beta ^i0`$ as $`u0`$, after the transformation
$`\mathrm{e}^{2ku}=1{\displaystyle \frac{2k}{\overline{d}r^{\overline{d}}}},\overline{d}\stackrel{\mathrm{def}}{=}d_01`$ (38)
the metric (18) for BHs and the corresponding scalar fields may be written as
$`ds_D^2=\left({\displaystyle \underset{\omega }{}}H_\omega ^{A_\omega }\right)\left[dt^2\left(1{\displaystyle \frac{2k}{\overline{d}r^{\overline{d}}}}\right){\displaystyle \underset{\omega }{}}H_\omega ^{2/Y_\omega ^2}+\left({\displaystyle \frac{dr^2}{12k/(\overline{d}r^{\overline{d}})}}+r^2d\mathrm{\Omega }^2\right)+{\displaystyle \underset{i=2}{\overset{n}{}}}ds_i^2{\displaystyle \underset{\omega }{}}H_\omega ^{A_\omega ^i}\right];`$
$`A_\omega \stackrel{\mathrm{def}}{=}{\displaystyle \frac{2}{Y_\omega ^2}}{\displaystyle \underset{\mu 𝒮_\omega }{}}{\displaystyle \frac{a_\mu d(I_\mu )}{D2}}\stackrel{\mathrm{OS}}{=}{\displaystyle \frac{2}{Y_s^2}}{\displaystyle \frac{d(I_s)}{D2}};A_\omega ^i\stackrel{\mathrm{def}}{=}{\displaystyle \frac{2}{Y_\omega ^2}}{\displaystyle \underset{\mu 𝒮_\omega }{}}a_\mu \delta _{iI_\mu }\stackrel{\mathrm{OS}}{=}{\displaystyle \frac{2}{Y_s^2}}\delta _{iI_s};`$ (39)
$`\phi ^a={\displaystyle \underset{\omega }{}}{\displaystyle \frac{1}{Y_\omega ^2}}\mathrm{ln}H_\omega {\displaystyle \underset{\mu 𝒮_\omega }{}}a_\mu \lambda _{\mu a}\stackrel{\mathrm{OS}}{=}{\displaystyle \underset{s}{}}{\displaystyle \frac{\lambda _{sa}}{Y_s^2}}\mathrm{ln}H_s,`$ (40)
where $`\stackrel{\mathrm{OS}}{=}`$ means “equal for OS, with $`\omega s`$”, and $`H_\omega `$ are harmonic functions in $`_+\times S^{d_0}`$:
$$H_\omega (r)=1+p_\omega /(\overline{d}r^{\overline{d}}),p_\omega \stackrel{\mathrm{def}}{=}\sqrt{k^2+\widehat{q}_\omega Y_\omega ^2}k.$$
(41)
The subfamily (37), (39)–(41) exhausts all BOS BH solutions; the OS ones are obtained in the special case of each block $`𝒮_\omega `$ consisting of a single element $`s`$.
The above relations describe the so-called non-extremal BHs. Extremal ones, corresponding to minimum black hole mass for given charges (the so-called BPS limit), are obtained in the limit $`k0`$. The same solutions follow directly from (36)–(35) under the conditions $`h_\omega =k=c^A=0`$. For $`k=0`$, the solution is defined in the whole range $`r>0`$, while $`r=0`$ in many cases corresponds to a naked singularity rather than an event horizon, so that we no more deal with a black hole. However, in many other important cases $`r=0`$ is an event horizon of extremal Reissner-Nordström type, with an AdS near-horizon geometry; some examples are mentioned in the Appendix.
Other families of solutions, mentioned at the end of the previouis section, also contain BH subfamilies. The most general BH solutions are considered in Ref. .
## 4. Perturbation equations
### 4.1. Truncated target space $`\overline{𝕍}`$
Consider now nonstatic spherically symmetric configurations corresponding to the action (1) with the metric (2) and all field variables depending on $`u`$ and $`t`$. As before, we are dealing with true electric and magnetic forms $`F_s`$, so that their $`I_s1`$, or
$$I_s=1J_s,J_s\{2,\mathrm{},n\}.$$
(42)
As in Refs. , it is helpful to pass to the Einstein frame in the physical $`(d_0+2)`$-dimensional space-time $`𝕄{}_{\mathrm{phys}}{}^{}=_u\times 𝕄{}_{0}{}^{}\times 𝕄_1`$. The action (1) is then rewritten in terms of the metric $`g_{\mu \nu }`$, the $`d_0+2`$-dimensional part of $`g_{MN}`$, and is transformed to the Einstein frame in $`𝕄_{\mathrm{phys}}`$ with the metric
$$\overline{g}_{\mu \nu }=\mathrm{e}^{2\sigma _2/d_0}g_{\mu \nu }.$$
(43)
The electric $`n_s`$-forms are re-parametrized as follows:
$`F_{utM_3\mathrm{}M_{n_s}}=\underset{s}{𝐹}{}_{ut}{}^{},{\displaystyle \frac{1}{n_s!}}F_s^2={\displaystyle \frac{1}{2}}\underset{s}{𝐹}{}_{\mu \nu }{}^{}\underset{s}{𝐹}{}_{}{}^{\mu \nu }\mathrm{e}_{}^{2\sigma (J_s)}=Q_s^2\mathrm{e}^{2\lambda _{sa}\phi ^a},`$ (44)
where the indices $`M_3,\mathrm{}M_{n_s}`$ belong to $`J_s`$; here and henceforth the indices $`\mu ,\nu `$ are raised and lowered using the metric $`\overline{g}_{\mu \nu }`$; in the last equality the solution (10) and the positive energy assumption (7) are taken into account.
The action (1) is written in terms of $`\overline{g}_{\mu \nu }`$ and $`\underset{s}{𝐹}_{\mu \nu }`$ as follows (up to a constant prefactor, connected with the volume of extra dimensions, and a subtracted total divergence):
$`S`$ $`=`$ $`{\displaystyle \underset{𝕄_{\mathrm{phys}}}{}}d^{d_0+2}z\sqrt{|\overline{g}|}\{[\overline{g}]{\displaystyle \frac{1}{d_0}}(\sigma _2)^2{\displaystyle \underset{i=2}{\overset{n}{}}}d_i(\beta ^i)^2`$ (45)
$`\delta _{ab}(\phi ^a,\phi ^b){\displaystyle \frac{1}{2}}{\displaystyle \underset{s𝒮}{}}\underset{s}{𝐹}{}_{\mu \nu }{}^{}\underset{s}{𝐹}{}_{}{}^{\mu \nu }\mathrm{e}_{}^{2\sigma _2/d_02\sigma (J_s)+2\lambda _{sa}\phi ^a}\}`$
$`=`$ $`{\displaystyle \underset{𝕄_{\mathrm{phys}}}{}}d^{d_0+2}z\sqrt{|\overline{g}|}\left\{[\overline{g}]H_{KL}(x^K,x^L){\displaystyle \frac{1}{2}}{\displaystyle \underset{s𝒮}{}}\underset{s}{𝐹}{}_{\mu \nu }{}^{}\underset{s}{𝐹}{}_{}{}^{\mu \nu }\mathrm{e}_{}^{2Z_{s,K}x^K}\right\}`$
where $`(f,g)=\overline{g}^{\mu \nu }_\mu f_\nu g`$, $`(f)^2=(f,f)`$; the non-degenerate symmetric matrix
$$(H_{KL})=\left(\begin{array}{cc}d_id_j/d_0+d_i\delta _{ij}& 0\\ 0& \delta _{ab}\end{array}\right)$$
(46)
defines a positive-definite metric in the vector space $`\overline{𝕍}`$ (truncated target space) parametrized by the variables $`(x^K)=(\beta ^2,\mathrm{},\beta ^n;\phi ^a)`$; the constant vectors $`\overline{Z}{}_{s}{}^{}\overline{𝕍}`$ are characterized by the components<sup>5</sup><sup>5</sup>5We will use the indices $`K,L`$ for quantities specified in $`\overline{𝕍}`$ to distinguish them from those in $`𝕍`$ where the indices $`A,B`$ are used; vectors in $`\overline{𝕍}`$ are marked with overbars, those in $`𝕍`$ by arrows. Scalar products are written as $`\stackrel{}{Y}\stackrel{}{Z}=G_{AB}Y^AZ^B`$ (as before) and $`\overline{Y}\overline{Z}=H_{KL}Y^KZ^L`$.
$$(Z_{s,K})=(d_i\delta _{iJ_s}\frac{d_i}{d_0},\lambda _{sa}),(Z_s{}_{}{}^{K})=(H^{KL}Z_{s,L})=(\delta _{iJ_s}\frac{d(I_s)}{D2},\lambda _{sa})$$
(47)
where the matrix $`(H^{KL})`$ is inverse to $`(H_{KL})`$:
$$(H^{KL})=\left(\begin{array}{cc}\delta ^{ij}/d_i1/(D2)& 0\\ 0& \delta ^{ab}\end{array}\right).$$
(48)
The truncated target space $`\overline{𝕍}`$ may be considered as the hyperplane $`x^1=\sigma _2/d_0`$ in $`𝕍`$, with the metric $`H_{KL}`$ induced by $`G_{AB}`$ (21). The components $`H^{KL}`$ turn out to be the same as $`G^{AB}`$ for $`i1`$; the components $`Y_s^A`$ and $`Z_s^K`$ coincide in the same manner. It is easy to find that for vectors whose $`I_s`$ satisfy (42) (which is always the case in the present paper),
$$\stackrel{}{Y}{}_{s}{}^{}\stackrel{}{Y}{}_{s^{}}{}^{}=\overline{Z}{}_{s}{}^{}\overline{Z}{}_{s^{}}{}^{}+\frac{d_01}{d_0},$$
(49)
whence it follows that, first, when different $`\stackrel{}{Y}_s`$ are mutually orthogonal in $`𝕍`$, the corresponding $`\overline{Z}_s`$ are never mutually orthogonal in $`\overline{𝕍}`$; second, for any $`\stackrel{}{Y}_s`$ whose $`I_s1`$ one has
$$\stackrel{}{Y}{}_{s}{}^{}\stackrel{}{Y}{}_{s}{}^{}Y_s^2>\frac{d_01}{d_0}.$$
(50)
### 4.2. Wave equations
The action (45) may be used to obtain the equations governing small spherically symmetric perturbations of static solutions. The metric (43) in $`𝕄_{\mathrm{phys}}`$ will be written in the form
$$ds_\mathrm{E}^2=\overline{g}_{\mu \nu }dz^\mu dz^\nu =\mathrm{e}^{2\alpha }du^2+\mathrm{e}^{2\beta }d\mathrm{\Omega }^2\mathrm{e}^{2\gamma }dt^2$$
(51)
where “E” stands for the Einstein frame and $`\alpha (u,t),\beta (u,t),\gamma (u,t)`$ are connected with the corresponding quantities from (2) as follows:
$$\alpha =\alpha ^0+\sigma _2/d_0,\beta =\beta ^0+\sigma _2/d_0,\gamma =\beta ^1+\sigma _2/d_0.$$
(52)
Since the field equations for the $`F`$-forms have been integrated — see (10) and (44), the remaining unknowns are $`\alpha ,\beta ,\gamma `$ and $`x^K`$ (that is, $`\beta ^i,i2`$, and $`\phi ^a`$). In what follows we will write
$$\alpha (u,t)=\alpha (u)+\delta \alpha (u,t),$$
where $`\delta \alpha `$ is a small perturbation, and similarly for other unknowns. We accordingly preserve only terms linear in $`\delta \alpha `$ and similar quantities and in time derivatives. The field equations may be written in the form
$`[\overline{g}]x^K={\displaystyle \underset{s}{}}Z_s{}_{}{}^{K}Q_{s}^{2}\mathrm{e}^{2d_0\beta +2\overline{Z}{}_{s}{}^{}\overline{x}},`$ (53)
$`_\mu ^\nu =T_\mu ^\nu {\displaystyle \frac{1}{d_0}}\delta _\mu ^\nu T_\lambda ^\lambda \stackrel{\mathrm{def}}{=}\stackrel{~}{T}_\mu ^\nu `$ (54)
where $`[\overline{g}]=\overline{g}^{\mu \nu }_\mu _\nu `$ is the D’Alembert operator, while for the nonzero components of the EMT $`T_\mu ^\nu `$ corresponding to (45) one has (no summing in $`\mu `$)
$`\stackrel{~}{T}{}_{\mu }{}^{\mu }=\mathrm{e}^{2\alpha }H_{KL}x_u^Kx_u^Ldiag(0,1,0,\mathrm{},0){\displaystyle \underset{s}{}}Q_s^2\mathrm{e}^{2\alpha +2\gamma +2\overline{Z}{}_{s}{}^{}\overline{x}}diag(1{\displaystyle \frac{1}{d_0}},{\displaystyle \frac{1}{d_0}},{\displaystyle \frac{1}{d_0}}\mathrm{},{\displaystyle \frac{1}{d_0}}),`$
$`\stackrel{~}{T}{}_{ut}{}^{}=H_{KL}x_u^Kx_t^L`$ (55)
where $`x_u=_ux`$ and $`x_t=_tx`$; the first and second places under the symbol “diag” belong to $`t`$ and $`u`$, respectively.
As in our previous papers on stability, we use the coordinate freedom in the perturbed space-time and put
$$\delta \beta 0$$
(56)
but preserve the harmonic $`u`$ coordinate condition in the unperturbed (static) space-time<sup>6</sup><sup>6</sup>6This coordinate is harmonic for both metrics (2) and (51); in the latter the coordinate condition has the form $`\alpha =d_0\beta +\gamma `$.. Then Eqs. (53) and (54) give
$`\widehat{L}\delta x^K+x_u^K(\delta \gamma _u\delta \alpha _u)2x_{uu}^K\delta \alpha `$ $`=`$ $`2{\displaystyle \underset{s}{}}Q_s^2\mathrm{e}^{2\gamma +2\overline{Z}{}_{s}{}^{}\overline{x}}Z_s{}_{}{}^{K}Z_{s,L}^{}\delta x^L;\widehat{L}\stackrel{\mathrm{def}}{=}\mathrm{e}^{2d_0\beta }_{tt}+_{uu};`$ (57)
$`d_0\delta \alpha _t`$ $`=`$ $`\overline{x}_u\delta \overline{x}_t;`$ (58)
$`d_0\beta _u(\delta \alpha _u\delta \gamma _u)`$ $`=`$ $`2\overline{x}_{uu}\delta \overline{x}2\beta _{uu}\delta \alpha ,`$ (59)
where (57) follows from the $`({}_{ut}{}^{})`$ component of (54) and (59) from one of the angular components of (54); we have also used the equations valid for static systems, in particular, (25), where, according to the definitions of $`\stackrel{}{Y}`$ and $`\overline{Z}`$, $`y_s(u)=\stackrel{}{Y}{}_{s}{}^{}\stackrel{}{x}=\gamma +\overline{Z}{}_{s}{}^{}\overline{x}`$. Integrating (58) in $`t`$ and omitting the emerging arbitrary function of $`u`$ (since we neglect static perturbations), we obtain
$$d_0\delta \alpha =\overline{x}_u\delta \overline{x}.$$
(60)
Substituting $`\delta \alpha `$ from (60) and $`\delta \alpha _u\delta \gamma _u`$ from (59) into (57), we finally arrive at the set of wave equations for the dynamical degrees of freedom in our system, represented by $`\delta x^K`$:
$`\widehat{L}\delta x^K=2P^K{}_{L}{}^{}\delta x^L,P^K{}_{L}{}^{}={\displaystyle \frac{1}{d_0}}{\displaystyle \frac{d}{du}}\left({\displaystyle \frac{x_u^Kx_{L,u}}{\beta _u}}\right)+{\displaystyle \underset{s}{}}Q_s^2\mathrm{e}^{2y_s(u)}Z_s^KZ_{s,L}.`$ (61)
The stability problem is now reduced to a boundary-value problem for $`\delta x^K(u,t)`$. Namely, if there exists a nontrivial solution to Eqs. (61) satisfying some physically reasonable conditions at the ends of the range of $`u`$, such that $`|\delta x^K|`$ (at least some of them) grow unboundedly with $`t`$, then the static system is unstable. Otherwise it is stable in the linear approximation.
The condition at infinity, $`u=0`$, is evident: the perturbations must vanish,
$$dx^K0\mathrm{as}u0.$$
(62)
It is less evident at $`u=u_{\mathrm{max}}`$ since some of the background static solutions are singular there. As in Refs. and others, dealing with minimally coupled or dilatonic scalar fields, we will use the minimal requirement providing the validity of the perturbation scheme, namely
$$|\delta x^K/x^K|<\mathrm{}.$$
(63)
When the background is regular, this condition requires that the perturbation should be regular as well.
## 5. Stability properties of single-brane solutions
### 5.1. Decoupling cases
Eqs. (61) in general do not decouple. Even in the simplest case when there is only one antisymmetric form $`F`$ (that is, one $`p`$-brane), so that Eqs. (27)–(30) yield the general static solution to the field equations, Eqs. (61) contain various linear combinations of $`\delta x^K`$ with $`u`$-dependent coefficients.
There is, however, an important case when Eqs. (61) do decouple for any configuration of $`𝕄`$ with the metric (2), namely, the single-brane solution (27)–(30) under the condition that the vector $`\overline{c}=(c^K)`$ is parallel to $`\overline{Z}`$ in $`\overline{𝕍}`$: <sup>7</sup><sup>7</sup>7Curiously, they are parallel in $`\overline{𝕍}`$ although the corresponding vectors $`\stackrel{}{c}`$ and $`\stackrel{}{Y}`$ are mutually orthogonal in the surrounding target space $`𝕍`$. For a clear picture, imagine two vectors in 3-dimensional space whose projections onto a plane lie on the same ray.
$$c^K=BY^K/Y^2,B=\mathrm{const}$$
(64)
(here and henceforth in this section we omit the index $`s`$ since, by our assumption, it takes only one value). This condition is automatically valid for the case of utmost interest, BHs with one $`p`$-brane (“a black $`p`$-brane”), which, by (37), corresponds to $`B=k=h0`$.
Due to the collinearity condition (64) and the constraint (29), the constants are now connected by the relation
$$N^{}(h^2signhB^2)=k^2signkB^2$$
(65)
with $`N^{}=Y^2(d_01)/d_0<1`$. It turns out that, besides BHs, the condition (64) is satisfied for some singular solutions whose behaviour is quite generic for the system under study:
1. $`k>0`$, $`h0`$, such that $`u_{\mathrm{max}}=\mathrm{}`$ and a singularity at the centre of symmetry is attractive at least in terms of the metric (51), $`\mathrm{e}^\gamma 0`$;
2. $`h<0`$, so that the solution behaviour is determined by the function $`s(h,u+u_1)=h^1\mathrm{sin}h(u+u_1)`$ in (27) where $`u_1=\mathrm{const}(0,\pi /|h|)`$. In this case the central singularity is repulsive, $`\mathrm{e}^\gamma \mathrm{}`$, of Reissner-Nordström type.
Due to (64), Eqs. (61) take the form
$`\widehat{L}\delta x^K=2Z^K\left[{\displaystyle \frac{1}{d_0}}\left({\displaystyle \frac{f_u^2}{\beta _u}}\right)_uf_{uu}\right](\overline{Z}\delta \overline{x}),f(u)\stackrel{\mathrm{def}}{=}{\displaystyle \frac{y+Bu}{Y^2}}`$ (66)
with $`y(u)`$ determined by (27); the area function $`\beta `$ has the form
$$\beta =\frac{1}{d_01}\left\{\mathrm{ln}[(d_01)s(k,u)]+\frac{1}{N^{}}f(u)Bu+\mathrm{const}\right\}$$
(67)
where the value of the constant is inessential.
Since $`\overline{𝕍}`$ is an $`l`$-dimensional Euclidean space ($`l=n1+|𝒜|`$), there are $`l1`$ linearly independent vectors $`\overline{Z}_{}`$ such that $`\overline{Z}{}_{}{}^{}\overline{Z}=0`$. Therefore the set of wave equations (66) decouples into one equation for $`\overline{Z}\delta \overline{x}`$ and $`l1`$ equations for different $`\overline{Z}{}_{}{}^{}\delta \overline{x}`$:
$`\widehat{L}(\overline{Z}\delta \overline{x})`$ $`=`$ $`U(u)(\overline{Z}\delta \overline{x}),U(u)=2Z^2\left[{\displaystyle \frac{1}{d_0}}\left({\displaystyle \frac{f_u^2}{\beta _u}}\right)_uf_{uu}\right];`$ (68)
$`\widehat{L}(\overline{Z}{}_{}{}^{}\delta \overline{x})`$ $`=`$ $`0.`$ (69)
The static nature of the background solution makes it possible to separate the variables:
$$\overline{Z}\delta \overline{x}=\psi (u)\mathrm{e}^{\mathrm{\Omega }t},\overline{Z}{}_{}{}^{}\delta \overline{x}=\psi ^{}(u)\mathrm{e}^{\mathrm{\Omega }^{}t},$$
(70)
so that Eqs. (68) and (69) lead to
$`\psi _{uu}`$ $`=`$ $`[\mathrm{e}^{2d_0\beta }\mathrm{\Omega }^2+U(u)]\psi ,`$ (71)
$`\psi _{uu}^{}`$ $`=`$ $`\mathrm{e}^{2d_0\beta }\mathrm{\Omega }_{}^{}{}_{}{}^{2}\psi ^{}.`$ (72)
The existence of an admissible solution of any of these equations with a real value of $`\mathrm{\Omega }`$ or $`\mathrm{\Omega }^{}`$ would mean that the perturbation can grow exponentially with time, hence the instability.
It is hard to solve Eqs. (71), (72) in their full range but it is rather easy to assess the asymptotic behaviour of their solutions near $`u=0`$ and $`u_{\mathrm{max}}`$, and this will be sufficient for making stability conclusions.
In particular, for $`u0`$, which corresponds to spatial infinity, one has
$$U(u)0\mathrm{and}\mathrm{e}^{d_0\beta }c_0u^{d_0/\overline{d}},c_0=\overline{d}{}_{}{}^{d_0/\overline{d}}.$$
(73)
The general asymptotic form of solutions to (71) and (72) at small $`u`$ for all cases under study may be written as follows:
$$\psi (\mathrm{or}\psi ^{})=u^{d_0/(2\overline{d})}\left[c_1\mathrm{exp}(c_0\overline{d}\mathrm{\Omega }u^{1/\overline{d}})+c_2\mathrm{exp}(c_0\overline{d}\mathrm{\Omega }u^{1/\overline{d}})\right],c_1,c_2=\mathrm{const}.$$
(74)
The boundary condition (62) then requires that in (74) $`c_1=0`$, and it remains to look at the other end of the $`u`$ range, $`uu_{\mathrm{max}}`$.
### 5.2. Instability of naked singularities
Consider Case 1 of the previous subsection, a “scalar type” singularity. As $`u\mathrm{}`$, the relevant functions of the static solution behave as follows:
$$y=hu+O(1),\beta \frac{u}{Y^2d_0}\frac{N_11}{h+k}(kB)(hB)$$
(75)
where $`N_1=d_0Y^2/\overline{d}>1`$. Since, due to (65), in the present case
$$|B|>k|B|>h\mathrm{and}|B|<k|B|<h,$$
one sees that $`\mathrm{e}^{\beta (u)}0`$ exponentially. The same happens to $`U(u)`$, therefore the asymptotic form of (71) or (72) is simply $`\psi _{uu}=0`$ or $`\psi _{uu}^{}=0`$, so that
$$\psi (\mathrm{or}\psi ^{})=c_3u+c_4$$
(76)
with constants $`c_3`$ and $`c_4`$. On the other hand, the background functions $`x^K`$ also behave as $`\mathrm{const}u`$ as $`u\mathrm{}`$, therefore the second boundary condition (63) is satisfied for any solution (76), including the one joining the solution (74) with $`c_1=0`$ at small $`u`$. We conclude that there are growing modes of perturbations for any $`\mathrm{\Omega }`$, hence the singular solution is catastrophically unstable.
In Case 2, $`h<0`$, we have $`u_{\mathrm{max}}=\pi /|h|u_1<\mathrm{}`$ and the relevant functions in the static solution approach $`u_{\mathrm{max}}`$ in the following way:
$$y(u)\mathrm{ln}(|h|\mathrm{\Delta }u),x^K\frac{Z^K}{Y^2}\mathrm{ln}(|h|\mathrm{\Delta }u),\beta \frac{1}{d_0Y^2}\mathrm{ln}\mathrm{\Delta }u\mathrm{}$$
(77)
where $`\mathrm{\Delta }u=u_{\mathrm{max}}u`$. One can make sure that $`U(u)`$ does not affect the asymptotic behaviour of solutions to Eq. (71) as $`uu_{\mathrm{max}}`$ as compared with that of Eq. (72), and for both one can write:
$$\psi (\mathrm{or}\psi ^{})=c_5+c_6\mathrm{\Delta }u$$
(78)
while the condition (63) only requires $`|\delta x^K/\mathrm{ln}\mathrm{\Delta }u|<\mathrm{}`$. Thus the solution satisfies (63) for any choice of the constants $`c_5,c_6`$, and, as in Case 1, this leads to the instability of the background singular solution.
### 5.3. Stability of black holes
In the BH case it is again hard to solve Eqs. (71), (72), but for our purpose it is sufficient to note that, to realize an instability, a solution should begin with a zero value at $`u=0`$ and tend to a finite limit as $`u\mathrm{}`$. This is evidently impossible for a solution to (72) since $`\psi _{uu}^{}/\psi ^{}>0`$. We conclude that at least the $`\psi ^{}`$ modes of BH perturbations are stable. The same reasoning works for the $`\psi `$ mode provided $`U(u)0`$ for all $`u>0`$. Let us pass to the variable $`R=r^{\overline{d}}/\overline{d}`$ in the expression for $`U`$ in (68), so that $`0<u<\mathrm{}`$ corresponds to $`\mathrm{}>R>2k`$:
$`U`$ $`=`$ $`{\displaystyle \frac{2Z^2pR(R2k)}{Y^2(R+p)^2}}\left\{2k+p\left[1{\displaystyle \frac{N^{}R^2}{(R+p^{})^2}}\right]+pN^{}{\displaystyle \frac{4kR+2k(p+p^{})+pp^{}}{(R+p^{})^2}}\right\},`$ (79)
$`N^{}\stackrel{\mathrm{def}}{=}{\displaystyle \frac{1}{Y^2}}{\displaystyle \frac{d_01}{d_0}}<1,p^{}\stackrel{\mathrm{def}}{=}p(1N^{}),`$
where we have used the explicit form of the single-brane BH solution (27)–(37) and the substitution (38) with $`R=r^{\overline{d}}/\overline{d}`$, replacing $`Y_sY`$, $`p_sp`$, $`u_su_1`$. Note that $`N^{}<1`$ due to (50), so that, in particular, $`p^{}>0`$. The expression (79) is manifestly positive for $`\mathrm{}>R>2k`$, therefore the $`\psi `$ mode also does not lead to an instability. Thus linear stability of all single-brane BH solutions under spherically symmetric perturbations has been established.
Our consideration did not apply to extremal BHs since in this case the behaviour of the background functions $`x^K(u)`$ is generically singular as $`u\mathrm{}`$ ($`R=1/u0`$):
$$x^K=\underset{s}{}\frac{Z_s^K}{Y^2}\mathrm{ln}\left(1+\frac{u}{u_s}\right)$$
(80)
and so there is no reason to require $`|\delta x^K|<\mathrm{}`$. In some cases it is regular (see Sec. 3.2 and examples in the Appendix). One can see, however, that again, as $`u\mathrm{}`$, Eqs. (71) and (72) for a single-brane extremal BH take the form $`\psi _{uu}=0`$; the linearly growing solution is discarded since it grows faster than $`x^K`$ in (80), so we are left with a constant and have to require $`|\psi <\mathrm{}|`$ for both regular and singular backgrounds. Then the same reasoning with $`\psi _{uu}/\psi >0`$ makes us conclude that such allowed solutions with $`\mathrm{\Omega }>0`$ do not exist and extremal BHs are stable as well. Indeed, an explicit form of $`U(u)`$ is
$$U(u)=\frac{2Z^2}{Y^2(u+u_1)^2}\left[1+N^{}\frac{u_1^2+N^{\prime \prime }u^2}{(u_1+N^{\prime \prime }u)^2}\right]$$
(81)
where $`N^{}<1`$ was defined in (79) and $`N^{\prime \prime }=1N^{}`$. The reasoning works since $`U>0`$ for all $`u>0`$ and $`U0`$ as $`u\mathrm{}`$.
We can now formulate the following result, to be used in the further consideration:
Proposition 2. If a decoupled linear perturbation mode $`\xi `$ of a static, spherically symmetric BH solution obeys the equation
$$\widehat{L}\xi =U(u)\xi $$
(82)
with $`U(u)0`$ (including the case $`U0`$), this mode is stable.
## 6. Some black holes with multiple branes
### 6.1. Two-brane black holes
We have seen that one-brane singular background solutions are catastrophically unstable; we would not like to treat more complex singular solutions since there is no reason to believe that interaction of modes can prevent the instability. We instead consider some multi-brane BH solutions for which the perturbation equations decouple and show that they are stable.
Suppose there is a BH background solution (37)–(41) with two branes, so that $`s`$ takes two values, $`s=1,2`$. The solution is characterized by two charges $`Q_s`$, two vectors $`\stackrel{}{Y}{}_{s}{}^{}𝕍`$ and their counterparts $`\overline{Z}{}_{s}{}^{}\overline{𝕍}`$, which we assume to be non-collinear (if they are collinear, the consideration simplifies and the result is the same as for a single brane).
The matrix $`P^K_L`$ in the perturbation equations (61) may be written in the form
$$P^K{}_{L}{}^{}=\underset{s}{}Q_s^2\mathrm{e}^{2y_s(u)}Z_s^KZ_{s,L}+\underset{ss^{}}{}Z_s^KZ_{s^{},L}f_{ss^{}}(u),$$
(83)
with
$$f_{ss^{}}(u)=\frac{1}{Y_s^2Y_s^{}^2}\left[\frac{(y_{s,u}+k)(y_{s^{},u}+k)}{d_0\beta _u}\right]_u.$$
(84)
Just as in the one-brane case, one easily separates the “transversal” degrees of freedom: for vectors $`\overline{Z}{}_{}{}^{}\overline{𝕍}`$ such that $`\overline{Z}{}_{}{}^{}\overline{Z}{}_{s}{}^{}=0`$ (they fill a $`(dim\overline{𝕍}2)`$-dimensional plane), the function $`\xi =\overline{Z}{}_{}{}^{}\delta \overline{x}`$ obeys the wave equation (82) with $`U0`$.
However, Eqs. (61) with the matrix (83) in general do not decouple. An exception is the special case when the two functions $`y_s`$ coincide,
$$y_1=y_2=y(u)=k\mathrm{sinh}(ku_1)/\mathrm{sinh}[k(u+u_1)]$$
(85)
although the vectors $`\stackrel{}{Y}_s`$ are different; in the BOS terminology (Sec. 3.3) the two vectors $`\stackrel{}{Y}_s`$ form a block and in our case this single block exhausts the whole system. If we suppose for simplicity that the norms of $`\stackrel{}{Y}_s`$ coincide, $`Y_1^2=Y_2^2=Y^2`$, then the charges coincide as well, $`Q_1^2=Q_2^2=Q^2`$, and one obtains for the two modes $`\xi _\pm =(oZ_1\pm \overline{Z}{}_{2}{}^{})\delta x`$:
$`\widehat{L}\xi _+`$ $`=`$ $`2Z^2(1+\mathrm{cos}\theta )\left[Q^2\mathrm{e}^{2y}+2F\right]\xi _+;`$ (86)
$`\widehat{L}\xi _{}`$ $`=`$ $`2Z^2Q^2\mathrm{e}^{2y}(1\mathrm{cos}\theta )\xi _{}`$ (87)
where $`Z^2=Y^2(d_01)/d_0`$ (see (49), $`\theta `$ is the angle between the vectors $`\overline{Z}_1`$ and $`\overline{Z}_2`$ in $`\overline{𝕍}`$ and $`F`$ is the function (84) with $`Y_1^2=Y_2^2`$ and $`y_1=y_2=y`$.
A direct substitution of $`y`$ and $`\beta _u`$ into (86) shows, as before, that the coefficient by $`\delta x_+`$ at the r.h.s. is nonnegative; for $`\delta x_{}`$ in (87) this is manifestly so. Hence the previous reasoning works and we conclude that such BHs (including extremal ones) with two branes are stable.
The case $`Y_1^2Y_2^2`$ is covered in the next section.
### 6.2. Single-block black holes
A natural question arises, whether or not the stability conclusion of the previous section extends to an arbitrary multi-brane BH described by a single function $`y(u)`$, in other words, to any single-block BH. Note that any set of linearly independent vectors $`\stackrel{}{Y}_s`$ may be treated as a BOS-block, hence a special static solution of this kind (and hence a BH solution) may always be obtained; the only restriction is $`a_\mu >0`$ for the charge factors obeying the consistency conditions (33).
Consider such a system: let there be a BOS BH solution with $`m`$ linearly independent vectors $`\stackrel{}{Y}{}_{s}{}^{}𝕍`$, $`s𝒮=𝒮_\omega `$, and the charge factors $`a_s`$ satisfy (33). The following relations are valid:
$$\stackrel{}{Y}{}_{\omega }{}^{}=\underset{s}{}a_s\stackrel{}{Y}{}_{s}{}^{};\stackrel{}{Y}{}_{s}{}^{}\stackrel{}{Y}{}_{\omega }{}^{}=Y_\omega ^2,s;\underset{s}{}a_s=1.$$
(88)
It is easy to see that, due to (49), similar relations hold for the corresponding vectors $`\overline{Z}{}_{s}{}^{}\overline{𝕍}`$:
$$\overline{Z}{}_{\omega }{}^{}=\underset{s}{}a_s\overline{Z}{}_{s}{}^{}s;\overline{Z}{}_{s}{}^{}\overline{Z}{}_{\omega }{}^{}=\overline{Z}{}_{\omega }{}^{2}.$$
(89)
For certainty we suppose that $`\overline{Z}_s`$ are linearly independent; if they are not, the consideration is slightly modified without changing the results.
The wave equations (61) take the form
$$\frac{1}{2}\widehat{L}\delta x^K=\widehat{q}\mathrm{e}^{2y(u)}\underset{s}{}a_sZ_s{}_{}{}^{K}Z_{sL}^{}\delta x^L+F(u)Z_\omega {}_{}{}^{K}Z_{\omega L}^{}\delta x^L,F(u)=\frac{1}{Y_\omega ^4}\left[\frac{(y_u+k)^2}{d_0\beta _u}\right]_u,$$
(90)
where $`\widehat{q}=\widehat{q}_\omega =_sQ_s^2`$ and $`y(u)`$ is given by (85). As before, the perturbations $`\stackrel{}{Z}{}_{}{}^{}\delta \overline{x}`$ (where $`\stackrel{}{Z}_{}`$ belong to the plane $`\overline{𝕍}_{}`$ orthogonal to all $`\stackrel{}{Z}_s`$, whose dimension is $`dim\overline{𝕍}m0`$) are decoupled and obey the equation (69), giving no unstable modes.
Multiplying (90) by $`\overline{Z}_\omega `$, one obtains a decoupled equation for $`\xi _\omega =\overline{Z}{}_{\omega }{}^{}\delta \overline{x}`$:
$$\widehat{L}\xi _\omega =2U_\omega (u)\xi _\omega ,U_\omega (u)=(\widehat{q}\mathrm{e}^{2y}+F)Z_\omega ^2$$
(91)
where $`Z_\omega ^2=\overline{Z}_\omega ^2`$. Since, as is directly verified, $`U_\omega (u)0`$, this mode is also stable.
The remaining $`(m1)`$ degrees of freedom may be described in terms of the vectors $`\overline{W}{}_{s}{}^{}=\overline{Z}{}_{s}{}^{}\overline{Z}_\omega `$ and the functions $`\xi _s=\overline{W}{}_{s}{}^{}\delta \overline{x}`$, such that
$$\overline{W}{}_{s}{}^{}\overline{Z}{}_{\omega }{}^{}=0;\underset{s}{}a_s\overline{W}{}_{s}{}^{}=0;\underset{s}{}a_s\xi _s=0.$$
(92)
Using (89) and (92), one obtains the following $`m`$ equations, coupled due to (92), for $`(m1)`$ independent variables:
$$\widehat{L}\xi _s=2\widehat{q}\mathrm{e}^{2y}\underset{s^{}=1}{\overset{m}{}}K_{ss^{}}\xi _s^{},K_{ss^{}}=a_s^{}\overline{W}{}_{s}{}^{}\overline{W}{}_{s^{}}{}^{}.$$
(93)
Excluding one of the unknowns, say, $`\xi _m`$, by virtue of (92), we arrive at a determined set of wave equations for $`\eta _s=\xi _s\xi _m`$, $`s=1,\mathrm{},m1`$:
$$\widehat{L}\eta _s=2\widehat{q}\mathrm{e}^{2y}\underset{s^{}=1}{\overset{m1}{}}F_{ss^{}}\eta _s^{},F_{ss^{}}=a_s^{}\overline{W}{}_{s^{}}{}^{}(\overline{W}{}_{s}{}^{}+\frac{1}{a_m}\underset{s^{\prime \prime }=1}{\overset{m1}{}}a_{s^{\prime \prime }}\overline{W}{}_{s^{\prime \prime }}{}^{}).$$
(94)
This is a good way of studying specific models. In the general case, however, the situation looks more transparent if we consider, instead, an auxiliary system with $`m`$ independent unknowns, described by Eqs. (93) where $`\overline{W}_m`$ is slightly shifted from its true value by some $`\mathrm{\Delta }\overline{W}`$, so that all $`\overline{W}_s`$ become linearly independent; the relation among $`\xi _s`$ in (92) is then cancelled as well. Our system is restored when $`\mathrm{\Delta }\overline{W}0`$.
For the auxiliary system the matrix $`(\overline{W}{}_{s}{}^{}\overline{W}{}_{s^{}}{}^{})`$ is symmetric and positive-definite; if all $`a_s`$ are equal, the same is true for the matrix of coefficients in (93), $`(K_{ss^{}})`$, hence there is a similarity transformation bringing it to a diagonal form with its positive eigenvalues along the diagonal. Such a transformation applied to Eqs. (93) decouples them into $`m`$ separate wave equations like (91), with some positive function replacing $`U(u)`$. In the limit $`\mathrm{\Delta }\overline{W}0`$, the worst thing that can happen is that some of the eigenvalues tend to zero, giving for some combinations of $`\xi _s`$ the equation $`\widehat{L}\xi =0`$ which, as we know, does not lead to an instability. One can assert “by continuity” that this picture is generic, at least for $`a_s`$ close enough to being equal, and stability is again concluded according to Proposition 2.
On the contrary, when the numbers $`a_m>0`$ are different, one cannot guarantee that the non-symmetric matrix $`K_{ss^{}}`$ is similar to a diagonal one . A failure in its diagonalization can be connected with the occurrence of a pair (or pairs) of complex roots $`\lambda _s`$ of the characteristic equation $`det|K_{ss^{}}\lambda \delta _{ss^{}}|=0`$. In this case there is at least one pair of coupled perturbations for which a special investigation is necessary. An inspection of the characteristic equation shows that the matrix $`K_{ss^{}}`$ cannot have negative eigenvalues, therefore a separate unstable mode cannot occur and the only possible instability can be connected with coupling between modes.
In particular, in an arbitrary OS BH solution there is a subfamily where all $`y_s(u)`$ coincide (i.e., the constants $`u_s`$ are the same for all $`s`$), so that the branes form a BOS block, and it turns out that all $`a_s`$ are also equal, as well as the squared charges $`Q_s^2`$. The above reasoning shows that such solutions are stable.
If $`rank(K_{ss^{}})<m1`$, that is, there are additional linear dependences among $`\overline{Z}_s`$, then some combinations of $`\xi _s`$ decouple leading to equations of the form $`\widehat{L}\xi =0`$, and for the remaining modes the above discussion can be repeated with slight modifications.
This is what can be said about the general case of single-block BOS BH solutions. If there is a block of only two branes ($`m=2`$), one can make a common stability conclusion generalizing the one made in Sec. 6.1. Indeed, for $`\xi _\omega =a_1\xi _1+a_2\xi _2`$ there is Eq. (91), whereas for $`\xi _{}=\xi _1\xi _2`$ one obtains
$$\widehat{L}\xi _{}=2\widehat{q}\mathrm{e}^{2y}a_1(1a_1)(\overline{Z}{}_{1}{}^{}\overline{Z}{}_{2}{}^{})^2\xi _{}.$$
(95)
In the special case $`Z_1^2=Z_2^2=Z^2`$ one recovers (86), (87).
For $`m3`$ one has to study specific models individually.
## 7. Concluding remarks
We have shown that all static single-brane BH solutions with the metric (2) are stable under linear spherically symmetric perturbations, whereas non-BH solutions possessing naked singularities of different types are unstable. Very probably other singular solutions, for which perturbation equations do not decouple, are unstable as well, since, as known from vibration theory, coupling between modes can hardly stabilize them. On the contrary, coupled modes can be unstable even when single ones are stable. It is therefore of interest to study the stability properties of more complex BH solutions; this work is in progress.
We have also shown that the BH stability conclusion can be extended to some BHs with multiple intersecting $`p`$-branes, namely, for the BOS case, characterized by a single function $`y(u)`$ (i.e., all $`y_s`$ coincide). It turns out that for such backgrounds the wave equations for perturbations also generically decouple and the absence of unstable modes can be proved. Though, such a general proof is available only in two cases: (i) two-brane BH solutions ($`m=2`$) and (ii) equal-charge subfamilies of arbitrary OS solutions. Nontrivial brane systems with $`m>2`$ should be studied individually to see whether or not the corresponding matrix $`(K_{ss^{}})`$ in (93) can be diagonalized over the field of real numbers. If yes, the solution is stable, otherwise a further study of coupled modes is necessary.
It should be stressed that a BOS-block solution exists for an arbitrary set of linearly independent vectors $`\stackrel{}{Y}_s`$. In particular, if any multi-brane static BH solution for a certain set of input parameters with independent vectors $`\stackrel{}{Y}_s`$ is known, e.g., any OS or BOS solution (see Sec. 3), then the additional requirement that all the functions $`y_s`$ coincide selects from it a special BOS-block solution, for which a stability study can be performed as described above. The only restriction is the requirement $`a_\mu (\omega )>0`$ for the charge factors obeying Eqs. (33).
Some technical points are worth mentioning. First, in gravitational stability studies it is sometimes rather hard to separate real physical perturbations from purely gauge degrees of freedom. We avoid this problem by obtaining the set of wave equations (61) where the number of equations is precisely the number of dynamical degrees of freedom, represented by scalars in the physical space-time $`𝕄_{\mathrm{phys}}`$. Due to the latter circumstance, one more complication is avoided: when dealing with vector and tensor perturbations of BHs, one has to take into account the apparent singularity of the metric on the horizon; to properly formulate the boundary conditions, it is then necessary to pass to Kruskal-like coordinates; to be admissible, and the perturbations are required to be finite on the future horizon . In our case the perturbations are scalars, so the finiteness requirement can be imposed in any coordinates. The choice of gauge only remains important for making the treatment more transparent.
To conclude, we would like to emphasize that our consideration did not depend on the number and dimensions of the factor spaces in the original space-time $`𝕄`$, on the number of scalar fields $`\phi ^a`$ and on the particular values of their coupling constants $`\lambda _{sa}`$.
## Appendix
Consider, for illustration, some solutions of $`D=11`$ supergravity, representing the low-energy limit of M-theory, as examples of systems to which our stability results apply.
The action (1) for this theory does not contain scalar fields ($`\phi ^a=\lambda _{sa}=0`$) and the only $`F`$-form is of rank 4, whose various nontrivial components $`F_s`$ (elementary $`F`$-forms according to Sec. 2, to be called simply $`F`$-forms) are associated with electric 2-branes \[for which $`d(I_s)=3`$\] and magnetic 5-branes \[such that $`d(I_s)=6`$\] (see and references therein). The orthogonality conditions (26) are satisfied if the following intersection rules hold:
$$33=1,36=2,66=4.$$
(A.1)
(the notations are evident); for all $`F`$-forms $`Y_s^2=2`$.
We will designate the branes by figures labelling their world volume coordinates (covered by $`I_s`$), beginning with “1” which corresponds to the time axis. Thus, e.g., (123) is an electic 2-brane whose world volume includes the time axis $`𝕄{}_{1}{}^{}=_t`$ and two extra dimensions. The number of dimensions where branes can be located is $`D1d_0=10d_0`$.
1. Single-brane BH solutions are described by (39), (41) where all sums and products in $`s`$ consist of a single term. These solutions are well known; the metric (2) can be presented as
$$ds_{11}^2=H^{d(I)/9}\left[\frac{12k/(\overline{d}r^{\overline{d}})}{H}dt^2+\left(\frac{dr^2}{12k/(\overline{d}r^{\overline{d}})}+r^2d\mathrm{\Omega }^2\right)+H^1ds_{\mathrm{on}}^2+ds_{\mathrm{off}}^2\right]$$
(A.2)
where $`H=H(r)=1+p/(\overline{d}r^{\overline{d}})`$, $`p=\sqrt{k^2+2Q^2}k`$, $`\overline{d}=d_01`$; $`ds_{\mathrm{on}}^2`$ and $`ds_{\mathrm{off}}^2`$ are the “on-brane” and “off-brane” extra-dimension line elements, respectively; the dimension $`d_0`$ of the sphere $`𝕄_0`$ varies from 2 to 7 for $`d(I)=3`$ (an electric brane) and from 2 to 4 for $`d(I)=6`$ (a magnetic brane). In particular, the cases of maximum $`d_0`$, when off-brane extra dimensions are absent, correspond in the extremal near-horizon limits to the famous structures $`AdS_4\times S^7`$ (electric) and $`AdS_7\times S^4`$ (magnetic). All these solutions are stable under linear spherically symmetric perturbations.
2. Some examples of orthogonal systems (OS), whose stability in the general case is yet to be studied, are:
$`(𝐢)`$ (123), (145), (167) — 3 electric branes; $`d_0=2`$ or 3.
$`(\mathrm{𝐢𝐢})`$ (123), (124567) — 1 electric and 1 magnetic branes; $`d_0=2`$ or 3.
$`(\mathrm{𝐢𝐢𝐢})`$ $`(123),(145),(124678),(135678);d_0=2.`$
The metrics are easily found from (39) with $`D2=9`$, $`Y_s^2=2`$ and the equalities marked $`\stackrel{\mathrm{OS}}{=}`$.
The systems (i) with $`d_0=3`$ and (iii) are remarkable in that their extremal limits have regular horizons and the near-horizon geometries are, respectively, $`AdS_2\times S^3\times T^6`$ and $`AdS_2\times S^2\times T^7`$ if the remaining extra dimensions are compactified on tori.
When the orthogonal systems form BOS-blocks (i.e., in the special case of equal charges and a unique function $`y(u)`$), the solutions are stable according to Sec. 6.2.
3. All two-brane BOS-block BHs are stable according to Sec. 6.1, for instance,
$$(𝐢)(123),(123456);\stackrel{}{Y}{}_{1}{}^{}\stackrel{}{Y}{}_{2}{}^{}=1.(\mathrm{𝐢𝐢})(123),(145678);\stackrel{}{Y}{}_{1}{}^{}\stackrel{}{Y}{}_{2}{}^{}=1.$$
The norms are equal ($`Y_s^2=2`$), and the angle $`\theta `$ is 60 in case (i) and 120 in case (ii).
4. Many seemingly possible three-brane blocks turn out to be forbidden due to a zero value of a certain $`a_s`$ (see Sec. 3.3). Consider, e.g.,
$$(123),(145),(123678);(\stackrel{}{Y}{}_{s}{}^{}\stackrel{}{Y}{}_{s^{}}{}^{})=\left(\begin{array}{ccc}2& 0& 1\\ 0& 2& 1\\ 1& 1& 2\end{array}\right).$$
(A.3)
One easily finds from (33) that $`(a_1,a_2,a_3)=(0,1/2,1/2)`$, so one of the charges should be zero, which means that such a system cannot exist.
5. The following is an example of a single-block BH whose stability can be established by an individual study as described in Sec. 6.2:
$$(123),(145),(123467);d_0=2\mathrm{or}3;(\stackrel{}{Y}{}_{s}{}^{}\stackrel{}{Y}{}_{s^{}}{}^{})=\left(\begin{array}{ccc}2& 0& 1\\ 0& 2& 0\\ 1& 0& 2\end{array}\right).$$
(A.4)
From (33) it follows
$$(a_1,a_2,a_3)=(2/7,3/7,2/7);\stackrel{}{Y}{}_{\omega }{}^{}=\underset{s=1}{\overset{3}{}}a_s\stackrel{}{Y}{}_{s}{}^{};\stackrel{}{Y}{}_{\omega }{}^{2}=6/7.$$
(A.5)
Suppose for certainty $`d_0=3`$. Then, in agreement with (49),
$$(\overline{Z}{}_{s}{}^{}\overline{Z}{}_{s^{}}{}^{})=\left(\begin{array}{ccc}4/3& 2/3& 1/3\\ 2/3& 4/3& 2/3\\ 1/3& 2/3& 4/3\end{array}\right),Z_\omega ^2=\frac{4}{21}.$$
(A.6)
The next stage is to separate the perturbation $`\overline{Z}{}_{\omega }{}^{}\delta \overline{x}`$, after which the remaining two degrees of freedom obey Eqs. (94) of the form
$$\widehat{L}\eta _s=2\widehat{q}\mathrm{e}^{2y}\underset{s^{}=1}{\overset{2}{}}F_{ss^{}}\eta _s^{},F_{ss^{}}=\left(\begin{array}{cc}2/7& 0\\ 2/7& 6/7\end{array}\right).$$
(A.7)
The characteristic equation $`det|F_{ss^{}}\lambda \delta _{ss^{}}|=0`$ has the form $`(2/7\lambda )(6/7\lambda )=0`$, and, according to Proposition 2, the positivity of its roots proves the stability of the background configuration.
### Acknowledgement
V.M. is grateful to Dept. Maths., University of the Aegean, Greece, and to Dept. Phys., Nara University, Japan, for their hospitality during his stay there in October and November-December 1999, respectively. K.B. acknowledges the hospitality of the colleagues of DFis-UFES, Vitória, ES, Brazil during his stay there in November-December 1999. The work was supported in part by the Russian Basic Research Foundation and by the Russian Ministry of Science and Technologies. |
warning/0002/hep-ph0002267.html | ar5iv | text | # Features of particle multiplicities and strangeness production in central heavy ion collisions between 1.7𝐴 and 158𝐴 GeV/𝑐.
## 1 Introduction
After scouring results from relativistic heavy ion collisions at many different energies over several years some common traits are starting to emerge. Indeed, statistical-thermal models have proved to be able to reproduce relative particle multiplicities in a satisfactory manner by using two or three relevant parameters: temperature, baryon chemical potential and a possible strange-quark suppression parameter, $`\gamma _\mathrm{s}`$. Such an analysis has been performed by many authors for heavy ion collisions data from CERN SPS, from Brookhaven AGS and also from GSI SIS. In this paper we present a simultaneous analysis of data from several different collisions, with emphasis on the similarity of the colliding system in order to study the behaviour of parameters as a function of centre-of-mass energy within one framework. Hence, we have focussed our attention on central Au–Au collisions at beam momenta of 1.7$`A`$ GeV/c (SIS) , 11.6$`A`$ GeV/c (AGS) and on central Pb–Pb collisions at 158$`A`$ GeV/c (SPS) . As far as the choice of data (and, consequently, colliding system) is concerned, our leading rule is the availability of full phase space integrated multiplicity measurements because a pure statistical-thermal model analysis of particle yields, without any consideration of dynamical effects, may apply only in this case . Such data, however, exist only in a few cases and whenever legitimate we have extrapolated spectra measured in a limited rapidity window to full phase space. The use of extrapolations is more correct than using data over limited intervals of rapidity, especially in the framework of a purely statistical-thermal analysis without a dynamical model. Moreover, the usual requirement of zero strangeness ($`S=0`$) demands fully integrated multiplicities because strangeness does not need to vanish in a limited region of phase space.
A point of considerable interest in heavy ion collisions is the enhanced production of strange quarks per u, d quark with respect to elementary collisions like $`\mathrm{e}^+\mathrm{e}^{}`$, $`\mathrm{pp}`$ , $`\mathrm{p}\overline{\mathrm{p}}`$. This could be related to properties of the system at the parton level prior to hadronisation . In order to further study strangeness production and enhancement at low energy, we also present a new analysis of Si–Au collisions at 14.6$`A`$ GeV (AGS) using only multiplicities obtained from fully integrated phase space distributions. This also allows to cross-check results of previous analyses performed using limited rapidity interval data. In particular, we have included the $`4\pi `$ pion multiplicity and results presented in . In order to assess the consistency of the results obtained, we have performed the statistical-thermal model analysis by using two completely independent numerical algorithms whose outcomes turned out to be in close agreement throughout.
Similar analyses have been recently made by other authors (see e.g. ); however, both the model and the used data set differ in several important details, such as the assumption of full or partial equilibrium for some quark flavours, the number of included resonances, the treatment of resonance widths, inclusion or not of excluded volume corrections, treatment of flow, corrections due to limited rapidity windows etc. Because of these differences it is difficult to trace the origin of discrepancies between different results. We hope that the present analysis, covering a wide range of beam energies using a consistent treatment, will make it easier to appreciate the energy dependence of the various parameters such as temperature and chemical potential.
## 2 Data set and model description
As emphasized in the introduction, in the present analysis we use the most recent available data, concentrating on fully integrated particle yields and discarding data that have been obtained in limited kinematic windows. The only exceptions to this rule are the $`\overline{\mathrm{\Lambda }}/\mathrm{\Lambda }`$ and $`\overline{\mathrm{p}}/\mathrm{K}^{}`$ ratios in Si–Au collisions which were not available in full phase space. It has been decided to keep them as they are the only available recent measurements involving antibaryons.
We have derived integrated multiplicities of $`\pi ^+`$, $`\mathrm{\Lambda }`$ and proton in Au–Au collisions at AGS by extrapolating published rapidity distributions with constrained mid-rapidity value ($`y_{\mathrm{NN}}`$=1.6). For proton and $`\mathrm{\Lambda }`$ we have fitted the data to Gaussian distributions, whilst for $`\pi ^+`$ we have used a symmetric flat distribution at midrapidity with Gaussian-shaped wings on each side; the point at which the Gaussian wing and the plateau connect is a free parameter of the fit. The fits yielded very good $`\chi ^2`$’s/dof: 0.27, 1.24 and 1.00 for $`\pi ^+`$, proton and $`\mathrm{\Lambda }`$ respectively. The integrated multiplicities have been taken as the area under the fitted distribution between the minimal $`y_{\mathrm{min}}`$ and maximal $`y_{\mathrm{max}}`$ values of rapidities for the reactions $`\mathrm{NN}\pi \mathrm{NN}`$, $`\mathrm{NN}\mathrm{\Lambda }\mathrm{K}`$ for pions and $`\mathrm{\Lambda }`$’s respectively; the difference between these areas and the total area has been taken as an additional systematic error. The area between $`y_{\mathrm{min}}`$ and $`y_{\mathrm{max}}`$ amounts to practically 100% of the total area for pions and about 95% for $`\mathrm{\Lambda }`$’s. Ref. quotes an additional experimental systematic error of 10% on $`\mathrm{\Lambda }`$ multiplicity that we have added in quadrature. Hence we obtain:
$`\pi ^+`$ $`=`$ $`133.7\pm 9.93`$
$`\mathrm{\Lambda }`$ $`=`$ $`20.34\pm 1.36\pm 1.23\pm 2.03`$ (1)
where the first error is the fit error, the second is the systematic error due to the variation of integration region and the third is the experimental systematic error. As to protons, the extracted rapidity interval corresponding to the reaction N N $``$ N N is only 79% of the total Gaussian area. The difference between the two areas is too large to be considered as an additional error; thus, in order to reduce the uncertainty, we have decided to take the ratio $`\mathrm{p}/\pi ^+`$ extracted in the above rapidity interval rather than the proton multiplicity itself. This yields:
$$\mathrm{p}/\pi ^+=1.234\pm 0.126$$
(2)
where the error includes both the fit error and an error stemming from a 10% systematic uncertainty quoted in ref. .
We have not included data on deuteron production because of the possible inclusion of fragments in the measured yields. This is particularly dangerous at low (SIS) energies where inclusion or not of deuterons modifies thermodynamic quantities like $`ϵ/n`$ .
The data analysis has been performed within an ideal hadron gas grand-canonical framework for Pb–Pb at SPS and Au–Au at AGS whereas for Au-Au at SIS and Si–Au at AGS we have required the exact conservation of strangeness instead of using a strangeness chemical potential (see the discussion later in the text); in both cases we have used a supplementary strange quark fugacity $`\gamma _\mathrm{s}`$. In the grand-canonical approach, the overall average multiplicities of hadrons and hadronic resonances are determined by an integral over a statistical distribution:
$$n_i=(2J_i+1)\frac{V}{(2\pi )^3}\mathrm{d}^3p\frac{1}{\gamma _\mathrm{s}^{s_i}\mathrm{exp}[(E_i𝝁\text{ }𝐪_i)/T]\pm 1}$$
(3)
where $`𝐪_i`$ is a three-dimensional vector with electric charge, baryon number and strangeness of hadron $`i`$ as components; $`𝝁`$ the vector of relevant chemical potentials; $`J_i`$ the spin of hadron $`i`$ and $`s_i`$ the number of valence strange quarks in it; the $`+`$ sign in the denominator is relevant for fermions, the $``$ for bosons. This formula holds in case of many different statistical-thermal systems (i.e. clusters or fireballs) having common temperature and $`\gamma _\mathrm{s}`$but different arbitrary momenta, provided that the probability of realizing a given distribution of quantum numbers among them follows a statistical rule . In this case $`V`$ must be understood as the sum of all cluster volumes measured in their own rest frame. Furthermore, since both volume and participant nucleons may fluctuate on an event by event basis, $`V`$ and $`𝝁`$ (and maybe $`T`$) in Eq. (3) should be considered as average quantities .
The overall abundance of a hadron of type $`i`$ to be compared with experimental data is determined by the sum of Eq. (3) and the contribution from decays of heavier hadrons and resonances:
$$n_i=n_i^{\mathrm{primary}}+\underset{j}{}\mathrm{Br}(ji)n_j$$
(4)
where the branching ratios Br$`(ji)`$ have been taken from the 1998 issue of the Particle Data Table .
It must be stressed that the unstable hadrons contributing to the sum in Eq. (4) may differ according to the particular experimental definition. This is a major point in the analysis procedure because quoted experimental multiplicities may or may not include contributions from weak decays of hyperons and K$`{}_{S}{}^{}{}_{}{}^{0}`$. We have included all weak decay products in our computed multiplicities except in Pb–Pb collisions on the basis of relevant statements in ref. and about antiproton production in refs. . It must be noted that switching this assumption in Au–Au at SIS and AGS does not affect significantly the resulting fit parameters whereas it does in Si–Au.
The overall multiplicities of hadrons depend on several unknown parameters (see Eq. (3)) which are determined by a fit to the data. The free parameters in the fit are $`T`$, $`V`$, $`\gamma _\mathrm{s}`$and $`\mu _B`$ (the baryon chemical potential) whereas $`\mu _S`$ and $`\mu _Q`$, i.e. the strangeness and electric chemical potentials, are determined by using the constraint of overall vanishing strangeness and forcing the ratio between net electric charge and net baryon number $`Q/B`$ to be equal to the ratio between participant protons and nucleons. The latter is assumed to be $`Z/A`$ of the colliding nucleus in Au–Au and Pb–Pb while it has been calculated to be 0.43 for central Si–Au collisions by means of a geometrical model.
As we have mentioned before, for SIS Au-Au and AGS Si–Au data we have required the exact conservation of strangeness instead of using a strangeness chemical potential. This gives rise to slightly more complex calculations which are necessary owing to either very small strange particle production (Au–Au) or a relatively small system size (Si–Au). The difference between the strangeness-canonical and pure grand-canonical calculationis of multiplicities of K and $`\mathrm{\Lambda }`$ for the final set of thermal parameters (see Table 1) turns out to be around 2-3% for K and $`\mathrm{\Lambda }`$ in Si–Au but it is as large as a factor 15 in Au–Au at 1.7$`A`$ GeV/c.
Owing to few available data points in SIS Au–Au collisions, we have not fitted the volume $`V`$ nor the $`\gamma _\mathrm{s}`$therein. The volume has been assumed to be $`4\pi r^3/3`$ where $`r=7`$ fm (approximately the radius of a Au nucleus) while $`\gamma _\mathrm{s}`$has been set to 1, the expected value for a completely equilibrated hadron gas. Since we have performed a strangeness-canonical calculation here, the yield ratios involving strange particles are not independent of the chosen volume value as in the grand-canonical framework. Thus, in this particular case, $`V`$ is meant to be the volume within which strangeness is conserved (i.e. vanishing) and not the global volume defining overall particle multiplicities as in Eq. (3). Also, in order to test the dependence of this assumption on our results, we have repeated the fit by varying $`V`$ by a factor 2 and 0.5 in turn.
The yields of resonances have been calculated by integrating Eq. (3) times a relativistic Breit-Wigner distribution over an interval $`[m\delta m_l,m+\delta m_u]`$, where $`\delta m_l=\mathrm{min}[mm_{\mathrm{threshold}},2\mathrm{\Gamma }]`$ and $`\delta m_u=2\mathrm{\Gamma }`$. The minimum mass $`m_{\mathrm{threshold}}`$ is required to open all decay modes<sup>1</sup><sup>1</sup>1In fact, in analysis A (see below) the integration interval has been taken symmetric $`[m\delta m_l,m+\delta m_l]`$. The relativistic Breit-Wigner distribution has been renormalised within the integration interval. The non-vanishing width of resonances plays a major role especially at low energies (e.g. SIS); for instance, the $`\mathrm{\Delta }(1232)`$ resonance creates pions more effectively than in the case of a vanishing width.
We have not used proper volume corrections in a Van der Waals type fashion which have been considered previously .
A major problem in Eq. (4) is where to stop the summation over hadronic states. Indeed, as mass increases, our knowledge of the hadronic spectrum becomes less accurate; starting from $`1.7`$ GeV many states are possibly missing, masses and widths are not well determined and so are the branching ratios. For this reason, it is unavoidable that a cut-off on hadronic states be introduced in Eq. (4). If the calculations are sensitive to the value of this cut-off, then the reliability of results is questionable. We have performed all our calculations with two cut-offs, one at 1.8 GeV (in the analysis algorithm A) and the other one at 2.4 GeV (in the analysis algorithm B). The contribution of missing heavy resonances is expected to be very important for temperatures $`200`$ MeV making thermal models inherently unreliable above this temperature.
## 3 Results
As mentioned in the introduction, we have performed two analyses (A and B) by using completely independent algorithms. In the analysis A all light-flavoured resonances up to 1.8 GeV have been included. The production of neutral hadrons with a fraction $`f`$ of $`\mathrm{s}\overline{\mathrm{s}}`$content has been suppressed by a factor $`(1f)+f\gamma _s^2`$. In the analysis B the mass cut-off has been pushed to 2.4 GeV and neutral hadrons with a fraction $`f`$ of $`\mathrm{s}\overline{\mathrm{s}}`$content have been suppressed by a factor $`\gamma _s^{2f}`$. Both algorithms use masses, widths and branching ratios of hadrons taken from the 1998 issue of Particle Data Table . However, it must be noted that differences between the two analyses exist in dealing with poorly known heavy resonance parameters, such as assumed central values of mass and width, where the Particle Data Table itself gives only a rough estimate. Moreover, the two analyses differ by the treatment of mass windows within which the relativistic Breit-Wigner distribution is integrated.
The results of the $`\chi ^2`$ fits are shown in Tables 1 and 2 for both analyses A and B. The agreement is indeed very good and confirms the reliability of the results obtained. The $`\chi ^2`$ minimisation in Au–Au collisions at AGS in analysis B did not converge to a reliable minimum; however, the $`\chi ^2`$ computed in analysis B by fixing the values of the parameters to the ones obtained in analysis A is approximately as large as in analysis A itself, thus confirming the good agreement between the two calculations.
We have investigated in detail the lack of convergence of analysis B in Au–Au collisions. The main reason of the fragility of the fit is the absence of measured antibaryon yields, which are very effective in fixing the baryon-chemical potential, in the main set of full phase space data. That shortage brings about a shallowness of $`\chi ^2`$ minima in four dimensions, and, consequently, a difficult convergence in both analyses. Notwithstanding, in analysis A the absolute minimum turned out to be deep enough, whereas in analysis B the convergence to a sufficiently nearby point was spoiled and the minimum drifted to $`T150`$ MeV with a nearly flat descent from the minimum found in A. This indicated a possible model dependence of the fit outcome. In order to check our result in analysis A and make it robust we have repeated the fits in Au–Au collisions at AGS by using an additional measurement of $`\overline{\mathrm{p}}/\mathrm{p}`$ ratio in the limited phase space region $`1.0<y<2.2`$ around midrapidity. The use of a ratio of particles measured by the same experiment under the same conditions reduces the involved systematic errors due to slightly different centrality definitions (with respect to the other data set) and other possible sources. However, the actual ratio in full phase space might be different owing to different shapes of $`\overline{\mathrm{p}}`$ and p rapidity distribution and this effect has been taken into account by conservatively assigning a 20% additional systematic error. The fit results are shown in Table 3; the two analyses are in very good agreement and, on top of that, the results for analysis A are in excellent agreement with those in Table 1 obtained without using $`\overline{\mathrm{p}}/\mathrm{p}`$ ratio, thus confirming the good quality of the calculation.
For each analysis an estimate of systematic errors on fit parameters have been obtained by repeating the fit
* assuming vanishing widths for all resonances
* varying the mass cut-off to 1.7 GeV in analysis A and to 1.8 GeV in analysis B
* for Au–Au at 1.7$`A`$ GeV/c, the volume $`V`$ has been varied to $`V/2`$ and to $`2V`$ (see discussion in Sect. 2)
The differences between the new fit parameters and the main parameters have been conservatively taken as uncorrelated systematic errors to be added in quadrature for each variation (see Table 1). The effect of errors on masses, widths and branching ratios of inserted hadrons has been studied in analysis A according to the procedure described in ref. and found to be negligible.
Finally, the results of the two analyses have been averaged according to a method suggested in ref. , well suited for strongly correlated measurements. Firstly, a simple no-correlation weighted average has been calculated as the central value of each parameter. Secondly, the error on it has been estimated by conservatively assuming that the results A and B are fully correlated, i.e. with a covariance matrix:
$$C=\left(\begin{array}{cc}\sigma _1& \sigma _1\sigma _2\\ \sigma _1\sigma _2& \sigma _2\end{array}\right)$$
(5)
yielding an error:
$$\sigma ^2=\frac{\frac{1}{\sigma _1^2}+\frac{1}{\sigma _2^2}+\frac{2}{\sigma _1\sigma _2}}{\left(\frac{1}{\sigma _1^2}+\frac{1}{\sigma _2^2}\right)^2}$$
(6)
The correlation between analyses A and B clearly arises from the use of the same set of hadronic data and theoretical model.
In Table 1 we also list the values of the Wroblewski factor $`\lambda _\mathrm{s}`$ measuring the number of newly created primary valence $`\mathrm{s}\overline{\mathrm{s}}`$pairs in comparison to the newly created non-strange primary valence quark pairs
$$\lambda _\mathrm{s}=\frac{2\mathrm{s}\overline{\mathrm{s}}}{\mathrm{u}\overline{\mathrm{u}}+\mathrm{d}\overline{\mathrm{d}}}$$
(7)
along with fit and systematic errors. The $`\mathrm{s}\overline{\mathrm{s}}`$and light quark pairs are computed on the basis of primary multiplicities of all hadron species, i.e. before particle decays take place. The behaviour of $`\lambda _\mathrm{s}`$as a function of collision type and centre-of-mass energy is shown in Fig. 1 including elementary and S–S, S–Ag collisions. Values for S–S, S–Ag and $`\mathrm{e}^+\mathrm{e}^{}`$, pp, $`\mathrm{p}\overline{\mathrm{p}}`$collisions have been taken from ref. .
## 4 Discussion and conclusions
From the results obtained, it emerges that a statistical-thermal description of multiplicities in a wide range of heavy ion collisions is indeed possible to a satisfactory degree of accuracy, for beam momenta ranging from 1.7$`A`$ GeV/c to 158$`A`$ GeV/c per nucleon. Furthermore, the fitted parameters show a remarkably smooth and consistent dependence as a function of centre-of-mass energy. The fit quality is generally good with the exception of Au–Au collisions at 1.7$`A`$ GeV/c where the large $`\chi ^2`$ is due to an underestimation of one ratio $`\eta /\pi ^0`$ (see Table 2).
The temperature varies considerably between the lowest and the highest beam energy, namely, between 50 MeV at SIS and 160 MeV at SPS. Similarly, the baryon chemical potential changes appreciably, decreasing from about 820 MeV at SIS to about 240 MeV at SPS. However, since the changes in temperature and chemical potential are opposite, the resulting energy per particle shows little variation and remains practically constant at about 1 GeV per particle; this is shown in Fig. 2.
The supplementary $`\gamma _\mathrm{s}`$factor, measuring the deviation from a completely equilibrated hadron gas, is around 0.7 – 0.8 at all energies where it has been considered a free fit parameter. At the presently found level of accuracy, a fully equilibrated hadron gas (i.e. $`\gamma _\mathrm{s}`$=1) cannot be ruled out in all examined collisions except in Pb–Pb, where $`\gamma _\mathrm{s}`$deviates from 1 by more than $`4\sigma `$. This result does not agree with a recent similar analysis of Pb–Pb data imposing a full strangeness equilibrium. The main reason of this discrepancy is to be found in the different data set used; whilst in ref. measurements in different limited rapidity intervals have been collected, we have used only particle yields extrapolated to full phase space. The temperature values that we have found essentially agree with previous analyses in Au–Au collisions and Si–Au collisions and estimates in Au–Au collision at 11.7 $`A`$ GeV/c .
The $`T`$ value in Pb–Pb is significantly affected by the multiplicity value of the heaviest particles measured, namely $`\varphi `$ and $`\mathrm{\Xi }`$, as they are almost entirely directly produced and provide a major lever arm on the slope of production vs. mass function. A recent 40% lowering of the $`\mathrm{\Xi }^{}`$ yield measured by NA49 with respect to a previous measurement results in a decrease of estimated temperature value from about 180 MeV to the actual 160 MeV. However, the removal of these two particles from the data set yields fitted parameter values which are in fair agreement with the main fit, as shown in Table 3. In particular, it is worth remarking that this exclusion does not bring significant change to $`\gamma _\mathrm{s}`$whose outcome is very sensitive to particles with multiple strange quark content and this confirms again the robustness of the main fit.
In order to further investigate strangeness production in Pb-Pb we have performed a consistency test between our fitted parameters, based on NA49 measurements, and the multiplicities of multi-strange hadrons measured by the experiment WA97 in central Pb–Pb collisions in a rapidity window $`\pm 0.5`$ around mid-rapidity . By fixing $`T`$, $`\gamma _\mathrm{s}`$and $`\mu _B`$ to the averaged values in Table 1 and adjusting the volume (i.e. an overall normalisation), we obtain a $`\chi ^2/\mathrm{dof}=28.9/6`$. Calculated $`\mathrm{\Lambda },\overline{\mathrm{\Lambda }}`$ multiplicities (see Table 4) do not include a residual feeding from $`\mathrm{\Xi }`$ decays in the experiment, estimated to be $`<5\%,<10\%`$ respectively. The high value of the $`\chi ^2`$ indicates that the statistical-thermal analysis is not able to reproduce data in a limited phase space region and in full phase space at the same time without resorting to a more detailed dynamical model. In particular, the parameters determined by the fit to NA49 data underestimate the yields of $`\mathrm{\Xi }`$ and $`\mathrm{\Omega }`$ baryons.
The parameter $`\gamma _\mathrm{s}`$as a function of centre-of-mass energy in heavy ion collision (including S–S and S–Ag ) is shown in Fig. 3. Again, the values for S–S, S–Ag and $`\mathrm{e}^+\mathrm{e}^{}`$, pp, $`\mathrm{p}\overline{\mathrm{p}}`$collisions have been taken from ref. . As can be seen from the Fig. 3 $`\gamma _\mathrm{s}`$is fairly constant, however, given the large error bars, it is quite difficult to exclude different behaviours. Also the behaviour of $`\lambda _\mathrm{s}`$factor (see Fig. 1) as a function of energy (provided that there is little dependence on system size at fixed $`\sqrt{s}`$, as the approximate equality of $`\lambda _\mathrm{s}`$in S–S and S–Ag confirms) is still unclear due to large experimental uncertainties. The line shape is either compatible with a monotonically increasing curve, saturating at $`\lambda _\mathrm{s}`$$`0.45`$, or with a curve having a maximum around Si–Au collisions, then decreasing and settling at an asymptotic $`0.45`$ value or maybe decreasing further to the characteristic value of elementary collisions.
Forthcoming lower energy Pb–Pb and high energy Au–Au data at RHIC should allow to clarify the behaviour of strangeness production in heavy ion collision. In order to easily compare our results with new measurements from RHIC experiments we also show in figure 4 the values of various particle-antiparticle ratios as a function of $`\overline{\mathrm{p}}/\mathrm{p}`$ ratio for different values of the temperature ($`T`$=160, 165 and 170 MeV) and a fixed value of the charge to baryon ratio of 0.401. RHIC results, however, will only be available for very limited kinematical region, while this kind of thermal model approach is largerly tied to full phase space ratios (see Introduction).
## Acknowledgements
We are very grateful to N. Carrer, U. Heinz, M. Morando, C. Ogilvie for useful suggestions and discussions about the data. We especially thank H. Oeschler for his help with the GSI SIS data and R. Stock for his help with NA49 data. |
warning/0002/gr-qc0002088.html | ar5iv | text | # 1. Introduction
## 1. Introduction
The SEE (Satellite Energy Exchange) concept of a space-based gravitational experiment was suggested in the early 90s and was aimed at precisely measuring the gravitational interaction parameters: the gravitational constant $`G`$, possible violations of the equivalence principle measured by the Eötvös parameter $`\eta `$, time variations of $`G`$, and hypothetical non-Newtonian gravitational forces (parametrized by the Yukawa strength $`\alpha `$ and range $`\lambda `$). Such tests are intended to fill gaps left by current methods of ground-based experimentation and observation of astronomical phenomena. The significance of new measurements is quite evident since nearly all modified theories of gravity and unified theories predict some violations of the Equivalence Principle (EP), either by deviations from the Newtonian law (inverse-square-law, ISL) or by composition-dependent (CD) gravity accelerations, due to the appearance of new possible massive particles (partners); time variations of $`G`$ are also generally predicted .
The idea of the SEE method is to study the relative motion of two bodies on board a drag-free Earth satellite using horseshoe-type trajectories, previously known in planetary satellite astronomy: if the lighter body (the ”Particle”) is moving along a lower orbit than the heavier one (the ”Shepherd”) and approaching from behind, then the Particle almost overtakes the Shepherd, but it gains energy due to their gravitational interaction, passes therefore to a higher orbit and begins to lag behind. The interaction phase can be studied within a drag-free capsule (a cylinder up to 20 m long, about 1 m in diameter) where the Particle can loiter as long as $`10^5`$ seconds. It was claimed that the SEE method exceeded in accuracy all other suggestions, at least with respect to $`G`$ and $`\alpha `$ for $`\lambda `$ of the order of metres. Some design features were considered, making it possible to reduce various sources of error to a negligible level. It was concluded, in particular, that the most favourable orbits are the sun-synchronous, continuous sunlight orbits situated at altitudes between 1390 and 3330 km.
Since the origination of the SEE concept, the development has focused on critical analyses and critical hardware requirements. All indications from this work are that the SEE concept is feasible and practicable . A “blue ribbon” Theory Advisory Group, formed two years ago to critique Project-SEE activities and goals, has concluded that they are sound.
This paper presents some new evaluations concerning the opportunities of the SEE concept and its yet-unresolved difficulties. In Sec. 2, for comparison, we briefly outline the current status of terrestrial and astronomical determination of the gravitational interaction parameters to be measured by the SEE method. In Sec. 3, on the basis of computer simulations of Particle trajectories, we estimate the requirements for the Shepherd quadrupole moment uncertainty. Sec. 4 shows the results of simulations of the measurement procedure itself, which enables us to estimate the possible measurement accuracy with respect to $`G`$ and $`\alpha `$ for $`\lambda `$ of the order of either metres or the Earth’s radius. Sec. 5 discusses a spurious effect of test body electric charging when the satellite orbit passes through the Van Allen radiation belts, rich in high-energy protons. Sec. 6 is a conclusion.
In what follows, the term “orbit” applies to satellite (or Shepherd) motion around the Earth, while the words “trajectory” or “path” apply to Particle motion with respect to the Shepherd inside the capsule.
## 2. State of the art: a brief survey
Since gravitational forces are so very small, precision-measurement techniques have been at the core of terrestrial gravity research for two centuries. However, evidence is increasingly accumulating which indicates that terrestrial methods have plateaued in accuracy and are unlikely to achieve significant accuracy gains in the future . For example, the uncertainty in the gravitational constant $`G`$ had been accepted as 128 ppm for nearly two decades, and the actual uncertainty in $`G`$ — as indicated by the scatter of results among recent experiments which claim high accuracy — is roughly the same (about 140 ppm). We discuss below the situation with respect to several key measurements.
### 2.1. Terrestrial determinations of $`G`$
The Luther & Towler determination of $`G`$ in 1982 , with the result $`(6.6726\pm 0.0005)\mathrm{\hspace{0.17em}10}^{11}`$ $`\mathrm{N}\mathrm{m}^2/\mathrm{kg}^2`$ and other, less precise experiments gave rise to the current official CODATA value of $`G`$, viz. 6.67259$`\mathrm{\hspace{0.17em}10}^{11}`$ $`\mathrm{N}\mathrm{m}^2/\mathrm{kg}^2`$ with an error of 128 ppm. Several still other experiments which also claimed high precision were ignored by CODATA because of inadequate documentation of systematic errors.
There is considerable evidence that the uncertainty in $`G`$ has plateaued at about 100 ppm. At a recent (November 23–24, 1998) conference in London, several new (1998) determinations of $`G`$ were reported. The obtained values for $`G`$ (in units of $`10^{11}`$ $`\mathrm{N}\mathrm{m}^2/\mathrm{kg}^2`$) and their estimated error $`\delta G/G`$ in ppm are as follows:
| Fitzgerald and Armstrong | 6.6742 | 90 |
| --- | --- | --- |
| (New Zealand) | 6.6746 | 134 |
| Nolting et al. (Zurich) | 6.6749 | 210 |
| Meyer et al. (Wuppethal) | 6.6735 | 240 |
| Karagioz et al. (Moscow) | 6.6729 | 75 |
| CODATA (1986) | 6.67259 | 128 |
Obviously, most of the stated errors are of the order 100 ppm. Moreover, the scatter (1-sigma) about the mean is about 140 ppm. Some of the investigators still hope for accuracy of 10 ppm. It remains to be seen whether they will be able to report error estimates of this size or, more importantly, whether their respective values for $`G`$ will actually agree within 10 ppm.
We note that this analysis is perhaps unduly optimistic since it excludes one extremely bad outlier: the very careful and well documented experiment by the Physikalisch-Technische Bundesanstalt in the late 1980s and early 1990s, which obtained a series of values for $`G`$ that were consistenly above the results of most experimenters by about 6000 ppm (0.6 % !), while claiming an error of about 100 ppm . No explanation for such a large discrepancy has been found.
It might seem that the problems of terrestrial apparatus must inexorably yield to new technologies — that the promise of ever increasing sensitivities would also lead to ever improving accuracy. However, this may not be true, since it is various systematic errors which limit the ultimate attainable accuracy in terrestrial experiments .
### 2.2. Terrestrial tests of the equivalence principle (EP) and search for Yukawa forces
The EP may be tested by searching for either violations of the inverse-square law (ISL) or composition-dependent (CD) effects in gravitational free fall.
In the watershed year of 1986, Fischbach startled the physics community by showing that Eötvös’s famous turn-of-the-century experiment is much less decisive as a null result than was generally believed . Prior to this time, experiments by Dicke and Braginsky had demonstrated the universality of free fall (UFF) to very high accuracy with respect to several metals falling in the gravitational field of the Sun (the Eötvös parameter $`\eta `$ was ultimately found to be smaller than $`10^{12}`$). The interpretation of these results at the time was that they validated UFF.
It was implicit that any violation would have infinite range, like gravity . During the 1970s and early 1980s there was also a flurry of activity concerning possible ISL violations, which eventually led to null results at the levels of precision then available (Fujii, , Long ).
Since 1986 it has become customary to parametrize possible apparent EP violations as if due to a Yukawa particle with a Compton wavelength $`\lambda `$. This approach unites both ISL and CD effects very naturally, while the parameter values in the Yukawa potential suggest which experimental conditions are required to detect the new interaction.
Following Fischbach’s conjecture, ISL and CD tests were undertaken by many investigators. Although a number of anomalies were initially reported, nearly all of these were eventually explained in terms of overlooked systematic errors or extreme sensitivity to models, while most investigators obtained null results. By far the tightest bounds are those obtained by Adelberger and his “Eot-Wash” group at the University of Washington . This group expects a further improvement of at least an order of magnitude . A positive result for a deviation from the Newtonian law (ISL) was obtained (and interpreted in terms of a Yukawa-type potential) in the range of 20 to 500 m by Achilli and colleagues ; this needs to be verified in other independent experiments.
For reviews of terrestrial searches for non-Newtonian gravity, see . The opportunities of the SEE concept in this respect are discussed in Refs. and in the present paper.
The UFF is still in the scope of the current experimental projects, and the SEE concept suggests here a progress of 3 to 4 orders of magnitude as compared with Ref. . Only one project, STEP (Satellite Test of the EP) promises a greater progress but meets some significant problems of its own , connected, in particular, with the radiation belts.
## 3. Simulations of Particle trajectories and the Shepherd quadrupole moment
In the previous studies of the SEE project it was assumed that the capsule was about 20 m long and the initial Shepherd-Particle separation $`x_0`$ along the capsule axis was as great as 18 m; some estimations were also made for $`5`$ m $`x_010`$ m. The Shepherd mass was taken to be $`M=500`$ kg and the Particle mass $`m=0.1`$ kg. The present study retains these values.
In what follows we describe some characteristic features of Particle trajectories with a goal to determine their sensitivity to the uncertainty of the Shepherd quadrupole moment $`J_2`$ for $`x_05`$ m. As in our previous studies, the capsule diameter is supposed to be 1 m.
The reason for considering the quadrupole moment uncertainty is technological by origin. Namely, it is hard to produce a spherically symmetric Shepherd to a required accuracy and, instead, it has been suggested to use a Cook-Marussi stack of cylinders with $`J_2>0`$, which may be manufactured more easily. A slow rotation of the Shepherd with $`J_2>0`$ will stabilize its position and orientation.
The value of $`J_2`$ can be provided with some uncertainty $`\delta J_2`$. To avoid the inclusion of $`\delta J_2`$ in the set of parameters to be determined in the experiment, it is useful to know which values of $`\delta J_2`$ will be negligible, since the growth of the number of parameters leads to serious problems in data processing.
### 3.1. Equations of motion and the initial data
Assuming that the relative motion of the test bodies inside the capsule occurs in the satellite orbital plane, the reduced Lagrangian of the Particle motion reads
$`L={\displaystyle \frac{M}{2}}(\dot{R}^2+R^2\dot{\phi }^2)`$
$`+{\displaystyle \frac{m}{2}}\left[\dot{r}^2+r^2(\dot{\phi }+\dot{\psi })^2\right]+G{\displaystyle \frac{M_{}m}{r}}`$
$`+G{\displaystyle \frac{Mm}{s}}\left\{1+J_2\left({\displaystyle \frac{r_s}{s}}\right)^2P_2(\mathrm{cos}\theta )\right\}\left(1+\alpha e^{s/\lambda }\right)`$ (1)
where $`(R,\phi )`$ are the Earth-centred polar coordinates of the Shepherd in the orbital plane; $`r=\sqrt{(R+y)^2+x^2}`$ and $`\psi `$ are the Earth-centred polar coordinates of the Particle; $`x`$ and $`y`$ are the Shepherd-centred Particle coordinates, where $`x`$ is the “horizontal” one, i.e., along the orbit and simultaneously along the capsule and $`y`$ is the “vertical” one, along the Earth-Shepherd radius vector; $`s=\sqrt{x^2+y^2}`$ is the Particle-Shepherd separation; $`M_{}`$, $`M`$ and $`m`$ are the Earth, Shepherd and Particle masses, respectively; $`J_2`$ is the quadrupole moment of the Shepherd, $`r_s`$ is its radius and $`P_2`$ is the Legendre polynomial
$$P_2(\mathrm{cos}\theta )=\frac{3\mathrm{cos}^2\theta 1}{2},$$
where $`\theta `$ is the angle between the line connecting the centres of the test bodies and the Shepherd equatorial plane. It is easy to see that if the Shepherd symmetry axis is in its orbital plane, then $`\theta =\theta _0=\mathrm{arctan}(y/x)+\phi `$. If the symmetry axis of Shepherd is orthogonal to its orbital plane, then $`\theta =0`$. In general, if $`\chi `$ is the angle between the Shepherd symmetry axis and its orbital plane, then $`\theta =\theta _0\mathrm{cos}\chi `$. Hence the influence of $`J_2`$ on the Particle motion is minimum if the Shepherd symmetry axis lies in its orbital plane and is maximum if they are mutually orthogonal.
For simplicity (and taking into account the corresponding estimate) we neglect the influence of the Particle on the Shepherd, so the Shepherd trajectory is considered to be given. Then, varying the above Lagrangian with respect to $`x`$ and $`y`$, taking into account that $`Mm`$ and $`Rs`$, we arrive at the following equations of Particle motion with respect to the Shepherd:
$`{\displaystyle \frac{d^2x}{dt^2}}=2\dot{y}\dot{\phi }+x\left\{\dot{\phi }^2{\displaystyle \frac{GM_{}}{r^3}}\right\}{\displaystyle \frac{2\dot{R}\dot{\phi }y}{R}}`$
$`{\displaystyle \frac{G\overline{M}}{s^3}}x\left\{1+J_2\left({\displaystyle \frac{r_s}{s}}\right)^2P_2(\mathrm{cos}\theta )\right\}`$
$`\alpha x{\displaystyle \frac{G\overline{M}}{s^2}}\{1+J_2\left({\displaystyle \frac{r_s}{s}}\right)^2P_2(\mathrm{cos}\theta )\}\left({\displaystyle \frac{1}{s}}+{\displaystyle \frac{1}{\lambda }}\right)\mathrm{e}^{s/\lambda }`$
$`{\displaystyle \frac{G\overline{M}r_0^2}{2s^5}}J_2(1+\alpha e^{s/\lambda })\times `$
$`\times \left[x(1+3\mathrm{cos}2\theta )+3y\mathrm{sin}2\theta \mathrm{cos}\chi \right]`$ (2)
$`{\displaystyle \frac{d^2y}{dt^2}}=2\dot{x}\dot{\phi }+(R+y)\left\{\dot{\phi }^2{\displaystyle \frac{GM_{}}{r^3}}\right\}+{\displaystyle \frac{2\dot{R}\dot{\phi }x}{R}}`$
$`{\displaystyle \frac{G\overline{M}}{s^3}}y\left\{1+J_2\left({\displaystyle \frac{r_s}{s}}\right)^2P_2(\mathrm{cos}\theta )\right\}`$
$`\alpha y{\displaystyle \frac{G\overline{M}}{s^2}}({\displaystyle \frac{1}{s}}+{\displaystyle \frac{1}{\lambda }})\left\{1+J_2\left({\displaystyle \frac{r_s}{s}}\right)^2P_2(\mathrm{cos}\theta )\right\}e^{s/\lambda }`$
$`+{\displaystyle \frac{G\overline{M}r_0^2}{s^5}}J_2(1+\alpha e^{s/\lambda })\times `$
$`\times \left[3x\mathrm{sin}2\theta \mathrm{cos}\chi +y(13\mathrm{cos}\theta )\right]`$ (3)
where $`\overline{M}=M+m`$.
Two kinds of initial conditions for Eqs. (3.1.) and (3.1.) were used during the simulations. First, we used the so-called “standard” initial conditions, taking the Particle velocity components $`\dot{x}(0)`$ and $`\dot{y}(0)`$ corresponding to its unperturbed (i.e., without the $`Mm`$ interaction) orbital motion distinguished from the Shepherd’s orbit only by its radius (for circular orbits) or major semiaxis (for elliptic orbits). Assuming that the Particle motion begins right at the moment when the Shepherd passes its perigee, these conditions have the form
$`x(0)=x_0,y(0)=y_0,`$
$`\dot{x}(0)={\displaystyle \frac{\omega e^{}y_0}{2(1e)^2}},\dot{y}(0)={\displaystyle \frac{\omega ex_0}{e^{}(1e)}}`$ (4)
where $`\omega ^2=GM_{}/R_0^3`$, $`R_0`$ is the Shepherd orbital radius (at the perigee), $`e`$ is the orbital eccentricity and $`e^{}=\sqrt{1e^2}`$.
For clearness, the relations (3.1.) are written in the linear approximation in the variables $`x`$ and $`y`$. Higher-order approximations were used in the simulation process as well.
The second kind of initial conditions correspond to small variations of initial velocities with respect to their “standard” values.
The set of equations (3.1.)–(3.1.) was solved numerically using the software developed previously to analyze the SEE project.
On the basis of numerical solution of Eqs. (3.1.) and (3.1.), we considered two types of Particle trajectories, corresponding to different choices of the initial data: (i) approximately U-shaped ones and (ii) cycloidal ones, containing loops (see more details on the trajectories in ), for orbital altitudes $`H_{\mathrm{orb}}=500`$, 1500 and 3000 km. The uncertainty $`\delta J_2`$ ranged in the interval $`10^3÷10^5`$; and initial Particle position changed in the range $`6`$ m $`x_018`$ m, $`25`$ cm $`y_05`$ cm.
For $`H_{\mathrm{orb}}=500`$ km, all U-shaped paths contained a sinusoidal component (with the orbital frequency), starting at $`x_08`$ m, while for $`H_{\mathrm{orb}}=1500`$ and $`3000`$ km it was present in paths starting at $`x_010`$ m. In other families of U-shaped trajectories a sinusoidal component was present only in the case $`|y_0|20`$ cm.
### 3.2. Restrictions on the Shepherd quadrupole moment uncertainty
Small Shepherd quadrupole moment uncertainties $`\delta J_2`$ create small displacements $`\delta \stackrel{}{r}`$ of a Particle trajectory with respect to unperturbed one, $`\stackrel{}{r}_0(t)`$:
$$\delta \stackrel{}{r}=\stackrel{}{r}_j(t)\stackrel{}{r}_0(t)$$
where $`\stackrel{}{r}_j`$ is the perturbed path. Instead of the full displacement $`\delta \stackrel{}{r}`$, a displacement $`\delta x`$ along the $`x`$ axis may be considered since, by numerical simulations, displacements along the $`y`$ axis are an order of magnitude smaller than $`\delta x`$.
Numerical simulations show that in the whole range of the above initial conditions the displacement $`\delta x`$ is (as it should naturally be) a linear function of $`\delta J_2`$; for $`\delta J_2=10^4`$. For the case when the Shepherd symmetry axis is located in the orbital plane, the maximum values of $`\delta x`$ for U-shaped trajectories are given in Table 1.
One can conclude that, if the distance measurement error is $`10^6`$ m, for most of the trajectories the uncertainty $`\delta J_2=10^4`$ is admissible.
The increase of $`\delta x`$ for small values of $`|y_0|`$ is explained by a large displacement of the turning point towards the Shepherd.
When the Shepherd symmetry axis is orthogonal to the orbital plane, these estimations change as shown in Table 2.
For cycloidal trajectories these values are approximately an order of magnitude smaller than those for the U-shaped ones.
It was also found that $`\delta x`$ decreases with increasing orbital altitude $`H_{\mathrm{orb}}`$. Table 3 shows, as an example, the displacements of U-shaped trajectories with $`x_0=18`$ m for $`H_{\mathrm{orb}}=500`$, 1500 and 3000 km and the Shepherd symmetry axis located in its orbital plane.
The above results show that the Shepherd quadrupole moment uncertainty $`\delta J_2`$ may be neglected in the SEE experiment with circular Shepherd orbits at $`H_{\mathrm{orb}}=1500`$ or $`3000`$ km and U-shaped Particle trajectories if $`\delta J_210^5`$ and the position measurement error $`\delta l`$ is $`10^6`$ cm, or $`\delta J_210^7`$ for $`\delta l=10^8`$ cm. For cycloidal Particle trajectories or elliptic Shepherd orbits these estimates become $`\delta J_210^4`$ and $`\delta J_210^6`$, respectively. For low orbits, $`H_{\mathrm{orb}}=500`$ km, the resitrictions on $`\delta J_2`$ become more stringent: $`\delta J_210^6`$ for $`\delta l=10^6`$ cm and $`\delta J_210^8`$ for $`\delta l=10^8`$ cm. However, as is evident from the above tables, these requirements may be relaxed by an order of magnitude if one discards some trajectories.
The influence of $`\delta J_2`$ on the accuracy of $`G`$ measurement may be now estimated as follows. Let some value of $`\delta J_2`$ produce the trajectory displacement $`|\delta \stackrel{}{r}|\delta l_j`$ while the variation $`\delta G_0`$ of $`G`$ with the same initial conditions gives the trajectory displacement $`|\delta \stackrel{}{r}|\delta l_G`$. Then, keeping in mind the linear dependence of trajectory displacements on $`\delta J_2`$ and $`\delta G`$, the accuracy of $`G`$ measurement under the uncertainty $`\delta J_2`$ may be estimated as
$$\frac{\delta G}{G}\frac{\delta l_j}{\delta l_G}\frac{\delta G_0}{G}.$$
Using this inequality and the results of trajectory simulations, we obtain the following estimates for U-shaped Particle trajectories in circular orbits with $`H_{\mathrm{orb}}=1500`$ km:
Table 4. Estimates of $`\delta G/G`$ in ppm for $`\delta J_2=10^4`$, when the symmetry axis of the Shepherd lies in ($`\chi =0`$) or is ortogonal to ($`\chi =\pi /2`$) its orbital plane. The second line shows $`x_0`$.
| $`y_0,`$ | $`\chi =0`$ | | $`\chi =\pi /2`$ | |
| --- | --- | --- | --- | --- |
| cm | 18 m | 6 m | 18 m | 6 m |
| -25 | $`0.88`$ | $`0.7`$ | $`3.57`$ | $`6.7`$ |
| -20 | $`0.27`$ | $`0.3`$ | $`1.05`$ | $`1.73`$ |
| -15 | $`0.06`$ | $`0.1`$ | $`0.24`$ | $`0.44`$ |
One can conclude that the uncertainties $`\delta J_2\underset{}{<}10^5`$ do not create substantial $`G`$ errors for most of the trajectories.
## 4. Simulations of experimental procedures
This section describes the results of computer simulations of the whole measurement procedures aimed at obtaining the sought-after gravitational interaction parameters. These simulations assumed the Shepherd mass $`M=500`$ kg, a circular orbit with $`H_{\mathrm{orb}}=1500`$ km under a spherical gravitational potential of the Earth, and a Particle mass of $`100`$ g. Where relevant, it is assumed that both the Shepherd and the Particle are made of tungsten. Their identical compositions are assumed for simplicity since this work is performed only for estimation purposes.
### 4.1. Equations of motion with Yukawa terms
We will begin with a presentation of the Particle equations of motion in the relevant approximation, including the contributions from hypothetical Yukawa forces, taking into account the finite size of the Yukawa field sources.
Let the interaction potential for two elementary masses $`m_1`$ and $`m_2`$ be described by the potential
$$dV^{\mathrm{Yu}}=\frac{Gdm_1dm_2}{r}\alpha \mathrm{e}^{r/\lambda }$$
(5)
where $`r`$ is the masses’ separation, $`\alpha `$ and $`\lambda `$ are the strength parameter and the range of the Yukawa forces. Then for two massive bodies with the radii $`R_1`$ and $`R_2`$ after integration over their volumes we obtain
$$V^{\mathrm{Yu}}=\frac{Gm_1m_2\beta _1\beta _2}{r}\alpha \mathrm{e}^{r/\lambda }$$
(6)
where
$$\beta _i=3\left(\frac{\lambda }{R_i}\right)^3\left[\frac{R_i}{\lambda }\mathrm{cosh}\frac{R_i}{\lambda }\mathrm{sinh}\frac{R_i}{\lambda }\right].$$
(7)
When $`R_i/\lambda 1`$, we have $`\beta _i1`$. This may be the case when we consider the interaction between the Shepherd and the Particle at a distance of the order of a few metres. The radii of the Shepherd and the Particle are small: $`R_118`$ cm for the Shepherd and $`R_21.1`$ cm for the Particle. If the range $`\lambda `$ is of the order of the Earth radius, $`\lambda R_{}`$, we have $`\beta _{}=1.10`$ and $`\beta _{1,2}=1`$ where the indices 1 and 2 label the Shepherd and the Particle, respectively.
The equations of motion are obtained under the following assumptions. There are two Yukawa interactions with the parameters $`\lambda _0`$ and $`\alpha _0`$ referring to the Earth-Shepherd and Earth-Particle interactions which are the same (due to the assumed identical composition for the Shepherd and the Particle), while $`\lambda `$ and $`\alpha `$ determine the Shepherd-Particle interaction. The equations of motion in the frame of reference connected with the Shepherd, with the same notations $`x`$, $`y`$, $`s`$ as previously, are
$`\ddot{x}+2\omega ^2\dot{y}+G(m_1+m_2){\displaystyle \frac{x}{s^3}}3\omega ^2{\displaystyle \frac{xy}{s}}`$
$`+G(m_1+m_2){\displaystyle \frac{x}{s^3}}\alpha \left(1+{\displaystyle \frac{s}{\lambda }}\right)\mathrm{e}^{s/\lambda }=0;`$
$`\ddot{y}2\omega \dot{x}3\omega ^2y+G(m_1+m_2){\displaystyle \frac{y}{s^3}}+{\displaystyle \frac{3\omega ^2}{r_{01}}}\left(y^2{\displaystyle \frac{x^2}{2}}\right)`$
$`+G(m_1+m_2){\displaystyle \frac{y}{s^3}}\alpha \left(1+{\displaystyle \frac{s}{\lambda }}\right)\mathrm{e}^{s/\lambda }`$
$`\omega ^2\beta _0\alpha _0\mathrm{e}^{r_{01}/\lambda _0}y=0`$ (8)
where $`\omega `$ is the orbital frequency:
$$\omega ^2=\frac{GM_{}}{r_{01}^3}\left[1+\beta _0\alpha _0\left(1+\frac{r_{01}}{\lambda _0}\right)\mathrm{e}^{r_{01}/\lambda _0}\right].$$
(9)
We have neglected the terms quadratic in $`s/r_{01}`$ times $`\alpha `$ or $`\alpha _0`$ due to their manifestly small contributions.
If we set $`\alpha _0=0`$ in Eqs. (8), we obtain the equations used to describe only the Shepherd-Particle Yukawa interaction. One can notice that Yukawa terms are roughly proportional to the gradients of the corresponding Newtonian accelerations, namely, $`Gm_1/s^3`$ for the Shepherd-Particle interaction and $`GM_{}/r_{01}^3\omega ^2`$ for (say) the Earth-Shepherd interaction. In our case these quantities are estimated as
$`{\displaystyle \frac{Gm_1}{s^3}}2.7\mathrm{\hspace{0.17em}10}^{10}\mathrm{s}^2\mathrm{for}s=5\mathrm{m},`$
$`\omega ^28.16\mathrm{\hspace{0.17em}10}^7\mathrm{s}^2.`$ (10)
Thus, given the same strength parameter, the Earth’s Yukawa force is three orders of magnitude greater than that between the Shepherd and the Particle, therefore one might expect some significant progress in an ISL test for $`\lambda `$ of the order of the Earth’s radius.
The effect of the Earth’s Yukawa force is proportional to the displacements of the statellite along the direction of the Earth’s radius. Therefore the sensitivity of the SEE method will increase if one uses orbits with eccentricities of the order of 0.01, following Nordtvedt’s suggestion . (Larger eccentricities would too much disturb the qualitative picture of a SEE encounter.) Tentative estimates show that in this way one can achieve sensitivities to $`\alpha 10^{10}`$, and more thourough studies are in progress.
Eqs. (8) were used to simulate the measurement procedures.
### 4.2. Simulations of an experiment for measuring $`G`$
The constant $`G`$ is determined from the best fitting condition between the “theoretical” ($`\stackrel{}{r}{}_{}{}^{\mathrm{th}}(t_i)=\stackrel{}{r}_i^{\mathrm{th}}`$) and “empirical” ($`\stackrel{}{r}_i`$) Particle trajectories near the Shepherd. The fitting quality is evaluated by minimizing a functional characterizing a “distance” between the trajectories. We have considered the following functionals for such “distances”:
$`S={\displaystyle \underset{i=1}{\overset{N}{}}}[(x_ix{}_{i}{}^{\mathrm{th}})^2+(y_iy{}_{i}{}^{\mathrm{th}})^2],`$ (11)
$`S_x={\displaystyle \underset{i=1}{\overset{N}{}}}(x_ix{}_{i}{}^{\mathrm{th}})^2,S_y={\displaystyle \underset{i=1}{\overset{N}{}}}(y_iy{}_{i}{}^{\mathrm{th}})^2,`$ (12)
$`S^{}={\displaystyle \underset{i=1}{\overset{N}{}}}[|x_ix{}_{i}{}^{\mathrm{th}}|+|y_iy{}_{i}{}^{\mathrm{th}}|],`$ (13)
$`S_x^{}={\displaystyle \underset{i=1}{\overset{N}{}}}|x_ix{}_{i}{}^{\mathrm{th}}|,S_y^{}={\displaystyle \underset{i=1}{\overset{N}{}}}|y_iy{}_{i}{}^{\mathrm{th}}|.`$ (14)
The theoretical trajectory depends on the gravitational constant $`G`$, on the initial coordinates $`x_0,y_0`$ and on the initial velocities $`v_{x0},v_{y0}`$. To estimate $`G`$, one chooses the value for which a “distance” functional in the space of the five variables ($`G,x_0,y_0,v_{x0},v_{y0}`$) reaches its minimum.
We carried out a computer simulation of the SEE experiment and estimated $`\delta G`$ for a given coordinate measurement error ($`\sigma =1\mathrm{\hspace{0.17em}10}^6`$ m). As “empirical” trajectories, we took computed trajectories, with specified values of the above five variables, where a Gaussian noise was introduced from a random number generator. Independent “empirical trajectories were created by non-intersecting random number sequences. The functional was minimized using the gradient descent method and the consecutive descent method. The starting value of the “vertical” (along the Earth’s radius) coordinate, $`y_0`$, was taken to be 0.25 m, while the horizontal one, $`x_0`$, varied between 2 and 18 m. Fig 1 shows the dependence of the errors $`\delta G/G=R_{\mathrm{grad}}`$, obtained by the gradient descent method and $`\delta G/G=R_\mathrm{s}`$, obtained by the consecutive descent method. All the errors are estimated by confidence intervals corresponding to a confidence of 0.95. The mean values of these errors are as follows:
$`R_{\mathrm{grad}}=4.69\mathrm{\hspace{0.17em}10}^8,R_\mathrm{s}`$ $`=`$ $`5.24\mathrm{\hspace{0.17em}10}^8.`$
Thus the errors estimated by the gradient and consecutive descent methods are close to each other and are about an order of magnitude smaller than the error from one-trajectory data. It has been discovered that the simulation results strongly depend on the random number generator, so that ordinary generators are not perfect.
The use of truncated functionals like (2) has shown that a functional incorporating the more informative “horizontal” coordinate $`x`$ leads to estimates close to those obtained from the total functional, whereas the use of $`y`$ alone substantially decreases the sensitivity. Therefore in practice, to determine $`G`$, it is sufficient to measure only one of the two coordinates, viz. $`x`$.
Since the “empirical” trajectory is built on the basis of a computed one, with a known value of the gravitational constant $`G_0`$, it appears possible to estimate a possible systematic error inherent in the data processing method. The latter has turned out to be in most cases much smaller than the random error. This result shows the correctness of the methods used.
As is evident from the results, the best accuracy is achieved at values of $`x_0`$ ($``$ the capsule size) about 4–5 metres.
### 4.3. Sensitivity to Yukawa forces with $`\lambda 1`$ m
In an experiment for finding a Yukawa interaction between the Shepherd and the Particle with the potential (6) with $`\beta _{1,2}=1`$, one computes two theoretical trajectories: one ignoring the Yukawa forces ($`x^0(t_i),y^0(t_i)`$) and another taking them into account $`(x^\alpha (t_i),y^\alpha (t_i))`$. These two computed curves are compared with the empirical trajectory using the functional $`S_k`$ ($`k=0,\alpha `$) according to (11) which may be considered as a dispersion characterizing a scatter of the “empirical” coordinates with respect to the fitting trajectory. This is true when the theoretical model is adequate to the real situation. In the case $`k=\alpha `$ the functional $`S_k=s_\alpha `$ has a $`\chi ^2`$ distribution with $`n_2=2N1`$ degrees of freedom. With $`k=0`$ the parameter $`\alpha `$ is absent, therefore $`S_0`$ is distributed according to the $`\chi ^2`$ law with $`N_1=2N`$ degrees of freedom. Then their ratio $`S_0/S_\alpha =F_{n_2,n_1}`$ will be distributed according to the Fischer law with $`n_2`$ and $`n_1`$ degrees of freedom. If an experiment shows that, on a given significance level $`q`$, the relation
$$S_0/S_\alpha F_{n_1,n_2,q}$$
(15)
is valid, one should conclude that a Yukawa force has been detected. An equality sign shows a minimum detectable force on the given significance level $`q`$. We have assumed $`q=0.95`$. The results of a sensitivity computation for different values of the space parameter $`\lambda `$ are presented in Fig. 2.A maximum sensitivity of $`\alpha =2.1\mathrm{\hspace{0.17em}10}^7`$ has been observed for $`\lambda =1.25`$ m. This value is 3 to 4 orders of magnitude better than the sensitivity of terrestrial experiments in the same range.
These results are based on the measurement method which was proposed in the original SEE paper ; as already mentioned, a method involving an eccentric orbit , is much more sensitive and, by our tentative estimates, can give an error $`\delta \alpha \underset{}{<}10^{10}`$.
### 4.4. Sensitivity to Yukawa forces with $`\lambda R_{}`$
To estimate the parameter $`\alpha _0`$ in Eqs. (8), computer simulations were carried out using the method as described above for $`\alpha `$, based on the Fischer criterion for the significance level 0.95. The range parameter $`\lambda _0`$ varied from $`(1/32)R_{}`$ to $`32R_{}`$. Two trajectories with the initial Shepherd-Particle separations $`x_0`$ of 2 and 5 m were calculated. In both cases the impact parameter $`y_0`$ was chosen to be 0.25 m. We used Eqs. (8) with $`\alpha =0`$, i.e., excluding the non-Newtonian interaction between the Shepherd and the Particle. As is evident from Eqs. (8), the Particle trajectory depends on the ratio $`r_{01}/\lambda _0`$ in the product $`(r_{01}/\lambda _0)\mathrm{e}^{r_{01}/\lambda _0}`$. This quantity reaches its maximum at $`\lambda _0=r_{01}/2`$. Our calculations have confirmed that a maximum sensitivity of the SEE method ($`3.4\mathrm{\hspace{0.17em}10}^8`$ for $`x_0=5`$ m) is indeed observed at this value of $`\lambda _0`$. This is about an order of magnitude better than the estimates obtained by other methods. Hopefully this estimate may be further improved by about an order of magnitude by optimisation of the orbital parameters. However, there is a factor which can, to a certain extent, spoil these results, namely, the uncertainty in the parameter $`\omega `$ which, in this calculation, was assumed to be precisely known.
The simulation results are shown in Fig. 3 for two trajectories with initial Shepherd-Particle separations of 2 and 5 metres.
## 5. A possible effect of the Earth’s radiation belt
Charged particles, penetrating into the SEE capsule from space and captured by the test bodies, create electrostatic forces that could substantially distort the experimental results. Among the sources of such particles one should mention (i) cosmic-ray showers, (ii) solar flares and (iii) the Earth’s radiation belts (Van Allen belts). The effect of cosmic-ray showers was estimated in Ref. and shown to be negligible. Solar flares are more or less rare events and, although they create very significant charged particle fluxes, sometimes even exceeding those in the most dense regions of the radiation belts, one can assume that the SEE measurements (except those of $`\dot{G}`$) are stopped for the period of an intense flare. On the contrary, the effect of the Van Allen belts is permanent as long as the satellite orbit passes, at least partially, inside them.
We will show here that the charging is unacceptably high at otherwise favourable satellite orbits, so that some kind of charge removal technique is necessary, but this problem may be solved rather easily by presently available technology.
The range of the most favourable SEE orbital altitudes, roughly 1400 to 3300 km , coincides with the inner region of the so-called inner radiation belt , situated presumably near the plane of the magnetic equator. This region is characterized by a considerable flux of high-energy protons and electrons. For a SEE satellite at altitudes near 1500 km the duration of the charging periods is about 12 minutes. Maximum charging rates occur in the central Atlantic. It should be noted that the South Atlantic Anomaly (SAA) — a region of intense Van Allen activity which results from the low altitude of the Earth’s magnetic field lines over the South Atlantic Ocean — cannot cause additional problems for the SEE experiments. The reason is that the SAA mostly contains low-energy protons which cannot penetrate into the SEE capsule.
Electrons are known to be stopped by even a thin metallic shell, so only protons are able to induce charges on the test bodies. Proton-induced charges on the test bodies can create considerable forces. The inner radiation belt contains protons with energies of 20 to 800 MeV, and their maximum fluxes at an altitude of 3000 km over the equator are as great as about $`3\mathrm{\hspace{0.17em}10}^6\mathrm{cm}^2\mathrm{s}^1`$ for energies $`E\underset{}{>}10^6`$ eV and about $`2\mathrm{\hspace{0.17em}10}^4\mathrm{cm}^2\mathrm{s}^1`$ for $`E\underset{}{>}10^7`$ eV. At 1500 km altitude these numbers are a few times smaller; the fluxes gradually decrease with growing latitude $`\phi `$ and actually vanish at $`\phi 40^{}`$.
It is thus necessary to have some estimates taking into account that (i) the capsule walls have a considerable thickness and stop the low-energy part of the proton flux and (ii) among the protons that penetrate the capsule and hit the Particle, the most energetic ones, whose path in the Particle material is longer than the Particle diameter, fly it through and hit the capsule wall again. As for the Shepherd, its size is large enough to stop the overwhelming majority of protons which hit it.
In what follows, we will assume a Shepherd radius of 20 cm and a Particle radius of 2 cm and estimate the captured charges for some satellite orbits in a capsule whose walls of aluminium are 2, 4, 6 and 8 cm thick. The SEE satellite must actually involve several coaxial cylinders for thermal-radiation control, and the combined thickness of their walls must amount to several cm. We will assume, in addition, that the Particle also consists of aluminium and stops all protons whose path is shorter than 4 cm (thus a little overestimating the charge since most of protons will cover a smaller path through the Particle material). A 100 g Particle of aluminium will have a radius of $`2.07`$ cm.
It is advisable to determine first which charges (and fluxes that create them) might be regarded negligible.
### 5.1. Admissible charges
Let us estimate the Coulomb interaction both between the Shepherd and the Particle and between each test body and its image in the capsule walls. To estimate the spurious effects on the Particle trajectory, it is reasonable to calculate its possible displacements due to the Coulomb forces from the growing captured charges. We assume that the test bodies are discharged by grounding to the capsule before launching the motion.
Criterion. We will call the induced charges, or the fields they create, admissible if they cause a displacement of the Particle with respect to the Shepherd smaller than a prescribed coordinate measurement error $`\delta l`$ (we take here $`\delta l=10^6`$ m) for a prescribed measurement time (we take $`t10^4`$ s).
A charge on the Shepherd can be estimated as
$$q_MeS_MJ(t,x)𝑑t=eS_MF(t,x)$$
(16)
where $`e`$ is the elementary charge, $`x`$ is the capsule wall thickness in cm; $`J(t,x)`$ is the integral proton flux in $`\mathrm{cm}^2\mathrm{s}^1`$ after passing through the wall, that is, the flux of protons with energies $`E_p>E_p(x)`$ where $`E_p(x)`$ is such an energy that the proton path in aluminium equals $`x`$ cm; $`S_M1256`$ cm<sup>2</sup> is the Shepherd’s cross-section; $`F(t,x)`$ is the fluence, i.e., the total number of protons of relevant energies that crosses a square centimeter of area for a certain period $`t`$.
In a similar way, the charge captured by the Particle may be found as
$`q_m`$ $`\underset{}{<}`$ $`eS_m{\displaystyle [J(t,x)J(t,x+4)]𝑑t}`$ (17)
$`=`$ $`eS_m[F(t,x)F(t,x+4)]`$
where $`S_m12.56`$ cm<sup>2</sup> is the Particle cross-section. The subtraction in the square brackets takes into account the protons which fly through the Particle without stopping there. The sign $`\underset{}{<}`$ is used since the effective Particle cross-section is smaller than its equatorial section.
The Coulomb acceleration $`a_Q(t)=q_Mq_m/(r^2m)`$ (in the Gaussian system of units) depends on the Shepherd-Particle separation $`r`$ and on the form of the function $`J(t)`$, which in turn depends on the satellite orbital motion.
The charge-induced Particle displacement is approximately
$$\mathrm{\Delta }l=𝑑t\left[𝑑ta_Q(t)\right]$$
(18)
since the acceleration is almost unidirectional. If, for estimation purposes, we suppose that the flux is time-independent, $`J=J_0=`$const, and take into account that in Eq. (17) the difference $`J(t,x)J(t,x+4)\frac{2}{5}J(t,x)`$ (or even smaller; see particular values in the next section), then the resulting displacement is about
$$\mathrm{\Delta }l\frac{1}{30}\frac{e^2S_MS_mJ_0^2t^4}{r^2m}.$$
(19)
The strong time dependence is explained by the rapid growth of the Coulomb force with capturing the charge. Numerically, with the above values of $`S_M`$ and $`S_m`$, taking $`m=100`$ g and $`r=1`$ m (the latter leads to an overestimated force since the Particle spends most of time at greater distances), one gets:
$$J_0^2t^4\underset{}{<}0.83\mathrm{\hspace{0.17em}10}^{18}\mathrm{s}^2\mathrm{cm}^4.$$
(20)
For $`t=10^4`$ s an admissible flux is only within 9 $`\mathrm{cm}^2\mathrm{s}^1`$ .
Another undesired effect is that the Particle, being charged by the belt protons, will interact with the capsule walls. This is well approximated as an interaction with the Particle’s mirror image in the wall, while the latter may be roughly imagined as a conducting plane. Then, assuming that the Particle is at average at about 25 cm from the capsule wall and using the same kind of reasoning as above, we obtain instead of (20)
$$J_0^2t^4\underset{}{<}2.07\mathrm{\hspace{0.17em}10}^{19}\mathrm{s}^2\mathrm{cm}^4$$
(21)
and an admissible proton flux within 45 $`\mathrm{cm}^2\mathrm{s}^1`$ for $`t=10^4`$ s.
Some more estimates are of interest: if a charge can be kept smaller than a certain value, then how great may it be to create only negligible displacements? Suppose that there are constant charges on both the Shepherd ($`q=q_M`$) and the Particle ($`q=q_m,m=100`$ g), then they are admissible according to the above criterion as long as
$`q_Mq_m`$ $`<2\mathrm{\hspace{0.17em}10}^6\mathrm{CGSE}_q^2=\frac{2}{9}\mathrm{\hspace{0.17em}10}^{24}\mathrm{C}^2,`$ (22)
$`q_m^2`$ $`<\frac{1}{2}\mathrm{\hspace{0.17em}10}^6\mathrm{CGSE}_q^2.`$ (23)
These inequalities follow, respectively, from considering the Shepherd-Particle interaction and the interaction between the Particle (located at 25 cm from the wall) and its image. Thus the maximum admissible Particle charge is about $`7\mathrm{\hspace{0.17em}10}^4\mathrm{CGSE}_q1.5\mathrm{\hspace{0.17em}10}^6e`$; assuming this value, it follows from (22) that the maximum Shepherd charge is about $`3\mathrm{\hspace{0.17em}10}^3\mathrm{CGSE}_q5.5\mathrm{\hspace{0.17em}10}^6e`$. With these charge values the electric potentials on the test body surfaces are
$`U_M1.5\mathrm{\hspace{0.17em}10}^4\mathrm{CGSE}_q/\mathrm{cm}=45\mathrm{mV};`$
$`U_m3.5\mathrm{\hspace{0.17em}10}^4\mathrm{CGSE}_q/\mathrm{cm}=105\mathrm{mV}.`$ (24)
If by any means the requirements (22), (23) are satisfied (e.g., the potentials are kept smaller than the values (24)), the electrostatic effect on the Particle trajectory may be neglected.
The Shepherd’s interaction with its image charge induced in its nearest bottom of the SEE experimental chamber does not lead to appreciable Particle displacements. A very demanding requirements on the Shepherd’s charge emerges, however, if the SEE satellite is used for G-dot determination (whose detailed discussion is postponed to future papers). One obtains then
$$U_M\underset{}{<}12\mathrm{mV}.$$
(25)
Evidently, in this case the Shepherd-Particle interaction per se is not the determining factor with respect to charge limits on the test bodies.
### 5.2. Evaluation of charges captured in some orbits
Let us now estimate the charges captured by the Shepherd and the Particle on board a satellite in various circular orbits for a single revolution around the Earth, a period of about two hours. Actual measurement times may exceed this period, but not too much.
Approximate values of time-averaged proton fluxes are presented for some circular orbits in Table 5.
The fluxes in Table 5 have been obtained using the computation software worked out at Nuclear Physics Institute (NPI) of Moscow State University, called SEE2 (Space Environment Effects 2) and SEREIS (Space Environment Radiation Effects Information System) . This software made use of the NASA models AP8-max and AP8-min for calculating the proton fluxes ; however, the latter rest on measurements performed in the solar maximum of 1970 and minimum of 1964, while the NPI software uses some modern models of the Earth’s magnetosphere, taking into account its evolution on the scale of decades.
The high-energy particle fluxes in the radiation belts are strongly time-dependent; they vary between maxima and minima of solar activity, being, at least at low altitudes relevant for a SEE mission, greater at solar minima \[2–5\]. Table 5 shows the fluxes at a solar minimum; similar calculations for a solar maximum show smaller values by at average 20-25 per cent; the difference exceeds this value only for $`x=0`$ (being thus greater for low-energy protons than for higher-energy ones).
The solar activity varies from one maximum or minimum to another, the Earth’s magnetic field is sensitive to all these variations and also varies due to certain terrestrial phenomena. Another source of uncertainty, probably not a very strong one, is that the shielding effect is calculated by SEE2 software for a detector placed at the centre of a spherical shell of shielding material, whereas the capsule is cylindrical and the angular distribution of the proton flow is also uncertain. It is thus clear that any values like those presented in Table 5 may only serve as a guide, giving correct orders of magnitude.
The above data make it possible to evaluate the captured charges. The results for two orbits, namely, 1500b and 1500c, are presented in Table 6. These are the worst and the best variants of orbits at 1500 km among those analyzed (see the caption of Table 5).
These and other similar data lead to some conclusions of importance for the SEE experiments.
First, the models show zero proton fluxes in equatorial orbits of 500 — 800 km altitudes but indicate considerable fluxes at the same altitudes due to crossing the SAA. It turns out, however, that the SAA is overwhelmingly a low-energy phenomenon and almost does not affect fluxes on the relevant energy scale beginning with approximately 65 MeV. Even more, there is a very small proton flux due to the SAA even at energies over 10 MeV, so that behind a layer of 1 mm the SAA influence is already negligible. Therefore, behind a thicker metal layer there are actually no secondary particles due to SAA protons.
Second, at 1500 km altitude the fluxes substantially depend on the orbit orientation but remain on the same scale of a few million protons per cm<sup>2</sup> at energies over 65 MeV.
Third, evidently, at 3000 km altitude both the total flux (for $`x=0`$) and especially its high-energy part are a few times greater than at 1500 km.
Fourth, and most important: for all orbits in the desirable range of altitudes the charges are quite large as compared with their admissible values; they remain large even behind rather thick walls. It is thus quite necessary to have means to detect and remove the charges during the measurements. Moreover, as seen from the peak values in Table 6 and from time scans of Van Allen charging in orbits of interest (also obtained using the above-mantioned software), at a charging peak when crossing the magnetic equator the time required for the charge on the test bodies to reach its maximum allowable values, as listed above, is a matter of seconds, not minutes. Therefore the charge must be detected and removed as it builds up, on a time scale of seconds.
The detection and measurement of the charge on the test bodies can probably be achieved relatively easily by an array of minute microvoltmeters attached to the inner wall of the experimental chamber.
Several methods for removing positive charge are now being evaluated. A simple and promising method may be to shoot electron beams directly at test bodies. The number of electrons needed is on the order of $`10^8/`$sec. Although this approach has the inherent drawback that it requires that an active system must perform correctly for many years, it is simple in principle and will accomplish the goal.
## 6. Conclusion
Space offers the prospect of quantum leaps in the accuracy of gravitational experiments. Although space is a challenging environment for research, the inherent quiet of space can be exploited to make very accurate determinations of $`G`$ and other gravitational parameters, providing that care is taken to understand the many physical phenomena in space which have the potential to vitiate accuracy. A distinctive feature of a SEE mission is its capability to perform such determinations simultaneously on multiple parameters, making it one of the most promising proposals.
To conclude, we enumerate different SEE tests and measurements and show their expected accuracy as currently estimated:
| Test/measurement | Expected accuracy |
| --- | --- |
| EP/ISL at a few metres | $`2\mathrm{\hspace{0.17em}10}^7`$ |
| EP/CD at a few metres | $`<10^7(\alpha <10^4)`$ |
| EP/ISL at $`R_{}`$ | $`<10^{10}`$ |
| EP/CD at $`R_{}`$ | $`<10^{16}(\alpha <10^{13})`$ |
| $`G`$ | $`3.3\mathrm{\hspace{0.17em}10}^7`$ |
| $`\dot{G}/G`$ | $`<10^{13}`$ in one year |
The last estimate is only tentative; the subject is under study.
### Acknowledgement
This work was supported in part by NASA grant # NAG 8-1442. K.A.B. wishes to thank Nikolai V. Kuznetsov for helpful discussions and for providing an access to the SEE2 and SEREIS software. |
warning/0002/astro-ph0002127.html | ar5iv | text | # MeV-scale Reheating Temperature and Thermalization of Neutrino Background
## I Introduction
In the standard big bang cosmology it had been assumed tacitly that the universe was dominated by the thermal radiation at the early epoch. Even in the paradigm of the modern cosmology it is commonly believed that thermal radiation was produced by the reheating process after the primordial inflation and they dominated the energy of the universe at sufficiently early epoch. Here we ask, “How early should the universe be dominated by radiation in order to success the standard big bang cosmology?”. We could say that the energy of the universe should be dominated by the radiation at least before the beginning of the big bang nucleosynthesis (BBN) epoch. In this paper we answer the above question.
The various models of the modern particle physics beyond the standard model predicts a number of unstable massive particles which have long lifetimes and decays at about BBN epoch. The energy density of the non-relativistic particles or the oscillation energy density of the scalar fields (inflaton and so on) decreases as $`\rho _{NR}(t)a(t)^3`$, where $`a(t)`$ is a scale factor. On the other hand since the radiation energy density decreases more rapidly $`\rho (t)a(t)^4`$, if the energy density of the massive non-relativistic particles or the oscillating scalar fields is large enough, it immediately dominates the universe as it expands, and the universe necessarily becomes matter-dominated until the cosmic time reaches to their lifetime. When the particles decay into standard particles (e.g. photon and electron), they produce the large entropy and the universe becomes radiation-dominated again. It is expected that such process would change the initial condition for the standard big bang scenario. We call the process “late-time entropy production”.
Now we have some interesting candidates for late-time entropy production in models based on supersymmetry (SUSY). It is known that gravitino and Polonyi field which exist in local SUSY (i.e. supergravity ) theories have masses of $`𝒪(100\mathrm{G}\mathrm{e}\mathrm{V}10\mathrm{T}\mathrm{e}\mathrm{V})`$ . In addition they have long lifetimes because they interact with the other particle only through gravity. For example, since Polonyi field which has a heavy mass $`10`$ TeV cannot be diluted by the usual inflation, it immediately dominates the universe and decays at the BBN epoch. Moreover it is also known that in the superstring theories there exist many light fields called dilaton and moduli which have similar properties to the Polonyi field.
Recently Lyth and Stewart considered a mini-inflation called “thermal inflation” which dilutes the above dangerous scalar fields. In the thermal inflation scenario, however, the flaton field which is responsible for the thermal inflation decays at late times. In particular, if Polonyi (moduli) mass is less than $`1`$ GeV which is predicted in the framework of gauge-mediated SUSY breaking models , the sufficient dilution requires that the flaton decays just before BBN . Thus, in thermal inflation models, one should take care of the late-time entropy production.
To keep the success of BBN, any long-lived massive particles or the coherent oscillation of any scalar fields which dominate the universe at that time must finally decay into the standard particles before the beginning of BBN. Moreover the decay products would have to be quickly thermalized through scatterings, annihilations, pair creations and further decays and make the thermal bath of photon, electron and neutrinos. Concerning photons and electrons which electromagnetically interact, the interaction rate is much more rapid than the Hubble expansion rate at that time. Therefore it is expected that the photon and electron which are produced in the decay and subsequent thermalization processes are efficiently thermalized. The problem is that neutrinos can interact only through the weak interaction. In the standard big bang cosmology the neutrinos usually decouple from the electromagnetic thermal bath at about $`T23`$MeV. Therefore it is approximately inferred that they can not be sufficiently thermalized at the temperature $`T`$ a few MeV. Namely the reheating temperature after the entropy production process should be high enough to thermalize the neutrinos. Though people had ever used the rough constraints on the reheating temperature between 1MeV - 10MeV, in the previous paper we pointed out that the neutrino thermalization is the most crucial for the successful BBN. In this paper we describe the detail of the method to obtain the neutrino spectrum and the formulations to integrate a set of Boltzmann equations numerically , and we study the constraint on the reheating temperature using the obtained neutrino spectrum and the full BBN network calculations with the revised observational light element abundances.
The above constraint is almost model-independent and hence conservative because we only assume that the massive particle decay produces the entropy. However, a more stringent constraint can be obtained if we assume a decay mode into quarks or gluons. In this case some modifications are needed for the above description. When the high energy quark-antiquark pairs or gluons are emitted, they immediately fragment into a lot of hadrons (pions , kaons, protons, neutrons, etc.). It is expected that they significantly influence the freeze-out value of neutron to proton ratio at the beginning of BBN through the strong interaction with the ambient protons and neutrons. In the previous paper we did not consider such hadron injection effects on BBN. Therefore we carefully treat the hadron injection effects in the present paper.
For another constraint, the late-time entropy production may induce the significant effects on the anisotropies of the cosmic microwave background radiations (CMB). Lopez et al. pointed out that the CMB anisotropies are very sensitive to the equal time of matter and radiation. When the reheating temperature is so low that neutrinos do not be sufficiently thermalized, the radiation density which consists of photon and neutrinos becomes less than that in the standard big bang scenario. It may give distinguishable signals in the CMB anisotropies as a modification of the effective number of neutrino species $`N_\nu ^{\mathrm{eff}}`$. With the present angular resolutions and sensitivities of COBE observation it is impossible to set a constraint on $`N_\nu ^{\mathrm{eff}}`$ but it is expected that future satellite experiments such as MAP and PLANCK will gives us a useful information about $`N_\nu ^{\mathrm{eff}}`$. In addition the above effect may also induce the signals in the observed power spectrum of the density fluctuation for the large scale structure as a modification of the epoch of the matter-radiation equality.
The paper is organized as follows. In Sec. II we introduce the formulation of the basic equations and the physical parameters. In Sec. III we briefly review the current status of the observational light element abundances. In Sec. IV we study the spectra of the electron neutrino and the mu(tau)-neutrino by numerically solving the Boltzmann equations, and the constraints from BBN are obtained there. In Sec V we investigate the additional effects in the hadron injection by the massive particle decay. In Sec. VI we consider the another constraints which come from observations for large scale structures and anisotropies of CMB. Sec VII is devoted to conclusions. In Appendix we introduce the method of the reduction for the nine dimension integrals into one dimension.
## II Formulation of Neutrino Thermalization
### A Reheating Temperature
In order to discuss the late-time entropy production process, we should formulate the equations which describe the physical process. Here the reheating temperature $`T_R`$ is an appropriate parameter to characterize the late-time entropy production. We define the reheating temperature $`T_R`$ by
$$\mathrm{\Gamma }3H(T_R),$$
(1)
where $`\mathrm{\Gamma }`$ is the decay rate (=$`\tau ^1`$) and $`H(T_R)`$ is the Hubble parameter at the decay epoch (t = $`\tau `$). <sup>*</sup><sup>*</sup>*Since the actual decay is not instantaneous, the matter-dominated universe smoothly changes into radiation-dominated one. Thus it is rather difficult to clearly identify the reheating temperature by observing the evolution of the cosmic temperature. Instead we “define” the reheating temperature formally by Eq. (1) The Hubble parameter is expressed by
$$H=\left(\frac{g_{}\pi ^2}{90}\right)^{1/2}\frac{T_R^2}{M_G},$$
(2)
where $`g_{}`$ is the statistical degrees of freedom for the massless particles and $`M_G`$ is the reduced Plank mass ($`=2.4\times 10^{18}`$GeV). Then the reheating temperature is given by
$$T_R=0.554\sqrt{\mathrm{\Gamma }M_G}.$$
(3)
Here we have used $`g_{}=43/4`$. From Eq. (3), we can see that the reheating temperature has the one to one correspondence with the lifetime of the parent massive particle.
Here we define the effective number of neutrino species $`N_\nu ^{\mathrm{eff}}`$ as a parameter which characterize the time evolution of the energy density of neutrinos. Here $`N_\nu ^{\mathrm{eff}}`$ is defined by
$$N_\nu ^{\mathrm{eff}}\frac{\rho _{\nu _e}+\rho _{\nu _\mu }+\rho _{\nu _\tau }}{\rho _{\mathrm{std}}},$$
(4)
where $`\rho _{\mathrm{std}}`$ is the total neutrino energy density in the standard big bang model (i.e. no late-time entropy production and three neutrino species).
### B Basic Equations
When a massive particle $`\varphi `$ which is responsible for the late-time entropy production decays, all emitted particles except neutrinos are quickly thermalized and make a thermal bath with temperature $`T_R`$. For relatively low reheating temperature $`T_R10`$MeV neutrinos are slowly thermalized. Since in realistic situations the decay branching ratio into neutrinos is very small, we assume that neutrinos are produced only through annihilations of electrons and positrons, i.e. $`e^++e^{}\nu _i+\overline{\nu }_i(i=e,\mu ,\tau )`$. The evolution of the distribution function $`f_i`$ of the neutrino $`\nu _i`$ is described by the momentum dependent Boltzmann equation :
$$\frac{f_i(𝒑,t)}{t}H(t)𝒑\frac{f_i(𝒑,t)}{𝒑}=C_{i,\mathrm{coll}}$$
(5)
where the right hand side is the total collision term. The integrated Boltzmann equation is not adequate in the present problem. As we show in Sec. IV, the spectral shape of the momentum distribution obtained by our scheme is much different from the equilibrium one. It should be noticed that the integrated Boltzmann equation assumes that the shape of the momentum distribution is the same as the equilibrium one. Thus we should solve the momentum dependent Boltzmann equation. When the reaction is two bodies scattering $`1+23+4`$, it is given by the expression,
$`C_{i,\mathrm{coll}}={\displaystyle \frac{1}{2E_1}}{\displaystyle }`$ $`{\displaystyle \frac{d^3p_2}{2E_2(2\pi )^3}}{\displaystyle \frac{d^3p_3}{2E_3(2\pi )^3}}{\displaystyle \frac{d^3p_4}{2E_4(2\pi )^3}}`$ (7)
$`\times (2\pi )^4\delta ^{(4)}(p_1+p_2p_3p_4)\mathrm{\Lambda }(f_1,f_2,f_3,f_4)S|M|_{1234}^2,`$
where $`|M|^2`$ is the scattering amplitude summed over spins of all particles, $`S`$ is the symmetrization factor which is 1/2 for identical particles in initial and final states, $`\mathrm{\Lambda }=f_3f_4(1f_1)(1f_2)f_1f_2(1f_3)(1f_4)`$ is the phase space factor including Pauli blocking of the final states. Then the total collision term $`C_{i,\mathrm{coll}}`$ is expressed by,
$$C_{i,\mathrm{coll}}=C_{i,\mathrm{ann}}+C_{i,\mathrm{scat}},$$
(8)
where $`C_{i,\mathrm{ann}}`$ is the collision term for annihilation processes and $`C_{i,\mathrm{scat}}`$ is collision term for elastic scattering processes. Here we consider the following processes:
$`\nu _i+\nu _i`$ $``$ $`e^++e^{},`$ (9)
$`\nu _i+e^\pm `$ $``$ $`\nu _i+e^\pm .`$ (10)
In this paper we have treated neutrinos as Majorana ones (i.e., $`\nu =\overline{\nu }`$). It should be noted that there are no differences between Majorana neutrinos and Dirac ones as long as they are massless, and since the temperature is $`𝒪(\mathrm{MeV})`$ at least in this situation, we could have treated them as if they were massless particles. The relevant reactions are presented in Table I for $`\nu _e`$ and Table II for $`\nu _\mu `$ and $`\nu _\tau `$Here we neglect the neutrino self-interactions. It may lead to underestimate the kinetic equilibrium rate for high reheating temperatures. However, we think that this effect does not change the results very much. The interactions between electrons and neutrinos are the most important because they transfer the energy of the thermal bath to neutrinos. The self-interactions of the neutrinos cannot increase the energy density of neutrinos but mainly change their momentum distribution. Furthermore, the neutrino number densities are much smaller than the electron number density at low reheating temperature with which we are concerned. Thus differences caused by the neutrino self interactions are expected to be small.
The collision terms are quite complicated and expressed by nine dimensional integrations over momentum space. However, if we neglect electron mass and assume that electrons obey the Boltzmann distribution $`e^{p/T}`$, the collision terms are simplified to one dimensional integration form. <sup>§</sup><sup>§</sup>§The errors due to neglecting the electron mass is small and the deviation is just a few percent. We show the reasons as follows. The difference between Fermi-Dirac and Maxwell-Boltzmann distribution “$`df`$” is less than one at most $`df<1.0`$. The week interaction rate is almost expressed by $`\sigma vn_e/H(t)`$, where $`\sigma vG_F^2m_e^2`$ and $`n_e`$ is an electron number density. Then the error is at most estimated by, $`\sigma _Wvn_e/H(t)\times df10^2`$ (for $`T`$ 0.5MeV). Therefore the deviation is a few percent and the neglecting the electron mass does not change the results. The other methods of the approximation to reduce the integral from nine to two dimensions in which the electron mass is not neglected are presented in ref. Then $`C_{i,\mathrm{ann}}`$ is given by
$$C_{i,\mathrm{ann}}=\frac{1}{2\pi ^2}p_i^2𝑑p_i^{}(\sigma v)_i(f_i(p_i)f_i(p_i^{})f_{eq}(p_i)f_{eq}(p_i^{})),$$
(11)
where $`f_{eq}(=1/(e^{p_i/T}+1))`$ is the equilibrium distribution and $`(\sigma v)_i`$ is the differential cross sections given by
$`(\sigma v)_e`$ $`=`$ $`{\displaystyle \frac{4G_F^2}{9\pi }}(C_V^2+C_A^2)pp^{},`$ (12)
$`(\sigma v)_{\mu ,\tau }`$ $`=`$ $`{\displaystyle \frac{4G_F^2}{9\pi }}(\stackrel{~}{C_V}^2+\stackrel{~}{C_A}^2)pp^{},`$ (13)
where we take $`C_V=\frac{1}{2}+2\mathrm{sin}^2\theta _W`$, $`C_A=\frac{1}{2}`$, $`\stackrel{~}{C_V}=C_V1`$ ($`\stackrel{~}{C_A}=C_A1`$) and $`\theta _W`$ is Weinberg angle ($`\mathrm{sin}^2\theta _W0.231`$.
As for elastic scattering processes, $`C_{i,scat}`$ is also simplified to one dimensional integration (see Appendix), and it is expressed as
$`C_{i,scat}`$ $`=`$ $`{\displaystyle \frac{G_F^2}{2\pi ^3}}(C_V^2+C_A^2)[{\displaystyle \frac{f_i}{p_i^2}}({\displaystyle _0^{p_i}}dp_i^{}F_1(p_i,p_i^{})(1f_i(p_i^{}))+{\displaystyle _{p_i}^{\mathrm{}}}dp_i^{}F_2(p_i,p_i^{})(1f_i(p_i^{}))`$ (16)
$`+{\displaystyle \frac{1f_i(p_i)}{p_i^2}}({\displaystyle _0^{p_i}}dp_i^{}B_1(p_i,p_i^{})f_i(p_i^{})+{\displaystyle _{p_i}^{\mathrm{}}}dp_i^{}B_2(p_i,p_i^{})f_i(p_i^{}))],`$
where $`(C_V^2+C_A^2)`$ should be replaced by $`(\stackrel{~}{C_V}^2+\stackrel{~}{C_A}^2)`$ for $`i=\mu ,\tau `$, and the functions $`F_1,F_2,B_1,B_2`$ are given by
$`F_1(p,p^{})`$ $`=`$ $`D(p,p^{})+E(p,p^{})e^{p^{}/T},`$ (17)
$`F_2(p,p^{})`$ $`=`$ $`D(p^{},p)e^{(pp^{})/T}+E(p,p^{})e^{p^{}/T},`$ (18)
$`B_1(p,p^{})`$ $`=`$ $`F_2(p^{},p),B_2(p,p^{})=F_1(p^{},p),`$ (19)
where
$`D(p,p^{})`$ $`=`$ $`2T^4(p^2+p^2+2T(pp^{})+4T^2),`$ (20)
$`E(p,p^{})`$ $`=`$ $`T^2[p^2p^2+2pp^{}(p+p^{})T`$ (22)
$`+2(p+p^{})^2T^2+4(p+p^{})T^3+8T^4].`$
Together with the above Boltzmann equations, we should solve the energy-momentum conservation equation in the expanding universe:
$$\frac{d\rho (t)}{dt}=3H(t)(\rho (t)+P(t)),$$
(23)
where $`\rho (t)=\rho _\varphi +\rho _\gamma +\rho _e+\rho _\nu `$ is the total energy density of $`\varphi `$, photon, electron and neutrinos and it is given by
$$\rho (t)=\rho _\varphi (t)+\frac{\pi ^2T_\gamma ^4}{15}+\frac{2}{\pi ^2}\frac{dqq^2E_e}{\mathrm{exp}(E_e/T_\gamma )+1}+\frac{1}{\pi ^2}𝑑qq^3f_{\nu _e}(q)+\frac{2}{\pi ^2}𝑑qq^3f_{\nu _\mu }(q),$$
(24)
where $`E_e=\sqrt{q^2+m_e^2}`$ is the electron energy. $`P(t)P_\gamma (t)+P_{e^\pm }(t)+P_\nu (t)`$ is the total pressure,
$$P(t)=\frac{\pi ^2T_\gamma ^4}{45}+\frac{2}{\pi ^2}\frac{dqq^4}{3E_e[\mathrm{exp}(E_e/T_\gamma )+1]}+\frac{1}{3\pi ^2}𝑑qq^3f_{\nu _e}(q)+\frac{2}{3\pi ^2}𝑑qq^3f_{\nu _\mu }(q).$$
(25)
$`H(t)`$ is the Hubble parameter,
$$H(t)=\frac{\dot{a}(t)}{a(t)}=\frac{1}{\sqrt{3}M_G}\sqrt{\rho (t)}.$$
(26)
The time evolution equation of $`\rho _\varphi `$ is given by
$$\frac{d\rho _\varphi }{dt}=\mathrm{\Gamma }\rho _\varphi 3H\rho _\varphi .$$
(27)
Practically we solve the time evolution of the photon temperature instead of Eq. (23),
$`{\displaystyle \frac{dT_\gamma }{dt}}`$ $`=`$ $`{\displaystyle \frac{\rho _\varphi /\tau _\varphi +4H\rho _\gamma +3H(\rho _{e^\pm }+P_{e^\pm })+4H\rho _\nu +d\rho _\nu /dt}{\rho _\gamma /T_\gamma |_{a(t)}+\rho _{e^\pm }/T_\gamma |_{a(t)}}},`$ (28)
together with Eqs. (5), (26) and (27).
## III Observational light element abundances
In this section we briefly show the current status of the observational light element abundances. Concerning the deuterium abundance, the primordial D/H is measured in the high redshift QSO absorption systems. For the most reliable D abundance, we adopt the following value which is obtained by the clouds at z = 3.572 towards Q1937-1009 and at z = 2.504 towards Q1009+2956 ,
$$\mathrm{D}/\mathrm{H}=(3.39\pm 0.25)\times 10^5.$$
(29)
On the other hand, recently the high deuterium abundance is reported in relatively low redshift absorption systems at z = 0.701 towards Q1718+4807 , $`\mathrm{D}/\mathrm{H}=(2.0\pm 0.5)\times 10^4`$. Another group also observes the clouds independently . However, because they do not have full spectra of the Lyman series, the analyses would be unreliable. More recently Kirkman et al. observed the quasar absorption systems at z = 2.8 towards Q0130-4021 and they obtain the upper bound, D/H $`6.7\times 10^5`$. Moreover Molaro et al. reported D/H $`1.5\times 10^5`$ which was observed in the absorber at z = 3.514 towards APM 08279+5255 although it has the large systematic errors in the hydrogen column density . Considering the current situation, we do not adopt the high deuterium value in this paper.
The primordial <sup>4</sup>He mass fraction $`Y_p`$ is observed in the low metalicity extragalactic HII regions. Since <sup>4</sup>He is produced with Oxygen in the star, the primordial value is obtained to regress to the zero metalicity O/H $`0`$ for the observational data. Using the 62 blue compact galaxies (BCG) observations, it was reported that the primordial $`Y`$ is relatively “ low”, $`Y_p0.234`$ . However, recently it is claimed that HeI stellar absorption is an important effect though it was not included in the previous analysis properly. They found the relatively “high” primordial value, $`Y_p=0.245\pm 0.004`$. More recently Fields and Olive also reanalyze the data including the HeI absorption effect and they obtain
$$Y_p=0.238\pm (0.002)_{stat}\pm (0.005)_{syst},$$
(30)
where the first error is the statistical uncertainty and the second error is the systematic one. We adopt the above value as the observational $`Y_p`$.
The primordial <sup>7</sup>Li/H is observed in the Pop II old halo stars. In general a halo star whose surface effective temperature is low (the mass is small), has the deep convective zone. For such a low temperature star, the primordial <sup>7</sup>Li is considerably depleted in the warm interior of the star. On the other hand for the high temperature stars ($`T_{eff}5500`$K), it is known that the primordial abundance is not changed and they have a “plateau”of the <sup>7</sup>Li as a function of the effective temperature. In addition, though it is also known that <sup>7</sup>Li/H decreases with decreasing Fe/H, <sup>7</sup>Li still levels off at lower metalicity, \[Fe/H\]$`1.5`$, in the plateau stars. We adopt the recent measurements which are observed by Bonifacio and Molaro . They observed 41 old halo stars which have the plateau. We take the additional larger systematic error, because there may be underestimates in the stellar depletion and the production by the cosmic ray spallation. Then we obtain
$$\mathrm{log}_{10}(^7\mathrm{Li}/\mathrm{H})=9.76\pm (0.012)_{stat}\pm (0.05)_{syst}\pm (0.3)_{add}.$$
(31)
## IV Neutrino Thermalization and BBN
### A Time evolution of Neutrino spectrum
The evolution of the cosmic temperature $`T`$ is shown in Fig. 1 (a) for $`T_R=10`$ MeV, and (b) for $`T_R=2`$ MeV. In Fig. 1 (a), it is seen that the temperature decreases slowly as $`t^{1/4}`$, i.e. $`a^{3/8}`$ before the decay epoch, $`t\mathrm{\Gamma }^1(5\times 10^2\mathrm{sec})`$ which corresponds to $`T_R=10`$ MeV. This is because the actual decay is not instantaneous and $`\varphi `$ decays into radiation continuously at the rate $`\mathrm{\Gamma }`$ . Then The universe is still in M.D. After the decay epoch $`t\mathrm{\Gamma }^1`$, all $`\varphi `$-particles decay and the temperature decreases as $`a^1`$ and $`t^{1/2}`$. Then the universe becomes radiation-dominated epoch. Since at the temperature $`T0.5`$MeV ($`t3\mathrm{sec}`$), electrons and positrons annihilate into photons $`e^+e^{}2\gamma `$, the temperature is slightly heated. From Fig. 1 (b) we can see that the temperature decreases as $`t^{1/4}`$ until the decay epoch ($`t0.1\mathrm{sec})`$ which corresponds to $`T_R=2`$ MeV. After the decay epoch, the temperature decreases as $`t^{1/2}`$ (R.D.). In the actual computation we take a initial condition that there exists the net radiation energy density though the universe is in M.D. This represents the situation that the massive particle necessarily dominated the universe as it expands. On the other hand even if there are at first no radiation $`\rho _R0,i.e.T0`$ which corresponds to the initial condition of the oscillation epoch after the primordial inflation or thermal inflation, the cosmic temperature immediately tracks the same curve $`t^{1/4}`$ and then their decay establish the radiation dominated universe $`Tt^{1/2}`$. Therefore our treatment is quite a general picture for each entropy production scenario and it does not depend on whether the net initial radiation energy exists or not, only if once the unstable non-relativistic particles dominate the energy density of the universe.
In Fig. 2 we show the evolutions of $`\rho _{\nu _e}`$ and $`\rho _{\nu _\mu }`$ (=$`\rho _{\nu _\tau }`$) (a) for $`T_R=10`$ MeV and (b)$`2`$ MeV. From Fig. 2(a) we can see that if $`T_R=10`$ MeV, cosmic energy density is as same as the case of standard big bang cosmology. As shown in Fig. 2(b), however, the energy density of each neutrino species for $`T_R=2`$ MeV is smaller than the case of standard scenario. Since the electron neutrinos interact with electrons and positrons through both charged and neutral currents, they are more effectively produced from the thermal bath than the other neutrinos which have only neutral current interactions. The final distribution functions $`f_e`$ and $`f_\mu (=f_\tau )`$ are shown in Fig. 3 (a) for $`T_R=10`$ MeV and (b) $`2`$ MeV. For $`T_R=10`$ MeV, each neutrino is thermalized well and the perfect Fermi-Dirac distribution is established. For $`T_R=2`$ MeV, however, the distributions are not thermal equilibrium forms. As we noted in Sec. II, we must not use the integrated Boltzmann equation instead of the momentum dependent Boltzmann equation in the present problem because the former assumes the equilibrium distribution. To see this, let us define the ratio $`R_E`$ for a neutrino species by $`R_E=(\rho _\nu /n_\nu )/(3.151\stackrel{~}{T_\nu })`$, where $`\rho _\nu `$ is the neutrino energy density, $`n_\nu `$ is the neutrino number density, $`\stackrel{~}{T_\nu }`$ is the effective neutrino temperature which is defined by the neutrino number density as, $`\stackrel{~}{T_\nu }\left(2\pi ^2/(3\zeta (3))n_\nu \right)^{1/3}`$. Here both $`\rho _\nu `$ and $`n_\nu `$ are computed by integrating the neutrino distribution function which is obtained by solving the momentum dependent Boltzmann equation. Approximately $`R_E`$ represents a ratio of the mean energy per a neutrino to the thermal equilibrium one. If the neutrino is in thermal equilibrium, $`R_E`$ is unity. In the case of the integrated Boltzmann equation, because it is assumed that the shape of the neutrino distribution is the same as the equilibrium one at any times, $`R_E`$ is necessarily unity. On the other hand, in the case of our scheme, i.e. the momentum dependent Boltzmann equation, $`R_E`$ can not be unity. We have computed the ratio $`R_E`$ in some representative reheating temperatures for electron neutrino and have found that they deviated from unity more at the lower reheating temperatures, $`R_E`$ = 1.00, 1.03 and 1.50 (for $`T_R`$ = 10 MeV, 3 MeV and 1 MeV). Moreover at the lower reheating temperature than 1 MeV, the deviation is much larger. This result tells us that the neutrino distribution deviates from the thermal equilibrium shape considerably at the low reheating temperatures and we should solve the momentum dependent Boltzmann equation. $`R_E`$ has a tendency to increase as the reheating temperature decreases. This is because neutrinos are produced by the annihilation of electrons-positron pairs whose mean energy per one particle is larger than that of neutrinos.
In Fig. 4 we can see the change of the effective number of neutrino species $`N_\nu ^{\mathrm{eff}}`$ as a function of the reheating temperature $`T_R`$. If $`T_R7`$ MeV, $`N_\nu ^{\mathrm{eff}}`$ is almost equal to three and neutrinos are thermalized very well. We can regard that it corresponds to the initial condition which has ever been used for the standard big bang cosmology. On the other hand, if $`T_R7`$ MeV, $`N_\nu ^{\mathrm{eff}}`$ becomes smaller than three.
### B Neutrino thermalization and neutron to proton ratio
If the neutrinos are not thermalized sufficiently and do not have the perfect Fermi-Dirac distribution, i.e. in this case there is the deficit of the neutrino distribution due to the low reheating temperature, it considerably influences the produced light element abundances. In particular, the abundance of the primordial <sup>4</sup>He is drastically changed. The change of the neutrino distribution function influences the neutrino energy density and the weak interaction rates between proton and neutron. At the beginning of BBN (T $``$ 1 MeV - 0.1 MeV) the competition between the Hubble expansion rate $`H`$ and the weak interaction rates $`\mathrm{\Gamma }_{np}`$ determines the freeze-out value of neutron to proton ratio $`n/p`$. After the freeze-out time, neutrons can change into protons only through the free decay with the life time $`\tau _n`$. Since <sup>4</sup>He is the most stable light element and the almost all neutrons are synthesized into <sup>4</sup>He, the abundance of the primordial <sup>4</sup>He is sensitive to the freeze-out value of neutron to proton ratio.
If the neutrino energy density gets smaller than that of the standard BBN (SBBN), Hubble expansion rate which is proportional to the square of the total energy density is also decreased. Then the freeze out time becomes later and the $`\beta `$ equilibrium between neutrons and protons continues for longer time. As a result less neutrons are left. In this case the predicted <sup>4</sup>He is less than the prediction of SBBN. The effect due to the speed-down expansion is approximately estimated by
$$\mathrm{\Delta }Y0.1(\mathrm{\Delta }\rho _{tot}/\rho _{tot}),$$
(32)
where $`Y`$ is the mass fraction of <sup>4</sup>He and $`\rho _{tot}`$ is the total energy density of the universe.
Moreover, when the electron neutrino is not thermalized, there is interesting effect by which more <sup>4</sup>He are produced. The weak reaction rates are computed by integrating neutrino distribution functions which are obtained by solving Boltzmann equations numerically. Using the neutrino distribution functions, the six weak interaction rates between neutron and proton are represented by
$`\mathrm{\Gamma }_{npe^{}\overline{\nu _e}}`$ $`=`$ $`K{\displaystyle _0^{Qm_e}}𝑑p_{\nu _e}\left[\sqrt{(p_{\nu _e}Q)^2m_e^2}(Qp_{\nu _e}){\displaystyle \frac{p_{\nu _e}^2}{1+e^{(p_{\nu _e}Q)/T_\gamma }}}\left(1f_{\nu _e}(p_{\nu _e})\right)\right],`$ (33)
$`\mathrm{\Gamma }_{ne^+p\overline{\nu _e}}`$ $`=`$ $`K{\displaystyle _{Q+m_e}^{\mathrm{}}}𝑑p_{\nu _e}\left[\sqrt{(p_{\nu _e}Q)^2m_e^2}(p_{\nu _e}Q){\displaystyle \frac{p_{\nu _e}^2}{e^{(p_{\nu _eQ})/T_\gamma }+1}}\left(1f_{\nu _e}(p_{\nu _e})\right)\right],`$ (34)
$`\mathrm{\Gamma }_{n\nu _epe^{}}`$ $`=`$ $`K{\displaystyle _0^{\mathrm{}}}𝑑p_{\nu _e}\left[\sqrt{(p_{\nu _e}+Q)^2m_e^2}(p_{\nu _e}+Q){\displaystyle \frac{p_{\nu _e}^2}{1+e^{(p_{\nu _e}+Q)/T_\gamma }}}f_{\nu _e}(p_{\nu _e})\right],`$ (35)
$`\mathrm{\Gamma }_{pe^{}\overline{\nu _e}n}`$ $`=`$ $`K{\displaystyle _0^{Qm_e}}𝑑p_{\nu _e}\left[\sqrt{(p_{\nu _e}Q)^2m_e^2}(Qp_{\nu _e}){\displaystyle \frac{p_{\nu _e}^2}{e^{(p_{\nu _eQ})/T_\gamma }+1}}f_{\nu _e}(p_{\nu _e})\right],`$ (36)
$`\mathrm{\Gamma }_{pe^{}n\nu _e}`$ $`=`$ $`K{\displaystyle _0^{\mathrm{}}}𝑑p_{\nu _e}\left[\sqrt{(p_{\nu _e}+Q)^2m_e^2}(Q+p_{\nu _e}){\displaystyle \frac{p_{\nu _e}^2}{e^{(p_{\nu _e+Q})/T_\gamma }+1}}\left(1f_{\nu _e}(p_{\nu _e})\right)\right],`$ (37)
$`\mathrm{\Gamma }_{p\overline{\nu _e}ne^+}`$ $`=`$ $`K{\displaystyle _{Q+m_e}^{\mathrm{}}}𝑑p_{\nu _e}\left[\sqrt{(p_{\nu _e}Q)^2m_e^2}(Qp_{\nu _e}){\displaystyle \frac{p_{\nu _e}^2}{1+e^{(p_{\nu _eQ})/T_\gamma }}}f_{\nu _e}(p_{\nu _e})\right],`$ (38)
where $`Q=m_nm_p=1.29`$ MeV and $`K`$ is a normalization factor which is determined by the neutron life time $`\tau _n`$ as $`K(1.636\tau _n)^1`$ and $`\tau _n`$ is obtained by the experiments . From the above equations we can see that if neutrino and anti-neutrino distribution functions are decreased, both $`\beta `$ decay rates $`\mathrm{\Gamma }_{np}=\mathrm{\Gamma }_{npe^{}\overline{\nu _e}}+\mathrm{\Gamma }_{ne^+p\overline{\nu _e}}+\mathrm{\Gamma }_{n\nu _epe^{}}`$ and $`\mathrm{\Gamma }_{pn}=\mathrm{\Gamma }_{pe^{}\overline{\nu _e}n}+\mathrm{\Gamma }_{pe^{}n\nu _e}+\mathrm{\Gamma }_{p\overline{\nu _e}ne^+}`$ are simultaneously decreased by the following reasons. The dominant effects by the deficit of the distribution functions are to decrease the rates $`\mathrm{\Gamma }_{n\nu _epe^{}}`$, $`\mathrm{\Gamma }_{pe^{}\overline{\nu _e}n}`$ and $`\mathrm{\Gamma }_{p\overline{\nu _e}ne^+}`$ which have the neutrino or anti-neutrino in the initial state. On the other hand, though the other rates $`\mathrm{\Gamma }_{npe^{}\overline{\nu _e}}`$, $`\mathrm{\Gamma }_{ne^+p\overline{\nu _e}}`$ and $`\mathrm{\Gamma }_{pe^{}n\nu _e}`$ which have the neutrino or anti-neutrino in the final state are slightly increased due to Fermi-blocking factor $`(1f_\nu )`$, the ratio of the difference $`\mathrm{\Delta }f_\nu `$ to $`(1f_\nu )`$ is much smaller than that of $`\mathrm{\Delta }f_\nu `$ to $`f_\nu `$, i.e.
$$|\mathrm{\Delta }f_\nu /(1f_\nu )||\mathrm{\Delta }f_\nu /f_\nu |\mathrm{for}f_\nu 1.$$
(39)
Therefore, the enhancement is small and the latter effect is not important. In total, both weak interaction rates $`\mathrm{\Gamma }_{np}`$ and $`\mathrm{\Gamma }_{pn}`$ are decreased and become smaller than those of SBBN. In Fig. 5 the weak interaction rates $`\mathrm{\Gamma }_{np}`$ and $`\mathrm{\Gamma }_{pn}`$ are plotted. The solid lines denote the case of $`T_R=10`$ MeV which corresponds to the standard big bang scenario. The dotted lines denote the case of $`T_R=1`$ MeV. In the plot we can see that the insufficient thermalization of the neutrino distributions derives the changes of the weak interaction rates.
The decrease of weak interaction rates gives significant effects on the abundance of $`{}_{}{}^{4}\mathrm{He}`$. When the weak interaction rate $`\mathrm{\Gamma }_{np}`$ decreases, Hubble expansion rate becomes more rapid than that of the weak interaction rate earlier. Namely the freeze-out time becomes earlier. Then the freeze-out value of neutron to proton ratio becomes larger than in SBBN and it is expected that the predicted <sup>4</sup>He abundance becomes larger. The above effect is approximately estimated by
$$\mathrm{\Delta }Y+0.19(\mathrm{\Delta }\mathrm{\Gamma }_{np}/\mathrm{\Gamma }_{np})$$
(40)
In Fig. 6 we plot the time evolution of the neutron to proton ratio. In Fig. 6(a) we change only the number of neutrino species in SBBN. The dotted line denoted the curve of $`N_\nu ^{\mathrm{eff}}=1.37`$ which corresponds to the effective number of neutrino species in the case of $`T_R=2`$ MeV in the late-time entropy production scenario. Then we find that the predicted $`n/p`$ curve is lower than that of $`N_\nu ^{\mathrm{eff}}=3`$ due to only the speed down effects or the later decoupling. In Fig. 6(b) we plot the time evolution of $`n/p`$ when we change the reheating temperature in the late-time entropy production scenario. The dotted line denotes the case of $`T_R=2`$ MeV. Comparing to the case of $`N_\nu ^{\mathrm{eff}}=1.37`$ in Fig. 6(a), the $`n/p`$ ratio becomes larger. It is because the weak interaction rates are decreased by the deficit of the distribution function. Moreover in the case of $`T_R=1`$ MeV the $`n/p`$ ratio becomes much larger.
### C Neutrino thermalization and light element abundances
Next we perform Monte Carlo simulation and the maximum likelihood analysis to discuss how the theoretical predictions with the low reheating temperature scenario agree with the observational light element abundances.
In Fig. 7 we plotted <sup>4</sup>He mass fraction Y as a function of $`T_R`$ at $`\eta =5\times 10^{10}`$ (solid line). The dashed line denotes the virtual <sup>4</sup>He mass fraction computed by including only the speed down effect due to the change of the effective number of neutrino species which is shown in Fig. 4. The dotted line denotes the predicted value of Y in SBBN at $`\eta =5\times 10^{10}`$. For $`T_R7`$MeV, the solid line and dashed line are quite equal to the value in SBBN. As $`T_R`$ decreases, both the solid and dashed lines gradually decrease because of the speed down effect due to the change of $`N_\nu ^{\mathrm{eff}}`$. The dashed line continues to decrease as the reheating temperature decreases.
On the other hand, for T$`{}_{R}{}^{}2`$ MeV the effect that the weak interaction rates are weakened due to the deficit of the neutrino distribution function begins to become important and the predicted value of $`Y`$ begins to increase as $`T_R`$ decreases. For $`T_R1`$ MeV, since it is too late to produce enough electrons whose mass is about $`m_e`$ = 0.511 MeV, the weak interaction rates are still more weakened and $`Y`$ steeply increases as $`T_R`$ decreases.
In Fig. 8 we plot the contours of the confidence level in the $`\eta `$-$`T_R`$ plane. The solid line denotes 95 $`\%`$ C.L. and the dotted line denotes 68 $`\%`$ C.L. The filled square is the best fit point between the observation and theoretical predictions. The observational data are consistent with the high baryon to photon ratio, $`\eta (36)\times 10^{10}`$. From Fig. 8 we find that $`T_R0.7`$ MeV is excluded at 95 $`\%`$ C.L. In other wards $`T_R`$ as low as $`0.7\mathrm{MeV}`$ is consistent with BBN. Then $`N_\nu ^{\mathrm{eff}}`$ can be as small as $`0.1`$ and it definitely influences the formation of the large scale structure and CMB anisotropy as is seen in Sec.VI.
## V Hadron injection by massive particle decay
### A Hadron Jets and $`e^+e^{}`$ collider experiments
In the previous section we discussed only the case in which the parent massive particle $`\varphi `$ decays into photons or the other electro-magnetic particles. In this section we consider the entropy production process along with the hadron injection, i.e. the case in which the massive particle has some decay modes into quarks or gluons. Then the emitted quark-antiquark pairs or gluons immediately fragment into hadron jets and as a result a lot of mesons and baryons, e.g. pions, kaons, nucleons (protons and neutrons ) are emitted into the electro-magnetic thermal bath which is constituted by photon, electron and nucleons.
For example, if the gravitino $`\psi _\mu `$ is the parent particle which produces the large entropy, it could have a hadronic decay mode (e.g. $`\psi _\mu \stackrel{~}{\gamma }q\overline{q}`$ ) with the branching ratio $`B_h𝒪`$($`\alpha `$) at least even if the main decay mode is only $`\psi _\mu \stackrel{~}{\gamma }\gamma `$($`\stackrel{~}{\gamma }`$ : photino) . Then about 0.6 - 3 hadrons are produced for $`m_\varphi 1100\mathrm{TeV}`$. In addition the emitted high energy photons whose energy is about $`m_\varphi /2`$ scatter off the background photons and could also produce the quark-antiquark pairs through the electromagnetic interaction. For the cosmic temperature $`𝒪(\mathrm{MeV})`$, the energy in the center of mass frame is $`\sqrt{s}220\mathrm{GeV}`$ for $`m_\varphi 1100\mathrm{TeV}`$. Then the number of the produced hadrons is about 2 - 7 which effectively corresponds to the hadron branching ratio $`B_h10^2`$ if we assume that the hadron fragmentation is similar to the results of $`e^+e^{}`$ collider experiments. Thus $`B_h`$ should not become less than about $`10^2`$ for gravitino decay. If the decay mode $`\psi _\mu \stackrel{~}{g}g`$ ($`\stackrel{~}{g}`$ : gluino) is kinematically allowed, the hadronic branching ratio becomes close to one. For the other candidate, if the “flaton” is the parent particle as in thermal inflation model, it would also have a hadronic decay mode ($`\varphi gg`$ if the flaton mass is larger than 1 GeV.
If once such hadrons are emitted to the electro-magnetic thermal bath in the beginning of BBN epoch (at $`T10\mathrm{MeV}0.1\mathrm{MeV})`$, they quickly transfer all the kinetic energy into the thermal bath through the electro-magnetic interaction or the strong interaction. Through such thermalization processes the emitted high energy hadrons scatter off the background particles, and then they induce some effects on BBN. Especially, the emitted hadrons extraordinarily inter-convert the ambient protons and neutrons each other through the strong interaction even after the freeze-out time of the neutron to proton ratio $`n/p`$. For the relatively short lifetime ($`\tau _\varphi 10^2\mathrm{sec}10^2\mathrm{sec}`$) in which we are interested, the above effect induces the significant change in the previous discussion. Namely protons which are more abundant than neutrons are changed into neutrons by the hadron-proton collisions and the ratio $`n/p`$ increases extremely. Because $`{}_{}{}^{4}\mathrm{He}`$ is the most sensitive to the freeze out value of $`n/p`$, the late-time hadron injection scenario tends to increase $`Y_p`$.
Reno and Seckel investigated the influences of the hadron injection on the early stage of BBN. They constrained the lifetime of the parent particle and the number density comparing the theoretical prediction of the light element abundances with the observational data. Here we basically follow their treatment and apply it to the scenario of late-time entropy production with hadron injections.
The emitted hadrons do not scatter off the background nucleons directly. At first hadrons scatter off the background photons and electrons because they are much more abundant than nucleons. For $`t200\mathrm{sec}`$, the emitted high energy hadrons are immediately thermalized through the electro-magnetic scattering and they reach to the kinetic equilibrium before they interact with the ambient protons and neutrons. Then we use the threshold cross section $`\sigma v_{NN^{}}^{H_i}`$ for the strong interaction process $`N+H_iN^{}+\mathrm{}`$ between hadron $`H_i`$ and the ambient nucleon $`N`$, where $`N`$ denotes proton $`p`$ or neutron $`n`$. The strong interaction rate is estimated by
$`\mathrm{\Gamma }_{NN^{}}^{H_i}`$ $`=`$ $`n_N\sigma v_{NN^{}}^{H_i}`$ (41)
$``$ $`10^8\mathrm{sec}^1f_N\left({\displaystyle \frac{\eta }{10^9}}\right)\left({\displaystyle \frac{\sigma v_{NN^{}}^{H_i}}{10\mathrm{mb}}}\right)\left({\displaystyle \frac{T}{2\mathrm{MeV}}}\right)^3,`$ (42)
where $`n_N`$ is the number density of the nucleon species $`N`$, $`\eta `$ is the baryon to photon ratio ($`=n_B/n_\gamma `$), $`n_B`$ denotes the baryon number density ($`=n_p+n_n`$) and $`f_N`$ is the nucleon fraction ($`n_N/n_B`$). This equation shows that every hadron whose lifetime is longer than $`𝒪(10^8)`$ sec contributes to the inter-converting interaction between neutron and proton at the beginning of BBN. Hereafter we will consider only the following long-lived hadrons, (mesons, $`\pi ^\pm `$, $`K^\pm `$ and $`K_L`$, and baryons, $`p`$, $`\overline{p}`$, $`n`$, and $`\overline{n}`$). For the relevant process ( $`N+\pi ^\pm N^{}\mathrm{}`$, and $`N+K^{}N^{}\mathrm{}`$, etc.), we can obtain the cross sections in . Here we ignore $`K^+`$ interaction because $`n+K^+p+K^0`$ is the endothermic reaction which has $`Q=2.8`$ MeV.
We estimate the average number of emitted hadron species $`H_i`$ per one $`\varphi `$ decay as
$$N^{H_i}=B_hN_{jet}f_{H_i}\frac{N_{ch}}{2},$$
(43)
where $`N_{ch}`$ is the averaged charged-particle multiplicity which represents the total number of the charged particles emitted per two hadron jets, $`f_{H_i}`$ is the number fraction of the hadron species $`H_i`$ to all the emitted charged particles, $`B_h`$ is the branching ratio of the hadronic decay mode and $`N_{jet}`$ is the number of the produced jets per one $`\varphi `$ decay.
Here it is reasonable to assume that the averaged charged particle multiplicity $`N_{ch}`$ is independent of the the source because the physical mechanism which governs the production of hadron jets is quite similar and does not depend on the detail of the origin only if the high energy quark-antiquark pairs or gluons are emitted. We adopt the data which are obtained by the $`e^+e^{}`$ collider experiments. LEPII experiments (ALEPH, DELPHI, L3 and OPAL) recently give us the useful data for $`\sqrt{s}`$ = 130 $``$ 172 GeV . We adopt the following fitting function for $`\sqrt{s}`$ = 1.4 $``$ 172 GeV ,
$$N_{ch}=1.73+0.268\mathrm{exp}\left(1.42\sqrt{\mathrm{ln}(s/\mathrm{\Lambda }^2)}\right),$$
(44)
where $`\sqrt{s}`$ denotes the center of mass energy, the functional shape is motivated by the next-to-leading order perturbative QCD calculations, $`\mathrm{\Lambda }`$ is the cut-off parameter in the perturbative calculations and we take $`\mathrm{\Lambda }=1`$ GeV. In Fig. 9 we plot the charged particle multiplicity for $`\sqrt{s}=1\mathrm{GeV}100\mathrm{TeV}`$. The error of the fitting is about 10$`\%`$. Using the available data , we obtain the emitted hadron fraction $`f_{H_i}n^{H_i}/N_{ch}`$,
$`f_{\pi ^+}=0.64,f_\pi ^{}=0.64,`$ (45)
$`f_{K^+}=0.076,f_K^{}=0.076,f_{K_L}=0.054`$ (46)
$`f_p=f_{\overline{p}}=0.035,f_n=f_{\overline{n}}=0.034,`$ (47)
where $`n^{H_i}`$ is the number of the emitted hadron species $`H_i`$ which is defined as the value after both $`K_S`$ and $`\mathrm{\Lambda }^0`$ had completely finished to decay. <sup>\**</sup><sup>\**</sup>\**Although the summation of $`f_{H_i}`$ is obviously more than one, it is because the experimental fitting of $`N_{ch}`$ is defined as a value before $`K_S`$ and $`\mathrm{\Lambda }^0`$ begin to decay . Here we assume that $`f_{H_i}`$ do not change significantly in the energy range $`\sqrt{s}`$ 10 GeV - 100 TeV. Since we do not have any experimental data for the high energy region more than about 200 GeV, we extrapolate $`N_{ch}`$ to the higher energy regions and we take $`f_{H_i}`$ as a constant. As we find easily, almost all the emitted particles are pions which are the lightest mesons. To apply the data of the $`e^+e^{}`$ collider experiments, we take $`\sqrt{s}=2E_{jet}`$ in $`N_{ch}`$ where $`E_{jet}`$ denotes the energy of one hadron jet because the $`N_{ch}`$ is obtained by the result for two hadron jets.
### B Formulation in Hadron Injection Scenario
In this section we formulate the time evolution equations in the late-time hadron injection scenario. As we mentioned in the previous section, the hadron injection at the beginning of BBN enhances the inter-converting interactions between neutron and proton equally and the freeze out value of $`n/p`$ can be extremely increased. Then the time evolution equations for the number density of a nucleon $`N(=p,n)`$ is represented by
$$\frac{dn_N}{dt}+3H(t)n_N=\left[\frac{dn_N}{dt}\right]_{weak}\mathrm{\Gamma }_\varphi n_\varphi \left(K_{NN^{}}K_{N^{}N}\right),$$
(48)
where $`H(t)`$ is Hubble expansion rate, $`[dn_N/dt]_{weak}`$ denotes the contribution from the weak interaction rates which are obtained by integrating the neutrino distribution functions as discussed in Sec. IV, see Eqs.(33 \- 38), $`n_\varphi =\rho _\varphi /m_\varphi `$ is the number density of $`\varphi `$, $`K_{NN^{}}`$ denotes the average number of the transition $`NN^{}`$ per one $`\varphi `$ decay.
The average number of the transition $`NN^{}`$ is expressed by
$$K_{NN^{}}=\underset{H_i}{}N^{H_i}R_{NN^{}}^{H_i},$$
(49)
where $`H_i`$ runs the hadron species which are relevant to the nucleon inter-converting reactions, $`N^{H_i}`$ denotes the average number of the emitted hadron species $`H_i`$ per one $`\varphi `$ decay which is given by Eq. (43) and $`R_{NN^{}}^{H_i}`$ denotes the probability that a hadron species $`H_i`$ induces the nucleon transition $`NN^{}`$,
$$R_{NN^{}}^{H_i}=\frac{\mathrm{\Gamma }_{NN^{}}^{H_i}}{\mathrm{\Gamma }_{dec}^{H_i}+\mathrm{\Gamma }_{abs}^{H_i}},$$
(50)
where $`\mathrm{\Gamma }_{dec}^{H_i}=\tau _{H_i}^1`$ is the decay rate of the hadron $`H_i`$ and $`\mathrm{\Gamma }_{abs}^{H_i}\mathrm{\Gamma }_{NN^{}}^{H_i}+\mathrm{\Gamma }_{N^{}N}^{H_i}+\mathrm{\Gamma }_{NN}^{H_i}+\mathrm{\Gamma }_{N^{}N^{}}^{H_i}`$ is the total absorption rate of $`H_i`$.
### C Hadron injection and BBN
In this subsection we compare the theoretical prediction of the light element abundances in the hadron injection scenario to the observational light element abundances. In the computations we assume that the massive particle decays into three bodies ($`E_{jet}=m_\varphi /3`$) and two jets are produced at the parton level ($`N_{jet}=2`$<sup>††</sup><sup>††</sup>††The above choice of a set of model parameters $`E_{jet}`$ and $`N_{jet}`$ is not unique in general and is obviously model dependent. However, since $`N_{ch}`$ has the logarithmic dependence of $`E_{jet}`$, we should not be worried about the modification of $`E_{jet}`$ by just a factor of two so seriously. On the other hand in Eq. (48), the second term in the right hand side scales as $`N_{jet}/m_\varphi `$. For the modification of $`N_{jet}`$, therefore, we only translate the obtained results according to the above scaling rule and push the responsibility off onto $`m_\varphi `$. In the computing we take the branching ratio of the hadronic decay mode $`B_h=𝒪(10^21)`$.
As we noted in the previous subsections, it is a remarkable feature that the predicted $`Y_p`$ tends to increase in the hadron injection scenario because $`{}_{}{}^{4}\mathrm{He}`$ is the most sensitive to the freeze-out value of the neutron to proton ratio. Since protons which are more abundant than neutrons are changed into neutrons through the strong interactions rapidly, the freeze out value of $`n/p`$ increase extremely if once the net hadrons are emitted. In Fig. 10 we plot the predicted $`{}_{}{}^{4}\mathrm{He}`$ mass fraction $`Y_p`$ as a function of $`T_R`$ for (a) $`m_\varphi `$=100 TeV and (b) $`m_\varphi `$ =10 GeV. The solid curve denotes the predicted $`Y_p`$. Here we take the branching ratio of the hadronic decay mode as $`B_h`$ = 1 (right one) and $`B_h`$ = 0.01 (left one). The dot-dashed line denotes $`B_h=0`$. The dashed line denotes the virtual value of $`Y_p`$ computed by including only the speed down effect due to the change of the effective number of neutrino species. The dotted line denotes the prediction in SBBN.
As we mentioned in the previous section, the speed down effect due to deficit of the electron neutrino distribution function are not important for $`T_R7`$MeV. In addition since it is high enough to keep $`n/p`$ $``$ 1 for the cosmic temperature $`T7`$MeV, the enhancements of the inter-converting interaction between n and p due to the hadron emission do not induce any changes on the freeze-out value of $`n/p`$. As $`T_R`$ decreases ($`T_R7`$MeV), $`Y_p`$ also decreases gradually because the speed down effect on the freeze-out value of $`n/p`$ begins to be important. On the other hand, if a lot of hadrons are emitted when the cosmic temperature is $`T`$ 6 - 7 MeV and the ratio $`n/p`$ is less than one, they enhance the inter-converting interactions more rapidly. As a result, the ratio $`n/p`$ attempts to get closer to one again although the cosmic temperature is still low. Thus the above effects extremely increase the freeze-out value of $`n/p`$ and is much more effective than the speed down effects. Namely the produced $`Y_p`$ becomes larger very sensitively only if $`T_R`$ is just a little lower than 6 - 7 MeV. One can obviously find that this effect becomes more remarkable for the larger $`B_h`$.
To understand how it depends on mass, it is convenient to introduce the yield variable $`Y_\varphi `$ which is defined by
$$Y_\varphi n_\varphi /s,$$
(51)
where $`s`$ denotes the entropy density in the universe. Because $`Y_\varphi `$ is a constant only while the universe expands without any entropy production, it represents the net number density of $`\varphi `$ per comoving volume. For the simplicity let’s consider the instantaneous decay of $`\varphi `$ and assume that the reheating process has been completed quickly. Because the radiation energy in the thermal bath or entropy $`s=2\pi ^2g_{}/45T_R^3`$ is produced only from the decay products of $`\varphi `$, $`Y_\varphi `$ is approximately estimated using $`T_R`$ and $`m_\varphi `$ by
$$Y_\varphi 0.28\frac{T_R}{m_\varphi }.$$
(52)
From the above equation, we can see that for the fixed value of $`T_R`$ the net number of $`\varphi `$, i.e. the net number of the emitted hadrons, becomes larger for the smaller mass. Comparing Fig. 10 (a) with Fig. 10 (b), we find that the theoretical curve of $`Y_p`$ for the case of $`m_\varphi `$ = 10 GeV is enhanced more steeply and the starting point to increase $`Y_p`$ becomes higher than for the case of $`m_\varphi `$ = 100 TeV.
Since the other elements (D and $`{}_{}{}^{7}\mathrm{Li}`$) are not so sensitive as $`{}_{}{}^{4}\mathrm{He}`$, it is expected that the observational value of $`Y_p`$ constrains $`T_R`$ most strongly. In order to discuss how a low reheating temperature is allowed by comparing the theoretical predictions with observational values (D, $`{}_{}{}^{4}\mathrm{He}`$ and $`{}_{}{}^{7}\mathrm{Li}`$), we perform the Monte Carlo simulation and maximum likelihood analysis as discussed in Sec. IV. In addition to the case of Sec. IV we take account of the following uncertainties, the error for the fitting of $`N_{ch}`$ as 10$`\%`$ and the experimental error for each cross section of the hadron interaction as 50$`\%`$. Because there are not any adequate experimental data for the uncertainties of cross sections , here we take the larger values to get a conservative constraint.
In Fig. 11 we plot the contours of the confidence level for $`m_\varphi =100`$ TeV in the ($`\eta `$-$`T_R`$) plane for (a) $`B_h`$ = 1 and (b) $`B_h=10^2`$. The solid line denotes 95 $`\%`$ C.L., the dotted line denotes 68 $`\%`$ C.L. and the filled square is the best fit point between the observation and theoretical prediction for D, $`{}_{}{}^{4}\mathrm{He}`$ and $`{}_{}{}^{7}\mathrm{Li}`$. The baryon to photon ratio which is consistent with the observational data is restricted in the narrow region, $`\eta (46)\times 10^{10}`$. From Fig. 11(a), we find that $`T_R3.7`$ MeV is excluded at 95 $`\%`$ C.L. for $`B_h`$= 1. On the other hand, from Fig. 11(b) we obtain the milder constraint that $`T_R2.5`$ MeV is excluded at 95 $`\%`$ C.L. for $`B_h=10^2`$. In Fig. 12 we plot the contours of the confidence level for $`m_\varphi =10`$ GeV in the same way as Fig. 11. Compared to Fig. 11, as we mentioned above, we find that the lower bound on the reheating temperature becomes higher for a smaller mass. From Fig. 12 we get the lower bound on the reheating temperature that $`T_R5.0`$ MeV (4.0 MeV) at 95 $`\%`$ C.L. for $`B_h`$= 1 ($`B_h=10^2`$)
In Fig. 13 the lower bound on $`T_R`$ as a function of $`m_\varphi `$ are plotted for (a) $`B_h`$ = 1 and (b) $`B_h=10^2`$. The solid line denotes 95 $`\%`$ C.L. and the dotted line denotes 68 $`\%`$ C.L. As it is expected, the curve of the lower bound on $`T_R`$ is a gentle monotonic decreasing function of $`m_\varphi `$. In Fig. 13(a), we can see that $`T_R`$ should be higher than 4 MeV at 95 $`\%`$ C.L. for $`B_h`$ = 1 in $`m_\varphi `$ = 10 GeV - $`10^2`$ TeV<sup>‡‡</sup><sup>‡‡</sup>‡‡Though we have adopted the experimental error of the each hadron interaction cross section as 50$`\%`$ in the Monte Carlo simulation because of no data, the lower bound on $`T_R`$ might become about 10$`\%`$ higher than the above values if we adopt the more sever experimental error as 10$`\%`$ instead of 50$`\%`$.. On the other hand, in Fig. 13(b) we find that the constraint gets milder for $`B_h=10^2`$. It is shown that $`T_R2.5`$ MeV is excluded at 95 $`\%`$ C.L. for $`B_h=10^2`$. In Fig. 4 we find that $`N_\nu ^{\mathrm{eff}}`$ can be allowed as small as 2.8 for $`B_h`$ = 1 ( 1.9 for $`B_h=10^2`$).
### D Summary of hadron injection
In this section we have seen that the BBN constraint on the reheating temperature becomes much more stringent if a massive particle has a branching to hadrons. For successful BBN the reheating temperature should be higher than $`2.54`$ MeV for the branching ratio $`B_h=110^2`$. The hadron injection generally occurs if the late-time reheating is caused by the heavy particle with mass larger than $`1`$ GeV. Many candidates for the late-time reheating such as SUSY particles and flatons have such large masses and hence the constraint obtained here is crucial in constructing particle physics models based on SUSY or thermal inflation models.
For the lower limit of the reheating temperature, the effective number of the neutrino species $`N_\nu ^{\mathrm{eff}}`$ is given by 2.8 and 1.9 for $`B_h`$ = 1 and $`10^2`$, respectively. Since the limiting temperature is close to the neutrino decoupling temperature, the deviation of $`N_\nu ^{\mathrm{eff}}`$ from the standard value (i.e. 3) is small and hence the detection may not be easy.
However, from more general point of view, it is possible that light particles with mass $`1`$ GeV are responsible for the late-time reheating. In this case, as seen in the previous section, the reheating temperatures as low as $`0.7`$ MeV are allowed. For such low reheating temperature, neutrinos cannot be produced sufficiently. Thus the effective number of the neutrino species $`N_\nu ^{\mathrm{eff}}`$ becomes much less than $`3`$. This leads to very interesting effects on the formation of large scale structures and CMB anisotropies, which we discuss in the next section.
## VI Constraints from Large Scale Structure and CMB anisotropy
In this section, we discuss possibility to set constraints on the late-time entropy production from the large scale structure and CMB anisotropies. Hereafter, we only consider flat universe models with cosmological constant which are suggested by recent distant Supernovae (SNe) surveys and measurements of CMB anisotropies .
The late-time entropy production influences formation of the large scale structure and CMB anisotropies since the matter-radiation equality epoch is shifted if the effective number of neutrino species changes. The ratio of neutrino density to black-body photon density is $`\rho _\nu /\rho _\gamma =(7/8)(4/11)^{4/3}N_\nu `$. Therefore the redshift of matter-radiation equality can be written as a function of $`N_\nu `$:
$$1+z_{\mathrm{eq}}=4.03\times 10^4\mathrm{\Omega }_0h^2\left(1+\frac{7}{8}\left(\frac{4}{11}\right)^{4/3}N_\nu \right)^1,$$
(53)
where $`\mathrm{\Omega }_0`$ is the density parameter and $`h`$ is the non-dimensional Hubble constant normalized by $`100\mathrm{k}\mathrm{m}/\mathrm{s}/\mathrm{Mpc}`$.
Let us now discuss distribution of galaxies on large scales. For a quantitative analysis, we define the matter power spectrum in Fourier space as $`P(k)<|\delta _k|^2>`$, where $`\delta _k`$ is the Fourier transform of density fluctuations and $`<>`$ denotes the ensemble average. Hereafter, we assume the Harrison-Zel’dovich power spectrum, which is motivated by the inflation scenario, as an initial shape of the power spectrum, i.e., $`P(k)k`$. As fluctuations evolve in the expanding universe, the shape of the power spectrum is changed. One often introduces the transfer function $`T(k)`$ to describe this modification of the initial power spectrum as $`P(k)=AkT(k)^2`$, where $`A`$ is an arbitrary constant. In case of standard cold dark matter (CDM) dominated models, Bardeen et al. found a fitting formula:
$$T(k)=\frac{\mathrm{ln}(1+2.34q)}{2.34q}[1+3.89q+(16.1q)^2+(5.46q)^3+(6.71q)^4]^{1/4},$$
(54)
where $`q=k/\mathrm{\Omega }_0h^2\mathrm{Mpc}^1`$ when the baryon density is negligible small compared to the total density. It is easy to explain why $`q`$ is parameterized by $`\mathrm{\Omega }_0h^2`$. This is because CDM density fluctuations cannot evolve and stagnate during a radiation dominated era. Only after the matter-radiation equality epoch, fluctuations can evolve. Therefore the CDM power spectrum has a peak which corresponds to the horizon scale of the matter-radiation equality epoch. In fact, the wave number of the horizon scale at the equality epoch can be written as $`k_{\mathrm{eq}}=\sqrt{2\mathrm{\Omega }_0(1+z_{\mathrm{eq}})}H_0`$, where $`H_0`$ is the Hubble constant at present, that is proportional to $`\mathrm{\Omega }_0h^2`$. In the actual observations, distances in between galaxies are measured in the units of $`h^1\mathrm{Mpc}`$. Therefore to fit the observational data by the CDM type power spectrum, we usually introduce so-called shape parameter $`\mathrm{\Gamma }_\mathrm{s}=\mathrm{\Omega }_0h`$. It is known that we can fit the galaxy distribution if $`\mathrm{\Gamma }_\mathrm{s}0.25\pm 0.05`$ which suggests a low density universe. If the late-time entropy production takes place, however, we need to take into account $`N_\nu ^{\mathrm{eff}}`$ dependence of the matter-radiation equality epoch (Eq. (53)). Therefore $`\mathrm{\Gamma }_\mathrm{s}`$ should be written as
$$\mathrm{\Gamma }_\mathrm{s}=1.68\mathrm{\Omega }_0h/(1+0.227N_\nu ^{\mathrm{eff}}).$$
(55)
We plot contours of $`\mathrm{\Gamma }_\mathrm{s}`$ on $`\mathrm{\Omega }_0N_\nu ^{\mathrm{eff}}`$ plane in Fig. 14. It is shown that smaller $`\mathrm{\Omega }_0`$ is preferable for $`N_\nu ^{\mathrm{eff}}<3`$ with the same value of $`\mathrm{\Gamma }_\mathrm{s}`$. We also plot the power spectra for $`\mathrm{\Omega }_0=0.3`$ and $`h=0.7`$ with different $`N_\nu ^{\mathrm{eff}}`$’s in Fig. 15. Here we do not simply employ the fitting formula but numerically solve the evolution of density fluctuations . It is shown that the peak location of a model with smaller $`N_\nu ^{\mathrm{eff}}`$ shifts to the smaller scale (larger in $`k`$) since smaller $`N_\nu ^{\mathrm{eff}}`$ makes the equality epoch earlier which means the horizon scale at the equality epoch becomes smaller. We have hope that current larege scale structrue surveys such as 2DF and Sloan Digital Sky Survey (SDSS) may determine the precise value of $`\mathrm{\Gamma }_\mathrm{s}`$.
Besides the shape of the power spectrum, the amplitude is another important observational quantities to test models. On very large scales, the amplitude of the power spectrum is determined by CMB anisotropies which are measured by COBE/DMR . Since COBE/DMR scales are much larger than the horizon scale of the matter-radiation equality epoch, however, it is not sensitive to the transfer function $`T(k)`$ but the overall amplitude $`A`$. In order to compare the expected amplitude of the power spectrum from each CDM model with large scale structure observations, we employ the specific mass fluctuations within a sphere of a radius of $`8h^1\mathrm{Mpc}`$, i.e., $`\sigma _8`$ which is defined as
$`\sigma _8^2`$ $`=`$ $`<(\delta M/M(R))^2>_{R=8h^1\mathrm{Mpc}}`$ (56)
$`=`$ $`{\displaystyle \frac{1}{2\pi ^2}}{\displaystyle 𝑑kk^2P(k)W(kR)^2}|_{R=8h^1\mathrm{Mpc}},`$ (57)
where $`W(kR)`$ is a window function for which we employ a top hat shape as $`W(kR)3\left(\mathrm{sin}(kR)kR\mathrm{cos}(kR)\right)/(kR)^3`$. Eke et al. obtained the observational value of $`\sigma _8`$ which is deduced from the rich cluster abundance at present as
$$\sigma _8=(0.52\pm 0.04)\mathrm{\Omega }_0^{0.52+0.13\mathrm{\Omega }_0}.$$
(58)
Other estimates of $`\sigma _8`$ are agreed with their result. For CDM models with standard thermal history, the value of $`\sigma _8`$ is a function of $`\mathrm{\Omega }_0`$ and $`h`$. With the late-time reheating, however, $`\sigma _8`$ for fixed $`\mathrm{\Omega }_0`$ and $`h`$ becomes larger. The reason is following. Since we fix $`\mathrm{\Omega }_0`$ and $`h`$, the normalization factor $`A`$ is same regardless of the value of $`N_\nu ^{\mathrm{eff}}`$. As is shown in Fig. 15, the amplitude of the power spectrum on $`8h^1\mathrm{Mpc}`$, i.e., $`\sigma _8`$, is larger for smaller $`N_\nu ^{\mathrm{eff}}`$. In Fig. 16, we show the allowed region on the $`\mathrm{\Omega }_0h`$ plane for $`N_\nu ^{\mathrm{eff}}=0.5,2`$ and $`3`$ for COBE-normalized flat CDM models with the Harrison-Zel’dovich spectrum. The shaded region satisfies the matching condition with the cluster abundance Eq. (58). For fixed $`h`$, models with smaller $`N_\nu ^{\mathrm{eff}}`$ prefer lower $`\mathrm{\Omega }_0`$. Recently, the HST key project on the Extragalactic Distance Scale has reported that $`h=0.71\pm 0.06`$ (1$`\sigma `$) by using various distant indicators . From SNe measurements, $`\mathrm{\Omega }_0=0.28\pm 0.8`$ for flat models (see Fig. 7 of ). CDM models with $`N_\nu ^{\mathrm{eff}}=0.53`$ are still consistent with above value of $`h`$ and $`\mathrm{\Omega }_0`$. However we expect further precise determination of $`\mathrm{\Omega }_0`$, $`h`$ (from distant SNe surveys and measurements of CMB anisotropies) and $`\sigma _8`$ (from 2DF or SDSS) will set a stringent constraint on $`N_\nu ^{\mathrm{eff}}`$ and $`T_R`$ in near future.
Finally we discuss the CMB constraint on $`T_R`$. Let us introduce temperature angular power spectrum $`C_{\mathrm{}}`$ where $`\mathrm{}`$ is the multipole number of the spherical harmonic decomposition. The rms temperature anisotropy of CMB can be written as $`<|\mathrm{\Delta }T/T|^2>=_{\mathrm{}}(2\mathrm{}+1)C_{\mathrm{}}/4\pi `$. Using $`C_{\mathrm{}}`$, we can extract various important information of cosmology, such as the curvature of the universe, $`\mathrm{\Omega }_0`$, cosmological constant, $`h`$ and so on (see, e.g., ). In fact, we can measure the matter radiation equality epoch by using the height of peaks of $`C_{\mathrm{}}`$. The peaks are boosted during the matter-radiation equality epoch. If the matter-radiation equality is earlier, the correspondent horizon scale is smaller. Therefore we expect lower heights for first one or two peaks since these peaks are larger than the horizon scale at the equality epoch and do not suffer the boost as is shown in Fig. 17. With the present angular resolutions and sensitivities of COBE observation or current balloon and ground base experiments, however, it is impossible to set a constraint on $`N_\nu ^{\mathrm{eff}}`$. It is expected that future satellite experiments such as MAP and PLANCK will give us a useful information about $`N_\nu ^{\mathrm{eff}}`$. From Lopez et al.’s analysis , MAP and PLANCK have sensitivities that $`\delta N_\nu ^{\mathrm{eff}}0.1`$ (MAP) and $`0.03`$ (PLANCK) including polarization data, even if all cosmological parameters are determined simultaneously (see also Fig. 17). From such future observations of anisotropies of CMB, it is expected that we can precisely determine $`T_R`$.
## VII Conclusion
In this paper we have investigated the various cosmological effects induced by the late-time entropy production due to the massive particle decay. The neutrino distribution functions have been obtained by solving the Boltzmann equations numerically. We have found that if the large entropy are produced at about $`t1`$ sec, the neutrinos are not thermalized very well and hence do not have the perfect Fermi-Dirac distribution. The deficits of the neutrino distribution functions due to the insufficient thermalization decrease the Hubble expansion rate and weakens the weak interaction rates between proton and neutron. The above two effects changes the freeze-out value of $`n/p`$ significantly. Especially the produced $`{}_{}{}^{4}\mathrm{He}`$ mass fraction $`Y`$ is so sensitive to $`n/p`$ that the predicted value of $`Y`$ is changed drastically. Comparing the theoretical predictions of D, $`{}_{}{}^{4}\mathrm{He}`$ and $`{}_{}{}^{7}\mathrm{Li}`$ to the observational data, we have estimated the lower bound on the reheating temperature $`T_R`$ after the entropy production. We have found that $`T_R0.7`$ MeV is excluded at 95 $`\%`$ C.L. In other wards, $`T_R`$ can be as low as 0.7 MeV. Then the effective number of neutrino species $`N_\nu ^{\mathrm{eff}}`$ can be as small as $`0.1`$. It is enough sensitive for the ongoing large scale structure observations such as 2DF and SDSS or future satellite experiments (MAP and PLANCK) of CMB anisotropies to detect such modifications on $`N_\nu ^{\mathrm{eff}}`$ and we can find out the vestige of the late-time entropy production.
Furthermore, we have also studied the case in which the massive particle has some decay modes into quarks or gluons. In this scenario, a lot of hadrons, e.g. pions, kaons, protons and neutrons, which are originated by the fragmentation of the high energy quarks and gluons are injected into thermal bath. The emitted hadrons extraordinarily inter-convert the ambient protons and neutrons each other through the strong interaction even after the freeze-out time of the neutron to proton ratio $`n/p`$. Then the predicted value of $`Y`$ increase extremely and we can constrain $`T_R`$ and the branching ratio of the hadronic decay mode $`B_h`$ comparing to the observational light element abundances. We have found $`T_R`$ should be higher than 2.5 MeV \- 4 MeV at 95 $`\%`$ C.L. for $`B_h`$ = $`10^2`$ \- 1. The above results tells us that $`N_\nu ^{\mathrm{eff}}`$ can be as small as 1.9 - 2.8 even in the hadron injection scenario for $`B_h=10^2`$ \- 1. Then it still may be possible to detect the modifications on $`N_\nu ^{\mathrm{eff}}`$ by MAP and PLANCK.
## Acknowledgment
K.K. wish to thank J. Yokoyama and T. Asaka for useful discussions. This work was partially supported by the Japanese Grant-in-Aid for Scientific Research from the Monbusho, Nos. 10640250 (MK), 10-04502 (KK), 9440106 (NS) and “Priority Area: Supersymmetry and Unified Theory of Elementary Particles(#707)”(MK) and by the Sumitomo Foundation (NS).
## Reduction of collision integral
This appendix shows how we can reduce the nine dimensional integrals in Eq. (7) of the collision term $`C_{i,\mathrm{scat}}`$ for the scattering process into one dimensional integrals. Notice that, since we treat the massless neutrino, the norm of the neutrino momentum equals to its energy $`|𝒑_𝒊|=E_i`$. Here we divide the collision term into two parts:
$$C_{i,\mathrm{scat}}=F+B,$$
(59)
where $`F`$ represents the forward process and $`B`$ represents the backward process. They are given by
$$F=\frac{g_e}{2E_1}\frac{dp_2^3}{2E_2(2\pi )^3}\frac{dp_3^3}{2E_3(2\pi )^3}\frac{dp_4^3}{2E_4(2\pi )^3}(2\pi )^4\delta ^4(p_1+p_2p_3p_4)S|M|^2\mathrm{\Lambda }_F,$$
(60)
$$B=\frac{g_e}{2E_1}\frac{dp_2^3}{2E_2(2\pi )^3}\frac{dp_3^3}{2E_3(2\pi )^3}\frac{dp_4^3}{2E_4(2\pi )^3}(2\pi )^4\delta ^4(p_1+p_2p_3p_4)S|M|^2\mathrm{\Lambda }_B,$$
(61)
where $`g_e`$ = 2 and the phase space factors are given by
$`\mathrm{\Lambda }_F`$ $`=`$ $`f_1(E_1)f_2(E_2)\left(1f_3(E_3)\right)\left(1f_4(E_4)\right),`$ (62)
$`\mathrm{\Lambda }_B`$ $`=`$ $`\left(1f_1(E_1)\right)\left(1f_2(E_2)\right)f_1(E_3)f_2(E_4).`$ (63)
The integral over $`d^3p_4`$ is immediately done using $`\delta ^3(𝒑_\mathrm{𝟏}\mathbf{+}𝒑_\mathrm{𝟐}\mathbf{}𝒑_\mathrm{𝟑}\mathbf{}𝒑_\mathrm{𝟒})`$. From the momentum conservation, $`\mathbf{|}𝒑_\mathrm{𝟒}\mathbf{|}`$ is given by
$$|𝒑_\mathrm{𝟒}|^2=E_4^2=E_2^2+2E_2R\mathrm{cos}\eta +R^2,$$
(64)
where $`𝑹`$ $``$ $`𝒑_\mathrm{𝟏}\mathbf{}𝒑_\mathrm{𝟑}`$, $`R`$ = $`|𝑹|`$ and $`\mathrm{cos}\eta `$ $``$ $`𝑹\mathbf{}𝒑_\mathrm{𝟐}`$/($`|𝒑_\mathrm{𝟐}|`$R).
The remaining delta function $`\delta (E_1+E_2E_3E_4)`$ shows the energy conservation low which is given by
$$E_4^2=E_1^2+E_2^2+E_3^2+2(E_1E_2E_1E_3E_2E_3).$$
(65)
We can generally take the momentum axes as
$`𝑹`$ $`=`$ $`(0,0,R)`$ (66)
$`𝒑_\mathrm{𝟐}`$ $`=`$ $`(E_2\mathrm{sin}\eta \mathrm{sin}\varphi ,E_2\mathrm{sin}\eta \mathrm{cos}\varphi ,E_2\mathrm{cos}\eta ),`$ (67)
$`𝒑_\mathrm{𝟑}`$ $`=`$ $`(E_3\mathrm{sin}\xi ,0,E_3\mathrm{sin}\xi ),`$ (68)
where
$$\mathrm{cos}\xi =\frac{E_1^2E_3^2R^2}{2E_3R}.$$
(69)
Then $`|\mathrm{cos}\xi |1`$ demands
$$|E_1E_3|RE_1+E_3.$$
(70)
The volume element of $`𝒑_\mathrm{𝟐}`$ is given by $`dp_2^3=E_2^2d\mathrm{cos}\eta d\varphi `$ and from Eqs. (64) and ( 65) the azimuthal angle is obtained by
$`\mathrm{cos}\eta ={\displaystyle \frac{R^2(E_1E_3)^22E_2(E_1E_3)}{2E_2R}}.`$ (71)
Then $`|\mathrm{cos}\eta |1`$ demands
$$|E_1E_3|RE_1+2E_2E_3.$$
(72)
From Eqs. (70) and (72), we can obtain the allowed region of R,
$$|E_1E_3|R\mathrm{Inf}[E_1+E_3,E_1+2E_2E_3].$$
(73)
Since the volume element of $`𝒑_\mathrm{𝟑}`$ is given by $`dp_3^2=2\pi E_3^2dE_3d\mathrm{cos}\theta `$ where $`\mathrm{cos}\theta =𝒑_\mathrm{𝟏}\mathbf{}𝒑_\mathrm{𝟑}/(E_3E_1)`$, the differential angle element is evaluated by
$$d\mathrm{cos}\theta =\frac{R}{2E_1E_3}dR.$$
(74)
From Eq. (73) we can see that the integration can be performed in the four allowed intervals,
$`F+B`$ $`=`$ $`{\displaystyle \frac{1}{128E_1^2}}{\displaystyle _0^{\mathrm{}}}𝑑E_3{\displaystyle _0^{\mathrm{}}}𝑑E_2{\displaystyle 𝑑R_0^{2\pi }\frac{d\varphi }{2\pi }|M|^2(\mathrm{\Lambda }_F+\mathrm{\Lambda }_B)}`$ (75)
$`=`$ $`{\displaystyle \frac{1}{128E_1^2}}[{\displaystyle _0^{E_1}}dE_3{\displaystyle _0^{E_3}}dE_2{\displaystyle _{E_1E_3}^{E_1+2E_2E_3}}dR+{\displaystyle _0^{E_1}}dE_3{\displaystyle _{E_3}^{\mathrm{}}}dE_2{\displaystyle _{E_1E_3}^{E_1+E_3}}dR`$ (78)
$`+{\displaystyle _{E_1}^{\mathrm{}}}dE_3{\displaystyle _{E_1+E_3}^{E_3}}dE_2{\displaystyle _{E_1+E_3}^{E_1+2E_2E_3}}dR+{\displaystyle _{E_1}^{\mathrm{}}}dE_3{\displaystyle _{E_3}^{\mathrm{}}}dE_2{\displaystyle _{E_1+E_3}^{E_1+E_3}}dR]`$
$`\times {\displaystyle _0^{2\pi }}{\displaystyle \frac{d\varphi }{2\pi }}S|M|^2(\mathrm{\Lambda }_F+\mathrm{\Lambda }_B).`$
Even though we only show the case of of $`\nu _e`$ here, we can get the same procedure for $`\nu _\mu `$ and $`\nu _\tau `$ if $`C_V`$ and $`C_A`$ are replaced by $`\stackrel{~}{C_V}`$ and $`\stackrel{~}{C_A}`$. As we also noted in Sec II, we assume that electrons obey the Boltzmann distribution function $`e^{E/T}`$. In addition, since neutrinos are massless, the energy momentum conservation gives $`p_1p_4=p_2p_3`$ in the elastic scattering process. The above assumptions simplify the integrations still more.
For the forward reaction, $`\nu (p_1)+e^\pm (p_2)\nu (p_3)+e^\pm (p_4)`$, the phase space factor is given by
$$\mathrm{\Lambda }_F=f_\nu (E_1)\left(1f_\nu (E_3)\right)\mathrm{exp}[\frac{E_2}{T}].$$
(79)
Then $`F_1`$ and $`F_2`$ in Eq. (16) are analytically estimated as
$`F_1`$ $``$ $`\left[{\displaystyle _0^{E_3}}𝑑E_2{\displaystyle _{E_1E_3}^{E_1+2E_2E_3}}𝑑R+{\displaystyle _{E_3}^{\mathrm{}}}𝑑E_2{\displaystyle _{E_1E_3}^{E_1+E_3}}𝑑R\right]{\displaystyle _0^{2\pi }}{\displaystyle \frac{d\varphi }{2\pi }}{\displaystyle \frac{S|M|^2e^{\frac{E_2}{T}}}{256(C_V^2+C_A^2)G_F^2}}`$ (80)
$`=`$ $`2T^4[E_1^2+E_3^2+2T(E_1E_3)+4T^4]T^2[E_1^2E_3^2+2E_1E_3(E_1+E_3)T`$ (82)
$`+2(E_1+E_3)^2T^2+4(E_1+E_3)T^3+8T^4]e^{\frac{E_3}{T}},`$
$`F_2`$ $``$ $`\left[{\displaystyle _{E_1+E_3}^{E_3}}𝑑E_2{\displaystyle _{E_1+E_3}^{E_1+2E_2E_3}}𝑑R+{\displaystyle _{E_3}^{\mathrm{}}}𝑑E_2{\displaystyle _{E_1+E_3}^{E_1+E_3}}𝑑R\right]{\displaystyle _0^{2\pi }}{\displaystyle \frac{d\varphi }{2\pi }}{\displaystyle \frac{S|M|^2e^{\frac{E_2}{T}}}{256(C_V^2+C_A^2)G_F^2}}`$ (83)
$`=`$ $`2T^4(E_1^2+E_3^22T(E_1E_3)+4T^4)e^{\frac{E_1E_3}{T}}T^2[E_1^2E_3^2+2E_1E_3(E_1+E_3)T`$ (85)
$`+2(E_1+E_3)^2T^2+4(E_1+E_3)T^3+8T^4]e^{\frac{E_3}{T}}.`$
On the other hand, for the backward reaction, $`\nu (p_1)+e^\pm (p_2)\nu (p_3)+e^\pm (p_4)`$, the phase space factor is given by,
$$\mathrm{\Lambda }_B=\left(1f_\nu (E_1)\right)f_\nu (E_3)\mathrm{exp}(\frac{E_1+E_2+E_3}{T}).$$
(86)
Then we can analytically obtain $`B_1`$ and $`B_2`$ in Eq. (16) as
$`B_1`$ $``$ $`\left[{\displaystyle _0^{E_3}}𝑑E_2{\displaystyle _{E_1E_3}^{E_1+2E_2E_3}}𝑑R+{\displaystyle _{E_3}^{\mathrm{}}}𝑑E_2{\displaystyle _{E_1E_3}^{E_1+E_3}}𝑑R\right]{\displaystyle _0^{2\pi }}{\displaystyle \frac{d\varphi }{2\pi }}{\displaystyle \frac{S|M|^2e^{\frac{E_1+E_2+E_3}{T}}}{256(C_V^2+C_A^2)G_F^2}},`$ (87)
$`=`$ $`2T^4(E_1^2+E_3^2+2T(E_1E_3)+4T^4)e^{\frac{E_1E_3}{T}}T^2[E_1^2E_3^2+2E_1E_3(E_1+E_3)T`$ (89)
$`+2(E_1+E_3)^2T^2+4(E_1+E_3)T^3+8T^4]e^{\frac{E_1}{T}},`$
$`B_2`$ $``$ $`\left[{\displaystyle _{E_1+E_3}^{E_3}}𝑑E_2{\displaystyle _{E_1+E_3}^{E_1+2E_2E_3}}𝑑R+{\displaystyle _{E_3}^{\mathrm{}}}𝑑E_2{\displaystyle _{E_1+E_3}^{E_1+E_3}}𝑑R\right]{\displaystyle _0^{2\pi }}{\displaystyle \frac{d\varphi }{2\pi }}{\displaystyle \frac{S|M|^2e^{\frac{E_1+E_2+E_3}{T}}}{256(C_V^2+C_A^2)G_F^2}}`$ (90)
$`=`$ $`2T^4(E_1^2+E_3^22T(E_1E_3)+4T^4)T^2[E_1^2E_3^2+2E_1E_3(E_1+E_3)T`$ (92)
$`+2(E_1+E_3)^2T^2+4(E_1+E_3)T^3+8T^4]e^{\frac{E_1}{T}}.`$ |
warning/0002/quant-ph0002024.html | ar5iv | text | # Entanglement and quantum computation with ions in thermal motion
## I Introduction
Quantum computing relies on the ability to perform a collection of unitary evolutions of a quantum system, composed of a number of two-level systems (the qubits), and it is a key result that a small set of so-called universal gates exists, which may form the basis for the entire computation . The development of proposals for physical implementation of quantum computing have followed different routes, according to the various views one may have on the quantum dynamical processes. (i): one may view a gate operation on a single or on several qubits as a controlled transition from the initial to the final states, and one may implement it by a Hamiltonian, or a sequence of Hamiltonians, that couple these states directly. (ii): one may consider Hamiltonians that couple quite many states, but where unwanted operations are dynamically suppressed by resonance conditions or by ’bang-bang’ Hamiltonians . (iii): one may depart from a more systematic analysis of the Lie algebra generated (by commutation) from a given set of basic Hamiltonians; If one has access to Hamiltonians $`H_1`$ and $`H_2`$ with variable strength parameters $`\kappa _1`$ and $`\kappa _2`$, subsequent application over short time intervals $`dt`$ of $`\kappa _1H_1`$, $`\kappa _2H_2`$, $`\kappa _1H_1`$ and $`\kappa _2H_2`$, leads to the evolution operator ($`\mathrm{}=1`$)
$`\mathrm{e}^{i\kappa _2H_2dt}\mathrm{e}^{i\kappa _1H_1dt}\mathrm{e}^{i\kappa _2H_2dt}\mathrm{e}^{i\kappa _1H_1dt}`$ (1)
$`=\mathrm{e}^{\kappa _1\kappa _2[H_1,H_2]dt^2}+O(dt^3),`$ (2)
so that effectively the Hamiltonian $`i[H_1,H_2]`$ is obtained. As expressed by Lloyd : ‘By going forward and backing up a sufficiently small distance a large enough number of times, it is possible to parallel park in a space only $`\epsilon `$ longer than the length of the car’. If $`H_1`$ and $`H_2`$ commute with the commutator $`[H_1,H_2]`$, the higher order terms in $`dt`$ vanish exactly and one may apply $`H_1`$ and $`H_2`$ for arbitrarily large $`dt`$ and ’make a round trip in the parking lot and park in one single operation’.
The different proposals for quantum computing with trapped ions can be roughly categorized according to the lines above: In their original proposal, Cirac and Zoller noted that lasers resonant with sideband excitation of the trapped ions couple the ground and first vibrational state conditioned on the internal state of the irradiated ion, and subsequent irradiation of a second ion can couple its internal states conditioned on the vibrational state. We have formulated a proposal for two-bit and multi-bit gates in the ion trap, which makes use of resonance conditions to couple certain states of the two-particle system. In our proposal we apply bichromatic light which selects certain virtually excited intermediate states, and by choosing appropriate parameters we show that the desired internal state dynamics of the ions may be perfectly achieved, even if the vibrational degrees of freedom, used to couple the ions, are not in their ground state. Recently, Milburn has proposed a realization of a multi-bit quantum gate in the ion trap, which also operates when the ions are vibrationally excited: Adjusting the phases of laser fields resonant with side band transitions, one may couple internal state operators to different quadrature components, e.g., position and momentum, $`X`$ and $`P`$, of the oscillatory motion. In Ref. it is proposed to use the two Hamiltonians $`H_1=\lambda _1J_zP`$ and $`H_2=\lambda _2J_zX`$, expressed in terms of the collective spin operators $`J_\xi =_kj_{\xi k}`$ $`(\xi =x,y,z)`$, where the sum is over the ions irradiated by the lasers, and where $`j_{\xi k}`$ is the spin operator for the atom $`k`$, which may be defined by the Pauli spin matrices $`j_{\xi k}=\sigma _{\xi k}/2`$ $`(\mathrm{}=1)`$. By alternating application of the Hamiltonians $`H_1`$ and $`H_2`$ we may obtain the exact propagator
$`e^{iH_2\tau }e^{iH_1\tau }e^{iH_2\tau }e^{iH_1\tau }=e^{i\lambda _1\lambda _2J_z^2\tau ^2}`$ (3)
because the commutator of the oscillator position and momentum is a number. The interaction contained in $`J_z^2`$ between the ions has been established via the vibrational degrees of freedom, but after the gate this motion is returned to the initial state and is not in any way entangled with the internal state dynamics. Milburn also considers the possibility of coupling different individual internal state operators successively to X and P, so that the commutator term provides the product of such operators.
In this paper, we shall demonstrate that our bichromatic excitation scheme is in fact already a realization of the proposal by Milburn, and that gate operation more rapid than concluded in is possible. We show that our bichromatic scheme implements a propagator of the form $`\mathrm{e}^{iA(\tau )J_y^2}`$ which is analogous to the one obtained by Milburn (3). In Ref. it was shown that this propagator can be used to prepare maximally entangled states $`\frac{1}{\sqrt{2}}(|gg\mathrm{}g+\mathrm{e}^{i\varphi }|ee\mathrm{}e)`$ of any number of ions ($`N`$), where the $`k`$’th letter denotes the internal state $`e`$ or $`g`$ of the $`k`$’th ion. These maximally entangled states, which have an interest in the their own right , are produced by applying the unitary operator $`\mathrm{e}^{i\pi /2J_y^2}`$ to a string of ions initially in the state $`|gg\mathrm{}g`$, and they may be produced even without experimental access to individual ions in the trap.
In this paper we focus on the preparation of maximally entangled states. This is both of convenience for the theoretical presentation and to emphasize results which are most easily verified experimentally. However, the procedures described here also apply to quantum computation. With two ions illuminated by laser light the bichromatic scheme produces $`\frac{1}{\sqrt{2}}(|ggi|ee)`$ and together with single qubit rotation this evolution forms a universal set of gates which may be used to constuct a quantum computer. The control-not operation for example may be obtained by applying single ion operations on each ion before and after the bichromatic pulse which creates the state $`\frac{1}{\sqrt{2}}(|ggi|ee)`$ from $`|gg`$.
In section II, we recall our proposal for a two-qubit gate operation and we show that it is equivalent to the proposal of Milburn, with a harmonic rather than a stroboscopic application of Hamiltonian coupling terms. In experiments it may be difficult to fulfill the requirements for the analysis of sec. II to be precise, and in section III we address the fidelity with which certain entangled states may be engineered when we take into account the off-resonant couplings and the finite value of the Lamb-Dicke parameter. In sec. IV we study the influence of the environment on the system. We analyse the role of spectator vibrational modes and energy exchange between the ionic motion and thermal surroundings. A summary of our results and a conclusion are presented in section V.
## II Gate operation under ideal conditions
Ions in a linear trap interacting with a laser field of frequency $`\omega `$ may be described by the Hamiltonian
$`H=`$ $`H_0+H_{\mathrm{int}}`$ (4)
$`H_0=`$ $`\nu (a^{}a+1/2)+\omega _{eg}{\displaystyle \underset{i}{}}\sigma _{zi}/2`$ (5)
$`H_{\mathrm{int}}=`$ $`{\displaystyle \underset{i}{}}{\displaystyle \frac{\mathrm{\Omega }_i}{2}}(\sigma _{+i}\mathrm{e}^{i(\eta _i(a+a^{})\omega t)}+h.c.),`$ (6)
where $`\nu `$ is the frequency of the vibration, $`a^{}`$ and $`a`$ are the ladder operators of the quantized oscillator, $`w_{eg}`$ is the energy difference between the internal states $`e`$ and $`g`$, and $`\mathrm{\Omega }_i`$ is the resonant Rabi frequency of the $`i`$’th ion in the laser field. The exponentials account for the position dependence of the laserfield, and the recoil of the ions upon absorption of a photon. The positions of the ions $`x_i`$ are replaced by ladder operators $`kx_i=\eta _i(a+a^{})`$, where the Lamb-Dicke parameter $`\eta _i`$ represents the ratio between the ionic excursions within the vibrational ground state wavefunction and the wavelength of the exciting radiation. In Eq. (6) we have assumed that the laser is close to a sideband $`\omega \omega _{eg}\pm \nu `$ for a single vibrational mode and that we may neglect the contribution from all other vibrational modes. We tune the lasers close to the center-of-mass vibrational mode where all ions participate equally in the vibration, so that the coupling of the recoil to the vibration is identical for all ions, i.e., $`\eta _i=\eta `$ for all $`i`$. For simplicity we also assume the same Rabi frequency for all ions participating in the gate $`\mathrm{\Omega }_i=\mathrm{\Omega }`$. In this section we will consider an ion trap operating in the Lamb-Dicke limit, i.e. the ions are cooled to a regime with vibrational numbers $`n`$ ensuring that $`(n+1)\eta ^2<<1`$, so that we may perform the expansion $`e^{i\eta (a+a^{})}1+i\eta (a+a^{})`$.
### A Weak field coupling
In our previous article we assumed that two ions in the string were both illuminated with two lasers of opposite detunings $`\omega \omega _{eg}=\pm \delta `$. With this choice of laser detunings the only energy conserving transitions are from $`|ggn`$ to $`|een`$ and from $`|gen`$ to $`|egn`$, where $`n`$ is the quantum number for the relevant vibrational mode of the trap, cf. Fig. 1. We considered the weak field regime $`\eta \mathrm{\Omega }<<\nu \delta `$, where only a negligible population is transfered to the intermediate levels with vibrational quantum numbers $`n\pm 1`$. In this regime the effective Rabi frequency $`\stackrel{~}{\mathrm{\Omega }}`$ for the transition from $`|ggn`$ to $`|een`$ may be determined in second order perturbation theory
$$\stackrel{~}{\mathrm{\Omega }}=2\underset{m}{}\frac{een|H_{int}|mm|H_{int}|ggn}{E_m(E_{ggn}+\omega _m)}=\frac{(\mathrm{\Omega }\eta )^2}{\nu \delta },$$
(7)
where we have used the intermediate states $`|m`$=$`|egn\pm 1`$ and $`|gen\pm 1`$, and where $`\omega _m`$ is the frequency of the laser exciting the intermediate state $`|m`$. For the transition $`|egn`$ to $`|gen`$ we get the same effective Rabi frequency.
The remarkable feature in Eq. (7) is that it contains no dependence on the vibrational quantum number $`n`$. This is due to interference between the paths shown in Fig. 1. If we take a path where an intermediate state with vibrational quantum number $`n+1`$ is excited, we have a factor of $`n+1`$ appearing in the numerator ($`\sqrt{n+1}`$ from raising and $`\sqrt{n+1}`$ from lowering the vibrational quantum number). In paths involving the vibrational state $`n1`$ we obtain a factor of $`n`$. Due to the opposite detunings, the denominators have opposite signs and the $`n`$ dependence disappears when the two terms are subtracted. The coherent evolution of the internal atomic state is thus insensitive to the vibrational quantum numbers, and it may be observed with ions in any superposition or mixture of vibrational states. The coherent evolution may even be seen if the vibrational quantum number $`n`$ changes during the gate due to heating .
### B General field coupling
We now consider the interaction without restricting the parameters to a regime where no population is transfered to states with different $`n`$. For this purpose it is convenient to change to the interaction picture with respect to $`H_0`$. In the Lamb-Dicke limit with lasers detuned by $`\pm \delta `$ the Hamiltonian becomes
$`H_{\mathrm{int}}=`$ $`2\mathrm{\Omega }J_x`$ $`\mathrm{cos}\delta t\sqrt{2}\eta \mathrm{\Omega }J_y`$ (10)
$`\times [x(\mathrm{cos}(\nu \delta )t+\mathrm{cos}(\nu +\delta )t)`$
$`+p(\mathrm{sin}(\nu \delta )t+\mathrm{sin}(\nu +\delta )t)],`$
where we have introduced the dimensionless position and momentum operators, $`x=\frac{1}{\sqrt{2}}(a+a^{})`$ and $`p=\frac{i}{\sqrt{2}}(a^{}a)`$, and the collective spin operators discussed above Eq. (3).
Choosing not too strong laser intensities $`\mathrm{\Omega }<<\delta `$ and tuning close to the sidebands $`\nu \delta <<\delta `$ we may neglect the $`J_x`$ term and the terms oscillating at frequency $`\nu +\delta `$ in Eq. (10), and our interaction is a special case of the Hamiltonian
$$H_{\mathrm{int}}=f(t)J_yx+g(t)J_yp.$$
(11)
The exact propagator for the Hamiltonian (11) may be represented by the ansatz
$$U(t)=\mathrm{e}^{iA(t)J_y^2}\mathrm{e}^{iF(t)J_yx}\mathrm{e}^{iG(t)J_yp},$$
(12)
and the Schrödinger equation $`i\frac{d}{dt}U(t)=H_{\mathrm{int}}U(t)`$ then leads to the expressions
$`F(t)=`$ $`{\displaystyle _0^t}f(t^{})𝑑t^{}`$ (13)
$`G(t)=`$ $`{\displaystyle _0^t}g(t^{})𝑑t^{}`$ (14)
$`A(t)=`$ $`{\displaystyle _0^t}F(t^{})g(t^{})𝑑t^{}.`$ (15)
With $`f(t)=\sqrt{2}\eta \mathrm{\Omega }\mathrm{cos}(\nu \delta )t`$ and $`g(t)=\sqrt{2}\eta \mathrm{\Omega }\mathrm{sin}(\nu \delta )t`$ following from (10) we get
$`F(t)`$ $`=`$ $`{\displaystyle \frac{\sqrt{2}\eta \mathrm{\Omega }}{\nu \delta }}\mathrm{sin}\left((\nu \delta )t\right)`$ (16)
$`G(t)`$ $`=`$ $`{\displaystyle \frac{\sqrt{2}\eta \mathrm{\Omega }}{\nu \delta }}\left[1\mathrm{cos}((\nu \delta )t)\right]`$ (17)
$`A(t)`$ $`=`$ $`{\displaystyle \frac{\eta ^2\mathrm{\Omega }^2}{\nu \delta }}\left[t{\displaystyle \frac{1}{2(\nu \delta )}}\mathrm{sin}(2(\nu \delta )t)\right].`$ (18)
In the $`xp`$ phase space the operator $`U`$ performs translations $`(x,p)(x+J_yG(t),pJ_yF(t))`$ entangled with the internal state of the ions.
Apart from a change of basis from $`J_z`$ to $`J_y`$ the interaction considered by Milburn may also be put in this form, with $`f(t)`$ and $`g(t)`$ alternating between zero and non-vanishing constants. Within the present formulation, the trick in Ref. is to use functions $`f(t)`$ and $`g(t)`$ such that $`F(t)`$ and $`G(t)`$ both vanish after a period $`\tau `$. At this instant the vibrational motion is returned to its original state and the propagator reduces to $`U(\tau )=\mathrm{e}^{iA(\tau )J_y^2}`$, i.e., we are left with an internal state evolution which is independent of the external vibrational state. This decoupling is possible because the effective internal state transition is completed in the same amount of time for all vibrational components and because the AC Stark shift of the atomic levels due to the laser fields are independent of the value of $`n`$. In the weak field case these properties are ensured by the interfering coupling amplitude in Fig. 1, see detailed discussion in Ref. . In the general case it follow from the formal structure of Eq. (12). According to (13) the acquired factor $`A(\tau )`$ is equal to the area swept by the line segment between $`(G(t),0)`$ and $`(G(t),F(t))`$, as shown in Fig. 2. If $`(G(t),F(t))`$ forms a closed path, $`A(t)`$ is plus (minus) the enclosed area if the path is traversed in the (counter) clockwise direction. In the proposal by Milburn successive constant Hamiltonians proportional to $`x`$ and $`p`$ are applied and the area enclosed by $`(G(t),F(t))`$ is rectangular. In our proposal the area is a circle of radius $`\sqrt{2}J_y\eta \mathrm{\Omega }/(\nu \delta )`$, as illustrated in Fig. 2.
With the propagator in Eq. (12) we may calculate the time evolution of the system. Suppose that the ions are initially in the internal ground state and an incoherent mixture of vibrational state as described by the density matrix $`\rho ^{tot}=_nP_n|g..gng..gn|`$. The time evolution of the internal state density operator $`\rho =\mathrm{Tr}_n(\rho ^{tot})`$ with any number of ions $`N`$ may be found from $`\rho _{a_1\mathrm{}a_N,b_1..b_N}(t)=_nP_ng..gn|U^{}(t)|b_1..b_Na_1..a_N|U(t)|g..gn`$ ($`a_j,b_j=e`$ or $`g`$), where we have used $`_n|nn|=1`$ to remove one of the summations over vibrational states. Here we list the density matrix elements for the case of two ions $`N=2`$:
$`\rho _{gg,gg}`$ $`=`$ $`{\displaystyle \underset{n}{}}P_n\text{[}{\displaystyle \frac{3}{8}}+{\displaystyle \frac{1}{2}}\mathrm{e}^{\frac{F(t)^2+G(t)^2}{4}}`$ (21)
$`\times L_n\left({\displaystyle \frac{F(t)^2+G(t)^2}{2}}\right)\mathrm{cos}\left(A(t)+{\displaystyle \frac{1}{2}}F(t)G(t)\right)`$
$`+{\displaystyle \frac{1}{8}}\mathrm{e}^{(F(t)^2+G(t)^2)}L_n\left(2(F(t)^2+G(t)^2)\right)\text{]}`$
$`\rho _{ee,ee}`$ $`=`$ $`{\displaystyle \underset{n}{}}P_n\text{[}{\displaystyle \frac{3}{8}}{\displaystyle \frac{1}{2}}\mathrm{e}^{\frac{F(t)^2+G(t)^2}{4}}`$ (24)
$`\times L_n\left({\displaystyle \frac{F(t)^2+G(t)^2}{2}}\right)\mathrm{cos}\left(A(t)+{\displaystyle \frac{1}{2}}F(t)G(t)\right)`$
$`+{\displaystyle \frac{1}{8}}\mathrm{e}^{(F(t)^2+G(t)^2)}L_n(2(F(t)^2+G(t)^2))\text{]}`$
$`\rho _{gg,ee}`$ $`=`$ $`{\displaystyle \underset{n}{}}P_n\text{[}{\displaystyle \frac{1}{8}}(1\mathrm{e}^{(F(t)^2+G(t)^2)}`$ (28)
$`\times L_n(2(F(t)^2+G(t)^2)))`$
$`{\displaystyle \frac{i}{2}}\mathrm{e}^{\frac{F(t)^2+G(t)^2}{4}}L_n\left({\displaystyle \frac{F(t)^2+G(t)^2}{2}}\right)`$
$`\times \mathrm{sin}\left(A(t)+{\displaystyle \frac{1}{2}}F(t)G(t)\right)\text{]},`$
where $`L_n`$ is the $`n`$’th order Laguerre polynomium.
These expressions can be evaluated in different regimes. In the weak field regime, $`\eta \mathrm{\Omega }<<\nu \delta `$, the $`xp`$ phase space trajectory is a very small circle, which is traversed several times. $`F(t)`$ and $`G(t)`$ are negligible for all times, and $`\mathrm{e}^{iF(t)J_yx}\mathrm{e}^{iG(t)J_yp}`$ is approximately unity, such that we have an internal state preparation which is disentangled from the vibrational motion throughout the gate. Since $`A(t)\eta ^2\mathrm{\Omega }^2t/(\nu \delta )`$ if $`(\nu \delta )t>>1`$ the time evolution corresponds to the one obtained from an effective Hamiltonian $`H=\stackrel{~}{\mathrm{\Omega }}J_y^2`$, and Eq. (21) describes simple Rabi oscillations between the states $`|gg`$ and $`|ee`$. This is demonstrated in Fig. 3 (a) which shows the time evolution described by Eq. (21). The curves show sinusoidal Rabi oscillation from $`|gg`$ to $`|ee`$ superimposed by small oscillations due to the weak entanglement with the vibrational motion.
Outside the weak field regime the internal state is strongly entangled with the vibrational motion in the course of the gate. For successful gate operation we have to ensure that we return to the initial vibrational state at the end of the gate by choosing parameters such that $`G(\tau )=F(\tau )=0`$, corresponding to $`(\nu \delta )\tau =K2\pi `$, where $`K`$ is an integer. A maximally entangled state is created if we adjust our parameters so that $`A(\tau )=\pi /2`$. This is achieved if the parameters are chosen in accordance with
$$\frac{\eta \mathrm{\Omega }}{\nu \delta }=\frac{1}{2\sqrt{K}},K=1,2,3,\mathrm{}.$$
(29)
By choosing a low value of $`K`$ such that an entangled state is created after a few rounds in phase space we may perform a faster gate than considered in the weak field case. See Fig. 3 (b), where we have used $`K=2`$, and where a maximally entangled state $`\frac{1}{\sqrt{2}}(|ggi|ee)`$ is created at the time $`\nu t250`$.
By combining the requirement $`(\text{29})`$ with the condition $`(\nu \delta )\tau =K2\pi `$ we may express the time for the state preparation as
$$\tau =\frac{\pi }{\eta \mathrm{\Omega }}\sqrt{K}.$$
(30)
In order to avoid off-resonant excitations of the ions we must require $`\frac{\mathrm{\Omega }^2}{\nu ^2}<<1`$ and $`\eta ^2`$ must be much less than unity to fulfill the Lamb-Dicke approximation (see subsec. III A and III B). For a given trap and/or laser intensity Eq. (30) sets a bound on the speed of the gate. In tabel I we give some numerical examples for the time of the gate for some typical experimental parameters. The control-not operation which plays a central role in quantum computation may be created by a combination of single particle rotations and a bichromatic pulse with the duration described by Eq. (30). The single particle operations may be performed much faster than the two qubit gates, so the time required to perform a control-not operation is also given by (30).
## III Non ideal conditions
In the previous section we used the Lamb-Dicke and the rotating wave approximations to arrive at an exactly solvable model. In this section we perform a more detailed analysis of the validity of the approximations and we estimate the effect of deviations from the ideal situation in an actual experiment. The general procedure in the section, is to change to the interaction picture with respect to the simple Hamiltonian (11) using the exact propagator in Eq. (12), and to treat the small deviations from the ideal situation by perturbation theory. The figure of merit for the performance of the gate is taken to be the fidelity $`F`$ of creation of the maximally entangled $`N`$-particle state $`|\mathrm{\Psi }_{max}=1/\sqrt{2}(|gg..gi|ee..e)`$, which in the ideal case is created at the time when $`A(\tau )=\pi /2`$, if the ions are initially in the $`|gg..g`$ state , i. e.,
$$F=\mathrm{\Psi }_{max}|\rho _{int}(\tau )|\mathrm{\Psi }_{max}.$$
(31)
### A Direct coupling
Going from Eq. (10) to Eq. (11) we neglected a term $`H_d=2\mathrm{\Omega }J_x\mathrm{cos}(\delta t)`$. This term describes direct off resonant coupling of $`g`$ and $`e`$ without changes in the vibrational motion. For high laser power this term has a detrimental effect on the fidelity, which we calculate below.
Changing to the interaction picture, we may find the propagator $`U_I(t)`$ from the Dyson series
$`U_I(t)`$ $`=`$ $`1i{\displaystyle _0^t}𝑑t^{}H_{d,I}(t^{})`$ (33)
$`{\displaystyle _0^t}{\displaystyle _0^t^{}}𝑑t^{}𝑑t^{\prime \prime }H_{d,I}(t^{})H_{d,I}(t^{\prime \prime })+\mathrm{},`$
where the interaction Hamiltonian is given by $`H_{d,I}(t)=U^{}(t)H_d(t)U(t)`$. Since $`H_d(t)`$ is oscillating at a much higher frequency than the propagator $`U(t)`$, we may treat $`U(t)`$ as a constant during the integration and we obtain
$`U_I(t)`$ $`=`$ $`1i{\displaystyle \frac{2\mathrm{\Omega }}{\delta }}\mathrm{sin}(\delta t)U^{}(t)J_xU(t)`$ (35)
$`{\displaystyle \frac{\mathrm{\Omega }^2}{\delta ^2}}(1\mathrm{cos}(2\delta t))U^{}(t)J_x^2U(t)+\mathrm{}.`$
Near the endpoint, $`U(t)\mathrm{e}^{i(\pi /2)J_y^2}`$ and we obtain the fidelity
$$F1\frac{N\mathrm{\Omega }^2}{2\delta ^2}(1\mathrm{cos}(2\delta \tau )),$$
(36)
where $`N`$ is the number of ions participating in the gate. We plot in Fig. 4 the product of the fidelity due to the carrier (36) and the population of the EPR-state $`\frac{1}{\sqrt{2}}(|ggi|ee)`$ expected from the time evolution in Eq. (21). The result agrees well with the result of a numerical integration of the Schrödinger equation with the Hamiltonian (10).
If the duration of the laser pulses can be controlled very accurately in the experiment, so that one fulfills both (29) and $`2\delta \tau =2K^{}\pi `$ the effect of the direct coupling vanishes. If one cannot perform such an accurate control, the net effect of the direct coupling is to reduce the average fidelity by $`\frac{N\mathrm{\Omega }^2}{2\delta ^2}`$ (=0.03 for the parameters of Fig. 4).
### B Lamb-Dicke approximation
In section II we used the Lamb-Dicke approximation $`\mathrm{e}^{i\eta (a+a^{})}1+i\eta (a+a^{})`$ to simplify our calculations. Now we investigate the validity of this approximation.
In the weak field case, we can use the exact matrix elements $`n|\mathrm{e}^{i\eta (a+a^{})}|n+1=i\eta \frac{\mathrm{e}^{\eta ^2/2}}{\sqrt{n+1}}L_n^1(\eta ^2)`$, to obtain the effective Rabi frequency between $`|ggn`$ and $`|een`$
$`\stackrel{~}{\mathrm{\Omega }}_n`$ $`=`$ $`\stackrel{~}{\mathrm{\Omega }}\mathrm{e}^{\eta ^2}\left[{\displaystyle \frac{\left(L_n^1(\eta ^2)\right)^2}{n+1}}{\displaystyle \frac{\left(L_{n1}^1(\eta ^2)\right)^2}{n}}\right]`$ (37)
$``$ $`\stackrel{~}{\mathrm{\Omega }}\left[1\eta ^2(2n+1)+\eta ^4\left({\displaystyle \frac{5}{4}}n^2+{\displaystyle \frac{5}{4}}n+{\displaystyle \frac{1}{2}}\right)\right],`$ (38)
where $`\stackrel{~}{\mathrm{\Omega }}`$ is given by Eq. (7), and where $`L_n^1`$ are the generalized Laguerre polynomials
$$L_n^\alpha (x)=\underset{m=0}{\overset{n}{}}(1)^m\left(\begin{array}{cc}n+\alpha & \\ nm& \end{array}\right)\frac{x^m}{m!}.$$
(39)
The effective Rabi frequency is no longer independent of the vibrational quantum number $`n`$, and the internal state becomes entangled with the vibrational motion, resulting in a non-ideal performance of the gate .
To illustrate the effect of deviations from the Lamb-Dicke approximation, we consider again the production of an EPR-state $`\frac{1}{\sqrt{2}}(|ggi|ee)`$. With an $`n`$-dependent coupling strength the fidelity is
$$F=\frac{1}{2}+\frac{1}{2}\underset{n=0}{\overset{\mathrm{}}{}}P_n\mathrm{sin}(\stackrel{~}{\mathrm{\Omega }}_nt),$$
(40)
where $`P_n`$ is the initial population of the vibrational state $`n`$. We show in Fig. 5 the evolution of the fidelity predicted by Eq. (40) and obtained by a direct integration of the full Hamiltonian in Eq. (6). Due to the deviation from the Lamb-Dicke approximation the effective Rabi frequency is reduced, cf., Eq. (37), and the optimal gate performance is achieved with a duration that is longer than $`\pi /(2\stackrel{~}{\mathrm{\Omega }})`$. The spreading of the values of $`\stackrel{~}{\mathrm{\Omega }}_n`$, causes entanglement with the vibrational motion which reduces the fidelity. With the parameters in Fig. 5 the maximally obtainable fidelity is 0.92 obtained after a pulse of duration $`\tau 1.9/\stackrel{~}{\mathrm{\Omega }}`$.
With more than two ions, the time evolution of the system may be obtained by expanding the initial state $`|gg\mathrm{}g`$ on eigenstates of the $`J_y`$ operator:
$`|gg\mathrm{}g={\displaystyle \frac{(i)^N}{2^{N/2}}}{\displaystyle \underset{k=0}{\overset{N}{}}}(1)^k\sqrt{\left(\begin{array}{cc}N& \\ k& \end{array}\right)}|M_y=N/2k.`$ (43)
In the $`J_y`$ basis the propagator (12) is diagonal and in the weak field regime ($`F(t)`$, $`G(t)0`$) with $`n`$ dependent coupling strengths we get the fidelity
$$F=\underset{n=0}{\overset{\mathrm{}}{}}P_n\left|\frac{1}{2^N}\underset{k=0}{\overset{N}{}}\left(\begin{array}{cc}N& \\ k& \end{array}\right)\mathrm{e}^{i(N/2k)^2(\pi /2\stackrel{~}{\mathrm{\Omega }}_nt)}\right|^2.$$
(44)
In the limit of many ions $`(N>>1)`$ and near the optimum ($`\stackrel{~}{\mathrm{\Omega }}_nt\pi /2`$) we may approximate this expression by assuming that $`k`$ is a continuous variable and by replacing the binominal coefficient by a Gaussian distribution with the same width. In this limit the fidelity becomes
$$F=\underset{n=0}{\overset{\mathrm{}}{}}P_n\frac{1}{\sqrt{1+\frac{N(N1)(\pi /2\stackrel{~}{\mathrm{\Omega }}_nt)^2}{4}}}.$$
(45)
Expanding this expression to lowest order in $`\eta `$ and adjusting the pulse duration to take into account the reduction in the coupling strength we find to lowest order in $`\eta `$
$$F=1\frac{\pi ^2N(N1)}{8}\eta ^4\mathrm{Var}(n)$$
(46)
at the optimum time
$$\tau _{opt}=\frac{\pi }{2\stackrel{~}{\mathrm{\Omega }}}(1+\eta ^2(2\overline{n}+1)),$$
(47)
where $`\overline{n}`$ and Var$`(n)`$ are the mean and variance of the vibrational quantum number.
In Eq. (45) and (46) we have replaced a quantity $`N^2`$ following from the Gaussian approximation to (44) by $`N(N1)`$. With this substitution (45) and (46) describe the fidelity well for all values of $`N`$. With the parameters of Fig. 5, Eq. (46) yields $`F=0.88`$ which is in good agreement with the numerical result in the figure.
The equations (37-46) were derived for weak fields, but they also provide an accurate description of the system outside this regime. To show this we note, that with bichromatic light, $`H_{\mathrm{int}}`$ in Eq. (6) may be written as
$`H_{\mathrm{int}}`$ $`=`$ $`2\mathrm{\Omega }\mathrm{cos}(\delta t)[J_x\mathrm{cos}\left(\eta \sqrt{2}(x\mathrm{cos}(\nu t)+p\mathrm{sin}(\nu t))\right)`$ (49)
$`J_y\mathrm{sin}\left(\eta \sqrt{2}(x\mathrm{cos}(\nu t)+p\mathrm{sin}(\nu t))\right)]`$
in the interaction picture with respect to $`H_0`$. An expansion of the trigonometric functions in this Hamiltonian leads to Eq. (10) which formed the basis of the discussion in section II. The term proportional to $`J_x`$ is suppressed because it is far off resonance. The lowest order contribution of this term was treated in the previous section, and we shall now consider corrections to the $`J_y`$ term which may have significant effects. In the interaction picture with respect to the lowest order Hamiltonian (11), $`x`$ and $`p`$ are changed into $`x+J_yG(t)`$ and $`pJ_yF(t)`$ and to lowest non-vanishing order in $`\eta `$ the interaction picture Hamiltonian is
$`H_3`$ $`=`$ $`\eta ^3J_y{\displaystyle \frac{\sqrt{2}\mathrm{\Omega }}{12}}[\mathrm{cos}((\nu \delta )t)h_1(x,p)`$ (51)
$`+\mathrm{sin}((\nu \delta )t)h_2(x,p)],`$
where
$`h_1(x,p)`$ $`=`$ $`3x^3+xp^2+pxp+p^2x`$ (52)
$`h_2(x,p)`$ $`=`$ $`3p^3+px^2+xpx+x^2p,`$ (53)
and where we have used that $`F(t)`$ and $`G(t)`$ are proportional to $`\eta `$. To calculate the effect of the Hamiltonian in (51) we note that the propagator
$`U_{3,\mathrm{int}}(t)=\mathrm{e}^{\left[i\frac{\mathrm{sin}((\nu \delta )t)}{\nu \delta }h_1(x,p)\right]}\mathrm{e}^{\left[i\frac{1\mathrm{cos}((\nu \delta )t)}{\nu \delta }h_2(x,p)\right]}`$ (54)
is consistent with the Hamiltonian (51) until order $`\eta ^5`$, i.e. $`i\frac{dU_{3,\mathrm{int}}(t)}{dt}=(H_3+O(\eta ^6))U_{3,\mathrm{int}}`$. (But the full Hamiltonian contains terms of order $`\eta ^4`$ and $`\eta ^5`$ which are not taken into account in $`U_{3,\mathrm{int}}`$. These terms are included below). We are interested in the propagator at times $`\tau =K2\pi /(\nu \delta )`$ where the vibrational motion is returned to the initial state. At these instants the exponents in Eq. (54) vanish and the propagator reduces to $`U_3(\tau )=1`$ such that it has no influence on the internal state preparation.
Expanding the Hamiltonian to order $`\eta ^6`$ we obtain the propagator to the same order in $`\eta `$ in the interaction picture with respect to $`H_0`$ in (6)
$`U_6(\tau )`$ $`=`$ $`\mathrm{e}^{i\stackrel{~}{\mathrm{\Omega }}\tau J_y^2\left[1\eta ^2(2n+1)+\eta ^4\left(\frac{5}{4}n^2+\frac{5}{4}n+\frac{1}{2}\right)\right]}`$ (56)
$`\times \mathrm{e}^{i\eta ^5J_y^3\frac{\sqrt{8}\mathrm{\Omega }^3}{(\nu \delta )^2}x\tau }\mathrm{e}^{i\eta ^6J_y^4\frac{5\mathrm{\Omega }^4}{2(\nu \delta )^3}\tau }`$
valid at times $`\tau =K2\pi /(\nu \delta )`$. The first exponential provides the time evolution with the modified effective Rabi frequency in Eq. (37). If we evaluate the propagator (56) in the weak field regime, the last two exponentials both vanish in the limit of large K when the requirement (29) is inserted, and the time evolution in (56) is consistent with Eq. (37-46). The last two exponentials are also of minor importance for a different reason: In Eq. (37) $`\eta ^2`$ appears in the combination $`\eta ^2n`$, whereas it appears as $`\eta ^2`$ in the last two exponentials of (56) when the condition (29) is inserted. In situations where deviations from the Lamb-Dicke approximation are important $`\eta ^2n1`$, the deviation is typically caused by a high value of $`n`$ rather than a large value of $`\eta `$ ($`\eta ^2<<1`$). In this case one may neglect the last two exponentials and the effect of the non-Lamb-Dicke terms are the same as in the case of weak fields as described by Eqs. (37-46). To achieve the optimum operation of the gate with the parameters of Fig. 5 we have to ensure $`\stackrel{~}{\mathrm{\Omega }}\tau 1.9`$ and there is a small correction to the condition in Eq. (29).
## IV External disturbances
So far we have considered a system described by the Hamiltonian (6), where only the center of mass motion is present in the ion trap and where the coupling of this mode to the surroundings is neglible. In this section we shall remove these two assumptions and consider the decrease in fidelity due to the presence of other modes in the trap and due to heating of the center of mass vibrational motion.
### A Spectator vibrational modes
With $`N`$ ions in the trap, the motional state is described by $`3N`$ non degenerate vibrational modes. With a proper laser geometry or if the transverse potential is much steeper than the longitudinal potential, the coupling of the laser to transverse modes will be neglible and the only contribution is from the $`N`$ longitudinal modes. With $`N`$ vibrational modes the ion trap may be described by the Hamiltonian
$`H=`$ $`H_0+H_{\mathrm{int}}`$ (57)
$`H_0=`$ $`{\displaystyle \underset{l=1}{\overset{N}{}}}\nu _l(a_l^{}a_l+1/2)+\omega _{eg}{\displaystyle \underset{i}{}}\sigma _{zi}/2`$ (58)
$`H_{\mathrm{int}}=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \frac{\mathrm{\Omega }_i}{2}}(\sigma _{+i}\mathrm{e}^{i(_{l=1}^N\eta _{i,l}(a_l+a_l^{})\omega t)}+h.c.),`$ (59)
where $`\nu _l`$ and $`a_l^{}`$ and $`a_l`$ are the frequency and ladder operators of the $`l`$’th mode. The excursion of the $`i`$’th ion in the $`l^{}th`$ mode is described by the Lamb-Dicke parameter $`\eta _{i,l}`$ which may be represented as $`\eta _{i,l}=\eta \frac{\sqrt{N}b_i^l}{\sqrt{\nu _l/\nu }}`$, where $`\eta `$ and $`\nu `$ refer to the center of mass mode as in the previous sections, and where $`b_i^l`$ obeys the orthogonality conditions $`_{i=1}^Nb_i^lb_i^l^{}=\delta _{l,l^{}}`$ and $`_{l=1}^Nb_i^lb_i^{}^l=\delta _{i,i^{}}`$ .
The center of mass mode ($`l=1`$), which is used to create the entangled states of the ions, has $`b_i^1=1/\sqrt{N}`$ for all ions and is well isolated from the remaning $`N1`$ vibrational modes $`\nu _{l>1}\sqrt{3}\nu `$, so that we could neglect the contribution from the other modes in the previous sections. In this section we shall extimate the effect of the presence of the spectator modes. They have both a direct effect, due to the off resonant coupling to the other modes, and an indirect ’Debye-Waller’ effect because the coupling strength of the center of mass mode is reduced due to the oscilations in the spectator modes. Below we shall calculate the direct and indirect effects separately.
The lowest order contribution of the direct coupling to the spectator modes may be found by expanding the exponentials as in Eq. (10).
$`H_{\mathrm{int}}=`$ $`2\mathrm{\Omega }J_x`$ $`\mathrm{cos}\delta t+{\displaystyle \underset{l=1}{\overset{N}{}}}\mathrm{\Theta }_l[x_lf_l(t)+p_lg_l(t)],`$ (60)
where $`f_l(t)=\sqrt{2}\eta \mathrm{\Omega }\sqrt{\nu /\nu _l}[\mathrm{cos}(\nu _l\delta )t+\mathrm{cos}(\nu _l+\delta )t]`$ and $`g_l(t)=\sqrt{2}\eta \mathrm{\Omega }\sqrt{\nu /\nu _l}[\mathrm{sin}(\nu _l\delta )t+\mathrm{sin}(\nu _l+\delta )t]`$, and where the internal and external state operators are defined by $`\mathrm{\Theta }_l=_{i=0}^Nb_i^lj_{y,i}`$ and $`x_l=\frac{1}{\sqrt{2}}(a_l+a_l^{})`$ and $`p_l=\frac{i}{\sqrt{2}}(a_l^{}a_l)`$. Since the ladder operators for different modes commute, we may find the propagator for this Hamiltonian using the steps that lead to Eq. (12)
$$U(t)=\underset{l=1}{\overset{N}{}}U_l(t),$$
(61)
where
$$U_l(t)=\mathrm{e}^{iA_l(t)\mathrm{\Theta }_l^2}\mathrm{e}^{iF_l(t)\mathrm{\Theta }_lx_l}\mathrm{e}^{iG_l(t)\mathrm{\Theta }_lp_l}$$
(62)
with the functions $`F_l`$, $`G_l`$ and $`A_l`$ defined analogously to Eq. (13). Note, that this is an exact solution of the Hamiltonian (60) without the $`J_x`$ term, so that to lowest order in the Lamb-Dicke parameter it includes all effects of the coupling to the other modes.
From the definition of $`\mathrm{\Theta }_l`$ it is seen that $`\mathrm{\Theta }_1=J_y`$ and the propagator $`U_1`$ reduces to Eq. (12) in the rotating wave approximation. The other $`N1`$ propagators in (61) cause a reduction of the fidelity due to the excursion into the $`x_lp_l`$ phase space of these modes. Expanding the exponentials, using $`gg\mathrm{}g|\mathrm{\Theta }_l\mathrm{\Theta }_l^{}|gg\mathrm{}g=\delta _{l,l^{}}N/4`$ and $`\delta \nu `$, and averaging over time we find
$$F=1\eta ^2N\frac{\mathrm{\Omega }^2}{\nu ^2}\underset{l=2}{\overset{N}{}}\frac{\nu }{\nu _l}(2\overline{n}_l+1)\frac{\nu _l^2/\nu ^2+1}{(\nu _l^2/\nu ^21)^2},$$
(63)
where $`\overline{n}_l`$ is the mean vibrational excitation of the $`l^{}th`$ mode.
In addition to the direct coupling to the spectator vibrational mode, the fidelity is also reduced because the coupling strength is dependent on the vibration of the other modes. Unlike the direct coupling discussed above, this effect is not suppressed by the other modes being far off-resonant, and it may have an effect comparable to the direct coupling.
Due to the vibration of the ions the coupling of the $`i`$’th ion to the sideband is reduced from $`i\eta \sqrt{n+1}`$ to $`n_1n_2\mathrm{}n_N|\mathrm{e}^{i_{l=1}^N\eta _{i,l}(a_l+a_l^{})}|n_1+1n_2\mathrm{}n_Ni\eta \sqrt{n+1}(1_{l=1}^N\eta _{i,l}^2(n_l+1/2))`$. With this reduced coupling strength the effective propagator at times $`\tau =K2\pi /(\nu \delta )`$ may be described by
$$U(\tau )=\mathrm{e}^{iA(\tau )\mathrm{\Lambda }^2},$$
(64)
where $`\mathrm{\Lambda }=_{i=1}^Nj_{y,i}(1_{l=1}^N\eta _{i,l}^2(n_l+1/2))`$. In the Cirac-Zoller scheme , the $`n`$-dependent AC Stark shifts caused by coupling to other vibrational modes lead to decoherence, unless these modes are cooled to the ground state. In our bichromatic scheme, these internal state level shifts depend much less on the vibrational excitation. By expanding (64) around the optimum $`A(t)\pi /2`$ we calculate the lowest order reduction in the fidelity
$`F`$ $`=`$ $`1{\displaystyle \frac{\pi ^2N(N1)}{8}}\eta ^4{\displaystyle \underset{l=1}{\overset{N}{}}}{\displaystyle \frac{\mathrm{Var}(n_l)}{(\nu _l/\nu )^2}}`$ (66)
$`{\displaystyle \frac{\pi ^2(N2)}{16}}\eta ^4{\displaystyle \underset{i,l,l^{}=1}{\overset{N}{}}}{\displaystyle \frac{(b_i^l)^2(b_i^l^{})^21/N^2}{\nu _l\nu _l^{}/\nu ^2}}\overline{n_ln_l^{}}.`$
The expressions in Eqs. (63) and (66) may be simplified if the vibrational motion is in a thermal equilibrium at a given temperature. In a thermal state Var$`(n_l)=\overline{n}_l^2+\overline{n}_l`$, $`\overline{n_ln_l^{}}=\overline{n}_l\overline{n}_l^{}`$ for $`ll^{}`$, and $`\overline{n}_l\overline{n}_1\nu /\nu _l`$, and using these expressions we find the lower estimate for the fidelity
$$F1\eta ^2N\frac{\mathrm{\Omega }^2}{\nu ^2}(\overline{n}_1\sigma _1(N)+\sigma _2(N))$$
(67)
for the direct coupling (63) and
$`F`$ $``$ $`1{\displaystyle \frac{\pi ^2N(N1)}{8}}\eta ^4(\overline{n}_1^2\sigma _3(N)+\overline{n}_1\sigma _4(N))`$ (69)
$`{\displaystyle \frac{\pi ^2(N2)}{16}}\eta ^4(\overline{n}_1^2\sigma _5(N)+\overline{n}_1\sigma _6(N))`$
for the Debye-waller coupling (66), where the sums $`\sigma _1\mathrm{}\sigma _6`$ may be derived from Eqs. (63) and (66). For example $`\sigma _3(N)=_{l=1}^N\frac{\nu ^4}{\nu _l^4}`$. With the mode functions and frequencies of Ref. these sums are readily evaluated, and the results are shown in Fig. 6. From the figure it is seen that $`\sigma _5,\sigma _6<<\sigma _3,\sigma _4`$, so that the last term in Eq. (69) may be neglected. All the sums have a very rapid convergence and we may estimate the fidelity by replacing the sums with their large $`N`$ values, i.e.
$$F1\eta ^2N\frac{\mathrm{\Omega }^2}{\nu ^2}0.8(\overline{n}_1+1)$$
(70)
for the direct coupling (63) and
$`F`$ $``$ $`1{\displaystyle \frac{\pi ^2N(N1)}{8}}\eta ^4(1.2\overline{n}_1^2+1.4\overline{n}_1)`$ (71)
for the Debye-Waller coupling (66).
We note that Eq. (71) is derived from terms beyond the Lamb-Dicke expansion and it incorporates the reduction of fidelity due to deviations from the Lamb-Dicke approximation in the center of mass mode, cf. the formal similarity of Eq. (71) and Eq. (46).
### B Heating of the vibrational motion
An ion trap cannot be perfectly isolated and the vibration of the ions will be subject to heating due to the interaction with the environment. Relaxation due to the interaction between the vibration and a thermal reservoir may be described by the master equation
$$\frac{d}{dt}\rho =i[H,\rho ]+(\rho ),$$
(72)
where $`(\rho )`$ is of the Lindblad form
$$(\rho )=\frac{1}{2}\underset{m}{}\left(C_m^{}C_m\rho +\rho C_m^{}C_m\right)+\underset{m}{}C_m\rho C_m^{}$$
(73)
with relaxation operators $`C_1=\sqrt{\mathrm{\Gamma }(1+n_{th})}a`$ and $`C_2=\sqrt{\mathrm{\Gamma }(n_{th})}a^{}`$, where $`\mathrm{\Gamma }`$ characaterizes the strength of the interaction, and $`n_{th}`$ is the mean vibrational number in thermal equilibrium.
We calculate the effect of heating assuming that the ions remain in the Lamb-Dicke limit. Changing to the interaction picture with respect to the Hamiltonian (11), the time evolution of $`\rho `$ is entirely due to the heating, i.e., the Lindblad terms which are transformed using the propagator (12)
$`\stackrel{~}{C_1}`$ $`=`$ $`U^{}C_1U=\sqrt{\mathrm{\Gamma }(1+n_{th})}\left(a+J_y{\displaystyle \frac{G(t)iF(t)}{\sqrt{2}}}\right)`$ (74)
$`\stackrel{~}{C_2}`$ $`=`$ $`U^{}C_2U=\sqrt{\mathrm{\Gamma }n_{th}}\left(a^{}+J_y{\displaystyle \frac{G(t)+iF(t)}{\sqrt{2}}}\right).`$ (75)
The density matrix is most conveniently expressed in the basis of $`J_y`$ eigenstates, and by tracing over the vibrational states we find the time derivative of the internal state density matrix in the interaction picture
$`{\displaystyle \frac{d}{dt}}\rho _{M_y,M_y^{}}`$ $`=`$ $`(M_yM_y^{})^2\mathrm{\Gamma }(1+2n_{th})`$ (77)
$`\times {\displaystyle \frac{G(t)^2+F(t)^2}{4}}\rho _{M_y,M_y^{}}.`$
This equation is readily integrated, and at times $`\tau =K2\pi /(\nu \delta )`$ we get
$$\rho _{M_y,M_y^{}}(\tau )=\rho _{M_y,M_y^{}}(0)\mathrm{e}^{(M_yM_y^{})^2\frac{\mathrm{\Gamma }(1+2n_{th})}{4K}\tau }.$$
(78)
The initial state is expanded on the $`J_y`$ eigenstates as in Eq. (43) and the population of the initial state (which is ideally constant in the interaction picture) equals
$$F=\frac{1}{2^{2N}}\underset{j=0}{\overset{N}{}}\underset{k=0}{\overset{N}{}}\left(\begin{array}{cc}N& \\ j& \end{array}\right)\left(\begin{array}{cc}N& \\ k& \end{array}\right)\mathrm{e}^{(jk)^2\frac{\mathrm{\Gamma }(1+2n_{th})}{4K}\tau }.$$
(79)
For two ions this expressions can be readily evaluated
$$F=\frac{3}{8}+\frac{1}{2}\mathrm{e}^{\frac{\mathrm{\Gamma }(1+2n_{th})}{4K}\tau }+\frac{1}{8}\mathrm{e}^{4\frac{\mathrm{\Gamma }(1+2n_{th})}{4K}\tau }.$$
(80)
In the limit of many ions ($`N>>1`$) and short times ($`\frac{\mathrm{\Gamma }(1+2n_{th})}{4K}\tau <<1`$) we may again approximate the expression in Eq. (79) by assuming that $`j`$ and $`k`$ are continuous variables and by replacing the binomial coefficients by Gaussian distributions with the same width. In this limit the fidelity becomes
$$F=\frac{1}{\sqrt{1+N\frac{\mathrm{\Gamma }(1+2n_{th})}{4K}\tau }}.$$
(81)
For 2 ions the deviation between (80) and (81) is less than 0.02 for all values of $`F`$ larger than $`0.5`$.
In the above expressions we have assumed the Lamb-Dicke approximation. This corresponds to a situation where the heating is counteracted for example by lasercooling on some ions reserved for this purpose. If the ions are not cooled the heating will proceed towards high vibrational numbers with a heating rate $`\mathrm{\Gamma }n_{th}`$ and the heating will eventually take the ions out of the Lamb-Dicke limit. With strong fields ($`K1`$) the reduction in the fidelity described by Eq. (81) will ruin the entangled state before the heating has made a significant change to the vibrational state ($`\mathrm{\Gamma }n_{th}\tau 1`$). For weak fields ($`K>>1`$) however, the situation is different. With weak fields one may produce an entangled state even though the time required to entangle the ions is much longer than the decoherence time of the vibrational motion which is used to communicate between the ions, i.e. if $`K>N\mathrm{\Gamma }n_{th}\tau `$ the effect of heating is small even though the change in the average vibrational number $`\mathrm{\Gamma }n_{th}\tau `$ is larger than unity . Since the effective Rabi-frequency has a small dependence on the vibrational quantum number $`n`$ as described in Eq. (37) the heating will have an indirect effect on the internal state preparation. This can be modelled by changing the probabilities in Eqs. (40-45) into time dependent functions $`P_n(t)`$ reflecting the change in the vibrational motion occurring during the internal state preparation.
## V Conclusion
We have in this paper evaluated the possibility for preparation of entangled states of ions by illumination with bichromatic light. We have identified two regimes: (i) a weak field regime where single photon absorption is suppressed and where two-photon processes interfere in a way that makes the internal state dynamics insensitive to the vibrational state, and (ii) a strong field regime where the individual ions are coherently excited and the motional state is highly entangled with the internal state until all undesirable excitations are deterministically removed towards the end of the interaction.
We have presented analytical estimates for the fidelity of the internal state preparation. These expressions are summarized in table II. The expressions for the fidelity may be readily applied to experimental parameters and they show that several ion trap experiments today are in a position to apply our proposal directly. In fact, using our proposal the NIST group at Boulder has been able to produce the maximally entangled state $`\frac{1}{\sqrt{2}}(|ggggi|eeee)`$ with four ions . In this experiment the heating of the center-of-mass mode was so strong that this mode could not be used to communicate between the ions. Instead the experiment used an asymmetric mode where all ions have the same amplitude but a different sign, i.e. $`|\eta _i|`$ are the same for all ions $`i`$. Apart from the center-of-mode such modes only exist in ion traps containing two or four ions, and the experiment could not go beyound four ions. In other existing traps the heating is much less significant , and these traps may be employed to produce entangled states with more particles.
The use of ancillary degrees of fredom (center-of-mass position and momentum) to communicate between two or more quantum systems is a key ingredient of quantum information processing. The algebraic property (3) which allows coupling and temporary entanglement with such an ancilla may find wide applications in many different systems for quantum computation with different ancillae (photons, phonons, Cooper-pairs, etc.). However, operators with a constant non-vanishing commutator (which allows the formal step from Eq. (2) to Eq. (3)) only exist in infinite-dimensional Hilbert spaces . In addition to the implementation in cavity QED realizations of quantum computing where quantized cavity fields play the role of the vibrational modes, it thus seems very relevant to investigate to which extent the ideas underlying Eq. (3) can be generalized to ancillae with a finite number of states and, e.g., for communication across a linear qubit register by only nearest neighbour interaction.
## Acknowledgments
We thank B. E. King, C. Monroe, D. J. Wineland, R. Blatt, D. Leibfried, and F. Schmidt-Kahler for fruitfull discussions and for enlightening us on details of their trapping experiments. We also thank D. F. V. James for providing the eigenfrequencies and modes for the collective vibrations which were used to produce Fig. 6. This work is supported by Thomas B. Thriges Center for Kvanteinformatik and by the Danish National Research Council. |
warning/0002/math-ph0002032.html | ar5iv | text | # Freiburg THEP-99/14 To appear in Rep. on Math. Phys. 47, 2001 A vertical exterior derivative in multisymplectic geometry and a graded Poisson bracket for nontrivial geometries
## 1 Introduction
In a geometrical framework to handle field theories over manifolds in a finite dimensional geometry is proposed. This mathematical setting appears under the name multisymplectic geometry, De Donder-Weyl theory, Hamilton-Cartan formalism, and covariant field theory in the literature (, further ). The basic idea is to treat the space coordinates of a given field theory as additional evolution parameters. Thus, there is a finite number of variables (the field and its first derivatives) that evolve in space-time rather than a curve in an infinite-dimensional vector space of field configurations. As shown in one can incorporate the field equations and the Noether theorem in that formulation, but in order to find a corresponding quantum field theory – at least in the sense of a formal deformation – one has to formulate the dynamics of the classical theory in terms of Poisson brackets first.
Kanatchikov () has constructed such a bracket for trivial vector bundles over orientable manifolds. In the nontrivial case the used “vertical exterior derivative” which plays a central rôle in the construction is not globally defined (the resulting bracket, however, does not depend on the coordinate system used). What is needed is a derivative in vertical directions that in particular has square zero. A first guess would be to use a connection and take an expression like $`dv^A_A`$ with $``$ being a covariant derivative and $`dv^A`$ being vertical. The condition that its square gives zero is then equivalent to the flatness of $``$ along fibres. As the fibres under consideration are vector spaces one would indeed expect that it is possible to construct such a covariant derivative. This construction constitutes the main part of this paper.
The remaining part of this article is organised as follows. In the first section a short overview over the multisymplectic approach is given. Then, with the help of a covariant derivative that is flat along the fibres of phase space, the already mentioned vertical exterior derivative is constructed and discussed. Then the Poisson structure is given and the defining properties are proved. Finally, mechanics as the case of a trivial (vector) bundle over a one-dimensional base manifold (i.e., the time axis $``$) is recovered and the scalar field case is considered.
The appendix contains some well known facts about connections viewed as sections of jet bundles and the construction of the already mentioned covariant derivative on the multisymplectic phase space.
## 2 From variational principles to multisymplectic geometry
In field theory, solutions of the field equations are stationary points of the action functional
$$L[\phi ]=_{}(\phi (x),\phi (x))d^{n+1}x,$$
where $``$ is some $`(n+1)`$-dimensional parameter space (e.g. space-time), $`\phi `$ is the gradient of the field $`\phi `$ and $``$ is the Lagrange density.
In general, $`\phi `$ is a section of a vector bundle $`\pi :𝒱`$. $`T\phi :TT𝒱`$ fulfils $`T\pi T\phi =T\text{id}_{}`$ and thus<sup>1)</sup><sup>1)</sup>1)Usually, the first jet bundle of a vector bundle $`(,\pi ,𝒱)`$ is defined to be the set of all equivalence classes at a point of $``$ of local sections, where equivalence means equal function value and first derivatives. But this can be viewed as a tangent map from $`T`$ to $`T𝒱`$ having the stated property. Further, such a tangent map defines how to (horizontally) lift $`T`$ at every point of $`𝒱`$, which is equivalent to having a connection. Hence, a connection defines a map $`𝒱𝒥^1𝒱`$, which turns the affine bundle $`𝒥^1𝒱`$ into a vector bundle over $`𝒱`$. defines an element of $`𝔍^1𝒱`$, the first jet bundle of $`𝒱`$ (). Using a linear connection
$$\mathrm{\Gamma }:𝒱𝔍^1𝒱$$
we obtain an isomorphism
$$i_\mathrm{\Gamma }:\left(𝔍^1𝒱\right)_v\left(𝒱T^{}\right)_{\pi (v)}$$
for all $`v`$ in $`𝒱`$, where in addition we have used $`\left(𝔙𝒱\right)_v𝒱_{\pi (v)}`$ for vector bundles $`𝒱`$ and their vertical tangent bundles $`𝔙𝒱`$. In particular, we find
$$i_\mathrm{\Gamma }T_x\phi \xi (x)=_\xi \phi (x),$$
(1)
for $``$ denoting the covariant derivative corresponding to $`\mathrm{\Gamma }`$ and $`\xi `$ being a tangent vector on $``$. This will be needed in section 4.
Now the Lagrange density can be interpreted as a mapping
$$:𝔍^1𝒱\mathrm{\Lambda }^{n+1}T^{},L[\phi ]=_{}j^1\phi ,$$
where $`j^1\phi (x)=T_x\varphi \left(𝔍^1𝒱\right)_{\phi (x)}`$ is the first jet prolongation of $`\phi \mathrm{\Gamma }(𝒱)`$. Stationary points of $`L`$ correspond to solutions of the Euler-Lagrange equations, which in local coordinates<sup>2)</sup><sup>2)</sup>2)When working in local coordinates of $`𝔍^1𝒱`$ we will use the following convention. Small Latin indices sum over the base manifold directions, that is $`i,j,k`$ run from $`1`$ to $`n+1`$ if not specified otherwise. Capital Latin characters as $`A,B,C,D`$ run from $`1`$ to $`N`$ which is the dimension of a fibre of $`𝒱`$. $`(x^i,v^A,v_i^A)`$ of $`𝔍^1𝒱`$ read (cf. )
$$\frac{}{v^A}j^1\phi \frac{}{x^i}\left(\frac{}{v_i^A}j^1\phi \right)=0.$$
(2)
Now we want to formulate the theory on what we shall call phase space. Since $`𝔍^1𝒱`$ is not a vector bundle but an affine bundle, one chooses the dual $`\left(𝔍^1𝒱\right)^{}`$ to be the bundle of affine mappings from $`𝔍^1𝒱`$ to $`\mathrm{\Lambda }^{n+1}T^{}`$. Thus, coordinates $`(x^i,v^A,v_i^A)`$ on $`𝔍^1𝒱`$ induce coordinates $`(x^i,v^A,p,p_A^i)`$ on $`\left(𝔍^1𝒱\right)^{}`$. One can show (see , ch. 2B) that $`\left(𝔍^1𝒱\right)^{}`$, being a vector bundle over $`𝒱`$ (it inherits a vector space structure from the target space $`\mathrm{\Lambda }^{n+1}T^{}`$), is canonically isomorphic to $`𝒵\mathrm{\Lambda }^{n+1}T^{}𝒱`$, where
$$𝒵_v=\{z\mathrm{\Lambda }^{n+1}T^{}𝒱_v|i_Vi_Wz=0V,W\left(𝔙𝒱\right)_v\},𝒵=\underset{v𝒱}{}𝒵_v.$$
Furthermore, on $`\mathrm{\Lambda }^{n+1}T^{}𝒱`$ there is a canonical $`(n+1)`$-form $`\mathrm{\Theta }_\mathrm{\Lambda }`$, defined by
$$\mathrm{\Theta }_\mathrm{\Lambda }(z)(u_1,\mathrm{},u_{n+1})=z(T\pi _{𝒱\mathrm{\Lambda }}u_1,\mathrm{},T\pi _{𝒱\mathrm{\Lambda }}u_{n+1}),$$
where $`z\mathrm{\Lambda }^{n+1}T^{}𝒱`$, $`u_1,\mathrm{},u_{n+1}T_z\mathrm{\Lambda }^{n+1}T^{}𝒱`$, $`\pi _{𝒱\mathrm{\Lambda }}:\mathrm{\Lambda }^{n+1}T^{}𝒱𝒱`$. Using the embedding $`i_{\mathrm{\Lambda }𝒵}:𝒵\mathrm{\Lambda }^{n+1}T^{}𝒱`$, we obtain an $`(n+1)`$-form on $`𝒵`$,
$$\mathrm{\Theta }=i_{\mathrm{\Lambda }𝒵}^{}\mathrm{\Theta }_\mathrm{\Lambda },$$
(3)
which will be called canonical $`(n+1)`$-form thereafter. There is a canonical $`(n+2)`$-form $`\mathrm{\Omega }`$ on $`𝒵`$, too,
$$\mathrm{\Omega }=d\mathrm{\Theta }.$$
Using coordinates $`(x^i,v^A,p,p_A^i)`$, one finds
$$\mathrm{\Theta }=p_A^idv^A\left(_{x^i}\text{}\text{ }\text{ }d^{n+1}x\right)+pd^{n+1}x,\mathrm{\Omega }=dv^Adp_A^i\left(_{x^i}\text{}\text{ }\text{ }d^{n+1}x\right)dpd^{n+1}x,$$
where $`d^{n+1}x=dx^1\mathrm{}dx^{n+1}`$. Now we are in the position to reformulate (2). As a first step we define a covariant Legendre transform for $``$:
$$\begin{array}{cc}\hfill 𝔽:𝔍^1𝒱\gamma & 𝔽(\gamma )\left(𝔍^1𝒱\right)^{}𝒵,\hfill \\ \hfill 𝔽(\gamma ):𝔍^1𝒱\gamma ^{}& (\gamma )+\frac{d}{dϵ}\text{ }ϵ=0\left(\gamma +ϵ(\gamma ^{}\gamma )\right)\mathrm{\Lambda }^{n+1}T^{}.\hfill \end{array}$$
(4)
In coordinates as above it takes the form
$$=L(x^i,v^A,v_i^A)d^{n+1}x,p_A^i=\frac{L}{v_i^A},p=L\frac{L}{v_i^A}v_i^A.$$
(5)
Using $`𝔽`$ we can pull back the canonical $`(n+1)`$-form $`\mathrm{\Omega }`$ to obtain the so-called Cartan form $`\mathrm{\Theta }_{}`$,
$$\mathrm{\Theta }_{}=\left(𝔽\right)^{}\mathrm{\Theta }.$$
One can show (, theorem 3.1) that the Euler-Lagrange equations (2) are equivalent to
$$\left(j^1\phi \right)^{}(i_W\mathrm{\Omega }_{})=0WT𝔍^1𝒱,$$
where
$$\mathrm{\Omega }_{}=d\mathrm{\Theta }_{}=\left(𝔽\right)^{}\mathrm{\Omega }.$$
## 3 A vertical exterior derivative
Let us denote the multisymplectic phase space $`\left(𝔍^1𝒱\right)^{}`$ by $`𝒫`$ to simplify notation. In what follows we will need a mapping that is in some sense the vertical part of the exterior derivative on $`𝒫`$. In particular, it must have square zero. Whereas the derivation along fibres of $`𝒫`$ can be defined without additional data, the space of vertical forms as a subspace of arbitrary forms cannot<sup>3)</sup><sup>3)</sup>3)One can, however, define the space of vertical forms canonically, but in what follows we need the wedge product of a vertical form and an arbitrary one. For this, one needs an embedding of vertical forms in the space of forms, which in turn requires the use of a connection. .This is due to the fact that one needs to specify what is not vertical if one is looking for the dual of vertical vectors. For this, one needs a connection in the bundle $`𝒫`$ over $``$. This is dealt with in appendix A. With the help of this connection we can split $`T_p𝒫`$ into horizontal and vertical components for each point $`p`$ of $`𝒫`$. In local coordinates<sup>4)</sup><sup>4)</sup>4)When working in coordinates of $`𝒫`$, we will use the following convention which is similar to the one for coordinates on $`𝔍^1𝒱`$ . Small Latin indices sum over the base manifold directions, that is $`i,j,k`$ run from $`1`$ to $`n+1`$ if not specified otherwise. Capital Latin characters as $`A,B,C,D`$ run from $`1`$ to $`N`$ which is the dimension of a fibre of $`𝒱`$. Small Greek indices can be both base manifold and $`𝒱`$-fibre and dual jet bundle indices, i.e. $`\rho ,\sigma ,\tau =i,A,{\scriptscriptstyle \genfrac{}{}{0.0pt}{}{i}{A}}`$. Finally, capital letters from $`M`$ onwards stand for both $`A,B\mathrm{}`$ and $`{\scriptscriptstyle \genfrac{}{}{0.0pt}{}{i}{A}},{\scriptscriptstyle \genfrac{}{}{0.0pt}{}{j}{B}},\mathrm{}`$. $`(x^i,v^A,p_A^i,p)`$ we have a basis $`(𝔢_{(p)}^\alpha ,𝔢)`$, $`\alpha =i,A,{\scriptscriptstyle \genfrac{}{}{0.0pt}{}{i}{A}}`$ of $`T_p^{}𝒫`$ that is dual to a basis $`(𝔢_\alpha (p),𝔢)`$ of $`T_p𝒫`$. The detailed definition of the latter is explained in the appendix. In coordinates as above,
$$𝔢_{(p)}^i=dx^i,𝔢_{(p)}^A=dv^A+\mathrm{\Gamma }_{iB}^A(\pi (p))v^Bdx^i,𝔢_{(p)}^{\stackrel{i}{A}}=dp^{\genfrac{}{}{0.0pt}{}{i}{A}}+(\mathrm{\Lambda }_{kj}^i\delta _A^B\mathrm{\Gamma }_{kA}^B\delta _j^i)p^{\genfrac{}{}{0.0pt}{}{k}{B}}dx^j,𝔢_{(p)}^{}=dp.$$
(6)
Using the duality between $`T𝒫`$ and $`T^{}𝒫`$, we obtain a covariant derivative $`D^{}`$ on $`T𝒫`$, in particular
$$\begin{array}{cccccccccc}\hfill \left(D_{}^{}{}_{𝔢_M}{}^{}𝔢^N\right)(𝔢_\rho )(p)& =& 𝔢^N\left(D_{𝔢_M}𝔢_\rho \right)(p)& =& 0,\hfill & \hfill \left(D_{}^{}{}_{𝔢_M}{}^{}𝔢^i\right)(𝔢_\rho )(p)& =& 𝔢^i\left(D_{𝔢_M}𝔢_\rho \right)(p)& =& 0\hfill \end{array}$$
for all fibre indices $`M,N=A,{\scriptscriptstyle \genfrac{}{}{0.0pt}{}{i}{A}}`$ and all indices $`\rho `$. Thus for every $`\alpha (p)=\frac{1}{l!}\alpha _{\rho _1\mathrm{}\rho _l}^{}{}_{(p)}{}^{}𝔢_{(p)}^{\rho _1}\mathrm{}𝔢_{(p)}^{\rho _l}\mathrm{\Omega }^l𝒫=\mathrm{\Gamma }(\mathrm{\Lambda }^lT^{}𝒫)`$ the mapping<sup>5)</sup><sup>5)</sup>5)This mapping is a globally defined version of the vertical differential used by Kanatchikov in .
$$d^V=\left(𝔢_{(p)}^MD_{}^{}{}_{𝔢_M}{}^{}\right):\mathrm{\Omega }^l𝒫\mathrm{\Omega }^{l+1}𝒫$$
fulfils ($`M,N=A,{\scriptscriptstyle \genfrac{}{}{0.0pt}{}{i}{A}}`$ for $`i=1,\mathrm{},n`$, $`A=1,\mathrm{},N`$, $`\rho _l=i,A,{\scriptscriptstyle \genfrac{}{}{0.0pt}{}{i}{A}}`$)
$$\begin{array}{cc}\hfill \left(d^V\right)^2\alpha (p)& =\left(d^V\right)^2\frac{1}{l!}\alpha _{\rho _1\mathrm{}\rho _l}^{}{}_{(p)}{}^{}𝔢_{(p)}^{\rho _1}\mathrm{}𝔢_{(p)}^{\rho _l}\hfill \\ & =\left(𝔢_{(p)}^MD_{}^{}{}_{𝔢_M}{}^{}\right)\left(𝔢_{(p)}^ND_{}^{}{}_{𝔢_N}{}^{}\right)\frac{1}{l!}\alpha _{\rho _1\mathrm{}\rho _l}^{}{}_{(p)}{}^{}𝔢_{(p)}^{\rho _1}\mathrm{}𝔢_{(p)}^{\rho _l}\hfill \\ & =\frac{1}{l!}\left(𝔢_{(p)}^MD_{}^{}{}_{𝔢_M}{}^{}\right)\left(𝔢_N\alpha _{\rho _1\mathrm{}\rho _l}\right)_{(p)}𝔢_{(p)}^N𝔢_{(p)}^{\rho _1}\mathrm{}𝔢_{(p)}^{\rho _l}\hfill \\ & +\frac{1}{l!}\left(𝔢_{(p)}^MD_{}^{}{}_{𝔢_M}{}^{}\right)\underset{k=1}{\overset{l}{}}\alpha _{\rho _1\mathrm{}\rho _l}^{}{}_{(p)}{}^{}𝔢_{(p)}^N𝔢_{(p)}^{\rho _1}\mathrm{}\underset{=0}{\underset{}{D_{}^{}{}_{e_N}{}^{}𝔢^{\rho _k}}}\mathrm{}𝔢_{(p)}^{\rho _l}\hfill \\ & =\frac{1}{l!}\left(𝔢_M𝔢_N\alpha _{\rho _1\mathrm{}\rho _l}\right)_{(p)}𝔢_{(p)}^M𝔢_{(p)}^N𝔢_{(p)}^{\rho _1}\mathrm{}𝔢_{(p)}^{\rho _l}\hfill \\ & +\frac{1}{l!}\left(𝔢_N\alpha _{\rho _1\mathrm{}\rho _l}\right)_{(p)}𝔢_{(p)}^M\underset{=0}{\underset{}{D_{}^{}{}_{𝔢_M}{}^{}𝔢^N}}𝔢_{(p)}^{\rho _1}\mathrm{}𝔢_{(p)}^{\rho _l}\hfill \\ & +\frac{1}{l!}\underset{k=1}{\overset{l}{}}\left(𝔢_N\alpha _{\rho _1\mathrm{}\rho _l}\right)_{(p)}𝔢_{(p)}^M𝔢_{(p)}^N𝔢_{(p)}^{\rho _1}\mathrm{}\underset{=0}{\underset{}{D_{}^{}{}_{e_N}{}^{}𝔢^{\rho _k}}}\mathrm{}𝔢_{(p)}^{\rho _l}\hfill \\ & =\frac{1}{2l!}\left([𝔢_M,𝔢_N]\alpha _{\rho _1\mathrm{}\rho _l}\right)_{(p)}𝔢_{(p)}^M𝔢_{(p)}^N𝔢_{(p)}^{\rho _1}\mathrm{}𝔢_{(p)}^{\rho _l}\hfill \\ & =0,\hfill \end{array}$$
that is, $`\left(d^V\right)^2=0`$. This justifies the name vertical exterior derivative.
### 3.1 Poincaré lemma for $`d^V`$
###### Lemma 3.1 (Poincaré lemma for $`d^V`$)
Let $`\alpha \mathrm{\Omega }^r𝒫`$ with $`d^V\alpha =0`$. Then for every $`p𝒫`$ there exists a neighbourhood $`U_p`$ and a $`(r1)`$-Form $`\beta `$ such that $`\alpha \text{ }𝒰_p=d^V\beta `$.
Proof: As fibres of $`𝒫`$ are contractible and $`d^V`$, restricted to such a fibre, acts like the exterior derivative, this is a consequence of the Poincaré lemma. In detail, let $`m=\pi (p)`$ and $`𝒰`$ be a neighbourhood of $`m`$ such that $`𝒫\text{ }𝒰`$ is trivial. Now let $`𝒰_p=\pi ^1(𝒰)`$. On $`𝒰_p`$, we can choose a basis $`(𝔢_{(p)}^\alpha ,𝔢_{(p)}^i)`$ of $`T^{}𝒫\text{ }𝒰_p`$ as above (in what follows we will omit the point $`p`$ when writing a covector). Then we have
$$\alpha (p)=\underset{l=0}{\overset{r}{}}\alpha _l(p),$$
where $`\alpha _l`$ is of the form
$$\alpha _l(p)=\frac{1}{r!}\alpha _{M_1\mathrm{}M_li_{l+1}\mathrm{}i_r}(p)𝔢^{M_1}\mathrm{}𝔢^{M_l}𝔢^{i_{l+1}}\mathrm{}𝔢^{i_r}.$$
As $`𝔢^{M_1}\mathrm{}𝔢^{M_l}𝔢^{i_{l+1}}\mathrm{}𝔢^{i_r}`$ and $`𝔢^{M_1}\mathrm{}𝔢^{M_j}𝔢^{i_{j+1}}\mathrm{}𝔢^{i_r}`$ are linearly independent for $`jl`$, $`d^V\alpha =0`$ implies
$$d^V\alpha _l=0l=1,\mathrm{},r.$$
Furthermore, we see that
$$d^V\alpha _l(p)=0d^V\alpha _{l,i_{l+1}\mathrm{}i_r}(p)=0i_{l+1},\mathrm{},i_r=1,\mathrm{},n,$$
where
$$\alpha _l(p)=\frac{1}{(rl)!}\alpha _{l,i_{l+1}\mathrm{}i_r}(p)𝔢^{i_{l+1}}\mathrm{}𝔢^{i_r}.$$
Now, if we restrict the $`\alpha _{l,i_{l+1}\mathrm{}i_r}`$ to a fixed fibre $`𝒫_m`$ of $`𝒫`$, applying $`d^V`$ corresponds to the exterior derivative on that space. As the fibre under consideration is a vector space, it follows that
$$\alpha _{l,i_{l+1}\mathrm{}i_r}\text{ }𝒵_m=d^V\beta _{(l1),i_{l+1}\mathrm{}i_r}^m,$$
and hence
$$\begin{array}{cc}\hfill \alpha (p)=\underset{l=0}{\overset{r}{}}\alpha _l(p)& =\underset{l=0}{\overset{r}{}}\frac{1}{(rl)!}\alpha _{l,i_{l+1}\mathrm{}i_r}(p)𝔢^{i_{l+1}}\mathrm{}𝔢^{i_r}\hfill \\ & =\underset{l=0}{\overset{r}{}}\frac{1}{(rl)!}\left(d^V\beta _{(l1),i_{l+1}\mathrm{}i_r}^{\pi (p)}\right)𝔢^{i_{l+1}}\mathrm{}𝔢^{i_r}\hfill \\ & =\underset{l=0}{\overset{r}{}}\frac{1}{(rl)!}d^V\left(\beta _{(l1),i_{l+1}\mathrm{}i_r}^{\pi (p)}𝔢^{i_{l+1}}\mathrm{}𝔢^{i_r}\right)\hfill \\ & =d^V\beta (p),\hfill \end{array}$$
where
$$\beta (p)=\underset{l=0}{\overset{r}{}}\frac{1}{(rl)!}\beta _{(l1),i_{l+1}\mathrm{}i_r}^{\pi (p)}𝔢^{i_{l+1}}\mathrm{}𝔢^{i_r}.$$
$`\mathrm{}`$
## 4 Field equations
As already mentioned the multisymplectic phase space $`𝒫`$ of a given field theory is chosen to be the affine dual of the first jet bundle $`𝒥^1𝒱`$, but the field equations (2) are formulated on $`𝒥^1𝒱`$ itself. Hence, similar to ordinary mechanics, one uses the covariant Legendre transformation (4) to reformulate the theory. For this, let us assume that the middle equation of (5) can be rearranged so that the variables $`v_i^A`$ can be expressed in terms of $`(x^i,v^A,p^{\genfrac{}{}{0.0pt}{}{i}{A}})`$. In other words, we require
$$det\left(\frac{^2L}{v_i^Av_j^B}\right)0,v_i^A=\phi _i^A(x^i,v^A,p^{\genfrac{}{}{0.0pt}{}{i}{A}}).$$
Then the Lagrange density $`L`$, (5), becomes a function over phase space,
$$\stackrel{~}{L}(x^i,v^A,p^{\genfrac{}{}{0.0pt}{}{i}{A}})=L(x^i,v^A,\phi _i^A(x^i,v^A,p^{\genfrac{}{}{0.0pt}{}{i}{A}}))$$
and we obtain the covariant Hamiltonian
$$H(x^i,v^A,p^{\genfrac{}{}{0.0pt}{}{i}{A}})=\stackrel{~}{L}(x^i,v^A,p^{\genfrac{}{}{0.0pt}{}{i}{A}})p^{\genfrac{}{}{0.0pt}{}{i}{A}}\phi _i^A(x^i,v^A,p^{\genfrac{}{}{0.0pt}{}{i}{A}}).$$
(7)
Using this, the generalised Hamiltonian equations
$$\frac{H}{v^A}=\frac{p^{\genfrac{}{}{0.0pt}{}{i}{A}}}{x^i},\frac{H}{p^{\genfrac{}{}{0.0pt}{}{i}{A}}}=\frac{v^A}{x^i},$$
(8)
are equivalent to the Euler-Lagrange equations (2), (, ch. 4.2). Note, however, that $`H`$ is not a function but (7) rather describes a subset of $`𝒫`$ which is the image of $`𝒥^1𝒱`$ under $`𝔽`$. The coordinates we have used up to now have arisen in a natural way from coordinates on $``$ and $`𝒱`$; they simply are the components of the tangent map of a given section. If one uses the connection $`\mathrm{\Gamma }`$ as a zero section of $`𝔍^1𝒱𝒱`$ one turns $`𝔍^1𝒱`$ into a vector space $`𝔙𝒱T^{}`$, and $`𝒫`$ splits into the direct sum of a line bundle and the bundle of linear mappings of the former vector bundle to $`\mathrm{\Lambda }^{n+1}T^{}`$ (cf. ). In coordinates this corresponds to the change
$$\mathrm{\Psi }:(x^i,v^A,v_i^A)(x^i,v^A,\stackrel{~}{v}_i^A=v_i^A+\mathrm{\Gamma }_{iB}^Av^B).$$
(9)
Using
$$\begin{array}{cc}\hfill \frac{}{v_i^A}\mathrm{\Psi }^1& =\frac{\mathrm{\Psi }^1}{\stackrel{~}{v}_i^A},\hfill \\ \hfill \frac{}{v^A}\mathrm{\Psi }^1& =\frac{\mathrm{\Psi }^1}{\stackrel{~}{v}^A}\mathrm{\Gamma }_{iA}^B\frac{\mathrm{\Psi }^1}{\stackrel{~}{v}_i^B}\hfill \end{array}$$
(10)
equation (2) becomes (for $`_\mathrm{\Gamma }=\mathrm{\Psi }^1`$)
$$\frac{_\mathrm{\Gamma }}{v^A}j^1\phi _i\left(\frac{_\mathrm{\Gamma }}{\stackrel{~}{v}_i^A}j^1\phi \right)=0$$
(11)
For the affine bundle $`𝒫`$ the change of coordinates induces a mapping $`\mathrm{\Psi }^{}:(x^i,v^A,p,p^{\genfrac{}{}{0.0pt}{}{i}{A}})(x^i,v^A,p+\mathrm{\Gamma }_{iB}^Ap^{\genfrac{}{}{0.0pt}{}{i}{A}}v^B,p^{\genfrac{}{}{0.0pt}{}{i}{A}})`$. Let $`_\mathrm{\Gamma }=H(\mathrm{\Psi }^{})^1`$. As we have a global splitting of $`𝒫`$ induced by the connection $`\mathrm{\Gamma }`$, this is a function on $`(𝔙𝒱T^{})^{}`$. Differentiating $`_\mathrm{\Gamma }`$ as in (8) with respect to $`v^A`$ and $`p_A^i`$ one obtains on solutions $`j^1\phi `$ of (2)
$$\frac{_\mathrm{\Gamma }}{v^A}=_i\stackrel{~}{p}^{\genfrac{}{}{0.0pt}{}{i}{A}},\frac{_\mathrm{\Gamma }}{\stackrel{~}{p}^{\genfrac{}{}{0.0pt}{}{i}{A}}}=_iv^A.$$
(12)
For the last equation we have used that in the coordinates introduced the first jet prolongation has the form (1). A similar result can be found in .
Now we are going to formulate the equations of motion in a coordinate free manner. Let solutions of (2) be described by $`(n+1)`$-vector fields $`\stackrel{n+1}{X}\mathrm{\Gamma }(\mathrm{\Lambda }^{n+1}T\left(𝔍^1𝒱\right)^{})`$ with $`T\overline{\pi }\stackrel{n+1}{X}0`$. Further, let $`\stackrel{n+1}{X}{}_{}{}^{V}=\stackrel{n+1}{X}(T\overline{\pi }\stackrel{n+1}{X})^h`$ be the vertical component of $`\stackrel{n+1}{X}`$, where $`(T\overline{\pi }\stackrel{n+1}{X})^h`$ is the horizontal lift according to the splitting induced by the mapping (43) in the appendix B. If $`\mathrm{\Omega }^{(2,n)}=d^V\mathrm{\Theta }^{(1,n)}`$, where $`\mathrm{\Theta }^{(1,n)}`$ denotes the vertical component of $`\mathrm{\Theta }`$ (so that in the splitting above $`\mathrm{\Omega }^{(2,n)}`$ has two vertical and $`n`$ horizontal components),
$$\mathrm{\Theta }^{(1,n)}=\mathrm{\Theta }\mathrm{\Theta }^H,(X)^h\text{}\text{ }\text{ }\mathrm{\Theta }^{(1,n)}=0X\mathrm{\Lambda }^{n+1}T,X\text{}\text{ }\text{ }\mathrm{\Theta }^H=0X𝔙𝒫.$$
the generalised Hamilton equations (8) are equivalent to
$$\left(X{}_{}{}^{V}\text{}\text{ }\text{ }\mathrm{\Omega }^{(2,n)}\right)^{(1,0)}=()^{n+1}d^VH.$$
## 5 Hamiltonian forms and a graded Poisson structure
With the help of the vertical exterior derivative we can define the graded vertical Lie derivative by an $`r`$-vector field by
$$_{\stackrel{r}{X}}\mathrm{\Phi }=\stackrel{r}{X}\text{}\text{ }\text{ }d^V\mathrm{\Phi }+()^{r+1}d^V\left(\stackrel{r}{X}\text{}\text{ }\text{ }\mathrm{\Phi }\right)$$
(13)
for every form $`\mathrm{\Phi }`$ on $`T𝒫`$.
An $`r`$-vector field $`\stackrel{r}{X}`$ is called a Hamiltonian multi-vector field iff there is a horizontal $`(n+1r)`$-form $`\stackrel{(n+1r)}{F}`$ that satisfies
$$\stackrel{r}{X}\text{}\text{ }\text{ }\mathrm{\Omega }^{(2,n)}=d^V\stackrel{(n+1r)}{F}.$$
(14)
The set of all such forms will be called the set of Hamiltonian forms and denoted by $``$. Not every horizontal form is automatically Hamiltonian. Indeed, if we write in local coordinates
$$\stackrel{(n+1r)}{F}=\frac{1}{r!}F^{i_1\mathrm{}i_r}(𝔢_{i_1\mathrm{}i_r}\text{}\text{ }\text{ }\omega ),$$
(15)
where $`\omega `$ is the horizontally lifted volume form of $``$ and $`𝔢_{i_1\mathrm{}i_r}=𝔢_{i_1}\mathrm{}𝔢_{i_r}`$, we find for $`n+1>r`$ ()
$$\begin{array}{cc}\hfill rX^{A[j_1\mathrm{}j_{r1}}\delta _j^{i]}& =_{\stackrel{i}{A}}F^{j_1\mathrm{}j_{r1}i}\hfill \\ \hfill rX^{\stackrel{i}{A}j_1\mathrm{}j_{r1}}& =_AF^{j_1\mathrm{}j_{r1}i}\hfill \end{array}$$
(16)
which puts a restriction on the admissible horizontal forms $`F`$ with $`r<n+1`$, namely
$$_{\stackrel{k}{B}}F^{j_1\mathrm{}j_r}=0$$
(17)
for all $`k\{j_1,\mathrm{},j_r\}`$. For $`r=n+1`$ the first equation in (16) does not lead to any restriction, since $`j`$ has to be in $`\{j_1,\mathrm{},j_n,i\}`$ in any case. Moreover, from $`d^V\stackrel{r}{X}\text{}\text{ }\text{ }\mathrm{\Omega }^{(2,n)}=(d^V)^2\stackrel{(n+1r)}{F}=0`$ we deduce in particular
$$\underset{i=1}{\overset{n+1}{}}\underset{A,B=1}{\overset{N}{}}_{\stackrel{i}{A}}X^{Bi_1\mathrm{}i_r}𝔢^{\stackrel{i}{A}}𝔢^{\stackrel{j}{B}}𝔢_{i_1\mathrm{}i_rj}\text{}\text{ }\text{ }\omega =0,$$
which implies
$$\left(_{\stackrel{j_1}{A}}\right)^2F^{j_1\mathrm{}j_r}=r_{\stackrel{j_1}{A}}X^{Bj_1\mathrm{}j_{r1}}=0\text{ (No summation over }j_1\text{.)}$$
(18)
Hence, as already remarked in , the coordinate expression of $`F`$ can depend on the coordinates of the fibre of $`𝒫`$ in a specific polynomial way only, where each coordinate $`p^{\stackrel{i}{A}}`$ appears at most to the first power.
If $`n=0`$ then $`\mathrm{\Omega }^{(2,0)}`$ does not contain any horizontal degree and the Hamiltonian forms are just functions on $`𝒫`$. For those, the conditions (16) become
$$X^A=_{\stackrel{1}{A}}F^1,X^{\stackrel{1}{A}}=_AF^1.$$
(19)
Hence, arbitrary functions $`F`$ are allowed.
###### Lemma 5.1
Let $`\stackrel{(n+1r)}{F}=\frac{1}{r!}F^{j_1\mathrm{}j_r}𝔢_{j_1\mathrm{}j_r}\text{}\text{ }\text{ }\omega `$ be a Hamiltonian form. If $`r<n+1`$, then the coefficient functions are of the following form.
$$F^{j_1\mathrm{}j_r}(x,v,p)=\frac{1}{r!}\underset{k=0}{\overset{r}{}}p^{\stackrel{j_1}{A_1}}\mathrm{}p^{\stackrel{j_k}{A_k}}f^{A_1\mathrm{}A_kj_{k+1}\mathrm{}j_r},$$
(20)
where the functions $`f`$ are antisymmetric in the upper indices.
If $`n+1=r`$, then the set of Hamiltonian forms consists of all functions on the phase space $`𝒫`$.
With that, we have the following observation.
###### Lemma 5.2
If $`\stackrel{r}{X}`$,$`\stackrel{s}{X}`$ are Hamiltonian multi-vector fields, then
$$\stackrel{r}{X}\text{}\text{ }\text{ }\stackrel{s}{X}\text{}\text{ }\text{ }\mathrm{\Omega }^{(2,n)}$$
(21)
is a Hamiltonian form.
Proof: This can be checked by a calculation using coordinates. Let us suppose $`n>0`$. (The case $`n=0`$ is easy because there is no additional restriction on Hamiltonian forms apart from having horizontal degree zero.) Firstly, the above expression (21) is horizontal. Since $`\stackrel{r}{X}`$ and $`\stackrel{s}{X}`$ are assumed to be Hamiltonian, there are horizontal forms $`F`$ and $`G`$ satisfying (14) respectively. We will show that $`\stackrel{r}{X}\text{}\text{ }\text{ }\stackrel{s}{X}\text{}\text{ }\text{ }\mathrm{\Omega }^{(2,n)}`$ is of the form (20).
$$\begin{array}{cc}\hfill \stackrel{r}{X}\text{}\text{ }\stackrel{s}{X}\text{}\text{ }\mathrm{\Omega }^{(2,n)}& =\frac{1}{(r1)!}\frac{1}{(s1)!}()^{(r1)}\stackrel{r}{X}\text{}^{Mi_1\mathrm{}i_{r1}}\stackrel{s}{X}\text{}^{Nj_1\mathrm{}j_{s1}}𝔢_M𝔢_N,𝔢^A𝔢^{\stackrel{i}{A}}\left(𝔢_{i_1\mathrm{}i_{r1}j_1\mathrm{}j_{s1}i}\text{}\text{ }\omega \right)\hfill \\ & =\frac{1}{(r+s1)!}H^{i_1\mathrm{}i_{r1}j_1\mathrm{}j_{s1}i}\left(𝔢_{i_1\mathrm{}i_{r1}j_1\mathrm{}j_{s1}i}\text{}\text{ }\omega \right).\hfill \end{array}$$
Because of the special form of $`\stackrel{r}{X}`$ and $`\stackrel{s}{X}`$ according to lemma 5.1 we find
$$_{\stackrel{i_1}{A}}H^{i_1\mathrm{}i_{r+s1}}=_{\stackrel{i_2}{A}}H^{i_1\mathrm{}i_{r+s1}}$$
(22)
and
$$_{\stackrel{i}{A}}H^{i_1\mathrm{}i_{r+s1}}=0\text{ for }i\{i_1,\mathrm{},i_{r+s1}\}.$$
(23)
This shows that $`\stackrel{r}{X}\text{}\text{ }\text{ }\stackrel{s}{X}\text{}\text{ }\text{ }\mathrm{\Omega }^{(2,n)}`$ fulfils the conditions derived from (16) and thus is Hamiltonian. $`\mathrm{}`$
Looking at equation (21) we can ask what the corresponding Hamiltonian multi-vector field might be. One calculates
$$\begin{array}{cc}\hfill d^V\left(\stackrel{r}{X}\text{}\text{ }\stackrel{s}{X}\text{}\text{ }\mathrm{\Omega }^{(2,n)}\right)& =d^V\left(\stackrel{r}{X}\text{}\text{ }\stackrel{s}{X}\text{}\text{ }\mathrm{\Omega }^{(2,n)}\right)+()^{r+1}\stackrel{r}{X}\text{}\text{ }d^V\left(\stackrel{s}{X}\text{}\text{ }\mathrm{\Omega }^{(2,n)}\right)\hfill \\ & =_{\stackrel{r}{X}}\stackrel{s}{X}\text{}\text{ }\mathrm{\Omega }^{(2,n)}\hfill \end{array}$$
As $`_{\stackrel{r}{X}}\mathrm{\Omega }^{(2,n)}=0`$ this looks like the Lie bracket of $`\stackrel{r}{X}`$ and $`\stackrel{s}{X}`$ being inserted in $`\mathrm{\Omega }^{(2,n)}`$. Now in symplectic mechanics the Lie bracket of two (locally) Hamiltonian vector fields is the vector field associated to the Poisson bracket of the Hamiltonian functions of the former. Hence, by analogy, we define a bracket as follows:
$$\{\stackrel{r}{F},\stackrel{s}{F}\}=()^{n+1r}\stackrel{n+1r}{X}\text{}\text{ }\text{ }\stackrel{n+1s}{X}\text{}\text{ }\text{ }\mathrm{\Omega }^{(2,n)},$$
(24)
where $`\stackrel{r}{F},\stackrel{s}{F}`$ are Hamiltonian forms and $`\stackrel{nr}{X},\stackrel{ns}{X}`$ denote the corresponding vector fields. Note that whereas there is some ambiguity in the choice of a Hamiltonian (multi-)vector field in eq. (14), this does not lead to an ambiguity of the above bracket. Indeed, a vector field $`X`$ that vanishes on $`\mathrm{\Omega }^{(2,n)}`$ must have vanishing coefficients $`X^{Mi_1\mathrm{}i_k}`$ but can have non vanishing components $`X^{M_1\mathrm{}M_ji_1\mathrm{}i_l}`$. The latter, however, do not contribute to the bracket since $`\mathrm{\Omega }^{(2,n)}`$ is of type $`(2,n)`$<sup>6)</sup><sup>6)</sup>6)The author wishes to thank the referees for pointing out the remaining ambiguity to him. .
###### Proposition 5.1
The bracket
$$\{,\}:\times $$
(25)
defined by (24) has the following properties:
1. It is graded antisymmetric,
$$\{\stackrel{r}{F},\stackrel{s}{F}\}=()^{(nr)(ns)}\{\stackrel{s}{F},\stackrel{r}{F}\}.$$
2. It fulfils a graded Jacobi identity,
$$()^{(nr)(nt)}\{\stackrel{r}{F},\{\stackrel{s}{F},\stackrel{t}{F}\}\}+()^{(ns)(nr)}\{\stackrel{s}{F},\{\stackrel{t}{F},\stackrel{r}{F}\}\}+()^{(nt)(ns)}\{\stackrel{t}{F},\{\stackrel{r}{F},\stackrel{s}{F}\}\}=0.$$
3. There is a product
$$\stackrel{r}{F}\stackrel{s}{F}=^1(\stackrel{r}{F}\stackrel{s}{F})=()^{(n+1r)(n+1s)}\stackrel{s}{F}\stackrel{r}{F},$$
(26)
where $``$ is the operation induced by the Hodge operator on $``$ that maps Hamiltonian functions to Hamiltonian functions. With respect to $``$, the above defined bracket shows a graded Leibniz rule,
$$\{\stackrel{r}{F},\stackrel{s}{F}\stackrel{t}{F}\}=\{\stackrel{r}{F},\stackrel{s}{F}\}\stackrel{t}{F}+()^{(nr)(n+1s)}\stackrel{s}{F}\{\stackrel{r}{F},\stackrel{t}{F}\}.$$
(27)
Proof. $`\mathit{1})`$ is an immediate consequence of the definition.
$`\mathit{2})`$ is a straightforward calculation if one uses
$$_{\stackrel{k}{B}}X^{\stackrel{i}{A}j_1\mathrm{}j_1}=_AX^{b[j_1\mathrm{}j_1}\delta _k^{i]},_BX^{\stackrel{i}{A}j_1\mathrm{}j_1}=_AX^{\stackrel{i}{B}j_1\mathrm{}j_1}$$
(28)
which can be deduced from changing the order of differentiation in (16).
As for $`\mathit{3})`$, using
$$(𝔢_{i_1\mathrm{}i_r}\text{}\text{ }\text{ }𝔢^1\mathrm{}𝔢^n)=𝔢^{i_1}\mathrm{}𝔢^{i_r}$$
we find
$$\stackrel{n+1q}{G}\stackrel{n+1r}{H}=\frac{1}{(q+r)!}G^{i_1\mathrm{}i_q}H^{i_{q+1}\mathrm{}i_{q+r}}\left(𝔢_{i_1\mathrm{}i_{q+r}}\text{}\text{ }\text{ }\omega \right)$$
(29)
and hence
$$\begin{array}{cc}\hfill \{\stackrel{n+1p}{F},\stackrel{n+1q}{G}\stackrel{n+1r}{H}\}& =()^p\frac{1}{(p1)!}X_F^{Mi_1\mathrm{}i_{p1}}\text{}\text{ }d^V\left(GH\right)\hfill \\ & =X_F^{Mi_1\mathrm{}i_{p1}}(_MG^{j_1\mathrm{}j_q}H^{j_{q+1}\mathrm{}j_{q+r}}𝔢_{i_1\mathrm{}i_{(p1)}j_1\mathrm{}j_{q+r}}\text{}\text{ }\omega \hfill \\ & +()^{(p1)q}G^{j_1\mathrm{}j_q}X_F^{Mi_1\mathrm{}i_{p1}}(_MH^{j_{q+1}\mathrm{}j_{q+r}})𝔢_{j_1\mathrm{}j_qi_1\mathrm{}i_{(p1)}j_1\mathrm{}j_{q+r}}\text{}\text{ }\omega \hfill \\ & =\{\stackrel{p}{F},\stackrel{q}{G}\}\stackrel{r}{H}+()^{(p1)q}\stackrel{q}{G}\{\stackrel{p}{F},\stackrel{r}{H}\}\hfill \end{array}$$
(30)
$`\mathrm{}`$
One might ask about the dependence of the bracket on the connections $`\mathrm{\Gamma }`$ and $`\mathrm{\Lambda }`$. As can be seen from (6), different choices of connections amount to differences in the horizontal terms of the vertical forms that have been used in the definition of $`d^V`$. But from (14) we learn that this change can have an effect on those terms of $`X`$ that have two or more vertical components only. Again, those terms do not contribute to the bracket. Hence the Poisson bracket does not depend on $`\mathrm{\Gamma }`$ nor $`\mathrm{\Lambda }`$.
## 6 Recovering mechanics
To recover Hamiltonian mechanics we proceed as follows. Let $`𝒬`$ be the coordinate space of the theory. Then, $`=`$ and $`𝒱`$ is trivial $`𝒱=\times 𝒬`$. Hence, $`T𝒱`$ decomposes into $`T𝒱=T𝒬`$. The condition for a mapping $`\phi \psi :T=T𝒱=T𝒬`$ to be in $`𝔍^1𝒱`$ is thus
$$T\pi (\phi \psi )=\psi =T\text{id}_{}=1.$$
(31)
As the mapping $`\phi `$ is defined by its value at $`1`$ we conclude $`𝔍^1𝒱=T𝒬\times `$ and, going to the dual we obtain the phase space,
$$𝒫\left(𝔍^1𝒱\right)^{}=(T^{}𝒬)\times .$$
(32)
The canonical $`1`$-form $`\mathrm{\Theta }`$ reads
$$\mathrm{\Theta }(t,v^A,p,p_A)=p_Adv^A+pdt$$
whereas $`\mathrm{\Omega }^{(2,0)}`$ is
$$\mathrm{\Omega }^{(2,0)}(t,v^A,p,p_A)=dp_Adv^A$$
which is just the canonical $`2`$-Form. As the base manifold is one-dimensional, horizontal forms are either functions or 1-forms on $`T^{}𝒬`$. Now in this case equation (14) admits the former case since $`\mathrm{\Omega }^{(2,0)}`$ does not contain any horizontal component. Therefore the Hamiltonian multi-vector fields can be ordinary vector fields on $`T^{}𝒬`$ only, and we have
$$X_F(t,v,p)=_{p^A}F(t,v,p)_{p^A}_{v^A}F(t,v,p)_{v^A}$$
(33)
There is no additional restriction to admissible Hamiltonian functions (cf. (19)) and we have arrived at the stage of Hamiltonian mechanics (cf. ). As the bundle $`𝒱`$ is trivial we do not need a connection really, so there is no need for $`𝒬`$ to be a vector bundle. As the base manifold is one-dimensional only, the product of two Hamiltonian forms always gives zero. This can be remedied if one includes horizontal $`1`$-forms in the set of observables in addition to functions<sup>7)</sup><sup>7)</sup>7)I. Kanatchikov, private communication.. This leads to the extension of the notion of Hamiltonian vector fields to form valued vector fields.
In , sec. 4, where a Poisson structure is defined on (de Rham) equivalence classes of forms on $`𝒫`$, the Poisson algebra consists of those functions only for which the dependence on the parameter is the physical time, i.e. which solve the equations of motion when differentiated with respect to this parameter. Here, in contrast, nothing can be said about the ”time” dependence of Hamiltonian forms.
## 7 The case of a scalar field
In the case of a scalar field, the fibre of $`𝒱`$ is isomorphic to $``$. Using a connection $`\mathrm{\Gamma }:𝒱𝔍^1𝒱`$, we obtain an isomorphism
$$𝔍^1𝒱\stackrel{\mathrm{\Gamma }}{}𝔙𝒱_𝒱T^{},𝔙𝒱\times .$$
(34)
Hence
$$𝔍^1𝒱\stackrel{\mathrm{\Gamma }}{}\text{pr}^{}(T^{}),$$
(35)
where pr denotes the canonical projection of the bundle $`𝒱`$. Using (14) one immediately verifies in local coordinates $`(x^i,v,p^i,p)`$ of $`𝒫`$ in this case (let $`e_i`$ denote the horizontal lifts of tangent vectors of $``$ and $`𝔢^i`$ be the vertical forms with respect to the splitting discussed in the appendix; the determinant comes from the volume element on $``$)
$`_v\text{}\text{ }\text{ }\mathrm{\Omega }^{(2,n)}=𝔢^i(e_i\text{}\text{ }\text{ }\omega )=d^Vp^i(e_i\text{}\text{ }\text{ }\omega ),`$
$`_{i=1}^{n+1}_{p^i}(()^i(\sqrt{detg})e_1\mathrm{}e_{i1}e_{i+1}\mathrm{}e_{n+1})\text{}\text{ }\text{ }\mathrm{\Omega }^{(2,n)}=d^Vv,`$
hence $`\mathrm{\Pi }(x,v,p)=p^i(e_i\text{}\text{ }\text{ }\omega )`$ satisfies
$$\{\mathrm{\Pi },\mathrm{\Phi }\}=1$$
for $`\mathrm{\Phi }(x,v)=v`$, but $`\mathrm{\Pi }1=0=\mathrm{\Phi }1`$. The unit with respect to $``$ is $`\omega `$, so one should look for solutions of
$`X\text{}\text{ }Y\text{}\text{ }\mathrm{\Omega }^{(2,n)}=\omega .`$
This cannot be solved, as $`\mathrm{\Omega }^{(2,n)}`$ contains $`n`$ horizontal components, whereas $`\omega `$ is a horizontal $`(n+1)`$-form. As before, one might have to include vector fields that are form valued, i.e. endomorphisms of $`\mathrm{\Lambda }^{}T^{}𝒫`$.
Note, however, that the connection $`\mathrm{\Gamma }`$ remains arbitrary: Although it is hidden in the expression for $`\mathrm{\Pi }`$,
$$\mathrm{\Pi }(x,v,p)=p^i(e_i\text{}\text{ }\text{ }\omega )=p^i(_i\text{}\text{ }\text{ }\omega ),$$
$`\mathrm{\Pi }`$ is in fact independent of it.
## 8 Conclusions
In multisymplectic geometry we take the phase space $`𝒫`$ to be the affine dual of the first jet bundle to a given vector bundle $`𝒱`$. It is then possible to define (graded) Poisson brackets (24) on $`𝒫`$ even for nontrivial vector bundles. For this one needs a covariant derivative on the $`(n+1)`$-dimensional base manifold $``$ (space-time) and a connection on the vector bundle of the fields under consideration.
Kanatchikov has proposed a similar construction by making use of equivalence classes of forms modulo forms of higher horizontal degree (). This is equivalent to the use of the construction elaborated in this article, as a vertical form, say $`𝔢^A`$, differs from the coordinate expression $`dv^A`$ by horizontal components only, cf. (6),
$$𝔢_{(p)}^A=dv^A+\mathrm{\Gamma }_{iB}^A(\pi (p))v^Bdx^i.$$
(36)
Hence, $`𝔢_{(p)}^A`$ and $`dv^A`$ define the same equivalence class, independent of the connections $`\mathrm{\Gamma }`$ and $`\mathrm{\Lambda }`$ used. The same applies to the bracket: Whereas the correspondence of Hamiltonian forms and multi-vector fields is ambiguous and does depend on the connections chosen, the (graded) Poisson bracket does not. Admissible observables are so-called Hamiltonian forms, horizontal forms that satisfy certain consistency relations, (16). It turns out that those Hamiltonian forms are polynomial in the momenta, i.e. coordinates of the fibres of $`𝒫𝒱`$, cf. (20).
In addition $``$ has to be orientable in order to define the multiplication (26) between Hamiltonian forms. For Hamiltonian forms of the same degree, this product is commutative but gives zero if the form degree is less than $`(n+1)/2`$.
If space-time is taken to be one-dimensional the whole formalism reduces to ordinary mechanics on a configuration space $`𝒬`$. Hamiltonian forms then are arbitrary functions on the extended phase space $`T^{}𝒬\times `$, and the Poisson bracket takes the standard form. However, the product $``$ of functions always gives zero in this case.
In the case of a scalar field, given a (local) field $`\mathrm{\Phi }`$ one can define a Hamiltonian form $`\mathrm{\Pi }`$ that satisfies $`\{\mathrm{\Pi },\mathrm{\Phi }\}=1`$, but the constant function $`1`$ is not the unit with respect to $``$. Rather, this rôle is played by $`\omega `$, the pulled back volume form from $``$. To obtain $`\{\mathrm{\Pi },\mathrm{\Phi }\}=\omega `$ one has to extend the notion of Hamiltonian vector fields in a way similar to that needed in the mechanical case (as mentioned above), namely one has to include form valued vector fields, i.e. endomorphisms of $`\mathrm{\Lambda }^{}T^{}𝒫`$.
The Poisson structure is graded in the following way. Let the degree of a (homogeneous) Hamiltonian form be its degree as an element of the exterior algebra. Then the degree of the Poisson bracket of two Hamiltonian forms is the sum of the respective degree minus $`n`$, the number of space directions, while the degree of the product of two Hamiltonian forms is the sum of the degrees minus $`n+1`$,
$$\mathrm{deg}\{\stackrel{r}{F},\stackrel{s}{F}\}=\mathrm{deg}\stackrel{r}{F}+\mathrm{deg}\stackrel{s}{F}n,\mathrm{deg}\stackrel{r}{F}\stackrel{s}{F}=\mathrm{deg}\stackrel{r}{F}+\mathrm{deg}\stackrel{s}{F}(n+1).$$
(37)
Looking at proposition 5.1 we find that the graded antisymmetry of the $`\{,\}`$, the graded Jacobi identity, the graded derivation property with respect to $``$ and the graded commutativity of $``$ all match with each other.
As already remarked in the examples, How to relate observables of physical fields and Hamiltonian forms. This point requires further investigation, especially the relation with the multiplicative structure. In particular, the notion of canonical conjugate momenta needs to be clarified.
Note added. As pointed out by one of the referees the above construction depends heavily on the vector space structure of fibres of $`𝒱`$. This might be sufficient for the study of such field theories where the fields take their values in a vector space. For classical mechanics on arbitrary configurations spaces, nevertheless, or in the case of string theory – whenever the target space is not Minkowski space – there’s is indeed a need for a generalisation of the construction. In this article, all that is used really is a splitting of the tangent space $`T𝔍^1𝒱`$ in horizontal and vertical subspaces with respect to the canonical projection onto $``$. Such a splitting does not exist canonically. There is, however, a natural way to split $`(\pi _1)_{0}^{1}{}_{}{}^{}\left(T𝔍^1𝒱\right)`$, the pull back of $`T𝔍^1𝒱`$ onto $`𝔍^1𝔍^1𝒱`$, the first jet bundle of $`𝔍^1𝒱`$. Now every connection $`\overline{\mathrm{\Gamma }}`$ on $`𝔍^1𝒱`$ (viewed as a bundle over $``$) defines a map $`\overline{\mathrm{\Gamma }}:𝔍^1𝒱𝔍^1𝔍^1𝒱`$ and hence induces a splitting of $`T𝔍^1𝒱`$. For $`𝒱`$ being a general fibre bundle, the connection $`\mathrm{\Gamma }`$ does not depend linearly on the fibre coordinates (cf. (39)). Rather, it takes the most general form
$$\mathrm{\Gamma }:𝒱(x^i,u^A)(x^i,u^A,\mathrm{\Gamma }_i^A).$$
In this case, in the local expression (42), one has to replace $`\mathrm{\Gamma }_{iB}^Au_j^B`$ by $`_{u^B}\mathrm{\Gamma }_i^Au_j^B`$ and $`\mathrm{\Gamma }_{kB}^Au^B`$ by $`\mathrm{\Gamma }_k^A`$.
Acknowledgements. The author’s interest in this subject was initiated by very elucidating discussions with H. Römer and M. Bordemann about quantisation schemes for field theories. In particular, the author thanks M. Bordemann for explaining to him and for critical remarks. Finally clarifying discussions with and valuable comments by I. Kanatchikov are gratefully acknowledged.
## Appendix A Connections and jet bundles
Given a bundle $`\pi :𝒱`$ over an $`n`$-dimensional base manifold $``$ every connection is defined by a section $`\mathrm{\Gamma }`$ of the first jet bundle $`𝔍^1𝒱`$ of $`𝒱`$, since it describes how to lift tangent vectors of the base manifold horizontally. If in addition $`𝒱`$ is a vector bundle (with fibre $`V`$) then as $`𝔍^1𝒱`$ is an affine bundle over $`𝒱`$ the connection $`\mathrm{\Gamma }`$ delivers an isomorphism
$$𝔍^1𝒱\stackrel{\mathrm{\Gamma }}{}𝔙𝒱_{}T^{},$$
(38)
where both sides ($`𝔙𝒱`$ being the vertical bundle to $`𝒱`$) are viewed as bundles over the base manifold $``$. Note in particular that the vertical bundle $`𝔙𝒱`$ is a vector bundle over $``$ (with typical fibre $`V\times V`$, , ch. II, 6.11.).
Now for $`𝒱`$ being a vector bundle we can form the covariant derivative $``$ that corresponds to the given connection $`\mathrm{\Gamma }`$. Then horizontal lifts of tangent vectors are represented by covariantly constant lifts of curves in the base manifold $``$. Therefore, in local coordinates $`(x^i)_{i=1,\mathrm{},n}`$ of $``$ and $`(x^i,v^A)_{i=1,\mathrm{},n,A=1,\mathrm{},N}`$ of $`𝒱`$ the map $`\mathrm{\Gamma }(v)\left(𝔍^1𝒱\right)_v`$, $`v𝒱`$, takes the form
$$\mathrm{\Gamma }(v):(x,\dot{c}^i(x))(x,v,\mathrm{\Gamma }_{iB}^A(x)v^B),$$
(39)
where $`\mathrm{\Gamma }_{iB}^A(x)`$ is the Christoffel symbol of $``$.
Now we are locking for a connection in $`𝔍^1𝒱`$, that is for a map
$$\overline{\mathrm{\Gamma }}:𝔍^1𝒱𝔍^1\left(𝔍^1𝒱\right).$$
For this, one needs a connection both in $`𝒱`$ and $``$ (, Prop. 4). If we use the isomorphisms
$$𝔍^1𝒱\stackrel{\mathrm{\Gamma }}{}𝔙𝒱T^{}\text{and}𝔍^1\left(𝔙𝒱T^{}\right)𝔍^1\left(𝔙𝒱\right)𝔍^1\left(T^{}\right),$$
the latter being natural, we see that all we need is a map $`𝔙𝒱𝔍^1𝔙𝒱`$, since a connection on $``$ defines a map $`\mathrm{\Lambda }^{}:T^{}𝔍^1\left(T^{}\right)`$. Now the desired map can be constructed by vertical prolongation if we make use of the isomorphism $`𝔙𝔍^1𝒱𝔍^1𝔙𝒱`$ (, eq. (1.4))<sup>8)</sup><sup>8)</sup>8)Let $`s_t`$ denote a one-parameter family of local sections of $`\pi :𝒱`$. Then
$$\frac{d}{dt}\text{ }t=0j^1(s_t)(x)j^1(\frac{d}{dt}\text{ }t=0s_t)(x)$$
gives the isomorphism. :
$$𝔙\mathrm{\Gamma }:𝔙𝒱𝔙𝔍^1𝒱𝔍^1𝔙𝒱.$$
Indeed,
$$𝔙\mathrm{\Gamma }\mathrm{\Lambda }^{}:𝔙𝒱T^{}𝔍^1𝔙𝒱𝔍^1T^{}𝔍^1(𝔙𝒱T^{})$$
gives a connection<sup>9)</sup><sup>9)</sup>9) In , p. 136, this construction is denoted by $`p(\mathrm{\Gamma },\mathrm{\Lambda })`$.
$$\overline{\mathrm{\Gamma }}:𝔍^1𝒱𝔍^1\left(𝔍^1𝒱\right).$$
(40)
In coordinates $`(x^i,v^A,v_i^A)`$ of $`𝔍^1𝒱`$ one calculates
$$\overline{\mathrm{\Gamma }}(x^i,v^A,v_i^A):(x^i,\dot{x}^i)(x^i,v^A,v_i^A,\dot{x}^i,\mathrm{\Gamma }_{jB}^A(x)v^B\dot{x}^j,\overline{\mathrm{\Gamma }}_{ij}^A\dot{x}^j),$$
(41)
where
$$\overline{\mathrm{\Gamma }}_{ij}^A(x^i,u^A,u_i^A)=\mathrm{\Gamma }_{jB}^A(u_i^B+\mathrm{\Gamma }_{iC}^Bu^C)\mathrm{\Lambda }_{ji}^k(u_k^A+\mathrm{\Gamma }_{kB}^Au^B)(_j\mathrm{\Gamma }_{iB}^A)u^B+\mathrm{\Gamma }_{iB}^A\mathrm{\Gamma }_{jC}^Bu^C.$$
(42)
Note that $`\mathrm{\Lambda }_{ij}^k`$ denote the Christoffel symbols of $`\mathrm{\Lambda }`$, not $`\mathrm{\Lambda }^{}`$.
## Appendix B A covariant derivative on $`T𝒫`$
Using a connection $`\mathrm{\Gamma }`$ of $`\pi :𝒱`$, which is a map
$$\mathrm{\Gamma }:𝒱𝔍^1𝒱,$$
the affine bundle $`\pi ^{}:𝔍^1𝒱𝒱`$ becomes a vector bundle,
$$𝔍^1𝒱\stackrel{\mathrm{\Gamma }}{}𝔙𝒱_𝒱\pi ^{}\left(T^{}\right),$$
where $`\mathrm{\Gamma }(𝒱)`$ is identified with the zero section.
If in addition $`\pi `$ is a vector bundle, then $`𝔙𝒱`$ is a vector bundle over $``$ as well (, ch. II, 6.11), and we have
$$𝔍^1𝒱\stackrel{\mathrm{\Gamma }}{}𝔙𝒱_{}T^{}.$$
Let $`\overline{𝒱}=𝔙𝒱_{}T^{}`$. In multisymplectic geometry the phase space $`\left(𝔍^1𝒱\right)^{}`$ consists of all with respect to $`\pi ^{}`$ fibre-wise affine mappings from $`𝔍^1𝒱`$ to $`\mathrm{\Lambda }^nT^{}`$. In order to simplify the notation, let us denote this bundle by $`𝒫:=\left(𝔍^1𝒱\right)^{}`$. Again, the connection $`\mathrm{\Gamma }`$ provides an isomorphism
$$𝒫\stackrel{\mathrm{\Gamma }}{}(\overline{𝒱}^{}\mathrm{\Lambda }^nT^{})_𝒱,$$
where $`\stackrel{~}{p}𝒫`$ is decomposed into a linear map $`\overline{p}:\overline{𝒱}\mathrm{\Lambda }^nT^{}`$ and a function $`p`$ on $`𝒱`$ in the following way:
$$\begin{array}{cc}\hfill \stackrel{~}{p}(\stackrel{~}{v})& =\stackrel{~}{p}(\stackrel{~}{v})\stackrel{~}{p}(\mathrm{\Gamma }(\pi ^{}(\stackrel{~}{v})))+\stackrel{~}{p}(\mathrm{\Gamma }(\pi ^{}(\stackrel{~}{v})))\hfill \\ & =\overline{p}(\overline{v})+p(v).\hfill \end{array}$$
Making use of the duality of $`\overline{𝒱}^{}`$ and $`\overline{𝒱}`$ , we obtain a connection $`\overline{\mathrm{\Gamma }}^{}`$ on $`\overline{𝒱}^{}`$ by
$$\overline{\mathrm{\Gamma }}^{}(v),\overline{v}=\overline{p},\overline{\mathrm{\Gamma }}(v),v𝒱,\overline{v}\overline{𝒱}_v,\overline{p}\overline{𝒱}_v^{}.$$
Here, $`\overline{\mathrm{\Gamma }}`$ is the connection on $`\overline{𝒱}`$ as explained in detail in (A). Further, this gives a connection on $`𝒫`$. In coordinates $`(x^i,v^A,p_A^i,p)`$ we calculate
$$\overline{\mathrm{\Gamma }}^{}(\overline{p}):T_x(x^i,\xi ^i)(x^i,\mathrm{\Gamma }_{iB}^Av^B\xi ^i,(\mathrm{\Lambda }_{ji}^kp_A^i\mathrm{\Gamma }_{jA}^Bp_B^k)\xi ^j,0)T𝒫.$$
Now $`\overline{\mathrm{\Gamma }}^{}`$ defines a covariant derivative $`\overline{}`$ on $`𝒫`$. With the help of this we define the connection mapping $`K`$ for $`[\alpha ]_p]T_p𝒫`$, represented by a curve $`\alpha (t)`$, by
$$K:T_p𝒫[\alpha ]_p\{\begin{array}{cc}\frac{d}{dt}\text{ }t=0\alpha (t)\hfill & \text{ if }T\overline{\pi }[\alpha ]=0\hfill \\ \left(\overline{}_{T\overline{\pi }[\alpha ]}\alpha \right)(0)\hfill & \text{ otherwise. }\hfill \end{array}$$
(43)
One easily verifies that $`K`$ is well defined. Let $`p`$ be a point in $`𝒫`$ and $`x`$ its image under the projection $`\overline{\pi }`$. For the tangent mapping of the canonical projection $`\overline{\pi }:𝒫`$, the map $`KT\overline{\pi }:T_p𝒫𝒫_xT_x`$ is bijective and hence provides a splitting of $`T𝒫_p`$. $`X_p^hT_p𝒫`$ is called the horizontal lift of $`HT_x`$ iff $`KT\overline{\pi }(X_p^h)=X`$. Similarly, $`q_p^vT_p𝒫`$ is called the vertical lift of $`q𝒫_x`$ iff $`KT\overline{\pi }(q_p^v)=q`$. Using this we define a covariant derivative $`D`$ on $`T𝒫`$ by<sup>10)</sup><sup>10)</sup>10)This method is inspired by the construction in . :
$$\begin{array}{cc}\hfill D_{X^h}Y^h\text{ }p& =\left(_X^{}Y\right)^h\text{ }p+\frac{1}{2}\left(\overline{R}(X,Y)p\right)^v\text{ }p\hfill \\ \hfill D_{X^h}\beta ^v\text{ }p& =\left(\overline{}_X\beta \right)^v\text{ }p\hfill \\ \hfill D_{\beta ^v}X^h\text{ }p& =0=D_{\beta ^v}\mathrm{\Gamma }^v\text{ }p,\hfill \end{array}$$
(44)
where $`p𝒫`$, $`\beta ^v,\mathrm{\Gamma }^v,X^h,Y^hT𝒫`$ are lifts as above, and $`^{}`$ is the (torsion free) covariant derivative on $`T`$. The curvature term $`\overline{R}`$ of $`\overline{}`$ is needed for $`D`$ to be torsion free.
Since at every point $`p`$ of $`𝒫`$ the tangent space $`T_p𝒫`$ decomposes into the direct sum of horizontal and vertical vectors, we can choose an appropriate basis as follows. If $`(x^i)`$ are coordinates of a neighbourhood $`𝒰`$ of $``$ that trivialises $`𝒫\text{ }𝒰`$ and $`(\xi ^i,v^A,p^{\genfrac{}{}{0.0pt}{}{i}{A}},p)`$ are coordinates on $`𝒫`$, we define for every $`p𝒫`$
$$\begin{array}{cc}\hfill 𝔢_i(p)& =\left(_{x^i}\right)^h\text{ }p=_{\xi ^i}\mathrm{\Gamma }_{iA}^Bv^A_{v^B}+\overline{\mathrm{\Gamma }}_{ij}^A_{p^{{\scriptscriptstyle \genfrac{}{}{0.0pt}{}{j}{A}}}}\hfill \\ \hfill 𝔢_A(p)& =_{v^A},𝔢_{\stackrel{i}{A}}(p)=_{p^{{\scriptscriptstyle \genfrac{}{}{0.0pt}{}{i}{A}}}},𝔢(p)=_p,i=1,\mathrm{},n,A=1,\mathrm{},N.\hfill \end{array}$$
we obtain a basis of $`T_p𝒫`$. From the definition of $`D`$ it follows in particular that
$$D_{𝔢_A}𝔢_\alpha =0,D_{𝔢_{\stackrel{i}{A}}}𝔢_\alpha =0,\alpha =i,A,{\scriptscriptstyle \genfrac{}{}{0.0pt}{}{j}{B}},A,B=1,\mathrm{},N,i,j=1,\mathrm{},n.$$ |
warning/0002/math0002102.html | ar5iv | text | # A 𝑊(𝐸₆)-equivariant projective embedding of the moduli space of cubic surfaces
## 1 The moduli space of cubic surfaces
### 1.1 The space $`M`$ and the action of the group $`\underset{¯}{G}`$
We first fix some notation and recall a few known facts on the moduli space of marked cubic surfaces. The moduli space of marked cubic surfaces, which we denote by $`M`$, is studied for example in and . Since any nonsingular cubic surface can be obtained by blowing up the projective plane $`𝐏^2`$ at six points, it can be represented by a $`3\times 6`$-matrix of which columns give homogeneous coordinates of the six points. In order to get a smooth cubic surface from six points, we assume that no three points are collinear and the six points are not on a conic. On the set of $`3\times 6`$ matrices, we have a cannical action of $`GL_3`$ on the left and the group $`𝐂^\times `$ acts naturally on homogeneous coordinates. By killing such ambiguity of coordinates, we get the following expression
$$x=\left(\begin{array}{cccccc}1& 0& 0& 1& 1& 1\\ 0& 1& 0& 1& x_1& x_2\\ 0& 0& 1& 1& x_3& x_4\end{array}\right);$$
in this paper we use local coordinates $`(x_1,x_2,x_3,x_4)`$ on $`M`$. The six points represented by the matrix above produces a non-singular cubic surface if and only if the following quantitiy does not vanish.
$`D(x)`$ $`:=`$ $`x_1x_2x_3x_4(x_11)(x_21)(x_31)(x_41)`$
$`\times (x_1x_2)(x_1x_3)(x_2x_4)(x_3x_4)D_1D_2Q,`$
where
$`D_1`$ $`:=`$ $`x_1x_4x_2x_3,`$
$`D_2`$ $`:=`$ $`x_1x_4x_4+x_2x_2x_3+x_3x_1,`$
$`Q`$ $`:=`$ $`x_2x_3x_1x_2x_3x_4+x_2x_3+x_1x_4x_2+x_1x_4x_3x_1x_4,`$
Thus we can identify the moduli space $`M`$ with the affine open set $`\{x=(x_1,\mathrm{},x_4)D(x)0\}.`$
Let us define as in six bi-rational transformations $`s_1,\mathrm{},s_6`$ in $`x=(x_1,\mathrm{},x_4):`$
$`s_1:(x_1,x_2,x_3,x_4)`$ $``$ $`({\displaystyle \frac{1}{x_1}},{\displaystyle \frac{1}{x_2}},{\displaystyle \frac{x_3}{x_1}},{\displaystyle \frac{x_4}{x_2}}),`$
$`s_2:(x_1,x_2,x_3,x_4)`$ $``$ $`(x_3,x_4,x_1,x_2),`$
$`s_3:(x_1,x_2,x_3,x_4)`$ $``$ $`({\displaystyle \frac{x_1x_3}{1x_3}},{\displaystyle \frac{x_2x_4}{1x_4}},{\displaystyle \frac{x_3}{x_31}},{\displaystyle \frac{x_4}{x_41}}),`$
$`s_4:(x_1,x_2,x_3,x_4)`$ $``$ $`({\displaystyle \frac{1}{x_1}},{\displaystyle \frac{x_2}{x_1}},{\displaystyle \frac{1}{x_3}},{\displaystyle \frac{x_4}{x_3}}),`$
$`s_5:(x_1,x_2,x_3,x_4)`$ $``$ $`(x_2,x_1,x_4,x_3),`$
$`s_6:(x_1,x_2,x_3,x_4)`$ $``$ $`({\displaystyle \frac{1}{x_1}},{\displaystyle \frac{1}{x_2}},{\displaystyle \frac{1}{x_3}},{\displaystyle \frac{1}{x_4}}).`$
If $`M`$ is regarded as the configuration space of six points in $`𝐏^2`$, the transformation $`s_1`$, for example, corresponds to the interchange of the two points represented by the first two column vectors of the matrix $`x`$. Each $`s_i`$ turns out to be a bi-regular involution on $`M`$, and they form a group $`\underset{¯}{G}`$ isomorphic to the Weyl group of type $`E_6`$; relation of the generators are given by the Coxeter graph
$$\begin{array}{ccccccccc}s_1& \text{}& s_2& \text{}& s_3& \text{}& s_4& \text{}& s_5\\ & & & & & & & & \\ & & & & s_6& & & & \end{array}$$
### 1.2 Root system $`\mathrm{\Delta }`$ of type $`E_6`$
We review the root system $`\mathrm{\Delta }`$ of type $`E_6`$, following . Consider an 8-dimensional Euclidean space $`\stackrel{~}{E}`$ with a standard basis $`\epsilon _1,\mathrm{},\epsilon _8`$. Let $`,`$ be the inner product on $`\stackrel{~}{E}`$ defined by $`\epsilon _j,\epsilon _k=\delta _{jk}`$ and let $`E`$ be the linear subspace of $`\stackrel{~}{E}`$ spanned by the six vectors
$$\epsilon _1,\mathrm{},\epsilon _5,\stackrel{~}{\epsilon }=\epsilon _6\epsilon _7\epsilon _8.$$
We introduce the 36 vectors:
$$r=\frac{1}{2}(\epsilon _1+\epsilon _2+\epsilon _3+\epsilon _4+\epsilon _5+\stackrel{~}{\epsilon }),$$
$$r_{1j}=\epsilon _{j1}+r_0,2j6$$
$$r_{jk}=\epsilon _{j1}\epsilon _{k1},2j<k6$$
$$r_{1jk}=\epsilon _{j1}\epsilon _{k1},2j<k6$$
$$r_{ijk}=\epsilon _{i1}\epsilon _{j1}\epsilon _{k1}+r_0,2i<j<k6$$
where $`r_{ij}=r_{ji},r_{ijk}=r_{jik}=r_{ikj},`$
$$r_0=\frac{1}{2}(\epsilon _1+\epsilon _2+\epsilon _3+\epsilon _4+\epsilon _5\stackrel{~}{\epsilon }).$$
Note that
$$rr_{ij},r_{ij}r_{kl},r_{ij}r_{ijk},r_{ij}r_{klm},r_{ijk}r_{ilm}.$$
The set
$$\mathrm{\Delta }=\{\pm r,\pm r_{ij},\pm r_{ijk}\}$$
forms a root system of type $`E_6`$. For example,
$$r_{12},r_{123},r_{23},r_{34},r_{45},r_{56}$$
can serve as a system of positive simple roots; its extended Dynkin diagram is given as
| $`r_{12}`$ | —— | $`r_{23}`$ | —— | $`r_{34}`$ | —— | $`r_{45}`$ | —— | $`r_{56}`$ |
| --- | --- | --- | --- | --- | --- | --- | --- | --- |
| | | | | $``$ | | | | |
| | | | | $`r_{123}`$ | | | | |
| | | | | $``$ | | | | |
| | | | | $`r`$ | | | | |
The set $`\{r`$, $`r_{jk}`$, $`r_{ijk}\}`$ is the totality of positive roots of $`\mathrm{\Delta }`$.
Let $`s_r,s_{ij}`$ and $`s_{ijk}`$ be the reflections on $`E`$ with respect to $`r,r_{ij}`$ and $`r_{ijk}`$. These reflections act on $`\mathrm{\Delta }`$ as
$$\begin{array}{cc}s_r:\hfill & r_{ij}r_{ij},r_{ijk}r_{lmn},\{i,j,k,l,m,n\}=\{1,\mathrm{},6\},\hfill \\ s_{ij}:\hfill & \mathrm{permutation}\mathrm{of}\mathrm{the}\mathrm{indices}i\mathrm{and}j,\hfill \end{array}$$
$$\begin{array}{cccccccccc}s_{123}:\hfill & r_{12}\hfill & \hfill & r_{12},\hfill & r_{14}\hfill & \hfill & r_{234},\hfill & r_{56}\hfill & \hfill & r_{56},\hfill \\ & r_{145}\hfill & \hfill & r_{145},\hfill & r\hfill & \hfill & r_{456},\hfill & & & \end{array}$$
modulo signs. Let us define two reflection group
$$G_1=s_{12},s_{23},s_{34},s_{45},s_{56}S_6,G=G_1,s_{123}W(E_6),$$
where $`S_6`$ is the symmetric group on six numerals $`\{1,\mathrm{},6\}`$, $`W(E_6)`$ the Weyl group of type $`E_6`$, and $`a,b,\mathrm{}`$ denotes the group generated by $`a,b,\mathrm{}`$ Note that
$$G_1G_1,s_r=G_1\times s_rG,$$
and $`G`$ acts transitively on $`\mathrm{\Delta }`$.
### 1.3 Naruki’s cross-ratio variety
A smooth compactification of $`M`$ known as Naruki’s cross-ratio variety $`𝒞`$ (, ), embedded in $`(𝐏^1)^{45}`$, is the union of $`M`$ and the 76 divisors. The 36 of them correspond to the positive roots of $`\mathrm{\Delta }`$ (they are said to be of the first kind), the other 40 divisors (said to be of the second kind, and are isomorphic to $`(𝐏^1)^3`$) can be represented as follows: Take three subsets $`\mathrm{\Delta }_1`$, $`\mathrm{\Delta }_2`$, $`\mathrm{\Delta }_3`$ of $`\mathrm{\Delta }`$ satisfying the following conditions:
* Each of $`\mathrm{\Delta }_1`$, $`\mathrm{\Delta }_2`$, $`\mathrm{\Delta }_3`$ is a root system of type $`A_2`$.
* $`\mathrm{\Delta }_1`$, $`\mathrm{\Delta }_2`$, $`\mathrm{\Delta }_3`$ are mutually orthogonal.
* The vectors in $`\mathrm{\Delta }_1\mathrm{\Delta }_2\mathrm{\Delta }_3`$ span $`E`$.
Note that each one of such three root systems determines the other two.
Such a triple $`\{\mathrm{\Delta }_1,\mathrm{\Delta }_2,\mathrm{\Delta }_3\}`$ dertermines a divisor. According to the naming of the roots, we must use two different expressions: The first one is of the form
$$\{\pm r_{12},\pm r_{23},\pm r_{13}\},\{\pm r_{45},\pm r_{56},\pm r_{46}\},\{\pm r,\pm r_{123},\pm r_{456}\},$$
(the corresponding divisor is denoted by $`Z_{123,456}`$ in ,) and the second one is of the form
$$\{\pm r_{12},\pm r_{234},\pm r_{134}\},\{\pm r_{34},\pm r_{356},\pm r_{456}\},\{\pm r_{56},\pm r_{125},\pm r_{126}\},$$
(the corresponding divisor is denoted by $`Z_{12,34,56}`$ in ). Note that
$$Z_{123,456}=Z_{456,123},Z_{12,34,56}=Z_{56,12,34}Z_{12,56,34}.$$
Thus, permuting the indices under $`G_1`$, we have 40 ($`=10+30`$) such divisors. The group $`G`$ acts transitively on these 40 divisors.
Remark: These divisors (of the second kind) are disjoint to each other, and can be blown-down to points. In fact they correspond bijectively to the cusps of the modular group studied in (see also and ).
## 2 Embedding $`\phi :M𝐏^{801}`$
### 2.1 Coordinates on $`𝐏^{801}`$ and the action of $`G`$
Let $`𝒜`$ be the set of 40 labels $`(123,456)`$ and $`(12,34,56)`$ with the following identification
$`(123,456)`$ $`=`$ $`(213,456)=(132,456)=(456,123),`$
$`(12,34,56)`$ $`=`$ $`(21,34,56)=(56,12,34).`$
Since $`G`$ acts on the set $`\mathrm{\Delta }`$ of roots, it acts also on the set of 40 divisors above, and so that it also acts on the set $`𝒜`$ of 40 labels.
We introduce 80 homogeneous coordinates
$$y_\alpha ,y_\alpha ,\alpha 𝒜$$
on $`𝐏^{801}`$. I define an action of $`G`$ on $`𝐏^{801}`$ by the following action of the generators $`s_{12},\mathrm{},s_{56}`$ and $`s_{123}`$ on the coordinates. Let $`s`$ be one of the generators and $`\alpha 𝒜`$; we assign
$`s(y_\alpha )`$ $`=`$ $`y_\alpha ,s(y_\alpha )=y_\alpha \text{if }s\alpha =\alpha \text{,}`$
$`s(y_\alpha )`$ $`=`$ $`y_\beta ,s(y_\alpha )=y_\beta \text{if }s\alpha =\beta \alpha \text{.}`$
### 2.2 Definition of $`\phi `$
In this section we define a map $`M𝐏^{801}`$. For a $`3\times 6`$ matrix $`x=(x_{ij})`$, we consider 80 polynomials of degree 18 as follows:
$`y_{(123,456)}(x)`$ $`=`$ $`D_{123}(x)D_{456}(x)Q(x),`$
$`y_{(12,34,56)}(x)`$ $`=`$ $`D_{134}(x)D_{234}(x)D_{356}(x)D_{456}(x)D_{512}(x)D_{612}(x),`$
and $`y_\alpha (x)=y_\alpha (\alpha 𝒜)`$, where $`Q(x)`$ is the determinant of the $`6\times 6`$-matrix with columns
$$(x_{1j}x_{2j},x_{2j}x_{3j},x_{3j}x_{1j},x_{1j}^2,x_{2j}^2,x_{3j}^2)j=1,\mathrm{},6.$$
Since we have
$$y_\alpha (gxh)=(detg)^6y_\alpha (x)(deth)^3,$$
the correswpondence above defines a map $`\phi :M𝐏^{801}`$. For later use, we present 40 polynomials $`y_\alpha (x)`$ in terms of the coordinates $`(x_1,x_2,x_3,x_4)`$ introduced in §1; the remaining 40 polynomoals are given by $`y_\alpha (x)=y_\alpha (x)`$. In the following table, $`y_\alpha (x)`$ is denoted simply by $`\alpha `$, and we number them as $`y_1,\mathrm{},y_{40}`$:
$`y_1=(156,234)`$ $`:=`$ $`D_1Q:y_2=(123,456):=D_2Q:`$
$`y_3=(124,356)`$ $`:=`$ $`(x_2x_1)Q:y_4=(145,236):=(x_3x_1)Q:`$
$`y_5=(146,235)`$ $`:=`$ $`(x_4x_2)Q:y_6=(134,256):=(x_4x_3)Q:`$
$`y_7=(135,246)`$ $`:=`$ $`x_1(x_41)Q:y_8=(136,245):=x_2(x_31)Q:`$
$`y_9=(125,346)`$ $`:=`$ $`x_3(x_21)Q:y_{10}=(126,345):=x_4(x_11)Q:`$
$`y_{11}=(12,56,34)`$ $`:=`$ $`D_1(x_11)(x_21)(x_4x_3):`$
$`y_{12}=(16,23,45)`$ $`:=`$ $`D_1(x_11)(x_31)(x_4x_2):`$
$`y_{13}=(15,23,46)`$ $`:=`$ $`D_1(x_21)(x_41)(x_3x_1):`$
$`y_{14}=(13,56,24)`$ $`:=`$ $`D_1(x_31)(x_41)(x_2x_1):`$
$`y_{15}=(15,24,36)`$ $`:=`$ $`D_1x_1(x_21)(x_31):`$
$`y_{16}=(16,24,35)`$ $`:=`$ $`D_1x_2(x_11)(x_41):`$
$`y_{17}=(15,34,26)`$ $`:=`$ $`D_1x_3(x_11)(x_41):`$
$`y_{18}=(16,34,25)`$ $`:=`$ $`D_1x_4(x_21)(x_31):`$
$`y_{19}=(12,36,45)`$ $`:=`$ $`D_2x_2x_3(x_11):`$
$`y_{20}=(12,35,46)`$ $`:=`$ $`D_2x_1x_4(x_21):`$
$`y_{21}=(13,26,45)`$ $`:=`$ $`D_2x_1x_4(x_31):`$
$`y_{22}=(13,25,46)`$ $`:=`$ $`D_2x_2x_3(x_41):`$
$`y_{23}=(13,24,56)`$ $`:=`$ $`D_2x_1x_2(x_4x_3):`$
$`y_{24}=(15,46,23)`$ $`:=`$ $`D_2x_1x_3(x_4x_2):`$
$`y_{25}=(16,45,23)`$ $`:=`$ $`D_2x_2x_4(x_3x_1):`$
$`y_{26}=(12,34,56)`$ $`:=`$ $`D_2x_3x_4(x_2x_1):`$
$`y_{27}=(13,46,25)`$ $`:=`$ $`x_1(x_21)(x_31)(x_4x_2)(x_4x_3):`$
$`y_{28}=(13,45,26)`$ $`:=`$ $`x_2(x_11)(x_41)(x_3x_1)(x_4x_3):`$
$`y_{29}=(12,46,35)`$ $`:=`$ $`x_3(x_11)(x_41)(x_2x_1)(x_4x_2):`$
$`y_{30}=(12,45,36)`$ $`:=`$ $`x_4(x_21)(x_31)(x_2x_1)(x_3x_1):`$
$`y_{31}=(14,35,26)`$ $`:=`$ $`x_1(x_11)(x_4x_2)(x_4x_3):`$
$`y_{32}=(14,36,25)`$ $`:=`$ $`x_2(x_21)(x_3x_1)(x_4x_3):`$
$`y_{33}=(14,25,36)`$ $`:=`$ $`x_3(x_31)(x_2x_1)(x_4x_2):`$
$`y_{34}=(14,26,35)`$ $`:=`$ $`x_4(x_41)(x_2x_1)(x_3x_1):`$
$`y_{35}=(16,25,34)`$ $`:=`$ $`x_2x_3(x_11)(x_4x_2)(x_4x_3):`$
$`y_{36}=(15,26,34)`$ $`:=`$ $`x_1x_4(x_21)(x_3x_1)(x_4x_3):`$
$`y_{37}=(16,35,24)`$ $`:=`$ $`x_1x_4(x_31)(x_2x_1)(x_4x_2):`$
$`y_{38}=(15,36,24)`$ $`:=`$ $`x_2x_3(x_41)(x_2x_1)(x_3x_1):`$
$`y_{39}=(14,56,23)`$ $`:=`$ $`D_1D_2:`$
$`y_{40}=(14,23,56)`$ $`:=`$ $`(x_2x_1)(x_3x_1)(x_4x_2)(x_4x_3):`$
### 2.3 $`G`$-Equivariance of $`\phi `$
Recall that the group $`\underset{¯}{G}`$ acts on $`M`$, and that $`G`$ acts on $`𝐏^{801}`$. Let us identify the groups $`\underset{¯}{G}`$ and $`G`$ by
$$\iota :s_{12}s_1,\mathrm{},s_{56}s_5,s_{123}s_6.$$
Then we have
###### Theorem 1
The map $`\phi :M𝐏^{801}`$ is $`G`$-equivariant:
$$g(\phi (x))=\phi (\iota (g)x),gG,xM,$$
that is,
$$(gy_\alpha )(x)=c_gy_\alpha (\iota (g)x),gG,\alpha \pm 𝒜,xM,$$
where $`c_g`$ is a rational function in $`(x_1,x_2,x_3,x_4)`$.
Convention: Once this theorem is established, we ignore the redundant ones $`y_\alpha (x)=y_\alpha (x)`$ and regard $`\phi `$ as the map
$$Mx:y_\alpha (x):𝐏^{401}.$$
The group $`G`$ still acts on $`𝐏^{401}`$ by the transformations given in§2.3.
In order to prove the theorem, we have only to check the identity for a set of gnerators of $`G`$. Under $`s_1`$, the fourty polynomials are transformed as follows:
$$\begin{array}{ccccc}y_1c_1y_6,\hfill & y_2c_1y_2,\hfill & y_3c_1y_3,\hfill & y_4c_1y_8,\hfill & y_5c_1y_7,\hfill \\ y_6c_1y_1,\hfill & y_7c_1y_5,\hfill & y_8c_1y_4,\hfill & y_9c_1y_9,\hfill & y_{10}c_1y_{10},\hfill \\ y_{11}c_1y_{11},\hfill & y_{12}c_1y_{28},\hfill & y_{13}c_1y_{27},\hfill & y_{14}c_1y_{40},\hfill & y_{15}c_1y_{32},\hfill \\ y_{16}c_1y_{31},\hfill & y_{17}c_1y_{35},\hfill & y_{18}c_1y_{36},\hfill & y_{19}c_1y_{19},\hfill & y_{20}c_1y_{20},\hfill \\ y_{21}c_1y_{25},\hfill & y_{22}c_1y_{24},\hfill & y_{23}c_1y_{39},\hfill & y_{24}c_1y_{22},\hfill & y_{25}c_1y_{21},\hfill \\ y_{26}c_1y_{26},\hfill & y_{27}c_1y_{13},\hfill & y_{28}c_1y_{12},\hfill & y_{29}c_1y_{29},\hfill & y_{30}c_1y_{30},\hfill \\ y_{31}c_1y_{16},\hfill & y_{32}c_1y_{15},\hfill & y_{33}c_1y_{38},\hfill & y_{34}c_1y_{37},\hfill & y_{35}c_1y_{17},\hfill \\ y_{36}c_1y_{18},\hfill & y_{37}c_1y_{34},\hfill & y_{38}c_1y_{33},\hfill & y_{39}c_1y_{23},\hfill & y_{40}c_1y_{14},\hfill \end{array}$$
where $`c_1=(x_1x_2)^3`$, under $`s_2`$,
$$\begin{array}{ccccc}y_1y_1,\hfill & y_2y_2,\hfill & y_3y_6,\hfill & y_4y_4,\hfill & y_5y_5,\hfill \\ y_6y_3,\hfill & y_7y_9,\hfill & y_8y_{10},\hfill & y_9y_7,\hfill & y_{10}y_8,\hfill \\ y_{11}y_{14},\hfill & y_{12}y_{12},\hfill & y_{13}y_{13},\hfill & y_{14}y_{11},\hfill & y_{15}y_{17},\hfill \\ y_{16}y_{18},\hfill & y_{17}y_{15},\hfill & y_{18}y_{16},\hfill & y_{19}y_{21},\hfill & y_{20}y_{22},\hfill \\ y_{21}y_{19},\hfill & y_{22}y_{20},\hfill & y_{23}y_{26},\hfill & y_{24}y_{24},\hfill & y_{25}y_{25},\hfill \\ y_{26}y_{23},\hfill & y_{27}y_{29},\hfill & y_{28}y_{30},\hfill & y_{29}y_{27},\hfill & y_{30}y_{28},\hfill \\ y_{31}y_{33},\hfill & y_{32}y_{34},\hfill & y_{33}y_{31},\hfill & y_{34}y_{32},\hfill & y_{35}y_{37},\hfill \\ y_{36}y_{38},\hfill & y_{37}y_{35},\hfill & y_{38}y_{36},\hfill & y_{39}y_{39},\hfill & y_{40}y_{40},\hfill \end{array}$$
under $`s_3`$,
$$\begin{array}{ccccc}y_1c_3y_1,\hfill & y_2c_3y_3,\hfill & y_3c_3y_2,\hfill & y_4c_3y_7,\hfill & y_5c_3y_8,\hfill \\ y_6c_3y_6,\hfill & y_7c_3y_4,\hfill & y_8c_3y_5,\hfill & y_9c_3y_9,\hfill & y_{10}c_3y_{10},\hfill \\ y_{11}c_3y_{11},\hfill & y_{12}c_3y_{16},\hfill & y_{13}c_3y_{15},\hfill & y_{14}c_3y_{39},\hfill & y_{15}c_3y_{13},\hfill \\ y_{16}c_3y_{12},\hfill & y_{17}c_3y_{17},\hfill & y_{18}c_3y_{18},\hfill & y_{19}c_3y_{29},\hfill & y_{20}c_3y_{30},\hfill \\ y_{21}c_3y_{34},\hfill & y_{22}c_3y_{33},\hfill & y_{23}c_3y_{40},\hfill & y_{24}c_3y_{38},\hfill & y_{25}c_3y_{37},\hfill \\ y_{26}c_3y_{26},\hfill & y_{27}c_3y_{32},\hfill & y_{28}c_3y_{31},\hfill & y_{29}c_3y_{19},\hfill & y_{30}c_3y_{20},\hfill \\ y_{31}c_3y_{28},\hfill & y_{32}c_3y_{27},\hfill & y_{33}c_3y_{22},\hfill & y_{34}c_3y_{21},\hfill & y_{35}c_3y_{35},\hfill \\ y_{36}c_3y_{36},\hfill & y_{37}c_3y_{25},\hfill & y_{38}c_3y_{24},\hfill & y_{39}c_3y_{14},\hfill & y_{40}c_3y_{23},\hfill \end{array}$$
where $`c_3=(1x_3)^3(1x_4)^3`$, under $`s_4`$,
$$\begin{array}{ccccc}y_1c_4y_5,\hfill & y_2c_4y_2,\hfill & y_3c_4y_9,\hfill & y_4c_4y_4,\hfill & y_5c_4y_1,\hfill \\ y_6c_4y_7,\hfill & y_7c_4y_6,\hfill & y_8c_4y_8,\hfill & y_9c_4y_3,\hfill & y_{10}c_4y_{10},\hfill \\ y_{11}c_4y_{29},\hfill & y_{12}c_4y_{12},\hfill & y_{13}c_4y_{40},\hfill & y_{14}c_4y_{27},\hfill & y_{15}c_4y_{33},\hfill \\ y_{16}c_4y_{35},\hfill & y_{17}c_4y_{31},\hfill & y_{18}c_4y_{37},\hfill & y_{19}c_4y_{19},\hfill & y_{20}c_4y_{26},\hfill \\ y_{21}c_4y_{21},\hfill & y_{22}c_4y_{23},\hfill & y_{23}c_4y_{22},\hfill & y_{24}c_4y_{39},\hfill & y_{25}c_4y_{25},\hfill \\ y_{26}c_4y_{20},\hfill & y_{27}c_4y_{14},\hfill & y_{28}c_4y_{28},\hfill & y_{29}c_4y_{11},\hfill & y_{30}c_4y_{30},\hfill \\ y_{31}c_4y_{17},\hfill & y_{32}c_4y_{38},\hfill & y_{33}c_4y_{15},\hfill & y_{34}c_4y_{36},\hfill & y_{35}c_4y_{16},\hfill \\ y_{36}c_4y_{34},\hfill & y_{37}c_4y_{18},\hfill & y_{38}c_4y_{32},\hfill & y_{39}c_4y_{24},\hfill & y_{40}c_4y_{13},\hfill \end{array}$$
where $`c_4=(x_1x_3)^3`$, under $`s_5`$,
$$\begin{array}{ccccc}y_1y_1,\hfill & y_2y_2,\hfill & y_3y_3,\hfill & y_4y_5,\hfill & y_5y_4,\hfill \\ y_6y_6,\hfill & y_7y_8,\hfill & y_8y_7,\hfill & y_9y_{10},\hfill & y_{10}y_9,\hfill \\ y_{11}y_{11},\hfill & y_{12}y_{13},\hfill & y_{13}y_{12},\hfill & y_{14}y_{14},\hfill & y_{15}y_{16},\hfill \\ y_{16}y_{15},\hfill & y_{17}y_{18},\hfill & y_{18}y_{17},\hfill & y_{19}y_{20},\hfill & y_{20}y_{19},\hfill \\ y_{21}y_{22},\hfill & y_{22}y_{21},\hfill & y_{23}y_{23},\hfill & y_{24}y_{25},\hfill & y_{25}y_{24},\hfill \\ y_{26}y_{26},\hfill & y_{27}y_{28},\hfill & y_{28}y_{27},\hfill & y_{29}y_{30},\hfill & y_{30}y_{29},\hfill \\ y_{31}y_{32},\hfill & y_{32}y_{31},\hfill & y_{33}y_{34},\hfill & y_{34}y_{33},\hfill & y_{35}y_{36},\hfill \\ y_{36}y_{35},\hfill & y_{37}y_{38},\hfill & y_{38}y_{37},\hfill & y_{39}y_{39},\hfill & y_{40}y_{40},\hfill \end{array}$$
and under $`s_6`$,
$$\begin{array}{ccccc}y_1c_6y_{39},\hfill & y_2c_6y_2,\hfill & y_3c_6y_{26},\hfill & y_4c_6y_{25},\hfill & y_5c_6y_{24},\hfill \\ y_6c_6y_{23},\hfill & y_7c_6y_{22},\hfill & y_8c_6y_{21},\hfill & y_9c_6y_{20},\hfill & y_{10}c_6y_{19},\hfill \\ y_{11}c_6y_{11},\hfill & y_{12}c_6y_{12},\hfill & y_{13}c_6y_{13},\hfill & y_{14}c_6y_{14},\hfill & y_{15}c_6y_{18},\hfill \\ y_{16}c_6y_{17},\hfill & y_{17}c_6y_{16},\hfill & y_{18}c_6y_{15},\hfill & y_{19}c_6y_{10},\hfill & y_{20}c_6y_9,\hfill \\ y_{21}c_6y_8,\hfill & y_{22}c_6y_7,\hfill & y_{23}c_6y_6,\hfill & y_{24}c_6y_5,\hfill & y_{25}c_6y_4,\hfill \\ y_{26}c_6y_3,\hfill & y_{27}c_6y_{27},\hfill & y_{28}c_6y_{28},\hfill & y_{29}c_6y_{29},\hfill & y_{30}c_6y_{30},\hfill \\ y_{31}c_6y_{35},\hfill & y_{32}c_6y_{36},\hfill & y_{33}c_6y_{37},\hfill & y_{34}c_6y_{38},\hfill & y_{35}c_6y_{31},\hfill \\ y_{36}c_6y_{32},\hfill & y_{37}c_6y_{33},\hfill & y_{38}c_6y_{34},\hfill & y_{39}c_6y_1,\hfill & y_{40}c_6y_{40},\hfill \end{array}$$
where $`c_6=(x_1x_2x_3x_4)^2`$. Maybe it is interesting to see what happens under the operation of the involution $`s_r`$ (classically called the association) which sends $`(x_1,x_2,x_3,x_4)`$ to
$$(\frac{(x_41)D_1}{(x_4x_2)(x_4x_3)},\frac{(x_31)D_1}{(x_3x_1)(x_4x_3)},\frac{(x_21)D_1}{(x_4x_2)(x_2x_1)},\frac{(x_11)D_1}{(x_3x_1)/(x_2x_1)}):$$
$$\begin{array}{ccccc}y_1c_ry_1,\hfill & y_2c_ry_2,\hfill & y_3c_ry_3,\hfill & y_4c_ry_4,\hfill & y_5c_ry_5,\hfill \\ y_6c_ry_6,\hfill & y_7c_ry_7,\hfill & y_8c_ry_8,\hfill & y_9c_ry_9,\hfill & y_{10}c_ry_{10},\hfill \\ y_{11}c_ry_{26},\hfill & y_{12}c_ry_{25},\hfill & y_{13}c_ry_{24},\hfill & y_{14}c_ry_{23},\hfill & y_{15}c_ry_{38},\hfill \\ y_{16}c_ry_{37},\hfill & y_{17}c_ry_{36},\hfill & y_{18}c_ry_{35},\hfill & y_{19}c_ry_{30},\hfill & y_{20}c_ry_{29},\hfill \\ y_{21}c_ry_{28},\hfill & y_{22}c_ry_{27},\hfill & y_{23}c_ry_{14},\hfill & y_{24}c_ry_{13},\hfill & y_{25}c_ry_{12},\hfill \\ y_{26}c_ry_{11},\hfill & y_{27}c_ry_{22},\hfill & y_{28}c_ry_{21},\hfill & y_{29}c_ry_{20},\hfill & y_{30}c_ry_{19},\hfill \\ y_{31}c_ry_{34},\hfill & y_{32}c_ry_{33},\hfill & y_{33}c_ry_{32},\hfill & y_{34}c_ry_{31},\hfill & y_{35}c_ry_{18},\hfill \\ y_{36}c_ry_{17},\hfill & y_{37}c_ry_{16},\hfill & y_{38}c_ry_{15},\hfill & y_{39}c_ry_{40},\hfill & y_{40}c_ry_{39},\hfill \end{array}$$
where
$$c_r=\left(\frac{D_1D_2}{(x_4+x_2)(x_4x_3)(x_1x_3)(x_1x_2)}\right)^3.$$
### 2.4 $`\phi `$ embeds $`M`$
It is known in that the map
$$Mxy_1(x):y_3(x):y_4(x):y_5(x):y_7(x)𝐏^4$$
is two-to-one, and induces an embedding of the quotient space $`M/s_r`$. Thus the composite of $`\phi `$ and the projection
$$M𝐏^{401}𝐏^4$$
is a two-to-one map. This fact together with the equivariance of $`\phi `$ under the involution $`s_r`$ shown just above implies
###### Theorem 2
$`\phi `$ embeds $`M`$ into $`𝐏^{401}`$.
### 2.5 Prolongation of $`\phi `$ to degenerate arrangements
Let us consider degenerate arrangements of six points on the plane. Since arrangements with three collinear points can be transformed under $`G`$ to those with six points on a conic, we assume, Without loss of generality, that our arrangements represented by $`x=(x_1,\mathrm{},x_4)`$ satisfies $`Q=0`$, that is, the six points are on a conic. Since a (nonsingular) conic is isomorphic to a line, such arrangements form the configuration space
$$X(2,6)=GL(2)\backslash \{Mat(2,6)\text{any }2\times 2\text{ minor}0\}/(𝐂^\times )^6$$
of six points on the projective line: if we represent a point of $`X(2,6)`$ by a matrix of the form
$$z=\left(\begin{array}{cccccc}1& 0& 1& 1& 1& 1\\ 0& 1& 1& z_1& z_2& z_3\end{array}\right),$$
where
$$\underset{i=1}{\overset{3}{}}z_i(z_i1)\underset{1i<j3}{}(z_iz_j)0,$$
then the degenerate arrangements in question can be parametrized by $`z=(z_1,z_2,z_3)`$ as
$$x_1=(1z_1)/(1z_2),x_2=(1z_1)/(1z_3),x_3=z_1/z_2,x_4=z_1/z_3.$$
Among the two $`S_6`$-equivariant projective embedding of $`X(2,6)`$ presented in , let us recall the following one given by the fifteen polynomials
$$D_{ij}(z)D_{kl}(z)D_{mn}(z),\{i,j,k,l,m\}=\{1,\mathrm{},6\},$$
where $`D_{ij}(z)`$ is the $`(i,j)`$-minor of the $`2\times 6`$-matrix $`z`$. Their actual forms are given by
$$(z_i1)(z_jz_k),z_jz_k,z_i(z_jz_k),z_i(z_j1).$$
It is known and easy to show that the image is projectively equivalent to the so-called Segre cubic defined by
$$t_0+\mathrm{},t_5=0,(t_0)^3+\mathrm{}+(t_5)^3=0.$$
On the other hand, let us prolong the domain of definition of the map $`\phi `$ on these degenerate arrangements by the same forty polynomials. Then the map $`\phi `$ in $`z`$-coodinates is given by $`y_1=\mathrm{}y_{10}=0`$ and
$$\begin{array}{ccc}cy_{11}=z_1(z_2z_3),\hfill & cy_{12}=z_2+z_1,\hfill & cy_{13}=z_1z_3,\hfill \\ cy_{14}=(1+z_1)(z_2z_3),\hfill & cy_{15}=(1+z_1)z_3,\hfill & cy_{16}=(1+z_1)z_2,\hfill \\ cy_{17}=z_1(1+z_3),\hfill & cy_{18}=z_1(1+z_2),\hfill & cy_{19}=(z_2+z_1)z_3,\hfill \\ cy_{20}=(z_1z_3)z_2,\hfill & cy_{21}=(z_2+z_1)(1+z_3),\hfill & cy_{22}=(z_1z_3)(1+z_2),\hfill \\ cy_{23}=(1+z_1)(z_2z_3),\hfill & cy_{24}=z_1z_3,\hfill & cy_{25}=z_2+z_1,\hfill \\ cy_{26}=z_1(z_2z_3),\hfill & cy_{27}=(z_1z_3)(1+z_2),\hfill & cy_{28}=(z_2+z_1)(1+z_3),\hfill \\ cy_{29}=(z_1z_3)z_2,\hfill & cy_{30}=(z_2+z_1)z_3,\hfill & cy_{31}=(1+z_3)z_2,\hfill \\ cy_{32}=(1+z_2)z_3,\hfill & cy_{33}=(1+z_2)z_3,\hfill & cy_{34}=(1+z_3)z_2,\hfill \\ cy_{35}=z_1(1+z_2),\hfill & cy_{36}=z_1(1+z_3),\hfill & cy_{37}=(1+z_1)z_2,\hfill \\ cy_{38}=(1+z_1)z_3,\hfill & cy_{39}=z_2z_3,\hfill & cy_{40}=z_2z_3,\hfill \end{array}$$
where
$$c=\frac{(z_11)(z_1z_3)(z_2z_3)(z_1z_2)z_1}{(1z_2)^2z_3^2(1z_3)^2z_2^2}.$$
This shows that the prolonged $`\phi `$ gives exactly the embedding of the arranged arragements given by $`\{xQ(x)=0\}`$, isomorphic to $`X(2,6)`$, given above.
Further put
$$x_1=t\xi _1,x_2=t\xi _2,x_3=t\xi _3,x_4=1+t\xi _4,$$
and let $`t`$ tends to zero in $`\phi (x)`$. We can easily see that the limit is a point whchi is independ of $`(\xi _1,\mathrm{},\xi _4)`$.
These results together with the facts on Naruki’s cross-ratio $`𝒞`$ variety reviewd in §1.3, we can readily show
###### Theorem 3
The closure of the image of $`M`$ under $`\phi `$ is isomorphic to the variety obtained from Naruki’s cross-ratio $`𝒞`$ by blowing down the 40 exceptional divisors of the second kind to points.
Remark: This variety is isomorphic to the Stake compactification of the modular variety obtained in . Note that the 40 cusps correspond to the 40 points obtained by blowing down the divisors of the second kind.
## 3 Defining equations
We will find a set of generators of the ideal defining the closure $`\overline{\phi (M)}`$ of the image of $`M`$ in $`𝐏^{401}.`$ Among the forty polynomials $`y_1(x),\mathrm{},y_{40}(x)`$, we can find some relations. The Prücker relations () of the $`3\times 3`$-minors of a $`3\times 6`$-matrix yield linear relations among $`y_1(x),\mathrm{},y_{10}(x)`$; for example,
$$(124,356)(145,236)+(146,235)(134,256)=y_3(x)y_4(x)+y_5(x)y_6(x)=0.$$
On the other hand, it is easy to check the cubic relation
$`(124,356)(15,23,46)(13,26,45)(145,236)(13,56,24)(12,35,46)`$
$`=y_3(x)y_{13}(x)y_{21}(x)y_4(x)y_{14}(x)y_{20}(x)=0.`$
Let $`V`$ be the subvariety of $`𝐏^{401}`$, coordinatized by $`y_1,\mathrm{},y_{40}`$, defined by the $`G`$-orbits of the linear and the cubic equations:
$$y_3y_4+y_5y_6=0,y_3y_{13}y_{21}y_4y_{14}y_{20}=0.$$
###### Theorem 4
$`\overline{\phi (M)}=V.`$
### 3.1 An outline of the proof
By operating the group $`G`$ to the linear equation, we get many linear relations among $`y_1,\mathrm{},y_{40}`$. These relations form a system of rank 30, that is, all $`y`$’s can be expressed linearly in terms of ten chosen ones. For example, in terms of
$$y_1,y_3,y_4,y_5,y_7,y_{11},y_{12},y_{13},y_{15},y_{19},$$
the remaining thirty ones are expressed as
$`y_2`$ $`=`$ $`y_1y_5+y_4,y_6=y_4+y_5+y_3,y_8=y_3y_1+y_7,`$
$`y_9`$ $`=`$ $`y_1+y_7y_4,y_{10}=y_7y_5y_3,y_{14}=y_{11}+y_{13}y_{12},`$
$`y_{16}`$ $`=`$ $`y_1+y_{12}y_{13}y_{11}+y_{15},y_{17}=y_{15}y_1y_{13},`$
$`y_{18}`$ $`=`$ $`y_{13}y_{11}+y_{15},y_{20}=y_{19}+y_3+y_{11},y_{21}=y_4+y_{19}+y_{12},`$
$`y_{22}`$ $`=`$ $`y_{19}+y_3+y_{11}+y_{13}+y_5,y_{23}=y_1y_{12}+y_{13}+y_3+y_{11},`$
$`y_{24}`$ $`=`$ $`y_4+y_1+y_{13},y_{25}=y_{12}+y_5y_1,`$
$`y_{26}`$ $`=`$ $`y_4+y_5y_1+y_3+y_{11},y_{27}=y_{19}+y_4+y_1+y_3y_7+y_{11}+y_{13},`$
$`y_{28}`$ $`=`$ $`y_7+y_{19}+y_{12}+y_5+y_3,y_{29}=y_3y_7+y_{11}+y_{19}+y_5,`$
$`y_{30}`$ $`=`$ $`y_3+y_1y_7+y_{19}+y_4,y_{31}=y_1y_{13}+y_{19}+y_{15}+y_{12},`$
$`y_{32}`$ $`=`$ $`y_{19}+y_3+y_{15},y_{33}=y_{19}+y_4+y_{15},`$
$`y_{34}`$ $`=`$ $`y_{19}y_{13}+y_{15}+y_{12}+y_5y_1+y_3,`$
$`y_{35}`$ $`=`$ $`y_5y_1y_3+y_7y_{11}+y_{15}y_{13},`$
$`y_{36}`$ $`=`$ $`y_1+y_7y_4y_{13}+y_{15},`$
$`y_{37}`$ $`=`$ $`y_3+y_7y_{11}+y_{15}y_1+y_{12}y_{13},`$
$`y_{38}`$ $`=`$ $`y_1+y_7+y_{15},y_{39}=y_1y_{12}+y_{13},`$
$`y_{40}`$ $`=`$ $`y_5+y_4+y_1+y_{13}y_{12}.`$
By operating the group $`G`$ to the cubic equation, we get many such relations. Substituting the expressions of the $`y`$’s obtained above, we get cubic relations in terms of the chosen ten $`y`$’s; let us here rename the ten coordinates as:
$`g_1=y_1,g_2=y_3,g_3=y_4,g_4=y_5,g_5=y_7,`$
$`g_6=y_{11},g_7=y_{12},g_8=y_{13},g_9=y_{15},g_0=y_{19}.`$
Among these cubic equations we can find exactly thirty linearly independent ones. Therefore the variety $`V`$ is isomorphic to a subvariety of $`𝐏^9`$, coodinatized by $`g_1,\mathrm{},g_9,g_0`$, defined by thirty cubic equations, say $`cub_1,\mathrm{},cub_{30}.`$ Some of them which are relatively simple are shown below:
$`cub_1`$ $`:=`$ $`g_2g_8g_0+g_2g_8g_7g_3g_6g_0g_2g_3g_6g_3g_6^2g_8g_0g_3g_3g_6g_8`$
$`+g_3g_7g_0+g_2g_3g_7+g_6g_7g_3,`$
$`cub_2`$ $`:=`$ $`g_0^2g_1+g_3g_1g_0+g_1^2g_0+g_2g_0g_1g_5g_1g_0+g_0g_1g_6+g_8g_1g_0+g_2g_3g_6`$
$`+g_2g_6g_1+g_2g_8g_6+g_3g_1g_6+g_1^2g_6+g_8g_1g_6g_5g_3g_6g_5g_1g_6g_5g_8g_6,`$
$`cub_3`$ $`:=`$ $`g_8g_0g_3+g_8g_0g_4+g_2g_8g_0g_3g_6g_0g_2g_3g_6g_3g_6^2g_3g_6g_8g_3g_6g_4,`$
$`cub_4`$ $`:=`$ $`g_2g_0g_7+g_2^2g_7+g_2g_6g_7+g_2g_8g_7+g_2g_7g_4g_0g_4g_6g_8g_0g_4+g_0g_4g_7,`$
$`cub_5`$ $`:=`$ $`g_5g_3g_2g_2g_5g_1g_5g_8g_2+g_2g_3g_0+g_2g_0g_1+g_2g_8g_0+g_2g_3g_7`$
$`+g_2g_1g_7+g_2g_8g_7+g_2g_3g_4+g_2g_1g_4+g_2g_8g_4+g_2^2g_3+g_2^2g_1+g_2^2g_8+g_5g_4g_0`$
$`g_5g_3g_0g_5g_1g_0g_5g_8g_0+g_5g_7g_0,`$
$`cub_6`$ $`:=`$ $`g_5g_8g_1g_8g_1g_4g_2g_8g_1g_5g_8g_7+g_4g_7g_8+g_2g_8g_7`$
$`+g_5g_8^2g_8^2g_4g_2g_8^2+g_5g_8g_2g_2g_8g_4g_2^2g_8+g_5g_8g_6g_8g_4g_6g_2g_8g_6`$
$`g_5g_7g_6g_5g_4g_6+g_5g_1g_6,`$
$`cub_7`$ $`:=`$ $`g_9g_2g_3+g_3g_9g_5g_3g_9g_6g_9g_3g_0g_9g_3g_4+2g_9g_2g_4g_4g_9g_5`$
$`+g_4g_9g_6+g_9g_0g_4+g_9g_4^2+g_9g_2^2g_9g_5g_2+g_9g_2g_6+g_9g_0g_2`$
$`g_5g_6g_0g_5g_3g_6g_5g_6g_9,`$
$`cub_8`$ $`:=`$ $`g_5g_3g_7+g_5g_4g_7+g_2g_5g_7+g_5g_7g_6g_2g_3g_1g_2g_0g_1g_2g_1g_7+g_5g_3g_1`$
$`+g_5g_1g_0g_3g_1g_6g_0g_1g_6g_1g_7g_6g_3g_1g_0g_0^2g_1g_1g_7g_0`$
$`g_1g_4g_3g_1g_4g_0g_1g_4g_7,`$
$`cub_9`$ $`:=`$ $`g_8g_1g_3g_3g_1g_6+g_3g_8g_7+g_6g_7g_3g_9g_3g_7g_3g_8^22g_3g_6g_8`$
$`+g_9g_3g_8g_3g_2g_8g_2g_3g_6+g_9g_2g_3g_3g_6^2+g_3g_9g_6g_3g_1g_0g_2g_3g_1`$
$`g_0^2g_1g_2g_0g_1g_9g_0g_1g_1g_7g_0g_2g_1g_7g_9g_1g_7,`$
$`cub_{10}`$ $`:=`$ $`g_5g_8g_0g_8g_0g_4g_2g_8g_0+g_5g_8g_9g_9g_8g_4g_2g_9g_8+g_3g_6g_8`$
$`+g_8g_0g_3+g_2g_3g_6g_5g_3g_6+g_3g_6^2+g_3g_6g_0+g_3g_6g_4g_9g_2g_3+g_3g_9g_5`$
$`g_3g_9g_6g_9g_3g_0g_9g_3g_4,`$
$`cub_{11}`$ $`:=`$ $`g_2g_5g_1+g_5g_1^2g_5^2g_1+g_5g_1g_0+g_5g_3g_1g_2g_5g_7g_5g_1g_7+g_5^2g_7`$
$`g_5g_7g_0g_5g_3g_7+g_5g_8g_2+g_5g_8g_1g_5^2g_8+g_5g_8g_0+g_5g_3g_8`$
$`g_3g_2g_8g_2g_8g_0g_2g_8g_7,`$
$`\mathrm{}`$
$`cub_{19}`$ $`:=`$ $`g_5g_1g_0g_2g_5g_1+g_0^2g_1+2g_2g_0g_1+g_0g_1g_6+g_1g_7g_0+g_2g_1g_7+g_1g_7g_6`$
$`+g_1g_4g_0+g_2g_1g_4+g_1g_4g_6+g_2^2g_1+g_2g_6g_1g_5g_7g_6g_5g_4g_6,`$
$`\mathrm{}`$
We first study this system over the field $`K:=𝐂(g_1,\mathrm{},g5)`$. Geometrically speaking, we project the variety $`V`$ onto the 4-dimensional projective space coordinatized by $`g_1:\mathrm{}:g_5`$, and study the generic fibre of the projection
$$\pi :𝐏^9Vg_1:\mathrm{}:g_0g_1:\mathrm{}:g_5𝐏^4.$$
We shall prove that $`\pi `$ is generically two-to-one. This implies $`\pi `$ is two-to-one on
$$V^{}:=VV_{j=1}^{40}\{y_j=0\}.$$
Thus the argument in §2.4 shows that $`\phi :MV^{}`$ is an isomorphism.
We next study the intersection $`V\{y_1=0\}`$, and prove that $`V_{j=1}^{40}\{y_j=0\}`$ is the totality of the $`G`$-orbit of the closure of the image of $`X(2,6)`$ under the prolonged $`\phi `$, stated in §2.5.
### 3.2 Computation over $`K`$
From $`cub_3=0`$ and $`cub_7=0,`$ we can solve $`g_8`$ and $`g_9`$ as:
$`g_8`$ $`=`$ $`g_3g_6(g_2+g_6+g_0+g_4)/(g_0g_2g_0g_3+g_0g_4g_3g_6),`$
$`g_9`$ $`=`$ $`g_5g_6(g_0+g_3)/(g_2g_3+g_5g_3g_3g_6g_0g_3g_4g_3+2g_2g_4g_5g_4+g_4g_6+g_0g_4`$
$`+g_4^2+g_2^2g_5g_2+g_2g_6+g_0g_2g_5g_6).`$
Substituting these into $`cub_{19}`$, we can solve $`g_6`$:
$$g_6=\frac{g_1(g_5g_0g_5g_2+g_0^2+2g_0g_2+g_0g_4+g_7g_0+g_2g_7+g_2^2+g_2g_4)}{g_0g_1+g_1g_4+g_1g_7+g_2g_1g_5g_7g_5g_4}.$$
Substituting these expressions into $`cub_1,\mathrm{},cub_{30}`$, we have $`cub_3=cub_7=cub_{19}=0`$ of course and $`cub_{10}=0`$; though most of the remaining ones are complicated, $`cub_8`$ is relatively simple:
$$cub_8=(g_1g_5)qq/(g_0g_1+g_1g_4+g_1g_7+g_2g_1g_5g_7g_5g_4),$$
where
$`qq`$ $`=`$ $`g_1g_4t^2+(g_2g_5g_5g_4+g_1g_4+g_5g_3)s^2+2g_1g_4st`$
$`+`$ $`(g_5g_3g_1+g_2g_1g_5+g_1g_4g_3g_4^2g_5+g_1g_2g_4+g_5g_3g_4g_2g_5g_4+g_4^2g_1)s`$
$`+`$ $`(g_4g_5g_1+g_4^2g_1+g_1g_4g_3+g_1g_2g_4)t+g_3g_2g_1g_4g_5g_3g_1g_4+g_1g_4^2g_3`$
is a quadratic form in
$$t:=g_0,s:=g_7.$$
By this quadratic equation, we reduce the degree with respect to $`t`$ of the numerators of the $`cub`$’s to 1. In this way we get 26 $`(=304)`$ equations of the form
$$cube_j:a_jt+b_j=0,j=1,\mathrm{},30,j3,7,10,19,$$
where $`a_j,b_jK[s].`$ The polynomials
$$D_{jk}=a_jb_ka_kb_jK[s],j<k$$
have a unique common factor $`dd`$, which is quadratic in $`s`$; its actual form is given by
$`dd`$ $`:=`$ $`(g_1^2g_4^2+2g_2g_1^2g_4+g_3^2g_5^2+g_2^2g_3^2+g_2^2g_1^2+g_5^2g_4^2+2g_3g_2g_1g_42g_2g_1g_5g_4`$
$`2g_3g_5^2g_42g_1g_5g_4^2+2g_2^2g_3g_12g_3^2g_5g_2+2g_5g_3g_2g_4+2g_5g_3g_1g_42g_5g_3g_1g_2)s^2`$
$`+`$ $`(g_2^2g_1^32g_1^3g_2g_4+2g_1^2g_2g_4^2+2g_1^2g_5g_4^2+g_1^2g_4^3g_1^3g_4^22g_1g_5g_4^3`$
$`2g_3g_5^2g_4^22g_3g_2^2g_1^2+2g_3g_2g_1g_4^2+g_3^2g_2^2g_4+g_5^2g_4^3+2g_1^2g_2g_5g_4`$
$`+2g_2^2g_3g_1g_44g_3g_5g_1g_2g_4+g_5^2g_3^2g_4+g_2^2g_1^2g_4g_1g_5^2g_4^2g_2^2g_3^2g_1`$
$`g_3^2g_5^2g_1+2g_5g_3g_2g_4^2+2g_5g_1^2g_2g_3+2g_3^2g_5g_1g_22g_1g_2g_5g_4^2`$
$`+2g_3g_5^2g_1g_4+2g_3g_1g_5g_4^22g_3g_5g_1^2g_42g_3^2g_5g_2g_42g_1^2g_2g_3g_4)s`$
$``$ $`g_1^2g_3^2g_4^2+g_3g_1^2g_4^3g_1^3g_3g_4^2g_1^2g_2g_3^2g_4g_1^3g_3g_2g_4g_3^2g_1g_2^2g_4`$
$`g_3^2g_1g_2g_4^2+g_3g_1g_2g_4^3+g_2^2g_3g_1g_4^2g_3g_5^2g_1^2g_4+g_1^3g_3g_5g_4`$
$`+g_3^2g_5g_1g_4^2g_3^2g_5^2g_1g_4+g_3g_5^2g_1g_4^2g_3g_5g_1g_4^3g_3g_1^2g_2^2g_4`$
$`+g_3^2g_5g_1^2g_42g_3g_1g_2g_5g_4^2+2g_3g_1^2g_2g_5g_4+2g_3^2g_1g_2g_5g_4.`$
So $`dd(s)=0`$ is the compatibility condition of the system of equations $`cube_j:a_j(s)t+b_j(s)=0`$, among which $`cube_{11}:a_{11}(s)t+b_{11}(s)=0`$ has the minimal degree of coefficients; actually, degrees of $`a_{11}`$ and $`b_{11}`$ with respect to $`s`$ are 3 and 4.
Now we get three equations $`dd=0,qq=0`$ and $`cube_{11}=0`$. Substituting the expression $`t=b_{11}/a_{11}`$ into $`qq`$, we get a polynomial in $`s`$. We can prove that this polynomial has $`dd`$ as a factor.
We have proved that for a given generic point $`g_0:\mathrm{}:g_5𝐏^4`$, the inverse under $`\pi `$ consists of two points. They are obtained as follows: solve the quadratic equation $`dd=0`$ with respect to $`s(=g_7)`$. For each root, $`t(=g_0)`$ is uniquely determined by the linear equation $`cube_{11}`$. And $`g_6,g_8,g_9`$ are uniquely determined by $`cub3,cub7,cub19`$ as stated above. Then all the relations are satisfied.
### 3.3 Intersection of $`V`$ and $`\{y_1=0\}`$
We should better work on $`V`$ living in $`𝐏^{401}`$. In terms of the forty coordinates $`y_1,\mathrm{},y_{40}`$, the cubic equations are 2-term equations. Thus the vanishing of $`y_1`$ implies that of some other coordinates. Thanks to $`G`$-action, we can assume that $`y_{10}=0`$. The vanishing of these two coordinates forces the vanishing of other coordinates. Tedious case-by-case study shows that every component of $`V\{y_1=0\}`$ is included in the $`G`$-orbit of
$$V\{y_1=\mathrm{}=y_9=y_0=0\},$$
which is the closure of the image of $`\{xQ(x)=0\}`$ under the prolonged $`\phi `$. |
warning/0002/astro-ph0002168.html | ar5iv | text | # The afterglows of gamma-ray bursts
## I Introduction
GRBs have mystified and fascinated astronomers since their discovery. Their brilliance and their short time variability clearly suggest a compact object (black hole or neutron star) origin. Three decades of high-energy observations, culminating in the definitive measurements of CGRO/BATSE, determined the spatial distribution to be isotropic yet inhomogeneous, suggestive of an extragalactic population (see fm95 for a review of the situation prior to the launch of the BeppoSAX mission). Further progress had to await the availability of GRB positions adequate for identification of counterparts at other wavelengths.
In the cosmological scenario, GRBs would have energy releases comparable to that of supernovae (SNe). Based on this analogy, Paczyńsk & Rhoads pr93 and Katz k94 predicted that the gamma-ray burst would be followed by long-lived but fading emission. These papers motivated systematic searches for radio afterglow, including our effort at the VLA fk95 . The broad-band nature of this “afterglow” and its detectability was underscored in later work mr97 ; v97 .
Ultimately, the detection of the predicted afterglow had to await localizations provided by the Italian-Dutch satellite, BeppoSAX b+97 . The BeppoSAX Wide Field Camera (WFC) observes about 3% of the sky, triggering on the low-energy (2 – 30 keV) portion of the GRB spectrum, localizing events to $``$ 5 – 10 arcminutes. X-ray afterglow was first discovered by BeppoSAX in GRB 970228, after the satellite was re-oriented (within about 8 hours) to study the error circle of a WFC detection with the 2 – 10 keV X-ray concentrators. The detection of fading X-ray emission, combined with the high sensitivity and the ability of the concentrators to refine the position to the arcminute level, led to the subsequent discovery of long lived emission at lower frequencies c+97 ; JvP+97 ; f+97 .
Optical spectroscopy of the afterglow of GRB 970508 led to the definitive demonstration of the extragalactic nature of this GRB mdk+97 . The precise positions provided by radio and/or optical afterglow observations have allowed for the identification of host galaxies, found in almost every case. Not only has this provided further redshift determinations, but it has been useful in tying GRBs to star formation through measurements of the host star formation rate (e.g. kdr+98 ; dkb+98 ). HST with its exquisite resolution has been critical in localizing GRBs within their host galaxies and thereby shed light on their progenitors (e.g. fpt+99 ; hh99 ; bod+99 ). Observations of the radio afterglow have directly established the relativistic nature of the GRB explosions f+97 and provided evidence linking GRBs to dusty star-forming regions. Radio observations are excellent probes of the circumburst medium and the current evidence suggests that the progenitors are massive stars with copious stellar winds. The latest twist is an apparent connection of GRBs with SNe bkd+99 . Separately, an important development is the possible association of a GRB with a nearby (40 Mpc) peculiar SN gvv+98 ; kfw+98 .
In this paper we review the primary advances resulting from afterglow studies. §II discusses the statistics of detections to-date, including possible causes for the lack of radio and optical afterglows from some GRBs. In §III we review constraints on the nature of the progenitor population(s), in particular evidence linking some classes of GRBs to SNe. §IV describes the status of current understanding of the physics of the afterglow emission. Here we compare observations to predictions of the basic spherically-symmetric model, and describe complications arising from deviations from spherical symmetry and non-uniform distribution of the circumburst medium. We conclude with speculations of the near and long-term advances in this field (§V).
We point out that this review has two biases. First, given the concentration of previous review articles on optical and X-ray observations, we emphasize the unique contributions of radio afterglow measurements. Second, this article is intended to also provide a summary of the efforts of the Caltech-NRAO-CARA GRB collaboration, and therefore details our work in particular. This review is in response to review talks given at the 1999 Maryland October meeting (SRK) and the 5th Huntsville GRB meeting (DAF and SRK).
## II Statistics of Afterglow Detections
Afterglow emission was first detected from GRB 970228, both at X-ray c+97 and optical frequencies JvP+97 , but not at radio wavelengths fksw98 . The first radio afterglow detection came following the localization of GRB 970508 f+97 . Figure 1 shows two examples of radio lightcurves. The radio afterglow of GRB 970508 is famous for several reasons: it was the first radio detection, it gave the first direct demonstration of relativistic expansion, and it remains the longest-lived afterglow fwk00 .
Afterglow emission is now routinely detected across the electromagnetic spectrum. BeppoSAX has been joined in studying the X-ray afterglows by the All Sky Monitor (ASM) aboard the X-ray Timing Explorer (XTE), the Japanese ASCA mission, and recently the Chandra X-ray observatory (CXO). A veritable armada of optical facilities (ranging from 1-m class telescopes to the 10-m Keck telescopes) routinely discover and study optical afterglows. The HST has been primarily used to make exquisite images of the host galaxies (see above) but in the near future we expect other uses such as UV spectroscopy and identification of underlying SNe. The VLA has led the detection in radio. However, other centimeter-wavelength facilities (the Australia Telescope National Facility, Westerbork Synthesis Radio Telescope, the Ryle Telescope) and millimeter wavelengths (James Clerk Maxwell Telescope, the Owens Valley Millimeter Array, IRAM and the Plateau de Bure Interferometer) are now regularly contributing to afterglow studies.
Figure 2 summarizes the statistics of afterglow detections. In almost all cases, X-ray emission has been detected, establishing the critical importance of prompt X-ray observations. Optical afterglow appears to be detected in about 2/3 of all well-localized events if sufficiently deep optical images are taken rapidly (i.e. within a day or so of the burst). Radio afterglows are detected in 40% of the cases – far more often than usually assumed. We refer the reader to the Frail et al. fkw+00 for a comprehensive summary of the X-ray/optical/radio afterglow detection statistics. The non-detections are, as discussed below, as interesting as the detections.
Radio Non-detection. The failure to find radio afterglow is most likely due to lack of sensitivity. The brightest radio afterglow to date is that from GRB 991208 (Frail GCN<sup>1</sup><sup>1</sup>1GCN refers to the GRB Coordinates Network Circular Services. This network is maintained by S. Barthelmy at the Goddard Space Flight Center;
see $`\mathrm{http}://\mathrm{lheawww}.\mathrm{gsfc}.\mathrm{nasa}.\mathrm{gov}/\mathrm{docs}/\mathrm{gamcosray}/\mathrm{legr}/\mathrm{bacodine}/\mathrm{gcn}\mathrm{\_}\mathrm{main}.\mathrm{html}`$ 451) with a peak flux of 2 mJy, a 60-$`\sigma `$ detection (at centimeter wavelengths) whereas the weakest afterglow is typically around 5$`\sigma `$. In contrast, at optical and X-ray wavelengths, afterglow emission is routinely detected at hundreds of sigma. If the VLA were to be upgraded by a factor of 10 in sensitivity, then we predict that radio afterglow emission would, like X-ray emission, be detected from most GRBs.
Optical Non-detection. Non-detection at optical wavelengths is more interesting, as it may result in some cases from extinction along the line of sight or within the source. Bad weather as well as rapid fading of the afterglow has certainly hindered some optical searches, which, due to notification delays, typically begin some hours after the event. Furthermore, low Galactic latitude events may be obscured, or hidden in crowded foregrounds. However, in some cases deep searches have been performed with no success. Here, non-detection likely results from extinction by dust in the burst host galaxy and/or absorption by the intergalactic medium. GRB 970828 ggv+98 is one example, as is the more dramatic case of GRB 980329. This burst was one of the brightest events in the WFCiaa+98 . Searches for optical afterglow emission failed to identify any counterpart. VLA observations identified an unusual radio variable in the field tfk+98 . Soon thereafter, a red afterglow and a bright IR afterglow were identified (Klose GCN 43, Larkin et al. GCN 44). Taylor et al. tfk+98 suggest that the GRB arose in a region with high extinction. Further optical and IR work on this interesting afterglow can be found in gcp+99 , ppm+98 , and rlm+99 .
Optically dim “red” but bright IR afterglows can also result from the GRB being located at high redshift. Intergalactic HI absorption will result in a wavelength cutoff below the Lyman limit, $`<912(1+z)`$ Å, where $`z`$ is the redshift of the source. This effect was originally invoked to explain the faint R-band but bright IR emission from GRB 980329 f99 . We now know, based on recent Keck observations, that the GRB host is blue, incompatible with a high-$`z`$ origin. Rather, it is more tenable that the host is a typical star-forming galaxy with dusty star-forming regions, and that the GRB occurred in one such regiontfk+98 . We are presently carrying out IR spectroscopy of this host to determine the redshift and the star formation rate (SFR). While searching for “R dropouts” may in the future provide an effective method for finding high-redshift events, it is clear that cross-calibrated multi-band photometry of higher quality than currently exists will be required to make this useful.
X-ray Non-detection. The spectra of most GRB events clearly extend into the X-ray band, as established by GINGA observationssfmy98 . How the X-ray emission observed during the burst connects to the X-ray afterglow is uncertain, due to sensitivity limitations of wide-field monitors. X-ray afterglow emission appears to be ubiquitous. Observations of the X-ray afterglow are important for two reasons: (i) the observations of the X-ray afterglow by sensitive imaging instruments (e.g. the concentrators aboard BeppoSAX) result in sufficiently precise (arcminute) localization and (ii) a significant (perhaps even a dominant) fraction of the explosion energy appears to be radiated in this band. Of all the SAX bursts, GRB 970111 is peculiar for the absence of X-ray afterglow (admittedly the data were obtained about 17 hours after the burst) fac+99 . In view of the critical role played by X-ray afterglow in localization of GRBs we regard this non-detection to be worthy of further investigation.
## III The Nature of the Progenitors
In almost all cases, a host galaxy has been identified at the location of the fading afterglow. GRB redshifts can be obtained either via absorption spectroscopy (when the transient is bright) or by emission spectroscopy of the host galaxy. In Figure 3 and Table 1 we summarize the measured redshifts and host galaxy magnitudes. While the distance scale debate is settled (at least for the class of long duration GRBs, see below) we remain relatively ignorant of the nature of the central engine. Currently popular GRB models fall into two categories: (i) the coalescence of compact objects (neutron stars, black holes and white dwarfs eic+89 ; ls74 ; moc93 ; npp92 ) and (ii) the collapse of the central iron core of a massive star to a spinning black hole, a “collapsar” woo93 ; mw99 . We now summarize the light shed on the progenitor problem by afterglow studies.
The Location of GRBs Within Hosts. A fundamental insight into the nature of SNe came from their location with respect to other objects within the host galaxy (specifically HII regions and spiral arms) and the morphology of the host galaxy itself (elliptical versus spiral). In a similar manner, we are now making progress in understanding GRB progenitors by measuring offsets with respect to other objects in the host galaxies. The rather good coincidence of GRBs with host galaxies already suggests that they are unlikely to be a halo population (as would be expected in the coalescence scenario bsp99 ). On the other hand, with the possible exception of GRB 970508pfb+99 , they are clearly not associated with galactic nuclei (i.e. massive central black holes). Typical offsets of GRBs from the centroid of their host galaxies are comparable to the half-light radii of field galaxies at comparable magnitudes, suggesting that GRBs originate from stellar populations.
Host Galaxies. Demonstrating a direct link between GRBs and (massive) star formation is more difficult. On the whole, the population of identified hosts seems typical in comparison to field galaxies in the same redshift and magnitude range. The hosts have average luminosities for field galaxies, modulo corrections due to evolution. Their emission line fluxes and equivalent widths are also statistically indistinguishable from the normal field galaxy population. The observed star formation rates, derived from recombination line fluxes (mostly the \[O II\] 3727$`\AA `$ line) and from the UV continuum flux range from less than $`1M_{}`$ yr<sup>-1</sup> to several tens of $`M_{}`$ yr<sup>-1</sup> – typical of normal galaxies at comparable redshifts (extinction corrections can increase these numbers by a factor of a few, but similar corrections apply to the comparison field galaxy population as well). It will probably be necessary to have a sample of several tens of GRB hosts before a correlation of GRBs with the (massive) star formation rate can be tested statistically. However, below we point to several specific examples which are suggestive of a link between GRBs and star-forming regions.
Association with Starforming Regions. There is evidence showing that GRBs arise from dusty regions within their host galaxies. In this respect, radio observations provide a unique tool for detecting events in regions of high ambient density (as was the case for GRB 980329). An even more extreme example is GRB 970828, where the host was identified based solely on the VLA discovery of a radio flare dfk+2000 . Interestingly enough, this is the dustiest galaxy in the sample of GRB hosts to-date.
Second, some GRBs appear to be located within identifiable star-forming regions. An example is GRB 990123 bod+99 ; ftm+99 ; hh99 . VLA observations of GRB 980703 fbk+2000 are perhaps more convincing. The radio observations can be sensibly interpreted by appealing to free-free absorption from a foreground HII region (which would dwarf the Orion complex). If this interpretation is correct then this would be strong evidence for a GRB being located within a starburst region.
The GRB–SN link. If GRBs arise from the collapse of a massive star, it is an unavoidable consequence that emission from the underlying supernova should be superimposed on the afterglow. Bloom et al. bkd+99 may have made the first detection of a possible SN component in the GRB 980326 lightcurve (Fig. 4). These authors noted that SNe, in contrast to afterglows, have distinctive temporal and spectral signatures: rising to a maximum at $`20(1+z)`$ days, with little emission blueward of about 4000 Å in the restframe (and certainly blueward of 3000 Å) owing to a multitude of resonance absorption lines. This discovery has led to other possible SN detections, most notably GRB 970228 gal+99 ; rei99 .
The suggestion of a GRB–SN connection is an intriguing one but it has yet to be placed on a firm footing. Important questions are: (i) are all long-duration GRBs accompanied by SNe? (ii) if so, are these SNe of type Ib/c? Ground-based observations are possible in those cases where the afterglow decays rapidly (e.g. GRB 980326) or if high quality optical and IR observations exist (e.g. GRB 970228).
We need more examples to test the GRB–SN link. Future progress will depend on a combination of ground and HST observations. For relatively nearby GRBs especially those with a rapidly decaying optical afterglow it would be attractive and feasible to obtain the spectrum of the SN around the time when the flux from the SN peaks. A moderate quality spectrum with SN-like features would have the singular advantage of definitively confirming the SN interpretation (as opposed to alternatives involving re-radiation by dust wd99 ). However, for most GRBs, we expect HST observations to play a critical role. HST’s widely recognized strengths in accurate photometry of sources embedded in galaxiesgkc+98 and photometric stability make the detection of a faint SN against the optical afterglow and the host galaxy possible.
Diversity of the Progenitor Population. As was the case with SNe, it is likely naive to think of a single progenitor population. Below, we discuss the two additional classes which show some promise: the mysterious short duration GRBs and a possible class of low luminosity GRBs associated with SNe.
Short Events. It has been known for some time that the distribution of the duration of GRBs appears to be bimodal fm95 ; see Figure 5. Furthermore, these two groups may have different spatial distributions kc96b , with the short bursts being detected out to smaller limiting redshifts. However, we know very little about this class of GRBs since, as noted earlier, all bursts localized by BeppoSAX and RXTE thus far are of long duration (Figure 5). Fortunately, improvements in BeppoSAX and the imminent launch of HETE-2 provide for the first time the opportunity to follow-up short GRBs.
The short duration bursts are difficult to accommodate in the collapsar model, given the long collapse time of the core. However, they find a natural explanation in the coalescence models. How would these bursts manifest themselves? Li & Paczyǹski lp98b speculate that if the short-duration bursts result from NS–NS mergers then they may leave a bright, but short-lived ($`<1`$ day) optical transient. Radio observations provide a complementary tool for determining the nature of the short duration bursts. The low ambient density would result in weak afterglows (since flux $`\rho ^{1/2}`$) which are potentially detectable. Radio observations have additional advantages of a longer lived afterglow, immunity from weather and freedom from the diurnal cycle.
Gamma-ray Bursts Associated with Supernovae. Observers and theorists alike have been intrigued by the possibility that the bright supernova, SN 1998bw, discovered by Galama et al. gvv+98 in the error circle of GRB 980425 paa+99 , is associated with the gamma-ray event (Figure 6). Kulkarni et al. kfw+98 discovered that the SN had an extremely bright radio counterpart; see Figure 7. We noted that the inferred brightness temperature exceeded the inverse Compton catastrophe limit of $`5\times 10^{11}`$ K and to avoid rapid cooling we postulated the existence of a relativistically expanding blastwave ($`\mathrm{\Gamma }>2`$). This relativistic shock is, of course, in addition to the usual sub-relativistic SN shock. This relativistic shock may have produced the GRB at early times. (We note here that we disagree with the much lower energy estimates of wl99 ; our recent calculations using the same assumptions as those made in wl99 result in an energy estimate similar to that obtained earlier kfw+98 from minimum-energy formulation). The optical modeling of the lightcurve and the spectra show that GRB 980425 was especially energetic imn+98 ; wes99 with an energy release of $`3\times 10^{52}`$ erg and Nickel production of nearly nearly a solar mass.
If GRB 980425 is associated with 1998bw, then this type of event is rare among the SAX localizations. GRB 980425 is most certainly not a typical GRB: the red-shift of SN 1998bw is 0.0085 and the $`\gamma `$-ray energy release in GRB 980425 is at least four orders of magnitude less than in other cosmologically located GRBs. For this reason, most astronomers (especially those in the GRB field; see Wheeler’s foray in experimental sociologyw99 ) do not believe the association between GRB 980425 and SN 1998bw. On the other hand, as evidenced by the intense interest in and modeling of the radio and optical data of SN 1998bw, this object is of considerable interest to the SN community. Indeed, we believe that the proposed GRB–SN association controversy has muddied the main issue: SN 1998bw is an interesting SN in its own right.
What is the true distinguishing feature of SN 1998bw that may connect it to a GRB event? Is it the large energy release, as suggested by several authorsimn+98 ; grs+99 ? We argue that in fact it is the energy coupled into relativistic ejecta that most closely connects SN 1998bw to a GRB. In a typical SN, about $`10^{51}`$ erg is coupled to the envelope of the star (a small fraction of the total SN energy release of $`10^{53}`$ erg). In a GRB, a similar amount of energy ($`10^{51}`$$`10^{52}`$ erg depending on the event) is coupled to a much smaller ejecta mass, resulting in relativistic outflow. For SN 1998bw, applying the minimum energy formulation to the radio observations we infer the relativistic shell to contain $`10^{50}`$ erg. Not only is this uncharacteristic of a typical SN (there exists no evidence for relativistic ejecta in ordinary SN), but it is not dissimilar from the energy implied for GRB outflows. One could therefore envisage a continuum of physical phenomenon between SN 1998bw and cosmological GRBs provided we use the energy in the relativistic ejecta as the basic underlying parameter and not the isotropic gamma-ray release.
## IV Afterglow: The Physics and Energetics of the Fireball
One can consider a GRB to be like a SN explosion with a central source releasing energy $`E_0`$ (comparable to the mechanical release of energy in an SN). This is the so-called fireball model. The difference between an SN and a GRB is primarily in ejecta mass: 1–10 $`M_{}`$ for SNe whereas only $`10^5M_{}`$ for GRBs. The evolution of a GRB is much faster than that of a SN due to two factors: the ejecta expand relativistically and, thanks to the smaller ejecta mass, the optical depth is considerably smaller.
As the ejecta encounter ambient gas, two shocks are produced: a short-lived reverse shock (traveling through the ejecta) and a long-lived forward shock (propagating into the swept-up ambient gas). Afterglow emission is identified with emission from the forward shock. In order to obtain significant afterglow emission, several conditions are necessary. (1) Rapid equipartition of electrons with the shocked protons (which hold most of the energy). (2) Acceleration of electrons to a power law spectrum (particle Lorentz factor distribution, $`dN/d\gamma \gamma ^p`$). (3) Rapid growth of the magnetic field with energy density in the range of $`10^2`$ of that of the protons. Under these circumstances, afterglow emission is dominated by synchrotron emission of the accelerated particles; see spn97 ; w97a . The weakness of this model is the assumption of growth in the magnetic field strength to the high values noted above (R. Blandford, pers. comm.).
The theoretically expected afterglow spectrum is shown in Figure 8. Three key frequencies can be identified: $`\nu _a`$, the synchrotron self-absorption frequency; $`\nu _m`$, the frequency of the electron with a minimum Lorentz factor (corresponding to the thermal energy behind the shock) and $`\nu _c`$, the cooling frequency. Electrons which radiate above $`\nu _c`$ cool on timescales equal to the age of the shock. The evolution of these three frequencies is determined by the hydrodynamical evolution of the shock which in turn is affected by two principal factors: the environment of the GRB and the geometry of the explosion.
The GRB environment. The earliest afterglow models made the simplifying assumption of expansion into a constant density medium. This is an appropriate assumption should the GRB progenitor explode into a typical location of the host galaxy. However, there is increasing evidence tying GRBs to massive stars (see §III). It is well known that massive stars lose matter throughout their lifetime and thus one expects the circumburst medium to exhibit a density profile, $`\rho r^2`$ where $`r`$ is the distance from the progenitor. Chevalier & Li cl99 refer to these two models as the ISM (interstellar medium) and the wind model respectively. As can be seen from Figure 8 these two models give rise to rather different evolution of the three critical frequencies.
Geometry: Jets versus Spheres. The hydrodynamics is also affected by the geometry of the explosion. Many powerful astrophysical sources have jet-like structure. There is evidence (from polarization observations) indicating asymmetric expansion in SNe w99 , so it is only reasonable to assume that GRB afterglows also have jet-like geometry as well. A clear determination of the geometry is essential in order to infer the true energy of the explosion. This is especially important for energetic bursts such as GRB 990123 whose isotropic energy release approaches $`M_{}c^2`$.
Let the opening angle of the jet be $`\theta _0`$. As long as the bulk Lorentz factor, $`\mathrm{\Gamma }`$, is larger than $`\theta _0^1`$, the evolution of the jet is exactly the same as that of a sphere (for an observer situated on the jet axis). However, once $`\mathrm{\Gamma }`$ falls below $`\theta _0^1`$ then two effects become important. First, for a well defined jet, the on-axis observer sees an edge and thus one expects to see a break in the afterglow emission. Second, the lateral expansion of the jet (due to heated and shocked particles) will start affecting the hydrodynamical explosion.
Wind or ISM? The two key diagnostics to distinguish these two models are the evolution of the cooling frequency (see Figure 8) and the early behavior of the radio emission. In the wind model, the radio emission rises rapidly (relative to the ISM model) and the synchrotron self-absorption frequency falls rapidly with time. Both these result from the fact that the ambient density decreases with radius (and hence in time) in the wind model.
Unfortunately, in general, the current data are not of sufficient quality to firmly distinguish the two models. For example in GRB 980519, the same optical and X-ray data appear to be adequately explained by the jet+ISM model sph99 and the sphere+wind model cl99 . Including the radio data tips the balance, but only slightly, in favor of the wind model fks+2000 . In our opinion, the best example for the wind model is that of GRB 980329 fks+99 ; see Figure 9. This afterglow exhibits the two unique signatures of the wind model: high $`\nu _a`$ and a rapid rise. Given the importance of making the distinction between the wind and the ISM model we urge early wide band radio observations (especially at high frequencies).
Energetics. Of all the physical parameters of the fireball, the most eagerly sought parameter is the total energy $`E_0`$. By analogy with supernovae, it is $`E_0`$ which sets the GRB phenomenon apart from other astrophysical phenomena. Classes of GRBs may eventually be distinguished and ranked by their energy budget; for example, long-duration events, short duration events and supernova-GRBs (see §III).
One approach has been to use the isotropic $`\gamma `$-ray energy as a measure of $`E_0`$; see Figure 3. There are three well known problems with such estimates. First, collimation of the ejecta (jets) will result in overestimation of the total energy release. For GRB 990510 where a good case for a jet has been established (Figure 9), the standard isotropic energy estimate is probably a factor of 300 more than the true energy hbf+99 . Second, even after accounting for a possible jet geometry, the efficiency of converting the shock energy into gamma-ray emission is very uncertain. For example, some authors k99 advocate low efficiency ($`1\%`$) which would result in an enormous upward correction to the usual isotropic estimates. Third, the bulk Lorentz factor is extremely high during the emission of $`\gamma `$-rays and thus the estimates critically depend on assumption of the geometry and granularity kp99 of the emitting region. In particular, if the emission is from small blobs kp99 then the inferred estimates are grossly in error.
In contrast to this highly uncertain situation, afterglows offer (in principle) more robust methods to evaluate $`E_0`$. In view of the importance of determining $`E_0`$ we summarize the different methods of determining $`E_0`$ from afterglow observations. One approach is to fit a “snapshot” broad-band afterglow spectrum (from radio to X-rays) to an afterglow model; this approach was pioneered by Wijers & Galama wg99 . The strength of this method is that the estimated $`E_0`$ is, in principle, robust. Specifically, the estimate does not depend on the usually unknown environmental factors (run of density). However, in practice, this method is very sensitive to the values of the critical frequencies (Figure 8) which are usually not well determined. This difficulty explains the wildly differing estimates of $`E_0`$ for GRB 970508 wg99 ; gps99 . Furthermore, this method uses measurements obtained at early times (when the afterglow at high frequencies is bright) with the result that the true source geometry is hidden by relativistic beaming.
A second approach is to model the light curves of the afterglow in a given band, specifically a radio band. The advantages of this method are the photometric stability of radio interferometers and the low Lorentz factor at the epoch of the peak of the radio emission. The disadvantages are two-fold: the sensitivity to the environmental parameters (density) and the assumption of the constancy of the microphysics parameters (electron and magnetic field equipartition factors). Application of this approach to GRB 980703 has resulted in seemingly accurate measures of the fireball parameters fbk+2000 .
Freedman & Waxman fw99 take yet another approach, and estimate the energy release from late time X-ray observations. They show that the X-ray flux is insensitive to the GRB environment, and obtain robust estimates of the fireball energy per unit solid angle: from $`3\times 10^{51}`$ erg to $`3\times 10^{53}`$ erg.
With all the above approaches, however, the possible collimation of the ejecta in jets is still a major uncertainty. This can be addressed by observing the evolution of the afterglow as the “edge” of the jet becomes visible. In most cases no evidence for jets has been seen, with the notable exceptions of GRB 990510 and possibly GRB 990123. In addition, a variety of statistical arguments (the absence of copious numbers of “orphan afterglows”)gri99 ; gvb+99 ; rho97 suggests that, on average, the collimation cannot be extreme, and that for most bursts the opening angle is not less than 0.1 radian. Thus the total energy for most bursts may be reduced to the range of $`10^{50}`$ erg to $`3\times 10^{51}`$ erg, but could easily be much higher in at least some cases.
Possibly the best approach to determining the energetics, which minimizes uncertainties due both to collimation (jets) and to the environment is to model the afterglow after it becomes non-relativistic. This method builds on the well established minimum energy formulation and the self-similarity of the Sedov solution. Not only are the ejecta truly non-relativistic, but they are also essentially spherical, as by this time jets will have had sufficient time to have undergone significant lateral expansion. Indeed, we can justifiably call this “fireball calorimetry” fwk00 . Applying this technique to the long-lived afterglow of GRB 970508 (Figure 1) led to the surprising result that $`E_05\times 10^{50}`$ erg – weaker than a standard SN! This is an astonishing result. If true, this result would suggest that it is not $`E_0`$ which is the prime distinction between GRBs and SNe but the ejecta mass. However, Chevalier & Li cl2000 interpret the same data in the wind framework and derive much larger $`E_0`$. Clearly, we need more well studied afterglows with sufficient observations to first distinguish the circumburst environment (wind versus ISM) and then radio observations over a sufficiently long baseline to undertake calorimetry. Nonetheless, one should bear in mind that the current evidence for large energy release in GRBs is not as strong as is usually assumed.
## V Epilogue and Future
Clearly, the GRB field is evolving rapidly. Along what direction\[s\] will this field proceed in the coming years? One way to anticipate the future is by considering analogies from the past.
In §III we already discussed the parallels between the SN field and the GRB field. Here we discuss the numerous parallels with quasar astronomy. First discovered at radio wavelengths, we now study quasars across the electromagnetic spectrum. Although still identified by their gamma-ray properties, we now recognize the tremendous value of pan-chromatic GRB and afterglow studies. In both cases, there was considerable controversy about the distance scale. However, once this issue was settled, it became clear that quasars are the most energetic objects (sustained power) whereas GRBs are the most brilliant. For both, the ultimate energy appears to be related to black holes (albeit of different masses).
The raging issues in GRB astronomy today are the same that fueled quasar studies in the 60’s: the spatial distribution, the extraction of energy from the central engine, the transfer of energy from stellar scales to parsec scales, and the geometry of the relativistic outflow (sphere or jet). Astronomers took decades to unify the seemingly diverse types of quasars, and to conclude that there are two types of central engines: radio loud and radio quiet. Likewise, there may well be two types of GRB engines: rapidly and slowly spinning black holes emerging respectively from collapse of a rotating core of a massive star or coalescence of compact objects and the collapse of a massive star. This picture could potentially explain both the cosmologically located GRBs and SN 1998bw. Finally, we can project that in the future, GRBs may be used to probe distant galaxies, just as quasars are used today to study the IGM.
There is a feeling in the astronomical community (outside the GRB community) that the GRB problem is “solved”. The truth is that the GRB problem is now getting defined! We now summarize our view of the major issues and anticipated near term advances. In our opinion the major issues are Diversity, Progenitors and Energy Generation.
As discussed earlier, high energy observations suggest the existence of two classes: short and long duration bursts. It is possible that afterglow observations may demarcate additional classes. If so, one can contemplate that within a year (assuming abundant localizations by HETE-2) that we will have new GRB designations such as sGRBs (GRBs with late time bump indicative of an underlying SN), wGRBs (GRBs whose afterglow clearly indicates a wind circumburst medium shaped by stellar winds), iGRBs (GRBs which explode in the interstellar medium) and so on.
The broad indications are that GRBs are associated with stars and most likely massive stars. However, we know little beyond this. Comparing the unbeamed GRB event rate of $`1.8\times 10^{10}`$ yr<sup>-1</sup> Mpc<sup>-3</sup> schmidt99 with $`3\times 10^5`$ Type Ibc SN yr<sup>-1</sup> Mpc<sup>-3</sup> and $`10^6`$ yr<sup>-1</sup> NS–NS merger Mpc<sup>-3</sup> lamb99 shows that GRBs events are extremely rare; here we note that the present data do not support a collimation correction in excess of 100. It will be quite some time before we will be in a position to identify the conditions necessary for a star to die as a GRB.
It is our opinion that SN 1998bw is a major development in the field of stellar collapse. The association (or lack) with GRB 980425 unfortunately has distracted our attention of this important development. The existence of a significant amount of mildly relativistic material, $`10^{50}`$ erg kfw+98 , is fascinating and it is ironic that none of the models can account for this inferred value whereas most of the theoretical effort has gone into explaining the gamma-ray burst itself (especially considering the uncertain association of GRB 980425 with SN 1998bw). Clearly, SN 1998bw is a rare event but we are convinced that more such events will be found and accordingly have mounted a major campaign to identify these SNe. The robust signatures of this class are high $`T_B`$ and prompt X-ray emission since these are necessary consequences of a relativistic ejecta. We note that if these future events are as bright as SN 1998bw then the energy in the relativistic ejecta can be directly measured by VLBI observations of the expanding radio shell.
It is vitally important to make quantitative progress in determining the energy release in GRBs. As discussed in §IV, firm estimates of the energy release require well sampled broad-band data at early times and densely sampled radio light curves out to late times. This will require a coordinated approach and necessarily involve many observatories around the world and in space. The same datasets will also help us understand a profound puzzle: if GRBs indeed arise from the death of massive stars then why do we not see signatures for a circumburst medium shaped by stellar winds in all long duration GRBs? Even ardent supporters of the wind model cl99 ; cl2000 concede that some GRBs (e.g. GRB 990123, 990510) are due to a jet expanding into a constant density medium.
We now discuss the anticipated returns. True to our tradition as observers, we order the discussion by wavelength regimes!
Radio Observations: Dusty galaxies, Circumstellar Edges and Reverse Shocks. Perhaps the most exciting use of radio afterglow is in identifying dusty star-forming host galaxies. Such host galaxies are not readily seen at optical wavelengths. Currently, such galaxies are eagerly sought and studied at sub-millimeter wavelengths. However, the sensitivity and localization of such galaxies by sub-millimeter telescopes is poor. In contrast, GRB host galaxies are identified at the sub-arcsecond level. The present radio afterglow detection rate of 40% already places an upper limit on the amount of star-formation in dusty regions, viz. this rate is not larger than that measured from optical observations. This result is entirely independent of the conclusion based on studies in the sub-millimeter regime, or the diffuse cosmic FIR background found in the COBE data. However, the result does rely on two assumptions: (i) GRBs trace star formation and (ii) the GRB explosion and its aftermath does not radically alter the ambient medium (i.e., with a prompt and complete destruction of dust grains along the line of sight).
Radio observations of SNe offer a probe of the distribution of the circumstellar matter. A spectacular example is SN 1980K whose radio flux dropped 14 yrs after the explosion mvw+98 . A progenitor star which suffered mass loss with variation in the wind speed could explain the observations. Indeed, one expects significant radial structure in the circumburst medium as the progenitor evolves from a blue star to a red supergiant and thence to possibly a blue supergiant etc. If GRBs come from binary stars which undergo a phase of common envelope envolution bllb99 then the structure would be even more complicated. Thus radio observations have the potential (in fortunate circumstances) to give us insight into the mass loss history of the progenitor star\[s\].
The prompt optical emission from GRB 990123 abb+99 has been interpreted to arise from the reverse shock sp99 . Far less discussed is the prompt radio emission – a radio flare – also seen from this burst kfs+99 . Sari & Piran sp99 suggest that the radio emission also originates from the reverse shock as the electrons cool. Observations related to the reverse shocks are important since it is only through these observations that we have a chance of studying the elusive ejecta. We now have four such examples of radio flares kf00 and this represents an order of magnitude better success rate than ROTSE+LOTIS. We urge theorists to pay attention to these new findings. More to the point, radio observations appear to be fruitful for the study of reverse shocks, especially when combined with observations of the prompt optical emission. This bodes well for the coming years given the efforts underway to increase the sensitivity of ROTSE abb+99 .
X-ray Observations: Diversity & Progenitors. GINGA identified a number of X-ray rich GRBs. BeppoSAX has found several such examples with some bursts lacking significant gamma-ray emission – the so-called X-ray flashes heise99 . We know very little about these X-ray transients. Could they be GRBs in a very dense environment (with red giant progenitors)? We need to take such transients more seriously and intensively followup on such bursts.
Another interesting finding from GINGA was the discovery of precursor soft X-ray emission m+91 . There is no simple explanation for this phenomenon in the current internal-external shock model. We suggest that the soft X-ray emission precursor is similar to the UV breakout of ordinary SNe. This hypothesis can be confirmed or rejected by obtaining the redshift to such bursts.
The X-ray rich GRB 981226 faa+2000 ; fkb+99 was marked with two additional peculiarities: a precursor emission and afterglow emission which is seemingly undetectable after about 12 hours but then rises rapidly before commencing decay. Above we alluded to the fact that massive stars do not have a single phase of mass loss but instead have a veritable history of mass loss (from birth to death). The X-ray observations of GRB 981226 could be accounted for in a model in which the progenitor has first a red supergiant wind followed by a blue supergiant wind.
Optical Observations: SN link, Short bursts & Geometry. The GRB–SN connection is best probed by optical observations. The value of optical observations has already been demonstrated by the current observations of GRB 980326 and 970228. Clearly, more observations are needed to establish this link. Once this link is established then one can undertake detailed spectroscopic studies of the SN with large ground-based telescopes and photometric studies with HST.
Offsets of GRBs and the morphology of the host galaxies will continue to be of great interest. Such observations will help us differentiate whether some GRBs come from nuclear regions or always from star-forming regions. Under the current paradigm, the discovery of GRBs coincident with elliptical galaxies would be a major surprise. On the other hand, one expects short bursts to arise in the halo of their galaxies and thus in this case no coincidence is expected. We expect HETE-2 to contribute significantly to these issues. Finally, polarization measurements offer a very convenient way to probe the geometry of the emitting region as has already been demonstrated from the discovery of polarization in GRB 990510 (e.g. lcg99 ; wvg+99 ).
Acknowledgments. Our research is supported by NASA and NSF. JSB holds a Fannie & John Hertz Foundation Fellowship, AD holds a Millikan Postdoctoral Fellowship in Experimental Physics, TJG holds a Fairchild Foundation Postdoctoral Fellowship in Observational Astronomy and RS holds Fairchild Foundation Senior Fellowship in Theoretical Astrophysics. The VLA is a facility of the National Science Foundation operated under cooperative agreement by Associated Universities, Inc. The W. M. Keck Observatory is operated by the California Association for Research in Astronomy, a scientific partnership among California Institute of Technology, the University of California and the National Aeronautics and Space Administration. It was made possible by the generous financial support of the W. M. Keck Foundation. |
warning/0002/physics0002006.html | ar5iv | text | # Spatiotemporally Localized Multidimensional Solitons in Self-Induced Transparency Media
## Abstract
”Light bullets” are multi-dimensional solitons which are localized in both space and time. We show that such solitons exist in two- and three-dimensional self-induced-transparency media and that they are fully stable. Our approximate analytical calculation, backed and verified by direct numerical simulations, yields the multi-dimensional generalization of the one-dimensional Sine-Gordon soliton.
The concept of multi-dimensional solitons that are localized in both space and time, alias ”light bullets” (LBs), was pioneered by Silberberg, and has since then been investigated in various nonlinear optical media, with particular emphasis on the question of whether these solitons are stable or not. For a second-harmonic generating medium, the existence of stable two- and three-dimensional (2D and 3D) solitons was predicted as early as in 1981, followed by studies of their propagation and stability against collapse, and of analogous 3D quantum solitons. In a nonlinear Schrödinger model both stable and unstable LBs were found and it was suggested that various models describing fluid flows yield stable 2D spatio-temporal solitons. Recently, the first experimental observation of a quasi-2D ”bullet” in a 3D sample was reported in Ref. .
In this letter we predict a new, hitherto unexplored type of LBs, obtainable by 2D or 3D self-induced transparency (SIT). SIT involves the solitary propagation of an electromagnetic pulse in a near-resonant medium, irrespective of the carrier-frequency detuning from resonance. The SIT soliton in 1D near-resonant media is exponentially localized and stable. In order to investigate the existence of ”light bullets” in SIT, i.e. solitons that are localized in both space and time, one has to consider a 2D or 3D near-resonant medium. Here we present an approximate analytical solution of this problem, which is checked by and in very good agreement with direct numerical simulations.
Our starting point are the two-dimensional SIT equations in dimensionless form
$`i_{xx}+_z𝒫`$ $`=`$ $`0`$ (2)
$`𝒫_\tau W`$ $`=`$ $`0`$ (3)
$`W_\tau +{\displaystyle \frac{1}{2}}(^{}𝒫+𝒫^{})`$ $`=`$ $`0.`$ (4)
Here $``$ and $`𝒫`$ denote the slowly-varying amplitudes of the electric field and polarization, respectively, $`W`$ is the inversion, $`z`$ and $`x`$ are respectively the longitudinal and transverse coordinates (in units of the effective absorption length $`\alpha _{\mathrm{eff}}`$), and $`\tau `$ the retarded time (in units of the input pulse duration $`\tau _p`$). The Fresnel number $`F`$ ($`F>0`$), which governs the transverse diffraction in 2D and 3D propagation, is incorporated in $`x`$ and the detuning $`\mathrm{\Delta }\mathrm{\Omega }`$ of the carrier frequency from the central atomic resonance frequency is absorbed in $``$ and $`𝒫`$. We have neglected polarization dephasing and inversion decay, considering pulse durations that are much shorter than the corresponding relaxation times. Eqs. (Spatiotemporally Localized Multidimensional Solitons in Self-Induced Transparency Media) are then compatible with the local constraint $`|𝒫|^2+W^2=1`$, which corresponds to conservation of the Bloch vector.
The first nontrivial question is to find a Lagrangian representation for these 2D equations, which is necessary for adequate understanding of the dynamics. To this end, we rewrite the equations in a different form, introducing the complex variable $`\varphi `$ defined as follows
$$\varphi \frac{1+W}{𝒫}=\frac{𝒫^{}}{1W}𝒫=\frac{2\varphi ^{}}{\varphi \varphi ^{}+1},W=\frac{\varphi \varphi ^{}1}{\varphi \varphi ^{}+1}.$$
(5)
Eqs. (3) and (4) can then be expressed as a single equation, $`\varphi _\tau +(/2)\varphi ^2+(1/2)^{}=0`$. Next, we define a variable $`f`$ so that $`\varphi 2f_\tau /(f)`$. In terms of $`f`$, the previous equation becomes $`f_{\tau \tau }(_\tau /)f_\tau +(1/4)||^2f=0.`$ This equation is equivalent to
$`f_\tau `$ $`=`$ $`{\displaystyle \frac{1}{2}}g`$ (7)
$`g_\tau `$ $`=`$ $`{\displaystyle \frac{1}{2}}^{}f,`$ (8)
with $`gf\varphi `$. Applying the same transformations to Eq. (2) yields
$$i_{xx}+_z2fg^{}=0.$$
(9)
The Lagrangian density corresponding to Eqs. (Spatiotemporally Localized Multidimensional Solitons in Self-Induced Transparency Media) and (9) can now be found in an explicit form,
$`(x,\tau )`$ $`=`$ $`{\displaystyle \frac{1}{4}}_x_x^{}+{\displaystyle \frac{i}{8}}(_z^{}_z^{}){\displaystyle \frac{i}{2}}\left(f^{}gfg^{}^{}\right)`$ (11)
$`{\displaystyle \frac{i}{2}}\left(f\dot{f}^{}\dot{f}f^{}\right){\displaystyle \frac{i}{2}}\left(g\dot{g}^{}\dot{g}g^{}\right).`$
Now we proceed to search for LB solutions. Before resorting to direct simulations, we obtain an analytical approximation of the solutions. The starting point for this approximation is the well-known soliton solution for 1D SIT (the Sine-Gordon soliton)
$`(\tau ,z)`$ $`=`$ $`\pm 2\alpha \text{sech}\mathrm{\Theta }`$ (13)
$`𝒫(\tau ,z)`$ $`=`$ $`\pm 2\text{sech}\mathrm{\Theta }\mathrm{tanh}\mathrm{\Theta }`$ (14)
$`W(\tau ,z)`$ $`=`$ $`\text{sech}^2\mathrm{\Theta }\mathrm{tanh}^2\mathrm{\Theta },`$ (15)
with $`\mathrm{\Theta }(\tau ,z)=\alpha \tau \frac{z}{\alpha }+\mathrm{\Theta }_0`$, and $`\alpha `$, $`\mathrm{\Theta }_0`$ arbitrary real parameters. Equation (13) is also called a $`2\pi `$-pulse, because its area $`_{\mathrm{}}^{\mathrm{}}(\tau ,z)𝑑\tau =\pm 2\pi `$.
Returning to the 2D SIT equations, we notice by straightforward substitution into Eqs. (Spatiotemporally Localized Multidimensional Solitons in Self-Induced Transparency Media) that a 2D solution with separated variables, in the form $`(\tau ,z,x)=_1(\tau ,z)_2(x)`$ (and similarly for $`f`$ and $`g`$), does not exist. To look for less obvious solutions, we first split equations (Spatiotemporally Localized Multidimensional Solitons in Self-Induced Transparency Media) into their real and imaginary parts, writing $`_1+i_2`$ and $`𝒫𝒫_1+i𝒫_2`$:
$`_{2xx}+_{1z}𝒫_1`$ $`=`$ $`0`$ (17)
$`_{1xx}_{2z}+𝒫_2`$ $`=`$ $`0`$ (18)
$`𝒫_{1\tau }_1W`$ $`=`$ $`0`$ (19)
$`𝒫_{2\tau }_2W`$ $`=`$ $`0`$ (20)
$`W_\tau +_1𝒫_1+_2𝒫_2`$ $`=`$ $`0.`$ (21)
In the absence of the $`x`$-dependence, these equations are invariant under the transformation $`(_1,𝒫_1)(_2,𝒫_2)`$. This suggests a 1D solution in which real and imaginary parts of the field and polarization are equal, $`_1=_2`$ and $`𝒫_1=𝒫_2`$, and such that the total field and polarization reduce to the SG solution (Spatiotemporally Localized Multidimensional Solitons in Self-Induced Transparency Media). Our central result is an approximate but quite accurate (see below) extension of this solution, applicable to the 2D SIT equations. In terms of the original physical variables it is given by
$`(\tau ,z,x)`$ $`=`$ $`\pm 2\alpha \sqrt{\text{sech}\mathrm{\Theta }_1\text{sech}\mathrm{\Theta }_2}\text{exp}(i\mathrm{\Delta }\mathrm{\Omega }\tau +i\pi /4)`$ (23)
$`𝒫(\tau ,z,x)`$ $`=`$ $`\pm \sqrt{\text{sech}\mathrm{\Theta }_1\text{sech}\mathrm{\Theta }_2}\{(\mathrm{tanh}\mathrm{\Theta }_1+\mathrm{tanh}\mathrm{\Theta }_2)^2+`$ (26)
$`{\displaystyle \frac{1}{4}}\alpha ^2C^4[(\mathrm{tanh}\mathrm{\Theta }_1\mathrm{tanh}\mathrm{\Theta }_2)^2`$
$`2(\text{sech}^2\mathrm{\Theta }_1+\text{sech}^2\mathrm{\Theta }_2)]^2\}^{1/2}\text{exp}(i\mathrm{\Delta }\mathrm{\Omega }\tau +i\mu )`$
$`W(\tau ,z,x)`$ $`=`$ $`[1\text{sech}\mathrm{\Theta }_1\text{sech}\mathrm{\Theta }_2\{(\mathrm{tanh}\mathrm{\Theta }_1+\mathrm{tanh}\mathrm{\Theta }_2)^2+`$ (29)
$`{\displaystyle \frac{1}{4}}\alpha ^2C^4[(\mathrm{tanh}\mathrm{\Theta }_1\mathrm{tanh}\mathrm{\Theta }_2)^2`$
$`2(\text{sech}^2\mathrm{\Theta }_1+\text{sech}^2\mathrm{\Theta }_2)]^2\}]^{1/2},`$
with
$`\mathrm{\Theta }_1`$ $`=`$ $`\alpha \tau {\displaystyle \frac{z}{\alpha }}+\mathrm{\Theta }_0+Cx`$ (30)
$`\mathrm{\Theta }_2`$ $`=`$ $`\alpha \tau {\displaystyle \frac{z}{\alpha }}+\mathrm{\Theta }_0Cx,`$ (31)
$`\mu `$ $``$ $`\mathrm{arctan}\left(𝒫_2/𝒫_1\right).`$ (32)
Here $`\alpha `$, $`\mathrm{\Theta }_0`$ and $`C`$ are real constants. Equations (Spatiotemporally Localized Multidimensional Solitons in Self-Induced Transparency Media) satisfy the two-dimensional SIT equations (17) and (18) and obey the normalization condition $`𝒫_1^2+𝒫_2^2+W^2=1`$. They reduce to the Sine-Gordon solution for $`C=0`$. The accuracy to which Eqs. (Spatiotemporally Localized Multidimensional Solitons in Self-Induced Transparency Media) satisfy Eqs. (19)-(21) is O($`\alpha C^2`$), which requires that $`|\alpha |C^21`$. This is the single approximation made. Numerical simulations discussed later on verify that Eq. (Spatiotemporally Localized Multidimensional Solitons in Self-Induced Transparency Media) indeed approximates the exact solution of Eq. (Spatiotemporally Localized Multidimensional Solitons in Self-Induced Transparency Media) to a high accuracy. In addition, we have checked that substitution of (Spatiotemporally Localized Multidimensional Solitons in Self-Induced Transparency Media) into the Lagrangian (11) and varying the resulting expression with respect to the parameters $`\alpha `$ and $`C`$ yields zero. This ”variational approach” is commonly used to obtain an approximate ”ansatz” solution to a set of partial differential equations in Lagrangian representation. Equations (Spatiotemporally Localized Multidimensional Solitons in Self-Induced Transparency Media) represents a light bullet, which decays both in space and time and is stable for all values of $`z`$. The latter follows directly from (23) and also from the Vakhitov-Kolokolov stability criterion.
Figs. 1-3 show the electric field and polarization, generated by direct numerical simulation of the 2D SIT equations (Spatiotemporally Localized Multidimensional Solitons in Self-Induced Transparency Media) at the point $`z=1000`$, using (Spatiotemporally Localized Multidimensional Solitons in Self-Induced Transparency Media) as an initial ansatz for $`z=0`$. To a very good accuracy (with a deviation $`<1\%`$), they still coincide with the initial configuration and analytic prediction (Spatiotemporally Localized Multidimensional Solitons in Self-Induced Transparency Media). The electric field has a typical shape of a 2D LB, localized in time and the transverse coordinate $`x`$, with an amplitude $`2\alpha `$ and a nearly sech-form cross-section in a plane in which two of the three coordinates $`\tau `$, $`z`$ and $`x`$ are constant. The ratio $`C/\alpha `$ determines how fast the field decays in the transverse direction. For $`|C/\alpha |1`$ (then $`|C|<1`$, as $`|\alpha |C^21`$), we have a relatively rapid decay in $`\tau `$ and slow fall-off in the $`x`$-direction, as is seen in Fig. 1. In the opposite case, $`|C/\alpha |1`$, the field decays more slowly in time and faster in $`x`$. The polarization field has the shape of a double-peaked bullet. Its cross-section at constant $`x`$ displays a minimum at $`\mathrm{\Theta }_{\text{min}}0`$, where $`|𝒫(\mathrm{\Theta }_{\text{min}})|0`$, and maxima at $`\mathrm{\Theta }_\pm =\pm \text{Arcosh}(\sqrt{2})`$, where $`|𝒫(\mathrm{\Theta }_\pm )|1`$. The field and polarization decay in a similar way, which is a characteristic property of SIT . Also the inversion decays both in time and in $`x`$, but to a value of $`1`$ instead of zero, corresponding to the atoms in the ground state at infinity. A numerical calculation of the field area at $`x=0`$ yields $`_{\mathrm{}}^{\mathrm{}}𝑑\tau |(\tau ,z,0)|=6.28\pm 0.052\pi `$, irrespective of $`z`$. By analogy with the SG soliton, one might thus name this a ”$`2\pi `$ bullet”.
We have also numerically obtained axisymmetric stable LBs in a 3D SIT medium, see Fig. 4. The 3D medium is described by Eqs. (Spatiotemporally Localized Multidimensional Solitons in Self-Induced Transparency Media) with the first one replaced by
$$i(_{rr}+r^1_r)+_z𝒫=0,$$
(33)
where $`r\sqrt{x^2+y^2}`$ is the transverse radial coordinate. Searching for an analytic 3D bullet solution in the transverse plane proves to be difficult. However, in the limit of either large or small $`r`$, an approximate analytic solution may be found. For large $`r`$, it again takes the form (Spatiotemporally Localized Multidimensional Solitons in Self-Induced Transparency Media), but now with $`\mathrm{\Theta }_1=\alpha \tau z/\alpha +\mathrm{\Theta }_0+Cr`$ and $`\mathrm{\Theta }_2=\alpha \tau z/\alpha +\mathrm{\Theta }_0Cr`$, where $`\alpha `$, $`\mathrm{\Theta }_0`$, and $`C`$ are constants, $`|\alpha |C^21`$, and it is implied $`r1/|C|`$. It is in sufficiently good agreement (deviations $`<5\%`$) with results of simulation of the 3D equations, using this solution as an initial ansatz. Comparison of Figs. 1 and 4 shows that the 2D and 3D bullets have similar shapes, but the 3D one decays faster in the radial direction for small $`r`$ than the 2D bullet in its transverse direction.
For constant $`\tau `$, the 2D and 3D bullets are localized in both the propagation direction $`z`$ and the transverse direction(s). One may also ask whether there exist SIT solitons which are traveling (plane) waves in $`z`$ and localized in $`x`$ (and $`y`$). Using a symmetry argument, it is straightforward to prove that they do not exist. Starting from the SIT equations (Spatiotemporally Localized Multidimensional Solitons in Self-Induced Transparency Media) (in 2D, the 3D case can be considered analogously) we adopt a plane-wave ansatz for $``$ and $`𝒫`$, changing variables as follows: $`x\sqrt{k}x`$ (assuming $`k>0`$), $`(\tau ,z,x)(\tau ,x)\mathrm{exp}(ikz)`$, $`𝒫(\tau ,z,x)k^1𝒫(\tau ,x)\mathrm{exp}(ikz)`$, and $`W(\tau ,z,x)k^1W(\tau ,x)`$. The equations for the real and imaginary parts of the field then become
$`_{2xx}_2𝒫_1`$ $`=`$ $`0`$ (35)
$`_{1xx}_1+𝒫_2`$ $`=`$ $`0,`$ (36)
with the equations for $`𝒫_\tau `$ and $`W_\tau `$ given by (19)-(21). Using the transformation $`(_1,𝒫_1)(_2,𝒫_2)`$, which leaves the last three equations invariant but changes the first two, one immediately finds that (Spatiotemporally Localized Multidimensional Solitons in Self-Induced Transparency Media) only admits the trivial solution $`_1=_2=𝒫_1=𝒫_2=0`$, $`W=1`$.
The observation of ”light bullets” in a SIT process requires high input power of the incident pulse and high density of the two-level atoms in the medium, in order to achieve pulse durations short compared to decoherence and loss times. These requirements are met e.g. for alkali gas media, with typical atomic densities of $`10^{11}`$ atoms/cm<sup>3</sup> and relaxation times $`50`$ ns, and for optical pulses generated by a laser with pulse duration $`\tau _p<0.1`$ ns. In order to include transverse diffraction, the incident pulse should be of uniform transverse intensity and satisfy $`\alpha _{\mathrm{eff}}d^2/\lambda <1`$ , where $`\lambda `$ and $`d`$ are its carrier wavelength and diameter respectively. The parameter $`\alpha `$ in the solution (Spatiotemporally Localized Multidimensional Solitons in Self-Induced Transparency Media), which determines the amplitude of the bullet and its decay in time, corresponds to $`\alpha \kappa _z\tau _pv_p`$, with $`\kappa _z`$ the wavevector component along the propagation direction $`z`$ and $`v_p`$ the velocity of the pulse in the medium, and can thus be controlled by the incident pulse duration and velocity. The parameter $`C\kappa _xL_x`$, where $`\kappa _x`$ is the transverse component of the wavevector and $`L_x`$ is the spatial transverse width of the pulse, is also controlled by the characteristics of the incident pulse and should satisfy the condition $`\kappa _z\kappa _x^2L_zL_x^21`$. For a homogeneous (atomic beam) absorber, the effective absorption length $`\alpha _{\mathrm{eff}}10^4`$ m<sup>-1</sup> and the Fresnel number $`F`$ can range from 1 to 100. The bullets then decay on a time scale of $`t110\tau _p10`$ ns and transverse length of $`x0.11`$ mm, which is well within experimental reach.
In conclusion, we predict the existence of fully stable ”light bullets” in 2D and 3D self-induced transparency media. The prediction is based on an approximate analytical solution of the multi-dimensional SIT equations and verified by direct numerical simulation of these PDE’s. Our results suggest an experiment aimed at detection of this ”bullet” in an SIT-medium and opens the road for analogous searches for ”light bullets” in other nonlinear optical processes, such as, e.g., stimulated Raman scattering, which is analogous to SIT.
M.B. acknowledges support from the Israeli Council for Higher Education. Support from ISF, Minerva and EU (TMR) is acknowledged by G.K. |
warning/0002/gr-qc0002030.html | ar5iv | text | # Untitled Document
Highly Excited Friedmann Universe <sup>a</sup><sup>a</sup>a Published in Physics of Atomic Nuclei, Vol. 62, No. 9, 1999, pp. 1524-1529. Translated from Yadernaya Fizika, Vol. 62, No.9, 1999, pp. 1625-1631.
V. V. Kuzmichev <sup>b</sup><sup>b</sup>b e-mail: vvkuzmichev@yahoo.com; specrada@bitp.kiev.ua
Bogolyubov Institute for Theoretical Physics, National Academy of Sciences of Ukraine, Metrolohichna St. 14b, Kiev, 03143 Ukraine
Abstract: A highly excited Friedmann universe filled with a scalar field and radiation has been considered. On the basis of a direct solution to the quantum-mechanical problem with a well-defined time variable, it has been shown that such a universe can have features (energy density, scale factor, Hubble constant, density parameter, matter mass, equivalent number of baryons, age, dimensions of large-scale fluctuations, amplitude of fluctuations of cosmic microwave background radiation temperature) identical to those of the currently observed Universe.
1. Available cosmological data suggest that, from the point of view of quantum theory, the currently observed Universe is likely to be in a highly excited state . This is confirmed by estimates of the number of the quantum state that corresponds to its averaged motion as a discrete unit \[1-4\]. In view of this, it is important to investigate cosmological systems featuring gravitational and matter fields and occurring in states with large quantum numbers by proceeding from a direct solution to the relevant quantum-mechanical problem.
A model of the Friedmann universe filled with a uniform scalar field has been proposed in . This model, which is appropriate for constructing a quantum theory, features a well-defined time variable. The reference frame was specified there with the aid of a subsidiary matter source in the form of radiation (relativistic matter of any nature) that was assumed to be initially present in the cosmological system along with a scalar field that forms a nonzero cosmological constant in the early universe. The evolution of the universe filled with not only scalar field but also with radiation differs from that which is realized in the absence of radiation. The main difference lies here in the emergence of a new region that is accessible to a classical motion and which is bounded by the potential barrier existing in the system of scalar and gravitational fields. A quantum universe involving a slowly varying scalar field and occurring in low-lying (quasistationary) states has been analyzed in , where it has been shown that the dynamical model proposed there is compatible with the currently prevailing ideas of the early Universe. In this study, we will consider a quantum Friedmann universe with large quantum numbers characterizing the possible physical states of the gravitational and matter fields involved. On the basis of a solution to the quantum mechanical problem, it is shown that the universe in highly excited states can have features (energy density, scale factor, Hubble constant, density parameter, matter mass, equivalent number of baryons) identical to those in the currently observed Universe. The problem of the age of the Universe and the possible new mechanism that could generate fluctuations of the metric due to a finite width of quasistationary states and to the anisotropy of cosmic microwave background radiation are discussed in the Appendices.
2. The wave function of the quantum Friedmann universe filled with radiation and a uniform scalar field $`\varphi `$ specified by the potential $`V(\varphi )`$ is determined by the Schrödinger equation
$$2i_T\mathrm{\Psi }=\left(_a^2\frac{2}{a^2}_\varphi ^2U\right)\mathrm{\Psi },$$
(1)
where $`T`$ is a privileged time coordinate related to the synchronous proper time $`t`$ by the coordinate condition $`dT/dt=1/a`$, $`a`$ being the scale factor, while
$$U=a^2a^4V(\varphi )$$
(2)
plays the role of the effective interaction potential in the cosmological system being considered. Here, we use the system of units in which $`l=\sqrt{2G\mathrm{}/3\pi c^3}=1`$ and $`\stackrel{~}{\varphi }=\sqrt{3c^4/8\pi G}=1.`$<sup>1</sup><sup>1</sup>1 In order to go over to the system of units where $`\mathrm{}=c=1`$, we must use the relation $`l\stackrel{~}{\varphi }^2=\sqrt{3/32\pi ^3}m_p`$ for the energy and the relation $`(\stackrel{~}{\varphi }/l)^2=(9/16)m_p^4`$ for the energy density, $`m_p`$ being the Planck mass.
Possible solutions to equation (1) are determined by the properties of the scalar field involved. If $`V(\varphi )`$ is everywhere positive definite, we can see that, in the variable $`a`$, the potential $`U`$ has the form of a barrier with height $`U_{max}=\frac{1}{4V}`$ and width $`\mathrm{\Delta }a=a_2(E)a_1(E)`$, where $`a_1(E)`$ and $`a_2(E)`$, $`a_1(E)<a_2(E)`$, are the classical turning points determined from the condition $`U=E`$. With the aid of the equation of motion for the field $`\varphi `$, it can be shown that, if the field satisfies the condition $`|_t^2\varphi ||\frac{dV}{d\varphi }|`$, the contribution of the operator $`\frac{2}{a^2}_\varphi ^2`$ can be approximated by the term $`\frac{a^4}{18}\left(\frac{1}{H}\frac{dV}{d\varphi }\right)^2`$, where $`H=\frac{_ta}{a}`$ is the Hubble constant; that is, this contribution is equivalent to the squared addition to the interaction $`U`$ of the gradient of the potential $`V(\varphi )`$. In this effective potential, stationary states cannot exist in the region $`aa_1`$. If, however, $`V(\varphi )1`$, quasistationary states with lifetimes exceeding the Planck time can exist within the barrier . The positions and widths of such states are determined by the solutions to equation (1) that satisfy the boundary condition in the form of a wave traveling toward greater values of $`a`$.
A general solution to equation (1) can be represented in the integral form
$$\mathrm{\Psi }(a,\varphi ,T)=_0^{\mathrm{}}𝑑E\text{e}^{\frac{i}{2}ET}C(E)\psi _E(a,\varphi ),$$
(3)
where the function $`C(E)`$ characterizes the $`E`$ distribution of the states of the universe at the instant $`T=0`$, while $`\psi _E(a,\varphi )`$ and $`E`$ are, respectively, the eigenfunctions and the eigenvalues for the equation
$`\left(_a^2+{\displaystyle \frac{2}{a^2}}_\varphi ^2+UE\right)\psi _E=0.`$ (4)
3. Let us now consider a quantum universe where $`\left|\frac{1}{V}\frac{dV}{d\varphi }\right|^21`$. A solution to equation (4) can then be represented as
$`\psi _E(a,\varphi )={\displaystyle _0^{\mathrm{}}}𝑑ϵ\phi _ϵ(a,\varphi )f_ϵ(\varphi ;E),`$ (5)
where $`\phi _ϵ`$ and $`ϵ`$ are, respectively, the eigenfunctions and eigenvalues for the operator $`\left[_a^2+U\right]`$ of the adiabatic approximation that correspond to continuum states at a fixed value of the field $`\varphi `$. The functions $`\phi _ϵ`$ can be normalized to the delta function $`\delta (ϵϵ^{})`$. Their form greatly depends on the value of $`ϵ`$. For $`ϵ<U_{max}`$, there are quasistationary states with $`ϵ=\stackrel{~}{ϵ}_nϵ_n+i\mathrm{\Gamma }_n`$ in the system , and the main contribution to the integral in (5) over the interval $`0<ϵ<U_{max}`$ comes from the values $`ϵϵ_n`$ and $`a<R`$, where $`Ra_2(ϵ)`$. Here, the wave function has the form
$$\phi _ϵ=A(ϵ)\phi _ϵ^{(0)},$$
(6)
where the function $`A(ϵ)`$ has a pole in the complex plane of $`ϵ`$ at $`ϵ=\stackrel{~}{ϵ_n}`$, while $`\phi _ϵ^{(0)}`$ is the solution over the interval $`0<a<R`$ that is regular at the point $`a=0`$, normalized to unity, and weakly dependent on $`ϵ`$. Proceeding in a way similar to that adopted in the theory of quasistationary states for short-range potentials , we can show that the quantity $`|A(ϵ)|^2`$ can be approximated by the delta function $`\delta (ϵϵ_n)`$ in the case of the potential (2) as well. In this approximation, expression (5) can eventually be reduced to the expansion
$$\psi _E(a,\varphi )=\underset{n}{}\phi _n(a,\varphi )f_n(\varphi ;E)+_{U_{max}}^{\mathrm{}}𝑑ϵ\phi _ϵ(a,\varphi )f_ϵ(\varphi ;E),$$
(7)
where $`\phi _n=\phi _{ϵ_n}^{(0)}`$ for $`0<a<R`$ and $`\phi _n=0`$ for $`a>R`$ (a state of this type was considered in ), while $`f_n(\varphi ;E)=_0^{\mathrm{}}𝑑a\phi _n^{}(a,\varphi )\psi _E(a,\varphi )`$. In the limit of an impenetrable barrier, the function $`\phi _n`$ reduces to the wave function of a stationary state with a definite value of $`ϵ_n`$. In the case of $`V1`$ considered here, the contribution of the integral to the expansion in (7) can be disregarded, and the quantities $`f_n(\varphi ;E)`$ can be interpreted as the amplitudes of the probability that the universe is in the state $`f_n(\varphi ;E)`$ with a given value of the field $`\varphi `$. They satisfy the set of differential equations
$$_\varphi ^2f_n+\underset{n^{}}{}K_{nn^{}}(\varphi ;E)f_n^{}=0,$$
(8)
where
$`K_{nn^{}}=\phi _n|_\varphi ^2|\phi _n^{}+2\phi _n|_\varphi |\phi _n^{}_\varphi +{\displaystyle \frac{1}{2}}\phi _n|a^2|\phi _n^{}\left(ϵ_n^{}E\right).`$ (9)
Over the time interval $`\mathrm{\Delta }T<\frac{1}{\mathrm{\Gamma }_n}`$, we can disregard the possibility of decay and consider the quasistationary state as a stationary state that arises in place of the quasistationary state when the decay probability tends to zero. The wave function $`\phi _n`$ of such a stationary state and the corresponding eigenvalue $`ϵ_n`$ can be found by perturbation theory by considering the interaction $`a^4V(\varphi )`$ as a small perturbation against $`a^2`$ (in the region $`a<a_1`$, we have $`a^2V<1`$). This yields
$`\phi _n`$ $`=`$ $`|n{\displaystyle \frac{V}{4}}[{\displaystyle \frac{1}{8}}\sqrt{N(N1)(N2)(N3)}|n2+`$ (10)
$`+`$ $`\sqrt{N(N1)}\left(N{\displaystyle \frac{1}{2}}\right)|n1`$
$``$ $`\sqrt{(N+1)(N+2)}\left(N+{\displaystyle \frac{3}{2}}\right)|n+1`$
$``$ $`{\displaystyle \frac{1}{8}}\sqrt{(N+1)(N+2)(N+3)(N+4)}|n+2]+O(V^2),`$
$`ϵ_n`$ $`=`$ $`ϵ_n^0{\displaystyle \frac{3}{4}}V\left[2N(N+1)+1\right]`$ (11)
$``$ $`{\displaystyle \frac{V^2}{4}}\left[{\displaystyle \frac{17}{2}}N^3+{\displaystyle \frac{51}{4}}N^2+{\displaystyle \frac{59}{4}}N+{\displaystyle \frac{21}{4}}\right]+O\left(V^3\right),`$
where $`N=2n+1`$, while $`ϵ_n^0=2N+1`$ and $`|n`$ are, respectively, an eigenvalue and the corresponding eigenfunction of the operator $`\left[_a^2+a^2\right]`$. Quasistationary states are realized for $`V<0.08=4.5\times 10^2m_p^4`$ and are characterized by values $`ϵ_n>2.6`$ and $`\mathrm{\Gamma }_n0.3`$ . The universe can undergo a tunnel transition to the region $`a>a_2(ϵ_n)`$ from any quasistationary state.
In the early universe, the quantity $`V(\varphi )`$ specifies the vacuum energy density, which determines the cosmological constant $`\mathrm{\Lambda }`$ at that era . In our Universe, the cosmological constant is very small. In the model featuring matter in the form of one scalar field, the reduction of the cosmological term can be described in terms of the potential $`V(\varphi (t))`$, which decreases with time . A similar behavior of $`V(\varphi (t))`$ is suggested by the results of investigations within inflationary models . It was shown in that, with a nonzero probability, a quantum universe initially filled with radiation and a scalar field whose potential $`V(\varphi (t))`$ decreases with time can evolve in the region bounded by the barrier. The expansion occurs here owing to transitions from lower to higher states via the interactions of the scalar and gravitational fields involved. In the zeroth approximation, the evolution of the universe can be described by considering transitions between unperturbed states $`|n`$ and $`|n+1`$ under the effect of the interaction $`a^4V`$ . In a more rigorous approach that takes into account variations in $`V(\varphi )`$, the universe expands through transitions between the states $`\phi _n`$ and $`\phi _{n+1}`$ due to the gradient of the potential $`V(\varphi )`$.
As the potential $`V(\varphi )`$ decreases, the number of quantum states in which the universe can occur increases. As a result, the competition between the tunneling processes and transitions from one state to another arises in the universe that has not had time to undergo tunneling. Since the decay probability decreases exponentially with decreasing $`V`$, the probability that the universe is excited to states with large quantum numbers $`n`$ in the course of time is nonzero. The probability that, in course of time, the universe in large-$`n`$ states will occur in the region outside the barrier is negligibly small, because small values of $`V`$ correspond to such states. In the limit $`V0`$, the universe is completely locked in the region within the barrier with $`ϵ_n`$ tending $`ϵ_n^0`$.
4. Let us investigate the properties of the universe in states with $`n1`$. In such states, $`|V|1`$ and $`\mathrm{\Gamma }_n0`$, and we can use the wave function (10) in calculating the quantities $`K_{nn^{}}`$ in (9). As before, we assume that the potential $`V(\varphi )`$ changes slowly in these states as well ($`\left|\frac{dV}{d\varphi }\right|1`$). The derivatives $`_\varphi \phi _n`$ in (9) can then be disregarded. Taking this into account and going over to the limit $`n1`$, we can then reduce equation (8) to the simplified form
$$_\varphi ^2f_n+\omega _n^2(\varphi ;E)f_n=0,$$
(12)
where
$$\omega _n^2(\varphi ;E)=2N^2\left[1\frac{E}{2N}2NV(\varphi )\right].$$
(13)
The potential of the field $`\varphi `$ will be chosen in the form $`V(\varphi )=\frac{m^2}{2}\varphi ^2`$. From the condition $`|V|1`$, it follows that the mass of the fields must be constrained by the condition $`m\frac{m_p^2}{|\varphi |}`$. In this case, only the region $`|\varphi |m_p`$ corresponds to extremely large masses $`mm_p`$. For masses satisfying the inequality $`mm_p`$, the field $`\varphi `$ can be determined over the entire interval $`0|\varphi |<\mathrm{}`$; in this case, equation (12) reduces to the Schrödinger equation for a harmonic oscillator. Under such conditions, the field $`\varphi `$ can oscillate near the minimum of the potential $`V(\varphi )`$, causing the production of particles.<sup>2</sup><sup>2</sup>2 At $`E=0`$, a similar mechanism leads to the production of particles by the inflaton field, which is identified with the scalar field $`\varphi `$ . In the case considered here, the field $`\varphi `$ for $`n1`$ states can be treated as an effective field obtained upon averaging over the internal degrees of freedom of real physical fields. Substituting the explicit expression for $`V(\varphi )`$ into (13) and introducing the new variable $`x=(2m^2N^3)^{1/4}\varphi `$, we arrive at the equation
$$_x^2f_n+(zx^2)f_n=0,$$
(14)
where $`z=\frac{\sqrt{2N}}{m}\left(1\frac{E}{2N}\right)`$. Equation (14) has solutions at $`z=2s+1`$, where $`s=0,1,2,\mathrm{}`$. Suppose that the universe occurs in a state with a definite quantum number $`s`$. The function $`f_n`$ can then be represented as
$$f_{ns}=\frac{1}{\sqrt{s!}}\left(B_n^{}\right)^sf_{n0},B_nf_{n0}=0,f_{n0}=\left(\frac{1}{\pi }\right)^{1/4}\text{e}^{x^2/2},$$
(15)
where $`f_{n0}`$ is the state vector of the universe in the $`n`$th state at $`s=0`$, while $`B_n^{}=\frac{1}{\sqrt{2}}(x_x)`$ and $`B_n=\frac{1}{\sqrt{2}}(x+_x)`$ are the operators that, respectively, create and annihilate particles in this state and which satisfy the conventional commutation relations $`[B_n,B_n^{}]=1`$ and $`[B_n,B_n]=[B_n^{},B_n^{}]=0`$. Since $`n1`$, small changes in $`n`$ do not affect the physical state of the universe where there are $`s`$ particles.
The condition of quantization of $`E`$ has the form
$$E=2N(2N)^{1/2}(2s+1)m.$$
(16)
For $`mm_p`$ and small $`s`$, the presence of the field $`\varphi `$ has no effect on the properties of the universe because, in this case, the approximate equality $`E2N`$ holds to a high precision, so that the universe is dominated by radiation. A transition from the radiation-dominated universe to the universe where matter (in the form of particles produced by the field $`\varphi `$) prevails occurs when, owing to an increase in the number of particles, the second term in (16) becomes commensurate with the first one. Let us consider the case where there are many particles in a highly excited $`n`$th state of the universe. For $`n1`$ and $`s1`$, the wave function has the form
$$\psi _{ns}(a,\varphi )=\phi _n(a)f_{ns}(\varphi ),$$
(17)
where
$$\phi _n(a)=\left(\frac{2}{N}\right)^{1/4}\mathrm{cos}\left(\sqrt{2N}a\frac{N\pi }{2}\right),$$
(18)
$$f_{ns}(\varphi )=\left(\frac{m(2N)^{3/2}}{4s}\right)^{1/4}\mathrm{cos}\left(\sqrt{2s}(2m^2N^3)^{1/4}\varphi \frac{s\pi }{2}\right).$$
(19)
This wave function corresponds to the semiclassical solution to equation (4) and is normalized to unity with allowance for the fact that the probability of finding the universe in the region $`a>a_2`$ is negligibly small.
Let us now calculate the energy density $`\rho _{tot}`$ for matter and radiation in the universe described by the wave function (17). The energy density for a classical field $`\varphi `$ and radiation is determined by the Einstein equation for the $`({}_{0}{}^{0})`$ component and is given by
$$\rho =\frac{2}{a^6}\pi _\varphi ^2+V+\frac{E}{a^4},$$
(20)
where $`\pi _\varphi `$ is the momentum canonically conjugate to the variable $`\varphi `$. Assuming that all the quantities in (20) are operator-valued, we set $`\rho _{tot}=\rho `$, where averaging is performed with the wave function (17). Disregarding the variances $`a^2a^2`$ and $`a^6a^6`$, we obtain
$$\rho _{tot}=\frac{2}{a^6}\pi _\varphi ^2+V+\frac{E}{a^4}.$$
(21)
Calculating the expectation values in (21), we arrive at the total energy density in the universe in the form of the sum of the energy densities for matter and radiation as in the general theory of relativity :
$$\rho _{tot}=\frac{193}{12}\frac{M_\varphi }{a^3}+\frac{E}{a^4}.$$
(22)
Here, $`M_\varphi =ms`$, and $`a=\sqrt{\frac{N}{2}}`$ is the scale factor in the universe occurring in an $`n1`$ state. The constant $`E`$ and the total mass $`M_\varphi `$ are related by the equation
$$E=4a\left[aM_\varphi \right].$$
(23)
If
$$a=M_\varphi ,$$
(24)
there is no radiation; we then have $`\rho _{tot}=\rho _{sub}`$ and
$$a=\left(\frac{193}{12}\frac{1}{\rho _{sub}}\right)^{1/2}.$$
(25)
5. Suppose that the system being considered occurs in an $`n1,s1`$ state and that it is characterized by the scale-factor value of $`a10^{28}\text{cm}`$. The condition in (24) will then be satisfied at $`M_\varphi 10^{56}\text{g}`$. This value coincides with the mass of matter in the observed part of our Universe . On the other hand, we set $`\frac{E}{a^4}`$ to the density of cosmic microwave background radiation energy in the present era, $`\rho _\gamma ^0=2.6\times 10^{10}`$ GeV/cm<sup>3</sup> . From (23), we then find that, for $`a=10^{28}\text{cm}`$, we have $`\frac{M_\varphi }{a}=10.7\times 10^5`$; that is, the equality in (24) holds to a high precision in this case as well. We note that, although the absolute value of $`E`$ is not small in the present era ($`E10^{117}`$ ), it is small in relation to $`a^210^{122}`$ (all estimates are presented here in the system of units where $`l=\stackrel{~}{\varphi }=1`$) and can be disregarded. It is the ratio of these quantities that determines the accuracy to which relation (24) holds.
The current values of the scale factor ($`aa_010^{28}\text{cm}`$) and of the mean energy density ($`\rho _{sub}\rho _010^5`$ GeV/cm<sup>3</sup>) satisfy relation (25). Associating the well-known equality $`a_0=\left(\frac{\mathrm{\Omega }_0}{\mathrm{\Omega }_01}\frac{1}{\rho _0}\right)^{1/2}`$ with relation (25), we find that the effective value of the density parameter is $`\mathrm{\Omega }_0=1.066`$; that is, the geometry of the universe with the above features is close to Euclidean geometry. It is characterized by the quantum-number values of $`na^210^{122}`$ and $`s\frac{a}{m}`$. Taking the proton mass for $`m`$, we obtain $`s10^{80}`$. The above value of $`n`$ complies with existing estimates for our Universe (see \[1-4\]), while $`s`$ is equal to the equivalent number of baryons . The value found for $`\mathrm{\Omega }_0`$ corresponds to the Hubble constant value of $`H_0^{theory}94`$ km/(s Mpc). This value of Hubble constant lies within the limits of experimental uncertainty of the value of $`H_0^{exp}=80\pm 17`$ km/(s Mpc), which was obtained with the aid of the Hubble cosmic telescope .
The approach developed here makes it possible to obtain realistic estimates for the age of the Universe (see Appendix 1), for the proper dimensions of the nonhomogeneities of matter which are consistent with those of the observed large-scale structure of the Universe, and for the amplitude of the fluctuations of the cosmic microwave background radiation temperature. The resulting estimate for the last quantity is close to the value extracted from experimental data (see Appendix 2).
6. The above numerical estimates of the parameters of the universe are of an illustrative character. If we nevertheless associate them with our Universe, it can be concluded that the values observed in the present era for the scale factor, the mass of matter in the Universe, and the density of the cosmic microwave background radiation energy satisfy relation (23). The zero-order approximation, which corresponds to setting $`\pi _\varphi ^2=0`$ in (20), leads to a very small value of $`\mathrm{\Omega }_00.1`$, but it is very close to the lower boundary of the uncertainty interval for $`\mathrm{\Omega }_0`$ . Although the potential $`V(\varphi )`$ undergoes only small variations in response to changes in the field $`\varphi `$, the field $`\varphi `$ itself changes fast, oscillating about the point $`\varphi =0`$, so that the approximation in which $`\pi _\varphi ^2=0`$ is invalid. The application of the present model in this approximation would result in the radiation-dominated universe; that is, it would not feature a mechanism capable of filling it with matter upon a slow descent of the potential $`V(\varphi )`$ to the equilibrium position, which corresponds to the true vacuum. In states of the universe that are characterized by large values of the quantum numbers, the kinetic term of the scalar field ensures the density parameter value close to unity.
Replacement of the entire set of actually existing massive fields by some averaged massive scalar field seems physically justified for states of the universe that have large values of $`n`$. We can see that, by and large, such an averaged field describes correctly the global features of our Universe. It effectively includes visible baryon matter and dark matter, which is globally manifested on cosmological scales via gravitational interaction. The status of the field $`\varphi `$ changes as we go over from one stage of universe evolution to another. In the early universe, the field $`\varphi `$ ensures a nonzero value of the vacuum-energy density (cosmological constant) due to $`V(\varphi )`$ values at which the equation for $`\phi _ϵ(a,\varphi )`$ admits nontrivial solutions in the form of quasistationary states. In a later era, when the field $`\varphi `$ descends to a minimum of the potential $`V(\varphi )`$ and begins to oscillate about this minimum, it appears to be a source of some averaged matter filling the visible volume of the universe, which has linear dimensions on the order of $`a`$.
APPENDIX 1
The time $`\frac{1}{H_0}`$ corresponding to the theoretical Hubble constant value $`H_0^{theory}`$ is equal to $`10^{10}`$ yr, and the age of the Universe calculated by the standard formula $`t=\frac{2}{3H}`$ is $`t_07\times 10^9`$ yr. This result is smaller than the expected value of $`t_0=(10÷20)\times 10^9`$ yr , but it is close to that which corresponds to $`H_0^{exp},t_08\times 10^9`$ yr, highlighting the problem of the age of the Universe. It should be borne in mind, however, that the age of the universe is calculated by the formula $`t=\frac{q}{H}`$, where $`q=\frac{1}{2}`$ (for the equation of state $`p=\frac{\rho }{3}`$ or $`\frac{2}{3}`$ (for $`p=0`$), which implies that $`a=bt^q`$, where the proportionality factor $`b`$ is independent of $`t`$. In the model of the universe filled with matter \[as represented by the field $`\varphi (t)`$\] and radiation, the factor $`b`$ depends on the variable $`\varphi (t)`$, which changes with time. The inclusion of this dependence leads to a greater value of $`t_0`$. By way of example, we indicate that, for $`2\sqrt{V}t1`$, the scale factor varies with t according to the law $`a(t)(2\sqrt{ϵ}t)^{1/2}`$, where $`ϵ=ϵ(\varphi (t))`$ . From the above, it follows that, in order to calculate the age of the universe, we must consider the transcendental equation $`t=\left[2H_t\mathrm{ln}ϵ\right]^1`$. Since $`ϵ>0`$ and since it increases with time, we have $`ϵ>0`$.
In order to perform a numerical estimation, we assume that, in the time interval being considered, $`ϵ`$ grows in proportion to a power of time: $`ϵt^\alpha `$. The age of the universe is then given by the expression $`t=\frac{1+\alpha }{2H}`$. Since we have $`ϵ10^{117}`$ and $`t10^{61}`$ (in $`l/c`$ units), the exponent $`\alpha `$ is $`\alpha 1.9`$ \[a similar power-law dependence follows from the relation $`a^2nϵ`$ and from the classical expression for $`a(t)`$\], whence we find that, if the Hubble constant set to its theoretical value $`H_0^{theory}`$, the age of the universe is $`t_015.2\times 10^9`$ yr. It should be noted that the above power-law dependence of $`ϵ`$ on $`t`$ ensures a correct relationship between the current value of the scale factor $`a_0`$ and the age $`t_0`$ of the universe: $`a_0t_010^{61}`$.
APPENDIX 2
A quasistationary state with a small, but finite value of the width $`\mathrm{\Gamma }`$ does not possess a definite value of $`ϵ`$. This uncertainty we denote it by $`\delta ϵ`$ can serve as source of fluctuations of the metric. Let us demonstrate this explicitly. By associating $`ϵ+\delta ϵ`$ with the scale factor $`a+\delta a`$ and by using the solution to the Einstein equation in the region $`aa_1`$ from , we find that the amplitude of fluctuations of the scale factor can be represented as
$$\frac{\delta a}{a}=\frac{1}{4}\frac{\delta ϵ/ϵ}{1\text{tanh}\sqrt{V}t/2\sqrt{Vϵ}}.$$
$`(A.1)`$
Since $`\delta ϵ\mathrm{\Gamma }`$, the fluctuations $`\delta a`$ that were generated at the early stage of the evolution of the Universe will take the greatest values. In order to estimate them, we adopt the values of $`V=0.08,ϵ=2.6`$, and $`\delta ϵ0.3`$, corresponding to the time $`t1`$ . We then have $`\frac{\delta a}{a}0.04`$. Since the dimension of large-scale fluctuations changed in direct proportion $`a(t)`$, this relation has remained valid up to the present time . Taking this into account, we find that $`\delta a130`$ Mpc for the current value of $`a10^{28}`$ cm. On the order of magnitude, the above value corresponds to the scale of superclusters of galaxies . Smaller values of $`\delta ϵ`$ are peculiar to quantum states with smaller $`V`$. The fluctuations $`\delta a`$ corresponding to them are smaller than those presented above and are expected to manifest themselves against the background of the large-scale structure. They can be associated with clusters of galaxies, galaxies themselves, and clusters of stars. Thus, the conclusions drawn on the basis of the model considered here are in line with the generally accepted concept that galaxies, their clusters, and other structures in the Universe are macroscopic manifestations of quantum fluctuations that have grown considerably .
Let us estimate the amplitude of fluctuations of the cosmic microwave background radiation temperature, $`\delta T/T`$. By using the relation $`ϵ=\rho _\gamma a^4`$, where $`\rho _\gamma =(\pi ^2/15)T^4`$ is the density of the cosmic microwave background radiation energy, we obtain
$$\frac{\delta T}{T}=\frac{1}{4}\frac{\delta ϵ}{ϵ}\frac{\delta a}{a}.$$
$`(A.2)`$
For $`\sqrt{V}t1`$, it follows from (A.1) and (A.2) that
$$\frac{\delta T}{T}\frac{t}{2\sqrt{ϵ}}\frac{\delta a}{a}.$$
$`(A.3)`$
For the time $`t10^5`$ yr ( $`t10^{56}`$ in $`l/c`$ units), which corresponds to the recombination of primary plasma (separation of radiation from matter), it can be found that, for the observed value of $`ϵ=5.2\times 10^{117}`$ (in which case $`\sqrt{V}t\frac{t}{2\sqrt{ϵ}}0.7\times 10^3`$), the sought amplitude of fluctuations of cosmic microwave background temperature radiation at $`\frac{\delta a}{a}0.04`$ can be estimated as
$$\frac{\delta T}{T}2.8\times 10^5.$$
$`(A.4)`$
Upon recombination, the fluctuations of the temperature undergo no changes; therefore, measurement of the quantity $`\delta T/T`$ for the present era furnishes information about the Universe at the instant of last interaction of radiation with matter. The estimate in (A.4) is in good agreement with experimental data on cosmic microwave background radiation (see ).
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warning/0002/gr-qc0002078.html | ar5iv | text | # Einstein-Yang-Mills Isolated Horizons: Phase Space, Mechanics, Hair and Conjectures
## I Introduction
The nonperturbative quantum geometry program also known as “loop quantum gravity”, has met recently with substantial success in obtaining a calculation of the statistical mechanical entropy of a non-rotating black hole that accounts for its phenomenological identification with $`\frac{1}{4}A`$ . In so doing it was necessary to introduce a complete classical Hamiltonian treatment of Black Holes. This was accomplished by generalizing and properly defining the sector of the theory that is going to be treated. This work was guided by the need to start with a well defined action that would be differentiable in the sector under consideration. This lead those authors to the specialization of the notion of Trapping Horizons of Hayward’s to that of ‘Isolated Horizons’. Physically the idea is to represent “horizons in internal equilibrium and decoupled from what is outside”.
The zeroth and first laws of black hole mechanics refer to equilibrium situations and small departures therefrom. Therefore, in the standard treatments one restricts oneself to stationary space-times admitting event horizons and perturbations off such space-times. The isolated horizons (IH) framework, which is tailored to more general physical situations was introduced in and the corresponding zeroth and first laws of black hole mechanics were established. . This framework generalizes the treatment of black hole mechanics in two directions. First, the notion of event horizons is replaced by that of ‘isolated horizons’ which can be defined quasi-locally, unlike the former which can only be defined retroactively, after having access to the entire space-time history. Second, the underlying space-time need not admit any Killing field; isolated horizons need not be Killing horizons. The static event horizons normally used in black hole mechanics are special cases of isolated horizons. Moreover, because one can now admit gravitational and matter radiation, there are many more examples. In particular, while the manifold $`𝒮`$ of static space-times admitting event horizons in the Einstein-Maxwell theory is finite dimensional, the manifold $``$ of space-times admitting isolated horizons is infinite dimensional .
The resulting formalism is then, not only a step in the construction of the quantum theory for the sector but also a new tool for studying classical aspects of black holes. When restricted to the Static sector of the theory, it leads, for example, to an improvement in the treatment of the physical process version of the first law of black hole mechanics . The formalism has so far only been applied to the Einstein-Maxwell system (with and without dilaton field), of which the known exact solutions (the Reissner Nordstrom solutions in EM and the so called Gibbons-Maeda solutions in EMD) are particular examples.
In this work we will explore the extent to which this formalism can be extended to the Einstein-Yang-Mills theory where, in the static spherically symmetric sector, one finds the so called Colored Black Hole solutions . The motivation is many-fold. First, we are interested in studying the robustness of the isolated horizons formalism in its ability to treat other theories, specially those where ‘hairy’ solutions are know to exist. From the isolated horizons perspective, this poses a special challenge, since now there is an apparent tension given by the mismatch between the number of conserved charges at infinity and at the horizon. In the colored black hole solutions, the only non-zero ‘charge’ at infinity is the ADM mass, but the gauge field is non-vanishing at the horizon and it contributes to the ‘horizon magnetic charge’. Is the isolated horizon framework robust enough to deal with this situation and to resolve this tension? Second, we would like to know whether the general formalism allows us to learn new facts about the Static Spherically Symmetric (SSS) solutions. In particular, we want to understand what differences, if any, appear when we treat static but unstable black hole solutions, which is the case for the Abelian magnetic solutions and the higher $`n`$ colored black holes. Finally, we want to investigate whether the formalism can –as can be expected given the fact that it yields a satisfactory treatment of the inner boundary in Hamiltonian terms– allow us to take the limit when the horizon area goes to zero and connect the black hole solutions with the regular, solitonic, solutions that are known to also exist in this theory.
We will in fact find answers to all of these questions and puzzles. We find however, some subtleties that need to be stressed. In particular, the fact that there are no Spherically Symmetric solutions for some values of the horizon parameters, poses a challenge to the isolated horizon framework, since there seems to be no canonical value for the Horizon Mass of the black hole in this situation. Thus, if the EYM system is to be in the same ‘status’ as Einstein-Maxwell system, some yet unknown part of the current scenario would have to yield to the tension mentioned above. The most natural resolution would entail the validity of a ‘uniqueness conjecture’ for Static solutions and a ‘completeness conjecture’ for the existence of Stationary solutions, that we will put forward. Their validity would guarantee the complete consistency of the formalism.
On the other hand, we will show some new results regarding SSS solutions as a nontrivial ‘application’ of the formalism as it allows us to predict a relation between the ADM mass of a static black hole solution, its Horizon mass and the ADM mass of the corresponding solitonic solutions. These relations can be corroborated by numerical computations (consistent with the reported results in the literature). Therefore, these ‘coincidences’ can be viewed in a sense, as a check on the formalism.
This paper is organized as follows. In Section II, we recall some basic facts about the EYM system and about the known SSS solutions, both Abelian and non-Abelian. In Sec. III we specify the isolated horizon boundary conditions that we impose, making the necessary adjustments to incorporate non-Abelian gauge fields. Section IV deals with the action principle of the theory in the presence of the horizon as an inner boundary, and the specification of the phase space of the theory. In Section V we consider the definition of surface gravity in the absence of a Killing field and we show the zeroth law of BH mechanics for general isolated horizons. The definition of horizon mass together with the first law are studied in Section VI. Section VII is devoted to the completeness conjecture and some of its implication for stationary solutions. In Section VIII, we return to the study of Static Spherically Symmetric solution to the EYM equations from the perspective of isolated black holes, and show some new results. Finally, we end with a discussion in Section IX.
Throughout the paper, we use units in which $`c=G=1`$, and the abstract index notation of Penrose’s , except in Sec. VI where we use differential forms.
## II Einstein-Yang-Mills and Static Black Holes
This section has two parts. In the first one, we recall some basic facts about the Einstein-Yang-Mills system. In the second part, we briefly review the Static Spherically Symmetric solutions to the EYM equations focusing on black hole solutions.
### A Einstein-Yang-Mills System
In the Einstein-Yang-Mills system, the gravitational part of the action, $`S_{\mathrm{Grav}}`$ is given by:
$$S_{\mathrm{Grav}}=\frac{1}{16\pi }_𝐌\sqrt{g}R\mathrm{d}^4x,$$
(1)
and the matter part of the action is given by:
$$S_{\mathrm{YM}}(𝐀)=\frac{1}{16\pi g_{\mathrm{YM}}^2}_𝐌\sqrt{g}[𝐅_{ab}^i𝐅_i^{ab}]\mathrm{d}^4x,$$
(2)
where the abstract indices $`a,b,\mathrm{}`$ denote space-time objects and the indices $`i,j,\mathrm{}`$ are internal indices in the Lie algebra of the gauge group $`G`$. In this paper we shall consider $`G=SU(2)`$. The field strength $`𝐅_{ab}`$ is given by $`𝐅_{ab}^i=2_{[a}𝐀_{b]}^i+ϵ_{}^{i}{}_{jk}{}^{}𝐀_a^j𝐀_b^k`$, that is, the curvature of the Lie algebra valued one form $`𝐀_a^i`$. The total action $`S_{\mathrm{Tot}}`$ is given by
$$S_{\mathrm{Tot}}=S_{\mathrm{Grav}}+S_{\mathrm{YM}}.$$
(3)
We remind the reader that in the case of a Non-Abelian Yang Mills theory there is dimension-full parameter $`g_{\mathrm{YM}}`$ that, unlike the Abelian case, can not be absorbed in the “gauge” fields. This endows the full theory with a natural scale at the classical level i.e. $`g_{\mathrm{YM}}^{}{}_{}{}^{2}`$ has dimensions of mass.
The equations of motion that follow from $`S_{\mathrm{Tot}}`$ are:
$`D_a𝐅^{iab}=0,`$ (4)
$`R_{ab}=2\left(𝐅_{ac}^i𝐅_{ib}^{}{}_{}{}^{c}{\displaystyle \frac{1}{4}}g_{ab}𝐅^2\right),`$ (5)
where $`𝐅^2=𝐅_{ab}^i𝐅_i^{ab}`$, and $`D_a`$ is the generalized covariant derivative defined by $`𝐀`$. Furthermore we have the Bianchi identity for the Yang Mills sector:
$`D_{[c}𝐅_{ab]}^i=0,`$ (6)
The dual field tensor is given by $`{}_{}{}^{}𝐅_{ab}^{i}=\frac{1}{2}ϵ_{ab}^{}{}_{}{}^{cd}𝐅_{cd}^i`$, where $`ϵ_{abcd}`$ is the canonical volume element associated with $`g_{ab}`$.
Note that in contrast with the Einstein-Maxwell system, where, under appropriate fall-of conditions, the conserved electric and magnetic charges can be defined at infinity, the naive expressions,
$$\widehat{Q}^i:=\frac{1}{4\pi }_S_{\mathrm{}}{}_{}{}^{}𝐅_{}^{i},\widehat{P}^i:=\frac{1}{4\pi }_S_{\mathrm{}}𝐅^i,$$
(7)
fail to be gauge invariant. It is only in the presence of a globally defined isometry that the conserved charges might be invariantly defined (i.e. a natural gauge might be chosen). We can nevertheless define new gauge invariant quantities, for any two-sphere $`S`$ as follows,
$$Q_S:=\frac{1}{4\pi }_S|{}_{}{}^{}𝐅|,P_S:=\frac{1}{4\pi }_S|𝐅|,$$
(8)
where $`|𝐅|_{ab}`$ is the two form defined in the following way: we take $`ϵ_{ab}`$ the area two form associated with the 2-sphere $`S`$ and define $`f^i=𝐅_{ab}^iϵ^{ab}`$. Then $`|𝐅|_{ab}=\sqrt{(f^i)^2}ϵ_{ab}`$. $`|{}_{}{}^{}𝐅|_{ab}`$ is analogously defined. In what follows, we shall refer to $`Q_S`$ and $`P_S`$ as the electric and magnetic charge contained ‘within’ $`S`$ respectively.
### B Static Solutions
Static Spherically Symmetric (SSS) solutions to the EYM equations representing black hole space-times are known to exist in different situations (for a recent review see and references therein).
A standard parameterization for the metric and gauge potential is given by,
$`\mathrm{d}s^2`$ $`=`$ $`\left(1{\displaystyle \frac{2m(r)}{r}}\right)e^{2\delta (r)}\mathrm{d}t^2+\left(1{\displaystyle \frac{2m(r)}{r}}\right)^1\mathrm{d}r^2+r^2\mathrm{d}\mathrm{\Omega }^2,`$ (9)
$`𝐀`$ $`=`$ $`a\tau _3\mathrm{d}t+b\tau _3\mathrm{d}r+(w\tau _1+d\tau _2)\mathrm{d}\theta +(\mathrm{cot}\theta \tau _3+w\tau _2d\tau _1)\mathrm{sin}\theta \mathrm{d}\varphi ,`$ (10)
where $`a,b,w`$ and $`d`$ are functions of $`(r,t)`$.
One can then look for either static regular solutions by requiring $`m(r)<r/2`$ for all $`r0`$ or for static black hole solutions with horizon at $`r=r_H`$ by requiring $`m(r_H)=r_H/2`$ and $`m(r)<r/2`$ for all $`rr_H`$. Solutions of both types are found in the purely magnetic sector, for which, with a further gauge choice one can set $`a=b=d=0`$, and $`w`$ becomes a function of $`r`$ only.
The regular or solitonic solutions compose a discrete set parameterized by the number of nodes of the function $`w(r)`$ and are characterized by their ADM mass whose scale is set by $`g_{\mathrm{YM}}^{}{}_{}{}^{2}`$ .
On the other hand, for every value of $`r_H`$, one finds also two classes of black hole solutions: Abelian and non-Abelian.
#### 1 Abelian Solutions
The first class of solutions is given by what are essentially Abelian solutions embedded in SU(2). Within the Abelian sector we have either electrically charged or magnetically charged solutions. Unlike Einstein-Maxwell theory where the well known Maxwell duality exists, in EYM there is no such duality and one is not allowed to treat them on ‘equal footing’.
The electrically charged solutions with electric charge $`Q`$ are given by the standard solutions that can be described by choosing $`a0`$ and $`b=d=w=0`$ in (10). The YM potential is of the form $`𝐀=\frac{Q}{r}\tau _3\mathrm{d}t`$ and the metric is given by (9) with the functions $`m(r)=M(Q^2)/2r`$ and $`\delta =0`$. Note that these represent a two parameter family of ‘RN solutions’ with parameters $`M`$ and $`Q`$.
The magnetically charged solutions with magnetic charge $`P`$ are precisely the Reissner-Nordstrom solutions given by $`w=0`$, $`m(r)=M(P^2)/2r`$ and $`\delta =0`$. Since we are considering the Abelian case, the magnetic charge $`P`$ is not arbitrary but can take only one value $`P=1`$. Note that the YM field strength, in the ‘magnetic’ sector of the EYM theory takes the form,
$$𝐅=w^{}\tau _1\mathrm{d}r\mathrm{d}\theta +w^{}\tau _2\mathrm{sin}\theta \mathrm{d}r\mathrm{d}\varphi (1w^2)\tau _3\mathrm{sin}\theta \mathrm{d}\theta \mathrm{d}\varphi .$$
(11)
Thus, for the RN solutions where $`w=0`$, we get from (11) that $`P=1`$ for any sphere containing the black hole. The magnetically charged solutions are then parameterized by only one charge, namely, the ADM mass $`M`$.
These two sectors share the $`Q=P=0`$ solution corresponding to the Schwarzschild solution. One can also construct dyonic solutions with both electric charge $`Q`$ and unit magnetic charge. In all these solutions, one has to satisfy the inequality $`r_H^2Q^2+P^2`$, in order to have black hole solutions with no naked singularities.
#### 2 Non-Abelian Colored Black Holes
These solutions correspond to the purely magnetic case, where, for each value of the horizon area the equations have a discrete number of solutions which are strictly Non-Abelian in nature (i.e. do not exit in the Abelian regime). These are labeled by an integer $`n`$ that represents the number of nodes of the function $`w(r)`$. The lowest mode, $`n=0`$, represents the Schwarzschild solution. Therefore, the solution can be completely parameterized by two numbers $`(a_\mathrm{\Delta },n)`$, the horizon area and the integer $`n`$. All this solutions, for $`n>0`$ are unstable under perturbations .
On the other hand it is known that there are no nontrivial dyonic solutions (i.e solutions with electric and magnetic fields) in the spherically symmetric sector.
From a historical perspective, these were the first examples of ‘hairy black holes’ . This is because the electric and magnetic charges (7) are both zero, so the only parameter at infinity is the ADM mass. If the no-hair conjecture were valid for the EYM system, the specification of $`M_{\mathrm{ADM}}`$ would suffice to characterize the solution completely. However, this is not the case, since for a given value of the ADM mass, there exist a countable number of different solutions, labeled by $`n`$. Equivalently we can label these solutions by the ADM mass $`M_{\mathrm{ADM}}`$ and $`n`$.
Even when the charges at infinity are not enough to specify static black holes uniquely, one might still hope to have quasi-local quantities defined at the horizon, that are in a sense, good coordinates for the manifold $`𝒮`$ of static solutions. Indeed, we shall put forward in the following sections a ‘quasi-local uniqueness conjecture’ ($`C1`$): All static BH solutions are characterized by its horizon parameters arising from the ‘isolated horizon’ framework. Let us refer to these quantities defined at the horizon as ‘quasi-local parameters’. In theories where no hair is present, as is the case of the Einstein-Maxwell-Dilaton system, the number of ‘quasi-local parameters’ equals the number of parameters at infinity labeling the static solutions . Thus, stating a uniqueness conjecture in this theory is insensitive as to whether one is postulating it in terms of quantities at infinity (the standard viewpoint), or in terms of ‘quasi-local charges’. Our proposal is that, for general theories, one should state the postulate in terms of purely quasi-local quantities. In the EYM system, the quasi-local charges are $`a_\mathrm{\Delta }`$, the horizon area, $`Q_\mathrm{\Delta }`$ and $`P_\mathrm{\Delta }`$, the horizon electric and magnetic charges respectively. In this case the first conjecture $`C1`$ reads: Given a triple of parameters $`(a_\mathrm{\Delta },Q_\mathrm{\Delta },P_\mathrm{\Delta })`$ for which a SSS solutions exists, then the solution is unique. However, note an apparent tension in this suggestion: Given the mismatch of the number of parameters at the horizon and at infinity, there is room for inconsistency when formulating, say, the laws of thermodynamics. This is because the number of ‘independent’ parameters is different, when one considers charges at infinity or quasi-local parameters for, say, SSS colored solutions. As we shall see in Section VI, the nature of the problem can be made precise within the isolated horizons framework, and some ideas can be put forward to understand the origin of the difficulty.
For the convenience of the reader, in the next section we shall recall the notion of isolated horizons as defined in and explore some of its consequences for EYM system of interest to this paper.
## III Boundary Conditions and Consequences
Let us recall the notion of isolated horizons $`\mathrm{\Delta }`$ in general, and include in its definition the relevant modifications to incorporate the Einstein-Yang-Mills system. The basic boundary conditions defining $`\mathrm{\Delta }`$ are the same as those introduced in .
Let us begin by recalling some notation. Fix any null surface $`𝒩`$, topologically $`S^2\times R`$, and consider foliations of $`𝒩`$ by families of spatial 2-spheres. Given a foliation, we parameterize its leaves by $`v=\mathrm{const}`$ such that $`v`$ increases to the future and set $`n_a=_av`$. Under a reparametrization $`vF(v)`$, we have $`n_aF^{}(v)n_a`$ with $`F^{}(v)>0`$. Thus, every foliation comes equipped with an equivalence class $`[n_a]`$ of normals $`n_a`$ related by rescalings which are constant on each leaf.<sup>*</sup><sup>*</sup>*These 1-form fields $`n_a`$ are defined intrinsically on $`𝒩`$. We can extend each $`n_a`$ to the full space-time uniquely by demanding that the extended 1-form be null. However, in this paper, we will not need this extension. Also recall that, given any one $`n_a`$, we can uniquely select a vector field $`\mathrm{}^a`$ which is normal to $`𝒩`$ and satisfies $`\mathrm{}^an_a=1`$. (Thus, $`\mathrm{}^a`$ is future-pointing.) If we change the parameterization, $`\mathrm{}^a`$ transforms via: $`\mathrm{}^a(F^{}(v))^1\mathrm{}^a`$. Thus, given a foliation, we acquire an equivalence class $`[\mathrm{}^a,n_a]`$ of pairs, $`(\mathrm{}^a,n_a)`$, of vector fields and 1-forms on $`𝒩`$ subject to the relation $`(\mathrm{}^a,n_a)(G^1\mathrm{}^a,Gn_a)`$, where $`G`$ is any positive function on $`𝒩`$ which is constant on each leaf of the foliation. Given a pair $`(\mathrm{}^a,n_a)`$ in the equivalence class, we introduce a complex vector field $`m^a`$ on $`𝒩`$, tangential to each leaf in the foliation, such that $`m^a\overline{m}_a=1`$. (By construction, $`m^a\mathrm{}_a=m^an_a=0`$ on $`𝒩`$.) The vector field $`m^a`$ is unique up to a phase factor. With this structure at hand, we now look at the main Definition.
Definition: The internal boundary $`\mathrm{\Delta }`$ of a space-time $`(𝐌,g_{ab})`$ will be said to represent a non-rotating isolated horizon provided the following conditions holdThroughout this paper, the symbol $`\widehat{}=`$ will denote equality at points of $`\mathrm{\Delta }`$. For fields defined throughout space-time, an under-arrow will denote pull-back to $`\mathrm{\Delta }`$. The part of the Newman-Penrose framework used in this paper is summarized in the Appendices A and B of .:
* (i) Manifold conditions: $`\mathrm{\Delta }`$ is a null surface, topologically $`S^2\times R`$.
* (ii) Dynamical conditions: All field equations hold at $`\mathrm{\Delta }`$.
* (iii) Main conditions: $`\mathrm{\Delta }`$ admits a foliation such that the Newman-Penrose coefficients associated with the corresponding direction fields $`[\mathrm{}^a,n_a]`$ on $`\mathrm{\Delta }`$ satisfy the following conditions:
(iii.a) $`\rho \widehat{}=\overline{m}^am^b_a\mathrm{}_b`$, the expansion of $`[\mathrm{}^a]`$, vanishes on $`\mathrm{\Delta }`$.
(iii.b) $`\lambda \widehat{}=\overline{m}^a\overline{m}^b_an_b`$ and $`\pi \widehat{}=\mathrm{}^a\overline{m}^b_an_b`$ vanish on $`\mathrm{\Delta }`$ and the expansion $`\mu :=m^a\overline{m}^b_an_b`$ of $`n_a`$ is negative For simplicity, in this paper we focus on black-hole-type horizons rather than cosmological ones. To incorporate interesting cosmological horizons, one has to weaken this condition and allow the possibility that $`\mu `$ is everywhere positive on $`\mathrm{\Delta }`$. See . and constant on each leaf of the foliation.
* (iv) Conditions on matter: The Yang-Mills field $`𝐅`$ is such that
$$|\mathrm{Re}\varphi _1|,\mathrm{and}|\mathrm{Im}\varphi _1|,$$
(12)
are constant on each leaf of the foliation introduced in condition (iii). (Recall that $`\varphi _1^i\widehat{}=\frac{1}{2}m^a\overline{m}^b(𝐅i{}_{}{}^{}𝐅)_{ab}^i`$) where $`|\mathrm{Re}\varphi _1|:=\sqrt{_i(\mathrm{Re}\varphi _1^\mathrm{i})(\mathrm{Re}\varphi _1^\mathrm{i})}`$, and $`|\mathrm{Im}\varphi _1|`$ is analogously defined.
The first two conditions are quite tame: (i) simply asks that $`\mathrm{\Delta }`$ be null and have appropriate topology while (ii) is completely analogous to the dynamical condition imposed at infinity. As the terminology suggests, (iii.a) and (iii.b) are the most important conditions. Note first that, if a pair $`(\mathrm{}^a,n_a)`$ in the equivalence class $`[\mathrm{}^a,n_a]`$ associated with the foliation satisfies these conditions, so does any other pair, $`((G(v))^1\mathrm{}^a,G(V)n_a)`$. Thus, the conditions are well-defined. They are motivated by the following considerations. Condition (iii.a) captures the idea that the horizon is isolated without having to refer to a Killing field. In particular, it implies that the area of each 2-sphere leaf in the foliation be the same. We will denote this area by $`a_\mathrm{\Delta }`$ and define the horizon radius $`r_\mathrm{\Delta }`$ via $`a_\mathrm{\Delta }=4\pi r_\mathrm{\Delta }^2`$.
Condition (iii.b) has three sets of implications. First, one can show that if, as required, one can find a foliation of $`\mathrm{\Delta }`$ satisfying (iii.b), that foliation is unique. (In the SSS family, as one might expect, this condition selects the foliation to which the rotational Killing fields are tangential.) Second, it implies that the imaginary part of (the Newman-Penrose Weyl component) $`\mathrm{\Psi }_2`$, which captures angular momentum, vanishes and thus restricts us to non-rotating horizons. Third, the requirement that the expansion $`\mu `$ of $`n^a`$ be negative implies that $`\mathrm{\Delta }`$ is a future horizon rather than past horizon . Finally, consider the spherical symmetry requirement on the Yang-Mills field component $`\varphi _1^i`$ given by condition (iv). While this condition is a strong restriction, it can be motivated by analogy with the Einstein-Maxwell case. (For further motivation and remarks on these conditions, see .)
Since these conditions are local to $`\mathrm{\Delta }`$, the notion of an isolated horizon is quasi-local; in particular, one does not need an entire space-time history to locate an isolated horizon. Furthermore, the boundary conditions allow for the presence of radiation in the exterior region, thus, space-times admitting isolated horizons need not admit any Killing field . Indeed, the manifold $``$ of solutions to field equations admitting isolated horizons is infinite dimensional ).
In spite of this generality, boundary conditions place strong restrictions on the structure of various fields at $`\mathrm{\Delta }`$. Let us begin with conditions on the Yang-Mills field. The stress-energy tensor $`T_{ab}`$ of $`𝐅`$ satisfies the dominant energy condition. Hence, on $`\mathrm{\Delta }`$, $`T_{ab}\mathrm{}^b`$ is a future directed, causal vector field. Now, using the Raychaudhuri equation and field equations at $`\mathrm{\Delta }`$ (condition (ii) of the Definition), we conclude $`T_{ab}\mathrm{}^a\mathrm{}^b\widehat{}=0`$. By expanding out this expression (see Eq (5)) we obtain
$$𝐅_{ab}^i\widehat{}=\varphi _1^i2(\mathrm{}_{[a}n_{b]}m_{[a}\overline{m}_{b]})+\varphi _2^i2(m_{[a}\mathrm{}_{b]})+\mathrm{CC},$$
(13)
for some complex algebra-valued functions $`\varphi _1^i`$ and $`\varphi _2^i`$ (the only non-vanishing Newman-Penrose components of $`𝐅_{ab}^i`$) on $`\mathrm{\Delta }`$, where $`\mathrm{CC}`$ stands for ‘the complex conjugate term’. These equations say that there is no flux of Yang-Mills radiation across $`\mathrm{\Delta }`$. Finally, condition (iv) in the Definition implies
$$|\mathrm{Re}\varphi _1|\widehat{}=\frac{2\pi }{a_\mathrm{\Delta }}Q_\mathrm{\Delta },|\mathrm{Im}\varphi _1|\widehat{}=\frac{2\pi }{a_\mathrm{\Delta }}P_\mathrm{\Delta },$$
(14)
where $`Q_\mathrm{\Delta }`$ is the electric charge and $`P_\mathrm{\Delta }`$ the magnetic charge at the horizon as defined by (8). Thus the boundary conditions severely restrict the form of matter fields at $`\mathrm{\Delta }`$. The component $`\varphi _0^i=\mathrm{}^am^b𝐅_{ab}^i`$ of the YM field vanishes and the gauge invariant components of $`\varphi _1^i`$ are completely determined by the electric and magnetic charges. However, the component $`\varphi _2^i`$ of the YM field is unconstrained.
Restrictions imposed on space-time curvature at $`\mathrm{\Delta }`$ are essentially the same as in Ref .<sup>§</sup><sup>§</sup>§This is because these restrictions were obtained assuming rather general conditions on the matter stress-energy which are satisfied in EYM. The derivation of some of these results involve long calculations and a topological result on the Chern-class of the $`SO(2)`$ connection associated with the dyad $`(m,\overline{m})`$. See . Results relevant to this paper can be summarized as follows. In the Newman-Penrose notation, for the Ricci tensor components, we have:
$`\mathrm{\Phi }_{00}`$ $`=`$ $`{\displaystyle \frac{1}{2}}R_{ab}\mathrm{}^a\mathrm{}^b\widehat{}=0,\mathrm{\Phi }_{01}={\displaystyle \frac{1}{2}}R_{ab}\mathrm{}^am^b\widehat{}=0,`$ (15)
$`\mathrm{\Phi }_{11}`$ $`=`$ $`{\displaystyle \frac{1}{4}}R_{ab}(\mathrm{}^an^b+m^a\overline{m}^b)\widehat{}=8\pi ^2{\displaystyle \frac{(Q_\mathrm{\Delta }^2+P_\mathrm{\Delta }^2)}{a_\mathrm{\Delta }^2}},R\widehat{}=0,`$ (16)
where $`R`$ is the scalar curvature. The Weyl tensor components satisfy
$`\mathrm{\Psi }_0`$ $`=`$ $`C_{abcd}\mathrm{}^am^b\mathrm{}^cm^d\widehat{}=0,\mathrm{\Psi }_1=C_{abcd}\mathrm{}^am^b\mathrm{}^cn^d\widehat{}=0,`$ (17)
$`\mathrm{\Psi }_2`$ $`=`$ $`C_{abcd}\mathrm{}^am^b\overline{m}^cn^d\widehat{}=\mathrm{\Phi }_{11}{\displaystyle \frac{2\pi }{a_\mathrm{\Delta }}}.`$ (18)
Furthermore,
$$\mathrm{\Psi }_3\widehat{}=\mathrm{\Phi }_{21},\text{that is}C_{abcd}\mathrm{}^an^b\overline{m}^cn^d\widehat{}=\frac{1}{2}R_{ab}\overline{m}^an^b.$$
(19)
As expected, for the SSS solutions discussed in Section II, these conditions are satisfied. We note that even when the curvature components $`\mathrm{\Psi }_2`$ and $`\mathrm{\Phi }_{11}`$ are the only ones different from zero in SSS solutions, for a general isolated horizon other curvature components (like $`\mathrm{\Psi }_3`$ and $`\mathrm{\Psi }_4`$) may be ‘dynamical’, i.e., vary along the integral curves of $`\mathrm{}`$.
In the Einstein-Maxwell-Dilaton system , we have a set of parameters $`(a_\mathrm{\Delta },Q_\mathrm{\Delta },P_\mathrm{\Delta },\varphi _\mathrm{\Delta })`$ at the horizon (and the same number at infinity), and those parameters were naturally selected as the horizon parameters. In the EYM case we have seen that we can define electric and magnetic charges at the horizon $`(Q_\mathrm{\Delta },P_\mathrm{\Delta })`$. The boundary conditions ensure that this quantities are constant along the horizon $`\mathrm{\Delta }`$ (and explicitly gauge invariant). Thus, it is natural to use the triplet $`a_\mathrm{\Delta },Q_\mathrm{\Delta },P_\mathrm{\Delta }`$ to parameterize general EYM isolated horizons. Let us denote by $``$ the space of horizon parameters with coordinates $`(a_\mathrm{\Delta },Q_\mathrm{\Delta },P_\mathrm{\Delta })`$, with the following restrictions: $`a_\mathrm{\Delta }>4\pi (Q_\mathrm{\Delta }^2+P_\mathrm{\Delta }^2)`$ and $`Q_\mathrm{\Delta },P_\mathrm{\Delta }[0,\mathrm{})`$.
## IV Action and Phase Space
The gravitational action has been shown to be differentiable, with respect to variations respecting the isolated horizons boundary conditions, in . We refer the reader to those papers for details. One important property of the variational principle is that one is varying, in the pure gravitational case, histories with a fixed value $`a_\mathrm{\Delta }^o`$ of the horizon area. We also need to ensure that the matter action is differentiable and work out the Hamiltonian framework for the matter sector as well. This could require imposition of additional boundary conditions on matter fields, but as we will see, the conditions already imposed on matter, described in Sec. III will be enough. The purpose of this section is concentrate on the variational principle for the Yang-Mills field and work out the Hamiltonian description.
Since we require that field equations hold on $`\mathrm{\Delta }`$, the gravitational boundary conditions already imply certain restrictions on the behavior of YM fields there. As noted in Section III, boundary conditions imply that several components of the Ricci tensor vanish on $`\mathrm{\Delta }`$, and that the curvature tensor $`𝐅`$ has the form (13). In particular $`|\mathrm{Re}(\varphi _1)|`$ and $`|\mathrm{Im}(\varphi _1)|`$ are spherically symmetric on the preferred cross-sections. In that case, $`|\varphi _1|`$ can be expressed in terms of the electric and magnetic charges, $`Q_\mathrm{\Delta }`$ and $`P_\mathrm{\Delta }`$ of the isolated horizon,
$`Q:\widehat{}={\displaystyle \frac{1}{4\pi }}{\displaystyle _{S_v}}|{}_{}{}^{}𝐅|\widehat{}={\displaystyle \frac{1}{2\pi }}{\displaystyle _S}|\mathrm{Re}\varphi _1|{}_{}{}^{2}ϵ,`$ (20)
$`P:\widehat{}={\displaystyle \frac{1}{4\pi }}{\displaystyle _{S_v}}|𝐅|\widehat{}={\displaystyle \frac{1}{2\pi }}{\displaystyle _{S_v}}|\mathrm{Im}\varphi _1|{}_{}{}^{2}ϵ,`$ (21)
(Here $`S_v`$ are the 2-spheres $`v=\mathrm{const}`$ in the preferred foliation. The minus sign in front of the first integrals in (20) and (21) arise because we have oriented $`S_v`$ such that the radial normal is in-going rather than outgoing.)
Let’s now discuss the action principle. For the same reasons that the area $`a_\mathrm{\Delta }`$ is kept fixed in the variational principle, we will now restrict ourselves to histories for which the values of electric and magnetic charges on the horizon are fixed to $`Q_\mathrm{\Delta }^o`$ and $`P_\mathrm{\Delta }^o`$ respectively. To make the action principle well-defined, we need to impose suitable boundary conditions on the YM fields. Conditions at infinity are the standard ones given in . To find boundary conditions on $`\mathrm{\Delta }`$, let us consider the standard YM bulk action:
$$S_{\mathrm{YM}}=\frac{1}{16\pi }_𝐌\sqrt{g}𝐅_{ab}^i𝐅_i^{ab}\mathrm{d}^4x.$$
(22)
Variation of $`S_{\mathrm{YM}}`$ yields
$$\delta \left(S_{\mathrm{YM}}\right)=\frac{1}{4\pi }_𝐌(D_a𝐅^{abi})\delta 𝐀_i\sqrt{g}\mathrm{d}^4x+\frac{1}{4\pi }_𝐌\delta 𝐀_{[a}^i{}_{}{}^{}𝐅_{bc]}^{i}\stackrel{~}{\eta }^{abc}\mathrm{d}^3x.$$
(23)
As usual, the bulk term provides the equations of motion provided the surface term vanishes. The boundary term (23) at infinity is the usual one and is dealt with in the standard manner . When evaluated at the horizon, the boundary term (23) does not automatically vanish. On $`\mathrm{\Delta }`$, the boundary term in (23) can be written as,
$$\frac{3}{4\pi }_\mathrm{\Delta }\delta 𝐀_a^i{}_{}{}^{}𝐅_{ibc}^{}\mathrm{}^{[c}\stackrel{~}{\eta }^{ab]}dv\mathrm{d}^2x.$$
(24)
Now, (13) implies that on $`\mathrm{\Delta }`$ the pull-back of $`𝐅_{ab}^i\mathrm{}^b`$ vanishes. Then, the horizon contribution reduces to
$$dv_{S_v}\delta (𝐀_{ia}\mathrm{}^a){}_{}{}^{}𝐅_{ab}^{i}\stackrel{~}{\eta }^{ab}\mathrm{d}^2x,$$
(25)
where, as before, $`v`$ is the affine parameter (with respect to the horizon metric) along the integral curves of $`\mathrm{}^a`$ such that $`v=\mathrm{const}`$ define the preferred foliation of $`\mathrm{\Delta }`$ and $`S_v`$ are the 2-spheres in this foliation. Now, since isolated horizons are to be thought of as “non-dynamical”, it is natural to ask that the gauge field $`𝐀_a^i`$ be invariant under the action of $`\mathrm{}`$. This is equivalent to ask that $`_{\mathrm{}}𝐀_a^i\widehat{}=D_a𝒱_{\mathrm{}}^i`$ be satisfied , for $`𝒱_{\mathrm{}}^i`$ the gauge generator. This condition is naturally satisfied since the pull-back of $`𝐅_{ab}^i\mathrm{}^b`$ to the horizon $`\mathrm{\Delta }`$ vanishes. Furthermore, the form of the boundary term (25) suggests that we fix gauge so that $`(𝐀^i\mathrm{}):=A_a^i\mathrm{}^a`$ is proportional to $`\mathrm{Re}(\varphi _1^i)`$, that is, $`(𝐀^i\mathrm{})=c\mathrm{Re}(\varphi _1^i)`$ for $`c`$ a constant on our space of histories This gauge condition was independently found by Ashtekar, Fairhurst and Krishnan in the more general context of distorted horizons. We thank S. Fairhurst for communicating their results prior to publication.. The norm of $`(𝐀^i\mathrm{})`$ as we shall see in following sections, is determined by the consistency of the Hamiltonian formalism (and will take its standard value in the static solutions, when they exist). Then the boundary term (25) arising in the variation of the action vanishes, i.e., the bulk action itself is differentiable and the action principle is well-defined. Note that the permissible gauge transformations are now restricted: If $`𝐀_a𝐀_a+D_af`$, the generating ‘function’ $`f^i`$ has to satisfy $`l^aD_af^i\widehat{}=0`$ on $`\mathrm{\Delta }`$ (as well as satisfy standard fall-off conditions at infinity).
Let us now provide a summary of the structure of the phase space of Yang-Mills fields. Fix a foliation of $`𝐌`$ by a family of space-like 3-surfaces $`M_t`$ (level surfaces of a time function $`t`$) which intersect $`\mathrm{\Delta }`$ in the preferred 2-spheres. Fix a normalization for the vector field $`\mathrm{}^a`$. Denote by $`t^a`$ a “time-evolution” vector field which is not tangential to the foliation with affine parameter $`t`$ which tends to a unit time-translation at infinity and to the vector field $`l^a`$ on $`\mathrm{\Delta }`$. The canonical conjugate moment is given by,
$$\stackrel{~}{𝚷}_i^a=\frac{\sqrt{h}}{4\pi }h^{ac}𝐧^b𝐅_{ibc},$$
(26)
where $`h_{ab}`$ is the intrinsic metric on $`M_t`$ induced from $`g_{ab}`$, and $`𝐧^a`$ is the normal to $`M_t`$. In terms of the lapse and shift fields $`N`$ and $`N^a`$ defined by $`t^a`$, the Legendre transform of the action yields:
$`S_{\mathrm{EM}}`$ $`=`$ $`{\displaystyle \frac{1}{4\pi }}{\displaystyle }\mathrm{d}t{\displaystyle _{M_t}}\mathrm{d}^3x(\stackrel{~}{𝚷}_i^a_a(𝐀^i\mathrm{})N^dF_{ad}^i\stackrel{~}{𝚷}^{ia}+`$ (28)
$`{\displaystyle \frac{N}{2\sqrt{h}}}h_{ab}\stackrel{~}{𝚷}^{ia}\stackrel{~}{𝚷}_i^b+{\displaystyle \frac{N}{2}}\sqrt{h}F_i^{ab}F_{ab}^i),`$
where the 2-form $`𝐄_{iab}`$ is the pull-back to $`M`$ of $`{}_{}{}^{}𝐅_{iab}^{}`$, $`{}_{}{}^{()}𝐄_{a}^{i}:=\frac{1}{2}ϵ_{a}^{}{}_{}{}^{bc}𝐄_{bc}^i`$, and $`{}_{}{}^{()}𝐅_{a}^{i}:=\frac{1}{2}ϵ_{a}^{}{}_{}{}^{bc}𝐅_{bc}^i`$, with $`ϵ_{abc}`$ the volume form defined by $`h_{ab}`$. The two-form $`𝐄_{ab}^i`$ is related to the canonically conjugate moment as follows: $`𝐄_{ab}^i=\frac{1}{4\pi }\stackrel{~}{}\eta {}_{abc}{}^{}\stackrel{~}{𝚷}_{}^{ic}`$.
Thus, as usual, the phase space consists of pairs $`(𝐀_a^i,𝐄_{abi})`$ on the 3-manifold $`M_t`$, subject to boundary conditions, where the connection $`𝐀_a^i`$ is now the pull-back to $`M_t`$ of the Yang-Mills 4-potential and the 2-form $`𝐄_{ab}^i`$ is the dual of the electric field vector density. These fields are subject to boundary conditions. On any horizon 2-sphere $`S_v`$, conditions (iv) must hold ensuring that the pull-backs of $`|𝐅|_{ab}`$ and $`|𝐄|_{ab}`$ are spherically symmetricat the horizon $`\mathrm{\Delta }`$, the natural parameter $`v`$ and the ‘time’ parameter $`t`$ coincide.. (Since $`(𝐀^i\mathrm{})`$ appears as a Legendre multiplier in (28) it is not part of the phase space variables at $`\mathrm{\Delta }`$.) At infinity, $`𝐀_a^i,𝐄_{iab}`$ are subject to the appropriate boundary conditions that ensure asymptotic flatness of the metric and are general enough to include the nontrivial EYM solutions mentioned in section III. The conditions for the YM sector are :
$$A_a^i=A_a^i(\theta ,\varphi )/r+o(r^1),F_{ab}^i=F_{ab}^i(\theta ,\varphi )/r+o(r^2).$$
(29)
The symplectic structure on this YM-sector of phase space can be read off from (28):
$$\mathrm{\Omega }|_{(𝐀,𝐄)}(\delta _1,\delta _2)=\frac{1}{4\pi }_M\left[\delta _1\stackrel{~}{𝐄}_i^a\delta _2𝐀_a^i\delta _2\stackrel{~}{𝐄}_i^a\delta _1𝐀_a^i\right].$$
(30)
(The asymptotic conditions ensure that the integrals converge.) As usual, there is one first class constraint, $`D_{[a}𝐄_{bc]}^i=0`$, which generates gauge transformations: Under the canonical transformation generated by $`f_iD_{[a}𝐄_{bc]}^i\stackrel{~}{\eta }^{abc}\mathrm{d}^3x`$, the canonical fields transform, as usual, via $`𝐀_a^i𝐀_a^i+D_af^i`$ and $`𝐄_{iab}`$ ‘rotates’. Note that our boundary conditions allow the generating function $`f_i`$ to be non-trivial on the (intersection of $`M_t`$ with) $`\mathrm{\Delta }`$; the smeared constraint function is still differentiable. Thus, as in the gravitational case, the gauge degrees of freedom do not become physical in this framework.
Let us summarize. By imposing a set of boundary conditions on the horizon, to be satisfied by all histories in the variational principle, we arrived at the phase space of EYM isolated horizons $`_0`$. This phase space can be seen as the (gauge equivalence class of) solutions to the equations of motion with constant and fixed quasi-local parameters $`a_\mathrm{\Delta }^o`$, $`Q_\mathrm{\Delta }^o`$ and $`P_\mathrm{\Delta }^o`$. For the purpose of the variational principle (in the sense that is well defined and yields the correct equation of motion), it is enough to consider configurations with fixed values of the horizon parameters. However, as we shall see in following sections, when considering the Hamiltonian formulation –essential for the formulation of the first law– one needs to extend the space of isolated horizons from $`_0`$ to $``$, where all possible values of the quasi-local parameters are considered .
## V Surface Gravity and the Zeroth Law
In each SSS solution there is a unique time-translational Killing field $`t^a`$ which is unit at infinity. As usual, surface gravity $`\kappa _{\mathrm{SSS}}`$ is defined in terms of its acceleration at the horizon: $`t^a_at^b\widehat{}=\kappa _{\mathrm{SSS}}t^b`$. Unlike the Einstein-Maxwell-Dilaton theory where the knowledge of the exact solutions allow us to write $`\kappa `$ in terms of the parameters of the solutions, for the YM field we do not have a closed form for $`\kappa `$. The only expression we have at our disposal is the general formula found by Visser for SSS space-times of the form (9), independently of the matter content of the theory, given by the expression ,
$$\kappa _{\mathrm{SSS}}=\frac{1}{2r_\mathrm{h}}e^{\delta (r_\mathrm{h})}\left[12m^{}(r_\mathrm{h})\right].$$
(31)
From the perspective of the isolated horizon framework, $`\kappa `$ is the acceleration of the properly normalized null normal $`\mathrm{}^a`$ to $`\mathrm{\Delta }`$ . In the SSS solutions, $`\mathrm{\Delta }`$ happens to be a Killing horizon and we can select a unique vector field $`\mathrm{}^a`$ from the the equivalence class $`[\mathrm{}^a]`$ simply by setting $`\mathrm{}^a\widehat{}=t^a`$. Then $`\kappa _{\mathrm{SSS}}`$ is the acceleration of this specific $`\mathrm{}^a`$. In the case of general isolated horizons, the challenge is to find a prescription to single out a preferred $`\mathrm{}^a`$, without reference to any Killing field. The strategy we would like to adopt has two steps and is motivated by the one used for isolated horizons in dilaton gravity . However, as we shall see below, we will encounter some difficulties which make the EYM case more subtle than the systems previously studied.
In the first step, we will normalize $`\mathrm{}^a`$ only up to a constant, leaving a rescaling freedom $`\mathrm{}^a\mathrm{}^a=c\mathrm{}^a`$, where $`c`$ is a constant on $`\mathrm{\Delta }`$ but may depend on the parameters $`r_\mathrm{\Delta },Q_\mathrm{\Delta },P_\mathrm{\Delta }`$ of the isolated horizon. For each such $`\mathrm{}^a`$, we can define the surface gravity $`\kappa _{\mathrm{}}`$ relative to that $`\mathrm{}^a`$ via $`\mathrm{}^a_a\mathrm{}^b\widehat{}=\kappa _{\mathrm{}}\mathrm{}^b`$. Rescaling of $`\mathrm{}^a`$ now induces to a ‘gauge transformation’ in $`\kappa `$: $`\kappa _{\mathrm{}}\kappa _{\mathrm{}^{}}=c\kappa _{\mathrm{}}`$. (Recall that in the general Newman-Penrose framework, $`\kappa `$ is a connection component and therefore undergoes the standard gauge transformations under a change of the null tetrad. By fixing $`\mathrm{}^a`$ up to a constant rescaling, we have reduced the general gauge freedom to that of a constant rescaling.) Since the zeroth law only says that the surface gravity is constant on $`\mathrm{\Delta }`$, if it holds for one $`\mathrm{}^a`$, it holds for every $`\mathrm{}^a=c\mathrm{}^a`$. Thus, for the zeroth law, it is in fact not essential to get rid of the rescaling freedom.
Recall that the isolated horizon is naturally equipped with equivalence classes $`[\mathrm{},n]`$ of vector and co-vector fields, subject to the relation: $`(\mathrm{},n)(G^1\mathrm{},Gn)`$ for any positive function $`GG(v)`$ on $`\mathrm{\Delta }`$. As we already mentioned, our first task is to reduce the freedom in the choice of $`G(v)`$ to that of a constant. We use the same strategy as in . (For motivation, see .) Recall that $`\mu `$, the expansion of $`n`$ is strictly negative and constant on each leaf of the preferred foliation; $`\mu \mu (v)<0`$. It is easy to verify that
$$n^aG(v)n^a\mathrm{implies}\mu G(v)\mu (v).$$
(32)
Hence, we can always use the $`G(v)`$ freedom to set $`\mu \widehat{}=\mathrm{const}`$. This condition restricts the family of $`(\mathrm{}^a,n_a)`$ pairs and reduces the equivalence relation to $`(\mathrm{}^a,n_a)(c\mathrm{}^a,c^1n_a)`$ where $`c`$ is any constant on $`\mathrm{\Delta }`$. We will denote the restricted equivalence class by $`[\mathrm{}^a,n_a]_R`$. In the second step, we can arbitrarily fix the numerical value of $`\mu `$ in terms of the parameters of the isolated horizon and eliminate the rescaling freedom altogether, thereby selecting, for each choice of $`\mu `$, a canonical pair $`(\mathrm{}^a,n_a)`$ on each isolated horizon.
With the equivalence class $`[\mathrm{}^a,n_a]_R`$ at our disposal, as discussed above, we can define a surface gravity $`\kappa _{\mathrm{}}`$ via $`\mathrm{}^a_a\mathrm{}^b\widehat{}=\kappa _{\mathrm{}}\mathrm{}^b`$. Constancy of $`\kappa _{\mathrm{}}`$ on $`\mathrm{\Delta }`$ follows from the same arguments that were used in . For completeness, let us briefly recall the structure of that proof. First, using conditions on derivatives of $`l^a,n_a`$ introduced in the Definition, one can express the self-dual part of the Riemann curvature in terms of $`\kappa _{\mathrm{}},_a\kappa _{\mathrm{}},\mu `$ (and another field which is not relevant to this discussion). Comparing this expression to the standard Newman-Penrose expansion of the self-dual curvature tensor in terms of curvature scalars , and using the fact that certain curvature scalars vanish on $`\mathrm{\Delta }`$ (see (15) and (17)), one can conclude
$$(_{[a}\kappa _{\mathrm{}})n_{b]}\widehat{}=0,\mathrm{and}\kappa _{\mathrm{}}\widehat{}=\frac{\mathrm{\Psi }_2}{\mu }.$$
(33)
The first equation implies that $`\kappa _{\mathrm{}}`$ is spherically symmetric. Hence, it only remains to show that $`_{\mathrm{}}\kappa _{\mathrm{}}\widehat{}=0`$. Since $`\mu `$ is now a constant on $`\mathrm{\Delta }`$, it suffices to show that $`_{\mathrm{}}\mathrm{\Psi }_2=0`$. Now, the (second) Bianchi identity implies that
$$_{\mathrm{}}\left(\mathrm{\Psi }_2\mathrm{\Phi }_{11}\right)\widehat{}=0.$$
(34)
Finally, using (13), we conclude: $`\mathrm{\Phi }_{11}=8\pi ^2\frac{(Q_\mathrm{\Delta }^2+P_\mathrm{\Delta }^2)}{a_\mathrm{\Delta }^2}`$. Thus, $`\mathrm{\Phi }_{11}`$ is constants on $`\mathrm{\Delta }`$. Combining these results, we conclude $`_{\mathrm{}}\kappa _{\mathrm{}}\widehat{}=0`$, whence $`\kappa _{\mathrm{}}`$ is constant on $`\mathrm{\Delta }`$. This establishes the zeroth law.
Let us now consider the second step in fixing the normalization of $`\mathrm{}`$. So far, we have only required that $`\mu `$ be a (negative) constant but not fixed its value. Under the rescaling $`\mu c^1\mu `$ we have: $`\mathrm{}c\mathrm{}`$, and $`\kappa _{\mathrm{}}c\kappa _{\mathrm{}}`$. Hence, the remaining rescaling freedom in $`\mathrm{}`$ and $`\kappa `$ can be exhausted simply by fixing the value of $`\mu `$ in terms of the isolated horizon parameters. We would like to single out a canonical choice, and the obvious strategy is to fix $`\mu `$ to the value $`\mu _{\mathrm{SSS}}`$ it takes on the SSS solutions. However, there are two difficulties: First, although $`\mu _{\mathrm{SSS}}`$ is a well-defined function of the isolated horizon parameters $`r_\mathrm{\Delta },Q_\mathrm{\Delta }`$ and $`P_\mathrm{\Delta }`$, where static solutions exist, there is no closed expression for it in terms of these parameters simply because the full set of solutions (including the colored black holes) is not known in closed form. The second and more serious problem is that, for an arbitrary point in parameter space, there might not be any SSS solution (recall that the colored black holes span only a countable number of points in the $`P_\mathrm{\Delta }`$ axis for a given value of $`r_\mathrm{\Delta }`$). As we shall see in next section, one can still have a consistent Hamiltonian formulation and a first law even for those isolated horizons lying on points of the parameter space where no SSS solutions exist. But, as should be clear from the discussion, at this point there is no canonical normalization of $`\mu `$ on the whole space $``$.
Even when we can not find an explicit functional form for $`\mu `$, we can still write the general form that $`\kappa _{\mathrm{}}`$ shall have. Using an identity coming from the isolated horizons boundary condition we have the following relation ,
$$\kappa _{\mathrm{}}\widehat{}=\frac{1}{\mu }\left(\frac{2\pi }{a_\mathrm{\Delta }}+\mathrm{\Phi }_{11}\right).$$
(35)
Using the expression for $`\mathrm{\Phi }_{11}`$ in terms of the charges (15) we have,
$$\kappa _{\mathrm{}}=\frac{1}{\mu }\frac{1}{2r_\mathrm{\Delta }^2}\left[1\frac{(Q_\mathrm{\Delta }^2+P_\mathrm{\Delta }^2)}{r_\mathrm{\Delta }^2}\right].$$
(36)
For those points of parameter space $``$ where an SSS static solution exists, we can go further and make use of the general form of the metric (9). One then finds that the expansion of the properly normalized $`n`$ is such that
$$\mu _{\mathrm{SSS}}=\frac{e^{\delta (r_\mathrm{\Delta })\delta (\mathrm{})}}{r_\mathrm{\Delta }}.$$
(37)
A coordinate transformation can always render $`\delta (\mathrm{})=0`$, so the properly normalized $`\kappa `$ for general isolated horizons takes the form,
$$\kappa =\frac{e^{\delta (r_\mathrm{\Delta },Q_\mathrm{\Delta },P_\mathrm{\Delta })}}{2r_\mathrm{\Delta }}\left[1\frac{(Q_\mathrm{\Delta }^2+P_\mathrm{\Delta }^2)}{r_\mathrm{\Delta }^2}\right].$$
(38)
It is in a sense remarkable that our ignorance about the explicit form of the solutions is encoded in the function $`\delta `$.
We end this section with a discussion. As we have emphasized in this section, any choice of $`\mu `$ in terms of the horizon parameters defines a vector field $`\mathrm{}^a`$ and a surface gravity on $``$. However, one would like to make contact with the space of static solutions in such a way that the surface gravity $`\kappa _{\mathrm{}}`$ coincides with the surface gravity of the properly normalized Killing field. In the case of Einstein-Maxwell-Dilaton this was indeed possible since the static solutions span the space of horizon parameters $``$. That is, for each value of the isolated horizon parameters, there exists a (unique) static solution which allows us to fix $`\mu `$ in a unique, canonical way. In the case of EYM, the space of SSS solutions does not span $``$. Thus, one is able to canonically fix $`\kappa `$ only in this subspace; for a general point, there is no preferred normalization. This seems to be a serious shortcoming of the formalism that might place the EYM system on a different status than the EM and EMD systems. As we shall see in next section, this ambiguity is also manifested in the definition of mass. However, this problem will also motivate a second conjecture that we shall put forward in Section VII.
## VI Mass and the First Law
The space-times that admit an isolated horizon are not necessarily stationary, therefore, it is no longer meaningful to identify the ADM mass $`M`$ with the mass $`M_\mathrm{\Delta }`$ of the isolated horizon. For the formulation of the first law, we must first introduce an appropriate definition of $`M_\mathrm{\Delta }`$. This definition should be general enough in the sense that, for each choice for the normalization of $`\mathrm{}^a`$, the Mass should be uniquely defined, and the first law should be valid for any choice of normalization. As we shall see, the Hamiltonian framework provides a natural strategy. In the Einstein-Maxwell-Dilaton case the total Hamiltonian consists of a bulk term and two surface terms, one at infinity and the other at the isolated horizon. As usual, the bulk term is a linear combination of constraints and the surface term at infinity yields the ADM energy. In a rest-frame adapted to the horizon it is then natural to identify the surface term at $`\mathrm{\Delta }`$ as the horizon mass, $`M_\mathrm{\Delta }`$. Indeed, there are several considerations that support this identification .
For the gravitational part of the action and Hamiltonian, the discussion of only assumed that the stress-energy tensor satisfies two conditions at $`\mathrm{\Delta }`$: i) $`T_{ab}\mathrm{}^a`$ is a future pointing causal vector field on $`\mathrm{\Delta }`$; and, ii) $`T_{ab}\mathrm{}^an^b`$ is spherically symmetric on $`\mathrm{\Delta }`$. Both these conditions are met in the present case. Therefore, we can take over the results of directly. For the matter part of the action and Hamiltonian, the overall situation is again analogous, although there are the obvious differences in the detailed expressions. As in the Einstein-Maxwell case, matter terms contribute to the surface terms in the Hamiltonian only because one has to perform one integration by parts to obtain the Gauss constraint in the bulk term.
The net result is the following. Consider a foliation of the given space-time region $`𝐌`$ by a 1-parameter family of (partial) Cauchy surfaces $`M_t`$, each of which extends from the isolated horizon $`\mathrm{\Delta }`$ to spatial infinity $`i^o`$ (see Figure). We will assume that $`M_t`$ intersects $`\mathrm{\Delta }`$ in a 2-sphere belonging to our preferred foliation and that the initial data induced on $`M_t`$ are asymptotically flat. Denote by $`S_\mathrm{\Delta }`$ and $`S_{\mathrm{}}`$ the 2-sphere boundary of $`M_t`$ at the horizon and infinity, respectively. Choose a time-like vector field $`t^a`$ in $`𝐌`$ which tends to the unit time-translation orthogonal to the foliation at spatial infinity and to the vector field $`\mathrm{}^a`$ on $`\mathrm{\Delta }`$, with the normalization fixed (or partially fixed) as in Section V. Then, the Hamiltonian $`H_t`$ generating evolution along $`t^a`$ is given by:
$`H_t`$ $`=`$ $`{\displaystyle _{M_t}}\mathrm{constraints}+\underset{r_o\mathrm{}}{lim}{\displaystyle _{S_{r_o}}}\left({\displaystyle \frac{r_o}{4\pi G}}\mathrm{\Psi }_2\right){}_{}{}^{2}ϵ\mathrm{\Phi }_{\mathrm{}}Q_{\mathrm{}}`$ (39)
$`+`$ $`{\displaystyle _{S_\mathrm{\Delta }}}\left({\displaystyle \frac{\mu ^1}{4\pi G}}\mathrm{\Psi }_2\right){}_{}{}^{2}ϵ+|(𝐀l)|Q_\mathrm{\Delta }+V,`$ (40)
where $`S_{r_o}`$ are large 2-spheres of radius $`r_o`$ and $`V`$ is a constant on $`_0`$, the space of isolated horizons with fixed values of the horizon parameters. (The calculation and the final result are completely analogous to those in the Einstein-Maxwell case .) Note that the surface terms depend only on the ‘Coulombic’ parts of the gravitational and Yang-Mills fields.
It is easy to check that the surface term at infinity is, as usual, the time component $`P_a^{\mathrm{ADM}}t^a`$ of the ADM 4-momentum $`P_a^{\mathrm{ADM}}`$, which in the present (-,+,+,+) signature is negative of the ADM energy, $`P_a^{\mathrm{ADM}}t^a=E^{\mathrm{ADM}}`$. It is natural to identify the surface term at $`S_\mathrm{\Delta }`$ as the energy of the isolated horizon. (There is no minus sign because $`S_\mathrm{\Delta }`$ is the inner boundary of $`M`$). Since $`t^a\widehat{}=\mathrm{}^a`$ and since $`\mathrm{}^a`$ represents the ‘rest frame’ of the isolated horizon, this energy can in turn be identified with the horizon mass $`M_\mathrm{\Delta }`$. Thus, we have:
$$M_\mathrm{\Delta }^{(\mathrm{})}=_{S_\mathrm{\Delta }}\left(\frac{\mu ^1}{4\pi G}\mathrm{\Psi }_2\right){}_{}{}^{2}ϵ+|(𝐀l)|Q_\mathrm{\Delta }+V_{(\mathrm{})}(a_\mathrm{\Delta },Q_\mathrm{\Delta },P_\mathrm{\Delta }).$$
(41)
Here, the so far undetermined function $`V`$, depends only on the horizon parameters (and coupling constants). In the variational principle, this term played no role, but in the Hamiltonian description it becomes essential, since we are now interested in variations along the full isolated horizons phase space $``$. Thus, one should be able to consider in the formalism displacements along directions in which the horizon parameters change. As we shall show below, requiring a consistent Hamiltonian formulation enables us to determine the function $`V`$ for the EYM system. Now, using the expression (33) of surface gravity in terms of the Weyl tensor (and $`\mu `$), and calling $`\mathrm{\Phi }_\mathrm{\Delta }:=|𝐀\mathrm{}|`$ on $`\mathrm{\Delta }`$, we can cast $`M_\mathrm{\Delta }`$ in a more familiar form:
$$M_\mathrm{\Delta }=\frac{1}{4\pi }\kappa a_\mathrm{\Delta }+\mathrm{\Phi }_\mathrm{\Delta }Q_\mathrm{\Delta }+V(a_\mathrm{\Delta },Q_\mathrm{\Delta },P_\mathrm{\Delta }),$$
(42)
where we have dropped the explicit $`\mathrm{}^a`$ dependence of the Mass for notational simplicity. Thus, as in the Einstein-Maxwell case, we obtain a Smarr formula. However, the meaning of various symbols in the equation is somewhat different. Since an isolated horizon need not be a Killing horizon, in general $`M_\mathrm{\Delta }`$ does not equal the ADM mass, nor is $`\kappa `$ or $`\mathrm{\Phi }_\mathrm{\Delta }`$ computed using a Killing field. Since the constraints are satisfied in any solution, the bulk term in (39) vanishes as well. Hence, in this case, $`H_t=M_\mathrm{\Delta }E^{\mathrm{ADM}}`$, the difference being the ‘radiative energy’ in the space-time. Finally, as emphasized in , the matter contribution to the mass formula (41) is subtle: while it does not include the energy in radiation outside the horizon, it does include the energy in the ‘Coulombic part’ of the field associated with the black hole hair. (Recall that the future limit of the Bondi energy has this property.) This fact is crucial to the analysis of the ‘physical process version’ of the first law. However, since this issue was discussed in detail in , we shall not discuss this here.
Now, a consistent Hamiltonian formulation (for a sector of a diffeomorphism invariant theory) requires that for an arbitrary vector $`\delta `$ tangent to the symplectic manifold, (i.e. the phase space $`\mathrm{\Gamma }`$), one has
$$\delta H=\mathrm{\Omega }(\delta ,X_H),$$
(43)
where $`X_H`$ is the vector field that corresponds to the equations of motion for a given choice of lapse and shift, and $`H`$ is the Hamiltonian function corresponding to the same choice of lapse and shift. In the preceding prescription one assumes that the evaluation of $`\delta H`$ is carried out by considering the change in $`H`$ associated with the displacement $`\delta `$ of the phase space point, but keeping the lapse and shift fixed. If we now let the choice of lapse and shift depend on the phase space point –as is the case when $`\mathrm{}`$ is normalized as in Sec. V– we would obtain a new variation $`\delta \stackrel{~}{H}`$. This might fail to satisfy
$$\delta \stackrel{~}{H}=\mathrm{\Omega }(\delta ,X_H),$$
(44)
with the same $`X_H`$ as in (43). It turns out that the necessary and sufficient condition to obtain the required consistency is the validity of the first law.
$$\delta M_\mathrm{\Delta }=\frac{1}{8\pi }\kappa \delta a_\mathrm{\Delta }+\mathrm{\Phi }_\mathrm{\Delta }\delta Q_\mathrm{\Delta }.$$
(45)
That is, the first law of black hole mechanics -for quantities defined only at the horizon- arises naturally as part of the requirements for a consistent Hamiltonian formulation in which, for every value of the horizon parameters, one has chosen a canonical lapse and shift functions making the latter dependent on the point in phase space. Note that in contrast with the above situation, when constructing a Hamiltonian to deal with the analogous problem at infinity, the canonical choice of normalization of lapse and shift at infinity is taken as independent of the phase space point, namely, they are chosen to correspond to a unit time translation (normal to the initial-data hyper-surface) at infinity.
In the cases of Einstein Vacuum, Einstein-Maxwell and Einstein-Maxwell-Dilaton theories, this consistency requirement translate into an identity that is automatically satisfied by the expressions of $`M_\mathrm{\Delta }`$, $`\kappa `$ etc, for all values of the parameters and variations there-off. In the case of Einstein-Yang-Mills Theory, -as well as in other theories where hair is present-, the only way to ensure the validity of the consistency requirement is to limit the class of variations $`\delta `$ allowed. This has two dramatic consequences: First, it defines a foliation of phase space by a collection of (symplectic) leaves over which the Hamiltonian formulation is consistent, a situation which puts the construction in the ‘non-standard’ class. In fact, some simple systems are described by a similar type of situation, as for example the (reduced) Hamiltonian description of a rotating body where the -3 dimensional- phase space is foliated by two spheres, each of which is a true symplectic manifold where the Hamiltonian motion is restricted . Second, the status of the first law changes from that of an identity, valid for all variations, to that of a specification of the class of variations that the formalism allows.
It is important to note that in the first law (45) only variations of the electric charge are involved, and not variations of the magnetic charge. On the other hand, the Horizon Mass (42) might depend on $`P_\mathrm{\Delta }`$ through $`V`$.
Let us now see that asking consistency of the formalism leads us to some conditions that the function $`V`$ should satisfy. In order to do this we shall follow, for completeness, Ref. closely. The first step in this direction is to regard (45) as an identity between one forms $`\mathrm{d}M_\mathrm{\Delta }=\frac{1}{8\pi }\kappa \mathrm{d}a_\mathrm{\Delta }+\mathrm{\Phi }_\mathrm{\Delta }\mathrm{d}Q_\mathrm{\Delta }`$, where $`\kappa ,\mathrm{\Phi }_\mathrm{\Delta }`$ and $`V`$ are functions on $``$. Thus one can consider differential forms on $``$ and take an exterior derivative of the ‘first law’ to arrive at,
$$0=\frac{1}{8\pi }\mathrm{d}\kappa \mathrm{d}a_\mathrm{\Delta }+\mathrm{d}\mathrm{\Phi }_\mathrm{\Delta }\mathrm{d}Q_\mathrm{\Delta }.$$
(46)
The first conclusion coming from (46) is that the variations on $``$ are restricted to sub-manifolds such that the pull-back of the form $`\mathrm{d}P_\mathrm{\Delta }\mathrm{d}a_\mathrm{\Delta }`$ vanishes. That is, $`P_\mathrm{\Delta }`$ is not free to vary independently of $`r_\mathrm{\Delta }`$ and $`Q_\mathrm{\Delta }`$. This is precisely what happens for SSS solutions (representing only a discrete set of curves in the plane $`(r_\mathrm{\Delta },P_\mathrm{\Delta })`$) and in the static axial-symmetric case (also covering a discrete set of curves in the plane). From now on, we restrict ourselves to the symplectic leaves where the formalism is well defined. On these sub-manifolds the magnetic charge becomes a function of the area and electric charge, $`P_\mathrm{\Delta }=P_\mathrm{\Delta }(r_\mathrm{\Delta },Q_\mathrm{\Delta })`$.
Second, the condition (46) gives us a relation between $`\kappa `$ and $`\mathrm{\Phi }_\mathrm{\Delta }`$,
$$\frac{\kappa }{Q_\mathrm{\Delta }}=8\pi \frac{\mathrm{\Phi }_\mathrm{\Delta }}{a_\mathrm{\Delta }}.$$
(47)
Thus, if we know the surface gravity $`\kappa `$ then we can derive an expression for the potential $`\mathrm{\Phi }_\mathrm{\Delta }`$.
Finally, taking the variation of (42) and comparing it to (45) we arrive at the following equations,
$`a_\mathrm{\Delta }{\displaystyle \frac{\beta }{a_\mathrm{\Delta }}}+8\pi r_\mathrm{\Delta }Q_\mathrm{\Delta }{\displaystyle \frac{\mathrm{\Phi }_\mathrm{\Delta }}{a_\mathrm{\Delta }}}+8\pi r_\mathrm{\Delta }{\displaystyle \frac{V}{a_\mathrm{\Delta }}}`$ $`=`$ $`0,`$ (48)
$`{\displaystyle \frac{r_\mathrm{\Delta }}{2}}{\displaystyle \frac{\beta }{Q_\mathrm{\Delta }}}+Q_\mathrm{\Delta }{\displaystyle \frac{\mathrm{\Phi }_\mathrm{\Delta }}{Q_\mathrm{\Delta }}}+{\displaystyle \frac{V}{Q_\mathrm{\Delta }}}`$ $`=`$ $`0,`$ (49)
where $`a_\mathrm{\Delta }=4\pi r_\mathrm{\Delta }^2`$ and, for convenience, we have defined $`\beta :=2r_\mathrm{\Delta }\kappa `$. Thus, given $`\kappa `$ and $`\mathrm{\Phi }_\mathrm{\Delta }`$ one can in principle, integrate equations (48) and (49) to find $`V`$. Note that these equations are defined over the horizon parameters space $``$, so they are completely local to the horizon $`\mathrm{\Delta }`$.
Recall that the general prescription for arriving at an explicit expression for surface gravity $`\kappa `$, for general isolated horizons, involves the fixing of the expansion $`\mu `$ as function of the horizon parameters. For this, one requires some input from the SSS solutions (where they exist). However, it is important to stress that the results of this section regarding a consistent Hamiltonian formulation and the validity of the first law, are independent of the particular choice of normalization $`\mu `$ (and $`\kappa `$) that one makes. Thus, there is a consistent Hamiltonian for each choice. This is particularly important for those points of horizon parameter space where no SSS solutions exist and therefore, no ‘canonical’ normalization is available. Nevertheless, Isolated Horizons still exist, and are well defined for those points of parameter space. It is when we want to have a canonical choice of $`\mu `$ and therefore, of $`\kappa `$ and $`M_\mathrm{\Delta }`$, that we are forced make contact with static solutions (for the allowed regions in $``$).
In the remaining of this section, we focus our attention to static spherically symmetric (SSS) solutions to the EYM equations described in Section II. We shall consider the three classes of SSS solutions described in Section II, and find expressions for their surface gravity $`\kappa `$ and Horizon Masses $`M_\mathrm{\Delta }`$. As discussed before, a study of the three sectors of SSS solutions serves two purposes. First, it provides us with a way of fixing the normalization of $`\mathrm{}^a`$ for general isolated horizons for those point of parameter space $``$ where SSS solutions exist, and second, it will allow us to find, in the next Section, new results regarding SSS solutions. This is because, even when the expressions we will find for $`\kappa `$ and $`M_\mathrm{\Delta }`$ are valid in the general framework, they are, in particular, also valid for SSS solutions. Some of the results of this and the next section, regarding SSS colored black holes, were already reported in .
Let us start with the electrically charged case, corresponding to the $`P_\mathrm{\Delta }=0`$ surface in $``$. Since these solutions are nothing but electrically charged Reissner Nordstrom solutions, the expansion of $`n`$ is given by ,
$$\mu =\frac{1}{r_\mathrm{\Delta }},$$
(50)
and the surface gravity is given by
$$\kappa =\frac{1}{2r_\mathrm{\Delta }}\left(1\frac{Q_\mathrm{\Delta }^2}{r_\mathrm{\Delta }^2}\right).$$
(51)
We can now consider Equation (47) and find that $`_{r_\mathrm{\Delta }}\mathrm{\Phi }_\mathrm{\Delta }=\frac{Q_\mathrm{\Delta }}{r_\mathrm{\Delta }^2}`$. Then, asking $`\mathrm{\Phi }_\mathrm{\Delta }`$ to vanish as $`r_\mathrm{\Delta }\mathrm{}`$, we have that
$$\mathrm{\Phi }_\mathrm{\Delta }=\frac{Q_\mathrm{\Delta }}{r_\mathrm{\Delta }},$$
(52)
which corresponds precisely to the value of the electric potential on RN solutions. The equations (48) now imply that $`_{r_\mathrm{\Delta }}V=0`$ and $`_{Q_\mathrm{\Delta }}V=0`$. Since the restriction of $`V`$ to the plane $`P_\mathrm{\Delta }=0`$ does not depend on $`P_\mathrm{\Delta }`$, the only possibility is that $`V=\mathrm{constant}`$ on the $`P_\mathrm{\Delta }=0`$ plane of $``$. In order to fix the value of $`V`$ we notice that the Schwarzschild one-parameter family of solutions –corresponding to zero electric field– are contained within the electric RN family. These solutions are purely gravitational since the gauge potential vanishes exactly, and it is known that for the pure Einstein theory one can set $`V=0`$ (see ).
We now have expressions of the mass $`M_\mathrm{\Delta }`$, surface gravity $`\kappa `$, area $`a_\mathrm{\Delta }`$ and the electric potential $`\mathrm{\Phi }_\mathrm{\Delta }`$ of any isolated horizon in terms of its fundamental parameters $`r_\mathrm{\Delta },Q_\mathrm{\Delta },`$:
$$M_\mathrm{\Delta }=\frac{r_\mathrm{\Delta }}{2}\left[1+\frac{Q_\mathrm{\Delta }^2}{r_\mathrm{\Delta }^2}\right].$$
(53)
This is precisely the same form as in the Einstein-Maxwell theory. The total energy of the system $`E`$, related to on-shell value of the Hamiltonian, is given by
$$E=H_t=M_{\mathrm{ADM}}M_\mathrm{\Delta },$$
(54)
which for the electric RN embedded family vanishes exactly.
Let us now consider embedded Abelian magnetic solutions. For this solutions, the electric charge $`Q_\mathrm{\Delta }`$ vanishes so the Horizon mass variation formula, when restricted to the purely magnetic sector of the SSS space takes the form,
$$\delta M_\mathrm{\Delta }=\frac{1}{8\pi }\kappa \delta a_\mathrm{\Delta }.$$
(55)
The formula for the mass (42) is given by,
$$M_\mathrm{\Delta }=\frac{1}{4\pi }\kappa a_\mathrm{\Delta }+V(r_\mathrm{\Delta },Q_\mathrm{\Delta }=0,P_\mathrm{\Delta }=1),$$
(56)
The normalization factor $`\mu `$ remains the same as in the electrically charged case given by (50) and the surface gravity is given by
$$\kappa =\frac{1}{2r_\mathrm{\Delta }}\left(1\frac{P_\mathrm{\Delta }^2}{r_\mathrm{\Delta }^2}\right).$$
(57)
Then equation (47) implies that $`\mathrm{\Phi }_\mathrm{\Delta }`$ is zero. The second set of equations (48) and (49) reduce to,
$$_{Q_\mathrm{\Delta }}V=0,$$
(58)
and
$$_{r_\mathrm{\Delta }}V=\frac{r_\mathrm{\Delta }}{2}_{r_\mathrm{\Delta }}\beta .$$
(59)
Using equation (57) we get the following equation that $`V`$ should satisfy,
$$_{r_\mathrm{\Delta }}V=\frac{P_\mathrm{\Delta }^2}{r_\mathrm{\Delta }^2}.$$
(60)
Now, recall that, for a given value of the magnetic charge $`P_\mathrm{\Delta }`$, one has black hole solutions for $`r_\mathrm{\Delta }|P_\mathrm{\Delta }|`$ (the extreme case corresponding to $`r_\mathrm{\Delta }=|P_\mathrm{\Delta }|`$). In order to integrate (60), one has to choose some ‘boundary conditions’ on the space $``$. Our choice, motivated by consistency with the colored black holes (see below), is to set $`V(r_\mathrm{\Delta }=P_\mathrm{\Delta })=0`$. With this choice, $`V`$ takes the form,
$$V=_{|P_\mathrm{\Delta }|}^{r_\mathrm{\Delta }}\frac{P_\mathrm{\Delta }^2}{\stackrel{~}{r}^2}d\stackrel{~}{r}=\frac{P_\mathrm{\Delta }^2}{r_\mathrm{\Delta }}|P_\mathrm{\Delta }|.$$
(61)
With this, the Horizon Mass $`M_\mathrm{\Delta }`$ is given by,
$`M_\mathrm{\Delta }`$ $`=`$ $`{\displaystyle \frac{r_\mathrm{\Delta }}{2}}\left(1+{\displaystyle \frac{P_\mathrm{\Delta }^2}{r_\mathrm{\Delta }^2}}\right)|P_\mathrm{\Delta }|,`$ (62)
$`=`$ $`M_{\mathrm{ADM}}|P_\mathrm{\Delta }|.`$ (63)
Let us now consider the case in which the Abelian solution has both electric and (unit) magnetic charge. The surface gravity is given by,
$$\kappa =\frac{1}{2r_\mathrm{\Delta }}\left[1\frac{(Q_\mathrm{\Delta }^2+P_\mathrm{\Delta }^2)}{r_\mathrm{\Delta }^2}\right].$$
(64)
Equation (47) leads us to conclude that $`\mathrm{\Phi }_\mathrm{\Delta }=\frac{Q_\mathrm{\Delta }}{r_\mathrm{\Delta }}`$, and equation (48) takes the form,
$$_{r_\mathrm{\Delta }}V=\frac{(Q_\mathrm{\Delta }^2+P_\mathrm{\Delta }^2)}{r_\mathrm{\Delta }^2}+\frac{Q_\mathrm{\Delta }^2}{r_\mathrm{\Delta }^2}=\frac{P_\mathrm{\Delta }^2}{r_\mathrm{\Delta }^2}.$$
(65)
Then, imposing again the boundary condition that $`V`$ vanishes on extremal solutions. we have that,
$$V=\frac{P_\mathrm{\Delta }^2}{r_\mathrm{\Delta }^2}\frac{P_\mathrm{\Delta }}{\sqrt{P_\mathrm{\Delta }^2+Q_\mathrm{\Delta }^2}}.$$
(66)
The Horizon Mass is now,
$$M_\mathrm{\Delta }=M_{\mathrm{ADM}}\frac{P_\mathrm{\Delta }}{\sqrt{P_\mathrm{\Delta }^2+Q_\mathrm{\Delta }^2}},$$
(67)
and the total energy is then $`E=\frac{P_\mathrm{\Delta }}{\sqrt{P_\mathrm{\Delta }^2+Q_\mathrm{\Delta }^2}}=\frac{1}{\sqrt{1+Q_\mathrm{\Delta }^2}}`$.
Finally, there is the most interesting case, i.e., the family of colored black holes labeled by $`r_\mathrm{\Delta }`$ and an integer $`n`$. Since these solutions correspond to the purely magnetic case, Equations (55), (56), (58) and (59) continue to hold. This last condition that the function $`V=V(r_\mathrm{\Delta })`$ should satisfy can be written as,
$$V^{}=\frac{r_\mathrm{\Delta }}{2}\beta ^{},$$
(68)
with ‘prime’ denoting differentiation with respect to $`r_\mathrm{\Delta }`$. (We are considering variation with fixed value of $`n`$.) Furthermore, by requiring that $`M_\mathrm{\Delta }0`$ as $`r_\mathrm{\Delta }0`$, -coming from physical considerations- we arrive at the following relation,
$$M_\mathrm{\Delta }=\frac{1}{2}_0^{r_\mathrm{\Delta }}\beta (\stackrel{~}{r})d\stackrel{~}{r},$$
(69)
where the integration is again performed over the space of parameters of the $`n`$-colored black hole, labeled by the horizon radius $`r_\mathrm{\Delta }`$, and not over space-time. Let us note that for the $`n=0`$ solution, where $`\beta `$ is known in closed form ($`\beta =1`$), we arrive at $`M_\mathrm{\Delta }^{(n=0)}=r_\mathrm{\Delta }/2=\kappa a_\mathrm{\Delta }/(4\pi )`$, as expected.
Several remarks are in order. First, we must emphasize that the determination of $`V`$, and thus of $`M_\mathrm{\Delta }`$ relied on considerations involving only variations of quantities associated with the horizon $`\mathrm{\Delta }`$. Thus, the horizon mass is a well defined quantity in the isolated horizons phase space $``$ (provided that a global normalization of $`\mu `$ exists). Second, the HHM defined by (42), when restricted to SSS configurations, does not agree with the usual definitions of mass that one finds in the literature (see for instance and ). It should be stressed that (69) comes from a consistent Hamiltonian formulation, and is not a definition as occurs in other treatments.
As it was discussed at the end of Section V, the fact that SSS solutions to the EYM equations do not span the space $``$ of horizon parameters is, in a sense, disturbing. It might seem that the Einstein-Yang-Mills system is in a different status than the Einstein-Maxwell-dilaton (EMD) system where the space of horizon parameters is in a one to one correspondence with the space of static solutions. One would like to have a similar result in the EYM case. However, the SSS solutions span only a subspace of $``$. Luckily, the Spherically Symmetric solutions in EYM do not exhaust all possible Static solutions (as occurs in EMD); there are static solutions with axial symmetry that are not spherically symmetric . With these results at hand, we propose a completeness conjecture in the next Section.
## VII The Canonical Normalization Problem: A Proposal
We have seen that in order for the isolated horizons scheme to define the surface gravity and the horizon mass of the colored black holes we needed to introduce a uniqueness conjecture $`C1`$ that guarantees that, given the isolated horizon parameters $`a_\mathrm{\Delta },P_\mathrm{\Delta },Q_\mathrm{\Delta }`$, there would be at most one SSS solution. This was needed for, otherwise, the normalization of $`\mu `$ would not be uniquely specified given those parameters. The existing numerical evidence does indeed strongly support this conjecture. However, as we have mentioned before, and as is evident from the previous discussion, this is not sufficient in order to have the Isolated Horizon framework working for the EYM system to the same extent that it works, say, for the Einstein Vacuum, Einstein-Maxwell, and Einstein-Maxwell-Dilaton systems. In order to achieve that, we would need to have a canonical normalization of $`\mu `$ for a ‘complete’ set of values of the Isolated Horizon parameters. In the previously mentioned cases this canonical choice is given by the existence of static (and spherically symmetric) Black Hole solutions for all isolated horizon values of the parameters.<sup>\**</sup><sup>\**</sup>\** In fact, although this is true for the Einstein Vacuum system, in the Einstein Maxwell case we already have a potential problem, because if $`Q>r_\mathrm{\Delta }`$ there are no such static solutions. The consistency of the whole scheme would require the impossibility of constructing a space-time containing an Isolated Horizon with such values of the parameters. This is a rather serious consequence of the present point of view, and one that should be testable. There is a strong correlation of these issues with the cosmic censorship conjecture, that prevents us from violating the inequality $`Q<r_\mathrm{\Delta }`$ for static solutions. It would seem that the IH formalism implies that one can not construct initial data for a solution containing a black hole with values of the parameters that violate this inequality.
There is strong numerical evidence against the validity of the analogous claim in the case of the EYM system. In fact in the regime of staticity and spherical symmetry there are, given a fixed value of $`a_\mathrm{\Delta }`$, only a discrete set of values of $`P_\mathrm{\Delta }`$ for which there are black hole solutions. Moreover, within this regime there are no Black Hole solutions for any value of $`P_\mathrm{\Delta }1,0`$ and $`Q_\mathrm{\Delta }0`$. Thus if we want to have any hope that any claim in that direction might be true, we must formulate it outside this restrictive regime. Indeed the fact that in EYM systems there are static Black Hole solutions that are not spherically symmetric, (indicating that that the analogous to Israel’s theorem is false in this case), already shows us that we must go beyond the SSS regime. In fact the solutions alluded above are axially symmetric, instead of spherically symmetric, but seem to share, with the SSS solutions, the discreetness of the allowed values of $`P_\mathrm{\Delta }`$ (at least to the extent that this issue has been studied ). Thus we have to go beyond this regime as well. In fact there are strong indications (see for example the discussion in ) that we must go beyond the static regime, and pose the conjecture in a broad enough setting that would still allow one to single out, for a given choice of IH parameters, a particular black hole solution and thus a canonical normalization of $`\mu `$. This would be of course the class of stationary black hole solutions, where we would have to keep track also of the angular momentum, both at infinity $`J_{\mathrm{}}`$ and at the horizon $`J_\mathrm{\Delta }`$. The completeness conjecture would thus be: $`C2`$: For every value of the Isolated Horizon parameters $`a_\mathrm{\Delta },P_\mathrm{\Delta },Q_\mathrm{\Delta },J_\mathrm{\Delta }`$ for which a space-time can be constructed, there exist also a stationary Black Hole Solution with the same value of the parameters, now characterizing the Killing Horizon.<sup>††</sup><sup>††</sup>††In this statement, a space-time ‘can be constructed’ whenever there exists an asymptotically flat solution to the EYM equation satisfying IH boundary conditions with the specified values of the horizon parameters.
Let us now consider some of the implications of this conjecture. First, a stationary black hole solution would be characterized by its parameters at infinity: $`M_{\mathrm{ADM}},P_{\mathrm{}},Q_{\mathrm{}},J_{\mathrm{}}`$ and therefore the conjecture would imply the existence of a well defined map $`\mathrm{\Psi }:(a_\mathrm{\Delta },P_\mathrm{\Delta },Q_\mathrm{\Delta },J_\mathrm{\Delta })(M_{\mathrm{ADM}},P_{\mathrm{}},Q_{\mathrm{}},J_{\mathrm{}})`$ The failure of the no hair conjecture would indicate that this map is not invertible. In fact we know that it would not be injective. Moreover the map would be nontrivally four dimensional, in the sense that fixing, say, $`J_\mathrm{\Delta }=0`$ would not fix $`J_{\mathrm{}}=0`$, as can be seen from the following expression,
$$4\pi \mathrm{\Omega }(J_\mathrm{\Delta }J_{\mathrm{}})=_\mathrm{\Sigma }\mathrm{d}^3x\left(t^b\stackrel{~}{𝐄}_i^a𝐅_{ab}^i+_t(𝐀_a^i)\stackrel{~}{E}_i^a\right),$$
(70)
valid for stationary black hole solutions in EYM theory. Here, $`\mathrm{\Sigma }`$ is a maximal hyper-surface intersecting the bifurcate horizon, and $`t^a`$ is the projection to $`\mathrm{\Sigma }`$ of the time translation Killing field. As it can be seen from the Eq. (70), there is a bulk contribution to $`J_{\mathrm{}}`$, the canonical angular momentum at infinity. Here, $`J_\mathrm{\Delta }`$ is a particular definition of ‘horizon angular momentum’ (given by a Komar integral), and $`\mathrm{\Omega }`$ stands for the angular velocity of the horizon, i.e, the expression appearing in the first law
$$\delta M_{\mathrm{ADM}}+\mathrm{\Phi }_{\mathrm{}}\delta Q_{\mathrm{}}\mathrm{\Omega }\delta J_{\mathrm{}}=\frac{1}{8\pi }\kappa \delta a_\mathrm{\Delta },$$
(71)
where $`\mathrm{\Phi }_{\mathrm{}}`$ is the ‘electric potential’ at infinity. In fact the EYM system is, in this respect, rather different from the Einstein-Maxwell system, because in the latter one can disentangle for example the expression for $`\mathrm{\Phi }_{\mathrm{}}Q`$ from the expression for $`\mathrm{\Omega }J_{\mathrm{}}`$, something that can not be done in the former, in which case the only relationship that can be obtained is given by the expression,
$$8\pi (\mathrm{\Omega }J_{\mathrm{}}\mathrm{\Phi }_{\mathrm{}}Q_{\mathrm{}})=_\mathrm{\Sigma }N(\pi _{ab}\pi ^{ab}+2𝐄_i^a𝐄_a^i)/h,$$
(72)
with $`\pi ^{ab}`$ the momentum conjugate to $`h_{ab}`$, the Riemannian metric on $`M`$ and $`N`$ the ‘lapse function’.
Assuming the validity of $`C2`$ one would have a canonical choice for the normalizations that could be used to uniquely define $`\kappa `$ and $`M_\mathrm{\Delta }`$, the one provided by the Killing field of the stationary solution that is null at the horizon and that is normalized so that at infinity is a unit time translation. In order to make all these considerations more precise from the isolated horizons point of view, one needs to consider the extension of the formalism given in and .
Now, let us concentrate for the moment in the SSS sector and see if we can understand the discreteness observed there in terms of this conjecture. In the analysis that lead to the discovery of SSS Black Holes in EYM theory, one is fixing $`P_{\mathrm{}}=0`$ because one is interested in solutions that are Abelian at large distances. Moreover spherical symmetry evidently requires $`J_{\mathrm{}}=0`$ and $`J_\mathrm{\Delta }=0`$. Moreover the mixture alluded to before would prevent us from achieving this (spherical symmetry) unless we also set $`Q_{\mathrm{}}=0`$. Thus we see that the (highly nonlinear) problem is given by four constraints in a four dimensional space, and thus that the set of solutions is expected to be given by a discrete set (i.e. the linearized problem about a given solution in the SSS sector has no nontrivial solution within the sector).
Next, let us consider some consequences of the completeness conjecture in the structure of the space $``$. We have seen in Sec. VI that the space $``$ (where we now have to include distortion and rotation ) is foliated by ‘symplectic leaves’ where the Hamiltonian formulation and the first law are valid. This foliation intersects the space $`𝒮`$ of stationary solutions and defines a one parameter foliation of it. Now, if the $`C2`$ conjecture is valid, we have an isomorphism between $`𝒮`$ and $``$, which then induces a canonical foliation of $``$.
On the other hand, we must point out the following heuristic argument against the conjecture. Consider a stationary black hole solution in EYM theory, in order to be asymptotically flat the YM field strength must fall off rather rapidly at large distances, thus the self interaction of the fields must be falling off faster than the fields themselves and thus the fields must behave in the large distance limit as free fields, i.e., as Abelian fields. This suggests that the only possible values of the magnetic charge at $`\mathrm{}`$, $`P_{\mathrm{}}`$ are the Abelian values $`0,1,..`$ (there is of course no such restriction on the allowed value of the electric charge). This view is supported by the experience with the static spherically symmetric solutions. Thus, any stationary solution would be characterized at infinity by the parameters $`M_{\mathrm{ADM}},Q_{\mathrm{}},J_{\mathrm{}}`$ (setting for the moment $`P_{\mathrm{}}=0`$ to simplify the discussion). On the other hand the solution will be characterized by its horizon parameters amongst them $`a_\mathrm{\Delta }`$. We know from the first law in its asymptotic infinity version Eq. (71) (see the discussion in ), that the stationary black hole solutions are extrema of $`M_{\mathrm{ADM}}`$ at fixed $`a_\mathrm{\Delta },Q_{\mathrm{}},J_{\mathrm{}}`$ within the constrained phase space (i.e., the space of allowed initial data, which as usual, can be identified with the space of solutions). Each such extrema is an isolated point in that space, so the manifold of stationary solutions is three dimensional. In an argument is given that indicates that there would be a countable infinity of solutions for each value of the parameters $`a_\mathrm{\Delta },Q_{\mathrm{}},J_{\mathrm{}}`$ which would generalize what happens in the case of the static solutions. This suggest that the manifold of stationary solutions is made of a countable infinity of connected three dimensional components. On the other hand, the conjecture would seem to indicate that such manifold must be four dimensional. The only way to avoid this would require the impossibility of constructing solutions of the equations representing asymptotically flat black hole space-times (not necessarily stationary) containing isolated horizons for all values of the isolated horizon parameters.
One would be tempted to take such position in view of the discretness of the SSS colored black holes with a fixed value of he horizon area $`a_\mathrm{\Delta }^0`$. Let $`\{P_\mathrm{\Delta }^i\}_{i=1}^{\mathrm{}}`$ be the values of the horizon magnetic charges of the SSS colored black holes with horizon area $`a_\mathrm{\Delta }^0`$. One might want to argue that given a value of $`P_\mathrm{\Delta }^{}(0,1)`$ that is not in that list, and is, say, between two of the values in the list, one can not construct a solution representing asymptotically flat black hole space-time with isolated horizon and, say $`a_\mathrm{\Delta }=a_\mathrm{\Delta }^0,P_\mathrm{\Delta }=P_\mathrm{\Delta }^{},Q_\mathrm{\Delta }=0,J_\mathrm{\Delta }=0`$. Unfortunately this claim is evidently false: The recipe for constructing such a space-time is to give initial data that, upon evolution would be static near the horizon and near infinity at least for a finite “time” interval. Take the equations for a the SSS (Eqs. (10),(11) and (12) in ) and set $`r_H=\sqrt{a_\mathrm{\Delta }^0/4\pi }`$, $`w(r_H)=\sqrt{1P_\mathrm{\Delta }^{}}`$, and evolve the elliptic equations up to say $`r=2r_H`$ (if we continue to evolve the equations attempting to obtain a static solution we would find that $`w`$ diverges so the solution would not be asymptotically flat). For, say, $`r5r_H`$ take the initial data for the Schwarzschild solution. In the intermediate region $`r(r_H,5r_H)`$ take any interpolating function for $`w`$, and set the time derivatives of the functions $`w,m,etc`$ (i.e the “momenta”) to satisfy the constraints. The point is that the evolved space-time will be static in a neighborhood of the horizon so, in particular, the Horizon will be isolated, at least during some finite time interval (until the radiation coming from the intermediate region arrives at the horizon). Note that, generically, in this case the event horizon will fail to coincide with the isolated horizon.
The argument above suggest that the manifold of isolated horizon parameters $``$ is indeed four dimensional and thus the conjecture would require the identification of this four dimensional manifold with the infinite set of three dimensional components that seem to constitute the manifold of stationary solutions. The best that can be hoped at this point is that the latter be dense in the former, a situation that would indicate that in the EYM theory there is much richer structure that what is found in, say, the Einstein-Maxwell system. In this case, the canonical normalization for $`\mathrm{}`$ would be given by an appropiate limit (within $`𝒮`$) of solutions where the normalization exists (assuming, of course, that there is such limit).
On the other hand the argument above is by not means a tight proof, particularly so in the case of the conclusion about the Abelian nature of the allowed values of $`P_{\mathrm{}}`$. As we have mentioned before, the validity of this conjecture, or some version of it (as for example a version based on the assumption that the manifold of stationary solutions is mapped into a dense subset of the isolated horizon parameters), seems to be the only reasonable way in which the Isolated Horizon scheme can be as successful in the general case as it has proven to be in the Einstein Vacuum, Einstein-Maxwell and Einstein-Maxwell-Dilaton theories.
Needless is to say that further research is required in order to elucidate whether one of the scenarios considered above is correct or whether, in fact, the Isolated Horizon scheme fails to achieve in EYM theory the same degree of success that is attained in the previously treated cases.
## VIII Spherically Symmetric Static Solutions: Mass and Hair
In this section we shall restrict our attention to the SSS sector of isolated horizons. One issue that has been considered for non-Abelian gauge theories is the relation that might exist between the existence of regular static, solitonic solutions and ‘hairy’ black hole solutions. This issue has been considered, for example, in from heuristic and dimensional arguments. In this section, two main issues are studied. First, by restricting the Hamiltonian formulation for Isolated black holes to the SSS sector, we can define the Hamiltonian Horizon Mass (HHM) of SSS black holes in EYM theory. We then use this expression to show that this quasi-local definition together with some basic properties of Hamiltonian Mechanics lead us to a formula relating HHM and ADM mass of the colored BH solutions with the ADM mass of the Solitons of the theory. We also conclude that the positivity of the ‘total energy’ spectrum of the colored black holes is related to their instability.
These results are quite surprising, because the IH formalism was developed to extend the notion of black holes to situations where radiation is present –and goes out to infinity– and one might have not expected to obtain new results already in the static sector of the theory.
As we have previously mentioned, the HHM Mass $`M_\mathrm{\Delta }`$ of a SSS BH does not correspond to any of the Quasi-local definitions of BH Mass considered in the literature. Furthermore, $`M_\mathrm{\Delta }`$ has the virtue of being constructed from a consistent Hamiltonian formulation which places it on a different status as the standard definitions.
To begin, let us calculate the value of the ‘total energy’ $`E`$ of the system, for the three sectors of SSS solutions. In order to do this, we use a general argument from symplectic geometry that states that, within each connected component of the space of static solutions $`𝒮`$ embedded in the space of isolated horizons, the value of the Hamiltonian $`H_t`$ remains constant . Let us review this argument since it is essential for our discussion. The Hamilton equations of motion can be written as $`\delta H=\mathrm{\Omega }(\delta ,X_H)`$, where $`\mathrm{\Omega }`$ is the symplectic form, $`\delta `$ is an arbitrary variation and $`X_H`$ is the Hamiltonian vector field. A static solution is one at which the Hamiltonian vector field either vanishes or generates pure gauge evolution. In either case, the symplectic structure evaluated on $`X_H`$ and any arbitrary vector field $`\delta `$ vanishes. Therefore, for this point of the phase space, $`\delta H=0`$ for any direction $`\delta `$. In particular $`\delta H=0`$ for variations relating two static solutions. Now, in the case of Einstein-Maxwell theory, the no-hair theorems ensure that all static solutions are given by the RN family. That is, the space of static solutions is in that case, connected. Furthermore, since there is no energy scale in the theory, the only ‘preferred’ value for $`H_t`$ is zero .
What is the situation in Einstein-Yang-Mills theory? First, there is the Abelian family of electrically charged solutions, that represent a connected component, parameterized by $`M,Q`$. For these solutions, the basic reasoning of applies and, as follows from (53) and (54), one has to conclude that for these solutions $`H=0`$. However, there is a subtle modification in the case of magnetic Abelian solutions. These solutions represent a disconnected component labeled by one parameter, namely, the mass $`M`$ (the magnetic charge $`P`$ is fixed to be unity). As discussed in Section II, the EYM system possesses an energy scale given by the YM coupling constant, so in principle, non-zero values of $`H`$ are allowed. In the one-dimensional component corresponding to Abelian magnetic solutions, the value of $`E`$ can be computed using (54) and (63) and is given by $`E=M_{\mathrm{ADM}}M_\mathrm{\Delta }=|P_\mathrm{\Delta }|=1`$. (We are taking the YM coupling constant $`g_{\mathrm{YM}}=1`$.)
Finally, let us consider colored black holes. Each connected component of the space of SSS colored black holes is one-dimensional (parameterized by $`a_\mathrm{\Delta }`$), and solutions corresponding to distinct values of $`n`$ belong to disconnected components. That is, the space SSS has a countable number of connected components. As we shall now show, for $`n1`$ the value of the Hamiltonian turns out to be different from zero: $`H_t^n0`$.
Recall that the general argument described above tells us that the (on shell) value of the Hamiltonian is constant for each family labeled by $`n`$. This in particular implies that its value is independent of the radius $`r_\mathrm{\Delta }`$ of the horizon. Thus one is allowed to take the limit
$$H^{(n)}=\underset{r_\mathrm{\Delta }0}{lim}[M_{\mathrm{ADM}}^{(n)}(r_\mathrm{\Delta })M_\mathrm{\Delta }^{(n)}(r_\mathrm{\Delta })].$$
(73)
Now, it is known that the colored black holes converge point-wise to the Bartnik-McKinnon soliton solutions and that the ADM mass satisfies $`M_{\mathrm{ADM}}^{(n)}M_{\mathrm{BK}}^{(n)}`$ when $`r_\mathrm{\Delta }0`$. Furthermore, the horizon mass of the black hole $`M_\mathrm{\Delta }`$ goes to zero in this limit, so we can conclude that
$$H^{(n)}=M_{\mathrm{BK}}^{(n)},$$
(74)
that is, the total value of the Hamiltonian equals the mass of the $`n`$th Bartnik-McKinnon soliton solution!
We now collect our results for colored black holes and arrive at the following unexpected relation,
$$M_{\mathrm{ADM}}^{(n)}(r_\mathrm{\Delta })=M_{\mathrm{BK}}^{(n)}+M_\mathrm{\Delta }^{(n)}(r_\mathrm{\Delta }).$$
(75)
Thus, we are in the position of writing an explicit formula for the ADM mass of the $`n`$ colored black hole as function of $`r_\mathrm{\Delta }`$,
$$M_{\mathrm{ADM}}^{(n)}(r_\mathrm{\Delta })=M_{\mathrm{BK}}^{(n)}+\frac{1}{2}_0^{r_\mathrm{\Delta }}\beta ^{(n)}(\stackrel{~}{r})d\stackrel{~}{r}.$$
(76)
In Figure 2, we show the values of the HHM as a function of the horizon radius $`r_\mathrm{\Delta }`$, for the $`n=1,2`$ families. Note that for a given value of the horizon area, the higher $`n`$ is, the lower the horizon mass of the corresponding black hole. In Figure 3 the value of the horizon magnetic charge $`P_\mathrm{\Delta }`$ is shown as function of the radius.
It is important to stress that, a priori, one would not expect to get the value of quantities defined at infinity, like the difference of ADM masses in terms of purely local quantities at $`\mathrm{\Delta }`$. At this point one might raise the following objection to the construction of $`M_\mathrm{\Delta }`$ for colored black holes: If we start by considering the equation $`\delta M_\mathrm{\Delta }=\kappa \delta a_\mathrm{\Delta }/4\pi `$, and try to integrate it along the one-dimensional curve defined for each $`n`$, one can “trivially” do so by using the usual form of the first law at infinity, which tells us that the general solution for $`M_\mathrm{\Delta }`$ is given by $`M_\mathrm{\Delta }=M_{\mathrm{ADM}}+c`$, with $`c`$ a constant. Then one might argue that the ‘only’ thing one is doing is to set $`c`$ such that $`M_\mathrm{\Delta }(r_\mathrm{\Delta }=0)=0`$. Thus, one would conclude, the derivation is trivial and even the notion of a Horizon mass would seem questionable. This argument would be perfectly valid had we postulated the equation $`\delta M_\mathrm{\Delta }=\kappa \delta a_\mathrm{\Delta }/4\pi `$. However, the non-trivial point here is that this equation comes as a consequence (and consistency requirement) from a Hamiltonian description respecting –physically motivated– boundary conditions. Therefore, even when the algebraic manipulations are simple, the final result is highly non-trivial.
We can now try to understand the physical meaning of the relation (76). Two facts are known about these solutions: first, we know that for fixed $`a_\mathrm{\Delta }`$ these solutions represent saddle points of the ADM mass function $`M`$ , and thus, as one can expect, for all values of $`n`$ these solutions are unstable under small perturbations . Let us now note that for the reported solutions in the literature (see, for instance, ), the BK Mass is a monotonic function of $`n`$, starting at $`M_{\mathrm{BK}}^10.828`$ and approaching $`1`$ as $`n`$ grows (in standard normalized units). The fact that the mass of the soliton, and therefore the total energy of the colored black holes is positive, confirms our expectation, coming from energetic considerations, that in general $`M^{\mathrm{ADM}}M_\mathrm{\Delta }`$. Indeed, since the difference between the Horizon and ADM masses can be seen as the energy that is available for radiation to fall both into the black hole and to infinity, one can understand the nonzero value of the Hamiltonian as an indication that there is a potentiality for instability of the solution. In static solutions there is, of course, no radiation. Thus, to be precise, a positive value of the energy $`E`$ means that, if we perturb slightly the initial data of a (unstable) static solution in such a way that the total energy is ‘very close’ to $`E`$, then the resulting space-time will approach Schwarzschild in the future, and the total radiated energy to both infinity and the horizon will be equal to $`E`$. In conclusion, a necessary condition for the solution to be unstable is for the value of the total energy on the solution in question to be positive.
Let us conclude this section with four remarks.
1. In the computation of the Horizon Mass for magnetically charged Abelian solution, we required that the function $`V`$ vanishes for the extremal black holes. Let us now motivate this choice. first, it is known that for $`r_\mathrm{\Delta }1`$, the $`n`$ colored black holes approach the Reissner Nordstrom magnetic solution (in a region around the horizon that expands unboundly with $`n`$) when $`n\mathrm{}`$. Let us abuse notation for a moment and refer to these limiting solution as the ‘$`n=\mathrm{}`$ colored black hole’. From the numerical solutions reported in the literature one can notice that the Horizon Mass $`M_\mathrm{\Delta }^{(n)}=M_{\mathrm{ADM}}^{(n)}M_{\mathrm{BK}}^{(n)}`$ tends to zero when $`n\mathrm{}`$ and $`r_\mathrm{\Delta }1`$ from above. Now, if we require continuity from above for $`M_\mathrm{\Delta }`$ on the space $``$ of SSS solutions, we should require that
$$\underset{r_\mathrm{\Delta }1^+}{lim}M_\mathrm{\Delta }^{(\mathrm{})}=0.$$
(77)
But the Horizon mass is given by $`M_\mathrm{\Delta }^{(\mathrm{})}=\frac{\kappa a_\mathrm{\Delta }}{4\pi }+V`$, and $`\kappa 0`$ as $`r_\mathrm{\Delta }1`$ (since this corresponds to the extremal RN case). Thus, we conclude that $`V0`$ when $`r_\mathrm{\Delta }1`$.
2. Let us now return to the puzzle that we posed in Section II: how can we reconcile the treatments of static solutions at infinity and at the horizon? That is, when formulating the first law, say at infinity, we only have two free parameters given by the ADM mass and electric charge $`Q_{\mathrm{}}`$. However, at the horizon we have three parameters, the horizon area $`a_\mathrm{\Delta }`$, the electric charge $`Q_\mathrm{\Delta }`$ and the magnetic charge $`P_\mathrm{\Delta }`$. Now, equation (45) tell us that their respective variations depend only on the horizon area and electric charge, which seems to be a contradiction. The proper setting of the problem again comes from the consistency of the Hamiltonian formulation. Recall that from Equation (46) we concluded that the magnetic charge $`P_\mathrm{\Delta }`$ is not free to vary at the horizon, but its variations are related to the variations of the horizon radius $`r_\mathrm{\Delta }`$. This means that the allowed values of $`P_\mathrm{\Delta }`$ and $`r_\mathrm{\Delta }`$ lie in sub-manifolds of co-dimension one embedded in the space $``$. Thus, even when we can in principle have three independent parameters at the horizons, in practice when employing the Hamiltonian formulation, one is restricted to values of the magnetic charge that are determined by the horizon area (within each family labeled by $`n`$). Thus, the variations of the parameters both at the horizon and at infinity are consistent. This can also be seen in the following way. The results of different formalisms at infinity have shown that the ADM mass varies as ,
$$\delta M_{\mathrm{ADM}}=\frac{1}{8\pi }\kappa \delta a_\mathrm{\Delta }+\mathrm{\Phi }_{\mathrm{}}\delta Q_{\mathrm{}},$$
(78)
but since we know that, at a static solution, an arbitrary variation $`\delta `$ satisfies $`\delta E=\delta (M_{\mathrm{ADM}}M_\mathrm{\Delta })=0`$, we have complete agreement with (45).
Now lets turn our view to the conflict that our completeness conjecture $`C2`$ seems to face in view of the fact that the colored black hole solutions, with different value of $`n`$, have different values of the Hamiltonian and the general argument, presented in Sec VI, ensures that $`H`$ must be constant over any sub-manifold of static (stationary) solutions. Thus, this general argument would lead us to conclude that $`H`$ is constant over $`𝒮`$. The solution of the apparent paradox lies in the fact, already mentioned in Sec VI that the consistency of the Hamiltonian formulation leads to a foliation of phase space by symplectic leaves over which the Hamiltonian formulation is valid. These leaves intersect the manifold of stationary configurations, which is embedded in $``$. The intersection results in hyper-surfaces of constant value of the Hamiltonian. In the case of SSS each one of those corresponds to a single family labeled by a fixed value of $`n`$. If we now restrict our consideration to the hyper-surface with $`Q=0`$ each of these families correspond to a curve in the $`(r_\mathrm{\Delta },P_\mathrm{\Delta })`$ plane. Namely, $`P_\mathrm{\Delta }`$ becomes a function of $`r_\mathrm{\Delta }`$ rather than an independent parameter. This again is similar to what happens in the case of a rotating body, where the manifold of stationary states does not correspond to a single value of the Hamiltonian, but the intersection of a symplectic leave with this manifold coincides with the curves of constant value of the Hamiltonian. Then, even when the space $``$ if foliated by a ‘continuum’ of leaves, the allowed Hamiltonian motions are restricted to lie within each of these level surfaces of the Hamiltonian $`H`$. This overall picture of the structure of the parameter space is in fact, consistent with the situation in EMD. In this case, there is only one ‘leave’ and thus, the intersection with the manifold of static solutions has also only one leave; the value of $`H`$ is indeed constant on the whole $`𝒮`$ (it is in fact zero).
It is intriguing to note that there seems to be a deep relation between the existence of non-trivial solitons and hairy black holes for which the charges are not independent (i.e., the possible SSS solutions are restricted to a discrete set of constant energy surfaces within $``$). The fact that the consistency requirement on the Hamiltonian formulation leads to the ‘nonstandard’ Hamiltonian framework for these cases (i.e the foliation of phase space by the (symplectic) leaves on which there is a truly Hamiltonian framework) together with the constancy of the full Hamiltonian on stationary solutions can be regarded as explaining such relation. That is, on the intersection of each leave with the manifold $`𝒮`$ of stationary solutions, the full Hamiltonian is a constant. In the limit $`a_\mathrm{\Delta }0`$ on each leave (provided such limit can be taken) we will find a soliton. The motion along the ‘leave’ in $`𝒮`$ determines the mutual dependence of the Isolated Horizon parameters.
3. It is now a general belief that the no-hair conjecture, even in its weakest form , is violated for some systems. One can thus hope that there be a uniqueness result for static solutions in terms of quasi-local parameters at the horizon. In Section VI we have put forward a ‘horizon parameters uniqueness conjecture’ stating that all static BH solutions are characterized by their horizons parameters (‘quasi-local charges’) in a unique way. In particular, it should be true that, given the horizon area $`a_\mathrm{\Delta }`$ and the horizon electric charge $`Q_\mathrm{\Delta }`$ and magnetic charge $`P_\mathrm{\Delta }`$, the Static solutions be uniquely determined. For instance, if we set $`Q_\mathrm{\Delta }`$, given an arbitrary value for the magnetic charge $`P_\mathrm{\Delta }[0,1]`$, there might be no solution, but if there is a solution, it should be unique. The numerical evidence available supports the conjecture (see Figure 3).
4. One other point already mentioned is the issue of the stability test provided by this type of analysis: It is only when $`M_{\mathrm{ADM}}>M_\mathrm{\Delta }`$ in (54) that the solution can be unstable. One very clear example of this is given by the magnetic RN solution, which can be considered within both the Einstein Maxwell (EM) theory and the EYM theory. This solution is stable within EM but unstable within EYM . We can now understand this situation in the following way: In the former case the gauge connection can not be globally described through a gauge field $`A_a`$ as it corresponds to a nontrivial bundle. In this case one can nevertheless apply the IH formalism through the use of the duality symmetry of Maxwell theory. This results in the appearance of a term $`\mathrm{\Phi }_MP_\mathrm{\Delta }`$ taking the place of $`\mathrm{\Phi }Q_\mathrm{\Delta }`$ in (42), and as it is well known the evaluation of $`E`$ as $`M_{\mathrm{ADM}}M_\mathrm{\Delta }`$ (in this case $`H^{(0)}=0`$ as there are no Abelian magnetically charged regular solitons) gives $`E=0`$, thus accounting for the stability of the solution. Let’s look at what happens in the EYM theory. In this case the solution can be described in terms of the gauge fields $`A_a^i`$ because it is associated with a trivial bundle (with larger group) and there is therefore no term of the form $`\mathrm{\Phi }_MP_\mathrm{\Delta }`$ in (42). The $`P_\mathrm{\Delta }`$ dependence of the Mass $`M_\mathrm{\Delta }`$ comes through $`V`$. As we have shown, the value of $`E`$ for this solutions is positive ($`E=1`$), thus allowing for the instability of these solutions. Note that this same argument is valid also for dyons with both electric and unit magnetic charge (See Equation (64) to (67)), and thus, it indicates a potential instability of these solutions.
Finally, let us conclude this remark by suggesting a ‘rule of thumb’ for finding potentially unstable solutions, suggested by the EYM system. In the static family of solutions, consider the limit $`r_\mathrm{\Delta }0`$. We have three possibilities: i) We arrive at a regular solution with zero energy (i.e. Minkowski). This indicates that the whole family, labeled by $`r_\mathrm{\Delta }`$, is stable; ii) There is a minimum allowed value of $`r_\mathrm{\Delta }`$ corresponding to zero surface gravity. In this case, we can not conclude anything, and; iii) In the limit one finds a regular solution with positive energy (a soliton different from the vacuum). In this case, the whole family of solutions (including the soliton) is potentially unstable. It would be interesting to re-examine, from this perspective, the (complete non-linear) stability of the Einstein-Skyrme Solitons and Black Holes .
## IX Discussion
In this note, we have studied the extension of the Isolated Horizon formalism to include the EYM system and found that it leads to a ‘nonstandard’ Hamiltonian formulation. The main feature of this formulation is that it provides a foliation of phase space into symplectic leaves in each of which we do get a standard Hamiltonian formulation. The framework nevertheless provides a powerful tool for studying some classical aspects of the theory already at the Static level. In particular, we found an unexpected relation between the the ADM mass of a static spherically symmetric black hole solution, its Horizon mass and the ADM mass of the corresponding solitonic solution. These relationships were checked numerically in terms of the known numerical results obtained in the process of finding those solutions and thus the agreement can be seen as a check on the whole formalism.
An apparent tension and a challenge for the formalism is given by the existence of hairy solutions, where the number of charges at infinity and at the horizon do not coincide. We have been able to pin-point the problem in a precise way using the Isolated Horizons formalism. This involves the analysis of the consistency of the Hamiltonian formulation and the nature of the first law. We have encountered difficulties in defining a canonical normalization for the vector $`\mathrm{}^a`$ and thus, for the Horizon Mass in general. We have proposed possible resolutions for this difficulty. Motivated by all these results, and in order to have a satisfactory treatment of the EYM system within the framework, we have put forward a ‘quasi-local uniqueness ’ and a ‘completeness’ conjectures for Stationary Black holes. In the case of the later we have put forward arguments both in favor and against it, but the main point is that some version of it seems to be the only possibility to have the Isolated Horizon framework working in EYM theory to the same extent that it does in, say, Einstein-Maxwell theory.
The present work can be generalized in several directions. First, our analysis allows us to propose a Isolated Horizon treatment for general theories containing non-trivial black holes and soliton solutions; it should be possible to apply the type of analysis presented here to these theories where nontrivial regular static solutions have been found. In particular, Einstein-Yang-Mills-Higgs, Einstein-Yang-Mills-Dilaton, and Einstein-Skyrme Theories, are examples in which there are both, solitonic and Black Hole solutions. In all these cases, formulas analogous to the EYM case can in principle be found by a straightforward application of the analysis carried out in the last Section. In particular, one should be able to to compute the ‘total energy’ of the ‘hairy’ black holes to test for a potential instability.
Second, one can use the very recent results of Ashtekar and collaborators who have been able to extend the isolated horizons framework to include distorted horizons as well as rotating horizons . Thus, one should be able to use the formalism for distorted horizons in order to study colored black holes which are static but not spherically symmetric . This analysis would be a first check for the ‘quasi-local uniqueness conjecture $`C1`$’ that we have proposed . The discussion of the preceding sections suggests that by restricting our attention to non-distorted, non-rotating horizons we were lead to a consistent but ‘incomplete’ formalism in the EYM system; a complete treatment (i.e., one providing a canonical choice of ‘normalization’ for all $``$, based on stationary solutions which are contained in the formalism) of EYM isolated horizons should be given within the context of distorted and rotating Isolated Horizons . It would be interesting to investigate this matter once the articles are made public.
Finally, the present analysis shows that the so called colored black hole solutions provide a nontrivial testing ground for the approach of to evaluate the statistical mechanical entropy of black holes. In particular it is known that by selecting the value $`\gamma =\mathrm{ln}2/(\pi \sqrt{3})`$ for the Immirzi parameter $`\gamma `$, (which amounts to selecting one among a continuous choice of unitarily inequivalent quantum theories corresponding to the same classical theory), the standard result $`S=A/4l_P^2`$ is obtained for the Einstein Vacuum and Einstein-Maxwell case. It is worth to point out that this choice can be made only once, and that it is conceivable that say, the choice needed in the case of Einstein-Maxwell black holes might have been different than the choice needed for pure gravity black holes. Now, it would be of interest to check whether these results are also valid when non-Abelian gauge fields are present.
## Acknowledgments
We would like to thank A. Ashtekar, C. Beetle, S. Fairhurst and R. Wald for discussions and correspondence. We are also grateful to the Center for Gravitational Physics and Geometry for its hospitality. This work was in part supported by DGAPA-UNAM grant No IN121298, by CONACyT grants J32754-E and 32272-E, by a NSF-CONACyT collaborative grant, by NSF grants INT9722514, PHY95-14240 and by the Eberly research funds of Penn State. |
warning/0002/gr-qc0002085.html | ar5iv | text | # 1 Introduction
## 1 Introduction
At present there exists some interest to $`M`$-theory (see, for example, -). This theory is “supermembrane” analogue of superstring models in $`D=11`$. The low-energy limit of $`M`$-theory after a dimensional reduction leads to models governed by a Lagrangian containing a metric, fields of forms and scalar fields. These models contain a large variety of the so-called $`p`$-brane solutions (see - and references therein).
In it was shown that after the dimensional reduction on the manifold $`M_0\times M_1\times \mathrm{}\times M_n`$ when the composite $`p`$-brane ansatz for fields of forms is considered the problem is reduced to the gravitating self-interacting $`\sigma `$-model with certain constraints imposed. (For electric $`p`$-branes see also .) This representation may be considered as a tool for obtaining different solutions with intersecting $`p`$-branes. In the Majumdar-Papapetrou type solutions were obtained (for non-composite case see ). These solutions correspond to Ricci-flat factor-spaces $`(M_i,g^i)`$, ($`g^i`$ is metric on $`M_i`$) $`i=1,\mathrm{},n`$.They were also generalized to the case of Einstein internal spaces . (Earlier some special classes of these solutions were considered in .) The solutions take place, when certain (block-)orthogonality relations (on couplings parameters, dimensions of ”branes”, total dimension) are imposed. In this situation a class of cosmological and spherically-symmetric solutions was obtained . The solutions with a horizon were studied in details in .
Here we present a family of $`p`$-brane black hole solutions with (next to) arbitrary intersections (see Sect. 2). These black hole solutions are governed by moduli functions $`H_s=H_s(R)`$ obeying a set of second order non-linear differential equations with some boundary relations imposed. Some general features of these black holes (e.g. “single-time” and “no-hair” theorems) were predicted earlier in . We suggested a conjecture: the moduli functions $`H_s`$ are polynomials when intersection rules correspond to semisimple Lie algebras. The conjecture was confirmed by special black-hole “block -orthogonal” solutions considered earlier in . An analogue of this conjecture for extremal black holes was considered earlier in .
In Sect. 3 explicit formulas for the solution corresponding to the Lie algebra $`A_2`$ are obtained. These formulas are illustrated by two examples of $`A_2`$-dyon solutions: a dyon in $`D=11`$ supergravity (with $`M2`$ and $`M5`$ branes intersecting at a point) and Kaluza-Klein dyon.
## 2 $`p`$-brane black hole solutions
We consider a model governed by the action
$`S={\displaystyle d^Dx\sqrt{|g|}\left\{R[g]h_{\alpha \beta }g^{MN}_M\phi ^\alpha _N\phi ^\beta \underset{a\mathrm{}}{}\frac{\theta _a}{n_a!}\mathrm{exp}[2\lambda _a(\phi )](F^a)^2\right\}}`$ (2.1)
where $`g=g_{MN}(x)dx^Mdx^N`$ is a metric, $`\phi =(\phi ^\alpha )\text{R}^l`$ is a vector of scalar fields, $`(h_{\alpha \beta })`$ is a constant symmetric non-degenerate $`l\times l`$ matrix $`(l\text{N})`$, $`\theta _a=\pm 1`$,
$$F^a=dA^a=\frac{1}{n_a!}F_{M_1\mathrm{}M_{n_a}}^adz^{M_1}\mathrm{}dz^{M_{n_a}}$$
(2.2)
is a $`n_a`$-form ($`n_a1`$), $`\lambda _a`$ is a 1-form on $`\text{R}^l`$: $`\lambda _a(\phi )=\lambda _{\alpha a}\phi ^\alpha `$, $`a\mathrm{}`$, $`\alpha =1,\mathrm{},l`$. In (2.1) we denote $`|g|=|det(g_{MN})|`$,
$$(F^a)_g^2=F_{M_1\mathrm{}M_{n_a}}^aF_{N_1\mathrm{}N_{n_a}}^ag^{M_1N_1}\mathrm{}g^{M_{n_a}N_{n_a}},$$
(2.3)
$`a\mathrm{}`$. Here $`\mathrm{}`$ is some finite set. In the models with one time all $`\theta _a=1`$ when the signature of the metric is $`(1,+1,\mathrm{},+1)`$.
We obtained a new family of (black hole) solutions to field equations corresponding to the action (2.1) (for derivation of these solutions see ). These solutions are defined on the manifold
$$M=(R_0,+\mathrm{})\times (M_1=S^{d_1})\times (M_2=\text{R})\times \mathrm{}\times M_n,$$
(2.4)
and have the following form
$`g=\left({\displaystyle \underset{sS}{}}H_s^{2h_sd(I_s)/(D2)}\right)\{f^1dRdR+R^2d\mathrm{\Omega }_{d_1}^2`$ (2.5)
$`\left({\displaystyle \underset{sS}{}}H_s^{2h_s}\right)fdtdt+{\displaystyle \underset{i=3}{\overset{n}{}}}\left({\displaystyle \underset{sS}{}}H_s^{2h_s\delta _{iI_s}}\right)g^i\},`$
$`\mathrm{exp}(\phi ^\alpha )={\displaystyle \underset{sS}{}}H_s^{h_s\chi _s\lambda _{a_s}^\alpha },`$ (2.6)
$`F^a={\displaystyle \underset{sS}{}}\delta _{a_s}^a^s,`$ (2.7)
where $`f=12\mu /R^{\overline{d}}`$,
$$^s=\frac{Q_s}{R^{d_1}}\left(\underset{s^{}S}{}H_s^{}^{A_{ss^{}}}\right)dR\tau (I_s),$$
(2.8)
$`sS_e`$,
$$^s=Q_s\tau (\overline{I}_s),$$
(2.9)
$`sS_m`$. Here $`Q_s0`$ ($`sS`$) are charges, $`R_0>0`$, $`R_0^{\overline{d}}=2\mu >0`$, $`\overline{d}=d_11`$. In (2.5) $`g^i=g_{m_in_i}^i(y_i)dy_i^{m_i}dy_i^{n_i}`$ is a Ricci-flat metric on $`M_i`$, $`i=3,\mathrm{},n`$ and
$$\delta _{iI}=\underset{jI}{}\delta _{ij}$$
(2.10)
is the indicator of $`i`$ belonging to $`I`$: $`\delta _{iI}=1`$ for $`iI`$ and $`\delta _{iI}=0`$ otherwise. Let $`g^2=dtdt`$, and $`g^1=d\mathrm{\Omega }_{d_1}`$ be a canonical metric on unit sphere $`M_1=S^{d_1}`$,
The $`p`$-brane set $`S`$ is by definition
$`S=S_eS_m,S_v=_a\mathrm{}\{a\}\times \{v\}\times \mathrm{\Omega }_{a,v},`$ (2.11)
$`v=e,m`$ and $`\mathrm{\Omega }_{a,e},\mathrm{\Omega }_{a,m}\mathrm{\Omega }`$, where $`\mathrm{\Omega }=\mathrm{\Omega }(n)`$ is the set of all non-empty subsets of $`\{2,\mathrm{},n\}`$, i.e. all $`p`$-branes do not “live” in $`M_1`$.
Any $`p`$-brane index $`sS`$ has the form $`s=(a_s,v_s,I_s)`$, where $`a_s\mathrm{}`$, $`v_s=e,m`$ and $`I_s\mathrm{\Omega }_{a_s,v_s}`$. The sets $`S_e`$ and $`S_m`$ define electric and magnetic $`p`$-branes correspondingly. In (2.6) $`\chi _s=+1,1`$ for $`sS_e,S_m`$ respectively. All $`p`$-branes contain the time manifold $`M_2=\text{R}`$, i.e.
$`2I_s,sS.`$ (2.12)
All the manifolds $`M_i`$, $`i>2`$, are assumed to be oriented and connected and the volume $`d_i`$-forms
$$\tau _i\sqrt{|g^i(y_i)|}dy_i^1\mathrm{}dy_i^{d_i},$$
(2.13)
are well–defined for all $`i=1,\mathrm{},n`$. Here $`d_i=\mathrm{dim}M_i`$, $`i=1,\mathrm{},n`$, $`d_1>1`$, $`d_2=1`$, and for any $`I=\{i_1,\mathrm{},i_k\}\mathrm{\Omega }`$, $`i_1<\mathrm{}<i_k`$, we denote
$`\tau (I)\tau _{i_1}\mathrm{}\tau _{i_k},d(I)={\displaystyle \underset{iI}{}}d_i.`$ (2.14)
Forms $`^s`$ correspond to electric and magnetic $`p`$-branes for $`sS_e,S_m`$ respectively. In (2.9)
$$\overline{I}=\{1,\mathrm{},n\}I.$$
(2.15)
The parameters $`h_s`$ appearing in the solution satisfy the relations
$$h_s=K_s^1,K_s=B_{ss},$$
(2.16)
where
$`B_{ss^{}}=d(I_sI_s^{})+{\displaystyle \frac{d(I_s)d(I_s^{})}{2D}}+\chi _s\chi _s^{}\lambda _{\alpha a_s}\lambda _{\beta a_s^{}}h^{\alpha \beta },`$ (2.17)
$`s,s^{}S`$, with $`(h^{\alpha \beta })=(h_{\alpha \beta })^1`$ and $`D=1+_{i=1}^nd_i`$. Here we assume that
$$(𝐢)B_{ss}0,$$
(2.18)
for all $`sS`$, and
$$(\mathrm{𝐢𝐢})\mathrm{det}(B_{ss^{}})0,$$
(2.19)
i.e. the matrix $`(B_{ss^{}})`$ is a non-degenerate one.
Let consider the matrix
$$(A_{ss^{}})=\left(2B_{ss^{}}/B_{s^{}s^{}}\right).$$
(2.20)
Here some ordering in $`S`$ is assumed.
Functions $`H_s=H_s(z)>0`$, $`z=2\mu /R^{\overline{d}}(0,1)`$ obey the equations
$$\frac{d}{dz}\left(\frac{(1z)}{H_s}\frac{dH_s}{dz}\right)=B_s\underset{s^{}S}{}H_s^{}^{A_{ss^{}}},$$
(2.21)
equipped with the boundary conditions
$`H_s(10)=H_{s0}(0,+\mathrm{}),`$ (2.22)
$`H_s(+0)=1,`$ (2.23)
$`sS`$. Here $`B_s=K_s\epsilon _sQ_s^2/(2\overline{d}\mu )^2`$ and
$$\epsilon _s=(\epsilon [g])^{(1\chi _s)/2}\epsilon (I_s)\theta _{a_s},$$
(2.24)
$`sS`$, $`\epsilon [g]signdet(g_{MN})`$. More explicitly (2.24) reads: $`\epsilon _s=\epsilon (I_s)\theta _{a_s}`$ for $`v_s=e`$ and $`\epsilon _s=\epsilon [g]\epsilon (I_s)\theta _{a_s}`$, for $`v_s=m`$.
Equations (2.22) are equivalent to Toda-type equations. First boundary condition guarantees the existence of a regular horizon at $`R^{\overline{d}}=2\mu `$. Second condition (2.23) ensures an asymptotical (for $`R+\mathrm{}`$) flatness of the $`(2+d_1)`$-dimensional section of the metric.
Due to (2.8) and (2.9), the dimension of $`p`$-brane worldsheet $`d(I_s)`$ is defined by
$`d(I_s)=n_{a_s}1,d(I_s)=Dn_{a_s}1,`$ (2.25)
for $`sS_e,S_m`$ respectively. For a $`p`$-brane: $`p=p_s=d(I_s)1`$.
The solutions are valid if the following restriction on the sets $`\mathrm{\Omega }_{a,v}`$ is imposed. (This restriction guarantees the block-diagonal structure of the stress-energy tensor, like for the metric, and the existence of $`\sigma `$-model representation , see also ). We denote $`w_1\{i|i\{2,\mathrm{},n\},d_i=1\}`$, and $`n_1=|w_1|`$ (i.e. $`n_1`$ is the number of 1-dimensional spaces among $`M_i`$, $`i=2,\mathrm{},n`$). It follows from (2.12) that $`2w_1`$.
Restriction. Let 1a) $`n_11`$ or 1b) $`n_12`$ and for any $`a\mathrm{}`$, $`v\{e,m\}`$, $`i,jw_1`$, $`ij`$, there are no $`I,J\mathrm{\Omega }_{a,v}`$ such that $`iI`$, $`jJ`$ and $`I\{i\}=J\{j\}`$.
These restriction is satisfied in the non-composite case : $`|\mathrm{\Omega }_{a,e}|+|\mathrm{\Omega }_{a,m}|=1`$, (i.e when there are no two $`p`$-branes with the same color index $`a`$, $`a\mathrm{}`$.) The restriction forbids certain intersections of two $`p`$-branes with the same color index for $`n_12`$.
The Hawking temperature corresponding to the solution is (see also for orthogonal case) found to be
$$T_H=\frac{\overline{d}}{4\pi (2\mu )^{1/\overline{d}}}\underset{sS}{}H_{s0}^{h_s},$$
(2.26)
where $`H_{s0}`$ are defined in (2.22)
This solution describes a set of charged (by forms) overlapping $`p`$-branes ($`p_s=d(I_s)1`$, $`sS`$) “living” on submanifolds of $`M_2\times \mathrm{}\times M_n`$.
## 3 Examples.
### 3.1 Orthogonal and block-orthogonal solutions
There exist solutions to eqs. (2.21)-(2.23) of polynomial type. The simplest example occurs in orthogonal case , when
$$B_{ss^{}}=0,$$
(3.1)
for $`ss^{}`$, $`s,s^{}S`$. In this case $`(A_{ss^{}})=\mathrm{diag}(2,\mathrm{},2)`$ is a Cartan matrix for semisimple Lie algebra $`A_1\mathrm{}A_1`$ and
$$H_s(z)=1+P_sz,$$
(3.2)
with $`P_s0`$, satisfying
$$P_s(P_s+1)=B_s,$$
(3.3)
$`sS`$. Relation (2.22) implies $`P_s>1`$. For $`B_s<0`$ parameters $`P_s>0`$ are uniquely defined.
In this solution was generalized to a block orthogonal case:
$`S=S_1\mathrm{}S_k,S_iS_j=\mathrm{},ij,`$ (3.4)
$`S_i\mathrm{}`$, i.e. the set $`S`$ is a union of $`k`$ non-intersecting (non-empty) subsets $`S_1,\mathrm{},S_k`$, and relation (3.1) should be satisfied for all $`sS_i`$, $`s^{}S_j`$, $`ij`$; $`i,j=1,\mathrm{},k`$. In this case (3.2) is modified as follows
$$H_s(z)=(1+P_sz)^{b_0^s},$$
(3.5)
where
$$b_0^s=2\underset{s^{}S}{}A^{ss^{}},$$
(3.6)
$`(A^{ss^{}})=(A_{ss^{}})^1`$ and parameters $`P_s`$ are coinciding inside blocks, i.e. $`P_s=P_s^{}`$ for $`s,s^{}S_i`$, $`i=1,\mathrm{},k`$. Parameters $`P_s0`$ satisfy the relations
$$P_s(P_s+2\mu )=\overline{B}_s/b_0^s,$$
$`b_0^s0`$, and parameters $`\overline{B}_s/b_0^s`$ are also coinciding inside blocks, i.e. $`\overline{B}_s/b_0^s=\overline{B}_s^{}/b_0^s^{}`$ for $`s,s^{}S_i`$, $`i=1,\mathrm{},k`$.
Let $`(A_{ss^{}})`$ be a Cartan matrix for a finite-dimensional semisimple Lie algebra $`𝒢`$. In this case all powers in (3.6) are natural numbers (coinciding with the components of twice the dual Weyl vector in the basis of simple coroots, see ) and hence, all functions $`H_s`$ are polynomials, $`sS`$.
Conjecture. Let $`(A_{ss^{}})`$ be a Cartan matrix for a semisimple finite-dimensional Lie algebra $`𝒢`$. Then the solution to eqs. (2.21)-(2.23) (if exists) is a polynomial
$$H_s(z)=1+\underset{k=1}{\overset{n_s}{}}P_s^{(k)}z^k,$$
(3.7)
where $`P_s^{(k)}`$ are constants, $`k=1,\mathrm{},n_s`$, integers $`n_s=b_0^s`$ are defined in (3.6) and $`P_s^{(n_s)}0`$, $`sS`$.
For certain series of simple finite-dimensional Lie algebras this conjecture will be proved in a separate publication. In the extremal case $`\mu =+0`$ an a analogue of this conjecture was suggested previously in .
### 3.2 Solutions for $`A_2`$ algebra
Here we consider some examples of solutions related to the Lie algebra $`A_2=sl(3)`$. In this case the Cartan matrix reads
$$\left(A_{ss^{}}\right)=\left(\begin{array}{cccccc}2& 1& & & & \\ 1& 2& & & & \end{array}\right)$$
(3.8)
According to the results of previous section we seek the solutions to eqs. (2.21)-(2.23) in the following form (see (3.7); here $`n_1=n_2=2`$):
$$H_s=1+P_sz+P_s^{(2)}z^2,$$
(3.9)
where $`P_s=P_s^{(1)}`$ and $`P_s^{(2)}0`$ are constants, $`s=1,2`$.
The substitution of (3.9) into equations (2.21) and decomposition in powers of $`z`$ lead us to the relations
$`P_s(P_s+1)+2P_s^{(2)}=B_s,`$ (3.10)
$`2P_s^{(2)}(P_s+2)=P_{s+1}B_s,`$ (3.11)
$`2P_s^{(2)}({\displaystyle \frac{1}{2}}P_s+P_s^{(2)})=P_{s+1}^{(2)}B_s,`$ (3.12)
corresponding to powers $`z^0,z^1,z^2`$ respectively, $`s=1,2`$. Here we denote $`s+1=2,1`$ for $`s=1,2`$ respectively. For $`P_1+P_2+20`$ the solutions to eqs. (3.10)-(3.12) read
$`P_s^{(2)}={\displaystyle \frac{P_sP_{s+1}(P_s+1)}{2(P_1+P_2+2)}},`$ (3.13)
$`B_s={\displaystyle \frac{P_s(P_s+1)(P_s+2)}{P_1+P_2+2}},`$ (3.14)
$`s=1,2`$. For $`P_1+P_2+2=0`$ there exist also a special solution with
$`P_1=P_2=1,2P_s^{(2)}=B_s>0,B_1+B_2=1.`$ (3.15)
Thus, in the $`A_2`$-case the solution is described by relations (2.5)-(2.9) with $`S=\{s_1,s_2\}`$, intersection rules following from (2.17), (2.20) and (3.8)
$`d(I_{s_1}I_{s_2})={\displaystyle \frac{d(I_{s_1})d(I_{s_2})}{D2}}\chi _{s_1}\chi _{s_2}\lambda _{a_{s_1}}\lambda _{a_{s_2}}{\displaystyle \frac{1}{2}}K,`$ (3.16)
$`d(I_{s_i}){\displaystyle \frac{(d(I_{s_i}))^2}{D2}}+\lambda _{a_{s_i}}\lambda _{a_{s_i}}=K,`$ (3.17)
where $`K=K_{s_i}0`$, and functions $`H_{s_i}=H_i`$ are defined by relations (3.9) and (3.13)-(3.15) with $`z=2\mu R^{\overline{d}}`$, $`i=1,2`$. Here $`\lambda \lambda ^{^{}}=\lambda _\alpha \lambda _\beta ^{^{}}h^{\alpha \beta }`$.
### 3.3 $`A_2`$-dyon in $`D=11`$ supergravity
Consider the “truncated” bosonic sector of $`D=11`$ supergravity (“truncated” means without Chern-Simons term). The action (2.1) in this case reads
$`S_{tr}={\displaystyle _M}d^{11}z\sqrt{|g|}\left\{R[g]{\displaystyle \frac{1}{4!}}F^2\right\}.`$ (3.18)
where $`\mathrm{rank}F=4`$. In this particular case, we consider a dyonic black-hole solutions with electric $`2`$-brane and magnetic $`5`$-brane defined on the manifold
$$M=(2\mu ,+\mathrm{})\times (M_1=S^2)\times (M_2=\text{R})\times M_3\times M_4,$$
(3.19)
where $`dimM_3=2`$ and $`dimM_4=5`$.
The solution reads,
$`g=H_1^{1/3}H_2^{2/3}\left\{f^1dRdR+R^2d\mathrm{\Omega }_2^2H_1^1H_2^1fdtdt+H_1^1g^3+H_2^1g^4\right\},`$ (3.20)
$`F={\displaystyle \frac{Q_1}{R^2}}H_1^2H_2dRdt\tau _3+Q_2\tau _1\tau _3,`$ (3.21)
where $`f=12\mu /R`$, metrics $`g^2`$ and $`g^3`$ are Ricci-flat metrics of Euclidean signature, and functions $`H_s`$ are defined by relations (3.9), (3.13) and (3.14) with $`z=2\mu R^1`$, $`B_s=2Q_s^2/(2\mu )^2`$, $`s=1,2`$; where $`\tau _1`$ is volume form on $`S^2`$.
The solution describes $`A_2`$-dyon consisting of electric $`2`$-brane with worldsheet isomorphic to $`(M_2=\text{R})\times M_3`$ and magnetic $`5`$-brane with worldsheet isomorphic to $`(M_2=\text{R})\times M_4`$. The “branes” are intersecting on the time manifold $`M_2=\text{R}`$. Here $`K_s=(U^s,U^s)=2`$, $`\epsilon _s=1`$ for all $`sS`$. The $`A_2`$ intersection rule reads (see (3.16))
$$25=1$$
(3.22)
Here and in what follows $`(p_1p_2=d)(d(I)=p_1+1,d(J)=p_2+1,d(IJ)=d)`$.
The solution (3.20), (3.21) satisfies not only equations of motion for the truncated model, but also the equations of motion for $`D=11`$ supergravity with the bosonic sector action
$`S=S_{tr}+c{\displaystyle _M}AFF`$ (3.23)
($`c=\mathrm{const}`$, $`F=dA`$), since the only modification related to “Maxwells” equations
$`dF=\mathrm{const}FF,`$ (3.24)
is trivial due to $`FF=0`$ (since $`\tau _i\tau _i=0`$).
This solution in a special case $`H_1=H_2=H^2`$ ($`P_1=P_2`$, $`Q_1^2=Q_2^2`$) was considered in . The 4-dimensional section of the metric (3.20) in this special case coincides with the Reissner-Nordström metric. For the extremal case, $`\mu +0`$, and multi-black-hole generalization see also .
### 3.4 $`A_2`$-dyon in Kaluza-Klein model
Let us consider $`4`$-dimensional model
$$S=_Md^4z\sqrt{|g|}\left\{R[g]g^{\mu \nu }_\mu \phi _\nu \phi \frac{1}{2!}\mathrm{exp}[2\lambda \phi ]F^2\right\}$$
(3.25)
with scalar field $`\phi `$, two-form $`F=dA`$ and $`\lambda =\sqrt{3/2}`$. This model originates after Kaluza-Klein (KK) reduction of $`5`$-dimensional gravity. The 5-dimensional metric in this case reads
$$g^{(5)}=\varphi g_{\mu \nu }dx^\mu dx^\nu +\varphi ^2(dy+𝒜)(dy+𝒜),$$
(3.26)
where $`𝒜=\sqrt{2}A=\sqrt{2}A_\mu dx^\mu `$ and $`\varphi =\mathrm{exp}(2\phi /\sqrt{6})`$.
We consider a dyonic black-hole solution carrying electric charge $`Q_1`$ and magnetic charge $`Q_2`$, defined on the manifold $`M=(2\mu ,+\mathrm{})\times (M_1=S^2)\times (M_2=\text{R})`$. This solution reads
$`g=\left(H_1H_2\right)^{1/2}\left\{{\displaystyle \frac{dRdR}{12\mu /R}}+R^2d\mathrm{\Omega }_2^2H_1^1H_2^1\left(1{\displaystyle \frac{2\mu }{R}}\right)dtdt\right\},`$ (3.27)
$`\mathrm{exp}(\phi )=H_1^{\lambda /2}H_2^{\lambda /2},`$ (3.28)
$`F=dA={\displaystyle \frac{Q_1}{R^2}}H_1^2H_2dRdt+Q_2\tau _1,`$ (3.29)
where functions $`H_s`$ are defined by relations (3.9), (3.13) and (3.14) with $`z=R^1`$, $`B_s=2Q_s^2/(2\mu )^2`$, $`s=1,2`$; where $`\tau _1`$ is volume form on $`S^2`$.
For 5-metric we obtain from (3.26)-(3.28)
$`g^{(5)}=H_2\left\{{\displaystyle \frac{dRdR}{12\mu /R}}+R^2d\mathrm{\Omega }_2^2H_1^1H_2^1\left(1{\displaystyle \frac{2\mu }{R}}\right)dtdt\right\}`$ (3.30)
$`+H_1H_2^1(dy+𝒜)(dy+𝒜),`$
$`d𝒜=\sqrt{2}F`$.
For $`Q_20`$ we get the black hole version of Dobiash-Maison solution from and for $`Q_10`$ we are led to the black hole version of Gross-Perry-Sorkin monopole solution from , (see ). The solution coincides with Gibbons-Wiltshire dyon solution . Our notations are related to those from ref. , as following : $`H_1R^2=B`$, $`H_2R^2=A`$, $`R^22\mu R=\mathrm{\Delta }`$, $`Q_1=\sqrt{2}q`$, $`Q_2=\sqrt{2}p`$, $`R\mu =rm`$, $`\mu ^2=m^2+d^2p^2q^2`$, $`(P_2P_1)/2(P_2+1)=d/(d\sqrt{3}m)`$). (For general spherically symmetric configurations see also ref. .)
We note that, quite recently, in the KK dyon solution was used for constructing the dyon solution in $`D=11`$ supergravity (3.20)-(3.21) for flat $`g^3`$ and $`g^4`$ and its rotating version. This dyon solution differs from $`M2M5`$ dyon from .
## 4 Conclusions
Thus here we presented a family of black hole solutions with intersecting $`p`$-branes with next to arbitrary intersection rules (see (2.18), (2.19) and Restriction). The metric of solutions contains $`n1`$ Ricci-flat “internal” space metrics. The solutions are defined up to a set of functions $`H_s`$ obeying a set of equations (equivalent to Toda-type equations) with certain boundary conditions imposed. Using a conjecture on polynomial structure of $`H_s`$ for intersections related to semisimple Lie algebras, we obtained explicit relations for the solutions in the $`A_2`$-case and considered two examples of $`A_2`$-dyon solutions: one in $`D=11`$ supergravity (with M2 and M5 branes intersecting at a point) and another in $`5`$-dimensional Kaluza-Klein theory.
Acknowledgments
This work was supported in part by the Russian Ministry for Science and Technology, Russian Foundation for Basic Research, and project SEE. |
warning/0002/math0002058.html | ar5iv | text | # Untitled Document
Espace des modules
des faisceaux de rang 2 semi-stables
de classes de Chern $`c_1=0`$, $`c_2=2`$ et $`c_3=0`$
sur la cubique de $`^4`$
Stéphane Druel
1. Introduction
(1.1) Soient $`X^4`$ une hypersurface cubique lisse et $`\mathrm{}X`$ une droite de $`^4`$. Soit $`X_{\mathrm{}}`$ la variété obtenue en éclatant $`\mathrm{}`$ dans $`X`$. La projection le long de $`\mathrm{}`$ induit un morphisme $`X_{\mathrm{}}\stackrel{𝑝}{}^2`$ dont les fibres sont les coniques qui sont coplanaires avec $`\mathrm{}`$. Lorsque la droite $`\mathrm{}`$ est générique les fibres de $`p`$ sont lisses ou réunion de deux droites distinctes. Le lieu de dégénérescence de $`p`$ est alors une courbe plane lisse et connexe $`C_0`$ de degré 5. Soit $`C`$ la variété des droites contenues dans $`X`$ et incidentes à $`\mathrm{}`$. Le morphisme $`CC_0`$ est un revêtement étale double connexe. Soit $`i`$ l’involution échangeant les deux feuillets dudit revêtement et notons encore $`i`$ l’automorphisme induit sur la jacobienne $`JC`$. La variété de Prym associée au revêtement $`(C,C_0)`$ est alors $`P=(Idi)JC`$. C’est une variété abélienne principalement polarisée de dimension 5. Soient $`A_1(X)`$ le groupe des 1-cycles algébriques modulo l’équivalence rationnelle et $`AA_1(X)`$ le sous-groupe des cycles algébriquement équivalents à zéro. L’application qui à $`tC`$ associe la classe de la droite correspondante $`z_tX_{\mathrm{}}`$ dans $`A`$ induit un isomorphisme de groupes $`PA`$. On démontre que pour toute variété lisse $`T`$ de dimension pure $`n1`$ et tout $`n+1`$-cycle $`z`$ sur $`X\times T`$ l’application d’Abel-Jacobi qui à $`tT`$ associe la classe du cycle $`z_tz_{t_0}`$ dans $`P`$, où $`t_0T`$ est fixé, est algébrique (\[Mu\]). La jacobienne intermédiaire de $`X`$ est définie par :
$$J(X)=(H^{2,1}(X))^{}/\alpha (H_3(X,))$$
$`\alpha `$ est l’application donnée par intégration sur les cycles. C’est une variété abélienne principalement polarisée de dimension 5 isomorphe à la variété de Prym. Via cet isomorphisme, l’image du cycle $`z_tz_{t_0}`$ par l’application d’Abel-Jacobi est la forme linéaire donnée par intégration sur le cycle $`\mathrm{\Gamma }`$ modulo le groupe $`\alpha (H_3(X,))`$$`\mathrm{\Gamma }=z_tz_{t_0}`$. On démontre enfin que l’application d’Abel-Jacobi induit un plongement de la surface de Fano de $`X`$ dans $`J(X)`$.
(1.2) Soient $`(X,𝒪_X(1))`$ une variété polarisée de dimension $`n1`$ et $`E`$ un faisceau cohérent sur $`X`$ de rang $`r`$. La pente $`\mu (E)`$ de $`E`$ est définie par la formule :
$$\mu (E)=\frac{c_1(E)c_1(𝒪_X(1))^{n1}}{r}$$
Le faisceau $`E`$ est dit $`\mu `$-semi-stable (resp. semi-stable) s’il est sans torsion et si pour tout sous-faisceau $`LE`$ de rang $`0<r^{}<r`$ on a $`\mu (L)\mu (E)`$ (resp. $`\frac{\chi (L(n))}{r^{}}\frac{\chi (E(n))}{r}`$ pour $`n0`$). Il est dit $`\mu `$-stable (resp. stable) s’il est sans torsion et si pour tout sous-faisceau $`LE`$ de rang $`0<r^{}<r`$ on a $`\mu (L)<\mu (E)`$ (resp. $`\frac{\chi (L(n))}{r^{}}<\frac{\chi (E(n))}{r}`$ pour $`n0`$). On a les implications suivantes :
$$\mu \text{-stable}\text{stable}\text{semi-stable}\mu \text{-semi-stable}$$
Supposons enfin $`\text{Pic}(X)`$ et soit $`F`$ un faisceau réflexif de rang 2 sur $`X`$, de première classe de Chern $`c_1(F)=0`$ ou $`c_1(F)=1`$. Alors $`F`$ est stable si et seulement si $`h^0(F)=0`$ et si $`c_1(F)=0`$ alors $`F`$ est semi-stable si et seulement si $`h^0(F(1))=0`$ (\[H2\] lemme 3.1).
(1.3) Soient $`X^4`$ une hypersurface cubique lisse et $`𝒪_X(1)`$ le générateur très ample de $`\text{Pic}(X)`$. Les $``$-modules $`H^2(X,)`$, $`H^4(X,)`$ et $`H^6(X,)`$ sont libres de rang 1. On identifie ainsi les classes de Chern d’un faisceau cohérent sur $`X`$ à des entiers relatifs. Nous étudions ici l’espace des modules des faisceaux semi-stables de rang 2 sur $`X`$. Nous démontrons le :
Théorème 1.4.$``$Soient $`X^4`$ une hypersurface cubique lisse et $`B`$ la surface de Fano de $`X`$. Alors l’espace des modules $`M_X`$ des faisceaux semi-stables de rang 2 sur $`X`$ de classe de Chern $`c_1=0`$, $`c_2=2`$ et $`c_3=0`$ est isomorphe à l’éclatement d’un translaté de la surface $`B`$ dans la jacobienne intermédiaire $`J(X)`$.
(1.5) Soit $`X`$ une variété projective lisse de dimension au moins 2 et $`E`$ un fibré vectoriel de rang 2 sur $`X`$. S’il existe une section globale dont le lieu des zéros $`Y`$ est de codimension pure 2 alors on a une suite exacte (\[H1\]) :
$$0𝒪_XEI_Y\text{det}(E)0$$
(1.6) Fibrés de rang 2 et construction de Serre.$``$ Supposons $`X`$ de dimension au moins 3. Soit $`L`$ un fibré inversible sur $`X`$ tel que $`h^1(L^1)=0`$ et $`h^2(L^2)=0`$ et soit $`YX`$ un sous-schéma fermé de codimension pure 2. On a un isomorphisme $`\text{Ext}_X^1(I_YL,𝒪_X)=H^0(𝒪_Y).`$ Le sous-schéma $`Y`$ est le lieu des zéros d’une section d’un fibré $`E`$ de rang 2 sur $`X`$ de déterminant $`L`$ si et seulement si $`Y`$ est localement intersection complète et $`\omega _Y=(\omega _XL)_{|Y}`$.
(1.7) Soit $`X^4`$ une hypersurface cubique lisse. Nous montrons que les fibrés vectoriels stables sont associés aux quintiques elliptiques normales tracées sur $`X`$ par la construction de Serre (2.4), au moyen du :
(1.8) Critère de Mumford-Castelnuovo.$``$Soit $`F`$ un faisceau cohérent sur une variété projective $`X`$ tel que $`h^i(F(i))=0`$ pour $`i1`$. Alors $`h^i(F(k))=0`$ pour $`i1`$ et $`ki`$ et $`F`$ est engendré par ses sections globales (\[Mum\] lect. 14).
(1.9) Soit $`X^4`$ une hypersurface cubique lisse. Nous montrons que les faisceaux stables non localement libres sont paramétrés par les coniques lisses tracées sur $`X`$ et que les faisceaux strictement semi-stables sont paramétrés par les couples de droites de $`X`$ (3.5). Nous montrons enfin que la seconde classe de Chern définit un morphisme vers la jacobienne intermédiaire $`J(X)`$. Ce morphisme est birationnel (\[I-M\]) et identifie $`M_X`$ à l’éclatement d’une surface lisse dans $`J(X)`$ (4.8).
(1.10) Soient $`F`$ un faisceau cohérent sur un schéma $`X`$ et $`YX`$ un sous-schéma fermé. La restriction de $`F`$ à $`Y`$ sera notée $`F_Y`$.
Remerciements.$``$Je tiens à exprimer toute ma gratitude à Arnaud Beauville pour m’avoir soumis ce problème et pour l’aide qu’il m’a apportée. Je remercie également le $`referee`$ pour ces nombreuses remarques et pour avoir relevé une erreur dans la preuve de la proposition 3.1.
2. Fibrés de rang 2 stables sur la cubique de $`^4`$
Lemme 2.1.$``$Soient $`X^N`$ une variété de dimension $`n2`$ et $`E`$ un fibré de rang 2 $`\mu `$-semi-stable de première classe de Chern $`c_1(E)=0`$. Si $`h^0(E)0`$ alors le lieu des zéros d’une section globale non nulle est de codimension pure 2 ou bien ladite section ne s’annule pas et $`c_2(E)=0`$.
Démonstration.$``$Le fibré $`E`$ est de rang 2 et toute section non triviale s’annule donc en codimension au plus 2 ou bien ne s’annule pas. S’il existe une section partout non nulle alors $`E`$ est extension de $`𝒪_X`$ par $`L`$ avec $`c_1(L)=0`$ et $`c_2(E)=0`$. Supposons qu’une section de $`E`$ s’annule en codimension 1 et soit $`D`$ la partie de codimension 1 du lieu des zéros de ladite section. On a ainsi $`h^0(E(D))0`$ et $`\mu (𝒪_X(D))0`$ puisque $`E`$ est semi-stable. Or $`D`$ est effectif et on a donc $`D=0`$. $`\mathrm{}`$
Lemme 2.2.$``$Soient $`S^3`$ une surface cubique lisse et $`E`$ un fibré de rang 2 $`\mu `$-semi-stable de classes de Chern $`c_1(E)=0`$ et $`c_2(E)=2`$. Si $`h^0(E)=0`$ alors $`h^1(E(n))=0`$ pour $`n`$ et $`h^2(E(n))=0`$ pour $`n1`$. Si $`h^0(E)0`$ alors $`h^0(E)=1`$, $`h^1(E(n))=0`$ pour $`n2`$ et $`n1`$, $`h^1(E(1))=h^1(E)=1`$ et $`h^2(E(n))=0`$ pour $`n0`$.
Démonstration.$``$Supposons $`h^0(E)=0`$.$``$On a $`h^1(E)=h^0(E)=0`$ puisque $`h^2(E)=h^0(E(1))=0`$ et $`\chi (E)=0`$. On a enfin $`h^2(E(1))=h^0(E)=0`$. Finalement $`h^i(E(1i))=0`$ pour $`i1`$ et le lemme est une conséquence de (1.8).
Supposons $`h^0(E)0`$.$``$Le fibré $`E`$ est semi-stable et le lieu des zéros d’une section globale non nulle est donc de codimension pure 2 (2.1). On a donc une suite exacte (1.5) :
$$0𝒪_SEI_Z0$$
$`Z`$ est un sous-schéma fermé de dimension 0 et de longueur 2. On en déduit en particulier $`h^0(E)=1`$ et $`h^1(E)=1`$ puisque $`h^2(E)=h^0(E(1))=0`$ et $`\chi (E)=0`$. L’application naturelle $`H^0(𝒪_S(1))H^0(𝒪_Z(1))`$ est surjective puisque $`\mathrm{}(Z)=2`$ et on a donc $`h^1(I_Z(1))=0`$. On en déduit $`h^1(E(1))=0`$. Finalement $`h^i(E(2i))=0`$ pour $`i1`$ et le lemme est une conséquence de (1.8). $`\mathrm{}`$
Lemme 2.3.$``$Soient $`S^3`$ une surface cubique lisse et $`E`$ un fibré de rang 2 $`\mu `$-semi-stable de classes de Chern $`c_1(E)=0`$ et $`c_2(E)=1`$. Si $`h^0(E)0`$ alors $`h^0(E)=1`$, $`h^1(E(n))=0`$ pour $`n`$ et $`h^2(E(n))=0`$ pour $`n0`$.
Démonstration.$``$Le fibré $`E`$ est semi-stable et le lieu des zéros d’une section globale non nulle est donc de codimension pure 2 (2.1). On a une suite exacte (1.5) :
$$0𝒪_SEI_Z0$$
$`Z`$ est un point de $`S`$. On en déduit $`h^0(E)=1`$. La suite exacte :
$$0𝒪_S(n)E(n)I_Z(n)0$$
donne $`h^1(E(n))=0`$ pour $`n0`$ puisque $`h^1(𝒪_S(n))=0`$ et $`h^1(I_Z(n))=0`$ pour $`n0`$. On en déduit $`h^1(E(n))=h^1(E(n1))=0`$ pour $`n<0`$. On a enfin $`h^2(E(n))=h^0(E(1n))=0`$ pour $`n0`$. $`\mathrm{}`$
Théorème 2.4.$``$Soient $`X^4`$ une cubique lisse et $`E`$ un fibré de rang 2 stable de classes de Chern $`c_1(E)=0`$ et $`c_2(E)=2`$. Alors $`E(1)`$ est engendré par ses sections globales.
Démonstration.$``$Soit $`S|𝒪_X(1)|`$ une section hyperplane générique de $`X`$ tel que le fibré $`E_S`$ soit $`\mu `$-semi-stable relativement à la polarisation $`𝒪_S(1)`$ (\[M\] thm. 3.1).
Supposons $`h^0(E_S)=0`$.$``$Il suffit de prouver $`h^i(E(1i))=0`$ pour $`i1`$ (1.8). Considérons la suite exacte :
$$0E(n1)E(n)E_S(n)0$$
On a $`h^1(E(n))h^1(E(n1))`$ puisque $`h^1(E_S(n))=0`$ pour $`n`$ (2.2). On en déduit $`h^1(E(n))=0`$ pour $`n`$ puisque $`h^1(E(n))=0`$ pour $`n0`$ puis $`h^2(E(n))=0`$ pour $`n`$. On a enfin $`h^3(E(2))=h^0(E)=0`$.
Supposons $`h^0(E_S)0`$ et montrons que nous aboutissons à une contradiction.$``$Le fibré $`E(2)`$ est alors engendré par ses sections globales. Il suffit en effet de prouver $`h^i(E(2i))=0`$ pour $`i1`$ (1.8). Considérons à nouveau la suite exacte :
$$0E(n1)E(n)E_S(n)0$$
On a $`h^1(E(n))h^1(E(n1))`$ pour $`n2`$ puisque $`h^1(E_S(n))=0`$ pour $`n2`$ (2.2). On en déduit $`h^1(E(n))=0`$ pour $`n2`$ puisque $`h^1(E(n))=0`$ pour $`n0`$. Calculons $`h^1(E(1))`$. On a $`h^2(E)=h^1(E(2))=0`$ et $`h^3(E)=h^0(E(2))=0`$. Puisque $`\chi (E)=0`$ on a donc $`h^1(E)=0`$ et la suite exacte :
$$0EE(1)E_S(1)0$$
entraîne $`h^1(E(1))=h^1(E_S(1))=0`$ (2.2). On a enfin $`h^3(E(1))=h^0(E(1))=0`$. Le fibré $`E(2)`$ est donc engendré par ses sections globales.
Si l’une des sections du fibré $`E(2)`$ est partout non nulle alors $`E(2)`$ est isomorphe au fibré $`𝒪_X(2)𝒪_X(2)`$ et $`c_2(E)=12`$ ce qui est absurde. On a donc une suite exacte (1.5) :
$$0𝒪_X(4)E(2)I_C0$$
$`CX`$ est une courbe lisse de degré $`c_2(E(2))=14`$. On a $`h^1(I_C)=0`$ et la courbe $`C`$ est donc connexe. On a $`\omega _C=𝒪_C(2)`$ (1.6) et $`g(C)=15`$. Enfin, la courbe $`C`$ est non dégénérée puisque le fibré $`E`$ est stable. Calculons $`h^0(𝒪_C(1))`$. La suite exacte :
$$0𝒪_X(3)E(1)I_C(1)0$$
entraîne l’égalité $`h^1(I_C(1))=h^1(E(1))`$ puisque $`h^1(𝒪_X(3))=0`$ et $`h^2(𝒪_X(3))=0`$. La suite exacte :
$$0E(2)E(1)E_S(1)0$$
donne $`h^1(E(1))=h^1(E_S(1))=1`$ (2.2) puisque $`h^1(E(2))=0`$ et $`h^2(E(2))=h^1(E)=0`$. On a donc $`h^1(I_C(1))=1`$. On déduit de la suite exacte :
$$0I_C(1)𝒪_X(1)𝒪_C(1)0$$
que $`h^0(𝒪_C(1))=6`$ puisque $`h^0(I_C(1))=0`$ et $`h^1(𝒪_X(1))=0`$. La courbe $`C`$ est donc la projection dans $`^4`$ d’une courbe de Castelnuovo de $`^5`$ et le lemme 2.5 fournit la contradiction cherchée. $`\mathrm{}`$
Lemme 2.5.$``$Soit $`C^5`$ une courbe non dégénérée de genre 15 et de degré 14 (courbe de Castelnuovo). Soit $`O^5`$ ($`OC`$) tel que la projection à partir de $`O`$ induise un plongement de $`C`$ dans $`^4`$. L’image de $`C`$ dans $`^4`$ n’est alors contenue dans aucune cubique lisse.
Démonstration.$``$La courbe $`C`$ est contenue dans une surface irréductible $`S^5`$ de degré 4. Ladite surface $`S`$ et la courbe $`C`$ sont (\[A-C-G-H\]) :
ou bien
$``$ la surface de Veronese et $`C`$ est l’image d’une courbe plane de degré 7 par le plongement de Veronese,
ou bien
$``$ l’image de $`S_{2k}=_^1(𝒪_^1𝒪_^1(2k))`$, $`k\{0,1,2\}`$, par le morphisme $`\phi _k`$ associé au système linéaire $`|C_0+(k+2)f|`$ et $`C|4C_0+(4k+6)f|`$, où $`C_0`$ est la section associée au fibré naturel $`𝒪_{S_{2k}}(1)`$ $`(C_0^2=2k)`$ et $`f`$ une génératrice de la surface réglée $`S_{2k}`$. Pour $`k\{0,1\}`$ le morphisme $`\phi _k`$ est un plongement fermé, $`S=\phi _2(S_4)`$ est un cône au-dessus d’une courbe rationnelle lisse de degré 4 et le morphisme $`\phi _2`$ s’identifie à l’éclatement de $`\phi _2(S_4)`$ en son sommet.
Notons $`\pi `$ la projection considérée et $`\pi (S)`$ l’image de $`S`$ par l’application rationnelle $`\pi `$. Si $`\pi (S)`$ est de dimension 1 alors $`S`$ est un cône au dessus de $`C`$ isomorphe à $`\phi _2(S_4)`$ ce qui absurde puisque $`g(C)1`$. La variété $`\pi (S)`$ est donc de dimension 2. Si $`S`$ est un cône alors son sommet et le point de projection sont donc distincts.
Supposons la courbe $`C^4`$ contenue dans une cubique lisse $`X`$ et notons $`\overline{X}^5`$ le cône de sommet $`O`$ et de base $`X`$.
Supposons que la cubique $`\overline{X}`$ ne contienne pas la surface $`S`$. L’hypersurface $`\overline{X}`$ découpe alors sur $`S`$ une courbe de degré 12 et ne peut donc pas contenir la courbe $`C`$. La cubique $`\overline{X}`$ contient donc la surface S. On en déduit en particulier que $`\pi (S)X`$.
Supposons $`OS`$. Si $`S`$ est l’une des deux surfaces $`\phi _k(S_{2k})`$ avec $`k\{0,1,2\}`$ alors la génératrice $`f`$ passant par $`O`$ est contractée par $`\pi `$. Or $`C.f=4`$ et $`\pi `$ ne peut donc pas induire un plongement de $`C`$ dans $`^4`$. Si $`S`$ est la surface de Veronese alors l’application rationnelle $`^2^4`$ obtenue est définie par le système linéaire des coniques passant par un point. Ce système linéaire induit un plongement de la surface de Hirzebruch $`𝔽_1`$ dans $`^4`$ dont l’image est une surface de degré 3. Or $`𝔽_1X`$ et ladite surface est un diviseur de Cartier associé au fibré $`𝒪_X(l)`$$`l`$ est un entier convenable. Son degré est donc $`3l`$. On en déduit que la surface $`\pi (S)`$ est une section hyperplane de $`X`$, ce qui est absurde.
Il nous reste à traiter le cas où $`OS`$. Notons $`d`$ le degré de $`\pi `$. La surface $`\pi (S)`$ est donc de degré $`\frac{4}{d}`$. C’est un diviseur de Cartier associé au fibré $`𝒪_X(l)`$$`l`$ est un entier convenable. Son degré est donc $`3l`$ ce qui constitue la contradiction cherchée puisque l’égalité $`3ld=4`$ est impossible avec $`l`$ et $`d`$ entiers. $`\mathrm{}`$
Corollaire 2.6.$``$Le fibré $`E`$ est associé à une quintique elliptique lisse non dégénérée par la construction de Serre.
Démonstration.$``$Il est donné par l’extension (1.5) :
$$0𝒪_X(2)E(1)I_C0$$
$`C`$ est une courbe lisse. On a en particulier $`h^1(I_C)=0`$ et la courbe $`C`$ est donc connexe. On a $`\omega _C=𝒪_C`$ (1.6) et la courbe $`C`$ est donc une courbe elliptique de degré $`c_2(E(1))=5`$. Enfin la courbe $`C`$ est linéairement normale puisque $`E`$ est stable. $`\mathrm{}`$
3. Faisceaux de rang 2 semi-stables sur la cubique de $`^4`$
Proposition 3.1.$``$Soient X une cubique lisse de $`^4`$ et $`E`$ un faisceau de rang 2 semi-stable de classes de Chern $`c_1(E)=0`$, $`c_2(E)=2`$ et $`c_3(E)=0`$. Soit $`F`$ le bidual de $`E`$. Alors ou bien $`E`$ est localement libre ou bien $`F`$ est localement libre de seconde classe de Chern $`c_2(F)=1`$ et $`h^0(F)=1`$ ou bien $`F=H^0(F)𝒪_X`$.
Démonstration.$``$Soit $`S|𝒪_X(1)|`$ une section hyperplane générique telle que $`E_S`$ soit $`\mu `$-semi-stable relativement à la polarisation $`𝒪_S(1)`$ (\[M\] thm 3.1) et telle que $`F_S`$ soit isomorphe au bidual de $`E_S`$. Le faisceau $`F`$ est $`\mu `$-semi-stable. Le faisceau $`F_S`$ est localement libre de rang $`2`$ et $`\mu `$-semi-stable de première classe de Chern $`c_1(F_S)=0`$ (\[H2\]). Notons $`R`$ le conoyau de l’inclusion canonique $`EF`$. Le faisceau $`E`$ est sans torsion et $`R`$ est de dimension au plus 1. On a les formules $`c_2(F_S)=c_2(E_S)+c_2(R_S)=2\mathrm{}(R_S)`$ et $`\chi (F_S)=\mathrm{}(R_S)`$. On en déduit la relation $`h^0(F_S)=h^1(F_S)+\mathrm{}(R_S)`$ puisque $`h^2(F_S)=h^0(F_S(1))=0`$. Supposons $`h^0(F_S)1`$. Le lieu des zéros d’une section non nulle est ou bien vide ou bien de codimension pure 2 (2.1). S’il est vide alors le fibré $`F_S`$ est trivial et s’il est de codimension pure 2 alors $`h^0(F_S)=1`$. On a donc $`\mathrm{}(R_S)\{0,1,2\}`$ et $`c_2(F_S)\{0,1,2\}`$.
Considérons la suite exacte de restriction à une section hyperplane :
$$0F(n1)F(n)F_S(n)0$$
On a $`h^1(F_S(n))=0`$ pour $`n2`$ et $`n1`$ (2.2 et 2.3). On en déduit $`h^1(F(n1))h^1(F(n))`$ pour $`n2`$ et $`h^2(F(n1))h^2(F(n))`$ pour $`n1`$. Or $`h^1(F(n))=0`$ pour $`n0`$ (\[H2\] thm. 2.5) et $`h^2(F(n))=0`$ pour $`n0`$ et on a donc $`h^1(F(n))=0`$ pour $`n2`$ $`h^2(F(n))=0`$ pour $`n0`$. On a enfin $`h^3(F)=h^0(F^{}(2))=h^0(F(2))=0`$ (\[H2\] prop. 1.10) et $`h^0(F)h^0(F_S)`$.
Supposons $`\mathrm{}(R_S)=0`$$``$Alors $`c_2(F)=2`$ et $`\chi (F)=\frac{c_3(F)}{2}`$. On en déduit la formule $`\frac{c_3(F)}{2}=h^0(F)h^1(F)`$. Or $`c_3(F)0`$ (\[H2\] prop. 2.6) et $`h^0(F)h^0(F_S)1`$ et on a donc $`c_3(F)=0`$ ou $`2`$.
Si $`c_3(F)=0`$ alors les faisceaux $`E`$ et $`F`$ sont canoniquement isomorphes et localement libres (\[H2\] prop. 2.6).
Si $`c_3(F)=2`$ alors $`h^0(F)=1`$ et $`h^1(F)=0`$. Le faisceau $`R`$ est donc de dimension 0 et $`\mathrm{}(R)=\chi (F)\chi (E)=1`$. On a donc $`R=k(p)`$ avec $`pX`$. Puisque $`h^0(F)=1`$ on a un morphisme non nul $`𝒪_XF`$. De plus, $`\chi (E(n))=n^3+3n^2+2n`$ et $`\chi (𝒪_X(n))=\frac{n^3}{2}+\frac{3n^2}{2}+2n+1`$ et on en déduit $`h^0(E)=0`$ puisque $`E`$ est semi-stable. Le morphisme induit $`𝒪_XR`$ est donc non nul. Il est surjectif et induit une inclusion $`I_pX`$. Or $`\chi (I_p(n))=\frac{n^3}{2}+\frac{3n^2}{2}+2n`$ ce qui est en contradiction avec la semi-stabilité de $`E`$.
Supposons $`\mathrm{}(R_S)1`$. On a donc $`h^0(F_S)1`$. Le lieu des zéros d’une section globale non nulle est ou bien vide ou bien de codimension pure 2 (2.1).
Supposons qu’il existe une section non nulle de $`F_S`$ dont le lieu des zéros est de codimension pure 2. Alors $`h^0(F_S)=1`$. On en déduit $`h^1(F_S)=0`$ puis $`\mathrm{}(R_S)=1`$ et $`c_2(F)=1`$. On en déduit l’inégalité $`\chi (F)=h^0(F)h^1(F)=1+\frac{c_3(F)}{2}1h^1(F)`$. Puis $`c_3(F)=0`$ puisque $`c_3(F)0`$ (\[H2\] prop. 2.6). Le faisceau $`F`$ est donc localement libre (\[H2\] prop. 2.6) de seconde classe de Chern $`c_2(F)=1`$ et $`h^0(F)=1`$.
Supposons enfin qu’il existe une section du fibré $`F_S`$ ne s’annulant pas auquel cas ledit fibré est isomorphe au fibré $`H^0(F_S)𝒪_S`$ et donc $`\mathrm{}(R_S)=2`$ et $`c_2(F)=0`$. On en déduit l’inégalité $`\chi (F)=\frac{c_3(F)}{2}+2=h^0(F)h^1(F)2h^1(F)`$. puis $`c_3(F)=0`$ puisque $`c_3(F)0`$ (\[H2\] prop. 2.6) et $`h^0(F)=2`$. Le faisceau $`F`$ est donc localement libre (\[H2\] prop. 2.6). Supposons qu’il existe une section globale non nulle de $`F`$ dont le lieu des zéros $`Z`$ est non vide. Le schéma $`Z`$ est de dimension pure 1 puisque $`h^0(F(1))=0`$ et $`F`$ est donc extension de $`I_Z`$ par $`𝒪_X`$. On en déduit $`h^0(F)=1`$ ce qui est absurde. Le faisceau $`F`$ est donc isomorphe au fibré $`H^0(F)𝒪_X`$. $`\mathrm{}`$
Lemme 3.2.$``$Soit $`R`$ un faisceau cohérent sur $`^n`$ $`(n1)`$ tel que $`h^0(R(1))=0`$ et $`\chi (R(n))=n+1`$. Il existe alors une droite $`\mathrm{}^n`$ telle que $`R=𝒪_{\mathrm{}}`$.
Démonstration.$``$Le faisceau $`R`$ est de dimension 1 et on a donc $`h^0(R)=h^1(R)+11`$. Soient $`sH^0(R)`$ une section non nulle et $`I_Z`$ le noyau de l’application induite $`𝒪_^nR`$. On a $`h^0(𝒪_Z(1))=0`$ et $`Z`$ est donc de dimension pure 1. Considérons une section hyperplane générique $`S|𝒪_^n(1)|`$. On a $`\mathrm{}(R_S)=1`$ et l’inclusion $`𝒪_{ZS}R_S`$ est donc un isomorphisme. On en déduit que $`Z_{\text{red}}`$ est une droite $`\mathrm{}^n`$ et que le schéma $`Z`$ est génériquement réduit le long de $`\mathrm{}`$. Le noyau de l’application surjective $`𝒪_Z𝒪_{\mathrm{}}`$ est de dimension zéro et donc trivial puisque $`h^0(𝒪_Z(1))=0`$. On a donc $`R=𝒪_{\mathrm{}}`$ puisque ces deux faisceaux ont même polynôme caractéristique. $`\mathrm{}`$
Lemme 3.3.$``$Soit $`R`$ un faisceau cohérent sur $`^n`$ $`(n1)`$ tel que $`h^0(R(1))=0`$ et $`\chi (R(n))=2n+2`$. Alors il existe deux droites $`\mathrm{}_1^n`$ et $`\mathrm{}_2^n`$ telles que $`R`$ soit extension de $`𝒪_\mathrm{}_2`$ par $`𝒪_\mathrm{}_1`$ ou bien $`R(1)`$ est une thêta-caractéristique sur une conique lisse $`C^n`$.
Démonstration.$``$Le faisceau $`R`$ est de dimension 1 et on a donc $`h^0(R)=h^1(R)+22`$. Soient $`sH^0(R)`$ une section non nulle et $`I_Z`$ le noyau de l’application induite $`𝒪_^nR`$. On a $`h^0(𝒪_Z(1))=0`$ et $`Z`$ est donc de dimension pure 1. Soit $`S|𝒪_^n(1)|`$ une section hyperplane générique. On a $`0<\mathrm{}(ZS)\mathrm{}(R_S)=2`$. Notons $`Q`$ le conoyau de l’inclusion $`𝒪_ZR`$.
Supposons $`\mathrm{}(ZS)=1`$.$``$Le support du schéma $`Z`$ est alors une droite $`\mathrm{}_1`$ et ledit schéma est génériquement réduit le long de $`\mathrm{}_1`$. On a donc une application surjective $`𝒪_Z𝒪_\mathrm{}_1`$ dont le noyau est de dimension zéro. Ledit noyau est en fait trivial puisque $`h^0(𝒪_Z(1))=0`$. Enfin, on a $`\chi (Q(n))=n+1`$ et $`h^0(Q(1))=0`$ et le lemme 3.2 permet de conclure.
Supposons $`\mathrm{}(ZS)=2`$ et $`Q`$ non trivial.$``$On a $`h^1(R(1))=0`$ et on a donc $`h^1(R(k))=0`$ pour $`k1`$ (1.8). On en déduit en particulier $`h^0(R)=2`$ et $`h^0(R(1))=4`$. Considérons la suite exacte :
$$0𝒪_Z(1)R(1)Q(1)0$$
$`Q`$ est de dimension zéro. Le faisceau $`R(1)`$ est engendré par ses sections globales (1.8) et l’application $`H^0(R(1))H^0(Q(1))`$ n’est donc pas identiquement nulle. On en déduit $`h^0(𝒪_Z(1))3`$ et $`h^0(I_Z(1))n2`$. Il existe donc un plan $`^2^n`$ contenant le schéma $`Z`$. Notons $`J_Z`$ l’idéal de $`Z`$ dans ledit plan. On a $`c_1(J_Z)=2`$ et on a donc une inclusion $`J_Z𝒪_^2(2)`$ qui induit une application surjective $`𝒪_Z𝒪_C`$$`C`$ est une conique. Son noyau est de dimension zéro et donc trivial puisque $`h^0(𝒪_Z(1))=0`$. On a donc une suite exacte :
$$0𝒪_CRk(p)0$$
Si $`pC`$ alors l’extension précédente est triviale ce qui est absurde puisque $`h^0(R(1))=0`$. On a donc $`pC`$. Montrons que $`R`$ est un $`𝒪_C`$\- module. Soit $`fH^0(I_C(k))`$ $`(k0)`$ l’équation d’une hypersurface de degré $`k`$ contenant $`C`$. Considérons le diagramme commutatif suivant :
$$\begin{array}{ccccccc}& & 0& & 0& & 0\\ & & & & & & & & \\ 0& & 𝒪_C(k)& & K(k)& & k(p)(k)\\ & & & & & & & & \\ 0& & 𝒪_C(k)& & R(k)& & k(p)(k)& & 0\\ & & \times f& & \times f& & \times f& & \\ 0& & 𝒪_C& & R& & k(p)& & 0\end{array}$$
où les complexes horizontaux sont exacts. Si l’application $`Kk(p)`$ est nulle alors on a une inclusion $`R(k)/K(k)R`$ avec $`R(k)/K(k)`$ de dimension zéro ce qui est impossible puisque $`h^0(R(1))=0`$. Ladite application est donc surjective et on en déduit que l’application $`R(k)R`$ est nulle. Le faisceau $`R`$ est donc un $`𝒪_C`$-module. On vérifie alors qu’on a une suite exacte :
$$0I_pH^0(R)𝒪_CR0$$
$`I_p`$ est l’idéal de $`p`$ dans $`C`$. Si $`C`$ est une conique lisse alors $`R(1)`$ est la thêta-caractéristique sur $`C`$. Supposons $`C`$ non lisse et soit $`\mathrm{}C`$ une droite contenant $`p`$. L’inclusion $`I_{\mathrm{}}H^0(R)𝒪_C`$ se factorise à travers l’inclusion $`I_{\mathrm{}}𝒪_C`$ et on obtient ainsi une inclusion $`𝒪_{\mathrm{}}R`$ dont le conoyau est également isomorphe au faisceau structural d’une droite (3.2).
Supposons $`\mathrm{}(ZS)=2`$ et $`Q=0`$.$``$Le schéma $`Z_{\text{red}}`$ est de degré au plus 2. S’il est de degré 2 ledit schéma est ou bien réunion de deux droites distinctes ou bien une conique lisse. On a alors une application surjective $`𝒪_Z𝒪_C`$ dont le noyau est trivial si $`C`$ est réunion de droites disjointes et supporté en un point sinon. Ce dernier cas est impossible puisque $`h^0(𝒪_Z(1))=0`$. Si $`Z_{\text{red}}`$ est de degré 1 alors $`Z_{\text{red}}`$ est une droite $`\mathrm{}^n`$ et on a une surjection $`𝒪_Z𝒪_{\mathrm{}}`$ dont le noyau $`K`$ vérifie $`\chi (K(n))=n+1`$ et $`h^0(K(1))=0`$. Ce noyau est donc isomorphe au faisceau $`𝒪_{\mathrm{}}`$ (3.2). $`\mathrm{}`$
Lemme 3.4.$``$Soient $`X^4`$ une cubique lisse et $`\theta `$ une thêta-caractéristique sur une conique lisse $`CX`$. On considére le faisceau $`E`$ noyau de l’application surjective $`H^0(\theta (1))𝒪_X\theta (1)`$. Alors $`E`$ est stable de classes de Chern $`c_1(E)=0`$, $`c_2(E)=2`$ et $`c_3(E)=0`$.
Démonstration.$``$Le calcul des classes de Chern de $`E`$ est immédiat. Soit $`FE`$ un sous-faisceau de rang 1 de $`E`$. Le faisceau $`F`$ est de la forme $`I_Z(a)`$$`ZX`$ est un sous-schéma fermé de dimension au plus 1 et $`a`$. On a un diagramme commutatif à lignes et colonnes exactes :
$$\begin{array}{ccccccccc}& & 0& & 0& & & & \\ & & & & & & & & & & \\ & & H^0(\theta (1))I_C& =& H^0(\theta (1))I_C& & & & \\ & & & & & & & & & & \\ 0& & E& & H^0(\theta (1))𝒪_X& & \theta (1)& & 0\\ & & & & & & & & & & \\ 0& & \theta & & H^0(\theta (1))𝒪_C& & \theta (1)& & 0\\ & & & & & & & & & & \\ & & 0& & 0& & & & \end{array}$$
Notons $`F_0`$ le noyau de l’application induite $`F\theta `$. On a une inclusion $`F_0H^0(\theta (1))I_C`$. Le faisceau $`H^0(\theta (1))I_C`$ est $`\mu `$-semi-stable de pente nulle et on a donc $`c_1(F)=c_1(F_0)c_1(H^0(\theta (1))I_C)=0`$ puisque $`\theta `$ est de dimension 1.
Si $`c_1(F)<0`$ on a $`\chi (F(n))<\frac{1}{2}\chi (E(n))`$ pour $`n0`$ par un calcul classique. Si $`c_1(F)=0`$ alors $`F=I_Z`$ avec $`\text{codim}(Z)2`$ et on a donc $`I_Z^{}=𝒪_X`$. L’inclusion $`I_ZH^0(\theta (1))𝒪_X`$ déduite de l’inclusion $`EH^0(\theta (1))𝒪_X`$ est donc donnée par une section non nulle $`sH^0(\theta (1))`$. L’application induite $`I_Z\theta (1)`$ associe donc la section $`f_{|C}s`$ à la fonction $`f`$. La section $`s`$ étant génériquement non nulle on en déduit $`I_ZI_C`$ et donc $`\chi (I_Z(n))\chi (I_C(n))<\frac{1}{2}\chi (E(n))`$ pour $`n0`$ puisque $`\chi (I_C(n))=\frac{n^3}{2}+\frac{3n^2}{2}`$ et $`\frac{\chi (E(n))}{2}=\frac{n^3}{2}+\frac{3n^2}{2}+n`$. $`\mathrm{}`$
Théorème 3.5.$``$Soient X une cubique lisse de $`^4`$ et $`E`$ un faisceau de rang 2 semi-stable de classes de Chern $`c_1(E)=0`$, $`c_2(E)=2`$ et $`c_3(E)=0`$. Si $`E`$ est stable alors ou bien $`E`$ est localement libre ou bien $`E`$ est associé à une conique lisse $`CX`$ (3.4). Si $`E`$ est semi-stable non stable alors le gradué de $`E`$ est somme directe des idéaux de deux droites de $`X`$.
Démonstration.$``$Soit $`F`$ le bidual de $`E`$ et $`R`$ le conoyau de l’injection canonique $`EF`$. Le faisceau $`E`$ est localement libre ou bien $`F`$ est localement libre de seconde classe de Chern $`c_2(F)=1`$ et $`h^0(F)=1`$ ou bien $`F=H^0(F)𝒪_X`$ (3.1). On a $`\chi (E(n))=n^3+3n^2+2n`$ et $`\chi (𝒪_X(n))=\frac{n^3}{2}+\frac{3n^2}{2}+2n+1`$ et on en déduit $`h^0(E)=0`$ puisque $`E`$ est semi-stable.
Supposons $`E`$ localement libre.$``$On a $`h^0(E)=0`$ puisque $`E`$ est semi-stable et le fibré $`E`$ est donc stable.
Supposons $`F`$ localement libre de seconde classe de Chern $`c_2(F)=1`$ et $`h^0(F)=1`$.$``$ Alors $`\chi (R(n))=n+1`$. Soit $`sH^0(F)`$ une section non nulle. Elle s’annule le long d’une droite $`\mathrm{}_2X`$ (2.1). On a $`h^0(E)=0`$ et la section $`s`$ de $`F`$ induit une application non nulle $`𝒪_XR`$. Notons $`I_Z`$ le noyau de l’application précédente. Le schéma $`Z`$ est de dimension 1. Sinon on aurait une inclusion $`I_ZE`$ avec $`\chi (I_Z(n))=\frac{n^3}{2}+\frac{3n^2}{2}+2n+1\mathrm{}(Z)`$ ce qui est impossible par semi-stabilité de $`E`$. Soit $`S|𝒪_X(1)|`$ une section hyperplane générique. L’inclusion $`𝒪_{ZS}R_S`$ est un isomorphisme puisque $`\mathrm{}(R_S)=1`$. Aussi la composante de dimension 1 du support de $`Z`$ est une droite $`\mathrm{}_1`$ et on a donc une application surjective $`𝒪_Z𝒪_\mathrm{}_1`$. On en déduit $`R=𝒪_\mathrm{}_1`$ puisque ces deux faisceaux ont même polynôme caractéristique. L’application $`𝒪_X𝒪_\mathrm{}_1`$ est non nulle et les droites $`\mathrm{}_1`$ et $`\mathrm{}_2`$ sont donc disjointes. Considérons le diagramme commutatif à lignes et colonnes exactes :
$$\begin{array}{ccccc}& & 0& & 0\\ & & & & & & & & & & & \\ 0& & I_\mathrm{}_1& & 𝒪_X& & 𝒪_\mathrm{}_1& & 0\\ & & & & & & & & \\ 0& & E& & F& & 𝒪_\mathrm{}_1& & 0\\ & & & & & & & & & & & \\ & & I_\mathrm{}_2& =& I_\mathrm{}_2& & & & & \\ & & & & & & & & & & & \\ & & 0& & 0& & & & & \end{array}$$
On en déduit que le faisceau $`E`$ est semi-stable non stable et que son gradué est $`I_\mathrm{}_1I_\mathrm{}_2`$.
Supposons $`F=H^0(F)𝒪_X`$.$``$ On a $`\chi (R(n))=2n+2`$ et $`h^0(E)=0`$ puisque $`E`$ est semi-stable. On en déduit en particulier $`h^0(R)2`$. Considérons une section hyperplane générique $`S|𝒪_X(1)|`$. On a la suite exacte :
$$0E_SH^0(X,F)𝒪_SR_S0$$
L’application $`H^0(H^0(X,F)𝒪_S)H^0(R_S)`$ n’est pas nulle puisque le morphisme $`H^0(F)𝒪_SR_S`$ est surjectif. On a donc $`h^0(E_S)1`$. Si $`h^0(E_S)=0`$ alors l’application $`H^0(F_S)H^0(R_S)`$ est un isomorphisme et l’application $`H^0(F_S)H^0(𝒪_S(n))H^0(R_S(n))`$ est surjective pour tout $`n0`$ puisqu’il existe une section hyperplane de $`S`$ évitant le support de $`R_S`$. Si $`h^0(E_S)=1`$ alors le quotient $`H^0(F_S)/H^0(E_S)`$ est de dimension 1 et on a une application surjective $`(H^0(F_S)/H^0(E_S))𝒪_SR_S`$. Or $`\mathrm{}(R_S)=2`$ et $`𝒪_S(1)`$ est très ample et l’application $`(H^0(F_S)/H^0(E_S))H^0(𝒪_S(n))H^0(R_S(n))`$ est donc surjective pour $`n1`$. Il en résulte que l’application $`H^0(F_S)H^0(𝒪_S(n))H^0(R_S(n))`$ est également surjective. On a donc finalement $`h^1(E_S(n))=0`$ pour $`n1`$ puisque $`h^1(𝒪_S(n))=0`$ pour $`n1`$. La suite exacte :
$$0E(n1)E(n)E_S(n)0$$
donne $`h^2(E(n1)h^2(E(n))`$ pour $`n1`$. On a donc $`h^2(E(n))=0`$ pour $`n0`$ puisque $`h^2(E(n))=0`$ pour $`n0`$. En particulier $`h^2(E)=0`$. On déduit de la suite exacte :
$$0EH^0(X,F)𝒪_XR0$$
l’égalité $`h^3(E)=0`$. Mais $`\chi (E)=0`$ et on a donc $`h^1(E)=0`$. On en déduit $`h^0(R)=2`$ et l’inclusion $`H^0(F)H^0(R)`$ est donc un isomorphisme. Montrons alors que l’application de restriction $`H^0(R)H^0(R_S)`$ est injective. Supposons le contraire. Il existe donc une section $`sH^0(R)`$ non nulle dont l’image dans $`H^0(R_S)`$ est nulle. Notons $`I_Z`$ le noyau de l’application $`𝒪_XR`$ définie par le section $`s`$ et $`Q`$ le conoyau de l’inclusion $`𝒪_ZR`$. Par hypothèse, l’application $`𝒪_{ZS}R_S`$ est nulle et on a donc $`R_S=Q_S`$. Le faisceau $`Q`$ est donc de dimension 1 et $`c_2(Q)=c_2(R)`$. Le schéma $`Z`$ est donc de dimension 0. Or $`\chi (I_Z(n))=\frac{n^3}{2}+\frac{3n^2}{2}+2n+1\mathrm{}(Z)`$ ce qui est en contradiction avec la semi-stabilité de $`E`$ puisqu’on a une inclusion $`I_ZE`$. L’application de restriction $`H^0(R)H^0(R_S)`$ est injective et on a donc $`h^0(R(1))=0`$. Le faisceau $`R(1)`$ est donc ou bien une thêta-caractéristique sur une conique lisse $`CX`$ auquel cas $`E`$ est stable (3.4) ou bien il existe deux droites $`\mathrm{}_1X`$ et $`\mathrm{}_2X`$ telles que $`R`$ soit extension de $`𝒪_\mathrm{}_1`$ par $`𝒪_\mathrm{}_2`$ (3.3) auquel cas on a un diagramme commutatif à lignes et colonnes exactes :
$$\begin{array}{ccccccc}& & 0& & 0& & 0\\ & & & & & & & & \\ 0& & I_\mathrm{}_1& & E& & I_\mathrm{}_2& & 0\\ & & & & & & & & \\ 0& & 𝒪_X& & H^0(F)𝒪_X& & 𝒪_X& & 0\\ & & & & & & & & \\ 0& & 𝒪_\mathrm{}_1& & R& & 𝒪_\mathrm{}_2& & 0\\ & & & & & & & & \\ & & 0& & 0& & 0\end{array}$$
On en déduit que le gradué de $`E`$ est le faisceau $`I_\mathrm{}_1I_\mathrm{}_2`$. $`\mathrm{}`$
4. Espace des modules des faisceaux semi-stables sur la cubique de $`^4`$
(4.1) Soit $`X^4`$ une hypersurface cubique lisse et soit $`(Def(E),0)`$ l’espace des déformations verselles d’un faisceau cohérent $`E`$ sur $`X`$. L’espace tangent à $`Def(E)`$ en $`0`$ s’identifie à l’espace vectoriel $`\text{Ext}_X^1(E,E)`$. Le germe analytique $`Def(E)`$ est lisse en $`0`$ si $`\text{Ext}_X^2(E,E)=0`$.
Lemme 4.2.$``$Soient $`X^4`$ une cubique lisse et $`\theta `$ une thêta-caratéristique sur une conique lisse $`CX`$. Soit $`E`$ le noyau de la surjection $`H^0(\theta (1))𝒪_X\theta (1)`$. Alors $`\text{Ext}_X^2(E,E)`$ est nul et l’espace vectoriel complexe $`\text{Ext}_X^1(E,E)`$ est de dimension 5.
Démonstration.$``$Soit $`F`$ le noyau de la surjection $`𝒪_^4𝒪_^4\theta (1)`$. On vérifie que le faisceau $`F(1)`$ est engendré par ses sections globales en utilisant le critère de Mumford-Castelnuovo (1.8). On en déduit que le faisceau $`E(1)`$ est également engendré par ses sections globales puisqu’on a un morphisme surjectif $`F_{|X}(1)E(1)`$. On a donc $`\text{Hom}_X(E,\theta (1))\text{Hom}_X(H^0(E(1))𝒪_X(1),\theta (1))=0`$. Considérons la suite exacte :
$$\text{Ext}_X^2(H^0(\theta (1))𝒪_X,E)\text{Ext}_X^2(E,E)\text{Ext}_X^3(\theta (1),E)$$
On a $`\text{Ext}_X^2(H^0(\theta (1))𝒪_X,E)H^0(\theta (1))H^2(E)=0`$ et $`\text{Ext}_X^3(\theta (1),E)\text{Hom}_X(E,\theta (1))^{}=0`$ et donc $`\text{Ext}_X^2(E,E)=0`$. Enfin, $`\text{Ext}_X^3(E,E)\text{Hom}_X(E,E(2))^{}=0`$ et $`\text{Hom}_X(E,E)`$. L’espace vectoriel complexe $`\text{Ext}_X^1(E,E)`$ est donc de dimension 5 puisque $`\chi (E,E)=(1)^i\text{Ext}_X^i(E,E)=4`$. $`\mathrm{}`$
Lemme 4.3.$``$Soit $`X^4`$ une cubique lisse et soient $`\mathrm{}_1X`$ et $`\mathrm{}_2X`$ deux droites. Le groupe $`\text{Ext}_X^2(I_\mathrm{}_1,I_\mathrm{}_2)`$ est nul et l’espace vectoriel complexe $`\text{Ext}_X^1(I_\mathrm{}_1,I_\mathrm{}_2)`$ est de dimension 1 si $`\mathrm{}_1\mathrm{}_2`$ et de dimension 2 si $`\mathrm{}_1=\mathrm{}_2`$.
Démonstration.$``$On a un isomorphisme $`\text{Ext}_X^3(𝒪_\mathrm{}_1,I_\mathrm{}_2)\text{Hom}_X(I_\mathrm{}_2,𝒪_\mathrm{}_1(2))^{}`$. Or le faisceau $`I_\mathrm{}_2(1)`$ est engendré par ses sections globales et on a donc $`\text{Hom}_X(I_\mathrm{}_2,𝒪_\mathrm{}_1(2))\text{Hom}_X(H^0(I_\mathrm{}_2(1))𝒪_X(1),𝒪_\mathrm{}_1(2))=0`$. Considérons la suite exacte :
$$\text{Ext}_X^2(𝒪_X,I_\mathrm{}_2)\text{Ext}_X^2(I_\mathrm{}_1,I_\mathrm{}_2)\text{Ext}_X^3(𝒪_\mathrm{}_1,I_\mathrm{}_2)$$
On en déduit $`\text{Ext}_X^2(I_\mathrm{}_1,I_\mathrm{}_2)=0`$ puisque $`\text{Ext}_X^2(𝒪_X,I_\mathrm{}_2)=0`$. On a $`\text{Ext}_X^3(I_\mathrm{}_1,I_\mathrm{}_2)\text{Hom}_X(I_\mathrm{}_2,I_\mathrm{}_1(2))^{}=0`$ et $`\chi (I_\mathrm{}_1,I_\mathrm{}_2)=(1)^i\text{Ext}_X^i(I_\mathrm{}_1,I_\mathrm{}_2)=\text{Ext}_X^0(I_\mathrm{}_1,I_\mathrm{}_2)\text{Ext}_X^1(I_\mathrm{}_1,I_\mathrm{}_2)=\chi (I_\mathrm{}_1,I_\mathrm{}_1)=1`$. L’espace vectoriel complexe $`\text{Ext}_X^1(I_\mathrm{}_1,I_\mathrm{}_2)`$ est donc de dimension 1 si $`\mathrm{}_1\mathrm{}_2`$ et de dimension 2 si $`\mathrm{}_1=\mathrm{}_2`$ puisque $`\text{Hom}_X(I_\mathrm{}_1,I_\mathrm{}_2)=0`$ si $`\mathrm{}_1\mathrm{}_2`$ et $`\text{Hom}_X(I_\mathrm{}_1,I_\mathrm{}_2)`$ si $`\mathrm{}_1=\mathrm{}_2`$. $`\mathrm{}`$
(4.4) Soient $`N1`$ un entier et $`V`$ un espace vectoriel complexe. Soient $`Q`$ le schéma de Hilbert-Grothendieck paramétrant les quotients $`V𝒪_X(N)E`$ sur $`X`$ de rang 2 et de classes de Chern $`c_1(E)=0`$, $`c_2(E)=2`$, $`c_3(E)=0`$ et $`L`$ la polarisation naturelle (\[S\]). Le groupe $`G=PGL(V)`$ agit sur $`Q`$ et une puissance convenable de $`L`$ est $`G`$-linéarisée. Soit $`Q_c^{ss}`$ l’ouvert des points $`PGL(V)`$-semi-stables correspondants à des quotients sans torsion et $`Q_c`$ l’adhérence de $`Q_c^{ss}`$ dans $`Q`$. Lorsque l’entier $`N`$ et l’espace vectoriel $`V`$ sont convenablement choisis les propriétés suivantes sont satisfaites. L’application $`V𝒪_XE(N)`$ induit un isomorphisme $`VH^0(E(N))`$ et $`h^i(E(k))=0`$ pour $`kN`$ et $`i1`$ et pour tout quotient $`E`$ de $`Q_c`$. Le point $`[E]Q_c`$ est $`PGL(V)`$-semi-stable si et seulement si le faisceau $`E`$ est semi-stable si et seulement si $`[E]Q_c^{ss}`$. Le stabilisateur de $`[E]`$ dans $`GL(V)`$ s’identifie au groupe des automorphismes du faisceau $`E`$ et l’espace des modules $`M`$ est alors le quotient de Mumford :
$$Q_c^{ss}//G$$
Lemme 4.5.$``$Sous les hypothèses précédentes, le schéma $`Q_c^{ss}`$ est lisse.
Démonstration.$``$L’espace tangent à $`Q_c^{ss}`$ en un point $`[E]`$ est isomorphe à $`\text{Hom}_X(F,E)`$$`F`$ est le noyau de l’application $`V𝒪_X(N)E`$. Le schéma $`Q_c^{ss}`$ est lisse en ce point si $`\text{Ext}_X^1(F,E)=0`$. Considérons la suite exacte :
$$\text{Ext}_X^1(V𝒪_X(N),E)\text{Ext}_X^1(F,E)\text{Ext}_X^2(E,E)$$
On en déduit une inclusion $`\text{Ext}_X^1(F,E)\text{Ext}_X^2(E,E)`$ puisque $`h^1(E(N))=0`$. Il suffit donc de prouver $`\text{Ext}_X^2(E,E)=0`$. Si $`E`$ est stable et localement libre alors le résultat est démontré par \[M-T\] (lemme 2.7). Si $`E`$ est stable non localement libre alors l’annulation cherchée est donnée par le lemme 4.2. Si $`E`$ est strictement semi-stable alors $`E`$ est extension de $`I_\mathrm{}_2`$ par $`I_\mathrm{}_1`$$`\mathrm{}_1X`$ et $`\mathrm{}_2X`$ sont deux droites. L’annulation cherchée résulte alors facilement du lemme 4.3. $`\mathrm{}`$
Théorème 4.6.$``$Soit $`X^4`$ une hypersurface cubique lisse. L’espace des modules $`M_X`$ des faisceaux semi-stables de rang 2 sur $`X`$ de classes de chern $`c_1=0`$, $`c_2=2`$ et $`c_3=0`$ est lisse de dimension 5.
Démonstration.$``$Soient $`xQ_c^{ss}`$ et $`E`$ le faisceau correspondant. Soit $`Q_c^sQ_c`$ l’ouvert des faisceaux stables et $`M_X^s`$ l’ouvert des classes d’isomorphismes de faisceaux stables. Le schéma $`Q_c^s`$ est un fibré principal sous $`G`$ au dessus de $`M_X^s`$ et $`M_X^s`$ est donc lisse (4.5). Il nous reste à étudier $`M_X`$ en $`E=I_\mathrm{}_1I_\mathrm{}_2`$$`\mathrm{}_1X`$ et $`\mathrm{}_2X`$ sont deux droites (3.5). L’orbite $`O(x)`$ du point $`x`$ sous $`G`$ est fermée. Son stabilisateur $`G_x`$ est un groupe réductif et il existe un sous-schéma affine $`WQ_c^{ss}`$ passant par $`x`$ localement fermé et stable sous l’action de $`G_x`$ tel que le morphisme $`W//G_xQ_c^{ss}//G`$ soit étale (\[L\]). Soit $`(W,x)`$ le germe de $`W`$ en $`x`$ et soit $`F`$ la restriction à $`X\times (W,x)`$ du quotient tautologique sur $`X\times Q`$. Alors $`((W,x),F)`$ est un espace de déformation verselles pour le faisceau $`E`$ (\[O\] prop. 1.2.3). Le germe $`W`$ est donc lisse en $`x`$ (4.3) et puisque le morphisme $`W//G_xQ_c^{ss}//G`$ est étale il suffit donc prouver que le quotient $`W//G_x`$ est lisse en $`[x]`$. Or il existe un morphisme $`G_x`$-linéaire $`(W,x)(T_xW,0)`$ étale en $`x`$ tel que le morphisme induit $`W//G_xT_xW//G_x`$ soit étale en $`[x]`$ (\[L\]). Il suffit donc de prouver que le quotient $`T_xW//G_x`$ est lisse en $`0`$.
Supposons les droites $`\mathrm{}_1`$ et $`\mathrm{}_2`$ distinctes. L’espace tangent $`T_xW=\text{Ext}_X^1(E,E)=_{i,j}\text{Ext}_X^1(I_\mathrm{}_i,I_\mathrm{}_j)`$ est de dimension 6 (4.3) et $`G_x=G_m\times G_m`$ agit sur ledit espace par la formule (\[O\] lemme 1.4.16):
$$(t,t^{}).(\underset{i,j}{}e_{i,j})=e_{1,1}+\frac{t}{t^{}}e_{1,2}+\frac{t^{}}{t}e_{2,1}+e_{2,2}$$
On vérifie facilement que le quotient $`T_xW//G_x`$ est isomorphe à l’espace affine $`𝔸^5`$ et en particulier lisse en $`0`$.
Supposons les droites $`\mathrm{}_1`$ et $`\mathrm{}_2`$ confondues et notons $`\mathrm{}`$ cette droite. L’espace tangent $`T_xW=\text{Ext}_X^1(E,E)`$ est de dimension 8 (4.3) et $`G_x=PGL(2)`$. Posons $`T=\text{Ext}_X^1(I_{\mathrm{}},I_{\mathrm{}})`$ et soit $`U`$ un espace vectoriel de dimension 2. Le groupe $`G_x`$ agit sur $`T_xW=T\text{End}(U)`$ par conjugaison sur $`\text{End}(U)`$ (\[O\] lemme 1.4.16). Le quotient $`T_x//G_x`$ est isomorphe à l’espace affine $`𝔸^5`$ (\[La\] III cas 2) et en particulier lisse en $`0`$. $`\mathrm{}`$
Lemme 4.7.$``$Soient $`X^4`$ une cubique lisse et $`M_X`$ l’espace des modules des faisceaux de rang 2 semi-stables de classes de Chern $`c_1=0`$, $`c_2=2`$ et $`c_3=0`$. Le sous-schéma localement fermé de $`M_X`$ des faisceaux stables non localement libres est irréductible de dimension 4 et le sous-schéma fermé de $`M_X`$ des faisceaux strictement semi-stables est également irréductible de dimension 4.
Démonstration.$``$Soient $`B`$ la surface de Fano de $`X`$ et $`ZX\times B`$ la variété d’incidence. La variété $`Z`$ est lisse et irréducible de dimension 3 et le morphisme $`ZB`$ induit par la seconde projection fait de $`Z`$ un fibré en droites projectives localement trivial pour la topologie de Zariski (\[T\]). Notons $`X_Z`$ la variété obtenue en éclatant $`Z`$ dans le produit $`X\times B`$ et notons $`p`$ et $`q`$ les morphismes induits sur $`X`$ et $`B`$ respectivement. La fibre du morphisme $`q`$ au-dessus d’un point $`[\mathrm{}]B`$ s’identifie à $`X_{\mathrm{}}`$ (1.1). Soit $`Q`$ le fibré de rang 3 sur $`B`$ dont la fibre au-dessus d’un point $`[\mathrm{}]B`$ est l’ensemble des équations des hyperplans de $`^4`$ contenant la droite $`\mathrm{}`$. Le morphisme surjectif naturel $`q^{}Qp^{}𝒪_X(1)`$ induit un morphisme propre et dominant $`X_Z\stackrel{𝑐}{}_B(Q)`$. Le morphisme $`p\times c`$ induit un plongement de $`X_Z`$ dans $`X\times _B(Q)`$ au-dessus de $`_B(Q)`$. L’ensemble des points de $`_B(Q)`$ est en bijection ensembliste avec l’ensemble des coniques tracées sur $`X`$. Soit $`U_B(Q)`$ l’ouvert des coniques lisses et soit $`\pi `$ la projection de $`_B(Q)`$ sur $`B`$. Soit $`xU`$ et $`C_x=c^1(x)X`$ la conique correspondante. La conique $`C_x`$ engendre un plan de $`^4`$ dont l’intersection résiduelle avec $`X`$ est la droite $`\pi (x)=[\mathrm{}_x]`$. Soit $`Y=(c^1(U)Z_U)_{\text{red}}c^1(U)U\times ^4`$$`Z_U=Z\times _BU`$. Le morphisme induit $`YU`$ est alors fini de degré 2. La fibre du morphisme précédent au dessus d’un point $`xU`$ est ensemblistement l’intersection $`C_x\mathrm{}_x`$. Supposons $`Y`$ irréductible. Le schéma $`c^1(U)\times _UY`$ posséde une section au dessus de $`Y`$. Ladite section détermine un morphisme quasi-fini $`YM_X`$ dont l’image est précisément le sous-schéma localement fermé des faisceaux stables non localement libres (3.4). Le lemme est donc démontré dans ce cas. Si $`Y`$ n’est pas irréductible alors le morphisme $`c`$ posséde une section au-dessus de l’ouvert $`U`$ et l’argument précédent s’applique directement.
Les faisceaux strictement semi-stables sont paramétrés par les couples de droites de $`X`$ (3.5). Or on a un morphisme quasi-fini $`B\times BM_X`$ dont l’image est le fermé des faisceaux strictement semi-stables. La surface $`B`$ est irréductible et ledit fermé l’est donc également. $`\mathrm{}`$
Théorème 4.8.$``$Soient $`X^4`$ une cubique lisse et $`M_X`$ l’espace des modules des faisceaux de rang 2 semi-stables de classes de Chern $`c_1=0`$, $`c_2=2`$ et $`c_3=0`$. Soit $`B`$ la surface de Fano de $`X`$. Alors $`M_X`$ est isomorphe à l’éclatement d’un translaté de la surface $`B`$ dans la jacobienne intermédiaire de $`X`$.
Démonstration.$``$Soit $`UM_X`$ l’ouvert des fibrés vectoriels stables ; $`M_XU`$ est de dimension 4 (3.5 et 4.7). L’espace des modules $`M_X`$ est lisse de dimension 5 (4.6) et l’ouvert $`U`$ est donc dense. La variété $`M_X`$ est donc irréductible (2.6 et \[I-M\] cor. 5.1).
L’espace des modules $`M_X`$ est le quotient de Mumford $`Q_c^{ss}//G`$$`Q_c^{ss}`$ est lisse (4.5). Soit $``$ une famille universelle sur $`Q_c^{ss}\times X`$. Fixons $`t_0Q_c^{ss}`$. L’application
$$\begin{array}{ccc}\hfill Q_c^{ss}& & J(X)\hfill \\ \hfill t& & c_2(_t)c_2(_{t_0})\hfill \end{array}$$
est algébrique (1.1) et équivariante sous l’action du groupe $`G`$. On en déduit un morphisme que nous noterons encore $`c_2`$ de $`M_X`$ vers $`J(X)`$. Ce morphisme est birationnel (\[I-M\] thm. 3.2) et induit un isomorphisme par restriction à l’ouvert des fibrés stables (2.6 et \[M-T\]).
Les variétés $`M_X`$ et $`J(X)`$ sont lisses et le lieu exceptionnel $`D`$ de $`c_2`$ est donc de codimension pure 1. Le diviseur $`D`$ a au plus deux composantes irréductibles. La restriction de $`c_2`$ au diviseur des faisceaux strictement semi-stables est génériquement finie et le morphisme $`c_2`$ est donc un isomorphisme au point générique du diviseur considéré.
La grassmanienne des plans de $`^4`$ est rationnelle et les cubiques planes tracées sur $`X`$ sont donc toutes rationnellement équivalentes. Soient $`C_0`$ et $`C_1`$ deux coniques tracées sur $`X`$ et $`\mathrm{}_0`$ et $`\mathrm{}_1`$ les intersections résiduelles respectives des plans des coniques avec $`X`$. On a donc $`C_1=C_0+\mathrm{}_0\mathrm{}_1`$ dans $`J(X)`$. Si $`E`$ est un faisceau associé à une conique $`CX`$ (3.4) alors $`c_2(E)=C`$. On en déduit que le diviseur adhérence des faisceaux stables non localement libres est contracté sur un translaté de $`B`$ dans $`J(X)`$. Ce diviseur est irréductible (4.7) et $`M_X`$ est donc isomorphe à l’éclatement d’un translaté de $`B`$ dans $`J(X)`$ (\[L\] thm. 2). $`\mathrm{}`$
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Stéphane Druel
DMA-École Normale Supérieure
45 rue d’Ulm
75005 PARIS
e-mail: druel@clipper.ens.fr |
warning/0002/cond-mat0002359.html | ar5iv | text | # Upper critical field in layered superconductors
## I Introduction
The theory of upper critical field in highly anisotropic quasi-two-dimensional superconductors for the field orientation parallel to conducting layers has been developed by A.Lebed’ and K.Yamaji . It was shown that like in quasi-one-dimentional case at low enough temperatures (see below) the upper critical field starts to diverge such that a superconductivity is conserved in arbitrary large magnetic fields. Moreover in high enough magnetic fields (see below) the critical temperature of a superconducting phase transition has a tendency to restore of its zero field value. These statements are literally true for the superconductors with triplet pairing. For a singlet pairing the existance of a superconductivity in high fields is restricted by the paramagnetic limit $`H_p`$. So, the low temperature stability of a superconducting state under magnetic field can serve as the indication on the superconductivity with triplet pairing.
The tendency for low temperature divergence of the upper critical field $`H_{c2}>H_p`$ in quasi-two-dimentional $`\kappa `$-type ET organic superconductors has been reported recently by T.Ishiguro . Similar behavior has been observed earlier in quasi-one-dimensional organic superconductor $`(TMTSF)_2PF_6`$ by the group of P.Chaikin . So, the observations being in correspondence with theoretical predictions say in favor of triplet type of superconductivity in both type of these materials.
Another popular layered superconductor is $`Sr_2RuO_4`$. It demonstrates the properties compatible with triplet superconductivity . Unlike isostructural high-$`T_c`$ cuprates $`La_{2x}Ba_xCuO_4`$ the normal state of layered perovskite oxide $`Sr_2RuO_4`$ conforms with the predictions of Fermi-liquid theory . Unconventional nature of the superconducting state in this material manifests itself through the strong supression of $`T_c`$ by nonmagnetic impurities as well as by the exibition of a sharp decrease without the coherence peak of $`{}_{}{}^{101}Ru`$ nucleare spin-relaxation rate $`1/T_1`$ followed by $`T^3`$ behavior down to $`0.15K`$ . The spin part of $`{}_{}{}^{17}O`$ Knight shift for the field $`H=0.65T`$ parallel to $`\mathrm{𝑎𝑏}`$ plane does not change down to $`15mK`$, much below $`T_c(H)1.2K`$ . Since the experiment has been performed on $`Sr_2RuO_4`$ in a clean limit this constitues the evidence for the triplet spin pairing with spin of paires lying in $`\mathrm{𝑎𝑏}`$ plane. On the other hand the measurements of the basal plane upper critical field down to $`0.2K`$ shows no tendency to divergency of $`H_{c2}(T)`$ ( see also results of measurements on less perfect crystals ). Moreover, the search for reentrance of superconductivity under magnetic fields up to $`33T`$ and temperatures down to 50 mK has given the negative results . The basal plane upper critical field saturates at low temperatures at $`1.5T`$ which is well below the paramagnetic limit $`H_p2.8T`$ and roughly corresponds to quasiclassical upper critical field value
$$H_{c2}=\frac{\mathrm{\Phi }_0}{2\pi \xi _{ab}\xi _c},$$
(1)
where $`\mathrm{\Phi }_0`$ is the flux quantum and $`\xi _{ab}`$, $`\xi _c`$ are the coherence lengths in basal plane and along the $`c`$ axis correspondingly.
These observations being certainly in contradiction with theory stimulate us to reinvestigate the problem of upper critical field in quasi-two-dimensional superconductors taking into account scattering of quasiparticles on the impurities. The main goal is to find a limits for crystal purity when one can hope to see the low temperature upper critical field divergency. In a few words the result can be described as follows. In a pure crystal the $`H_{c2}`$ divergency starts to be developed when the thermal coherence length $`\xi (T)=v_F/2\pi T`$ begins to be larger than the ”cyclotron radius” of the quasiparticles orbits $`R_c(H)=v_F/\omega _c`$. Here $`v_F`$ is the basal plane Fermi velocity, $`\omega _c=eHv_Fd/c`$ is the ”cyclotron frequency”, $`d`$ is the distance between conducting layers. The impurities does not prevent the upper critical field divergency if at the temperature determined by equality
$$\xi (T)=R_c(H_{c2})$$
(2)
the quasiparticles mean free path $`l`$ is still larger than $`\xi (T)`$. Otherwise the upper critical field is saturated at temperatures below
$$T_l=\frac{v_F}{2\pi l}$$
(3)
and its value is roughly described by formula (1).
The paper has following structure. The general formalism for the upper critical field problem in layered conventional and unconventional superconductors is introduced in the next Section. By the way of the derivation the several simplifications have been used: (i) The equations are written for axially symmetric crystal where the only one-dimentional, even or odd in respect to the reflections in ($`𝐇`$, crystal axis) plane superconducting states are realized; (ii) Among them the only even superconducting states are considered; (iii) The equations are derived with but solved in neglect Pauli paramagnetic interaction; (iv)The equations are solved for only unconventional superconducting states with zero anomalous Green function self-energy. The analytical solution of the equations accompanied by the discussion of limits of crystal purity sufficient for the upper critical field saturation at low temperature is containned in the third Section.
## II The order parameter equations
The electron spectrum of a layered crytals obeys basal plain anizotropy, $`z`$-dependent corrugation of the Fermi surface and several bands in general case. It seems however unimportant for our purposes to take into account all these complifications. So, we shall consider a metal with electron spectrum
$`ϵ`$ $`(𝐩)={\displaystyle \frac{1}{2m}}(p_x^2+p_y^2)2t\mathrm{cos}p_zd,`$ (4)
$`t`$ $`ϵ_F={\displaystyle \frac{mv_F^2}{2}},{\displaystyle \frac{\pi }{d}}<p_z<{\displaystyle \frac{\pi }{d}}`$ (5)
in the magnetic field $`𝐇=(0,H,0)`$, $`𝐀=(0,0,Hx)`$ parallel to the conducting layers with distance $`d`$ between them. In absence of impurities the normal state electron Green function $`G_{\omega _n,\sigma }(p_y,p_z,x,x^{})`$ is obtained as the result of solution of the equation
$$\left[i\omega \frac{1}{2m}\left(\frac{d^2}{dx^2}+p_y^2\right)+2t\mathrm{cos}\left(p_zd\frac{\omega _cx}{v_F}\right)+\sigma \mu _eH+\mu \right]G_{\omega ,\sigma }(p_y,p_z,x,x^{})=\delta (xx^{}),$$
(6)
where $`\omega _n=\pi T(2n+1)`$ is the Matsubara frequency, $`\omega _c=ev_FdH/c`$, $`v_F`$ is the Fermi velocity, $`\mathrm{}=1`$, $`\mu `$ is the chemical potential, $`\mu _e`$ is a magnetic moment of an electron in a crystal. To use the diagonal shape of the Green function matrix we have chosen the $`\widehat{y}`$ direction as the spin quantization axis, such that $`\sigma =\pm 1`$.
In a layered crystal with a singlet Cooper pairing the superconducting states obey the following order parameters
$$\widehat{\mathrm{\Delta }}^s=i\mathrm{\Delta }^s(p_x,p_y,p_z,x)\widehat{\sigma }_y,$$
(7)
Here, $`\widehat{\sigma }_y`$ is the Pauli matrix. For a triplet superconductivity we shall limit ourselves by the consideration of the so called equal spin pairing states with spin lying in the plane of the conducting layers. In neglect of small effects of spontaneous magnetism a vector wave function of such the states is
$$𝐝=\widehat{z}\mathrm{\Delta }^t(p_x,p_y,p_z,x).$$
(8)
As we have put the spin quantization axis along $`\widehat{y}`$ direction we will use the corresponding basis of Pauli matrices $`\stackrel{}{\sigma }=(\widehat{\sigma }_y,\widehat{\sigma }_z,\widehat{\sigma }_x)`$. So the order parameter for triplet pairing state in our case is
$$\widehat{\mathrm{\Delta }}^t=i(𝐝\stackrel{}{\sigma })\sigma _y=\mathrm{\Delta }^t(p_x,p_y,p_z,x)\widehat{\sigma }_z.$$
(9)
The order parameter function is represented as a linear combination of the basis functions $`\psi _i(\varphi ,p_z)`$ of one of the irreducible representations of crystal symmetry group
$$\mathrm{\Delta }^{s,t}(p_x,p_y,p_z,x)=\psi _i(\varphi ,p_z)\mathrm{\Delta }_i^{s,t}(x).$$
(10)
Here, $`\varphi `$ is the angle between basal plane vector of momentum $`𝐩_{}`$ and magnetic field $`𝐇\widehat{y}`$. There are only one and two dimensional representations ($`i=1,\mathrm{}d;d=1,2`$) in the crystals with axially symmetric spectrum (5).
The crystal with axial symmetry under magnetic field lying in the basal plane obeys the symmetry in respect of reflections in plane where the vectors of magnetic field and crystal axis lie, that is $`(y,z)`$ plane in our case <sup>*</sup><sup>*</sup>*For the uniaxial crystals with hexagonal or tetragonal symmetry this property takes place only for particular directions of magnetic field in the basal plane.. The two components of vector basis functions $`(\psi _1(\varphi ,p_z),\psi _2(\varphi ,p_z))`$ can be always chosen such that each of them will have definite (even or odd) and at the same time mutually opposite parity in respect to reflections in $`(y,z)`$ plane. As the consequence the set of the order parameter equations for two component superconductivity splits on two independent equation for each component of the order parameter. One of them corresponds to the higher value of the upper critical field. Thus we always deal with one-component superconductivity with definite parity. The simplest examples of even functios $`\psi (\varphi ,p_z)`$ are: $`1,\mathrm{sin}^2\varphi \mathrm{cos}^2\varphi `$ for singlet pairing and $`\mathrm{sin}(p_zd),\mathrm{cos}\varphi `$ for triplet pairing. As an examples of odd states one can pointed out on $`\mathrm{sin}(p_zd)\mathrm{sin}\varphi `$ for singlet pairing and $`\mathrm{sin}\varphi `$ for the triplet pairing.
The upper critical field is found from the equation on the order parameter which have the different shape for even and odd superconducting states. For determiness we shall consider just the even order parameter states. In this case the order parameter equation for a clean superconductor with singlet pairing is
$`\mathrm{\Delta }^s(\varphi ,p_z,x)`$ $`=g{\displaystyle 𝑑x^{}\frac{dp_y^{}}{2\pi }\underset{\pi /d}{\overset{\pi /d}{}}\frac{dp_z^{}}{2\pi }\psi (\varphi ,p_z)\psi ^{}(\varphi ^{},p_z^{})}`$ (11)
$`\times `$ $`T{\displaystyle \underset{n}{}}{\displaystyle \frac{1}{2}}{\displaystyle \underset{\sigma =\pm 1}{}}G_{\omega _n,\sigma }(p_y^{},p_z^{},x,x^{})G_{\omega _n,\sigma }(p_y^{},p_z^{},x,x^{})\mathrm{\Delta }^s(\varphi ^{},p_z^{},x^{})`$ (12)
and for the triplet pairing
$`\mathrm{\Delta }^t(\varphi ,p_z,x)`$ $`=g{\displaystyle 𝑑x^{}\frac{dp_y^{}}{2\pi }\underset{\pi /d}{\overset{\pi /d}{}}\frac{dp_z^{}}{2\pi }\psi (\varphi ,p_z)\psi ^{}(\varphi ^{},p_z^{})}`$ (13)
$`\times `$ $`T{\displaystyle \underset{n}{}}{\displaystyle \frac{1}{2}}{\displaystyle \underset{\sigma =\pm 1}{}}G_{\omega _n,\sigma }(p_y^{},p_z^{},x,x^{})G_{\omega _n,\sigma }(p_y^{},p_z^{},x,x^{})\mathrm{\Delta }^t(\varphi ^{},p_z^{},x^{})`$ (14)
In presence of the impurities the order parameter of a superconducting singlet pairing state $`\widehat{\mathrm{\Delta }}^s=i\mathrm{\Delta }^s(\varphi ,p_z,x)\widehat{\sigma }_y`$ acquires a self energy part
$$\widehat{\mathrm{\Sigma }}^s(\omega _n,x)=i\mathrm{\Delta }_{\omega _n}^s(x)\widehat{\sigma }_y+\mathrm{\Delta }_{\omega _n}^t(x)\widehat{\sigma }_x$$
(15)
consisting of singlet and triplet components Compare with the paper , where a similar theory in frame of quasiclasical approach has been developed.
For a triplet pairing states with an order parameter $`\widehat{\mathrm{\Delta }}^t=\mathrm{\Delta }^t(\varphi ,p_z,x)\widehat{\sigma }_z`$ the corresponding self energy part
$$\widehat{\mathrm{\Sigma }}^t(\omega _n,x)=\mathrm{\Delta }_{\omega _n}^t(x)\widehat{\sigma }_z+i\stackrel{~}{\mathrm{\Delta }}_{\omega _n}^t(x)\widehat{\sigma }_0$$
(16)
consists of two different triplet components. Here $`\widehat{\sigma }_0`$ is the two dimensional unit matrix. The singlet component of self energy part is absent for chosen basal plane orientation of magnetic field and equal spin triplet pairing with spin directions parallel to conducting layers.
The self-consistency equations for the order parameters and self energy parts have the form (see )
$`\mathrm{\Delta }_{\alpha \beta }^{s,t}`$ $`={\displaystyle 𝑑x^{}\frac{dp_y^{}}{2\pi }\underset{\pi /d}{\overset{\pi /d}{}}\frac{dp_z^{}}{2\pi }V_{\beta \alpha ,\lambda \mu }^{s,t}}`$ (17)
$`\times `$ $`T{\displaystyle \underset{n}{}}G_{\stackrel{~}{\omega }_n}^{\lambda \gamma }(p_y^{},p_z^{},x,x^{})G_{\stackrel{~}{\omega }_n}^{\mu \delta }(p_y^{},p_z^{},x,x^{})[\mathrm{\Delta }_{\gamma \delta }^{s,t}(\varphi ^{},p_z^{},x^{})+\mathrm{\Sigma }_{\gamma \delta }^{s,t}(\stackrel{~}{\omega }_n,x^{})],`$ (18)
$`\mathrm{\Sigma }_{\gamma \delta }^{s,t}`$ $`(\stackrel{~}{\omega }_n,x)=nu^2{\displaystyle 𝑑x^{}\frac{dp_y^{}}{2\pi }\underset{\pi /d}{\overset{\pi /d}{}}\frac{dp_z^{}}{2\pi }}`$ (19)
$`\times `$ $`G_{\stackrel{~}{\omega }_n}^{\gamma \alpha }(p_y^{},p_z^{},x,x^{})G_{\stackrel{~}{\omega }_n}^{\beta \delta }(p_y^{},p_z^{},x,x^{})[\mathrm{\Delta }_{\alpha \beta }^{s,t}(\varphi ^{},p_z^{},x^{})+\mathrm{\Sigma }_{\alpha \beta }^{s,t}(\stackrel{~}{\omega }_n,x^{})],`$ (20)
where $`V_{\beta \alpha ,\lambda \mu }^{s,t}=gg_{\beta \alpha }^{s,t}g_{\lambda \mu }^{s,t+}\psi (\varphi ,p_z)\psi ^{}(\varphi ^{},p_z^{})/2`$, $`\widehat{g}^s=i\widehat{\sigma }_y`$, $`\widehat{g}^t=\widehat{\sigma }_z`$;
$`G_{\stackrel{~}{\omega }_n}^{\lambda \gamma }=G_{\stackrel{~}{\omega }_n,1}(\sigma _0^{\lambda \gamma }+\sigma _z^{\lambda \gamma })/2+G_{\stackrel{~}{\omega }_n,1}(\sigma _0^{\lambda \gamma }\sigma _z^{\lambda \gamma })/2`$.
Here $`\stackrel{~}{\omega }_n=\omega _n+\mathrm{\Gamma }\mathrm{𝑠𝑖𝑔𝑛}\omega _n`$, $`\mathrm{\Gamma }=mnu^2/2d`$ and $`u`$, $`n`$ are impurity potential amplitude and concentration. We will use also a quasiparticle life time $`\tau `$ and a mean free path $`l`$ introduced by means $`\mathrm{\Gamma }=1/2\tau =v_F/2l`$.
The equations (20) are obtained in frame of procedure of averaging over an impurity positions. As it was shown in the paper one can use the field independent value of $`\mathrm{\Gamma }`$ so long
$$\frac{v_F}{\omega _c}>\frac{l}{\sqrt{k_F}l}.$$
(21)
The system of ”the scalar” self-consistency equations for the singlet pairing states following of equations (20) is
$`\mathrm{\Delta }^s`$ $`(\varphi ,p_z,x)=g{\displaystyle 𝑑x^{}\frac{dp_y^{}}{2\pi }\underset{\pi /d}{\overset{\pi /d}{}}\frac{dp_z^{}}{2\pi }\psi (\varphi ,p_z)\psi ^{}(\varphi ^{},p_z^{})}`$ (22)
$`\times `$ $`T{\displaystyle \underset{n}{}}{\displaystyle \frac{1}{2}}{\displaystyle \underset{\sigma =\pm 1}{}}G_{\stackrel{~}{\omega }_n,\sigma }(p_y^{},p_z^{},x,x^{})G_{\stackrel{~}{\omega }_n,\sigma }(p_y^{},p_z^{},x,x^{})[\mathrm{\Delta }^s(\varphi ^{},p_z^{},x^{})+\mathrm{\Delta }_{\stackrel{~}{\omega }_n}^s(x^{})+\sigma \mathrm{\Delta }_{\stackrel{~}{\omega }_n}^t(x^{})],`$ (23)
$`\mathrm{\Delta }`$ $`{}_{\stackrel{~}{\omega }_n}{}^{s}(x)=nu^2{\displaystyle }dx^{}{\displaystyle }{\displaystyle \frac{dp_y}{2\pi }}{\displaystyle \underset{\pi /d}{\overset{\pi /d}{}}}{\displaystyle \frac{dp_z}{2\pi }}{\displaystyle \frac{1}{2}}{\displaystyle \underset{\sigma =\pm 1}{}}G_{\stackrel{~}{\omega }_n,\sigma }(p_y,p_z,x,x^{})G_{\stackrel{~}{\omega }_n,\sigma }(p_y,p_z,x,x^{})`$ (24)
$`\times `$ $`[\mathrm{\Delta }^s(\varphi ,\widehat{p}_z,x^{})+\mathrm{\Delta }_{\stackrel{~}{\omega }_n}^s(x^{})+\sigma \mathrm{\Delta }_{\stackrel{~}{\omega }_n}^t(x^{})],`$ (25)
$`\mathrm{\Delta }`$ $`{}_{\stackrel{~}{\omega }_n}{}^{t}(x)=nu^2{\displaystyle }dx^{}{\displaystyle }{\displaystyle \frac{dp_y}{2\pi }}{\displaystyle \underset{\pi /d}{\overset{\pi /d}{}}}{\displaystyle \frac{dp_z}{2\pi }}{\displaystyle \frac{1}{2}}{\displaystyle \underset{\sigma =\pm 1}{}}G_{\stackrel{~}{\omega }_n,\sigma }(p_y,p_z,x,x^{})G_{\stackrel{~}{\omega }_n,\sigma }(p_y,p_z,x,x^{})`$ (26)
$`\times `$ $`[\sigma (\mathrm{\Delta }^s(\varphi ,\widehat{p}_z,x^{})+\mathrm{\Delta }_{\stackrel{~}{\omega }_n}^s(x^{}))+\mathrm{\Delta }_{\stackrel{~}{\omega }_n}^t(x^{})].`$ (27)
For the triplet pairing case the corresponding equations are
$`\mathrm{\Delta }^t`$ $`(\varphi ,p_z,x)=g{\displaystyle 𝑑x^{}\frac{dp_y^{}}{2\pi }\underset{\pi /d}{\overset{\pi /d}{}}\frac{dp_z^{}}{2\pi }\psi (\varphi ,p_z)\psi ^{}(\varphi ^{},p_z^{})}`$ (28)
$`\times `$ $`T{\displaystyle \underset{n}{}}{\displaystyle \frac{1}{2}}{\displaystyle \underset{\sigma =\pm 1}{}}G_{\stackrel{~}{\omega }_n,\sigma }(p_y^{},p_z^{},x,x^{})G_{\stackrel{~}{\omega }_n,\sigma }(p_y^{},p_z^{},x,x^{})[\mathrm{\Delta }^t(\varphi ^{},p_z^{},x^{})+\mathrm{\Delta }_{\stackrel{~}{\omega }_n}^t(x^{})i\sigma \stackrel{~}{\mathrm{\Delta }}_{\stackrel{~}{\omega }_n}^t(x^{})],`$ (29)
$`\mathrm{\Delta }`$ $`{}_{\stackrel{~}{\omega }_n}{}^{t}(x)=nu^2{\displaystyle }dx^{}{\displaystyle }{\displaystyle \frac{dp_y}{2\pi }}{\displaystyle \underset{\pi /d}{\overset{\pi /d}{}}}{\displaystyle \frac{dp_z}{2\pi }}{\displaystyle \frac{1}{2}}{\displaystyle \underset{\sigma =\pm 1}{}}G_{\stackrel{~}{\omega }_n,\sigma }(p_y,p_z,x,x^{})G_{\stackrel{~}{\omega }_n,\sigma }(p_y,p_z,x,x^{})`$ (30)
$`\times `$ $`[\mathrm{\Delta }^t(\varphi ,p_z,x^{})+\mathrm{\Delta }_{\stackrel{~}{\omega }_n}^t(x^{})i\sigma \stackrel{~}{\mathrm{\Delta }}_{\stackrel{~}{\omega }_n}^t(x_1)],`$ (31)
$`\stackrel{~}{\mathrm{\Delta }}`$ $`{}_{\stackrel{~}{\omega }_n}{}^{t}(x)=nu^2{\displaystyle }dx^{}{\displaystyle }{\displaystyle \frac{dp_y}{2\pi }}{\displaystyle \underset{\pi /d}{\overset{\pi /d}{}}}{\displaystyle \frac{dp_z}{2\pi }}{\displaystyle \frac{1}{2}}{\displaystyle \underset{\sigma =\pm 1}{}}G_{\stackrel{~}{\omega }_n,\sigma }(p_y,p_z,x,x^{})G_{\stackrel{~}{\omega }_n,\sigma }(p_y,p_z,x,x^{})`$ (32)
$`\times `$ $`(i\sigma (\mathrm{\Delta }^t(\varphi ,p_z,x^{})+\mathrm{\Delta }_{\stackrel{~}{\omega }_n}^t(x^{}))+\stackrel{~}{\mathrm{\Delta }}_{\stackrel{~}{\omega }_n}^t(x^{})].`$ (33)
Taking into account that in common the paramagnetic limit of superconductivity $`H_p`$ is much higher than an orbital upper critical field we shall rest the general problem of influence of paramagnetism on superconductivity for a future investigations. In the absence of the Pauli paramagnetic interaction the equations (23)-(33) are greatly simplified and we have the system of two equations with equivalent structure for singlet and triplet pairing states
$`\mathrm{\Delta }^{s,t}(\varphi ,p_z,x)`$ $`=g{\displaystyle 𝑑x^{}\frac{dp_y^{}}{2\pi }\underset{\pi /d}{\overset{\pi /d}{}}\frac{dp_z^{}}{2\pi }\psi (\varphi ,p_z)\psi ^{}(\varphi ^{},p_z^{})}`$ (34)
$`\times `$ $`T{\displaystyle \underset{n}{}}G_{\stackrel{~}{\omega }_n}(p_y^{},p_z^{},x,x^{})G_{\stackrel{~}{\omega }_n}(p_y^{},p_z^{},x,x^{})(\mathrm{\Delta }^{s,t}(\varphi ^{},p_z^{},x^{})+\mathrm{\Delta }_{\stackrel{~}{\omega }_n}^{s,t}(x^{})),`$ (35)
$`\mathrm{\Delta }_{\stackrel{~}{\omega }_n}^{s,t}(x)`$ $`=nu^2{\displaystyle 𝑑x^{}\frac{dp_y}{2\pi }\underset{\pi /d}{\overset{\pi /d}{}}\frac{dp_z}{2\pi }}`$ (36)
$`\times `$ $`G_{\stackrel{~}{\omega }_n}(p_y,p_z,x,x^{})G_{\stackrel{~}{\omega }_n}(p_y,p_z,x,x^{})(\mathrm{\Delta }^{s,t}(\varphi ,\widehat{p}_z,x^{})+\mathrm{\Delta }_{\stackrel{~}{\omega }_n}^{s,t}(x^{})).`$ (37)
Below we shall omit the supercripts $`s,t`$ using the common notation $`\mathrm{\Delta }^{s,t}(\varphi ,p_z,x)=\mathrm{\Delta }(\varphi ,p_z,x)`$,
$`\mathrm{\Delta }_{\stackrel{~}{\omega }_n}^{s,t}(x)=\mathrm{\Delta }_{\stackrel{~}{\omega }_n}(x)`$ for the order parameters and the self energy parts both in singlet and in triplet case.
On this stage it is useful to note that the normal metal electron Green function which we should find as a solution of the equation (6) depends of $`p_y`$ only through its square. Hence for some unconventional superconducting phases like $`\mathrm{\Delta }(\varphi ,\widehat{p}_z,x)=\sqrt{2}(\mathrm{sin}^2\varphi \mathrm{cos}^2\varphi )\mathrm{\Delta }(x)`$ for singlet pairing or $`\mathrm{\Delta }(\varphi ,\widehat{p}_z,x)=\sqrt{2}\mathrm{cos}\varphi \mathrm{\Delta }(x)`$ for triplet pairing, the following property takes place
$$\frac{dp_y}{2\pi }G_{\stackrel{~}{\omega }_n}(p_y,p_z,x,x^{})G_{\stackrel{~}{\omega }_n}(p_y,p_z,x,x^{})(\mathrm{\Delta }(\varphi ,\widehat{p}_z,x)=0.$$
(38)
We will discuss further the only unconventional superconducting states obeying the property (38). For such the states the self energy part is equal to zero and we deal only with the order parameter equation.
$`\mathrm{\Delta }(\varphi ,p_z,x)`$ $`=g{\displaystyle 𝑑x^{}\frac{dp_y^{}}{2\pi }\underset{\pi /d}{\overset{\pi /d}{}}\frac{dp_z^{}}{2\pi }\psi (\varphi ,p_z)\psi ^{}(\varphi ^{},p_z^{})}`$ (39)
$`\times `$ $`T{\displaystyle \underset{n}{}}G_{\stackrel{~}{\omega }_n}(p_y^{},p_z^{},x,x^{})G_{\stackrel{~}{\omega }_n}(p_y^{},p_z^{},x,x^{})\mathrm{\Delta }(\varphi ^{},p_z^{},x^{}).`$ (40)
The equality (38) is not valid for conventional superconducting state as well for many unconventional superconducting states where, as the consequence, there are nonzero selfenergy parts. In isotropic conventional superconducting state it prevents a supression of the superconductivity by the ordinary impurities. For unconventional superconducting states the presence of the self energy leads just to the mathematical complifications and does not change qualitatively the main results.
The only difference of (40) from the pure case consists of change $`\omega _n\stackrel{~}{\omega }_n`$ and one can use the expresion for the normal metal electron Green function found in the paper . It has nonzero value in the regions determined by the inequality $`\omega _n(xx_1)0`$ for positive value of $`x`$ component of electron momentum
($`\alpha =1`$) and by the inequality $`\omega _n(xx_1)0`$ for negative values of $`x`$ component momentum ( $`\alpha =1`$), where it is defined as
$`G_{\stackrel{~}{\omega }_n}`$ $`(\varphi ,p_z,x,x_1)={\displaystyle \frac{i\mathrm{𝑠𝑖𝑔𝑛}\omega _n}{v_F\mathrm{sin}\varphi }}\mathrm{exp}\left[{\displaystyle \frac{\alpha \stackrel{~}{\omega }_n(xx_1)}{v_F\mathrm{sin}\varphi }}\right]\mathrm{exp}[i\alpha p_F(xx_1)\mathrm{sin}\varphi ]`$ (41)
$`\times `$ $`\mathrm{exp}`$ $`\left\{{\displaystyle \frac{i\alpha \lambda }{\mathrm{sin}\varphi }}\mathrm{sin}\left[{\displaystyle \frac{\omega _c(xx_1)}{2v_F}}\right]\mathrm{cos}\left[p_zd{\displaystyle \frac{\omega _c(x+x_1)}{2v_F}}\right]\right\}.`$ (42)
Here $`\lambda =4t/\omega _c`$.
This expression is valid under the condition
$$|\mathrm{sin}\varphi |>\sqrt{\frac{t}{ϵ_F}}.$$
(43)
The disregard of the small intervals of the angle $`\varphi `$ where $`|\mathrm{sin}\varphi |<(t/ϵ_F)^{1/2}`$ means that the only open trajectories of quasiparticles on the Fermi surface are taken into account The numerical calculation taking into account both open and slosed trajectories, performed for clean case in the paper, just confirms the qualitative results of the article have been obtained in neglect of closed trajectories. They give the main singular part to the kernels of the integral equation (40).
The substitution of (42) to the equation (43) gives after summation over the Matsubara frequency
$`\mathrm{\Delta }(x)`$ $`=\stackrel{~}{g}{\displaystyle \underset{|xx_1|>asin\varphi }{}}𝑑x_1{\displaystyle \underset{0}{\overset{\pi }{}}}{\displaystyle \frac{d\varphi }{2\pi v_F\mathrm{sin}\varphi }}{\displaystyle \underset{\pi }{\overset{\pi }{}}}{\displaystyle \frac{d(p_zd)}{2\pi }}\psi ^{}(\varphi ,p_z)\psi (\varphi ,p_z){\displaystyle \frac{2\pi T\mathrm{exp}\left[\frac{|xx_1|}{l\mathrm{sin}\varphi }\right]}{\mathrm{sinh}\left[\frac{2\pi T|xx_1|}{v_F\mathrm{sin}\varphi }\right]}}`$ (44)
$`\times `$ $`\mathrm{exp}\left\{{\displaystyle \frac{2i\lambda }{\mathrm{sin}\varphi }}\mathrm{sin}\left[{\displaystyle \frac{\omega _c(xx_1)}{2v_F}}\right]\mathrm{sin}(p_zd)\mathrm{sin}\left[{\displaystyle \frac{\omega _c(x+x_1)}{2v_F}}\right]\right\}\mathrm{\Delta }(x_1).`$ (45)
Here $`\stackrel{~}{g}=gm/4\pi d`$ is dimensionless constant of pairing interaction, $`a`$ is the small distances cutoff. The equation (45) is the basic equation of the paper. The properties of its solution we shall discuss in the next Section.
## III The upper critical field
Let us make an unessential simplification and consider $`p_z`$ independent superconducting states determined by functions $`\psi (\varphi )`$ normalized as follows
$$\underset{0}{\overset{\pi }{}}\frac{d\varphi }{\pi }\psi ^{}(\varphi )\psi (\varphi )=1.$$
In this case one can perform the integration over $`p_z`$ in (45):
$`\mathrm{\Delta }(x)`$ $`=\stackrel{~}{g}{\displaystyle \underset{|xx_1|>asin\varphi }{}}𝑑x_1{\displaystyle \underset{0}{\overset{\pi }{}}}{\displaystyle \frac{d\varphi }{2\pi v_F\mathrm{sin}\varphi }}\psi ^{}(\varphi )\psi (\varphi ){\displaystyle \frac{2\pi T\mathrm{exp}\left[\frac{|xx_1|}{l\mathrm{sin}\varphi }\right]}{\mathrm{sinh}\left[\frac{2\pi T|xx_1|}{v_F\mathrm{sin}\varphi }\right]}}`$ (46)
$`\times `$ $`_0\left\{{\displaystyle \frac{2\lambda }{\mathrm{sin}\varphi }}\mathrm{sin}\left[{\displaystyle \frac{\omega _c(xx_1)}{2v_F}}\right]\mathrm{sin}\left[{\displaystyle \frac{\omega _c(x+x_1)}{2v_F}}\right]\right\}\mathrm{\Delta }(x_1).`$ (47)
Here $`_0(\mathrm{})`$ is the Bessel function. For the further purposes it is not important to work with the equation (45) or with (47) and for the determiness we will operate with the latter.
The equation (47) in the absence of a magnetic field
$$1=\stackrel{~}{g}\underset{\frac{2\pi aT}{v_F}}{\overset{\mathrm{}}{}}\frac{dz}{\mathrm{sinh}z}\mathrm{exp}\left(\frac{z}{2\pi T\tau }\right)$$
(48)
determines of the critical temperature $`T_c`$, which is expressed from here through the critical temperature $`T_{c0}`$ in a perfect crystal without impurities $`l=\mathrm{}`$
$$T_{c0}=\frac{v_F}{\pi a}\mathrm{exp}\left(\frac{1}{\stackrel{~}{g}}\right)$$
(49)
by means of well known relation
$$\mathrm{ln}\frac{T_c}{T_{c0}}=\psi \left(\frac{1}{2}\right)\psi \left(\frac{1}{2}+\frac{1}{4\pi \tau T_c}\right).$$
(50)
For the critical temperatures $`T_cT_{c0}`$ the suppresion of the superconductivity by the impurities is:
$$T_c=T_{c0}\frac{\pi }{8\tau }.$$
(51)
One can also point out the condition of complete suppresion of the superconductivity
$$\tau =\tau _c=\frac{\gamma }{\pi T_{c0}},$$
(52)
where $`\mathrm{ln}\gamma =C=0,577\mathrm{}`$ is the Euler constant.
The behavior of the upper critical field is determined by the relationship of the three spacial scales: $`v_F/2\pi T`$, $`l`$ and $`v_F/\omega _c`$. For the temperatures near to the zero field critical temperature $`TT_c(H=0)`$ the upper critical field $`H_{c2}(T)`$, or $`\omega _{c2}(T)=eH_{c2}(T)v_Fd/c`$ tends to zero and the inequality
$$min\{\frac{v_F}{2\pi T},l\}<\frac{v_F}{\omega _{c2}(T)}$$
(53)
always presents <sup>§</sup><sup>§</sup>§Near $`T_c`$ there is also formal solution of (47) with $`\omega _{c2}>t`$ ( see ). This solution, however, is related to the region of magnetic fields, out the region of validity of the present theory (21). .
At the temperature decrease the upper critical field increase, but so long the inequality (53) takes place the essential interval of integration over $`(xx_1)`$ in (47) determined by the $`min(v_F/2\pi T,l)`$. Hence, one can use the smallness of $`(xx_1)v_F/\omega _c`$ in the kernel of the equation (47):
$$\mathrm{\Delta }(x)=\stackrel{~}{g}\underset{|xx_1|>asin\varphi }{}𝑑x_1\underset{0}{\overset{\pi }{}}\frac{d\varphi }{2\pi v_F\mathrm{sin}\varphi }\psi ^{}(\varphi )\psi (\varphi )\frac{2\pi T\mathrm{exp}\left[\frac{|xx_1|}{l\mathrm{sin}\varphi }\right]}{\mathrm{sinh}\left[\frac{2\pi T|xx_1|}{v_F\mathrm{sin}\varphi }\right]}_0\left\{\frac{4t(xx_1)}{v_F\mathrm{sin}\varphi }\mathrm{sin}\frac{\omega _cx}{v_F}\right\}\mathrm{\Delta }(x_1).$$
(54)
The solution of this equation gives the correct value of $`\omega _{c2}(T)`$ so long the inequality (53) is truth. By the substitution of the $`\omega _{c2}(T)`$ obtained from this equation to the inequality (53) one can establish the upper limit of the superconductor purity at which the $`H_{c2}(T)`$ found from this equation represent the correct value of the upper critical field up to zero temterature.
For the pure enough samples and low enough temperatures it can be turn out that the opposite to the (53) relationship
$$\frac{v_F}{\omega _{c2}(T)}<min\{\frac{v_F}{2\pi T},l\}$$
(55)
breaking the correctness of transfering from (47) to (54) is realized. It should be noted that ultraclean case demands a special investigation (see condition (21) and discussion in the paper ). One can claim however, that at the temperatures
$$\frac{v_F}{2\pi l}<T<\frac{\omega _c(T)}{2\pi }$$
the magnetic field dependence of the critical temperature in the equation (47) starts disappear or, in other words, the tendency to the divergency of the upper critical field pointed out in the paper appears.
To solve of the equation (54) at arbitrary temperature and purity is possible only numerically. Here we shall discuss the case when it allows the analitical solution. If the legth scale $`\xi `$, on which the function $`\mathrm{\Delta }(x)`$ is lokalized, is larger than the essential distance of integration over $`(xx_1)`$
$$\xi >min\{\frac{v_F}{2\pi T},l\},$$
(56)
then one can expand $`\mathrm{\Delta }(x_1)\mathrm{\Delta }(x)+\mathrm{\Delta }^{}(x)(xx_1)+\mathrm{\Delta }^{\prime \prime }(x)(xx_1)^2/2`$ under the integral in (54). Taking into consideration that under this condition the argument of Bessel function turns to be small even on the upper boundary of effective interval of the integration over $`(xx_1)`$:
$$\frac{t\omega _{c2}(T)\xi min\{\frac{v_F}{2\pi T},l\}}{v_F^2}\frac{tmin\{\frac{v_F}{2\pi T},l\}}{ϵ_F\xi }<1,$$
(57)
one can also expand Bessel function $`_0(x)1x^2/4`$. As the result we get the differential equation
$$\left(\mathrm{ln}\frac{T_{c0}}{T}\psi \left(\frac{1}{2}+\frac{1}{4\pi \tau T}\right)+\psi \left(\frac{1}{2}\right)\right)\mathrm{\Delta }(x)=\frac{CI(\alpha )}{2}\left(\frac{v_F}{2\pi T}\right)^2\mathrm{\Delta }^{\prime \prime }(x)+I(\alpha )\left(\frac{t\omega _{c2}(T)x}{\pi v_FT}\right)^2\mathrm{\Delta }(x),$$
(58)
where $`\alpha =(2\pi T\tau )^1`$,
$$I(\alpha )=\underset{0}{\overset{\mathrm{}}{}}\frac{z^2dz}{\mathrm{sinh}z}\mathrm{exp}(\alpha z)=4\underset{n=0}{\overset{\mathrm{}}{}}\frac{1}{(2n+1+\alpha )^3},$$
(59)
$$C_\psi =\underset{0}{\overset{\pi }{}}\frac{d\varphi }{\pi }\psi ^{}(\varphi )\psi (\varphi )\mathrm{sin}^2\varphi .$$
(60)
In pure case $`\alpha \alpha _c=(2\pi T_c\tau )^11`$ the inequality (56) is valid only in vicinity of the critical temperature. Puting $`T=T_c`$ in the right hand side and taking its lowest eigen value we obtain
$$\mathrm{ln}\frac{T_{c0}}{T}\psi \left(\frac{1}{2}+\frac{1}{4\pi \tau T}\right)+\psi \left(\frac{1}{2}\right)=\frac{\sqrt{C_\psi }I(\alpha _c)t\omega _{c2}(T)}{2\sqrt{2}\pi ^2T_c^2}.$$
(61)
Summing this equation with equation (50) and leaving only linear on $`TT_c`$ and on the impurity concentration terms we get
$$\omega _{c2}(T)=\frac{ev_Fd}{c}H_{c2}(T)=\frac{4\sqrt{2}\pi ^2}{7\zeta (3)\sqrt{C_\psi }t}\left(T_{c0}\beta \frac{\pi }{8\tau }\right)(T_cT).$$
(62)
Here the coefficient
$$\beta =2\frac{90\zeta (4)}{7\pi ^2\zeta (3)}0.83,$$
(63)
shows that, the slop of $`H_{c2}(T)`$ at $`T=T_c`$ decreases with the increase of the impurity concentration somewhat slower than $`T_c`$ itself (see eqn. (51)) .
The equation (54) does not contain any divergency of $`H_{c2}(T)`$ and the expression (62) which is valid in Ginzburg-Landau region can be used at arbitrary temperature as the estimate of upper critical field from above. Hence, to establish the limits of a sample purity, at which the equation (54) works, one may substitute (62) at zero temperature into the inequality (53). Omitting the numerical factor of the order of unity we have
$$l<\frac{t}{T_c}\xi _{ab},$$
(64)
where $`\xi _{ab}=v_F/2\pi T_c`$ is the basal plane coherence length. We see, that there is the good reserve in sample purity in the limits of which one can not expect low temperature divergency of the upper critical field. In the $`Sr_2RuO_4`$ the mean free path should be approximately 10 times larger than the basal plane coherence length $`\xi _0`$ to go out limit (64).
Let us consider now the dirty case: $`T_cT_{c0}`$, $`\alpha 1`$ and $`I(\alpha )\alpha ^2`$ allowing analytical solution for $`H_{c2}(T)`$ at arbitrary temperature. Taking the lowest eigen value of the equation (58) we get
$$\mathrm{ln}\frac{T_{c0}}{T}\psi \left(\frac{1}{2}+\frac{1}{4\pi \tau T}\right)+\psi \left(\frac{1}{2}\right)=\sqrt{2C_\psi }\tau ^2t\omega _{c2}(T).$$
(65)
Summing of this equation with equation (50) yields
$$\omega _{c2}(T)=\frac{\mathrm{ln}\frac{T_c}{T}\psi \left(\frac{1}{2}+\frac{1}{4\pi \tau T}\right)+\psi \left(\frac{1}{2}+\frac{1}{4\pi \tau T_c}\right)}{\sqrt{2C_\psi }\tau ^2t}.$$
(66)
This expression is correct at any temperature. One can rewrite it approximatively in more simple form
$$\omega _{c2}(T)=\frac{\sqrt{2}\pi ^2}{3\sqrt{C_\psi }t}(T_c^2T^2),$$
(67)
where
$$T_c=\frac{1}{\pi \tau }\left(\frac{3}{2}\mathrm{ln}\frac{\pi T_{c0}\tau }{\gamma }\right)^{1/2}.$$
(68)
## IV Conclusion
The system of linear integral equations for the order parameter of conventional and unconventional superconducting state in a layered crystals with impurities under magnetic field parallel to the conducting layers is derived. It is shown that so long the purity of the samples does not exceed high enough level (64) there is no tendency to the low temperature divergency of the upper critical field. The analytical solution of the equations in the clean (Ginzburg-Landau region) and dirty (arbitrary temperature) limits are presented.
## Acknowledgements
I express my best gratitudes to Dr. Manfred Sigrist and all the members of Condensed matter theory group at Yukawa Institute for Theoretical physics for their kind hospitality and interest to my work during my stay in Kyoto autumn 1999 where the significant part of this work have been completed.
I also indebted to Dr. Yoshiteru Maeno and his collaborators having stimulated my interest to the problem of upper critical field in layered superconductors.
I appreciate the valuable remarques of prof. M.Walker, prof. P.Nozieres and prof. Yu.A.Bychkov have resulted in the improvement in the initial text of the article. |
warning/0002/astro-ph0002506.html | ar5iv | text | # The properties of cooling flows in X-ray luminous clusters of galaxies
## 1 Introduction
X-ray observations of clusters of galaxies show that in the central regions of most clusters the cooling time of the intracluster gas is significantly less than a Hubble time (e.g. Edge, Stewart & Fabian 1992; White, Jones & Forman 1997; Peres et al. 1998). This cooling is thought to lead to a slow net inflow of material towards the cluster centre; a process known as a cooling flow (see Fabian 1994 for a review). X-ray imaging data show that gas typically ‘cools out’ throughout the central few tens to hundreds of kpc in clusters, with $`\dot{M}(r)r`$, where $`\dot{M}(r)`$ is the integrated mass deposition rate within radius $`r`$ (e.g. Thomas, Fabian & Nulsen 1987). Spatially resolved X-ray spectroscopy has confirmed the presence of distributed cooling gas in cooling flows, with a spatial distribution and luminosity in good agreement with the predictions from the imaging data and cooling-flow models (e.g. Allen & Fabian 1997).
For some years, the primary uncertainty with the standard model for cooling flows was the fate of the cooled matter (e.g. see Fabian, Nulsen & Canizares 1991). However, the discovery of large column densities of intrinsic X-ray absorbing material associated with cooling flows observed with the Solid State Spectrometer (SSS) on the Einstein Observatory (White et al. 1991; Johnstone et al. 1992) opened one interesting possibility. Follow-up spatially-resolved X-ray spectroscopy with the Position Sensitive Proportional Counter (PSPC) on ROSAT (e.g. Allen et al. 1993; Irwin & Sarazin 1995; Allen & Fabian 1997) confirmed the presence of intrinsic X-ray absorbing material in cooling-flow clusters and further showed this material to be distributed throughout, but centrally-concentrated within, the cooling flows. The X-ray data thus identify the intrinsic X-ray absorbing material as a plausible sink for the cooled gas deposited by the cooling flows.
In this paper we examine the X-ray properties of the cooling flows in a sample of 30 of the most X-ray luminous ($`L_{\mathrm{Bol}}>10^{45}`$ $`\mathrm{erg}\mathrm{s}^1`$) clusters of galaxies known. Using ASCA spectra and ROSAT High Resolution Imager (HRI) data, we present independent determinations of the mass deposition rates in the cooling flows and measure the column densities of intrinsic X-ray absorbing material associated with these systems. We show that the cooling flow model provides a consistent description of the X-ray imaging and spectral data and that the observed masses of intrinsic X-ray absorbing material are in reasonable agreement with the masses expected to have been accumulated by the cooling flows over their lifetimes. We also discuss the evidence for mass deposition from the cooling flows in other wavebands.
This is the final paper in a series that has examined the impact of cooling flows on mass measurements (Allen 1998), the $`kT_\mathrm{X}L_{\mathrm{Bol}}`$ relation (Allen & Fabian 1998a) and metallicity measurements (Allen & Fabian 1998b) for X-ray luminous clusters. In this paper, we describe the data reduction and analysis procedures used in these works. A large number of other, previous studies have also analysed one or more of data sets included here (see references in the papers listed above), although these studies have not, in general, examined the properties of the cooling flows in the clusters. Exceptions are the deprojection analysis of Peres et al. (1998), and the combined ASCA/ROSAT studies of Allen et al. (1996), Schindler et al. (1997), Böhringer et al. (1998) and Rizza et al. (1998).
The cosmological parameters $`H_0`$=50 $`\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$, $`\mathrm{\Omega }=1`$ and $`\mathrm{\Lambda }=0`$ are assumed throughout.
## 2 Observations and data reduction
### 2.1 Sample selection
Our sample was identified from clusters contained in the ROSAT X-ray Brightest Abell-type Cluster Sample (XBACS; Ebeling et al. 1996) and ROSAT Brightest Cluster Sample (BCS; Ebeling et al. 1998) with X-ray luminosities in the $`0.12.4`$ keV band exceeding $`10^{45}`$ $`\mathrm{erg}\mathrm{s}^1`$. From this initial list, we selected for study those targets with ASCA X-ray spectra and ROSAT High Resolution Imager (HRI) data available on the Goddard Space Flight Centre (GSFC) public archive as of 1997 July 15. We supplemented this sample with a number of other, southern X-ray luminous clusters, known to exhibit strong gravitational lensing effects (PKS0745-191, RXJ1347.5-1145, MS2137.3-2353, AC114) and the distant, luminous cooling-flow clusters Abell 1068, 1704 and IRAS 09104+4109.
The primary goal of this project was to investigate the X-ray properties of cooling flows in X-ray luminous clusters and their effects on the integrated X-ray properties of their host systems. We have not included either the Perseus (Abell 426) or Coma (Abell 1656) clusters in our study, although these systems met our selection criteria, since both are too close ($`z=0.0183`$ and $`0.0232`$, respectively) to allow their integrated cluster properties to be studied using the techniques used in this paper. Detailed analyses of these clusters and other nearby cooling flows are discussed by Allen et al. (1999).
The clusters Abell 370, 115 and 1758 also met the selection criteria but were not included in our study since ROSAT HRI images show them not be single, coherent structures but to consist of a number of merging subunits. The integrated X-ray spectra for such clusters will relate to the virial properties of the individual subclusters rather than the systems as a whole, and the assumptions of spherical symmetry and hydrostatic equilibrium required by the X-ray modeling will not apply. The X-ray images for the other clusters included in the sample do not exhibit any dramatic substructure that would suggest these assumptions to be invalid. Our final sample consists of 30 clusters, spanning the redshift range $`0.056<z<0.451`$, with a mean redshift of 0.21.
### 2.2 The ASCA observations
The ASCA (Tanaka, Inoue & Holt 1994) observations were made over a three-and-a-half year period between 1993 April and 1996 December. The ASCA X-ray Telescope array (XRT) consists of four nested-foil telescopes, each focussed onto one of four detectors; two X-ray CCD cameras, the Solid-state Imaging Spectrometers (SIS; S0 and S1), and two Gas scintillation Imaging Spectrometers (GIS; G2 and G3). The XRT provides a spatial resolution of $`3`$ arcmin Half Power Diameter (HPD) in the energy range $`0.312`$ keV. The SIS detectors provide good spectral resolution \[$`\mathrm{\Delta }E/E=0.02(E/5.9\mathrm{keV})^{0.5}`$\] over a $`22\times 22`$ arcmin<sup>2</sup> field of view. The GIS detectors provide poorer energy resolution \[$`\mathrm{\Delta }E/E=0.08(E/5.9\mathrm{keV})^{0.5}`$\] but cover a larger circular field of view of $`50`$ arcmin diameter.
For our analysis we have used the screened event lists from the rev1 processing of the data sets available on the GSFC ASCA archive (for a detailed description of the rev1 processing see the GSFC ASCA Data Reduction Guide, published by GSFC.) The ASCA data were reduced using the FTOOLS software (version 3.6) issued by GSFC, from within the XSELECT environment (version 1.3). Further data-cleaning procedures as recommended in the ASCA Data Reduction Guide, including appropriate grade selection, gain corrections and manual screening based on the individual instrument light curves, were followed. A summary of the ASCA observations, including the individual exposure times after all screening procedures were carried out, is provided in Table 1.
Spectra were extracted from all four ASCA detectors (except in those few cases where the S1 data were lost due to saturation problems caused by flickering pixels in the CCDs). The spectra were extracted from circular regions, centred on the peaks of the X-ray emission. For the SIS data, the radii of the regions used were selected to minimize the number of chip boundaries crossed (thereby minimizing the systematic uncertainties introduced by such crossings) whilst covering as large a region of the clusters as possible. Data from the regions between the chips were masked out and excluded. The final extraction radii for the SIS data are summarized in Table 2. Also included in that Table are the chip modes used in the observations (whether 1,2 or 4 chip mode) and the number of chips from which the extracted data were drawn. For the GIS data a constant extraction radius of 6 arcmin was used. \[We note that for Abell 2142, the 2 arcmin (3 arcmin) radius region surrounding the X-ray bright Seyfert-1 galaxy 1556+274 (offset by $`4`$ arcmin from the X-ray centre of the cluster) was masked out and excluded from the analysis of the SIS (GIS) data.\]
For the GIS observations, and SIS observations of clusters in regions of low Galactic column density ($`N_\mathrm{H}<5\times 10^{20}`$ atom cm<sup>-2</sup>), background subtraction was carried out using the ‘blank sky’ observations of high Galactic latitude fields compiled during the performance verification stage of the ASCA mission. For such data sets, the blank-sky observations provide a reasonable representation of the cosmic and instrumental backgrounds in the detectors. The background data were screened and grade selected in the same manner as the target observations and background spectra were extracted from the same regions of the detectors as the cluster spectra. For the SIS observations of clusters in directions of higher Galactic column density, background spectra were extracted from regions of the chips that were relatively free from foreground cluster emission.
For the SIS data, response matrices were generated using the FTOOLS SISRMG software. Where the spectra covered more than one chip, response matrices were created for each chip, which were then combined to form a counts-weighted mean matrix. For the GIS analysis, the response matrices issued by GSFC on 1995 March 6 were used.
### 2.3 The ROSAT observations
The ROSAT HRI observations were carried out between 1991 November and 1995 June. The HRI provides a $`5`$ arcsec (FWHM) X-ray imaging facility covering a $`40\times 40`$ arcmin<sup>2</sup> field of view (David et al. 1996). Reduction of the data was carried out with the Starlink ASTERIX package. X-ray images were created on a $`2\times 2`$ arcsec<sup>2</sup> pixel scale, from which centres for the cluster X-ray emission were determined. Where more than one observation of a source was made, a mosaic was constructed from the individual observations. For the cooling flow and intermediate clusters, the X-ray centres were identified from the peaks of the X-ray surface brightness distributions (which are easily determined from the HRI images). For the non-cooling flow clusters, the X-ray emission is not as sharply-peaked and for these systems we have identified the X-ray centres with the results from iterative determinations of the centroids within a circular aperture of fixed radius. For Abell 665, 2163 and AC114 a 1 arcminute radius aperture was used. For Abell 2744, 773, 2218 and 2219 a 2 arcmin aperture was adopted, and for Abell 520 a 3 arcmin aperture was used. For two of the clusters included in the sample, Abell 665 and 1413, ROSAT HRI images were not available at the time of writing and for these clusters PSPC imaging data have been used instead. A summary of the ROSAT observations and the X-ray centers for the clusters is given in Table 3.
### 2.4 Classification of clusters as cooling flow (CF) and non-cooling flow (NCF) systems
For the purposes of this paper, we have classified the clusters into subsamples of cooling-flow (hereafter CF) and non-cooling flow (NCF) systems. CFs are those clusters for which the upper (90 percentile) limit to the central cooling time, as determined from the deprojection analysis of the ROSAT HRI X-ray images (Section 4), is less than $`10^{10}`$ yr. (The ‘central’ cooling time is the mean cooling time of the cluster gas in the innermost bin included in the deprojection analysis, which is of variable size. The use of a fixed physical size of 100 kpc for the central bin leads to similar results; Section 4.) Using this simple classification scheme we identify 21 CFs and 9 NCFs in our sample. The mean redshift for the subsamples of both CF and NCF clusters is $`\overline{z}=0.21`$. The fraction of CFs in our sample is $`70`$ per cent, the same as that determined by Peres et al. (1998) from a study of a larger, flux-limited sample of clusters (primarily at lower redshift; $`\overline{z}=0.056`$) using a similar classification scheme.
## 3 Spectral Analysis of the ASCA data
### 3.1 The spectral models
The modeling of the X-ray spectra has been carried out using the XSPEC spectral fitting package (version 9.0; Arnaud 1996). For the SIS data, only counts in pulse height analyser (PHA) channels corresponding to energies between 0.6 and 10.0 keV were included in the analysis (the energy range over which the calibration of the SIS instruments is best-understood). For the GIS data only counts in the energy range $`1.010.0`$ keV were used. The spectra were grouped before fitting to ensure a minimum of 20 counts per PHA channel, allowing $`\chi ^2`$ statistics to be used.
The spectra have been modeled using the plasma codes of Kaastra & Mewe (1993; incorporating the Fe L calculations of Liedahl, Osterheld & Goldstein 1995, in XSPEC version 9.0) and the photoelectric absorption models of Balucinska-Church & McCammon (1992). The data from all four ASCA detectors were analysed simultaneously (except where the S1 data were lost due to chip saturation problems) with the fit parameters linked to take the same values across the data sets. The exceptions to this were the emission measures of the ambient cluster gas in the four detectors which, due to the different extraction radii used and residual uncertainties in the flux calibration of the instruments, were allowed to fit independently.
The spectra were examined with a series of spectral models. Model A, consisted of an isothermal plasma in collisional equilibrium, at the optically-determined redshift for the cluster, and absorbed by the nominal Galactic column density (Dickey & Lockman 1990). The free parameters in this model were the temperature ($`kT`$) and metallicity ($`Z`$) of the plasma and the emission measures in the four detectors. (The metallicities are measured relative to the solar values of Anders & Grevesse (1989), with the different elements assumed to be present in their solar ratios.) The second model, model B, was the same as model A but with the absorbing column density $`(N_\mathrm{H})`$ also included as a free parameter in the fits. The third model, model C, was the same as model A but with an additional component explicitly accounting for the emission from the cooling flows. The material in the cooling flows is assumed to cool at constant pressure from the ambient cluster temperature, following the prescription of Johnstone et al. (1992). The normalization of the cooling-flow component was parameterized in terms of a mass deposition rate, $`\dot{M}_\mathrm{S}`$, which was a free parameter in the fits. The mass deposition rate is linked to take the same value in all four detectors, scaled by a normalization factor proportional to the total flux measured in that detector. The metallicity of the cooling gas was assumed to be equal to that of the ambient ICM. The emission from the cooling flows was also assumed to be absorbed by an intrinsic column density, $`\mathrm{\Delta }N_\mathrm{H}`$, of cold gas which was a further free parameter in the fits. The abundances of metals in the absorbing material were fixed to their solar values (Anders & Grevesse 1989) although the effects of varying these values were examined. Finally, a fourth spectral model, model D, was examined which was similar to model C but with the constant-pressure cooling flow replaced with a second isothermal emission component. The temperature and normalization of this second emission component were included as free parameters in the fits. (As with model C, the normalizations were linked to the same value in all four detectors, scaled by the appropriate normalization factors.) The second emission component was also assumed to be intrinsically absorbed by a column density of cold gas, $`\mathrm{\Delta }N_\mathrm{H}`$, which was a free parameter in the fits. The metallicities of the two emission components were linked to take the same values.
Fig. 1 shows the ASCA data and best-fitting spectral models for four of the clusters in the sample; Abell 1704, 2029, 2204 and 2219. Table 4 summarizes the fit results for all of the clusters using the four spectral models. For each cluster we list the best-fitting parameter values and 90 per cent ($`\mathrm{\Delta }\chi ^2=2.71`$) confidence limits. The mass deposition rates ($`\dot{M}_\mathrm{S}`$), $`210`$ keV X-ray luminosities ($`L_\mathrm{X}`$) and bolometric luminosities ($`L_{\mathrm{Bol}}`$) are the values measured in the G3 detector. Note that for the clusters at relatively low redshifts ($`z<0.1`$), we have corrected the $`L_\mathrm{X}`$ and $`L_{\mathrm{Bol}}`$ values for the emission arising at radii $`>6`$ arcmin by scaling from the David et al. (1993) results, which were based on GINGA and Einstein MPC observations.
### 3.2 The goodness-of-fit
Table 5 summarizes the goodness-of-fit measurements obtained with the different spectral models. The tabulated results are the probabilities of exceeding the $`\chi ^2`$ values obtained, assuming in each case that the model correctly describes the spectral properties of the clusters. Goodness-of-fit values $`<10^2`$ may be regarded as indicating a formally unacceptable fit to the data.
The results from the spectral analysis show that model A (the isothermal model with the absorbing column density fixed at the nominal Galactic value) provides a reasonable description of most (21/30) of the clusters. However, model A fails significantly for the massive CF clusters Abell 1068, 1795, 2029 and 2204. The significance of the improvements to the fits obtained with spectral model B over model A have been evaluated using the F-test for the introduction of an additional free parameter (Bevington 1969). These significances are listed in Column 6 of Table 5. We see that the improvements obtained when including the column density as a free parameter in the fits with the isothermal models are significant at $`90`$ per cent confidence for 23 of the 30 clusters (and would also be highly significant for Abell 478 if the nominal Dickey & Lockman 1990 value for the Galactic column density, rather than the X-ray value determined by Allen & Fabian 1997 were used in model A). Model B provides a statistically acceptable fit for 27 of the 30 clusters.
The results on the goodness-of-fit determined with spectral models C and D are summarized in columns 4 and 5 of Table 5. The only clusters that these models do not adequately describe are the two nearest, brightest systems; Abell 1795 and Abell 2029. Spectral models C and D include an absorption component acting on the entire cluster emission which is normalized to the appropriate Galactic column densities determined from HI studies (Dickey & Lockman 1990). We have examined the significance of the improvements to the fits obtained when allowing the Galactic column density to also be a free parameter in the fits with these models. The fits were not significantly improved for any of the clusters except PKS0745-191, for which the Galactic column density is most uncertain. (The Galactic column density to PKS0745-191 varies between 3.2 and $`4.4\times 10^{21}`$ atom cm<sup>-2</sup> in the surrounding regions studied by Dickey & Lockman 1990). We find $`\mathrm{\Delta }\chi ^2=7.9`$ with model C and $`\mathrm{\Delta }\chi ^2=20.2`$ with model D for PKS0745-191 (indicating improvements significant at $`>99`$ per cent confidence) when allowing the Galactic column density to be free. The preferred value for the Galactic column density for this cluster from the ASCA data is $`3.5\times 10^{21}`$ atom cm<sup>-2</sup>. For Abell 1795, we also find a marginal improvement to the fits with models C and D when allowing the Galactic column density to be a free parameter ($`\mathrm{\Delta }\chi ^2=4.5`$ with Model C and $`\mathrm{\Delta }\chi ^2=5.4`$ for model D, which is significant at the 96.5 and 98 per cent level, respectively) with a preferred value of $`8\times 10^{19}`$ atom cm<sup>-2</sup>. The agreement between the Galactic column densities inferred from the 21 cm and X-ray data for most of the clusters with spectral models C and D provides further support for the validity of these models.
### 3.3 The requirement for multiphase models
The statistical improvements obtained by introducing an additional emission component into the fits with the single-temperature model (model B) are summarized in Columns 7 and 8 of Table 5. We list the results obtained for both cases, where the extra emission component is modeled as a cooling flow (for the CF systems only) or a second isothermal emission component. (Essentially we use spectral models C and D, but also include the Galactic column density as a free parameter it the fits so as to permit a direct comparison with the goodness of fit measured with spectral model B). The significance of the improvements to the fits obtained with the multiphase models with respect to model B have been quantified using an F-test for the introduction of 2 or 3 extra parameters (for models C and D respectively). We find that the improvements with model D are significant at $`90`$ per cent confidence for 21 of the 30 clusters. The only clusters for which we find no significant improvement with the introduction of a second emission component are the NCF clusters Abell 520, 773, 2218 and AC114, and the more distant and/or less-luminous CF clusters Abell 963, MS1358.4+6245, MS1455.0+2232, MS2137.3-2353 and Abell 2390, which have amongst the lowest numbers of counts in their ASCA spectra. For the CF clusters (the only clusters to which model C was applied) the improvements obtained with the cooling-flow model over model B are significant at $`90`$ per cent confidence in 13 of 21 systems, with again those CF clusters with the lowest numbers of total counts in their ASCA spectra exhibiting the least-significant improvements.
We thus find that most CF clusters, and a number of NCF systems (in general those systems with the best signal-to-noise ratios in their ASCA spectra) exhibit significant improvements to their fits with the use of multiphase, over single-phase, models. These improvements generally indicate the presence of emission from material cooler than the mean cluster temperatures (see Table 4). For the CF clusters, this can be naturally understood as being due to emission associated with the cooling flows. For the NCF clusters, the improvements to the fits obtained with model D are likely to reflect the presence of merging subclusters with lower virial temperatures. This is supported by optical, radio and X-ray imaging (e.g. Edge et al. 1992; Buote & Tsai 1996; Markevitch 1996; Feretti, Giovannini & Böhringer 1997; Rizza et al. 1998) and gravitational lensing studies (e.g. Kneib et al. 1995; Smail et al. 1995, 1997; Squires et al. 1997; Allen 1998) of the NCF clusters, which show that these systems typically exhibit complex morphologies and centroid shifts indicative of merger events.
### 3.4 Two-temperature versus cooling-flow models
The results in Table 4 show that spectral model D typically provides at least as good a fit to the ASCA spectra for the CF clusters as the cooling-flow model (model C). This has sometimes been taken to indicate that the X-ray gas in these systems is distinctly two-phase, with the cooler phase being due to the dominant cluster galaxy (e.g. Makishima 1997; Ikebe et al. 1999). However, when interpreting these results it is important to recall that the two-temperature model provides a more flexible fitting parameterization, with an extra degree of freedom, when applied to ASCA observations. Simulated cooling-flow spectra constructed with spectral model C, including plausible levels of intrinsic absorption (Section 6) and observed at the spectral resolution and signal-to-noise levels typical of ASCA observations, are invariably well-described by two-temperature models. This is illustrated in Fig. 2 where we show a simulated ASCA SIS spectrum for a 7 keV cluster with a metallicity of 0.4$`Z_{}`$, containing a cooling flow (intrinsically absorbed by a column density of $`4\times 10^{21}`$atom cm<sup>-2</sup>) accounting for 30 per cent of the total $`210`$ keV flux. A count rate of 1ct s<sup>-1</sup>, an exposure time of 40ks and a Galactic column density of $`10^{20}`$atom cm<sup>-2</sup> have been assumed. Overlaid, we show the best-fitting two-temperature model, which provides a good fit to the simulated data.
Thus, even in the case where the constant-pressure cooling flow model provides an exact description of the data (i.e. in the simulations), the two-temperature model provides a similarly good fit. In a real cluster, where spectral model C will undoubtedly over-simplify the true situation, the more flexible two-temperature model (model D) is likely to provide a better match to the observations. Thus, our finding that spectral model D often provides a better fit to the observed cluster spectra than model C, only shows that the constant pressure cooling flow model over-simplifies the true spectra of the cooling flows in the clusters (see also Section 6.4). We note that such tendencies would be enhanced by the presence of additional cool components in the spectra, for example due to temperature gradients at large radii (e.g. Markevitch et al. 1998). Finally, we note that the very high luminosities associated with the cool emission components in clusters like RXJ1347.5-1145 (with $`L_{\mathrm{cool}}`$ a few $`10^{45}`$$`\mathrm{erg}\mathrm{s}^1`$) are difficult to explain with models in which the cool emission is due to the interstellar medium of the central cluster galaxies.
## 4 Deprojection analysis of the X-ray images
### 4.1 Method and primary results
The analysis of the HRI imaging data has been carried out using an extensively updated version of the deprojection code of Fabian et al. (1981). Azimuthally-averaged X-ray surface brightness profiles were determined for each cluster from the HRI images. These profiles were background-subtracted, corrected for telescope vignetting and re-binned to provide sufficient counts in each annulus to allow the analysis to be extended to radii of at least 500 kpc. (Bin sizes of 8-24 arcsec were used.)
With the X-ray surface brightness profiles as the primary input, and under assumptions of spherical symmetry and hydrostatic equilibrium, the deprojection technique can be used to study the basic properties of the intracluster gas (temperature, density, pressure, cooling rate) as a function of radius. The deprojection code uses a monte-carlo method to determine the statistical uncertainties on the results and incorporates the appropriate HRI spectral response matrix issued by GSFC. The cluster metallicities were fixed at the values determined from the spectral analysis (Table 4). The absorbing column densities were fixed at the appropriate Galactic values (Table 1).
The deprojection code requires the total mass profiles for the clusters (which define the pressure profiles) to be specified. We have iteratively determined the mass profiles that result in deprojected temperature profiles (which approximate the mass-weighted temperature profiles in the clusters) that are isothermal within the regions probed by the HRI data (the central $`0.51`$ Mpc) and which are consistent with the spectrally-determined temperatures (Section 3). The assumption of approximately isothermal mass-weighted temperature profiles in the central regions of the clusters is supported by the following evidence: firstly, ASCA observations of nearby cooling flows show that in the central regions of these systems the gas is multiphase, but that the bulk of the X-ray gas there has a temperature close to the cluster mean (e.g. Fabian et al. 1994b; Fukazawa et al. 1994; Matsumoto et al. 1996; Ikebe et al. 1999). Secondly, combined X-ray and gravitational lensing studies of CF clusters (e.g. Allen 1998) show that approximately isothermal mass-weighted temperature profiles in the cores of such clusters lead to excellent agreement between their X-ray and gravitational-lensing masses. Thirdly, the use of approximately constant mass-weighted temperature profiles implies a more plausible range of initial density inhomogeneities in the clusters than would be the case if the temperature profiles decreased within their cores (Thomas, Fabian & Nulsen 1987). Finally, the use of approximately isothermal mass-weighted temperature profiles in the deprojection analyses leads to independent determinations of the mass deposition profiles in the cooling flows, from the X-ray spectra and imaging data, in excellent agreement with each other (e.g. Allen & Fabian 1997). We note that the assumption of a constant mass-weighted deprojected temperature profile is consistent with measurements of a decreasing emission-weighted temperatures in the cores of many CF clusters (e.g. Waxman & Miralda-Escudé 1995).
The mass profiles for the clusters were parameterized as isothermal spheres (Equation 4-125 of Binney & Tremaine 1987) with adjustable core radii ($`r_\mathrm{c}`$) and velocity dispersions ($`\sigma `$). The core radii were adjusted until the temperature profiles determined from the deprojection code became isothermal. The velocity dispersions were then adjusted until the temperatures determined from the deprojection code came into agreement with the spectrally-determined values. Errors on the velocity dispersions are the range of values that result in isothermal deprojected temperature profiles that are consistent, at the 90 per cent confidence limit, with the spectrally-determined temperature results. Estimates of the thermal gas pressure in the outermost radial bins used in the analysis are also required by the deprojection code and were determined iteratively. (The uncertainties on the outer pressure estimates do not significantly affect the results presented here.) Although the deprojection method of Fabian et al. (1981) is essentially a single-phase technique, we note that it produces results in good agreement with more detailed multiphase treatments (Thomas, Fabian & Nulsen 1987) and, due to its simple applicability at large radii in clusters, is better-suited to the present project.
The mass distributions determined from the deprojection analysis are summarized in columns 3 and 4 of Table 6. These simple parameterizations permit direct comparisons with independent mass constraints from gravitational lensing and dynamical studies (e.g. Allen 1998). We note that for a few of the clusters (primarily the brightest CFs with the best data) the single-component mass models cannot adequately satisfy the requirement for approximately isothermal deprojected temperatures profiles and for these systems a significantly better match was obtained by introducing a second ‘linear’ mass component, truncated at a specified outer radius. The clusters requiring the additional mass components were Abell 586 ($`2\times 10^{11}`$ $`\mathrm{M}_{}`$kpc<sup>-1</sup> within the central 20 kpc), Abell 2029 ($`3\times 10^{11}`$ $`\mathrm{M}_{}`$kpc<sup>-1</sup> within the central 20 kpc), Abell 2219 ($`10^{11}`$ $`\mathrm{M}_{}`$kpc<sup>-1</sup> within the central 30 kpc), Abell 478 ($`4\times 10^{11}`$ $`\mathrm{M}_{}`$kpc<sup>-1</sup> within the central 20 kpc and an adjustment of the cluster parameters to $`\sigma =840_{60}^{+80}`$$`\mathrm{km}\mathrm{s}^1`$ and $`r_\mathrm{c}=100`$ kpc) and Abell 1795 ($`10^{11}`$ $`\mathrm{M}_{}`$kpc<sup>-1</sup> within the central 20 kpc, and an adjustment of the cluster parameters to $`\sigma =720_{20}^{+20}`$$`\mathrm{km}\mathrm{s}^1`$ and $`r_\mathrm{c}=60`$ kpc).
The core radii of the mass distributions determined from the deprojection analysis are similar to the values measured from simple ‘$`\beta `$model’ fits to the X-ray surface brightness profiles (e.g. Cavaliere & Fusco-Femiano 1976, 1978). It is well known that the $`\beta `$-model does not provide a good match to the X-ray surface brightness profiles of the brightest CF clusters when their central regions are included in the analysis (e.g. Jones & Forman 1984). However, where $`\beta `$-models are used to estimate the mass core radii in CF clusters, the central regions of these clusters should not simply be excised from the fits. The central surface brightness peaks in CF clusters trace the presence of sharp central density rises in these systems and are not only the result of excess emission due to cooling gas. (It is the high central densities that lead to the high cooling rates.) Although fully accounting for the effects of cooling on the central density measurements is difficult with current data, simply excising the inner regions of the clusters can lead to overestimates of the mass core radii and corresponding underestimates of the central gas and total mass densities.
The basic results on the cooling flows from the deprojection analysis are summarized in columns $`58`$ of Table 6. For those clusters in common, the results obtained are generally in good agreement with values reported from previous works (e.g. Edge et al. 1992, White et al. 1997; Peres et al. 1998).
## 5 The mass deposition rates from the cooling flows
### 5.1 Comparison of the spectral and imaging results
The spectral and image deprojection analyses discussed in Sections 3 and 4 provide essentially independent estimates of the mass deposition rates in the clusters. A comparison of the results obtained from these analyses therefore provides a test of the validity of the cooling flow model.
The deprojection method describes the X-ray emission from a cluster as arising from a series of concentric spherical shells. The luminosity in a particular shell, $`j`$, may be written as the sum of four components (Arnaud 1988).
$$L_j=\mathrm{\Delta }\dot{M}_jH_j+\mathrm{\Delta }\dot{M}_j\mathrm{\Delta }\mathrm{\Phi }_j+\left[\underset{i=1}{\overset{j1}{}}\mathrm{\Delta }\dot{M}_i(\mathrm{\Delta }\mathrm{\Phi }_j+\mathrm{\Delta }H_j)\right],$$
(1)
where $`\mathrm{\Delta }\dot{M}_j`$ is the mass deposited in shell $`j`$, $`H_j`$ is the enthalpy of the gas in shell $`j`$, and $`\mathrm{\Delta }\mathrm{\Phi }_j`$ is the gravitational energy released in crossing that shell. $`_{i=1}^{j1}\mathrm{\Delta }\dot{M}_i`$ is the mass flow rate through shell $`j`$, and $`\mathrm{\Delta }H_j`$ the change of enthalpy of the gas as it moves through that shell. The first term in equation 1 thus accounts for the enthalpy of the gas deposited in shell $`j`$. The second term is the gravitational work done on the gas deposited in shell $`j`$. The third and fourth terms respectively account for the gravitational work done on material flowing through shell $`j`$ to interior radii, and the enthalpy released by that material as it passes through the shell.
In any particular shell, the densest material in the cooling flow is assumed to cool out and be deposited. Since the cooling time of this material will be short compared to the flow time, the cooling can be assumed to take place at a fixed radius. Thus, the luminosity contributed by the first term in equation 1 should have a spectrum that can be approximated by gas cooling at constant pressure from the ambient cluster temperature i.e. the same spectrum as the cooling component incorporated into the spectral analysis with model C (Section 3).
For the bulk of the material continuing to flow inwards towards the cluster centre, the cooling via X-ray emission is assumed to be offset by the gravitational work done on the gas as it moves inwards. The emission accounted for in the second and third terms of equation 1 should therefore have a spectrum that can be approximated by an isothermal plasma at the appropriate ambient temperature for the cluster i.e. the spectrum of the isothermal emission component in model C. Since the mass-weighted temperature profiles in the clusters are assumed to remain approximately isothermal with radius, the luminosity contributed by the fourth term of equation 1 should be negligible.
The mass deposition rates, $`\dot{M}_\mathrm{I}`$, listed in Table 6 are the mass flow rates ($`_{i=1}^{j1}\mathrm{\Delta }\dot{M}_i`$) determined from the deprojection analysis, at the point where the mean cooling time of the cluster gas first exceeds the Hubble time ($`1.3\times 10^{10}`$ yr). Thus, if the cooling flow model is correct, and the cooling flows have been undisturbed for a significant fraction of a Hubble time, the mass deposition rates determined from the deprojection analysis should be similar to those measured independently from the spectral data.
Fig. 3 shows the mass deposition rates determined from the spectral analysis ($`\dot{M}_\mathrm{S}`$) versus the image deprojection results ($`\dot{M}_\mathrm{I}`$), for those clusters for which both quantities have been measured. We see that the results exhibit an approximately linear correlation; a fit to the data with a power-law model of the form $`\dot{M}_\mathrm{S}=P\dot{M}_\mathrm{I}^Q`$, using the Akritas & Bershady (1996) bisector modification of the ordinary least-squares statistic, gives a best-fitting power-law slope of $`1.0\pm 0.2`$ (where the error is the standard deviation determined by bootstrap re-sampling), although the spectrally-determined mass deposition rates typically exceed the deprojection results by a factor of $`23`$. However, the spectral analysis presented in Section 3 also requires that the cooling gas is intrinsically absorbed by equivalent hydrogen column densities of, typically, a few $`10^{21}`$ atom cm<sup>-2</sup> (model C). The mass deposition rates determined from the deprojection analysis must therefore also be corrected for the effects of this absorbing material.
The corrections for the effects of intrinsic absorption on the deprojection results have been carried out by re-running the deprojection analysis with the absorbing column densities set to the total values determined with spectral model C. We assume that the absorption is due to cold gas with solar metallicity. If, however, the intrinsic absorption were instead due to dust, a possibility examined in more detail in Section 7, then the required correction factors would likely be reduced by a few tens of per cent. (The introduction of a simple OIK absorption edge at $`E0.54`$keV, such as might be associated with oxygen-rich, silicate dust grains, generally provides as good a description of the intrinsic absorption detected in the $`0.610.0`$keV ASCA spectra as a cold, gaseous absorber.) Since the intrinsic column densities measured with model C are redshifted quantities, we set the total column densities used in the revised deprojection analysis, which assumes zero redshift for the absorber, to be $`N_\mathrm{H}+\mathrm{\Delta }N_\mathrm{H}/(1+z)^3`$, where $`N_\mathrm{H}`$ is the Galactic column density). The absorption-corrected mass deposition rates ($`\dot{M}_\mathrm{C}`$) so determined are summarized in Table 7. Fig. 4 compares the spectrally-determined mass deposition rates with the absorption-corrected deprojected values. We see that once the presence of intrinsic absorption has been accounted for in a consistent manner in the imaging and spectral analyses, the results on the mass deposition rates in the clusters show good agreement.
### 5.2 The fraction of flux from the cooling flows
Table 8 lists the fractions of the total $`210`$ keV X-ray fluxes from the CF clusters associated with cooling gas, as determined with spectral model C. The results range from $`10`$ per cent for systems like Abell 1689 and 2142 to $`>30`$ per cent for Zwicky 3146, Abell 1068, RXJ1347-1145 and Abell 2204. Such a range of measurements is consistent with the results of Peres et al. (1998), from a deprojection analysis of a flux-limited sample of the brightest clusters known.
The fraction of the X-ray flux from a cluster due to it’s cooling flow can be expected to increase with the age of the flow. It is therefore likely that the cooling flows in Abell 1689 and 2142 are younger than those in Abell 2204 and Zwicky 3146. We note that the X-ray data are not, in general, of sufficient quality to place firm constraints on the ages of the cooling flows in these clusters (for a discussion of cooling flow age measurements see Allen et al. 1999).
It is interesting that the clusters with the largest cooling-flow flux fractions (with the exception of Abell 1413 and 2261) are also amongst the most optically line-luminous cooling flows known (e.g. Crawford et al. 1999). The production of powerful optical emission lines and associated UV/blue continuum emission is known to be due, at least in part, to the formation of massive stars at the centres of cooling flows (e.g. Johnstone, Fabian & Nulsen 1987; Allen 1995; Cardiel et al. 1995, 1998; McNamara et al. 1996; Voit & Donahue 1997). The results presented here are consistent with the idea that such star formation is typical of undisturbed cooling flows. It is also interesting to note that the clusters with the largest cooling flow flux fractions and most optically line-luminous emission line nebulae also tend to have the shortest cooling times averaged over their central 100 kpc regions, with $`\overline{t_{100}}<3`$Gyr (see also Peres et al. 1998).
## 6 Intrinsic X-ray absorption in clusters
### 6.1 The measured column densities and model dependency of the results
The equivalent column densities of intrinsic X-ray absorbing material inferred from the spectral analysis are summarized in Table 9. All three of the spectral models incorporating a free-fitting absorption component (models B,C,D) reveal the presence of excess X-ray absorption over and above the nominal Galactic values along the lines of sight to the clusters (Dickey & Lockman 1990), although the measured column densities are sensitive to the spectral model used.
Using the simple isothermal model with free-fitting absorption (model B), we determine a mean excess column density for the whole sample of 30 clusters of $`3.8\pm 3.2\times 10^{20}`$ atom cm<sup>-2</sup>. When we examine the CF and NCF clusters separately, we find that both subsamples exhibit intrinsic absorption, at a similar level on average. For the CF clusters, the mean excess column density measured with spectral model B is $`3.5\pm 3.0\times 10^{20}`$ atom cm<sup>-2</sup>. For the NCF systems, the value is $`4.4\pm 3.9\times 10^{20}`$ atom cm<sup>-2</sup>. The application of a Students t-test (accounting for the possibility of unequal variances in the two distributions; Press et al. 1992) shows the mean excess column densities for the CF and NCF clusters, determined with spectral model B, to differ at only the 45 per cent confidence level. The application of a Kolmogorov-Smirnov test shows the two subsamples to be drawn from different parent populations at only the $`10`$ per cent confidence level. Thus, our analysis with spectral model B suggests that intrinsic absorption is not confined to CF clusters.
The results on the column densities are quite different, however, when the more-sophisticated multiphase spectral models are used. As discussed in Section 3, the multiphase models generally provide a better description of the X-ray properties of the CF clusters (and some of the brighter NCF systems). Using spectral model C, our preferred model for the CF clusters in that it provides a consistent description of the spectral and imaging X-ray data for these systems (Section 5.1), we determine a mean intrinsic column density acting on the cooling-flow components of $`3.4\pm 1.3\times 10^{21}`$ atom cm<sup>-2</sup>. Using the two-temperature spectral model (model D), we determine a mean intrinsic column density acting on the cooler emission components in the CF systems of $`6.3\pm 4.7\times 10^{21}`$ atom cm<sup>-2</sup>. Thus, the intrinsic column densities for the CF clusters measured with the multiphase models (models C and D) are similar and approximately an order of magnitude larger than the values inferred using the single-phase model B. These results demonstrate the need for adopting an appropriate spectral when attempting to measure the column densities of absorbing material in CF clusters. We also note that where spectral model D provided a significant improvement to the fits to the NCF clusters, with respect to model B, the measured column densities acting on the cooler emission components were also significantly larger than the values determined using the single-phase models.
We have also measured the intrinsic column densities in the CF clusters using one further spectral model, referred to in Table 9 as model C’. Model C’ is identical to model C except that it assumes that the excess absorption acts on the entire cluster spectrum, rather than just the cooling gas, and that the absorbing material lies at zero redshift. Model C’ has been used in a number of previous studies and is included here for comparison purposes. The mean excess column density for the CF clusters determined with spectral model C’ is $`7.2\pm 3.9\times 10^{20}`$ atom cm<sup>-2</sup>.
Allen & Fabian (1997) present results from an X-ray colour deprojection study of 18 clusters observed with the ROSAT PSPC. These authors determine intrinsic column densities across the central 30 arcsec (radius) regions of Abell 478, 1795 and 2029, using spectral model C’, of $`7.40\pm 0.63`$, $`2.20\pm 0.20`$ and $`0.83\pm 0.42\times 10^{20}`$ atom cm<sup>-2</sup>, respectively. The results for Abell 478 and Abell 1795 are in excellent agreement with those presented here, although the Allen & Fabian (1997) value for Abell 2029 is $`4`$ times smaller than our result. Abell 2029 is unusual in that it hosts a strong cooling flow without associated optical line emission. This may indicate that the central regions of the cluster have been disrupted (perhaps by a minor merger event or the strong central radio source in this system), complicating the distributions of cooling and absorbing material (cf. Section 5.2).
Four of the CF clusters studied here were also examined by White et al. (1991), using Einstein Observatory SSS data. These authors measured intrinsic column densities for Abell 478, 1795, 2029 and 2142 (in a 3 arcmin radius circular aperture, using spectral model C’) of $`1.7_{0.7}^{+0.8}\times 10^{21}`$ atom cm<sup>-2</sup>, $`0.8_{0.3}^{+0.3}\times 10^{21}`$ atom cm<sup>-2</sup>, $`1.8_{0.5}^{+0.5}\times 10^{21}`$ atom cm<sup>-2</sup>and $`1.3_{0.4}^{+0.3}\times 10^{21}`$ atom cm<sup>-2</sup>, respectively. (The SSS result for Abell 478 has been corrected to account for the different value of Galactic absorption assumed in that study). The intrinsic column densities determined by White et al. (1991) are $`25`$ times larger than the values measured from the ASCA data.
Allen & Fabian (1997) compared the results from their X-ray colour deprojection study of PSPC data with the White et al. (1991) SSS analysis and concluded that if the intrinsic absorption were due to cold gas, then agreement between the measured column densities could only be obtained if the covering fraction of the X-ray absorbing material were, in general, $`<0.5`$. We have therefore examined the constraints on the covering fraction, $`f`$, that can be obtained from the ASCA data using spectral model C. The results are also listed in Table 9. In all cases we determine a best-fit covering fraction for the intrinsic X-ray absorbing material of unity. For many clusters, covering fractions of $`<70`$ per cent can be firmly ruled out. (The physical significance of these results are further explored in Section 6.4). The results on the covering fractions, and the reasonable agreement of the ASCA and PSPC results for Abell 478 and 1795, suggest that the White et al. (1991) results may have systematically over-estimated the intrinsic column densities in clusters by a factor of a few.
Before considering in more detail the possible origin and nature of the absorbing material, it is pertinent to consider whether the intrinsic absorption, inferred to be present using a variety of spectral models, could in fact be an artifact due to an unconsidered emission process in the clusters. Fig. 5 shows the SIS0 spectra for Abell 1068 and 1795, which exhibit two of the clearest absorption signatures. The data between $`3.0`$ and $`10.0`$ keV have been fitted with a single-temperature emission model with Galactic absorption (i.e. spectral model A). The residuals to these fits, over the full $`0.610.0`$ keV band, are shown in the lower panels. We see that the residuals, which can be naturally explained by the introduction of cool emission components with associated intrinsic absorption (using spectral models C and D), take the form of a large excess of counts between $`0.82.5`$ keV. We therefore include the caveat that if some extra emission process which can account for such residuals in the X-ray spectra were identified, then the requirement for intrinsic absorption from the X-ray data would be greatly diminished. Finally, we note that the fits with spectral Model A over the restricted $`3.010.0`$ energy range lead to determinations of the cluster temperatures for Abell 1068 and 1795 of $`6.1_{1.6}^{+2.6}`$ and $`6.1_{0.4}^{+0.5}`$ keV, respectively, in excellent agreement with the results obtained for the full data sets with spectral model C.
### 6.2 The mass of absorbing material
The results presented in Section 6.1 show that the cooling flows in our sample are typically intrinsically absorbed by equivalent hydrogen column densities of a few $`\times 10^{21}`$ atom cm<sup>-2</sup> (taking spectral model C as our preferred model). Assuming, in the first case, that this absorption is due to cold gas with solar metallicity, the mass of absorbing gas, $`M_{\mathrm{abs}}`$, within a radius, $`r_{\mathrm{abs}}`$, may be approximated as
$$M_{\mathrm{abs}}3.2\times 10^7r_{\mathrm{abs}}^2\mathrm{\Delta }N_\mathrm{H}\mathrm{M}_{},$$
(2)
where $`r_{\mathrm{abs}}`$ is in units of kpc and $`\mathrm{\Delta }N_\mathrm{H}`$ in $`10^{21}`$ atom cm<sup>-2</sup>. For metallicities of $`0.4Z_{}`$ in the absorbing gas (the mean emission-weighted value for the X-ray emitting gas in the CF clusters; Allen & Fabian 1998b) the implied masses are a factor $`2`$ larger. For zero metallicity these masses are a factor $`3`$ times larger than the values given by equation 2. (The scaling factors for different metallicities in the absorbing material have determined from spectral simulations using the XSPEC code.)
We have adopted $`r_{\mathrm{abs}}`$ as the radius where the cooling time of the cluster gas first exceeds $`5\times 10^9`$ yr in the clusters (a plausible age for the cooling flows in most of the CF systems). These radii, and the masses of absorbing gas implied by equation 2 (using the column densities determined with spectral model C) for both solar metallicity and $`Z=0.4Z_{}`$, are listed in columns $`24`$ of Table 10.
The observed masses of absorbing material may be compared to the masses expected to have been accumulated by the cooling flows $`(M_{\mathrm{acc}})`$ within the same radii over the past $`5\times 10^9`$ yr. Assuming that from time $`t=0`$ to $`t=5\times 10^9`$ yr, the integrated mass deposition rate within radius $`r_{\mathrm{abs}}`$ increases approximately linearly with time, the total mass accumulated within radius $`r_{\mathrm{abs}}`$ after $`5\times 10^9`$ yr is approximately
$$M_{\mathrm{acc}}2.5\dot{M}(r<r_{\mathrm{abs}})\times 10^9\mathrm{M}_{}.$$
(3)
At times $`t>5\times 10^9`$ yr, the mass deposition rate within $`r_{\mathrm{abs}}`$ should remain approximately constant. If instead we assume $`\dot{M}(r<r_{\mathrm{abs}})`$ to be constant during the first 5 Gyr, the accumulated mass will be twice the value indicated by equation 3. (Exactly how the mass deposition rate within $`r_{\mathrm{abs}}`$ grows with time during the first 5 Gyr is unclear and will depend upon the evolution of the cluster and the detailed properties of the cluster gas.)
The mass deposition rates within radii $`r_{\mathrm{abs}}`$ (corrected for the effects of absorption due to cold gas, as discussed in Section 5.1) and the estimated masses of cooled material accumulated by the cooling flows within these radii over a 5 Gyr period, are summarized in Table 10. Fig. 6 compares the accumulated masses with the masses of absorbing material determined from the spectral data. The agreement between the results is reasonable and is maximized for a metallicity in the absorbing gas of $`0.40.6Z_{}`$ (depending on the growth of the cooling flow over the first 5Gyr). This supports the idea that the observed X-ray absorption may be due to material accumulated by the cooling flows. We note that if the X-ray absorption in the CF clusters were due to dust rather than cold gas, a possibility examined in more detail in Section 7, then the agreement between the observed and predicted masses of absorbing material shown in Fig. 6 might still be expected to hold if the dust were contained in the material deposited by the cooling flows. In this case, the mass of the absorber calculated with equation 2 would be the mass of cooled gas associated with the dust when it was deposited from the cooling flow. For a Galactic dust to gas ratio, the mass in dust would be $`100`$ times smaller than the associated gas mass.
### 6.3 The luminosity reprocessed in other wavebands
The luminosities absorbed at X-ray wavelengths must eventually be reprocessed in other wavebands. If the absorbing material is dusty (as is likely to be the case in the central regions of the clusters e.g. Voit & Donahue 1995; Fabian, Johnstone & Daines 1994; Allen et al. 1995) then the bulk of this reprocessed emission is likely to emerge in the infrared band (e.g. Dwek, Rephaeli & Mather 1990). Table 11 summarizes the reprocessed luminosities (i.e. the bolometric luminosities absorbed within the clusters) measured with spectral model C. The reprocessed luminosities range from $`2\times 10^{44}`$ $`\mathrm{erg}\mathrm{s}^1`$ for Abell 1795 to $`5\times 10^{45}`$ $`\mathrm{erg}\mathrm{s}^1`$for RXJ1347.5-1145, with a mean value of $`10^{45}`$ $`\mathrm{erg}\mathrm{s}^1`$.
Allen et al. (1999) report detections of spatially extended $`100\mu `$m emission coincident with the X-ray centroids in the nearby Centaurus cluster and Abell 2199. The infrared fluxes from those clusters were also shown to be in good agreement with the predicted values due to reprocessed X-ray emission from their cooling flows. We have also used the IPAC SCANPI software and archival IRAS data to measure the 60 and $`100\mu `$m fluxes within a four arcmin (radius) aperture centred on the X-ray centroids for the clusters in the present sample. (The medians of the co-added SCANPI results were used.) The results are listed in Table 11. The error bars associated with the measured fluxes are the root-mean-square deviations in the residuals, external to the source extraction regions, after baseline subtraction. Where no detection was made, an upper limit equal to three times the r.m.s. deviation in the residuals is listed. Where a detection was made, we also list the in-scan separations (in arcmin) between the peak of the $`100\mu `$m emission (or the $`60\mu `$m emission in the case of Abell 1835) and the X-ray centre.
Two of the clusters in our sample, IRAS 09104+4109 and Abell 1068 have known infrared point sources coincident with their X-ray centroids (Kleinmann et al. 1988; Moshir et al. 1989). The presence of these sources is confirmed by our analysis. The 60 and $`100\mu `$m data for Abell 1704 and $`60\mu `$m data for Abell 1835 also suggest the presence of unresolved sources, coincident with the X-ray centroids for the clusters. Abell 478 and Zwicky 3146 have $`100\mu `$m emission originating from close to their X-ray centroids, which appears spatially extended, although the infrared flux from Abell 478 is probably contaminated by Galactic cirrus. The data for Abell 2142, 2204, 2261 and 2390 also provide significant detections within the four arcmin (radius) source apertures, although the peaks of the infrared emission from these systems are spatially offset from their X-ray centroids, suggesting that the detected flux is likely to originate, at least in part, from some other source.
Several of the clusters included in our study (Abell 586, 963, 1413, 1795, 2029, 2142 and PKS0745-191) have previously been studied at infrared wavelengths, using IRAS data, by Wise et al. (1993) and/or Cox, Bregman & Schombert (1995). These authors also report no clear detections of infrared emission from these sources. Edge et al. (1999) report measurements of $`60\mu `$m emission from Abell 1835 and 60 and $`100\mu `$m emission from Abell 2390, using a similar analysis to that presented here. These authors also present detections of $`850\mu `$m emission from these clusters which, for Abell 1835, they suggest is likely to be due to dust heated by the vigorous star formation observed in the cluster core or an obscured active galactic nucleus.
Following Helou et al. (1988) and Wise et al. (1993), we can estimate the total infrared luminosities implied by the observed IRAS fluxes and flux limits using the relation
$$L_{11000\mu \mathrm{m}}2.8\times 10^{44}(\frac{z}{0.05})^2(2.58S_{60}+S_{100})\mathrm{erg}\mathrm{s}^1,$$
(4)
where $`S_{60}`$ and $`S_{100}`$ are the 60 and $`100\mu `$m IRAS fluxes in units of Jy. This relation assumes a dust temperature of $`30`$K (which is consistent with the observations; e.g. Allen et al. 1999; Edge et al. 1999) and an emissivity index, $`n`$, in the range $`02`$, where the emissivity is proportional to the frequency, $`\nu ^n`$. Where measurements, rather than upper limits, were obtained, we associate a systematic uncertainty of $`30`$ per cent with the estimated $`11000\mu `$m luminosities, which is combined in quadrature with the random errors. The $`11000\mu `$m luminosities calculated from this relation are summarized in Table 11. In general, the measurements and upper limits to the observed $`11000\mu `$m luminosities exceed the predicted luminosities due to reprocessed X-ray radiation from the cooling flows. Where measurements (rather than upper limits) are made, this may suggest the presence of an additional source of infrared flux, such as star formation or AGN associated with the central cluster galaxies, or some other contaminating (possibly Galactic) source. At some level, the massive starbursts at the centres of systems like Abell 1068, 1835, 2204 2390 and Zwicky 3146 (Allen 1995; Crawford et al. 1999; Edge et al. 1999) must contribute to the detected infrared flux.
### 6.4 Systematic uncertainties in the absorption measurements
As discussed in Section 6.1, the measured intrinsic equivalent hydrogen column densities of absorbing material inferred from the ASCA spectra exhibit significant systematic variation, depending upon whether the single-phase or multiphase emission models are used. We have adopted the results obtained with the cooling-flow model (model C) as our preferred values, since this model provides a consistent description of the spectral and imaging X-ray data for the clusters. However, a number of systematic uncertainties affecting the absorption results remain with the cooling flow model.
Spectral model C undoubtedly over-simplifies the situation in a real cooling flow. In particular, the model assumes that the absorbing material lies in a uniform screen in front of the cooling flows. If the absorber is intrinsically associated with the clusters, then the absorbing material must, to some extent, be distributed throughout the cooling flows. Allen & Fabian (1997; see also Wise & Sarazin 1999) discuss two different, limiting geometries for the absorbing material; partial covering and multilayer absorption. The partial covering model, with a covering factor $`f<1`$, can be expected to apply or where there are $`<1`$ individual absorbing clouds along each line of sight and/or where the absorbing material is strongly clumped on large scales. However, as discussed in Section 6.1, covering fractions significantly less than unity are found to provide a poor description of the ASCA spectra, suggesting that the absorbing material is likely to be more smoothly distributed throughout the cooling flows.
The multilayer absorption model discussed by Allen & Fabian (1997) is applicable where the absorbing material is made up of a large number of similarly sized absorbing clouds, each with a column density much less than the total intrinsic column density ($`\mathrm{\Delta }N_\mathrm{H}`$), homogeneously distributed throughout the X-ray emitting medium. For an intrinsic emission spectrum $`A(E)`$, the observed spectrum, $`A^{}(E)`$, emerging from the emitting/absorbing region may be written (where $`\sigma (E)`$ is the absorption cross-section) as
$$A^{}(E)=A(E)\left(\frac{1e^{\sigma (E)\mathrm{\Delta }N_\mathrm{H}}}{\sigma (E)\mathrm{\Delta }N_\mathrm{H}}\right).$$
(5)
We have simulated cooling flow spectra with multilayer absorption, with total column densities of between $`10^{21}`$ and $`5\times 10^{22}`$ atom cm<sup>-2</sup>, and fitted these with simple models which assume that the absorbing material acts as a uniform screen in front of the cooling flow (as with model C) . For the simulations we adopt an emission spectrum appropriate for a constant pressure cooling flow, with an upper temperature of 7 keV and a metallicity of $`0.4Z_{}`$. We use a response matrix appropriate for the ASCA SIS detectors. The comparison between the fitted (uniform screen) column densities and the true total (multilayer) values are shown in Fig. 7. We see that for multilayer column densities of $`10^{21}5\times 10^{22}`$atom cm<sup>-2</sup>, the fitted column densities underestimate the true values by factors of $`510`$. (These results are not sensitive to any reasonable choice of Galactic column density.)
If the multilayer model approximates the true distribution of X-ray emitting and absorbing material in cooling flows, then the results obtained with spectral model C would imply true equivalent column-densities of intrinsic absorbing gas in the cluster cores of, typically, a few $`10^{22}`$atom cm<sup>-2</sup>. Importantly, we also note that for signal-to-noise ratios appropriate for ASCA observations of bright clusters (with count rates of $`1`$ct s<sup>-1</sup> in the SIS detectors and exposure times of $`40`$ ks), the simple uniform-screen absorption model provides a good fit to the simulated spectra with multilayer absorption, for total column densities $`<3\times 10^{22}`$ atom cm<sup>-2</sup>. (For larger total column densities, complex residuals are detected in the spectra at energies $`<2.0`$keV.) The large intrinsic column densities implied by the multilayer models would have strong implications for the nature of the absorbing matter. In particular, column densities of $`>10^{22}`$atom cm<sup>-2</sup> of gaseous material with solar metallicity would be very difficult to reconcile with current HI and CO limits, even if this material were very cold and highly molecular. (We also note that in the case of a multilayer distribution, the implied masses of absorbing gas would often exceed the predicted accumulated masses due to the cooling flows and imply, at least in part, some other origin for the absorbing medium.) In the following section, we explore these issues in more detail and summarize the current observational constraints on the physical nature of the absorbing matter.
## 7 Searches for the X-ray absorbing material in other wavebands
The results presented in this paper provide a consistent picture for the X-ray properties of cooling flows in X-ray luminous clusters of galaxies. The ASCA spectra and ROSAT images provide independent determinations of the mass cooling rates in the clusters in good agreement with each other, when a consistent method of analysis is employed (Section 5.1). The spectral data require the presence of large column densities of intrinsic X-ray absorbing material associated with the cooling gas (Section 6). The implied masses of absorbing material are in reasonable agreement with the masses expected to have been accumulated by the cooling flows over their lifetimes (Section 6.2), identifying this material as a plausible sink for the bulk of the cooled gas from the flows. Despite the evidence for intrinsic X-ray absorption associated with cooling flows, however, searches for this material in other wavebands have, to date, proved largely unsuccessful (at least outside of the central $`520`$ kpc). We here briefly review these observations and comment on their implications for the physical state of the X-ray absorbing matter. More detailed discussions of a number of these issues are given by Ferland, Fabian & Johnstone (1994), Daines, Fabian & Thomas (1994), Fabian et al. (1994a), O’Dea et al. (1994a), Voit & Donahue (1995), O’Dea & Baum (1996), Arnaud & Mushotzky (1988) and Henkel & Wiklind (1998).
As mentioned in Section 5.2, the central ($`520`$ kpc) regions of CF clusters are commonly observed to be sites of ongoing star formation with powerful, associated optical line emission (e.g. Johnstone et al. 1987; Heckman et al. 1989; McNamara & O’Connell 1989; Crawford & Fabian 1992; Allen 1995; Cardiel et al. 1995, 1998; McNamara et al. 1996; Voit & Donahue 1997; Crawford et al. 1999). Nearly all CF clusters with short central cooling times ($`t_{\mathrm{cool}}<2`$Gyr) exhibit such phenomena (Peres et al. 1998), whereas systems with longer central cooling times (including all NCF clusters) generally do not. The observed levels of star formation imply the presence of at least moderate masses of cooled gas in the cores of CF clusters, sufficient in a few cases to account for the bulk of the mass deposited by the cooling flows in those regions (Allen 1995) and at least some causal connection with the cooling flows.
At radio wavelengths, extensive searches have been carried out for 21cm line-emission associated with warm atomic hydrogen in the cores of CF clusters (e.g. Valentijn & Giovanelli 1982; O’Dea & Baum 1996). The negative results obtained (with the exception of M87; Jaffe 1992) typically constrain the column densities of optically thin material to be less than a few $`10^{19}`$ atom cm<sup>-2</sup> and, for an assumed spin temperature, $`T_\mathrm{S}20`$K, the number of clouds along the line of sight to be $`<1`$ (O’Dea & Baum 1996).
Searches have also been made for the 21 cm absorption signature of cold atomic hydrogen (e.g. McNamara, Bregman & O’Connell 1990; Jaffe 1992; Dwarakanath, van Gorkom & Owen 1994; O’Dea, Gallimore & Baum 1995; O’Dea & Baum 1996; Johnstone et al. 1998). Such material has been detected across the central few kpc of a few CF clusters, including Abell 426 (Crane, Van der Hulst & Haschick 1982; Sijbring et al. 1989; Jaffe 1990) 2A0335+096 (McNamara et al. 1990), Hydra A (Dwarakanath, van Gorkom & Owen 1995), MKW3s (McNamara et al. 1990) and Abell 2597 (O’Dea, Baum & Gallimore 1994) with typical column densities of a few $`10^{20}(T_\mathrm{S}/100)`$ atom cm<sup>-2</sup>. However, for most clusters only negative results have been obtained, with upper limits to the column densities of 10K gas (assuming a covering fraction of unity and a velocity width of $`1.5`$ kms<sup>-1</sup>; Dwarakanath et al. 1994, O’Dea & Baum 1996) of $`10^{18}10^{19}`$ atom cm<sup>-2</sup>.
Laor (1997) presents tight limits on the column density of atomic hydrogen towards the core of Abell 426 ($`\mathrm{\Delta }N_\mathrm{H}<4\times 10^{17}`$ atom cm<sup>-2</sup>) from Ly$`\alpha `$ observations made with the Hubble Space Telescope (see also Johnstone & Fabian 1995). Together, the results on HI emission and absorption, and Ly$`\alpha `$ absorption, imply that the bulk of the material responsible for the observed X-ray absorption cannot be in the form of atomic hydrogen.
An issue of crucial importance in interpreting the 21cm HI results and the results from CO studies (see below) is the expected temperature of X-ray absorber, if in gaseous form, which has proved controversial. Johnstone, Fabian & Taylor (1998; updating earlier calculations by Ferland, Fabian & Johnstone 1994 and Fabian et al. 1994a) suggest that cooled gas clouds, in equilibrium with the cooling-flow environment, are likely to be very cold. Johnstone et al. (1998) suggest that dust-free clouds form an extended outer envelope, with a column density of $`5\times 10^{21}`$ atom cm<sup>-2</sup>and a temperature of 13-17K, beyond which the gas drops to the microwave background temperature. For Galactic dust/gas ratios, however, the results are significantly modified, such that the cloud temperature drops to the microwave background temperature after a column density of only $`10^{20}`$ atom cm<sup>-2</sup>. At such low temperatures, dusty gas is likely to be highly molecular. However, these predictions differ from those of O’Dea et al. (1994a) and Voit & Donahue (1995) who suggest a minimum temperature for such clouds of $`20`$K, even in the presence of dust. Braine et al. (1995) also suggest a minimum temperature of $`10`$K for the absorbing gas. Possible reasons for the discrepancies between these results (with contrasting views) are discussed by Voit & Donahue (1995) and Ferland et al. (in preparation).
Extensive searches have been carried out for CO associated with molecular gas in cooling flows. To date CO emission has only been detected in Abell 426 (Lazareff et al. 1989; Mirabel, Sanders D.B. & Kazés I. 1989). Braine et al. (1995) also present limits on CO absorption in the central regions of this cluster). The upper limits to the column density of molecular hydrogen in the inner regions of other clusters (for an assumed kinetic temperature of $`T_{\mathrm{CO}}^\mathrm{K}=20`$K and a covering fraction of unity) are typically $`<`$ a few $`10^{20}`$ atom cm<sup>-2</sup>(O’Dea et al. 1994a, Antonucci & Barvainis 1994, McNamara & Jaffe 1994; Braine & Dupraz 1995; O’Dea & Baum 1996) which is significantly below the X-ray inferred column densities. The CO results also constrain the number of clouds along the line of sight (for $`T_{\mathrm{CO}}^\mathrm{K}=20`$K) to be $`<10`$ (O’Dea et al. 1994a; O’Dea & Baum 1996). The CO emission limits appear to exclude the possibility of large column densities (i.e. values consistent with the X-ray measurements) of molecular gas, with a kinetic temperature $`>10`$K. However, if the bulk of the X-ray gas cools to the microwave background temperature, as suggested by Johnstone et al. (1998), the CO limits may be consistent with the X-ray data.
It should be noted that the determinations of molecular hydrogen masses from CO observations are affected by a number of additional uncertainties. Firstly, the CO/H<sub>2</sub> ratios in cluster cooling flows may differ from Galactic values. In particular, where the temperature and metallicity of the absorbing gas is low, the use of a standard Galactic CO/H<sub>2</sub> ratio may underestimate the mass of molecular gas (Maloney & Black 1988; Madden et al. 1997).
It is unclear whether X-ray absorbing gas accumulated by cooling flows could have substantially sub-solar metallicity. Allen & Fabian (1998b) determine a mean emission-weighted metallicity for the X-ray emitting gas in CF clusters of $`Z0.4Z_{}`$. These authors also show that CF clusters generally contain metallicity gradients, so that the metallicity of the X-ray gas within the cooling radii of CF systems will, on average, exceed this value. Thus, it seems unlikely that the material currently being deposited from the cooling flows will have substantially sub-solar metallicities. However, if the X-ray absorbing material is formed from the matter deposited by the cooling flows over their lifetimes, the metallicity of this material could be somewhat lower. Within an inhomogeneous cooling flow (e.g. Nulsen 1986, 1998; Thomas, Fabian & Nulsen 1987) the denser gas at any particular radius will tend to have a smaller volume filling factor. If the metals are evenly distributed throughout the volume of the cooling flow at any particular radius (whether this will occur is unclear), the least dense gas will contain most of the metals (Reisenegger, Miralda-Escudé & Waxman 1996). Since the densest material in cooling flows will have the shortest cooling time and be deposited first, the material accumulated by the cooling flows over their histories could have a significantly lower metallicity than the material being deposited today. Conversely, if the metals in cooling flows were concentrated in the densest material, the metallicity of the X-ray absorbing material could exceed that of the X-ray emitting gas.
Fabian et al. (1994a), have suggested that the masses of molecular gas inferred from CO observations could be significantly underestimated if most of the CO in the absorbing material has frozen onto dust grains. Such freezing should occur rapidly (on timescales $`10^5`$ yr) if even small amounts of dust are initially present in the absorbing gas (Voit & Donahue 1995). Daines et al. (1994) and Fabian et al. (1994a) discuss how dust is likely to form even in clouds that are initially free of dust, on timescales $`10^9`$ yr (although see also Voit & Donahue 1995). Optical, UV and sub-mm studies of the cores of cooling flows show that dust lanes (e.g. Sparks, Macchetto & Golombek 1989; McNamara & O’Connell 1992) and intrinsic reddening (Hu 1992; Allen et al. 1995; Edge et al. 1999) are common.
Elston & Maloney (1994) and Jaffe & Bremer (1997) report detections of H<sub>2</sub>(1-0)S(1) emission from warm ($`2000`$K) molecular hydrogen from K-band spectroscopy of the inner few kpc of a number of CF systems (including Abell 478 and and PKS0745-191). No detections of such emission have been made in NCF systems. The H<sub>2</sub>S emission is likely to arise from the heated skin of cold molecular clouds which could also be responsible for the optical emission-line phenomena discussed above.
A final possibility, first discussed in detail by Voit & Donahue (1995; see also Arnaud & Mushotzky 1998), is that the X-ray absorption could be due to dust, with little or no associated gas. The K-absorption edge of oxygen is a primary source of X-ray absorption and its signature is fairly insensitive to the physical state of the oxygen, so long as it is not highly ionized. The introduction of a simple OIK absorption edge at $`E0.54`$ keV in the ASCA analysis (such as might be associated with oxygen-rich, silicate dust grains) typically provides at least as good a model for the intrinsic absorption as a cold gaseous absorber. Arnaud & Mushotzky (1998) have shown that $`0.357.0`$ keV Broad Band X-ray Telescope data for the Perseus cluster require excess absorption which is significantly better-explained by a simple oxygen absorption edge (which they associate with dust) than by a cold, gaseous absorption model. Verification of this important result will be possible in the near future using observations made with the Chandra Observatory and XMM.
Optical and UV spectroscopy of the central emission-line nebulosities in cooling flows indicate the presence of significant amounts of dust (Hu 1992; Allen et al. 1995; Crawford et al. 1999). Allen et al. (1995) also show that the column densities of absorbing gas inferred from the optical reddening studies are in reasonable agreement with the X-ray values, for Galactic dust/gas ratios. The infrared observations reported here and elsewhere, together with the sub-mm results of Edge et al. (1999), require the presence of significant dust masses in the core regions of at least a few nearby and/or exceptionally massive cooling flows (for which the best data exist) and are consistent with dust being a common feature of CF clusters. Limits on background counts of galaxies and quasars behind clusters are consistent with extinctions ranging from $`A_\mathrm{B}=0.20.5`$ (e.g. Boyle, Fong & Shanks 1988; Romani & Maoz 1992), although measurements of galaxy colours in nearby clusters constrain the reddening on cluster-wide scales to be $`E(BV)<0.06`$ mag (Ferguson 1993).
Voit & Donahue (1995) argue that the levels of dust required to account for the observed X-ray absorption are unlikely to lie entirely within the cores of the clusters, and should be distributed more widely throughout the cluster gas. ROSAT results on the spatial distribution of the X-ray absorbing material (Allen & Fabian 1997) show it to be centrally concentrated within the cooling radii of clusters and possibly be confined within cooling flows. However, the X-ray data are also consistent with a scenario in which the absorbing material is distributed in a manner similar to the X-ray emitting gas. This is illustrated in Fig. 8(a), where we show the equivalent column densities of hydrogen ions in the X-ray emitting gas through two simulated clusters; one CF and one NCF system. The gas distributions have been modeled as isothermal spheres, with core radii of 50 and 300 kpc, respectively (and have been normalized to provide total X-ray gas masses within 3Mpc of the cluster centres of $`2\times 10^{14}`$$`\mathrm{M}_{}`$). We see that the distribution of X-ray emitting gas is much more centrally concentrated in the CF system, particularly within the central $`100`$ kpc.
Fig. 8(b) shows the corresponding X-ray absorbing column densities as a function of radius (as would be determined from observations with ASCA-like CCD detectors in the $`0.610.0`$ keV band) in the case where the intracluster gas has an equivalent metallicity in dust grains of $`0.2`$ solar. For simplicity, we assume that the dust has an absorption spectrum in the ASCA band similar to that of metal-rich cold gas (Morrison & McCammon 1983; qualitatively similar results are obtained for an oxygen edge at 0.54 keV). We also assume that the absorber follows a multilayer distribution within the X-ray emitting medium, but is modeled as a uniform absorbing screen in front of the emitting regions. The results imply mean emission-weighted column densities (where we assume the X-ray emissivity to be proportional to square of the gas density) of $`2\times 10^{21}`$ atom cm<sup>-2</sup> for the CF cluster, and a few $`\times 10^{20}`$ atom cm<sup>-2</sup> for the NCF system. (The absolute results on the excess column densities will depend upon the chemical composition of the dust). We thus see that if the X-ray absorption is due to dust, with a spatial distribution that follows the X-ray gas, then we can expect absorption signatures to be present in both CF and NCF clusters, but to be stronger in CF systems (assuming in each case that an appropriate spectral model is used in the analysis). We would also expect the absorption to be concentrated towards the cores of CF clusters, in agreement with the observations. We note that if the X-ray absorption were entirely due to a pre-existing distribution of dust grains distributed throughout the clusters, then the agreement between the observed masses of absorbing matter (calculated assuming that the absorption is due to cold gas confined within the cooling flows) and the masses predicted to have been accumulated by the cooling flows over their lifetimes (Fig. 6) must be regarded as coincidence. However, it remains possible that both material accumulated by the cooling flows and a more extended distribution of large dust grains (presumably due to supernovae enrichment at early epochs) could contribute to the X-ray absorption.
The lifetime of grains of radius $`a\mu `$m to sputtering in hot gas of density $`n`$ is $`2\times 10^6a/n\mathrm{yr}`$ (Draine & Salpeter 1979). Provided that individual grains exceed 10$`\mu `$m in radius, they should survive for a Hubble time or longer throughout NCF clusters and beyond the cooling radius in clusters with cooling flows. Within cooling flows, the gas density rises inward so the grains are increasingly sputtered, releasing the metals into the gas phase which thus becomes increasingly metal rich towards the cluster centre. This provides one possible explanation for the abundance gradients inferred to be present in most CF clusters (Allen & Fabian 1998b; see also Irwin & Bregman 1999). However, $`10\mu `$m grains are optically thick to soft X-rays and so some distribution of grain sizes extending to smaller values is required if dust is responsible for the observed X-ray absorption (e.g. Laor & Draine 1993). We note that the sputtering of small ($`a<0.1\mu `$m) grains will significantly modify the reddening law in clusters with respect to the standard Galactic relation, reducing the optical/UV reddening (Laor & Draine 1993).
In conclusion, the available data from other wavebands provide some support for the large column densities of intrinsic absorbing material inferred to be present from the X-ray data, at least in the innermost regions of cooling flows. Optical and sub-mm observations of star formation in the cores of CF clusters and K-band observations of emission lines from molecular hydrogen suggest that significant masses of dusty, molecular gas are present in the central few tens of kpc of many CF clusters. The constraints from 21cm observations suggest that the bulk of the material accumulated by cooling flows cannot remain in a long-lived reservoir of atomic hydrogen. Substantial masses of molecular gas may be distributed throughout cooling flows and as yet have avoided detection, but this material must be very cold ($`T3`$K) and dusty, a possibility that remains controversial. The very large X-ray column densities ($`>10^{22}`$ atom cm<sup>-2</sup>) required if the absorbing material has a multilayer distribution are probably incompatible with a gaseous absorber. It remains possible that dust grains, either present in the cooled material deposited by the cooling flows or distributed throughout the clusters in a manner similar to the X-ray gas, may be responsible for much of the observed X-ray absorption.
## 8 Conclusions
The main conclusions that may be drawn from this paper may be summarized as follows:
(i) We have demonstrated the need for multiphase models to consistently explain the spectral and imaging X-ray data for the CF clusters included in our study. The mass deposition rates from the cooling flows, independently inferred from multiphase analyses of the ASCA spectra and deprojection analyses of the ROSAT HRI images, exhibit good agreement, especially once the effects of intrinsic X-ray absorption have been accounted for in a consistent manner. The mass deposition rates from the largest cooling flows exceed 1000 $`\mathrm{M}_{}\mathrm{yr}^1`$, identifying these as some of the most massive cooling flows known.
(ii) We have confirmed the presence of intrinsic X-ray absorption in the cluster spectra using a variety of spectral models. The measured equivalent hydrogen column densities are sensitive to the spectral models used in the analysis, ranging from a few $`10^{20}`$ atom cm<sup>-2</sup> for a simple isothermal emission model to a few $`10^{21}`$ atom cm<sup>-2</sup> using our preferred cooling-flow models, assuming in each case that the absorber lies in a uniform foreground screen. Both the CF and NCF systems exhibit excess absorption at a similar level (on average) when analysed with the same, simple isothermal model.
(iii) The masses of X-ray absorbing material inferred to be present in the CF clusters (assuming the absorption to be due to cold gas with a covering fraction of unity) are in reasonable agreement with the masses expected to have been accumulated by the cooling flows over their lifetimes.
(iv) The ASCA spectra constrain the covering fraction of the absorbing material acting on the cooling flows to be close to (or exceed) unity. If the X-ray absorption is due to many small, similarly-sized clouds along each line of sight, intermixed with the X-ray emitting gas, then the column densities inferred from the spectral analysis may significantly underestimate the true column densities of absorbing material in the cooling flows.
(vii) We have summarized the constraints on the physical properties of the X-ray absorbing material from observations in other wavebands. Substantial 60 and 100$`\mu `$m fluxes, sufficient to account for the bolometric luminosities absorbed and reprocessed within the clusters, are detected from several of the largest cooling flows. Optical observations of star formation in the cores of cooling flows also suggests that significant masses of molecular gas are present in the central few tens of kpc of many CF systems. Constraints from 21cm observations suggest that the bulk of the material accumulated within the cooling radii cannot remain in a long-lived reservoir of atomic hydrogen. Substantial masses of molecular gas may be distributed throughout cooling flows and as yet have avoided detection, although this material must be very cold ($`T3`$K), a possibility that remains controversial. The very large X-ray column densities ($`>10^{22}`$ atom cm<sup>-2</sup>) required if the absorbing material has a multilayer distribution are probably incompatible with a gaseous absorber. It remains possible that dust grains, either present in the material accumulated by the cooling flows or distributed more widely in a manner similar to the X-ray gas, may be responsible for much of the X-ray absorption.
## Acknowledgments
I thank Andy Fabian, Roderick Johnstone and Dave White for many helpful discussions and Harald Ebeling, Alastair Edge and Carolin Crawford for their continuing efforts with the BCS project. I acknowledge the support of the Royal Society. |
warning/0002/hep-ph0002266.html | ar5iv | text | # Calculating three-loop diagrams in heavy quark effective theory with integration-by-parts recurrence relations
## 1 Introduction
Perturbative quantum field theory is progressing fast. New high-precision experiments require calculation of higher radiative corrections — multiloop Feynman diagrams. Recently, some calculations have been done which would seem impossible only a few years ago. This is due to the high degree of automation of the process of generation, analyses and calculation of Feynman diagrams, which is achieved via extensive use of computer algebra (see for review and references). All such calculations are performed in the framework of dimensional regularization , i.e. diagrams are calculated as analytical functions of the space-time dimension $`d=42ϵ`$.
The integration-by-parts method was invented for calculation of three-loop massless propagator diagrams. It is the most systematic method of those currently used, and the most appropriate for computer-algebra implementation. It was first implemented as a SCHOONSCHIP package MINCER , and later re-implemented in FORM (in fact, FORM was created mainly to run MINCER). Since then, MINCER has been the engine behind most of spectacular successes of perturbative field theory. Some of these calculations, with gigabyte-size intermediate expressions, are among the largest computer-algebra calculations ever undertaken.
Integration-by-parts was used for other classes of problems, too. Many interesting physical results have been obtained with packages for calculating two-loop on-shell massive diagrams and three-loop vacuum diagrams with a single mass (see, e.g., ). Reduction of two-loop propagator integrals with generic masses and momentum to a finite set of bases integrals has been achieved . First three-loop on-shell calculations have been done recently .
Several years ago, an interesting new approach to heavy-quark problems in Quantum Chromodynamics has been formulated — Heavy Quark Effective Theory (HQET), see, e.g., for review and references. In collaboration with David Broadhurst, I applied the integration-by-parts method for calculating two-loop propagator diagrams in HQET . Since then, the algorithm suggested was used in a large number of physics applications. A short review of the integration-by-parts method as applied in heavy quark physics is presented in .
In the present work, I apply this method for calculating three-loop propagator diagrams in HQET. Three-loop anomalous dimensions and spectral densities in HQET are necessary for a number of physics applications, such as improved extraction of the $`B`$ meson decay constant $`f_B`$ from lattice simulations and from QCD sum rules. HQET lagrangian does not involve mass in the leading order, and, therefore, the problem is quite similar to the massless one. My aim is to produce a reliable package for three-loop HQET calculations, which could play the same role as MINCER in massless theories. I call it Grinder. Some complication comes from the fact that there are two kinds of lines now — massless propagators and infinitely-heavy ones, and hence the number of diagram topologies is substantially larger.
In section 2, we consider three-loop HQET propagator diagrams with one- or two-loop massless or HQET propagator subdiagrams. To this end, we first recall well-known results for massless and HQET one- and two-loop diagrams. After that, diagrams with two-loop insertions, or with two one-loop insertions, are easily calculated. Two-loop HQET diagrams with a single one-loop insertion are dealt with in a manner similar to the plain two-loop diagrams. In section 3, we consider proper three-loop HQET propagator diagrams. Some details of implementation and testing are presented in section 4.
## 2 Diagrams with lower-loop propagator insertions
### 2.1 Two-loop massless propagator diagrams
The one-loop massless propagator integral (figure 1) can be easily calculated by using the Feynman parameterization or Fourier transform to the coordinate space and back:
$`{\displaystyle \frac{d^dk}{D_1^{n_1}D_2^{n_2}}}`$ $`=`$ $`i\pi ^{d/2}(p^2)^{d/2n_1n_2}G(n_1,n_2),`$
$`D_1`$ $`=`$ $`k^2,D_2=(k+p)^2,`$ (1)
$`G(n_1,n_2)`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }(n_1+n_2d/2)\mathrm{\Gamma }(d/2n_1)\mathrm{\Gamma }(d/2n_2)}{\mathrm{\Gamma }(n_1)\mathrm{\Gamma }(n_2)\mathrm{\Gamma }(dn_1n_2)}}.`$
The integral with a numerator can be written as a finite sum
$`{\displaystyle \frac{P_n(k)d^dk}{D_1^{n_1}D_2^{n_2}}}`$ $`=`$ $`i\pi ^{d/2}(p^2)^{d/2n_1n_2}\times `$ (2)
$`\times {\displaystyle \underset{m}{}}G(n_1,n_2;n,m){\displaystyle \frac{(p^2)^m}{m!}}({\displaystyle \frac{1}{4}}{\displaystyle \frac{}{k_\mu }}{\displaystyle \frac{}{k^\mu }})^mP_n(k)|_{kp},`$
$`G(n_1,n_2;n,m)`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }(n_1+n_2md/2)\mathrm{\Gamma }(d/2n_1+nm)\mathrm{\Gamma }(d/2n_2+m)}{\mathrm{\Gamma }(n_1)\mathrm{\Gamma }(n_2)\mathrm{\Gamma }(dn_1n_2+n)}},`$
where $`P_n(k)`$ is an arbitrary homogeneous polynomial: $`P_n(\lambda k)=\lambda ^nP_n(k)`$.
We write the two-loop propagator integral (figure 2*a*) as
$`{\displaystyle \frac{d^dk_1d^dk_2}{D_1^{n_1}D_2^{n_2}D_3^{n_3}D_4^{n_4}D_5^{n_5}}}=\pi ^d(p^2)^{d{\scriptscriptstyle n_i}}G(n_1,n_2,n_3,n_4,n_5),`$
$`D_1=k_1^2,D_2=k_2^2,D_3=(k_1+p)^2,D_4=(k_2+p)^2,`$
$`D_5=(k_1k_2)^2.`$ (3)
It is symmetric with respect to $`12`$, $`34`$, and also $`13`$, $`24`$. If one of the indices is zero, it can be easily calculated using (1) (figure 2*b*,*c*)
$`G(n_1,n_2,n_3,n_4,0)=G(n_1,n_3)G(n_2,n_4),`$ (4)
$`G(0,n_2,n_3,n_4,n_5)=G(n_3,n_5)G(n_2,n_4+n_3+n_5d/2)`$ (5)
(and symmetric relations).
Applying the operators $`_1(k_1k_2)`$ and $`_1k_1`$ (where $`_i=/k_i`$) to the integrand of (3), we obtain the recurrence relations for $`G(n_1,n_2,n_3,n_4,n_5)`$ (known as triangle relations )
$`\left[dn_1n_32n_5+n_1\mathrm{𝟏}^+(\mathrm{𝟐}^{}\mathrm{𝟓}^{})+n_3\mathrm{𝟑}^+(\mathrm{𝟒}^{}\mathrm{𝟓}^{})\right]G`$ $`=`$ $`0,`$ (6)
$`\left[dn_3n_52n_1+n_3\mathrm{𝟑}^+(1\mathrm{𝟏}^{})+n_5\mathrm{𝟓}^+(\mathrm{𝟐}^{}\mathrm{𝟏}^{})\right]G`$ $`=`$ $`0,`$ (7)
where, for example,
$$\mathrm{𝟏}^\pm G(n_1,n_2,n_3,n_4,n_5)=G(n_1\pm 1,n_2,n_3,n_4,n_5).$$
(8)
Of course, more relations are obtained by symmetry. Another interesting relation is derived by applying the operator $`\frac{}{p}(k_2+p)`$. Substituting the general form of the relevant vector integral, we arrive at
$`[{\displaystyle \frac{1}{2}}d+n_4n_1n_2n_5+({\displaystyle \frac{3}{2}}dn_1n_2n_3n_4n_5)(\mathrm{𝟒}^{}\mathrm{𝟐}^{})+`$ (9)
$`+n_3\mathrm{𝟑}^+(\mathrm{𝟒}^{}\mathrm{𝟓}^{})]G`$ $`=`$ $`0.`$
This formula was derived long ago by S.A. Larin in his M.Sc. thesis (again, similar relations follow by symmetries).
If indices of two adjacent lines are non-positive integers, the integral contains a no-scale vacuum subdiagram and hence vanishes. The cases with zero indices are given by (4), (5). When $`n_5<0`$ and $`n_31`$, $`n_5`$ can be raised by (9); the cases $`n_11`$, $`n_21`$, $`n_41`$ are symmetric. The case $`n_5<0`$, $`n_1=n_1=n_3=n_4=1`$ is handled by
$`\left[(d2n_54)\mathrm{𝟓}^++2(dn_53)\right]G(1,1,1,1,n_5)=`$
$`=2\mathrm{𝟏}^+(\mathrm{𝟑}^{}\mathrm{𝟐}^{}\mathrm{𝟓}^+)G(1,1,1,1,n_5),`$ (10)
which follows from (6) and (7) at $`n_1=n_2=n_3=n_4=1`$ (note that the terms in the right-hand side of (10) are trivial for any $`n_5`$; for $`n_5<0`$, they vanish). When $`n_2<0`$, it can be raised by (9); the cases $`n_1<0`$, $`n_3<0`$, $`n_4<0`$ aresymmetric.
We are left with the most important situation when all the indices are positive. Applying (6), we reduce $`n_2`$, $`n_4`$, $`n_5`$ until one of them vanishes. Then (4) and (5) apply. If $`\mathrm{max}(n_1,n_3)<\mathrm{max}(n_2,n_4)`$, it is more efficient to lower $`n_1`$, $`n_3`$, $`n_5`$.
All one-loop integrals (figure 1) with integer $`n_{1,2}`$ are proportional to $`G_1=G(1,1)`$, the coefficient being a rational function of $`d`$. All two-loop integrals with integer indices reduce to $`G_1^2`$ (figure 3*a*) and $`G_2=G(0,1,1,0,1)`$ (figure 3*b*), with rational coefficients. Here
$$G_n=\frac{1}{\left(n+1n\frac{d}{2}\right)_n\left((n+1)\frac{d}{2}2n1\right)_n}\frac{\mathrm{\Gamma }(1+nϵ)\mathrm{\Gamma }^{n+1}(1ϵ)}{\mathrm{\Gamma }(1(n+1)ϵ)},$$
(11)
where $`(x)_n=\mathrm{\Gamma }(x+n)/\mathrm{\Gamma }(x)`$ is the Pochhammer symbol.
### 2.2 Two-loop HQET propagator diagrams
The one-loop HQET propagator integral (figure 4) can be easily calculated by using the modified Feynman parameterization (see, e.g., ), or Fourier transform to the coordinate space and back:
$`{\displaystyle \frac{d^dk}{D_1^{n_1}D_2^{n_2}}}`$ $`=`$ $`i\pi ^{d/2}(2\omega )^{d2n_2}I(n_1,n_2),`$
$`D_1`$ $`=`$ $`{\displaystyle \frac{(k+p)v}{\omega }},D_2=k^2,`$
$`I(n_1,n_2)`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }(n_1+2n_2d)\mathrm{\Gamma }(d/2n_2)}{\mathrm{\Gamma }(n_1)\mathrm{\Gamma }(n_2)}}`$ (12)
(here $`v`$ is 4-velocity of the heavy quark, $`v^2=1`$, and $`\omega =pv`$ is the residual energy). Similarly to (2), we obtain
$`{\displaystyle \frac{P_n(k)d^dk}{D_1^{n_1}D_2^{n_2}}}`$ $`=`$ $`i\pi ^{d/2}(2\omega )^{d2n_2}\times `$
$`\times {\displaystyle \underset{m}{}}I(n_1,n_2;n,m){\displaystyle \frac{(2\omega )^{2m}}{m!}}({\displaystyle \frac{1}{4}}{\displaystyle \frac{}{k_\mu }}{\displaystyle \frac{}{k^\mu }})^mP_n(k)|_{k2\omega v},`$
$`I(n_1,n_2;n,m)`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }(n_1+2n_2nd)\mathrm{\Gamma }(d/2n_2+nm)}{\mathrm{\Gamma }(n_1)\mathrm{\Gamma }(n_2)}}.`$ (13)
There are two topologies of two-loop propagator HQET diagrams (figure 5*a*,*b*). We write the first of them as
$$\begin{array}{ccccccc}\hfill \frac{d^dk_1d^dk_2}{D_1^{n_1}D_2^{n_2}D_3^{n_3}D_4^{n_4}D_5^{n_5}}& =& \multicolumn{5}{c}{\pi ^d(2\omega )^{2(dn_3n_4n_5)}I(n_1,n_2,n_3,n_4,n_5),}\\ \hfill D_1& =& \frac{(k_1+p)v}{\omega },\hfill & & \hfill D_2& =& \frac{(k_2+p)v}{\omega },\hfill \\ \hfill D_3& =& k_1^2,\hfill & & \hfill D_4& =& k_2^2,D_5=(k_1k_2)^2.\hfill \end{array}$$
(14)
It is symmetric with respect to $`13`$, $`24`$. If one of the indices is zero, it can be easily calculated using (12) and (1) (figure 5*c*,*d*,*e*)
$`I(n_1,n_2,n_3,n_4,0)`$ $`=`$ $`I(n_1,n_3)I(n_2,n_4),`$ (15)
$`I(0,n_2,n_3,n_4,n_5)`$ $`=`$ $`G(n_3,n_5)I(n_2,n_4+n_3+n_5d/2),`$ (16)
$`I(n_1,n_2,0,n_4,n_5)`$ $`=`$ $`I(n_1,n_5)I(n_2+n_1+2n_5d,n_4)`$ (17)
(and symmetric relations).
Applying the operators $`_1(k_1k_2)`$, $`_1k_1`$ and $`_1v`$ to the integrand of (14), we obtain
$`\left[dn_1n_32n_5+n_1\mathrm{𝟏}^+\mathrm{𝟐}^{}+n_3\mathrm{𝟑}^+(\mathrm{𝟒}^{}\mathrm{𝟓}^{})\right]I`$ $`=`$ $`0,`$ (18)
$`\left[dn_1n_52n_3+n_1\mathrm{𝟏}^++n_5\mathrm{𝟓}^+(\mathrm{𝟒}^{}\mathrm{𝟑}^{})\right]I`$ $`=`$ $`0,`$ (19)
$`\left[2n_1\mathrm{𝟏}^++n_3\mathrm{𝟑}^+(\mathrm{𝟏}^{}1)+n_5\mathrm{𝟓}^+(\mathrm{𝟏}^{}\mathrm{𝟐}^{})\right]I`$ $`=`$ $`0.`$ (20)
Of course, more relations are obtained by symmetry. Applying $`\omega \frac{d}{d\omega }`$ and using homogeneity in $`\omega `$, we obtain
$$\left[2(dn_3n_4n_5)n_1n_2+n_1\mathrm{𝟏}^++n_2\mathrm{𝟐}^+\right]I=0,$$
(21)
which is nothing but the sum of (19) and its mirror-symmetric. Subtracting the $`\mathrm{𝟐}^{}`$ shifted version of (21) from (18), we obtain the most useful relation
$`[dn_1n_2n_32n_5+1`$
$`(2(dn_3n_4n_5)n_1n_2+1)\mathrm{𝟐}^{}+n_3\mathrm{𝟑}^+(\mathrm{𝟒}^{}\mathrm{𝟓}^{})]I`$ $`=`$ $`0,`$ (22)
which lowers $`n_2`$, $`n_4`$, $`n_5`$, and does not raise heavy-quark indices.
If indices of two adjacent lines are non-positive integers, the integral contains a no-scale vacuum subdiagram and hence vanishes. The cases with zero indices are given by (15), (16) and (17). When $`n_2<0`$, it can be raised by (22); the case $`n_1<0`$ is symmetric. Similarly, if $`n_3<0`$, it can be raised by (22); the case $`n_4<0`$ is symmetric. When $`n_5<0`$ and $`n_31`$, $`n_5`$ can be raised by (22) (the case $`n_41`$ is symmetric); when $`n_5<0`$ and $`n_11`$, $`n_5`$ can be raised by (19) (the case $`n_21`$ is symmetric). The case $`n_5<0`$, $`n_1=n_1=n_3=n_4=1`$ is handled by
$`\left[(d2n_54)\mathrm{𝟓}^+2(dn_53)\right]I(1,1,1,1,n_5)=`$
$`=\left[(2d2n_57)\mathrm{𝟏}^{}\mathrm{𝟓}^+\mathrm{𝟑}^{}\mathrm{𝟒}^+\mathrm{𝟓}^{}+\mathrm{𝟏}^{}\mathrm{𝟑}^+\right]I(1,1,1,1,n_5),`$ (23)
which follows from (18), (19) and (20) at $`n_1=n_2=n_3=n_4=1`$ (note that the terms on the right-hand side of (23) are trivial for any $`n_5`$; for $`n_5<0`$, they vanish).
We are left with the most important situation when all the indices are positive. Applying (22), we reduce $`n_2`$, $`n_4`$, $`n_5`$ until one of them vanishes. Then (15), (16) and (17) apply. If $`\mathrm{max}(n_1,n_3)<\mathrm{max}(n_2,n_4)`$, it is more efficient to lower $`n_1`$, $`n_3`$, $`n_5`$.
In the second topology (figure 5*b*), three heavy-quark denominators depend on only two variables $`k_{1,2}v`$, hence they are linearly dependent. Therefore, there is one scalar product which cannot be expressed via the denominators. Let’s define the integral
$$\begin{array}{ccccccccccc}\multicolumn{11}{c}{\frac{N^{n_0}d^dk_1d^dk_2}{D_1^{n_1}D_2^{n_2}D_3^{n_3}D_4^{n_4}D_5^{n_5}}=\pi ^d(2\omega )^{2(d+n_0n_4n_5)}J(n_1,n_2,n_3,n_4,n_5;n_0),}\\ \hfill D_1& =& \frac{(k_1+p)v}{\omega },\hfill & & \hfill D_2& =& \frac{(k_2+p)v}{\omega },\hfill & & \hfill D_3& =& \frac{(k_1+k_2+p)v}{\omega },\hfill \\ \hfill D_4& =& k_1^2,\hfill & & \hfill D_5& =& k_2^2,\hfill & & \hfill N& =& 2k_1k_2\hfill \end{array}$$
(24)
(it is symmetric with respect to $`12`$, $`45`$). Noting that $`D_1+D_2D_3=1`$, we immediately have
$$(1\mathrm{𝟏}^{}\mathrm{𝟐}^{}+\mathrm{𝟑}^{})J=0.$$
(25)
Applying $`_1k_1`$, $`_1k_2`$ and $`_1v`$ to the integrand, we have
$`\left[d+n_0n_1n_32n_4+n_1\mathrm{𝟏}^++n_3\mathrm{𝟑}^+\mathrm{𝟐}^{}\right]J`$ $`=`$ $`0`$ (26)
$`\left[n_1n_3n_1\mathrm{𝟏}^+\mathrm{𝟑}^{}+n_3\mathrm{𝟑}^+\mathrm{𝟏}^{}+n_4\mathrm{𝟒}^+\mathrm{𝟎}^+2n_0\mathrm{𝟎}^{}\mathrm{𝟓}^{}\right]J`$ $`=`$ $`0`$ (27)
$`\left[2n_1\mathrm{𝟏}^+2n_3\mathrm{𝟑}^++n_4\mathrm{𝟒}^+(\mathrm{𝟏}^{}1)+n_0\mathrm{𝟎}^{}(\mathrm{𝟐}^{}1)\right]J`$ $`=`$ $`0.`$ (28)
Homogeneity in $`\omega `$ gives $`[2(d+n_0n_4n_5)n_1n_2n_3+n_1\mathrm{𝟏}^++n_2\mathrm{𝟐}^++n_3\mathrm{𝟑}^+]J=0`$, which is nothing but the sum of (26) and its mirror-symmetric. The boundary values of the integral (figure 5*c*,*e*) are
$`J(n_1,n_2,0,n_4,n_5;n_0)`$ $`=`$ $`(\mathrm{𝟓}^{}\mathrm{𝟑}^{}\mathrm{𝟒}^{})^{n_0}I(n_1,n_2,n_4,n_5,0)`$ (29)
$`J(0,n_2,n_3,n_4,n_5;n_0)`$ $`=`$ $`(\mathrm{𝟒}^{}\mathrm{𝟑}^{}+\mathrm{𝟓}^{})^{n_0}I(n_3,n_2,0,n_5,n_4)`$ (30)
and the symmetric relation for $`n_2=0`$. If $`n_0=0`$,
$`J(n_1,n_2,0,n_4,n_5)`$ $`=`$ $`I(n_1,n_4)I(n_2,n_5)`$ (31)
$`J(0,n_2,n_3,n_4,n_5)`$ $`=`$ $`I(n_3,n_4)I(n_2+n_3+2n_4d,n_5).`$ (32)
If $`n_40`$, or $`n_50`$, or two adjacent heavy-quark indices are non-positive, the integral vanishes. If any of $`n_1`$, $`n_2`$, $`n_3`$ is negative, it can be raised by (25). If all of them are positive, we use (25) to lower $`n_1`$, $`n_2`$ or $`n_3`$, until one of these indices vanish.
Instead of using (29), (30) when $`n_0>0`$, we could proceed in another way. If $`n_4>1`$, we lower both $`n_0`$ and $`n_4`$ using (27) (the case $`n_5>1`$ is symmetric). We are left with $`J(n_1,n_2,0,1,1;n_0)`$ and $`J(0,n_2,n_3,1,1;n_0)`$. In the first case, if $`n_1>1`$, we lower it using (28) (the case $`n_2>1`$ is symmetric). In the second case, if $`n_3>1`$, we lower it using (28); if $`n_2>1`$, we lower it using the difference of (28) and its mirror-symmetric. The two remaining cases (figure 5*c*,*e*) are easily evaluated using (13):
$`J(1,1,0,1,1;n_0)`$ $`=`$ $`(1)^{n_0}I(1,1;n_0,0)I(1,1),`$
$`J(0,1,1,1,1;n_0)`$ $`=`$ $`I(1,1;n_0,0){\displaystyle \underset{l=0}{\overset{n_0}{}}}{\displaystyle \frac{(1)^ln_0!}{l!(n_0l)!}}I(42n_0+ld,1)`$
All one-loop integrals (figure 4) with integer $`n_{1,2}`$ are proportional to $`I_1=I(1,1)`$, the coefficient being a rational function of $`d`$. All two-loop integrals with integer indices reduce to $`I_1^2`$ (figure 6*a*) and $`I_2=I(0,1,1,0,1)`$ (figure 6*b*), with rational coefficients. Here
$$I_n=\frac{1}{(1n(d2))_{2n}}\mathrm{\Gamma }(1+2nϵ)\mathrm{\Gamma }^n(1ϵ).$$
(33)
### 2.3 Three-loop HQET diagrams with lower-loop insertions
Diagrams with a two-loop propagator insertion, or with two one-loop insertions, figure 7, are trivially calculated by multiplying the relevant insertion\[s\] by the one-loop HQET integral with non-integer indices, whose values are obvious by dimensionality.
Next we consider the diagrams obtained from the two-loop HQET diagram of figure 5*a* by adding a single one-loop propagator insertion (figure 8). They are equal to the product of the corresponding one-loop integral and the integral of figure 5*a* with one non-integer index. All relations derived for the diagram of figure 5*a* are valid; however, the non-integer index changes the strategy of their application.
Let’s consider the case of figure 8*a* with a non-integer $`n_3`$. When $`n_5<0`$, it can always be raised by (22). When $`n_2<0`$, it can be raised by (22), too; when $`n_1<0`$, it can be raised by the mirror-symmetric relation. When $`n_4<0`$, it can always be raised by (22), too. Finally, when all indices $`n_1`$, $`n_2`$, $`n_4`$, $`n_5`$ are positive, we use (22) to lower $`n_2`$, $`n_4`$ or $`n_5`$, until (15), (16) and (17) are reached.
In the case of figure 8*b*, $`n_1`$ is non-integer. When $`n_5<0`$ and $`n_31`$, we can raise $`n_5`$ by (22) (the case $`n_41`$ is symmetric); when $`n_5<0`$ and $`n_21`$, we can lower $`n_2`$ by the relation symmetric to (19); when $`n_5<0`$ and $`n_2=n_3=n_4=1`$, we can use the relation
$$\left[2(dn_53)\mathrm{𝟒}^+n_5\mathrm{𝟓}^+\mathrm{𝟏}^{}\right]I(n_1,1,1,1,n_5)=0,$$
(34)
which follows from (20) and (19), and then the term with $`\mathrm{𝟒}^+`$ can be treated as above. When $`n_2<0`$, it can be raised by (22). When $`n_3<0`$ and $`n_41`$, we can raise $`n_3`$ by the relation symmetric to (22); when $`n_3<0`$ and $`n_51`$, we can raise $`n_3`$ by (19); when $`n_3<0`$ and $`n_21`$, we lower $`n_2`$ by (21); finally, when $`n_3<0`$ and $`n_2=n_4=n_5=1`$, we use (22) to get a trivial term with $`n_2=0`$. When $`n_4<0`$ and $`n_31`$, we can raise $`n_4`$ by (22); when $`n_4<0`$ and $`n_51`$, we can raise $`n_4`$ by (19); when $`n_4<0`$ and $`n_21`$, we lower $`n_2`$ by (21); finally, when $`n_4<0`$ and $`n_2=n_3=n_5=1`$, we use (20) to raise $`n_3`$ or $`n_5`$, and proceed as above. Finally, when all indices $`n_2`$, $`n_3`$, $`n_4`$, $`n_5`$ are positive, we use (22) to lower $`n_2`$, $`n_4`$ or $`n_5`$, until (15), (16) and (17) are reached.
In the case of figure 8*c*, $`n_5`$ is non-integer. When $`n_2<0`$, it can be raised by (22); the case $`n_1<0`$ is symmetric. When $`n_3<0`$, it can be raised by (19); the case $`n_4<0`$ is symmetric. When $`n_1>1`$, it can be lowered by (19); the case $`n_2>1`$ is symmetric. When $`n_3>1`$, it can be lowered by (20), with a trivial additional term having $`n_1=0`$; the case $`n_4>1`$ is symmetric. We are left with $`n_1=n_2=n_3=n_4=1`$; the relation (23) can be used to lower or raise $`n_5`$, with trivial additional terms. An integral with some specific value of $`n_5`$ (of the form integer plus $`ϵ`$) has to be considered as a new basis element. This integral has been calculated, exactly in $`d`$ dimensions, in in terms of $`{}_{3}{}^{}F_{2}^{}`$ hypergeometric functions, using the Hegenbauer polynomial technique in the coordinate space .
Finally, we consider the diagrams obtained from the two-loop HQET diagram of figure 5*b* by adding a single one-loop propagator insertion (figure 9). The diagram of figure 9*a*, with a non-integer $`n_4`$ is calculated exactly as the two-loop one.
In the case of figure 9*b*, $`n_1`$ is non-integer. When $`n_3<0`$ or $`n_2<0`$, they can be raised by (25). In the case of positive indices, we use
$$\left[dn_32n_4+n_1\mathrm{𝟏}^+(\mathrm{𝟐}^{}\mathrm{𝟑}^{})+n_3\mathrm{𝟑}^+\mathrm{𝟐}^{}\right]J=0$$
(35)
(which follows from (26) and (25)). It either lowers $`n_2+n_3`$, or, at a fixed $`n_2+n_3`$, lowers $`n_2`$. Therefore, sooner or later, we reach (29) and (30).
In the case of figure 9*c*, $`n_3`$ is non-integer. If $`n_1<0`$ or $`n_2<0`$, they can be raised by (25). When $`n_4>1`$, we can lower it or $`n_1`$ by (28); the case $`n_5>1`$ is symmetric. When $`n_1>1`$, it can be lowered by (26); the case $`n_2>1`$ is symmetric. We are left with $`n_1=n_2=n_4=n_5=1`$; the relation (25) can be used to lower or raise $`n_3`$, with trivial additional terms. An integral with some specific value of $`n_3`$ (of the form integer plus $`2ϵ`$) has to be considered as a new basis element.
It is not difficult to calculate $`J(1,1,n_3,n_4,n_5)`$ for arbitrary $`n_{3,4,5}`$ (not necessarily integer) in the coordinate space:
$`J(1,1,n,n_1,n_2)`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }(n2(dn_1n_21))\mathrm{\Gamma }(d/2n_1)\mathrm{\Gamma }(d/2n_2)}{\mathrm{\Gamma }(n)\mathrm{\Gamma }(n_1)\mathrm{\Gamma }(n_2)}}J,`$
$`J`$ $`=`$ $`t^{2(dn_1n_2)n1}{\displaystyle \underset{0<t_1<t_2<t}{}}𝑑t_1𝑑t_2t_2^{2n_1d}(tt_1)^{2n_2d}(t_2t_1)^{n1}`$ (38)
$`=`$ $`{\displaystyle \frac{1}{n(2n_1+n+1d)}}_3F_2\left(\begin{array}{c}1,d2n_2,\mathrm{\hspace{0.33em}2}n_1+n+1d\\ n+1,\mathrm{\hspace{0.33em}2}n_1+n+2d\end{array}|\mathrm{\hspace{0.17em}1}\right).`$
All diagrams considered in this section are particular cases of the generic three-loop topologies, which will be discussed in section 3, when some lines are shrunk (i.e., some indices vanish). Therefore, we don’t consider diagrams with numerators here: numerators should be dealt with in the context of generic topologies, and the formulae of this section are used only as boundary values for the corresponding recurrence relations, after elimination of numerators.
All three-loop HQET propagator integrals with lower-loop propagator insertions are linear combinations of 7 basis integrals (figure 10*a**g*), coefficients being rational functions of $`d`$. The basis integrals of figure 10*a* ($`I_1^3`$), figure 10*b* ($`I_1I_2`$), figure 10*c* ($`I_3`$), figure 10*d* ($`I_3I_1^2/I_2`$), and figure 10*e* ($`I_3G_1^2/G_2`$) are known exactly in $`d`$ dimensions in terms of $`\mathrm{\Gamma }`$ functions. Those of figure 10*f*,*g* contain hypergeometric $`{}_{3}{}^{}F_{2}^{}`$ functions of the unit argument. Several terms of their expansion in $`ϵ`$ can be obtained using the methods which were recently developed in . As we shall see in section 3, there is only one additional basis integral, figure 10*h*.
## 3 Proper three-loop HQET propagator diagrams
In the massless case , there are only 3 topologies of proper three-loop propagator diagrams: Mercedez, Ladder, and Non-planar (plus their reduced forms obtained by shrinking some lines). Now we have 10 topologies instead (figure 11). Each of them has 8 propagators. In the diagrams with 4 heavy-quark lines (figure 11*e**g*), there is one linear dependence between their denominators; with 5 heavy-quark lines (figure 11*h**j*) — 2 dependences. There are 9 independent scalar products of 3 loop momenta $`k_{1,2,3}`$ and the 4-velocity $`v`$. Therefore, in the diagrams with two or three heavy-quark lines (figure 11*a**d*), there is one scalar product in the numerator which cannot be cancelled against the denominators; with 4 heavy-quark lines (figure 11*e**g*) — two scalar products; with 5 heavy-quark lines (figure 11*h**j*) — three scalar products.
When calculating these diagrams using recurrence relations, some indices may vanish. This corresponds to shrinking the corresponding lines. In some cases, this results in diagrams with lower-loop propagator insertions, which were calculated in section 2. The diagrams of figure 12 are still non-trivial.
### 3.1 Diagram with two heavy-quark lines
Let’s consider the diagram of figure 11*a* first. We define
$$\begin{array}{ccccccccccc}\multicolumn{7}{c}{\frac{N^{n_0}d^dk_1d^dk_2d^dk_3}{D_1^{n_1}D_2^{n_2}D_3^{n_3}D_4^{n_4}D_5^{n_5}D_6^{n_6}D_7^{n_7}D_8^{n_8}}=}& & & & \\ & =& \multicolumn{9}{c}{i\pi ^{3d/2}(2\omega )^{3d+2n_02_{i=3}^8n_i}I_a(n_1,n_2,n_3,n_4,n_5,n_6,n_7,n_8;n_0),}\\ \hfill N& =& \frac{k_3v}{\omega },\hfill & & \hfill D_1& =& \frac{(k_1+p)v}{\omega },\hfill & & \hfill D_2& =& \frac{(k_2+p)v}{\omega },\hfill \\ \hfill D_3& =& k_1^2,\hfill & & \hfill D_4& =& k_2^2,\hfill & & \hfill D_5& =& (k_1k_2)^2,\hfill \\ \hfill D_6& =& k_3^3,\hfill & & \hfill D_7& =& (k_3+k_1)^2,\hfill & & \hfill D_8& =& (k_3+k_2)^2.\hfill \end{array}$$
This integral is mirror-symmetric with respect to $`12`$, $`34`$, $`78`$. It vanishes when the indices of the following groups of lines are non-positive: 12, 67, 68, 78, 375, 485, 315, 425, 364, 137, 248, 157, or 258.
First we are going to get rid of the numerator. When $`n_0>0`$ and $`n_71`$, we can lower $`n_0`$ by
$$\left[2n_1\mathrm{𝟏}^++n_3\mathrm{𝟑}^+(\mathrm{𝟏}^{}1)+n_5\mathrm{𝟓}^+(\mathrm{𝟏}^{}\mathrm{𝟐}^{})+n_7\mathrm{𝟕}^+(\mathrm{𝟏}^{}1+\mathrm{𝟎}^+)\right]I_a=0,$$
(39)
which is obtained by applying $`_1v`$ to the integrand of (3.1); the case $`n_0>0`$, $`n_81`$ is symmetric. When $`n_0>0`$ and $`n_61`$, we can lower $`n_0`$ by
$`[2[3d+n_0n_1n_22(n_3+n_4+n_5+n_6+n_7+n_8)]+`$
$`+n_3\mathrm{𝟑}^+(\mathrm{𝟏}^{}1)+n_4\mathrm{𝟒}^+(\mathrm{𝟐}^{}1)n_6\mathrm{𝟔}^+\mathrm{𝟎}^+2n_0\mathrm{𝟎}^{}]I_a`$ $`=`$ $`0,`$ (40)
which is the $`(_1+_2_3)v`$ relation simplified using the homogeneity relation
$$\left[3d+n_0n_1n_22(n_3+n_4+n_5+n_6+n_7+n_8)+n_1\mathrm{𝟏}^++n_2\mathrm{𝟐}^+\right]I_a=0.$$
(41)
When $`n_0>0`$ and $`n_11`$, we can lower $`n_0`$ by
$`[dn_1n_3n_52n_7+n_1\mathrm{𝟏}^+(1\mathrm{𝟎}^+)+`$
$`+n_3\mathrm{𝟑}^+(\mathrm{𝟔}^{}\mathrm{𝟕}^{})+n_5\mathrm{𝟓}^+(\mathrm{𝟖}^{}\mathrm{𝟕}^{})]I_a`$ $`=`$ $`0,`$ (42)
which is obtained by applying $`_1(k_1+k_3)`$ to the integrand of (3.1); the case $`n_0>0`$, $`n_21`$ is symmetric. We are left with $`n_0>0`$, $`n_1=n_2=n_6=n_7=n_8=1`$; we use $`_3(k_3+k_1)`$ relation
$`[d+n_0n_6n_82n_7+n_0\mathrm{𝟎}^{}(\mathrm{𝟏}^{}1)+`$
$`+n_6\mathrm{𝟔}^+(\mathrm{𝟑}^{}\mathrm{𝟕}^{})+n_8\mathrm{𝟖}^+(\mathrm{𝟓}^{}\mathrm{𝟕}^{})]I_a`$ $`=`$ $`0`$ (43)
to lower $`n_0`$ or raise $`n_6`$ or $`n_8`$, and apply the method described above again.
Now we shall discuss the integral $`I_a(n_1,n_2,n_3,n_4,n_5,n_6,n_7,n_8)`$ without numerator ($`n_0=0`$). Applying $`_3k_3`$, $`_1k_1`$, $`_1(k_1k_2)`$, $`(_1+_2_3)k_1`$ to the integrand of (3.1), we obtain the recurrence relations
$`\left[dn_7n_82n_6+n_7\mathrm{𝟕}^+(\mathrm{𝟑}^{}\mathrm{𝟔}^{})+n_8\mathrm{𝟖}^+(\mathrm{𝟒}^{}\mathrm{𝟔}^{})\right]I_a`$ $`=`$ $`0,`$ (44)
$`[dn_1n_5n_72n_3+n_1\mathrm{𝟏}^++`$
$`+n_5\mathrm{𝟓}^+(\mathrm{𝟒}^{}\mathrm{𝟑}^{})+n_7\mathrm{𝟕}^+(\mathrm{𝟔}^{}\mathrm{𝟑}^{})]I_a`$ $`=`$ $`0,`$ (45)
$`[dn_1n_3n_72n_5+n_1\mathrm{𝟏}^+\mathrm{𝟐}^{}+`$
$`+n_3\mathrm{𝟑}^+(\mathrm{𝟒}^{}\mathrm{𝟓}^{})+n_7\mathrm{𝟕}^+(\mathrm{𝟖}^{}\mathrm{𝟓}^{})]I_a`$ $`=`$ $`0,`$ (46)
$`[2(dn_5n_7n_8)n_2n_4n_6+n_2\mathrm{𝟐}^+\mathrm{𝟏}^{}+`$
$`+n_4\mathrm{𝟒}^+(\mathrm{𝟑}^{}\mathrm{𝟓}^{})+n_6\mathrm{𝟔}^+(\mathrm{𝟑}^{}\mathrm{𝟕}^{})]I_a`$ $`=`$ $`0,`$ (47)
where the last relation was simplified using (41).
The cases $`n_1=0`$, $`n_2=0`$, $`n_6=0`$, $`n_7=0`$, $`n_8=0`$ are trivial. When $`n_1<0`$, it can be raised by
$`\left[3dn_1n_22(n_3+n_4+n_5+n_6+n_7+n_8)\right]I_a=`$ (48)
$`=\left[dn_1n_2n_4n_82n_5+n_4\mathrm{𝟒}^+(\mathrm{𝟑}^{}\mathrm{𝟓}^{})+n_8\mathrm{𝟖}^+(\mathrm{𝟕}^{}\mathrm{𝟓}^{})\right]\mathrm{𝟏}^+I_a,`$
which is the difference of (41) and $`\mathrm{𝟏}^+`$ shifted version of the relation symmetric to (46); the case $`n_2<0`$ is symmetric. When $`n_6<0`$ and $`n_71`$, we can raise $`n_6`$ by (45) (the case $`n_81`$ is symmetric); when $`n_6<0`$ and $`n_7=n_8=1`$, we can raise $`n_6`$ or $`n_8`$ by (43). When $`n_7<0`$ and $`n_81`$, we can raise $`n_7`$ by the relation symmetric to (46); when $`n_7<0`$ and $`n_61`$, we can raise $`n_7`$ by (47); when $`n_7<0`$ and $`n_6=n_8=1`$, we can raise $`n_7`$ or $`n_6`$ by the relation symmetric to (43). The case $`n_8<0`$ is symmetric. The case $`n_3=0`$ (figure 12*a*, $`J_a(n_1,n_2,n_4,n_5,n_6,n_7,n_8)`$) will be considered later in this section; the case $`n_4=0`$ is symmetric. When $`n_3<0`$ and $`n_61`$, we can raise $`n_3`$ by (43); when $`n_3<0`$ and $`n_71`$, we can raise $`n_3`$ by (44); when $`n_3<0`$ and $`n_6=n_7=1`$, we can raise $`n_6`$ or $`n_7`$ by the relation symmetric to (43). The case $`n_4<0`$ is symmetric. The case $`n_5=0`$ (figure 12*b*, $`J_b(n_1,n_2,n_3,n_4,n_6,n_7,n_8)`$) will be considered later in this section. When $`n_5<0`$ and $`n_81`$, we can raise $`n_5`$ by (43) (the case $`n_71`$ is symmetric); when $`n_5<0`$ and $`n_7=n_8=1`$, we can raise $`n_7`$ or $`n_8`$ by (44). When all the indices are positive, we can use (44) to kill one of the lines 3, 4, 6; or we can use (43) to kill one of the lines 3, 5, 7; or we can use the symmetric relation to kill one of the lines 4, 5, 8; or we can use (47) to kill one of the lines 1, 3, 5, 7; or we can use the symmetric relation to kill one of thelines 2, 4, 5, 8.
Now we consider $`J_a(n_1,n_2,n_4,n_5,n_6,n_7,n_8)=I_a(n_1,n_2,0,n_3,n_4,n_5,n_6,n_7,n_8)`$ (figure 12*a* ). This integral vanishes when the indices of the following groups of lines are non-positive: 12, 67, 68, 78, 57, 15, 17, 46, 485, 425, 248, or 258. If any index is zero, the integral becomes trivial. If $`n_1<0`$, $`n_2<0`$, $`n_6<0`$, $`n_7<0`$, or $`n_8<0`$, we just consider $`J_a`$ as $`I_a`$ with $`n_3=0`$ and proceed as usual; the integral reduces to trivial ones not including $`J_a`$. When $`n_4>0`$, we can kill one of the lines 2, 4, 8 using the relation symmetric to (47). When $`n_5>0`$, we can kill one of the lines 2, 5, 8 using (46). When both $`n_4`$ and $`n_5`$ are negative, we proceed as follows. If $`n_61`$, $`n_4`$ can be raised by the relation symmetric to (47); if $`n_71`$, $`n_5`$ can be raised by (46); $`n_2`$ can be lowered down to 1 using (41). The remaining integral can be calculated by a repeated use of (13):
$`J_a(n_1,1,|n_4|,|n_5|,1,1,n_8)=`$
$`=`$ $`{\displaystyle \underset{l_1,l_2,l_3,m}{}}{\displaystyle \frac{(1)^{l_1+l_3}|n_4|!|n_5|!}{(|n_5|l_1)!m!(l_12m)!(|n_4|l_2)!l_3!(l_2l_3)!}}\times `$
$`\times I(n_1,1;l_1,m)I(1,n_8|n_4||n_5|+l_1+l_2m;l_2,0)\times `$
$`\times I(2(d+|n_4|+|n_5|n_8)+n_1+l_3+3,1).`$
Finally, we consider $`J_b(n_1,n_2,n_3,n_4,n_6,n_7,n_8)=I_a(n_1,n_2,n_3,n_4,0,n_6,n_7,n_8)`$ (figure 12*b*). This integral is mirror-symmetric; it vanishes when the indices of the following groups of lines are non-positive: 12, 67, 68, 78, 37, 48, 13, 17, 24, 28, or 364. If any index is zero, the integral becomes trivial. If $`n_1<0`$, $`n_2<0`$, $`n_3<0`$, $`n_4<0`$, $`n_6<0`$, $`n_7<0`$, or $`n_8<0`$, we just consider $`J_b`$ as $`I_a`$ with $`n_5=0`$ and proceed as usual; the integral reduces to trivial ones not including $`J_b`$ (though, possibly, including $`J_a`$). When all the indices are positive, one of the lines 3, 4, 6 can be killed by (44).
### 3.2 Mercedez with three heavy-quark lines
Now let’s consider the diagram of figure 11*b*. We define
$$\begin{array}{ccccccccccc}\multicolumn{5}{c}{\frac{N^{n_0}d^dk_1d^dk_2d^dk_3}{D_1^{n_1}D_2^{n_2}D_3^{n_3}D_4^{n_4}D_5^{n_5}D_6^{n_6}D_7^{n_7}D_8^{n_8}}=}& & & & & & \\ & =& \multicolumn{9}{c}{i\pi ^{3d/2}(2\omega )^{3d+2n_02_{i=4}^8n_i}I_b(n_1,n_2,n_3,n_4,n_5,n_6,n_7,n_8;n_0),}\\ \hfill D_1& =& \frac{(k_1+p)v}{\omega },\hfill & & \hfill D_2& =& \frac{(k_2+p)v}{\omega },\hfill & & \hfill D_3& =& \frac{(k_3+p)v}{\omega },\hfill \\ \hfill D_4& =& k_1^2,\hfill & & \hfill D_5& =& k_2^2,\hfill & & \hfill D_6& =& (k_1k_3)^2,\hfill \\ \hfill D_7& =& (k_2k_3)^2,\hfill & & \hfill D_8& =& (k_1k_2)^2,\hfill & & \hfill N& =& k_3^2.\hfill \end{array}$$
This integral is mirror-symmetric with respect to $`12`$, $`45`$, $`67`$. It vanishes when the indices of the following groups of lines are non-positive: 36, 37, 67, 123, 148, 168, 416, 486, 258, 278, 527, or 587.
First we are going to get rid of the numerator. When $`n_41`$, we can lower $`n_0`$ by the $`_1(k_1k_3)`$ relation
$$\left[dn_1n_4n_82n_6+n_1\mathrm{𝟏}^+\mathrm{𝟑}^{}+n_4\mathrm{𝟒}^+(\mathrm{𝟎}^+\mathrm{𝟔}^{})+n_8\mathrm{𝟖}^+(\mathrm{𝟕}^{}\mathrm{𝟔}^{})\right]I_b=0;$$
(49)
the case $`n_50`$ is symmetric. When $`n_61`$, we can lower $`n_0`$ by the $`_1k_1`$ relation
$$\left[dn_1n_6n_82n_4+n_1\mathrm{𝟏}^++n_6\mathrm{𝟔}^+(\mathrm{𝟎}^+\mathrm{𝟒}^{})+n_8\mathrm{𝟖}^+(\mathrm{𝟓}^{}\mathrm{𝟒}^{})\right]I_b=0;$$
(50)
the case $`n_71`$ is symmetric. When $`n_31`$, we can lower $`n_0`$ or raise $`n_6`$ or $`n_7`$ by the $`_3v`$ relation
$$\left[2n_3\mathrm{𝟑}^++n_0\mathrm{𝟎}^{}(1\mathrm{𝟑}^{})+n_6\mathrm{𝟔}^+(\mathrm{𝟑}^{}\mathrm{𝟏}^{})+n_7\mathrm{𝟕}^+(\mathrm{𝟑}^{}\mathrm{𝟐}^{})\right]I_b=0.$$
(51)
Finally, we can lower $`n_0`$ or raise $`n_3`$ or $`n_7`$ by the $`_3(k_3k_1)`$ relation
$$\left[d+n_0n_3n_72n_6+n_3\mathrm{𝟑}^+\mathrm{𝟏}^{}+n_0\mathrm{𝟎}^{}(\mathrm{𝟔}^{}\mathrm{𝟒}^{})+n_7\mathrm{𝟕}^+(\mathrm{𝟖}^{}\mathrm{𝟔}^{})\right]I_b=0.$$
(52)
Now we shall discuss the integral $`I_b(n_1,n_2,n_3,n_4,n_5,n_6,n_7,n_8)`$ without numerator ($`n_0=0`$). Applying $`_1v`$, $`_1(k_1k_2)`$, $`(_1+_2+_3)k_1`$ to the integrand of (3.2), we obtain the recurrence relations
$`\left[2n_1\mathrm{𝟏}^++n_4\mathrm{𝟒}^+(\mathrm{𝟏}^{}1)+n_6\mathrm{𝟔}^+(\mathrm{𝟏}^{}\mathrm{𝟑}^{})+n_8\mathrm{𝟖}^+(\mathrm{𝟏}^{}\mathrm{𝟐}^{})\right]I_b`$ $`=`$ $`0,`$ (53)
$`[dn_1n_4n_62n_8+n_1\mathrm{𝟏}^+\mathrm{𝟐}^{}+`$
$`+n_4\mathrm{𝟒}^+(\mathrm{𝟓}^{}\mathrm{𝟖}^{})+n_6\mathrm{𝟔}^+(\mathrm{𝟕}^{}\mathrm{𝟖}^{})]I_b`$ $`=`$ $`0,`$ (54)
$`[2(dn_6n_7n_8)n_2n_3n_5+`$
$`+(n_2\mathrm{𝟐}^++n_3\mathrm{𝟑}^+)\mathrm{𝟏}^{}+n_5\mathrm{𝟓}^+(\mathrm{𝟒}^{}\mathrm{𝟖}^{})]I_b`$ $`=`$ $`0,`$ (55)
where the last relation was simplified using the homogeneity relation
$`[3dn_1n_2n_32(n_4+n_5+n_6+n_7+n_8)+`$
$`+n_1\mathrm{𝟏}^++n_2\mathrm{𝟐}^++n_3\mathrm{𝟑}^+]I_b`$ $`=`$ $`0.`$ (56)
Using it to simplify the sum of (53), its symmetric and (51), we get
$`[2[3dn_1n_2n_32(n_4+n_5+n_6+n_7+n_8)]+`$
$`+n_4\mathrm{𝟒}^+(\mathrm{𝟏}^{}1)+n_5\mathrm{𝟓}^+(\mathrm{𝟐}^{}1)]I_b`$ $`=`$ $`0.`$ (57)
The cases $`n_3=0`$, $`n_4=0`$, $`n_5=0`$, $`n_6=0`$, $`n_7=0`$ are trivial. When $`n_3<0`$ and $`n_61`$, we can raise $`n_3`$ by (53) (the case $`n_71`$ is symmetric); when $`n_3<0`$ and $`n_6=n_7=1`$, we can raise $`n_3`$ or $`n_7`$ by (52). When $`n_4<0`$ and $`n_51`$, we can raise $`n_4`$ by the relation symmetric to (54); when $`n_4<0`$ and $`n_5=1`$, we can raise $`n_4`$ or $`n_5`$ by (57). The case $`n_5<0`$ is symmetric. When $`n_6<0`$ and $`n_71`$, we can raise $`n_6`$ by (52); when $`n_6<0`$ and $`n_31`$, $`n_7=1`$, we can raise $`n_6`$ or $`n_7`$ by (51); when $`n_6<0`$ and $`n_3=n_7=1`$, we can raise $`n_6`$ or $`n_3`$ by the relation symmetric to (52). The case $`n_7<0`$ is symmetric. The case $`n_1=0`$ is $`J_a(n_3,n_2,n_5,n_7,n_4,n_6,n_8)`$ (figure 12*a*, section 3.1), $`n_2=0`$ is symmetric. When $`n_1<0`$ and $`n_81`$, we can raise $`n_1`$ by the relation symmetric to (53); when $`n_1<0`$ and $`n_61`$, we can raise $`n_1`$ by (51); when $`n_1<0`$ and $`n_41`$, we can raise $`n_1`$ by (57); when $`n_1<0`$ and $`n_21`$, we can raise $`n_1`$ by the relation symmetric to (54); when $`n_1<0`$ and $`n_31`$, we can raise $`n_1`$ by (52); when $`n_1<0`$ and $`n_2=n_3=n_4=n_6=n_8=1`$, we can raise $`n_1`$, $`n_2`$ or $`n_3`$ by (56). The case $`n_2<0`$ is symmetric. The case $`n_8=0`$ (figure 12*c*, $`J_c(n_1,n_2,n_3,n_4,n_5,n_6,n_7)`$) will be considered later in this section. When $`n_8<0`$ and $`n_71`$, we can raise $`n_8`$ by (52) (the case $`n_61`$ is symmetric); when $`n_8<0`$ and $`n_51`$, we can raise $`n_8`$ by (55) (the case $`n_41`$ is symmetric); when $`n_8<0`$ and $`n_11`$, we can raise $`n_8`$, $`n_4`$ or $`n_6`$ by (53) (the case $`n_21`$ is symmetric); when $`n_8<0`$ and $`n_31`$, we can raise $`n_6`$ or $`n_7`$ by (51); when $`n_8<0`$ and $`n_1=n_2=n_3=n_4=n_5=n_6=n_7=1`$, we can raise $`n_1`$, $`n_2`$ or $`n_3`$ by (56). When all the indices are positive, we can kill one of the lines 1, 6, 8 by (52), or one of the lines 2, 7, 8 by its mirror-symmetric relation.
Now we consider $`J_c(n_1,n_2,n_3,n_4,n_5,n_6,n_7)=I_b(n_1,n_2,n_3,n_4,n_5,n_6,n_7,n_8)`$ (figure 12*c*). It is mirror-symmetric; it vanishes if any two indices of $`n_4`$, $`n_5`$, $`n_6`$, $`n_7`$ are non-positive, or $`n_1`$, $`n_2`$, $`n_3`$ are all non-positive. If any index is zero, the integral becomes trivial. If any index is negative, we just consider $`J_c`$ as $`I_b`$ with $`n_8=0`$ and proceed as usual. Using the $`_1v`$ and $`_3v`$ relations
$`\left[2n_1\mathrm{𝟏}^++n_4\mathrm{𝟒}^+(\mathrm{𝟏}^{}1)+n_6\mathrm{𝟔}^+(\mathrm{𝟏}^{}\mathrm{𝟑}^{})\right]J_c`$ $`=`$ $`0,`$ (58)
$`\left[2n_3\mathrm{𝟑}^++n_6\mathrm{𝟔}^+(\mathrm{𝟑}^{}\mathrm{𝟏}^{})+n_7\mathrm{𝟕}^+(\mathrm{𝟑}^{}\mathrm{𝟐}^{})\right]J_c`$ $`=`$ $`0,`$ (59)
we can lower $`n_1`$, $`n_2`$ (symmetric case) and $`n_3`$ down to 1.
We are left with $`J_c(1,1,1,n_4,n_5,n_6,n_7)`$. It is rather difficult to apply the standard techniques to this integral, because each operator $`_ik_j`$ produces a scalar product which cannot be expressed via the denominators. For example, $`_1k_1`$ gives the term $`\frac{n_6}{D_6}2k_1k_3`$, and $`_1(k_1k_3)`$ gives the term $`\frac{n_4}{D_4}2k_1k_3`$. We can cancel $`2k_1k_3`$ by forming the difference
$$(n_41)\mathrm{𝟔}^{}_1k_1(n_61)\mathrm{𝟒}^{}_1(k_1k_3),$$
which results in
$$\left[(n_41)\left(dn_12n_4+n_1\mathrm{𝟏}^+\right)\mathrm{𝟔}^{}(n_61)\left(dn_12n_6+n_1\mathrm{𝟏}^+\mathrm{𝟑}^{}\right)\mathrm{𝟒}^{}\right]J_c=0.$$
(60)
Another useful combination is
$$(n_41)\mathrm{𝟓}^{}\left(_1+_2+_3\right)k_1(n_51)\mathrm{𝟒}^{}\left(_1+_2+_3\right)k_2$$
(it can be simplified by the homogeneity relation). Using (58) and (59), we obtain at $`n_1=n_2=n_3=1`$
$`[(n_41)[2(d2n_41)n_4\mathrm{𝟒}^+(1\mathrm{𝟏}^{})]\mathrm{𝟔}^{}`$
$`2(n_61)(d2n_61+\mathrm{𝟏}^+\mathrm{𝟑}^{})\mathrm{𝟒}^{}+(n_41)(n_61)(\mathrm{𝟏}^{}\mathrm{𝟑}^{})]J_c`$ $`=`$ $`0,`$ (61)
$`[(n_41)[2(dn_5n_6n_7)+(\mathrm{𝟐}^++\mathrm{𝟑}^+)\mathrm{𝟏}^{}]\mathrm{𝟓}^{}`$
$`(n_51)[2(dn_4n_6n_7)+(\mathrm{𝟏}^++\mathrm{𝟑}^+)\mathrm{𝟐}^{}]\mathrm{𝟒}^{}]J_c`$ $`=`$ $`0.`$ (62)
One more useful relation is obtained by adding (58), its mirror-symmetric, formula (59) and using the homogeneity relation:
$`[2[3dn_1n_2n_32(n_4+n_5+n_6+n_7)]+`$
$`+n_4\mathrm{𝟒}^+(\mathrm{𝟏}^{}1)+n_5\mathrm{𝟓}^+(\mathrm{𝟐}^{}1)]J_c`$ $`=`$ $`0.`$ (63)
When $`n_6>1`$, we can lower it down to 1 (raising $`n_4`$) by (61); the case $`n_7>1`$ is symmetric. Finally, we can lower $`n_4`$ and $`n_5`$ down to 1 by (62) together with (63).
The integral $`J_c(1,1,1,1,1,1,1)`$ cannot be reduced to simpler ones (in contrast to the massless case ), and should be considered as a new basis integral. Its value is currently unknown, and its calculation is highly non-trivial.
### 3.3 Ladder with three heavy-quark lines
Now let’s consider the diagram of figure 11*c*. We define
$$\begin{array}{ccccccccccc}\multicolumn{5}{c}{\frac{N^{n_0}d^dk_1d^dk_2d^dk_3}{D_1^{n_1}D_2^{n_2}D_3^{n_3}D_4^{n_4}D_5^{n_5}D_6^{n_6}D_7^{n_7}D_8^{n_8}}=}& & & & & & \\ & =& \multicolumn{9}{c}{i\pi ^{3d/2}(2\omega )^{3d+2n_02_{i=4}^8n_i}I_c(n_1,n_2,n_3,n_4,n_5,n_6,n_7,n_8;n_0),}\\ \hfill D_1& =& \frac{(k_1+p)v}{\omega },\hfill & & \hfill D_2& =& \frac{(k_2+p)v}{\omega },\hfill & & \hfill D_3& =& \frac{(k_3+p)v}{\omega },\hfill \\ \hfill D_4& =& k_1^2,\hfill & & \hfill D_5& =& k_2^2,\hfill & & \hfill D_6& =& (k_1k_3)^2,\hfill \\ \hfill D_7& =& (k_2k_3)^2,\hfill & & \hfill D_8& =& k_3^2,\hfill & & \hfill N& =& 2k_1k_2.\hfill \end{array}$$
This integral is mirror-symmetric with respect to $`12`$, $`45`$, $`67`$. It vanishes when the indices of the following groups of lines are non-positive: 14, 16, 46, 25, 27, 57, 132, 637, 368, 378, 687, 485, 487, 586.
First we are going to get rid of the numerator. When $`n_71`$, we can lower $`n_0`$ by the $`_3(k_3k_1)`$ relation
$`[dn_3n_7n_82n_6+n_3\mathrm{𝟑}^+\mathrm{𝟏}^{}+`$
$`+n_7\mathrm{𝟕}^+(\mathrm{𝟎}^++\mathrm{𝟒}^{}+\mathrm{𝟓}^{}\mathrm{𝟔}^{})+n_8\mathrm{𝟖}^+(\mathrm{𝟒}^{}\mathrm{𝟔}^{})]I_c`$ $`=`$ $`0;`$ (64)
the case $`n_60`$ is symmetric. When $`n_51`$, we can lower $`n_0`$ by
$`[2(dn_5n_6n_7)+n_0n_2n_3n_8+2n_0\mathrm{𝟎}^{}\mathrm{𝟒}^{}+`$
$`+(n_2\mathrm{𝟐}^++n_3\mathrm{𝟑}^+)\mathrm{𝟏}^{}n_5\mathrm{𝟓}^+\mathrm{𝟎}^++n_8\mathrm{𝟖}^+(\mathrm{𝟒}^{}\mathrm{𝟔}^{})]I_c`$ $`=`$ $`0,`$ (65)
which is the $`(_1+_2+_3)k_1`$ relation simplified by the homogeneity relation. the case $`n_41`$ is symmetric. When $`n_11`$, we can lower $`n_0`$ or raise $`n_4`$ or $`n_6`$ by the $`_1v`$ relation
$$\left[2n_1\mathrm{𝟏}^++n_0\mathrm{𝟎}^{}(\mathrm{𝟐}^{}1)+n_4\mathrm{𝟒}^+(\mathrm{𝟏}^{}1)+n_6\mathrm{𝟔}^+(\mathrm{𝟏}^{}\mathrm{𝟑}^{})\right]I_c=0;$$
(66)
the case $`n_21`$ is symmetric. Finally, we can raise $`n_1`$ or $`n_6`$ by the $`_1k_1`$ relation
$$\left[d+n_0n_1n_62n_4+n_1\mathrm{𝟏}^++n_6\mathrm{𝟔}^+(\mathrm{𝟖}^{}\mathrm{𝟒}^{})\right]I_c=0.$$
(67)
Now we shall discuss the integral $`I_c(n_1,n_2,n_3,n_4,n_5,n_6,n_7,n_8)`$ without numerator ($`n_0=0`$). Applying $`_3v`$, $`_1(k_1k_3)`$ to the integrand of (3.3), we obtain the recurrence relations
$`\left[2n_3\mathrm{𝟑}^++n_6\mathrm{𝟔}^+(\mathrm{𝟑}^{}\mathrm{𝟏}^{})+n_7\mathrm{𝟕}^+(\mathrm{𝟑}^{}\mathrm{𝟐}^{})+n_8\mathrm{𝟖}^+(\mathrm{𝟑}^{}1)\right]I_c`$ $`=`$ $`0,`$ (68)
$`\left[dn_1n_42n_6+n_1\mathrm{𝟏}^+\mathrm{𝟑}^{}+n_4\mathrm{𝟒}^+(\mathrm{𝟖}^{}\mathrm{𝟔}^{})\right]I_c`$ $`=`$ $`0.`$ (69)
Homogeneity in $`\omega `$ gives the relation identical with (56).
The cases $`n_1=0`$, $`n_2=0`$, $`n_4=0`$, $`n_5=0`$, $`n_6=0`$, $`n_7=0`$ are trivial. When $`n_1<0`$ and $`n_61`$, we can raise $`n_1`$ by (68); when $`n_1<0`$ and $`n_41`$, we can raise $`n_1`$ using the sum of (66) and (68); when $`n_1<0`$ and $`n_4=n_6=1`$, we can raise $`n_1`$ or $`n_6`$ by (67). The case $`n_2<0`$ is symmetric. When $`n_4<0`$ and $`n_61`$, we can raise $`n_4`$ by (67); when $`n_4<0`$ and $`n_11`$, $`n_6=1`$, we can raise $`n_4`$ or $`n_6`$ by (66); when $`n_4<0`$ and $`n_1=n_6=1`$, we can raise $`n_4`$ or $`n_1`$ by (69). The case $`n_5<0`$ is symmetric. When $`n_6<0`$ and $`n_41`$, we can raise $`n_6`$ by (69); when $`n_6<0`$ and $`n_11`$, $`n_4=1`$, we can raise $`n_6`$ or $`n_4`$ by (66); when $`n_6<0`$ and $`n_1=n_4=1`$, we can raise $`n_6`$ or $`n_1`$ by (67). The case $`n_7<0`$ is symmetric. The case $`n_3=0`$ is $`J_b(n_1,n_2,n_4,n_5,n_8,n_6,n_7)`$ (figure 12*a*, section 3.1). When $`n_3<0`$ and $`n_61`$, we can raise $`n_3`$ by (66) (the case $`n_71`$ is symmetric); when $`n_3<0`$ and $`n_81`$, we can raise $`n_3`$ by the relation
$`[2[3dn_1n_2n_32(n_4+n_5+n_6+n_7+n_8)]+`$
$`+n_4\mathrm{𝟒}^+(\mathrm{𝟏}^{}1)+n_5\mathrm{𝟓}^+(\mathrm{𝟐}^{}1)+n_8\mathrm{𝟖}^+(\mathrm{𝟑}^{}1)]I_c`$ $`=`$ $`0,`$ (70)
which is (66) plus its mirror-symmetric plus (68) simplified by the homogeneity relation; when $`n_3<0`$ and $`n_11`$, we can raise $`n_3`$ by (69) (the case $`n_21`$ is symmetric); when $`n_3<0`$ and $`n_1=n_2=n_6=n_8=1`$, we can raise $`n_1`$, $`n_2`$ or $`n_3`$ by the homogeneity relation. The case $`n_8=0`$ is $`J_c(n_1,n_2,n_3,n_4,n_5,n_6,n_7)`$ (figure 12*c*, section 3.2). When $`n_8<0`$ and $`n_61`$, we can raise $`n_8`$ by (67) (the case $`n_71`$ is symmetric); when $`n_8<0`$ and $`n_41`$, we can raise $`n_8`$ by (69) (the case $`n_51`$ is symmetric); when $`n_8<0`$ and $`n_31`$, $`n_4=n_6=n_7=1`$, we can raise $`n_8`$, $`n_6`$ or $`n_7`$ by (68); when $`n_8<0`$ and $`n_11`$, $`n_3=n_4=n_5=n_6=n_7=1`$, we can raise $`n_4`$ or $`n_6`$ by (66) (the case $`n_21`$ is symmetric); when $`n_8<0`$ and $`n_1=n_2=n_3=n_4=n_5=n_6=n_7=1`$, we can raise $`n_1`$, $`n_2`$ or $`n_3`$ the homogeneity relation. When all the indices are positive, we can kill one of the lines 3, 6, 8 by (69), or one of the lines 3, 7, 8 by its mirror-symmetric relation.
### 3.4 Non-planar diagram with three heavy-quark lines
Now let’s consider the diagram of figure 11*d*. We define
$$\begin{array}{ccccccccccc}\multicolumn{7}{c}{\frac{N^{n_0}d^dk_1d^dk_2d^dk_3}{D_1^{n_1}D_2^{n_2}D_3^{n_3}D_4^{n_4}D_5^{n_5}D_6^{n_6}D_7^{n_7}D_8^{n_8}}=}& & & & \\ & =& \multicolumn{9}{c}{i\pi ^{3d/2}(2\omega )^{3d+2n_02_{i=4}^8n_i}I_d(n_1,n_2,n_3,n_4,n_5,n_6,n_7,n_8;n_0),}\\ \hfill D_1& =& \frac{(k_1+p)v}{\omega },\hfill & & \hfill D_2& =& \frac{(k_2+p)v}{\omega },\hfill & & \hfill D_3& =& \frac{(k_1+k_2k_3+p)v}{\omega },\hfill \\ \hfill D_4& =& k_1^2,\hfill & & \hfill D_5& =& k_2^2,\hfill & & \hfill D_6& =& (k_1k_3)^2,\hfill \\ \hfill D_7& =& (k_2k_3)^2,\hfill & & \hfill D_8& =& k_3^2,\hfill & & \hfill N& =& 2k_1k_2.\hfill \end{array}$$
This integral is mirror-symmetric with respect to $`12`$, $`45`$, $`67`$. It vanishes when the indices of the following groups of lines are non-positive: 46, 57, 132, 485, 487, 586, 687, 314, 136, 325, 237, 148, 178, 258, 268, 417, 526, 637, 368, 378.
When $`n_10`$, $`n_20`$ or $`n_30`$, we use
$`I_d(|n_1|,n_2,n_3,n_4,n_5,n_6,n_7,n_8;n_0)=`$
$`=(\mathrm{𝟎}^++\mathrm{𝟏}^{})^{|n_1|}(\mathrm{𝟓}^{}+\mathrm{𝟔}^{}\mathrm{𝟑}^{}\mathrm{𝟖}^{})^{n_0}I_a(n_3,n_2,0,n_5,n_6,n_7,n_4,n_8),`$ (71)
$`I_d(n_1,n_2,|n_3|,n_4,n_5,n_6,n_7,n_8;n_0)=`$
$`=(\mathrm{𝟎}^++\mathrm{𝟏}^{}+\mathrm{𝟐}^{}1)^{|n_3|}(\mathrm{𝟓}^{}\mathrm{𝟑}^{}\mathrm{𝟒}^{})^{n_0}I_a(n_1,n_2,n_4,n_5,0,n_8,n_6,n_7),`$ (72)
and relation symmetric to (71). Applying $`_2v`$ and $`_3v`$ to the integrand of (3.4), we obtain
$`\left[2n_2\mathrm{𝟐}^++2n_3\mathrm{𝟑}^++n_5\mathrm{𝟓}^+(1\mathrm{𝟐}^{})+n_7\mathrm{𝟕}^+(\mathrm{𝟏}^{}\mathrm{𝟑}^{})\right]I_d=0,`$ (73)
$`[2n_3\mathrm{𝟑}^++n_6\mathrm{𝟔}^+(\mathrm{𝟐}^{}\mathrm{𝟑}^{})+n_7\mathrm{𝟕}^+(\mathrm{𝟏}^{}\mathrm{𝟑}^{})+`$
$`+n_8\mathrm{𝟖}^+(\mathrm{𝟏}^{}+\mathrm{𝟐}^{}\mathrm{𝟑}^{}1)]I_d`$ $`=`$ $`0.`$ (74)
When $`n_1>1`$, we can lower it by
$`[2[3dn_1n_2n_32(n_4+n_5+n_6+n_7+n_8n_0)]+`$
$`+n_0\mathrm{𝟎}^{}(\mathrm{𝟏}^{}1)+2n_1\mathrm{𝟏}^++n_5\mathrm{𝟓}^+(\mathrm{𝟐}^{}1)+n_7\mathrm{𝟕}^+(\mathrm{𝟑}^{}\mathrm{𝟏}^{})]I_d`$ $`=`$ $`0,`$ (75)
which is twice the homogeneity relation
$`[3dn_1n_2n_32(n_4+n_5+n_6+n_7+n_8n_0)+`$
$`+n_1\mathrm{𝟏}^++n_2\mathrm{𝟐}^++n_3\mathrm{𝟑}^+]I_d`$ $`=`$ $`0`$ (76)
minus (73); the case $`n_2>1`$ is symmetric. When $`n_3>1`$, we can lower it by (74).
We are left with $`I_d(1,1,1,n_4,n_5,n_6,n_7,n_8;n_0)`$, and now are going to get rid of the numerator. Applying $`_3(k_3k_2)`$ and $`_1(k_1k_3)`$ to the integrand of (3.4), we obtain
$`[dn_3n_6n_82n_7+n_3\mathrm{𝟑}^+\mathrm{𝟏}^{}+`$
$`+n_6\mathrm{𝟔}^+(\mathrm{𝟒}^{}+\mathrm{𝟓}^{}\mathrm{𝟕}^{}+\mathrm{𝟎}^+)+n_8\mathrm{𝟖}^+(\mathrm{𝟓}^{}\mathrm{𝟕}^{})]I_d`$ $`=`$ $`0,`$ (77)
$`[d+n_0n_3n_42n_6+n_0\mathrm{𝟎}^{}(\mathrm{𝟓}^{}\mathrm{𝟕}^{}+\mathrm{𝟖}^{})+`$
$`+n_1\mathrm{𝟏}^+(\mathrm{𝟐}^{}\mathrm{𝟑}^{})+n_3\mathrm{𝟑}^+\mathrm{𝟐}^{}+n_4\mathrm{𝟒}^+(\mathrm{𝟖}^{}\mathrm{𝟔}^{})]I_d`$ $`=`$ $`0.`$ (78)
When $`n_41`$, we can lower $`n_0`$ by
$`[2(dn_4n_6n_7)+n_0n_1n_2n_3n_8+1`$
$`\left[3dn_1n_2n_3+12(n_4+n_5+n_6+n_7+n_8n_0)\right]\mathrm{𝟐}^{}+`$
$`+2n_0\mathrm{𝟎}^{}\mathrm{𝟓}^{}n_4\mathrm{𝟒}^+\mathrm{𝟎}^++n_8\mathrm{𝟖}^+(\mathrm{𝟓}^{}\mathrm{𝟕}^{})]I_d`$ $`=`$ $`0,`$ (79)
which is the $`(_1+_2+_3)k_2`$ relation plus (76) minus $`\mathrm{𝟐}^{}`$ shifted (76); the case $`n_51`$ is symmetric. When $`n_61`$, we can lower $`n_0`$ by (77); the case $`n_71`$ is symmetric. Finally, when $`n_4=n_5=n_6=n_7=1`$, we can lower $`n_0`$ or raise $`n_4`$ by (78).
Now we shall discuss the integral $`I_d(n_1,n_2,n_3,n_4,n_5,n_6,n_7,n_8)`$ without numerator ($`n_0=0`$). Applying $`_1k_1`$, $`_3k_3`$, $`(_2+_3)k_3`$ to the integrand of (3.4), we obtain
$`\left[dn_1n_62n_4+n_1\mathrm{𝟏}^++n_3\mathrm{𝟑}^+(1\mathrm{𝟏}^{})+n_6\mathrm{𝟔}^+(\mathrm{𝟖}^{}\mathrm{𝟒}^{})\right]I_d`$ $`=`$ $`0,`$ (80)
$`[dn_3n_6n_72n_8+n_3\mathrm{𝟑}^+(\mathrm{𝟏}^{}+\mathrm{𝟐}^{}1)+`$
$`+n_6\mathrm{𝟔}^+(\mathrm{𝟒}^{}\mathrm{𝟖}^{})+n_7\mathrm{𝟕}^+(\mathrm{𝟓}^{}\mathrm{𝟖}^{})]I_d`$ $`=`$ $`0,`$ (81)
$`[dn_2n_5n_62n_8+n_2\mathrm{𝟐}^+(1+\mathrm{𝟑}^{}\mathrm{𝟏}^{})+`$
$`+n_5\mathrm{𝟓}^+(\mathrm{𝟕}^{}\mathrm{𝟖}^{})+n_6\mathrm{𝟔}^+(\mathrm{𝟒}^{}\mathrm{𝟖}^{})]I_d`$ $`=`$ $`0.`$ (82)
Some other useful relations are
$$\left[2n_2\mathrm{𝟐}^++n_5\mathrm{𝟓}^+(1\mathrm{𝟐}^{})+n_6\mathrm{𝟔}^+(\mathrm{𝟑}^{}\mathrm{𝟐}^{})+n_8\mathrm{𝟖}^+(1+\mathrm{𝟑}^{}\mathrm{𝟏}^{}\mathrm{𝟐}^{})\right]I_d=0,$$
(83)
$$\left[2(dn_5n_7n_8)n_2n_3n_6+n_2\mathrm{𝟐}^++n_3\mathrm{𝟑}^+\mathrm{𝟏}^{}+n_6\mathrm{𝟔}^+(\mathrm{𝟒}^{}\mathrm{𝟖}^{})\right]I_d=0,$$
(84)
the differences of (73) and (74), and of (76) and (80).
The case $`n_1=0`$ is $`J_a(n_3,n_2,n_5,n_6,n_7,n_4,n_8)`$ (figure 12*a*, section 3.1; $`n_2=0`$ is symmetric); the case $`n_3=0`$ is $`J_b(n_1,n_2,n_4,n_5,n_8,n_6,n_7)`$ (figure 12*b*, section 3.1); the case $`n_8=0`$ is $`J_c(n_1,n_2,n_3,n_4,n_5,n_7,n_6)`$ (figure 12*c*, section 3.2). The cases $`n_6=0`$ ($`J_d(n_1,n_3,n_2,n_4,n_7,n_5,n_8)`$, figure 12*d*; $`n_7=0`$ is symmetric) and $`n_4=0`$ ($`J_e(n_1,n_3,n_2,n_6,n_8,n_5,n_7)`$, figure 12*e*; $`n_5=0`$ is symmetric) will be discussed later in this section. When $`n_2<0`$, it can be raised using
$`\left[3dn_1n_2n_32(n_4+n_5+n_6+n_7+n_8)\right]I_d=`$
$`=\left[dn_2n_3n_42n_6n_1\mathrm{𝟏}^+\mathrm{𝟑}^{}+n_4\mathrm{𝟒}^+(\mathrm{𝟖}^{}\mathrm{𝟔}^{})\right]\mathrm{𝟐}^+I_d,`$ (85)
which is (76) minus $`\mathrm{𝟐}^+`$ shifted (78); the case $`n_1<0`$ is symmetric. When $`n_3<0`$ and $`n_11`$, we can raise $`n_3`$ by (78) ($`n_21`$ is symmetric); when $`n_3<0`$ and $`n_71`$, we can raise $`n_3`$ by (73) ($`n_61`$ is symmetric); when $`n_3<0`$ and $`n_81`$, we can raise $`n_3`$ using
$`[2[3dn_1n_2n_32(n_4+n_5+n_6+n_7+n_8)]+`$
$`+n_4\mathrm{𝟒}^+(\mathrm{𝟏}^{}1)+n_5\mathrm{𝟓}^+(\mathrm{𝟐}^{}1)+n_8\mathrm{𝟖}^+(\mathrm{𝟏}^{}+\mathrm{𝟐}^{}\mathrm{𝟑}^{}1)]I_d`$ $`=`$ $`0,`$ (86)
which is (73) plus its symmetric minus (74) simplified by (76); when $`n_3<0`$ and $`n_1=n_2=n_6=n_7=n_8=1`$, we can raise $`n_3`$, $`n_1`$ or $`n_2`$ by (76). When $`n_1>1`$, it can be lowered by (75); the case $`n_2>1`$ is symmetric. When $`n_3>1`$, it can be lowered by (74).
We are left with $`I_d(1,1,1,n_4,n_5,n_6,n_7,n_8)`$. When $`n_4<0`$, it can be raised by (78); the case $`n_5<0`$ is symmetric. When $`n_6<0`$ and $`n_41`$, we can raise $`n_6`$ by (78); when $`n_6<0`$ and $`n_4=1`$, we can raise $`n_6`$ or $`n_4`$ using
$`[2(dn_1n_62n_4)2n_3\mathrm{𝟑}^+\mathrm{𝟏}^{}+n_4\mathrm{𝟒}^+(\mathrm{𝟏}^{}1)+`$
$`+n_6\mathrm{𝟔}^+[2(\mathrm{𝟖}^{}\mathrm{𝟒}^{})+\mathrm{𝟑}^{}\mathrm{𝟐}^{}]]I_d`$ $`=`$ $`0,`$ (87)
which is twice (80) minus the relation symmetric to (73). The case $`n_7<0`$ is symmetric. When $`n_8<0`$ and $`n_41`$, we can raise $`n_8`$ by (78) ($`n_51`$ is symmetric); when $`n_8<0`$ and $`n_61`$, we can raise $`n_8`$ by (87) ($`n_71`$ is symmetric); when $`n_8<0`$ and $`n_1=n_2=n_3=n_4=n_5=n_6=n_7=1`$, we can raise $`n_8`$, $`n_4`$ or $`n_5`$ by (86). When all the indices are positive, we can kill the line 6 or 8 using (78).
Thus, the non-planar diagram reduces to planar ones, in contrast to the massless case where such a reduction is impossible .
Now we consider $`J_d(n_1,n_2,n_3,n_4,n_5,n_7,n_8)=I_d(n_1,n_2,n_3,n_4,n_5,0,n_7,n_8)`$ (figure 12*d*). This integral vanishes when indices of the following groups of lines are non-positive: 4, 13, 23, 57, 58, 78, 37, 38, 25, 28. It becomes trivial if any of the indices is zero. Applying $`(_2+_3)k_2`$ and $`_3(k_3k_2)`$ to the integrand of (3.4), we obtain
$`\left[dn_2n_82n_5+n_2\mathrm{𝟐}^++n_8\mathrm{𝟖}^+(\mathrm{𝟕}^{}\mathrm{𝟓}^{})\right]J_d`$ $`=`$ $`0,`$ (88)
$`\left[dn_3n_82n_7+n_3\mathrm{𝟑}^+\mathrm{𝟏}^{}+n_8\mathrm{𝟖}^+(\mathrm{𝟓}^{}\mathrm{𝟕}^{})\right]J_d`$ $`=`$ $`0.`$ (89)
When $`n_3<0`$ and $`n_21`$, we can raise $`n_3`$ by (82); when $`n_3<0`$ and $`n_71`$, we can raise $`n_3`$ by (75); when $`n_3<0`$ and $`n_81`$, we can raise $`n_3`$ by (83); when $`n_3<0`$ and $`n_2=n_7=n_8=1`$, we can raise $`n_3`$ or $`n_2`$ by (84). When $`n_2<0`$ and $`n_31`$, we can raise $`n_2`$ by (81); when $`n_2<0`$ and $`n_51`$, we can raise $`n_2`$ by (75); when $`n_2<0`$ and $`n_81`$, we can raise $`n_2`$ by (74); when $`n_2<0`$ and $`n_3=n_5=n_8=1`$, we can raise $`n_2`$ or $`n_3`$ by (84). When $`n_8<0`$ and $`n_71`$, we can raise $`n_8`$ by (81); when $`n_8<0`$ and $`n_51`$, we can raise $`n_8`$ by (82); when $`n_8<0`$ and $`n_31`$, $`n_5=n_7=1`$, we can raise $`n_8`$ or $`n_7`$ by (74); when $`n_8<0`$ and $`n_21`$, $`n_3=n_5=n_7=1`$, we can raise $`n_8`$ or $`n_5`$ by (83); when $`n_8<0`$ and $`n_2=n_3=n_5=n_7=1`$, we can raise $`n_8`$ or $`n_2`$ by (88). When $`n_1<0`$ and $`n_31`$, we can raise $`n_1`$ by (84); when $`n_1<0`$ and $`n_41`$, we can raise $`n_1`$ using the relation symmetric to (73); when $`n_1<0`$ and $`n_21`$, we can raise $`n_1`$ by (82); when $`n_1<0`$ and $`n_71`$, we can raise $`n_1`$ by (75); when $`n_1<0`$ and $`n_81`$, we can raise $`n_1`$ by (83); when $`n_1<0`$ and $`n_2=n_3=n_4=n_7=n_8=1`$, we can raise $`n_1`$, $`n_2`$ or $`n_3`$ by (76). When $`n_5<0`$ and $`n_81`$, we can raise $`n_5`$ by (89); when $`n_5<0`$ and $`n_71`$, we can raise $`n_5`$ by (81); when $`n_5<0`$ and $`n_21`$, we can raise $`n_5`$ or $`n_8`$ by (83); when $`n_5<0`$ and $`n_2=n_7=n_8=1`$, we can raise $`n_5`$ or $`n_2`$ by (82). When $`n_7<0`$ and $`n_81`$, we can raise $`n_7`$ by (89); when $`n_7<0`$ and $`n_51`$, we can raise $`n_7`$ by (82); when $`n_7<0`$ and $`n_31`$, $`n_5=n_8=1`$, we can raise $`n_7`$ or $`n_8`$ by (74); when $`n_7<0`$ and $`n_1=n_3=n_5=1`$, we can raise $`n_7`$ or $`n_3`$ by (81). When all the indices are positive, we can kill one of the lines 1, 5, 7 by (89).
Finally, we consider $`J_e(n_1,n_2,n_3,n_5,n_6,n_7,n_8)=I_d(n_1,n_2,n_3,0,n_5,n_6,n_7,n_8)`$ (figure 12*e*). This integral vanishes when indices of the following groups of lines are non-positive: 6, 13, 18, 17, 58, 78, 57, 325, 237. It becomes trivial if any of the indices is zero. Applying $`(_1+_3)(k_3k_2)`$ to the integrand of (3.4), we obtain
$$\left[dn_1n_82n_7+n_1\mathrm{𝟏}^+\mathrm{𝟑}^{}+n_8\mathrm{𝟖}^+(\mathrm{𝟓}^{}\mathrm{𝟕}^{})\right]J_e=0.$$
(90)
When $`n_1<0`$, it can be raised using the relation symmetric to (85). When $`n_2<0`$, it can be raised by (85). When $`n_3<0`$ and $`n_11`$, we can raise $`n_3`$ by (78); when $`n_3<0`$ and $`n_61`$, we can raise $`n_3`$ using the relation symmetric to (73); when $`n_3<0`$ and $`n_21`$, we can raise $`n_3`$ using the relation symmetric to (78); when $`n_3<0`$ and $`n_71`$, we can raise $`n_3`$ by (73); when $`n_3<0`$ and $`n_81`$, we can raise $`n_3`$ by (86); when $`n_3<0`$ and $`n_1=n_2=n_6=n_7=n_8=1`$, we can raise $`n_3`$, $`n_1`$ or $`n_2`$ by (76). When $`n_8<0`$ and $`n_71`$, we can raise $`n_8`$ using the relation symmetric to (82); when $`n_8<0`$ and $`n_51`$, we can raise $`n_8`$ using the relation symmetric to (78); when $`n_8<0`$ and $`n_11`$, $`n_5=n_7=1`$, we can raise $`n_8`$ or $`n_7`$ using the relation symmetric to (83); when $`n_8<0`$ and $`n_1=n_5=n_7=1`$, we can raise $`n_8`$ or $`n_1`$ by (90). When $`n_5<0`$ and $`n_81`$, we can raise $`n_5`$ by (90); when $`n_5<0`$ and $`n_71`$, we can raise $`n_5`$ using the relation symmetric to (84); when $`n_5<0`$ and $`n_11`$, we can raise $`n_8`$ or $`n_7`$ using the relation symmetric to (83); when $`n_5<0`$ and $`n_1=n_7=n_8=1`$, we can raise $`n_5`$, $`n_7`$ or $`n_1`$ by (75). When $`n_7<0`$ and $`n_81`$, we can raise $`n_7`$ by (90); when $`n_7<0`$ and $`n_51`$, we can raise $`n_7`$ using the relation symmetric to (78); when $`n_7<0`$ and $`n_11`$, we can raise $`n_7`$ or $`n_8`$ using the relation symmetric to (83); when $`n_7<0`$ and $`n_1=n_5=n_8=1`$, we can raise $`n_7`$ or $`n_1`$ using the relation symmetric to (82). When all the indices are positive, we can kill one of the lines 3, 5, 7 by (90).
### 3.5 Diagrams with four heavy-quark lines
Let’s define (figure 11*e*)
$$\begin{array}{ccccccccccc}\multicolumn{5}{c}{\frac{N_{13}^{n_{13}}N_{23}^{n_{23}}d^dk_1d^dk_2d^dk_3}{D_1^{n_1}D_2^{n_2}D_3^{n_3}D_4^{n_4}D_5^{n_5}D_6^{n_6}D_7^{n_7}D_8^{n_8}}=}& & & & & & \\ & =& \multicolumn{9}{c}{i\pi ^{3d/2}(2\omega )^{3d+2(n_{13}+n_{23})2_{i=5}^8n_i}I_e(n_1,n_2,n_3,n_4,n_5,n_6,n_7,n_8;n_{13},n_{23}),}\\ \hfill D_1& =& \frac{(k_3+p)v}{\omega },\hfill & & \hfill D_2& =& \frac{(k_1+k_3+p)v}{\omega },\hfill & & \hfill D_3& =& \frac{(k_1+p)v}{\omega },\hfill \\ \hfill D_4& =& \frac{(k_2+p)v}{\omega },\hfill & & \hfill D_5& =& k_3^2,\hfill & & \hfill D_6& =& k_1^2,\hfill \\ \hfill D_7& =& k_2^2,\hfill & & \hfill D_8& =& (k_1k_2)^2,\hfill & & & & \\ \hfill N_{13}& =& 2k_1k_3,\hfill & & \hfill N_{23}& =& 2k_2k_3.\hfill & & & & \end{array}$$
This integral vanishes when the indices of the following groups of lines are non-positive: 5, 12, 67, 68, 78, 47, 48, 234, 326, 238. The heavy-quark denominators are linearly dependent: $`D_1D_2+D_3=1`$, and therefore
$$\left[1\mathrm{𝟏}^{}+\mathrm{𝟐}^{}\mathrm{𝟑}^{}\right]I_e=0.$$
(91)
The cases $`n_10`$, $`n_20`$, $`n_30`$ reduce to $`I_c`$, $`I_d`$:
$`I_e(|n_1|,n_2,n_3,n_4,n_5,n_6,n_7,n_8;n_{13},n_{23})=(1+\mathrm{𝟏}^{}\mathrm{𝟑}^{})^{|n_1|}\times `$
$`\times (\mathrm{𝟔}^{}\mathrm{𝟒}^{}+\mathrm{𝟖}^{})^{n_{13}}(\mathrm{𝟓}^{}\mathrm{𝟕}^{}+\mathrm{𝟖}^{}+\mathrm{𝟎}^+)^{n_{23}}I_c(n_2,n_4,n_3,0,n_7,n_5,n_8,n_6),`$
$`I_e(n_1,|n_2|,n_3,n_4,n_5,n_6,n_7,n_8;n_{13},n_{23})=`$
$`=(\mathrm{𝟏}^{}+\mathrm{𝟑}^{}1)^{|n_2|}(\mathrm{𝟔}^{}\mathrm{𝟒}^{}\mathrm{𝟖}^{})^{n_{13}}(\mathrm{𝟎}^+)^{n_{23}}I_c(n_1,n_4,n_3,n_5,n_7,0,n_8,n_6),`$
$`I_e(n_1,n_2,|n_3|,n_4,n_5,n_6,n_7,n_8;n_{13},n_{23})=`$
$`=(1\mathrm{𝟏}^{}+\mathrm{𝟑}^{})^{|n_3|}(\mathrm{𝟒}^{}\mathrm{𝟔}^{}+\mathrm{𝟖}^{}+\mathrm{𝟎}^+)^{n_{13}}(\mathrm{𝟎}^+)^{n_{23}}I_d(n_1,n_4,n_2,n_5,n_7,0,n_6,n_8).`$
If $`n_{1,2,3}`$ are all positive, we can lower them by (91) until one of them vanish.
Let’s define (figure 11*f*)
$$\begin{array}{ccccccccccc}\multicolumn{7}{c}{\frac{N_{13}^{n_{13}}N_{23}^{n_{23}}d^dk_1d^dk_2d^dk_3}{D_1^{n_1}D_2^{n_2}D_3^{n_3}D_4^{n_4}D_5^{n_5}D_6^{n_6}D_7^{n_7}D_8^{n_8}}=}& & & & \\ & =& \multicolumn{9}{c}{i\pi ^{3d/2}(2\omega )^{3d+2(n_{13}+n_{23})2_{i=5}^8n_i}I_f(n_1,n_2,n_3,n_4,n_5,n_6,n_7,n_8;n_{13},n_{23}),}\\ \hfill D_1& =& \frac{(k_1+p)v}{\omega },\hfill & & \hfill D_2& =& \frac{(k_1+k_3+p)v}{\omega },\hfill & & \hfill D_3& =& \frac{(k_2+k_3+p)v}{\omega },\hfill \\ \hfill D_4& =& \frac{(k_2+p)v}{\omega },\hfill & & \hfill D_5& =& k_3^2,\hfill & & \hfill D_6& =& k_1^2,\hfill \\ \hfill D_7& =& k_2^2,\hfill & & \hfill D_8& =& (k_1k_2)^2.\hfill & & & & \end{array}$$
This integral is mirror-symmetric with respect to $`14`$, $`23`$, $`67`$. It vanishes when indices of the following groups of lines are non-positive: 5, 23, 67, 68, 78, 216, 128, 347, 438. The heavy-quark denominators are linearly dependent: $`D_1D_2+D_3D_4=0`$, and therefore
$$\left[\mathrm{𝟏}^{}\mathrm{𝟐}^{}+\mathrm{𝟑}^{}\mathrm{𝟒}^{}\right]I_f=0.$$
(92)
The cases $`n_10`$, $`n_20`$ reduce to $`I_c`$, $`I_d`$:
$`I_f(|n_1|,n_2,n_3,n_4,n_5,n_6,n_7,n_8;n_{13},n_{23})=(\mathrm{𝟏}^{}+\mathrm{𝟐}^{}\mathrm{𝟑}^{})^{|n_1|}\times `$
$`\times (\mathrm{𝟔}^{}\mathrm{𝟒}^{}+\mathrm{𝟖}^{})^{n_{13}}(\mathrm{𝟓}^{}\mathrm{𝟕}^{}+\mathrm{𝟖}^{}+\mathrm{𝟎}^+)^{n_{23}}I_d(n_2,n_4,n_3,0,n_7,n_5,n_8,n_6),`$
$`I_f(n_1,|n_2|,n_3,n_4,n_5,n_6,n_7,n_8;n_{13},n_{23})=(\mathrm{𝟏}^{}+\mathrm{𝟐}^{}\mathrm{𝟑}^{})^{|n_2|}\times `$
$`\times (\mathrm{𝟒}^{}\mathrm{𝟔}^{}+\mathrm{𝟖}^{}+\mathrm{𝟎}^+)^{n_{13}}(\mathrm{𝟕}^{}+\mathrm{𝟖}^{}\mathrm{𝟓}^{})^{n_{23}}I_c(n_1,n_3,n_4,n_6,0,n_8,n_5,n_7)`$
(the cases $`n_40`$, $`n_30`$ are symmetric). If $`n_{1,2,3,4}`$ are all positive, we can use (92) to raise, say, $`n_1`$ and kill one of the lines 2, 3, 4.
Let’s define (figure 11*g*)
$$\begin{array}{ccccccccccc}\multicolumn{7}{c}{\frac{N_{13}^{n_{13}}N_{23}^{n_{23}}d^dk_1d^dk_2d^dk_3}{D_1^{n_1}D_2^{n_2}D_3^{n_3}D_4^{n_4}D_5^{n_5}D_6^{n_6}D_7^{n_7}D_8^{n_8}}=}& & & & \\ & =& \multicolumn{9}{c}{i\pi ^{3d/2}(2\omega )^{3d+2(n_{13}+n_{23})2_{i=4}^8n_i}I_g(n_1,n_2,n_3,n_4,n_5,n_6,n_7,n_8;n_{13},n_{23})}\\ \hfill D_1& =& \frac{(k_3+p)v}{\omega },\hfill & & \hfill D_2& =& \frac{(k_1+k_3+p)v}{\omega },\hfill & & \hfill D_3& =& \frac{(k_2+k_3+p)v}{\omega },\hfill \\ \hfill D_4& =& \frac{(k_2+p)v}{\omega },\hfill & & \hfill D_5& =& k_3^2,\hfill & & \hfill D_6& =& k_1^2,\hfill \\ \hfill D_7& =& k_2^2,\hfill & & \hfill D_8& =& (k_1k_2)^2.\hfill & & & & \end{array}$$
This integral vanishes when the indices of the following groups of lines are non-positive: 5, 67, 68, 78, 26, 28, 123, 234, 347, 438. The heavy-quark denominators are linearly dependent: $`D_1D_3+D_4=1`$, and therefore
$$\left[1\mathrm{𝟏}^{}+\mathrm{𝟑}^{}\mathrm{𝟒}^{}\right]I_g=0.$$
(93)
The cases $`n_10`$, $`n_30`$, $`n_40`$ reduce to $`I_b`$, $`I_d`$:
$`I_g(|n_1|,n_2,n_3,n_4,n_5,n_6,n_7,n_8;n_{13},n_{23})=(1\mathrm{𝟐}^{}+\mathrm{𝟑}^{})^{|n_1|}\times `$
$`\times (\mathrm{𝟔}^{}\mathrm{𝟒}^{}+\mathrm{𝟖}^{})^{n_{13}}(\mathrm{𝟓}^{}\mathrm{𝟕}^{}+\mathrm{𝟖}^{}+\mathrm{𝟎}^+)^{n_{23}}I_d(n_2,n_4,n_3,0,n_7,n_5,n_8,n_6),`$
$`I_g(n_1,n_2,|n_3|,n_4,n_5,n_6,n_7,n_8;n_{13},n_{23})=(\mathrm{𝟏}^{}+\mathrm{𝟐}^{}1)^{|n_3|}\times `$
$`\times (\mathrm{𝟒}^{}\mathrm{𝟔}^{}+\mathrm{𝟖}^{}+\mathrm{𝟎}^+)^{n_{13}}(0+)^{n_{23}}I_d(n_1,n_4,n_2,n_5,n_7,0,n_6,n_8),`$
$`I_g(n_1,n_2,n_3,|n_4|,n_5,n_6,n_7,n_8;n_{13},n_{23})=(1\mathrm{𝟏}^{}+\mathrm{𝟐}^{})^{|n_4|}\times `$
$`\times (\mathrm{𝟒}^{}+\mathrm{𝟔}^{}\mathrm{𝟎}^+)^{n_{13}}(\mathrm{𝟒}^{}\mathrm{𝟓}^{}+\mathrm{𝟖}^{})^{n_{23}}I_b(n_1,n_3,n_2,n_5,0,n_6,n_8,n_7).`$
If $`n_{1,3,4}`$ are all positive, we can lower them by (93) until one of them vanish.
### 3.6 Diagrams with five heavy-quark lines
Let’s define (figure 11*h*)
$$\begin{array}{ccccccccccc}\multicolumn{5}{c}{\frac{N_{12}^{n_{12}}N_{13}^{n_{13}}N_{23}^{n_{23}}d^dk_1d^dk_2d^dk_3}{D_1^{n_1}D_2^{n_2}D_3^{n_3}D_4^{n_4}D_5^{n_5}D_6^{n_6}D_7^{n_7}D_8^{n_8}}}& =& \multicolumn{5}{c}{i\pi ^{3d/2}(2\omega )^{3d+2(n_{12}+n_{13}+n_{23})2_{i=6}^8n_i}\times }\\ & & & & & & \multicolumn{5}{c}{\times I_h(n_1,n_2,n_3,n_4,n_5,n_6,n_7,n_8;n_{12},n_{13},n_{23})}\\ \hfill D_1& =& \frac{(k_2+p)v}{\omega },\hfill & & \hfill D_2& =& \frac{(k_1+k_2+p)v}{\omega },\hfill & & \hfill D_3& =& \frac{(k_1+p)v}{\omega },\hfill \\ \hfill D_4& =& \frac{(k_1+k_3+p)v}{\omega },\hfill & & \hfill D_5& =& \frac{(k_3+p)v}{\omega },\hfill & & \hfill D_6& =& k_2^2,\hfill \\ \hfill D_7& =& k_3^2,\hfill & & \hfill D_8& =& k_1^2,\hfill & & & & \\ \hfill N_{12}& =& 2k_1k_2,\hfill & & \hfill N_{13}& =& 2k_1k_3,\hfill & & \hfill N_{23}& =& 2k_2k_3.\hfill \end{array}$$
This integral is mirror-symmetric with respect to $`15`$, $`24`$, $`67`$. It vanishes when the indices of the following groups of lines are non-positive: 6, 7, 8, 12, 45, 234. There are two linear relations among the heavy-quark denominators:
$$\left[1\mathrm{𝟏}^{}+\mathrm{𝟐}^{}\mathrm{𝟑}^{}\right]I_h=0,\left[1\mathrm{𝟑}^{}+\mathrm{𝟒}^{}\mathrm{𝟓}^{}\right]I_h=0.$$
(94)
The cases
$`I_h(n_1,|n_2|,n_3,|n_4|,n_5,n_6,n_7,n_8;n_{12},n_{13},n_{23})=`$
$`=K_a(n_1,n_5,n_3,n_6,n_7,n_8;|n_2|,|n_4|;n_{23},n_{12},n_{13}),`$
$`I_h(n_1,|n_2|,n_3,n_4,|n_5|,n_6,n_7,n_8;n_{12},n_{13},n_{23})=`$
$`=K_b(n_1,n_4,n_3,n_6,n_7,n_8;|n_2|,|n_5|,0,0;n_{12},n_{13},n_{23}),`$
$`I_h(n_1,|n_2|,|n_3|,n_4,n_5,n_6,n_7,n_8;n_{12},n_{13},n_{23})=`$
$`=K_b(n_1,n_4,n_5,n_6,n_8,n_7;0,|n_3|,0,|n_2|;n_{23},n_{13},n_{12}),`$
$`I_h(|n_1|,n_2,n_3,n_4,|n_5|,n_6,n_7,n_8;n_{12},n_{13},n_{23})=`$
$`=K_c(n_2,n_4,n_3,n_6,n_7,n_8;|n_1|,|n_5|,0;n_{12},n_{13},n_{23}),`$
$`I_h(|n_1|,n_2,|n_3|,n_4,n_5,n_6,n_7,n_8;n_{12},n_{13},n_{23})=`$
$`=K_e(n_2,n_5,n_4,n_7,n_8,n_6;|n_1|,0,|n_3|;n_{12},n_{23},n_{13}),`$
as well as the symmetric cases $`n_40`$, $`n_10`$; $`n_40`$, $`n_30`$; $`n_50`$, $`n_30`$, reduce to $`I_c`$, $`I_d`$, as will be discussed later in this section. When $`n_30`$ (figure 12*f*), we can use
$$\left[\mathrm{𝟏}^{}\mathrm{𝟐}^{}+\mathrm{𝟒}^{}\mathrm{𝟓}^{}\right]I_h=0$$
(95)
to raise, say, $`n_1`$ and kill one of the lines 2, 4, 5. When $`n_{1,2}`$ are both positive, we can use (94) to kill one of the lines 1, 2, 3; when $`n_{4,5}`$ are both positive, we can kill one of the lines 3, 4, 5.
Let’s define (figure 11*i*)
$$\begin{array}{ccccccccccc}\multicolumn{5}{c}{\frac{N_{12}^{n_{12}}N_{13}^{n_{13}}N_{23}^{n_{23}}d^dk_1d^dk_2d^dk_3}{D_1^{n_1}D_2^{n_2}D_3^{n_3}D_4^{n_4}D_5^{n_5}D_6^{n_6}D_7^{n_7}D_8^{n_8}}}& =& \multicolumn{5}{c}{i\pi ^{3d/2}(2\omega )^{3d+2(n_{12}+n_{13}+n_{23})2_{i=6}^8n_i}\times }\\ & & & & & & \multicolumn{5}{c}{\times I_i(n_1,n_2,n_3,n_4,n_5,n_6,n_7,n_8;n_{12},n_{13},n_{23}),}\\ \hfill D_1& =& \frac{(k_1+p)v}{\omega },\hfill & & \hfill D_2& =& \frac{(k_1+k_2+p)v}{\omega },\hfill & & & & \\ \hfill D_3& =& \frac{(k_1+k_2+k_3+p)v}{\omega },\hfill & & \hfill D_4& =& \frac{(k_1+k_3+p)v}{\omega },\hfill & & \hfill D_5& =& \frac{(k_2+p)v}{\omega },\hfill \\ \hfill D_6& =& k_1^2,\hfill & & \hfill D_7& =& k_2^2,\hfill & & \hfill D_8& =& k_3^2.\hfill \end{array}$$
This integral vanishes when the indices of the following groups of lines are non-positive: 6, 7, 8, 34, 123. There are two linear relations among the heavy-quark denominators:
$$\left[1\mathrm{𝟏}^{}+\mathrm{𝟑}^{}\mathrm{𝟒}^{}\right]I_i=0,\left[1\mathrm{𝟏}^{}+\mathrm{𝟐}^{}\mathrm{𝟓}^{}\right]I_i=0.$$
(96)
The cases
$`I_i(n_1,|n_2|,|n_3|,n_4,n_5,n_6,n_7,n_8;n_{12},n_{13},n_{23})=`$
$`=K_b(n_1,n_4,n_5,n_6,n_8,n_7;|n_2|,0,|n_3|,0;n_{12},n_{23},n_{13}),`$
$`I_i(|n_1|,n_2,|n_3|,n_4,n_5,n_6,n_7,n_8;n_{12},n_{13},n_{23})=`$
$`=K_c(n_2,n_4,n_5,n_6,n_8,n_7;|n_1|,0,|n_3|;n_{12},n_{23},n_{13}),`$
$`I_i(|n_1|,|n_2|,n_3,n_4,n_5,n_6,n_7,n_8;n_{12},n_{13},n_{23})=`$
$`=K_d(n_3,n_5,n_4,n_7,n_6,n_8;|n_1|,0,0,|n_2|;n_{23},n_{12},n_{13}),`$
$`I_i(|n_1|,n_2,n_3,|n_4|,n_5,n_6,n_7,n_8;n_{12},n_{13},n_{23})=`$
$`=K_d(n_3,n_5,n_2,n_7,n_8,n_6;0,|n_1|,0,|n_4|;n_{12},n_{23},n_{13}),`$
$`I_i(n_1,n_2,n_3,|n_4|,|n_5|,n_6,n_7,n_8;n_{12},n_{13},n_{23})=`$
$`=K_d(n_3,n_1,n_2,n_6,n_8,n_7;0,|n_5|,|n_4|,0;n_{12},n_{13},n_{23}),`$
$`I_i(n_1,n_2,|n_3|,n_4,|n_5|,n_6,n_7,n_8;n_{12},n_{13},n_{23})=`$
$`=K_e(n_4,n_1,n_2,n_6,n_7,n_8;0,|n_3|,|n_5|;n_{23},n_{13},n_{12}),`$
$`I_i(n_1,|n_2|,n_3,|n_4|,n_5,n_6,n_7,n_8;n_{12},n_{13},n_{23})=`$
$`=K_f(n_1,n_5,n_3,n_6,n_7,n_8;|n_2|,0,|n_4|;n_{13},n_{23},n_{12}),`$
$`I_i(|n_1|,n_2,n_3,n_4,|n_5|,n_6,n_7,n_8;n_{12},n_{13},n_{23})=`$
$`=K_g(n_2,n_4,n_3,n_6,n_8,n_7;|n_5|,|n_1|,0;n_{12},n_{23},n_{13})`$
reduce to $`I_c`$, $`I_d`$, as will be discussed later in this section. When $`n_10`$, we can use
$$\left[\mathrm{𝟐}^{}\mathrm{𝟑}^{}+\mathrm{𝟒}^{}\mathrm{𝟓}^{}\right]I_i=0$$
(97)
to raise, say, $`n_2`$ and kill one of the lines 3, 4, 5. When $`n_{2,5}`$ are both positive, we can use (96) to kill one of the lines 1, 2, 5; when $`n_{3,4}`$ are both positive, we can kill one of the lines 1, 3, 4.
Finally, let’s define (figure 11*j*)
$$\begin{array}{ccccccccccc}\multicolumn{5}{c}{\frac{N_{12}^{n_{12}}N_{13}^{n_{13}}N_{23}^{n_{23}}d^dk_1d^dk_2d^dk_3}{D_1^{n_1}D_2^{n_2}D_3^{n_3}D_4^{n_4}D_5^{n_5}D_6^{n_6}D_7^{n_7}D_8^{n_8}}}& =& \multicolumn{5}{c}{i\pi ^{3d/2}(2\omega )^{3d+2(n_{12}+n_{13}+n_{23})2_{i=6}^8n_i}\times }\\ & & & & & & \multicolumn{5}{c}{\times I_j(n_1,n_2,n_3,n_4,n_5,n_6,n_7,n_8;n_{12},n_{13},n_{23}),}\\ \hfill D_1& =& \frac{(k_1+p)v}{\omega },\hfill & & \hfill D_2& =& \frac{(k_1+k_2+p)v}{\omega },\hfill & & & & \\ \hfill D_3& =& \frac{(k_1+k_2+k_3+p)v}{\omega },\hfill & & \hfill D_4& =& \frac{(k_2+k_3+p)v}{\omega },\hfill & & \hfill D_5& =& \frac{(k_3+p)v}{\omega },\hfill \\ \hfill D_6& =& k_1^2,\hfill & & \hfill D_7& =& k_2^2,\hfill & & \hfill D_8& =& k_3^2.\hfill \end{array}$$
This integral is mirror-symmetric with respect to $`15`$, $`24`$, $`68`$. It vanishes when the indices of the following groups of lines are non-positive: 6, 7, 8, 123, 234, 345. There are two linear relations among the heavy-quark denominators:
$$\left[1\mathrm{𝟏}^{}+\mathrm{𝟑}^{}\mathrm{𝟒}^{}\right]I_j=0,\left[1\mathrm{𝟐}^{}+\mathrm{𝟑}^{}\mathrm{𝟓}^{}\right]I_j=0.$$
(98)
The cases
$`I_j(n_1,|n_2|,|n_3|,n_4,n_5,n_6,n_7,n_8;n_{12},n_{13},n_{23})=`$
$`=K_b(n_1,n_4,n_5,n_6,n_7,n_8;0,0,|n_3|,|n_2|;n_{13},n_{23},n_{12}),`$
$`I_j(|n_1|,|n_2|,n_3,n_4,n_5,n_6,n_7,n_8;n_{12},n_{13},n_{23})=`$
$`=K_d(n_3,n_5,n_4,n_8,n_6,n_7;|n_1|,0,|n_2|,0;n_{23},n_{13},n_{12}),`$
$`I_j(|n_1|,n_2,|n_3|,n_4,n_5,n_6,n_7,n_8;n_{12},n_{13},n_{23})=`$
$`=K_e(n_2,n_5,n_4,n_8,n_7,n_6;|n_1|,|n_3|,0;n_{12},n_{13},n_{23}),`$
$`I_j(n_1,|n_2|,n_3,|n_4|,n_5,n_6,n_7,n_8;n_{12},n_{13},n_{23})=`$
$`=K_f(n_1,n_5,n_3,n_6,n_8,n_7;0,|n_2|,|n_4|;n_{12},n_{23},n_{13}),`$
$`I_j(|n_1|,n_2,n_3,n_4,|n_5|,n_6,n_7,n_8;n_{12},n_{13},n_{23})=`$
$`=K_g(n_2,n_4,n_3,n_6,n_8,n_7;0,|n_1|,|n_5|;n_{12},n_{23},n_{13}),`$
as well as the symmetric cases $`n_40`$, $`n_30`$; $`n_50`$, $`n_40`$; $`n_50`$, $`n_30`$, reduce to $`I_c`$, $`I_d`$, as will be discussed later in this section. When $`n_30`$, we can use
$$\left[\mathrm{𝟏}^{}\mathrm{𝟐}^{}+\mathrm{𝟒}^{}\mathrm{𝟓}^{}\right]I_j=0$$
(99)
to raise, say, $`n_1`$ and kill one of the lines 2, 4, 5. When $`n_{1,4}`$ are both positive, we can use (98) to kill one of the lines 1, 3, 4; when $`n_{2,5}`$ are both positive, we can kill one of the lines 2, 3, 5.
The reduced forms of the integrals $`I_{h,i,j}`$ are (figure 13)
$`K_a(n_1,n_2,n_3,n_4,n_5,n_8;n_9,n_{10};n_{45},n_{48},n_{58})=(\mathrm{𝟏}^{}+\mathrm{𝟑}^{}1)^{n_9}(\mathrm{𝟐}^{}+\mathrm{𝟑}^{}1)^{n_{10}}\times `$
$`\times (\mathrm{𝟎}^+)^{n_{45}}(\mathrm{𝟔}^{}\mathrm{𝟒}^{}\mathrm{𝟖}^{})^{n_{48}}(\mathrm{𝟕}^{}\mathrm{𝟓}^{}\mathrm{𝟖}^{})^{n_{58}}I_c(n_1,n_2,n_3,n_4,n_5,0,0,n_8),`$
$`K_b(n_1,n_2,n_3,n_4,n_7,n_8;n_9,n_{10},n_{11},n_{12};n_{48},n_{78},n_{47})=(\mathrm{𝟏}^{}+\mathrm{𝟑}^{}1)^{n_9}\times `$
$`\times (\mathrm{𝟐}^{}\mathrm{𝟑}^{}+1)^{n_{10}}(\mathrm{𝟏}^{}+\mathrm{𝟐}^{}1)^{n_{11}}(\mathrm{𝟏}^{}+\mathrm{𝟐}^{}\mathrm{𝟑}^{})^{n_{12}}(\mathrm{𝟔}^{}\mathrm{𝟒}^{}\mathrm{𝟖}^{})^{n_{48}}\times `$
$`\times (\mathrm{𝟕}^{}\mathrm{𝟓}^{}+\mathrm{𝟖}^{})^{n_{78}}(\mathrm{𝟒}^{}\mathrm{𝟔}^{}+\mathrm{𝟖}^{}+\mathrm{𝟎}^+)^{n_{47}}I_c(n_1,n_2,n_3,n_4,0,0,n_7,n_8),`$
$`K_c(n_1,n_2,n_3,n_6,n_7,n_8;n_9,n_{10},n_{11};n_{68},n_{78},n_{67})=(\mathrm{𝟏}^{}\mathrm{𝟑}^{}+1)^{n_9}\times `$
$`\times (\mathrm{𝟐}^{}\mathrm{𝟑}^{}+1)^{n_{10}}(\mathrm{𝟏}^{}+\mathrm{𝟐}^{}\mathrm{𝟑}^{})^{n_{11}}(\mathrm{𝟔}^{}\mathrm{𝟒}^{}+\mathrm{𝟖}^{})^{n_{68}}(\mathrm{𝟕}^{}\mathrm{𝟓}^{}+\mathrm{𝟖}^{})^{n_{78}}\times `$
$`\times (\mathrm{𝟒}^{}+\mathrm{𝟓}^{}\mathrm{𝟔}^{}\mathrm{𝟕}^{}+\mathrm{𝟎}^+)^{n_{67}}I_c(n_1,n_2,n_3,0,0,n_6,n_7,n_8),`$
$`K_d(n_1,n_2,n_3,n_5,n_6,n_7;n_9,n_{10},n_{11},n_{12};n_{57},n_{56},n_{67})=(\mathrm{𝟏}^{}\mathrm{𝟑}^{}+1)^{n_9}\times `$
$`\times (\mathrm{𝟑}^{}\mathrm{𝟐}^{}+1)^{n_{10}}(\mathrm{𝟏}^{}\mathrm{𝟐}^{}+1)^{n_{11}}(\mathrm{𝟏}^{}+\mathrm{𝟐}^{}\mathrm{𝟑}^{})^{n_{12}}(\mathrm{𝟓}^{}+\mathrm{𝟕}^{}\mathrm{𝟖}^{})^{n_{57}}\times `$
$`\times (\mathrm{𝟓}^{}\mathrm{𝟕}^{}+\mathrm{𝟖}^{}+\mathrm{𝟎}^+)^{n_{56}}(\mathrm{𝟔}^{}+\mathrm{𝟕}^{}\mathrm{𝟒}^{}\mathrm{𝟓}^{}\mathrm{𝟎}^+)^{n_{67}}\times `$
$`\times I_c(n_1,n_2,n_3,0,n_5,n_6,n_7,0),`$
$`K_e(n_1,n_2,n_3,n_5,n_6,n_8;n_9,n_{10},n_{11};n_{68},n_{58},n_{56})=(\mathrm{𝟏}^{}+\mathrm{𝟐}^{}\mathrm{𝟑}^{})^{n_9}\times `$
$`\times (\mathrm{𝟏}^{}+\mathrm{𝟐}^{}1)^{n_{10}}(\mathrm{𝟑}^{}\mathrm{𝟐}^{}+1)^{n_{11}}(\mathrm{𝟔}^{}\mathrm{𝟒}^{}+\mathrm{𝟖}^{})^{n_{68}}(\mathrm{𝟕}^{}\mathrm{𝟓}^{}\mathrm{𝟖}^{})^{n_{58}}\times `$
$`\times (\mathrm{𝟓}^{}\mathrm{𝟕}^{}+\mathrm{𝟖}^{}+\mathrm{𝟎}^+)^{n_{56}}I_d(n_1,n_2,n_3,0,n_5,n_6,0,n_8),`$
$`K_f(n_1,n_2,n_3,n_4,n_5,n_8;n_9,n_{10},n_{11};n_{48},n_{58},n_{45})=(\mathrm{𝟏}^{}+\mathrm{𝟐}^{}1)^{n_9}\times `$
$`\times (\mathrm{𝟑}^{}\mathrm{𝟐}^{}+1)^{n_{10}}(\mathrm{𝟑}^{}\mathrm{𝟏}^{}+1)^{n_{11}}(\mathrm{𝟒}^{}\mathrm{𝟔}^{}+\mathrm{𝟖}^{})^{n_{48}}(\mathrm{𝟓}^{}\mathrm{𝟕}^{}+\mathrm{𝟖}^{})^{n_{58}}\times `$
$`\times (\mathrm{𝟎}^+)^{n_{45}}I_d(n_1,n_2,n_3,n_4,n_5,0,0,n_8),`$
$`K_g(n_1,n_2,n_3,n_6,n_7,n_8;n_9,n_{10},n_{11};n_{68},n_{78},n_{67})=(\mathrm{𝟏}^{}+\mathrm{𝟐}^{}\mathrm{𝟑}^{})^{n_9}\times `$
$`\times (\mathrm{𝟑}^{}\mathrm{𝟐}^{}+1)^{n_{10}}(\mathrm{𝟑}^{}\mathrm{𝟏}^{}+1)^{n_{11}}(\mathrm{𝟔}^{}\mathrm{𝟒}^{}+\mathrm{𝟖}^{})^{n_{68}}(\mathrm{𝟕}^{}\mathrm{𝟓}^{}+\mathrm{𝟖}^{})^{n_{78}}\times `$
$`\times (\mathrm{𝟒}^{}+\mathrm{𝟓}^{}\mathrm{𝟔}^{}\mathrm{𝟕}^{}+\mathrm{𝟎}^+)^{n_{67}}I_d(n_1,n_2,n_3,0,0,n_6,n_7,n_8).`$
## 4 Implementation and testing
The package Grinder consists of a set of mutually recursive procedures for Feynman integrals of various topologies, which reduce a given Feynman integral to simpler ones, until boundary-case integrals with known values are reached. The package is written in REDUCE . Remembering results of previous function calls may make the program run much faster, if there is enough memory (unfortunately, REDUCE uses linear look-up in remember-tables).
I also re-implemented it in Axiom . All expressions involved are linear combinations of basis integrals with coefficients which are rational functions of $`d`$. It is convenient to use Axiom domain Vector Fraction UnivariatePolynomial, which has all the necessary operations. This makes intermediate expressions shorter than in the case when multivariate rational functions are used for entire expressions, because, typically, not all basis integrals are accompanied by every possible denominator. The amount of GCD calculations is thus reduced. This improvement can be, in principle, back-propagated to the REDUCE implementation by using matrices. However, working with matrices in REDUCE is awkward, because there are no local matrix variables, and no easy way for a function to return a matrix. On the other hand, complete diagram calculations, including tensor and $`\gamma `$-matrix algebra, can be done in REDUCE.
The main method of testing was checking various recurrence relations (including those which were not directly used for construction of the algorithm) in nested loops over $`n_i`$. For each of two-loop and generalized two-loop integrals, which depend on 5 indices, a typical number of checks was about 20000; each test set of this size runs for a few hours. Integrals with 8 or 9 indices are more difficult to check. Some test sets required a few days of CPU time. A tool showing how many times each linear code segment has been executed would be invaluable for setting up test cases which check all branches at least once. Unfortunately, such a tool is not available in either programming system, and I had to emulate it by hand.
###### Acknowledgments.
I am grateful to P.A. Baikov, D.J. Broadhurst, K.G. Chetyrkin, S.A. Larin and J.A.M. Vermaseren for useful discussions, and to INTAS for the grant which allowed me to buy Axiom. |
warning/0002/cond-mat0002190.html | ar5iv | text | # Josephson coupling and plasma resonance in vortex crystal
## 1 Introduction
Josephson coupling characterizes the ability of layered superconductors to carry supercurrents across the layers. In very anisotropic superconductors this coupling is suppressed by magnetic field applied along the c-axis. Thermal fluctuations and uncorrelated pinning lead to misalignment of pancake vortices induced by the magnetic field (see Fig. 1). Misalignment results in nonzero phase difference and in the suppression of the Josephson interlayer coupling. This suppression is quantitatively characterized by the “local coherence factor” $`𝒞\mathrm{cos}\phi _{n,n+1}(𝐫)`$, where $`\phi _{n,n+1}(𝐫)`$ is the gauge-invariant phase difference between layers $`n`$ and $`n+1`$, $`\mathrm{}`$ means average over thermal disorder and pinning. Josephson plasma resonance (JPR) measurements in highly anisotropic layered superconductors probe directly the interlayer Josephson coupling and the effect of pancake vortices on this coupling, because the squared JPR frequency, $`\omega _p^2`$, in the most part of vortex phase diagram, is proportional to the average interlayer Josephson energy ,
$$\omega _p^2\omega _0^2𝒞J_0𝒞,$$
(1)
where $`\omega _0(T)=c/\sqrt{ϵ_0}\lambda _c(T)`$ is the zero field plasma frequency, $`\lambda _c(T)`$ is the c-component of the London penetration depth, $`ϵ_0`$ is dielectric constant, and $`J_0`$ is the Josephson critical current.
The JPR measurements performed in the liquid vortex phase at relatively high magnetic fields, $`B>B_J=\mathrm{\Phi }_0/\lambda _J^2`$, revealed that the plasma frequency drops approximately as $`1/\sqrt{B}`$ , where $`\lambda _J=\gamma s`$ is the Josephson length, $`\gamma `$ is the anisotropy ratio and $`s`$ is the interlayer distance. The above dependence is characteristic for the pancake liquid weakly correlated along the $`c`$ axis. In this phase many pancake vortices contribute to the suppression of the phase difference at a given point. In contrast, in the vortex solid pancake vortices form aligned stacks and suppression of coupling is caused by weak misalignments of the pancake vortices due to the thermal fluctuations and random pinning. JPR measurements in Bi<sub>2</sub>Sr<sub>2</sub>CaCu<sub>2</sub>O<sub>8-δ</sub> (Bi-2212) crystals have shown that the JPR frequency decreases approximately linearly with field in the vortex solid. In the fields above 20 Oe the interlayer phase coherence changes drastically at the transition line, implying the decoupling nature of the first-order melting transition in agreement with theoretical expectations (see, e.g., Ref. ). On the other hand, at smaller field phase coherence does not change considerably at the melting point.
In this paper we consider the Josephson coupling and JPR in the vortex lattice. We focus on the suppression of coupling due to thermal fluctuations of pancake vortices near the equilibrium crystal positions, and neglect influence of pinning potential. For real Bi-2212 crystals this approximation is justified at sufficiently high temperatures ($`40`$ K). This problem has been considered in the past in the simple limiting cases. However quantitative calculation suitable for comparison with existing JPR data in a wide field range is still absent. At small fields, when vortices act independently, $`\omega _p^2`$ decreases linear with $`B`$. The linear dependence was observed experimentally in Refs. in solid state in Bi-2212 crystals. In fields below 20 Oe near $`T_c`$ this linear dependence extends to the liquid state providing evidence for a line structure of the vortices in the liquid at low fields. The regime of independent vortices has been considered in Ref. . In this paper we extend our consideration to higher fields up to the melting field.
## 2 Low fields. Isolated vortex lines
Consider small magnetic fields $`BB_J,B_\lambda \mathrm{\Phi }_0/4\pi \lambda _{ab}^2,`$. At these fields regions of suppressed coupling are localized near the vortex lines (pancake stacks) and do not overlap (single vortex regime). The field-induced change in $`𝒞`$ in this regime is given by $`\delta 𝒞1𝒞=BI/\mathrm{\Phi }_0`$, where $`I=d^2𝐫\left(1\mathrm{cos}\left(\phi _{n,n+1}(𝐫)\right)\right)`$, and $`\phi _{n,n+1}(𝐫)`$ is phase difference induced by fluctuation displacements $`𝐮_n`$ in a single line. The same integral determines the tilt stiffness due to the Josephson coupling. We split integration domain in $`I`$ into two region, $`r<R`$ and $`r>R`$, where $`R`$ is the intermediate scale $`r_w<R<\gamma s`$, with $`r_w^2𝐮_{n,n+1}^2`$ and $`𝐮_{n,n+1}𝐮_{n+1}𝐮_n`$. At $`r<R`$ we can neglect screening due to the Josephson currents and take
$$\phi _{n,n+1}(𝐫)=\mathrm{arctan}\frac{yu_{yn+1}}{xu_{xn+1}}\mathrm{arctan}\frac{yu_{yn}}{xu_{xn}}.$$
At $`r>R`$ we can take $`\phi _{n,n+1}`$in linear approximation with respect to $`𝐮_n`$,
$$\phi _{n,n+1}(𝐫,k_z)\frac{sdk_z}{2\pi }\stackrel{~}{k}_z\left[𝐮(k_z)\times \right]\mathrm{K}_0(\frac{\stackrel{~}{k}_zr}{\gamma }),$$
where $`𝐮(k_z)=s_n\mathrm{exp}(isk_zn)𝐮_n`$, $`\stackrel{~}{k}_z(2/s)\mathrm{sin}(sk_z/2)`$, and $`\mathrm{K}_0\left(z\right)`$ is the modified Bessel function. At intermediate distances $`r_wr\lambda _J`$, both expressions give the same simple result $`\phi _{n,n+1}(𝐫)[𝐫\times 𝐮_{n,n+1}]/r^2`$. Using above asymptotics of $`\phi _{n,n+1}`$ we obtain
$`I`$ $`={\displaystyle \frac{\pi }{2}}{\displaystyle \frac{dk_z}{2\pi }\left(1\mathrm{cos}\left(sk_z\right)\right)\left|u(k_z)\right|^2}`$
$`\times \mathrm{ln}\left({\displaystyle \frac{3.72\lambda _J^2}{u_{n,n+1}^2\left(1\mathrm{cos}\left(sk_z\right)\right)}}\right).`$
Weak logarithmic dependence on $`u_{n,n+1}`$ leads to the nonharmonic tilt energy. In the following we will eliminate this nonharmonicity using the self consistent harmonic approximation (SCHA), which results in the substitution $`\mathrm{ln}(A/u_{n,n+1}^2)\mathrm{ln}(0.24A/r_w^2)`$. Approximate evaluation of the above integral gives a simple practical relation connecting the field-induced correction of the plasma frequency $`\omega _p(B,T)`$ with $`r_w`$ for the case $`r_w<\lambda _J<a`$:
$$\frac{\omega _0^2(T)\omega _p^2(B,T)}{\omega _0^2(T)}\frac{\pi Br_w^2}{2\mathrm{\Phi }_0}\mathrm{ln}\frac{0.8\lambda _J}{r_w}.$$
(2)
This relation allows one to extract $`r_w^2`$ from the plasma resonance measurements.
We now calculate $`r_w^2`$ when wandering of the vortex lines is caused by thermal fluctuations. In the single vortex regime $`r_w^2`$ is determined by the wandering energy consisting of the Josephson and magnetic contributions,
$$_w\frac{1}{2}\frac{dk_z}{2\pi }\left[\epsilon _{1J}\stackrel{~}{k}_z^2+w_M\right]\left|u(k_z)\right|^2,$$
(3)
where $`\epsilon _{1J}(\epsilon _0/\gamma ^2)\mathrm{ln}\left(1.33\gamma /(r_w\stackrel{~}{k}_z)\right)`$ is the line tension due to the Josephson coupling, $`w_M(\epsilon _0/\lambda _{ab}^2)\mathrm{ln}(1.5\lambda _{ab}/r_w)`$ is the effective cage potential, which appears due to nonlocal magnetic interactions between pancake vortices in different layers (it describes the magnetic tilt stiffness at wave vectors $`k_z>1/r_w`$), and $`\epsilon _0\mathrm{\Phi }_0^2/(4\pi \lambda _{ab})^2`$. Assuming Gaussian fluctuations of pancake vortices we obtain
$$r_{wT}^2=\frac{8T}{sw_M}\frac{1}{1+\zeta +\sqrt{1+\zeta }},$$
(4)
where the parameter $`\zeta (T)4\lambda _{ab}^2(T)/\lambda _J^2`$ describes the relative roles of the Josephson and magnetic interactions. Substituting this result into Eq. (2) we obtain
$$\frac{\omega _0^2\omega _p^2(B)}{\omega _0^2}\frac{4\pi \lambda _{ab}^2BT}{sϵ_0\mathrm{\Phi }_0}\frac{1}{1+\zeta +\sqrt{1+\zeta }}.$$
(5)
This result of the single vortex regime is valid in both solid and liquid states for $`BB_J`$, because in this field range wandering of lines at short scales does not change much at the melting point. Eq. (5) describes fairly well the suppression of the plasma frequency at small fields . With the data of Ref. for underdoped Bi-2212 with T$`{}_{c}{}^{}84.5`$ K Eq. (2) gives unexpectedly large wandering length $`r_w1`$$`\mu `$m at $`77`$ K, which is comparable with both $`\lambda _{ab}`$ and $`\lambda _J`$ at this temperature. However, we found that this estimate is a in good agreement with the theoretical calculation (4).
## 3 High fields, $`B>B_J`$, $`B_\lambda `$
As the field increases two competing effects start to influence pancake fluctuations and field dependence of the average Josephson energy. Vortex interactions diminish pancake fluctuations. On the other hand, collective suppression of the Josephson energy decreases tilt stiffness and enhances pancake fluctuations. Using general relations connecting the phase perturbations with the elastic lattice deformations (see, e.g. Ref. ), we obtain
$`\delta 𝒞={\displaystyle \frac{\left[\phi _{n,n+1}\right]^2}{2}}{\displaystyle \frac{\left(2\pi sn_v\right)^2}{2}}`$ (6)
$`\times {\displaystyle }{\displaystyle \frac{d^2qdk_z}{\left(2\pi \right)^3}}{\displaystyle \underset{Q<q_m}{}}{\displaystyle \frac{\stackrel{~}{k}_z^2\left(\left[𝐐\times 𝐪\right]^2𝐮_l^2+\left(q^2+\mathrm{𝐐𝐪}\right)^2𝐮_t^2\right)}{q^2\left(\left(𝐪+𝐐\right)^2+\stackrel{~}{k}_z^2/\gamma ^2\right)^2}},`$
where $`n_vB/\mathrm{\Phi }_0`$, $`𝐐`$ are the reciprocal lattice vectors (the cut off $`q_m2.2/r_w`$ in the summation over $`𝐐`$ is established by comparison with the single vortex regime), and $`𝐮_l(𝐪,k_z)`$ ($`𝐮_t(𝐪,k_z)`$) are the longitudinal (transversal) elastic lattice displacements. Integration with respect to the in-plane wave vector $`𝐪`$ is limited by the first Brillouin zone, which we approximate by the circle $`q<K_0`$, $`K_0^2=4\pi n_v`$. For thermal Gaussian fluctuations the mean squared averages $`𝐮_{t,l}^2|𝐮_{t,l}(𝐪,k_z)|^2`$ are determined by the corresponding components of the elastic matrix $`\mathrm{\Phi }_{t,l}(𝐪,k_z)`$
$$|𝐮_{t,l}(𝐪,k_z)|^2=T/\mathrm{\Phi }_{t,l}(𝐪,k_z).$$
At high fields $`B>B_J`$, $`B_\lambda `$, and at large $`k_z`$, $`k_z1/\lambda _{ab}`$, $`\sqrt{n_v}`$ we have
$`\mathrm{\Phi }_t(q,k_z)`$ $`=C_{66}q^2+\mathrm{\Phi }_{44}(q,k_z),`$
$`\mathrm{\Phi }_l(q,k_z)`$ $`=\mathrm{\Phi }_{11}(q)+\mathrm{\Phi }_{44}(q,k_z),`$
where $`\mathrm{\Phi }_{11}\frac{B^2}{4\pi }\left(1\frac{q^2}{4K_0^2}\right)`$ is the compression stiffness, $`C_{66}=A_{66}n_vϵ_0/4`$ is the shear modulus, parameter $`A_{66}<1`$ describes fluctuation suppression of $`C_{66}`$, which we approximate as $`A_{66}=10.4B/B_m`$ with $`B_m`$ being the melting field, and
$`\mathrm{\Phi }_{44}(q,k_z){\displaystyle \frac{n_v\epsilon _0k_z^2}{2\gamma ^2}}\mathrm{ln}{\displaystyle \frac{0.11a^2}{r_w^2\left(10.53q^2/K_0^2\right)^2}}`$
$`+{\displaystyle \frac{n_v\epsilon _0}{2\lambda _{ab}^2}}\mathrm{ln}\left(0.5+{\displaystyle \frac{0.13a^2}{r_w^2}}\right)+{\displaystyle \frac{B^2}{4\pi \lambda _{ab}^2}}{\displaystyle \frac{k_z^2}{k_z^2+\gamma ^2q^2}},`$
is the tilt stiffness, computed within SCHA, with $`a1/\sqrt{n_v}`$. The wandering length $`r_w`$ has to be determined self-consistently from the equation
$`r_w^2`$ $`=`$ $`2{\displaystyle \frac{d^2qdk_z}{\left(2\pi \right)^3}\left(1\mathrm{cos}(sk_z)\right)}`$
$`\times `$ $`\left(|𝐮_l(𝐪,k_z)|^2+|𝐮_t(𝐪,k_z)|^2\right).`$
Expression (6) for $`\delta 𝒞`$ can be naturally split into the collective contribution $`\delta 𝒞_{coll}`$, corresponding to $`Q=0`$ term in the $`Q`$-summation,
$$\delta 𝒞_{coll}\frac{\left(2\pi sn_v\right)^2}{2}\frac{d^2qdk_z}{\left(2\pi \right)^3}\frac{\stackrel{~}{k}_z^2q^2𝐮_t^2}{\left(𝐪^2+\stackrel{~}{k}_z^2/\gamma ^2\right)^2},$$
(8)
and the local contribution $`\delta 𝒞_{loc}`$ coming from $`Q>0`$ terms. At high fields $`B\mathrm{\Phi }_0/\lambda _J^2`$ we obtain approximate expression for $`\delta 𝒞_{loc}`$, which resembles the single-vortex result (2)
$$\delta 𝒞_{loc}\frac{\pi n_vr_w^2}{2}\mathrm{ln}\frac{0.58a}{r_w}.$$
(9)
If we consider suppression of coupling in the cylindrical Bravais cell near the chosen vortex line, then the local term determines suppression of coupling caused by this vortex line and the collective term determines suppression of coupling caused by all other vortex lines. In general, relative role of the collective term in $`\delta 𝒞`$ grows with field. We use above expressions to calculate the field dependence of $`𝒞`$ for comparison with JPR data.
Recently detailed measurements of field dependence of JPR frequency in the vortex crystal state of Bi-2212 have been done by M. Gaifullin et al. using frequency scan. To compare our calculations with JPR data we need to know $`\lambda _{ab}(T)`$ and $`\gamma =\lambda _c/\lambda _{ab}`$. $`\lambda _c`$ is extracted directly from JPR frequency at $`B=0`$, $`\lambda _c(T)=c/\sqrt{ϵ_0}\omega _0(T)`$ taking $`ϵ_0=11`$ and $`\gamma `$ is chosen as a fitting parameter. Fig. 2 compares the computed dependence $`𝒞(B)`$ with $`(\omega _p(B)/\omega _p(0))^2`$ for three values of temperature. We also show obtained values of $`\lambda \lambda _{ab}`$ and $`\gamma `$. We obtain $`\gamma `$, that slightly grows with temperature (from $`460`$ at $`40`$ K to $`510`$ at $`60`$ K). Such enhancement of $`\gamma `$ is expected due to the phase fluctuations.
In conclusion, we have calculated the field dependence of the JPR frequency in the vortex crystal. In the single vortex regime at low magnetic fields the JPR provides a direct probe for meandering of individual lines. The wandering length extracted from the JPR data is in agreement with the theoretical calculations. Our theory of pancake fluctuations gives a very good description of the field dependence of the plasma frequency up to the melting field. The authors thank M. Gaifullin, Y. Matsuda, T. Tamegai, and T. Shibauchi for providing their experimental data prior publication and V. Vinokur for constructive comments. |
warning/0002/hep-ph0002057.html | ar5iv | text | # A new determination of the Pomeron intercept in hard processes
## Abstract
A method allowing for a direct comparison of data with theoretical predictions is proposed for forward jet production at HERA. It avoids the reconstruction of multi-parton contributions by expressing the experimental cuts directly as correction factors on the QCD forward jet cross-section. An application to the determination of the effective Pomeron intercept in the BFKL-LO parametrization from $`d\sigma /dx`$ data at HERA leads to a good fit with a significantly higher effective intercept, $`\alpha _P=1.43\pm 0.025(stat.)\pm 0.025(syst.),`$ than for proton (total and diffractive) structure functions. It is however less than the value of the pomeron intercept using dijets with large rapidity intervals obtained at Tevatron. We also evaluate the rapidity veto contribution to the higher order BFKL corrections. The method can be extended to other theoretical inputs.
1. Introduction
The study of forward jets at colliders is considered as the milestone of QCD studies at high energies, since it provides a direct way of testing the perturbative resummations of soft gluon radiation. More precisely, the study of one forward jet (w.r.t. the proton) in an electron-proton collider seems to be a good candidate to test the energy dependence of hard QCD cross-sections. It is similar to the previous proposal of studying two jets separated by a large rapidity interval in hadronic colliders , for which only preliminary results are available . This test is also possible in $`\gamma ^{}`$-$`\gamma ^{}`$ scattering but here the statistics and the energy range are still insufficient to get a reliable determination of the physical parameters for hard QCD cross-sections. Indeed, the proposed (and favored for the moment being) set-up is to consider jets with transverse momentum $`k_T`$ of the order of the photon virtuality $`Q`$ allowing to damp the QCD evolution as a function of $`k_T`$ (DGLAP evolution ) in favor of the evolution in energy at fixed $`k_T`$ (BFKL evolution ).
Since proposal was made, a set of interesting studies have been performed to check its relevance. On the experimental ground, H1 and ZEUS have published useful results with appropriate cuts (to be displayed later on) at relatively small $`x.`$ On the theoretical ground, the general formulation and some quantitative estimates have been performed prior to experiments confirming the interest in such processes. The recent theoretical analyses have been mainly based on the use of Monte-Carlo simulations, including the multi-parton cross-sections and starting from the various frameworks in competition . Quite a few analyses arrive at a satisfactory description of the data, taking into account the specific parametrizations which are choosen. Indeed, the BFKL-based Monte-Carlo lead to quite satisfactory results, while those based on DGLAP evolution meet some difficulty to describe the data<sup>*</sup><sup>*</sup>*Note however, that some more refined versions of DGLAP evolution including contributions from the resolved off-mass-shell photon can describe the data ..
However, there still remains a problem in the interpretation of those results. Due to the difficulty in handling the experimental cuts without introducing in the simulation the whole set of theoretical n-parton contributions to the cross-section, it appears difficult to avoid the uncertainties of the reconstruction (with the parameters and constraints which are needed to define the scheme in practice). It seems thus difficult to determine unambiguously genuine theoretical parameters defining the cross-section one is looking at. One example is the dependence in parameters such as infra-red cut-offs, which are not a-priori required in the expression of the total $`d\sigma /dx`$ jet cross sections. Another illustration is the so-called “consistency constraint” which appears very useful in the expression of next-to-leading corrections to the BFKL formula coming from the n-parton contributions, but again is not expressed in terms of the $`d\sigma /dx`$ jet cross sections itself. In fact, it does not seem easy, in those schemes, to extract with some precision the value of the effective Pomeron intercept $`\alpha _P`$, i.e. the main theoretical parameter describing the theoretical energy dependence in this process. As we know, this parameter is of primordial importance to evaluate the amount of next-leading corrections in a BFKL framework and to confront its effective value with the recent theoretical determinations .
We want to address this problem in a quite different way, that is on focussing on the jet cross section $`d\sigma /dx`$ observable itself, by a consistent treatment of the experimental cuts and minimizing the uncertainties for that particular observable. Let us remark that our approach is not intended to provide a substitution to the other methods, since the Monte-Carlo simulations have the great merit of making a set of predictions for various observables. Hence, our method has to be considered as complementary to the others and dedicated to a better determination of the effective Pomeron intercept using the $`d\sigma /dx`$ data. As we shall see, it will fix more precisely this parameter, but it will leave less constrained other interesting parameters, such as the cross-section normalization.
One fruitful outcome of the method proposed in the present paper is the possibility of comparing the effective intercept with its determination in other processes involving QCD at high energy. In fact, using the parameters determined from forward jets at HERA, it is possible to compare with double jet production at Tevatron following Ref. for which preliminary experimental analyses have been performed and find a high value of the intercept ($`\alpha _P=1.7\pm .1\pm .1`$ in Ref. ). It can also be confronted with the effective BFKL analysis of proton structure functions at small-$`x_{Bj},`$ which give rather low valuesNote however that taking into account the full BFKL formula may lead to higher $`\alpha _P,`$ namely $`1.21.3,`$ see Ref. , for total structure functions and even reach $`.4,`$ see Ref. , for diffractive proton structure functions.($`\alpha _P1.11.21`$ see Ref. ). However, in those cases, the result may be different, since non-perturbative effects related to the “soft” proton scales are expected to influence the determination of parameters.
Thus, the comparison of the effective BFKL parameter $`\alpha _P`$ obtained for the forward jet production cross-section allows for a study of QCD at high energy, aiming at a better understanding of the corrections to the leading-order BFKL predictions .
The plan of our study is the following: in section 2, we introduce the QCD formalism and our method for determining $`\alpha _P.`$ In the following section 3 we determine the kinematic correction factors to the forward jet cross-section data on $`d\sigma /dx`$ due to the experimental cuts. In the following section 4, we perform and discuss a (separately and then common) fit to H1 and ZEUS data. This determines the BFKL parameter $`\alpha _P`$ which is subsequently used in section 4 for a comparison with the two-jet cross-section at Tevatron from the experimental D0 analysis. Discussions on these results and comparison with the BFKL study of (total and diffractive) structure functions are presented in section 5 and conclusions and outlook in section 6.
2. Formalism
The cross-section for forward jet production at HERA in the dipole model reads :
$`{\displaystyle \frac{d^{(4)}\sigma }{dxdQ^2dx_Jdk_T^2d\mathrm{\Phi }}}={\displaystyle \frac{\pi N_C\alpha ^2\alpha _S(k_T^2)}{Q^4k_T^2}}f_{eff}(x,\mu _f^2)\mathrm{\Sigma }e_Q^2{\displaystyle _{\frac{1}{2}i\mathrm{}}^{\frac{1}{2}+i\mathrm{}}}{\displaystyle \frac{d\gamma }{2i\pi }}\left({\displaystyle \frac{Q^2}{k_T^2}}\right)^\gamma \times `$ (1)
$`\times \mathrm{exp}\{ϵ(\gamma ,0)Y\}\left[{\displaystyle \frac{h_T(\gamma )+h_L(\gamma )}{\gamma }}(1y)+{\displaystyle \frac{h_T(\gamma )}{\gamma }}{\displaystyle \frac{y^2}{2}}\right]\mathrm{exp}\{ϵ(\gamma ,1)Y\}cos2\mathrm{\Phi }\left[{\displaystyle \frac{h_T(\gamma )}{\gamma }}{\displaystyle \frac{\gamma (1\gamma )}{(\gamma +1)(2\gamma )}}\right]`$ (2)
where
$`Y`$ $`=`$ $`\mathrm{ln}{\displaystyle \frac{x_J}{x}}`$ (3)
$`ϵ(\gamma ,p)`$ $`=`$ $`\overline{\alpha }\left[2\psi (1)\psi (p+1\gamma )\psi (p+\gamma )\right]`$ (4)
$`f_{eff}(x,\mu _f^2)`$ $`=`$ $`G(x,\mu _f^2)+{\displaystyle \frac{4}{9}}\mathrm{\Sigma }(Q_f+\overline{Q_f})`$ (5)
$`\mu _f^2`$ $``$ $`k_T^2,`$ (6)
are, respectively, $`Y`$ the rapidity interval between the photon probe and the jet, $`ϵ(\gamma ,p)`$ the BFKL kernel eigenvalues, $`f_{eff}`$ the effective structure function combination, and $`\mu _f`$ the corresponding factorization scale. The main BFKL parameter is $`\overline{\alpha },`$ which is the (fixed) value of the effective strong coupling constant in LO-BFKL formulae. Note that we gave for completion the full BFKL formula including the azimuthal dependence but we will stick to the azimuth-independent contribution with the dominant $`\mathrm{exp}\{ϵ(\gamma ,0)Y\}`$ factor.
The so-called “impact factors”
$`\left(\begin{array}{c}h_T\\ h_L\end{array}\right)={\displaystyle \frac{\alpha _S(k_T^2)}{3\pi \gamma }}{\displaystyle \frac{(\mathrm{\Gamma }(1\gamma )\mathrm{\Gamma }(1+\gamma ))^3}{\mathrm{\Gamma }(22\gamma )\mathrm{\Gamma }(2+2\gamma )}}{\displaystyle \frac{1}{1\frac{2}{3}\gamma }}\left(\begin{array}{c}(1+\gamma )(1\frac{\gamma }{2})\\ \gamma (1\gamma )\end{array}\right),`$ (11)
are obtained from the $`k_T`$ factorization properties of the coupling of the BFKL amplitudes to external hard probes. The same factors can be related to the photon wave functions within the equivalent context of the QCD dipole model .
Our goal is to compare as directly as possible the theoretical parametrization (2) to the data which are collected in experiments . The crucial point is how to take into account the experimentally defined kinematic cuts listed in Table I for the reported three sets of data (two for H1 with $`k_T=3.5`$ or $`5`$ GeV, and one for ZEUS).
| H1 cuts | ZEUS cuts |
| --- | --- |
| $`E_e^{^{}}>`$ 11 GeV | $`E_e^{^{}}>`$ 10 GeV |
| 160 $`\theta _e^{}173`$ deg. | |
| $`y>`$0.1 | $`y>`$0.1 |
| 7$`\theta _{jet}`$20 deg. | $`\theta _{jet}`$8.5 deg. |
| $`k_{Tjet}`$ 3.5 or 5 GeV | $`k_{Tjet}`$ 5 GeV |
| $`x_{jet}>`$0.035 | $`x_{jet}>`$0.036 |
| $`0.5<\frac{k_T^2}{Q^2}<2`$ | $`0.5<\frac{k_T^2}{Q^2}<2`$ |
| $`10^4<x<4.10^3`$ | $`4.510^4<x<4.510^2`$ |
Table I- Experimental cuts (H1/ZEUS)
The main problem to solve is to investigate the effect of these cuts on the determination of the integration variables leading to a prediction for $`d\sigma /dx`$ from the given theoretical formula for $`d^{(4)}\sigma `$ as given in formula (2). The effect is expected to appear as bin-per-bin correction factors to be multiplied to the theoretical cross-sections for average values of the kinematic variables for a given $`x`$-bin before comparing to data (e.g. fitting the cross-sections).
The idea of our method is threefold: i) for each $`x`$-bin, determining the average values of $`x`$, $`Q^2`$, $`E_J`$, $`k_T`$ from a known and reliable Monte-Carlo simulation of the cross-sections. For this sake, we use the Ariadne Monte-Carlo programme ; ii) choosing a set of integration variables over $`d^{(4)}\sigma `$ in (2) in such a way to match closely the experimental cuts and minimize the variation of the cross-sections over the bin size; iii) fixing the correction factors due to the experimental cuts for each $`x`$-bin, by a random simulation of the kinematic constraints with no dynamical input.
The point i) proposed already in allows a determination of which average values of the kinematic variables have to be taken in the theoretical formula (2) for each experimental $`x`$-bin. The point ii) comes from the crucial requirement to minimize the variation (over the $`x`$-bin) of the variables to be retained for the integration. Indeed, since the integration procedure multiplies the central value of the integrand by the size of the integration bins, it is compulsory to choose adequate variables which lead to a smooth dependence of the integrand and of the effect of the kinematic cuts.
This double stringent requirement can be solved for the forward jet $`d\sigma /dx.`$ For this sake we choose
$`{\displaystyle \frac{d\sigma }{dx}}={\displaystyle \left[Q^6\frac{d^{(4)}\sigma }{dxdQ^2dx_Jdk_T^2d\mathrm{\Phi }}\right]\times \mathrm{\Delta }\left(\frac{1}{Q^2}\right)\mathrm{\Delta }x_J\mathrm{\Delta }\left(\frac{k_T^2}{Q^2}\right)\mathrm{\Delta }\mathrm{\Phi }}.`$ (12)
The property of this non-trivial choice is the following. The integration variables are choosen in such a way that the expression in the square brackets $`[Q^6\mathrm{}]`$ in (12) is dependent on the ratio $`\frac{k_T^2}{Q^2}`$ and not on each scale separately. Looking at the experimental cuts (see Table I), it becomes clear that the choice of this scale-invariant integrand minimizes the variation of the observable on the bin, while each scale $`k_T^2`$ and $`Q^2`$ presents large variations and thus would generate large integration errors. Indeed, various numerical studies we have performed have demonstrated that it was a sine qua non stability condition for the fits. The overall azimuthal integration ($`\mathrm{\Delta }\mathrm{\Phi }=2\pi )`$ cancels the second term in (2).
3. Correction Factors
The experimental correction factors have been determined using a toy Monte-Carlo designed as follows. We generate flat distributions in the variables $`k_T^2/Q^2`$, $`1/Q^2`$, $`x_J,`$ using reference intervals which include the whole of the experimental phase-space (the $`\mathrm{\Phi }`$ variable is not used in the generation since all the cross-section measurements are $`\varphi `$ independent). In practice, we get the correction factors by counting the numbers of events which fulfill the experimental cuts given in Table I for each $`x`$-bin. The correction factor is obtained by the ratio to the number of events which pass the experimental cuts and the kinematic constraints, and the number of events which fullfil only the kinematic constraints,i.e. the so-called reference bin.
| x | $`\sigma `$ | $`Q^2`$ | $`E_{jet}`$ | $`k_T`$ | Corr. Factor (.$`10^3`$) |
| --- | --- | --- | --- | --- | --- |
| 0.00036 | 202.5 | 13.9 | 32.6 | 4.5 | 0.270 |
| 0.00073 | 342. | 21.5 | 34.4 | 5.0 | 0.993 |
| 0.0012 | 224. | 26.9 | 36.9 | 5.5 | 1.14 |
| 0.0017 | 138. | 31.4 | 38.1 | 5.8 | 1.11 |
| 0.0024 | 67. | 38.1 | 38.8 | 6.3 | 0.921 |
| 0.0035 | 32. | 47.0 | 37.9 | 6.9 | 0.711 |
Table IIa- Average values of kinematic quantities and correction factors - H1 $`k_T>3.5GeV`$
| x | $`\sigma `$ | $`Q^2`$ | $`E_{jet}`$ | $`k_T`$ | Corr. Factor (.$`10^3`$) |
| --- | --- | --- | --- | --- | --- |
| 0.00036 | 27.5 | 18.0 | 35.9 | 5.5 | 0.108 |
| 0.00073 | 126. | 27.0 | 36.5 | 5.8 | 0.695 * |
| 0.0012 | 132. | 32.2 | 37.9 | 6.3 | 0.895 |
| 0.0017 | 96. | 34.8 | 39.3 | 6.5 | 0.979 |
| 0.0024 | 55. | 40.1 | 39.4 | 6.7 | 0.870 |
| 0.0035 | 28. | 48.2 | 39.6 | 7.2 | 0.696 |
Table IIb- Average values of kinematic quantities and correction factors - H1 $`k_T>5GeV`$
| x | $`\sigma `$ | $`Q^2`$ | $`E_{jet}`$ | $`k_T`$ | Corr. Factor (.$`10^3`$) |
| --- | --- | --- | --- | --- | --- |
| 0.0006 | 114.0 | 28.0 | 36.5 | 6.3 | 0.304 * |
| 0.0011 | 96.2 | 39.0 | 38.0 | 6.9 | 0.656 |
| 0.0019 | 77.8 | 50.7 | 39.9 | 7.6 | 0.966 |
| 0.0033 | 34.4 | 75.6 | 43.8 | 8.7 | 0.996 |
| 0.006 | 14.1 | 113.6 | 49.6 | 10.4 | 0.995 |
| 0.01 | 6.53 | 176.4 | 58.5 | 12.9 | 0.896 * |
| 0.018 | 2.65 | 244.7 | 67.3 | 15.1 | 0.653 * |
| 0.031 | 0.65 | 366.8 | 78.8 | 18.8 | 0.373 * |
Table IIc- Average values of kinematic quantities and correction factors - ZEUS $`k_T>5GeV`$
The correction factors are given in Table IIa for H1 ($`k_T>3.5`$ GeV), Table IIb for H1 ($`k_T>5`$ GeV), and Table IIc for ZEUS bins together with the value of the bin centers determined with the full Monte-Carlo simulation , and the experimental values of the cross-sections Note that we did not use the full Monte-Carlo to get the correction factors in order to avoid any strong model dependence as these factors are only due to kinematic effects. It is however more difficult to use a toy Monte Carlo to get accurate values for the bin centers, and this is why we used a full Monte Carlo for this sake. However, the dependence of the theoretical cross-section on the bin centers is minimized by our specific choice of kinematic variables (see formula (7)). . We note that the correction factors are quite different from one $`x`$-bin to an other and much less than one (in $`10^3`$ units), explicitely showing that the experimental cuts play an important role in the cross-section measurement, and that these factors are compulsory to be taken into account if we want to get a direct comparison with the theoretical cross-sections. We also note that the correction factors are very much different from one another at very low $`x`$, showing that the acceptance of these bins is quite low. This is also why it is not so easy to be able to get a correct value of the measured cross-section after cuts in those bins. We also get the same order of magnitude for the correction factors for the H1 and ZEUS experiments because the experimental cuts are quite similar. The differences between both experiments are due mainly to the fact that the range in $`x`$ and $`Q^2`$ is much lower for H1 than for ZEUS (the reference bin for H1 goes to lower $`Q^2`$ compared to ZEUS).
4. Fits
Using the kinematic correction factors determined as described in the previous section, we perform a fit to the H1 and ZEUS data with only two free parameters. these are the effective strong coupling constant in LO BFKL formulae $`\overline{\alpha }`$ corresponding to the effective Lipatov intercept $`\alpha _P=1+4\mathrm{log}2\overline{\alpha }N_C/\pi `$, and the cross-section normalisation. The obtained values of the parameters and the $`\chi ^2`$ of the fit are given in Table III for a fit to the H1 and ZEUS data separately, and then to the H1 + ZEUS data together. Note that one H1 point at $`k_T>5`$ GeV ($`7.3`$ 10<sup>-4</sup>), and four ZEUS points ($`x=4.`$ 10<sup>-4</sup>, and the three highest-x points), were not taken into account in the fit and are distinguished in Tables II with a star. We will discuss this selection in a little while.
| fit | $`\overline{\alpha }`$ | $`\alpha _P`$ | Norm. | $`\chi ^2(/dof)`$ |
| --- | --- | --- | --- | --- |
| H1 | 0.17 $`\pm `$ 0.02 $`\pm `$ 0.01 | 1.44 $`\pm `$ 0.05 $`\pm `$ 0.025 | 29.4 $`\pm `$ 4.8 $`\pm `$ 5.2 | 5.7 (/9) |
| ZEUS | 0.20 $`\pm `$ 0.02 $`\pm `$ 0.01 | 1.52 $`\pm `$ 0.05 $`\pm `$ 0.025 | 26.4 $`\pm `$ 3.9 $`\pm `$ 4.7 | 2.0 (/2) |
| H1+ZEUS | 0.16 $`\pm `$ 0.01 $`\pm `$ 0.01 | 1.43 $`\pm `$ 0.025 $`\pm `$ 0.025 | 30.7 $`\pm `$ 2.9 $`\pm `$ 3.5 | 12.0 (/13) |
Table III- Fit results
The $`\chi ^2`$ of the fits have been calculated using statistical error only and are very satisfactory (about $`0.6perpoint`$ for H1 data, and $`1.perpoint`$ for ZEUS data). We give both statistical and systematic errors for the fit parameters. The values of the Lipatov intercept are close to one another and compatible within errors for the H1 and ZEUS sets of data, and indicate a preferable medium value ($`\alpha _P=1.41.5`$). We also notice that the ZEUS data have the tendency to favour a higher exponent, but the number of data points used in the fit is much smaller than for H1, and the H1 data are also at lower $`x`$. The normalisation is also compatible between ZEUS and H1. The fit results are shown in Figure 1 and compared with the H1 and ZEUS measurements.
Let us discuss our selection criterium for the fits. Both lowest $`x`$ points for H1 and ZEUS show large correction factors but only the lowest $`x`$ point for ZEUS lies a bit above the prediction, which shows the relevance of the correction factors we determined. On the other hand, the three highest $`x`$ points for ZEUS cannot be described by a BFKL fit probably because the $`x`$-value is too high ($`x>10^2`$). Consider now the second lower $`x`$ point at $`k_T>5`$ GeV for the H1 experiment that we suppressed from the fit (see Table IIb). If we include it in the fit the $`\chi ^2`$ value goes from 5.7 to 32, which is due to the small statistical error of this data point (the systematic error is on the contrary very large). By comparison, including the lowest x point for ZEUS changes the $`\chi ^2`$ from 2.0 to 7.9. In the same way, including the highest x points still increases the $`\chi ^2`$ to 67.4, showing clearly that these highest x points cannot be described using the BFKL formalism. It is interesting to note that all similar discrepancies appear also in other types of fitting procedures, e.g. in Ref. .
5. Comparison with other processes
The final result of our new determination of the effective pomeron intercept is $`\alpha _P=1.43\pm 0.025`$ (stat.) $`\pm 0.025`$ (syst.). This high value of the intercept leads to the following remarks. Our analysis confirms the trend observed using DGLAP based Monte-Carlo which have difficulties to reproduce the forward jet cross-section due to a low effective pomeron intercept when both $`k_T^2`$ and $`Q^2`$ scales are of the same order.
On the other hand, our method allows a direct comparison of the intercept values with those obtained in other experimental processes, i.e. $`\gamma ^{}\gamma ^{}`$ cross-sections at LEP , jet-jet cross-sections at Tevatron at large rapidity intervals , $`F_2`$ and $`F_2^D`$ proton structure function measurements . Let us first consider the known determinations of the effective intercepts in $`F_2`$ and $`F_2^D`$ measurements at HERA . It is known that the effective intercept determined in these measurements is rather low<sup>§</sup><sup>§</sup>§It is interesting to note that the “hard” Pomeron intercept obtained within the framework of two-Pomeron models fits with our determination. However our parametrization (2) corresponds to only one Pomeron.(1.2-1.3). This is the reason why these data can be both described by a DGLAP or a BFKL-LO fit Note that in the BFKL descriptions of these data , the effective intercept is taken to be constant, while the $`Q^2`$ dependence comes from the BFKL integration (see for instance formula (2)).
Now let us consider processes initiated by two hard probes which allow a more direct comparison between experiments and BFKL predictions. These processes suppress DGLAP evolution by selecting events with comparable hard scales for both hard probes. Recent data on $`\gamma ^{}\gamma ^{}`$ cross-section measurements at LEP lead to a BFKL description with a low effective intercept compatible with the one of $`F_2`$ and $`F_2^D`$ at HERA ($`\alpha _P`$=1.2-1.3 ) The statistics for these data is still very low. L3 and OPAL Collaborations have released the cuts used to enhance BFKL effects to get more statistics . These data can be both described by BFKL and DGLAP evolution equations.. The fact that similar values of the intercepts are found could be interpreted by sizeable higher order corrections to BFKL equation. On the other hand, it is interesting to note that our result based on forward jet measurement at HERA obtained in comparable $`Q^2`$ ($`Q^210`$ GeV<sup>2</sup>) and rapidity ($`Y`$ 3-4) domains is quite different. The value of the intercept is significantly higher.
It is also fruitful to compare our results with the effective intercept we obtain from recent preliminary dijet data obtained by the D0 Collaboration at Tevatron . The measurement consists in the ratio $`R=\sigma _{1800}/\sigma _{630}`$ where $`\sigma `$ is the dijet cross-section at large rapidity interval $`Y\mathrm{\Delta }\eta `$ for two center-of-mass energies (630 and 1800 GeV), $`\mathrm{\Delta }\eta _{1800}=4.6`$, $`\mathrm{\Delta }\eta _{630}=2.4.`$ The experimental measurement is $`R=2.9\pm 0.3`$ (stat.) $`\pm 0.3`$ (syst.). Using the Mueller-Navelet formula , this measurement allows us to get a value of the effective intercept for this process <sup>\**</sup><sup>\**</sup>\**Formula (13) is obtained after integration over the jet tranverse energies at 630 and 1800 GeV, $`E_{T_1}`$, $`E_{T_2}`$. We note that a non integrated formula shows a sizeable dependence on $`E_{T_1}/E_{T_2}`$, which could be confronted with experiment .
$`R={\displaystyle \frac{_{\frac{1}{2}i\mathrm{}}^{\frac{1}{2}+i\mathrm{}}\frac{d\gamma }{2i\pi \gamma (1\gamma )}e^{ϵ(\gamma ,0)\mathrm{\Delta }\eta _{1800}}}{_{\frac{1}{2}i\mathrm{}}^{\frac{1}{2}+i\mathrm{}}\frac{d\gamma }{2i\pi \gamma (1\gamma )}e^{ϵ(\gamma ,0)\mathrm{\Delta }\eta _{630}}}}.`$ (13)
We get $`\alpha _P`$=1.65 $`\pm `$ 0.05 (stat.) $`\pm `$ 0.05 (syst.), in agreement with the value obtained by D0 using a saddle-point approximation . This intercept is higher than the one obtained in the forward jet study.
The question arises to interpret the different values of the effective intercept. It could reasonably come from the differences in higher order QCD corrections for the BFKL kernel and/or in the impact factors depending on the initial probes ($`\gamma ^{}`$ vs. jets). In order to evaluate the approximate size of the higher order BFKL corrections, we will use their description in terms of rapidity veto effects . In formula (2), we make the following replacement
$`\mathrm{exp}(ϵ(\gamma ,0)Y)\mathrm{\Sigma }_{n=0}^{\mathrm{}}\theta (Y(n+1)b){\displaystyle \frac{\left[ϵ(\gamma ,0)(Y(n+1)b)\right]^n}{\mathrm{\Gamma }(n+1)}}.`$ (14)
The Heaviside function $`\theta `$ ensures that a BFKL ladder of $`n`$ gluons occupies $`(n+1)b`$ rapidity interval where $`b`$ parametrises the strength of NLO BFKL corrections. The value of the leading order intercept is fixed to $`\alpha _p=1.75(\alpha _S(Q^2=10)=0.28)`$, where $`Q^2=10`$ GeV<sup>2</sup> is inside the average range of $`Q^2`$ in the forward jet measurement. The fitted value of the $`b`$ parameter obtained using the forward jet data is found to be 1.28 $`\pm `$ 0.08 (stat.) $`\pm `$ 0.02 (syst.). Imposing the same value of $`\alpha _P`$ with Tevatron data gives $`b`$=0.21 $`\pm `$ 0.11 (stat.) $`\pm `$ 0.11 (syst.). Note that the theoretical value of $`b`$ for the NLO BFKL kernel is expected to be of the order 2.4, which is also compatible with the result obtained for the $`\gamma ^{}\gamma ^{}`$ cross-section. A contribution from the NLO impact factors is not yet known, and could perhaps explain the different values of $`b`$.
6. Conclusions
To summarize our results, using a new method to disantangle the effects of the kinematic cuts from the genuine dynamical values of the forward jet cross-sections at HERA, we find that the effective pomeron intercept is $`\alpha _P=1.43\pm 0.025`$ (stat.) $`\pm 0.025`$ (syst.). It is much higher than the soft pomeron intercept, and, among those determined in hard processes, it is intermediate between $`\gamma ^{}\gamma ^{}`$ interactions at LEP and dijet productions with large rapidity intervals at Tevatron.
Looking for an interpretation of our results in terms of higher order BFKL corrections expressed by rapidity gap vetoes $`b`$ between emitted gluons, we find a value of $`b=`$1.3, which is sizeable but less than the theoretically predicted value for the NLO BFKL kernel ($`b=`$2.4). The observed dependence in the process deserves further more precise studies .
Last but not least, the derivation of the correction factors given in Table II is independent of the theoretical input and could be used to test any model suitable for the jet cross-section.
Acknowledgments
We would like to thank Lev Lipatov for his fruitful remarks and suggestions. One of us (J.G.C.) acknowledges supports by CONACyT. |
warning/0002/hep-th0002005.html | ar5iv | text | # Multifractality of time and space, covariant derivatives and gauge invariance
## I Introduction
1. In the Riemannian Geometry and in the standard theory of fundamental particles (see for example, , )the covariant derivatives are used, instead of usual derivatives. In the standard theory of fundamental particles the covariant derivatives has the form
$`D^\mu =^\mu ig_1{\displaystyle \frac{Y}{2}}B^\mu ig_2{\displaystyle \frac{\tau _i}{2}}W_i^\mu ig_3{\displaystyle \frac{\lambda _a}{2}}G_a^\mu `$ (1)
$`i=1,2,3,a=1,2,\mathrm{..8}`$ (2)
where $`B^\mu `$, $`W_i^\mu `$, $`G_a^\mu `$ are gauge-invariant vector fields with symmetry groups $`U(1)`$,$`SU(2)`$,$`SU(3)`$, $`\tau _i`$,$`\lambda _a`$ \- are matrices of isospin and color fields, $`Y`$ \- hypercharge, $`\mu =0,1,2,3`$. The covariant derivatives with respect to tensor $`t^\mu \nu `$, in space with a metric tensor $`\gamma ^\mu \nu `$, are defined by relations
$$D_\alpha t^{\mu \nu }=_\alpha t^{\mu \nu }+\gamma _{\alpha \beta }^\nu t^{\mu \beta }$$
(3)
where $`\gamma _{\alpha \beta }^\nu `$ are Kristoffel’s coefficients
$$\gamma _{\alpha \beta }^\nu =\frac{1}{2}\gamma ^{\mu \sigma }(_\alpha \gamma _{\beta \sigma }+_\beta \gamma _{\alpha \sigma }+_\sigma \gamma _{\alpha \beta })$$
(4)
In this and other cases to $`^\mu `$ adds vector, tensor etc. functional terms usual derivatives (in last cases $`^\mu `$ is multiplied on a unit matrix of the according type).
2. For describing of dynamic characteristics of systems defined on multifractal time and space sets it is necessary to define the functionals (left-sided and right-sided) determined on the functions, given on a multifractal sets. These functionals are the elementary generalization of the Riemann - Liouville fractional derivatives and integrals (about the Riemann-Liouville fractional derivatives see ) for functions defined on multifractal sets and were introduced in . These functionals (generalized fractional derivatives (GFD)) read (for differentiation with respect to time)
$$D_{+,t}^df(t)=\left(\frac{d}{dt}\right)^n_a^t\frac{f(t^{})dt^{}}{\mathrm{\Gamma }(nd(t^{}))(tt^{})^{d(t^{})n+1}}$$
(5)
$`D_{,t}^df(t)=`$ (6)
$`=`$ $`(1)^n\left({\displaystyle \frac{d}{dt}}\right)^n{\displaystyle _t^b}{\displaystyle \frac{f(t^{})dt^{}}{\mathrm{\Gamma }(nd(t^{}))(t^{}t)^{d(t^{})n+1}}}`$ (7)
where $`\mathrm{\Gamma }`$-is a gamma function, $`a<b`$, $`a`$ and $`b`$ stationary values selected on an infinite axis (from $`\mathrm{}`$ to $`\mathrm{}`$), $`n1d_t<n`$, $`n=\{d_t\}+1`$, $`\{d_t\}`$\- integer part of $`d_t0`$ and $`n=0`$ for $`d_t<0`$. Generalized fractional derivatives (GFD) (5)-(6) coincide with fractional derivatives or fractional integrals of the Riemann - Liouville for a case $`d_t=const`$. At $`d_t=n+\epsilon (t)`$, $`\epsilon (t)0`$ GFD are represented through usual derivatives and integrals . The functions in integrals in (5)-(6) are generalized functions, given on set of a finitary functions . Similar (5)-(6) definitions GFD can be defined and with respect to space variables $`𝐫`$. The integral functionals (5)-(6) allow to describe dynamic of functions defined on multifractal sets and replace for such functions orderly or fractional (Riemann-Liouville) differentiation and integration. These functionals partially describe the memory about past (or future, if use right-hand GFD) time or spatial events. In it is shown, for $`d_\alpha =1\epsilon _\alpha (𝐫(t),t)`$, $`\alpha =t,𝐫`$, $`|\epsilon |1`$ where
$$\epsilon _\alpha =\underset{i}{}\beta _{\alpha ,i}L_{\alpha ,i}(\mathrm{\Phi }_i)$$
(8)
($`L_{t,i}`$ -are the Lagrangians densities of energy of physical fields which presents in a point $`x`$, $`\beta _{t,i}`$ are a numerical dimension factor, $`L_{𝐫,i}`$\- are the Lagrangians densities of energy of the new fields which appears be-cause of fractionality of space dimensions ) that GFD can be represented in that case by usual derivatives (scalar field):
$$D_{+,x^\mu }^{1+\epsilon _\alpha (x)}f(x)=^\mu f(x)+a^\mu [\epsilon _\alpha (x)f(x)]$$
(9)
where $`a`$ -is a factor of the regularization of integrals (5)-(6) or an analogies integrals for case of differentiation with respect to space coordinates. The relations (8), defining GFD for sets with almost integer dimension, have more composite structure than (1) and (3), nevertheless, they are very similar to definitions of the covariant derivatives. For a gravitational field the connections GFD (for the component $`\mathrm{\Phi }_00`$ of a gravity’s potentials) with covariant derivatives of effective Riemannian space was obtained in . The purpose of this paper is the establishment of the connections between GFD defined by (5)-(8) and covariant derivatives (1)-(3), for taking into account only the feeble physical fields and multifractal nature of time and using the relations for covariant derivatives (5)-(8) which are more complicate, than derivatives defined by relations (1) - (3).
## II Covariant derivatives and multifractal time
1. We shall be restricted by consideration of connection GFD and covariant derivatives defined on the multifractal time sets (i.e. $`\alpha =t`$ , for multifractal 3-dimensions coordinates space the consideration may be carried out by similarly methods) for a case of the small fractional corrections to integer dimension of $`t(|\epsilon |1)`$. From (9) follows
$$D_{+,x^\mu }^{1+\epsilon (x)}f(x)=[(1+\epsilon )^\mu ^\mu \epsilon (x)]f(x)$$
(10)
or, with the account of (8)
$`D_{+,x^\mu }^{1+\epsilon (x)}f(x)`$ $`=`$ $`[1{\displaystyle \underset{i}{}}\beta _iL_i(\mathrm{\Phi }_i(𝐫(t),t))]^\mu f(x)`$ (11)
$``$ $`{\displaystyle \underset{i}{}}\beta _i^\mu [L_i(\mathrm{\Phi }_i(𝐫(t),t))f(x)]`$ (12)
The relations (11) defines, for $`|\epsilon |1`$, connections between GFD (in the set of multifractal time with dimension $`d=1+\epsilon `$) and orderly derivatives with respect to time and coordinates in the time’s set with dimension $`d=1`$. Note that from the point of view of the multifractal theory, it is possible to treat the time’s set with dimension $`d=1`$, in which derivatives with respect to time and coordinates are replaced on covariant derivatives of the (11), as the effective time’s set corresponding to multifractal time with dimension $`d=1+\epsilon `$.
2. The comparison (11) and (1) allows to establish difference of general structure of derivatives (11) from covariant derivatives (1): except for presence more composite, than in (1), additive terms proportional derivatives of Lagrangians density with respect to time and coordinates, there are renormalization of the function’s factors before derivatives $`^\mu `$, which quantity depends both on time and from coordinates. The presence this renormalization allows to introduce, instead of space of time with fractional dimension for tensor fields an effective Riemannian space with the metric defined by dependence of $`\epsilon `$ from a metric tensor (see special case in ).
3. What is the physical sense of replacement of usual derivatives on the covariant derivatives (11)? The derivatives of the Lagrangians density with respect to time and coordinates can be interpreted as a birth or disappearance of energy (the signs of derivatives are plus or are minus, accordingly) if the time or the coordinates are changing. This energy is transmitted to a field of time or is taken from it by the carrier of a measure (by the set $`R^n`$). The conservation laws are fulfilled for closed system consisting of material fields of the time and the space and set of the carrier of a measure $`R^n`$. The fields of time and space without the carrier of measure set are open systems. The birth or annihilation of energy of the field of time is accompanied by changes of before derivatives quantity (factor before derivatives $`^\mu `$) depending not only from characteristics of function $`f(t)`$, with respect to which the operation of differentiation is applied, but also from characteristics of the field of time (defined by Lagrangians density in the given instant and given coordinates). The comparison (11) and (1) allows to state if the selection of a Lagrangians in the form of standard model of the theory of fundamental particles (with replacement in a Lagrangians the usual derivatives by covariant derivatives (1)) is made, that the covariant derivatives (11) contain the description of process of birth or disappearance all physical fields ( for example electromagnetic fields or fields of fundamental particles). The intensity of these processes depends on density of energy and, thus (11) are not reduced to (1) at the any selections of Lagrangians of standard model. The comparison (11) with the covariant derivatives general theory of relativity (GTR) (3) allows to state, as the components of a metric tensor are functions of an energy-momentum tensor of gravitational field $`t_{\mu \nu }`$, that covariant derivatives (11) contains (3) as a special case (according to presence only one gravitational field). This case was considered in and it was demonstrated the opportunity of introduction of effective Riemannian space (according to Riemannian space of GFR) for approximate describing of gravitational phenomenon in multi-fractal time (if the time is multifractal in reality) at small $`\epsilon `$.
## III Gauge invariance of the fields $`_\mu \underset{i}{}\beta _iL_i`$
Are the fields $`^\mu \underset{i}{}\beta _iL_i`$ gouge invariant? Whether is it necessary to require the gauge invariance of GFD (for small $`\epsilon `$) from the equations of theoretical physics which are wrote down with the help of GFD for multifractal sets, where time and space have fractional dimensions almost undistinguished from the integer dimensions ( these equations are wrote down in )? In the last case if to follow fields quantum theory, the requirement of the gauge invariance introduce the new ”charges”, that is the charges, defining birth or disappearance of fields $`L_i`$. This charge, is defined by both factors $`\beta _i`$ and factors included in Lagrangians (that define already known charges). The requirement of the gouge invariance enters as well the new massless physical fields: the field’s of production of all known physical fields $`\phi `$:
$$\phi ^\mu (𝐫(t),t),L_i)=^\mu \beta _iL_i$$
(13)
As the relations (11) for GFD are approximate and are valid only for small $`\epsilon `$, the gauge invariance (13) is also approximate and it is meaningful (if it at all is meaningful) only for small $`\epsilon `$ (though this case is most spread). The problem of validity of introduction of a gauge invariance (invariance of derivatives densities of Lagrangians) remains open.
## IV Is the fractal dimension of time influenced by characteristics of multifractal sets that arise in the internal sets of ”time intervals”?
The form of the covariant derivatives in the standard theories of fundamental particles (1), allow (within the framework of representations of multifractal time) to put the problem: is it possible to construct dimension of time $`d_t`$ in such a manner that covariant derivatives (1) will be appears in the theory by explicit form, instead of appears through Lagrangians? As will be shown below, it is possible, being non-essential for $`d_t`$. It is possible presents the covariant derivatives (11) in the form more similar to covariant derivatives (1), including the relations (1) as a special case. Till now ”time intervals” from which, on definition, the material field of time is consists, were treated as ”points” of time sets. It was considered possible to neglect by interior structure of time sets component those ”time intervals”. The covariant derivatives (1) describe characteristics of interior symmetry of fundamental particles and their gauge invariance, so, apparently, it is convenient in multifractal model investigate the characteristics of inner sets that consist the time and the space ”points”. So let’s add to density of Lagrangians (in $`\epsilon `$ that defined the multifractal dimensions of time and space sets) the terms, which origin can be connected with characteristics of sets constructing the ”points” by the vicinity of points $`x`$. These time and space sets ( intervals) near $`x`$ earlier approximately were described as the ”points” and characterized by fractional dimension $`d_t`$ or $`d_r`$, but that part of sets contents and carry the information about structural characteristics of fundamental particles (as the field of time, following , generates all material fields and defines their characteristics). If the sets components of ”time intervals” are multifractal, than to each point this ”interior” sets surrounding the medial ”point” with coordinate $`x`$ should be compared with the fractal dimensions or with their medial integral description. Most simply in this case for $`d=1+\epsilon (𝐫(t),t)`$ to write down $`d`$ as
$$d=1+\epsilon =1+\underset{i}{}\beta _iL_{0,i}(\mathrm{\Phi }_i)+\underset{i,\mu }{}\underset{x^\mu }{\overset{x_0^\mu }{}}\stackrel{~}{B}_i^\mu (x)𝑑x^\mu $$
(14)
where $`\stackrel{~}{B}^\mu `$ \- are vector quantifies (possessing by complicated interior symmetry) and define characteristics of sets in a vicinity of each point $`x`$ of the time (or space) intervals, in which this point is considered. In particular, for definition of dependence $`\stackrel{~}{B}^\mu `$ from physical fields and interior symmetries of fundamental particles, for example, relations (1) can be chosen. In (14) $`x_0^\mu `$ \- are stationary values proportional to ”size” of the according time intervals (or space intervals), $`L_{0,i}`$ -are densities of Lagrangians of free fields. As $`x_0^\mu `$ are very small, their contributions are essential only for derivatives $`D_{\pm ,x^\mu }^{1+\epsilon }`$, which in this case will accept the form
$`D_{+,x^\mu }^{1+\epsilon (x)}[1{\displaystyle \underset{i}{}}\beta _iL_{0,i}(\mathrm{\Phi }_i(𝐫(t),t))]^\mu `$ (15)
$``$ $`{\displaystyle \underset{i}{}}\beta _i^\mu L_{0,i}(\mathrm{\Phi }_i(𝐫(t),t)){\displaystyle \underset{i}{}}B_i^\mu (𝐫(t),t)`$ (16)
It is possible to rewrite covariant derivatives (15) as
$`D_{+,x^\mu }^{1+\epsilon (x)}[1{\displaystyle \underset{i}{}}\beta _iL_{0,i}(\mathrm{\Phi }_i(𝐫(t),t))]\times `$ (17)
$`\times `$ $`^\mu {\displaystyle \frac{\underset{i}{}[\beta _i^\mu L_{0,i}(\mathrm{\Phi }_i(𝐫(t),t))B_i^\mu (𝐫(t),t)]}{[1\underset{i}{}\beta _iL_{0,i}(\mathrm{\Phi }_i(𝐫(t),t))]}}`$ (18)
If $`B_i^\mu `$ in (15) determine by use of relations (1) and if neglect by the contributions from fractal dimensions in square brackets ($`\beta 0`$), the definition covariant derivatives (15) coincides with the covariant derivatives of standard theory of fundamental particles (at the according selection of stationary values and matrices). It is necessary to note, apparently, the equivalence of mathematical expositions of dynamic characteristics of fundamental particles in viewed model as with the help of use of densities of Lagrangians (into which the terms describing interactions of fields and particles are included), and exposition with the help of the introduction of covariant derivatives. In the last case and elimination, from the according densities of Lagrangians of terms describing interactions must be made.
## V Conclusions
In the presented theory of multifractal time and space the generalized fractional derivatives, defined by (5) - (6) are used instead of usual differentiation and integration for describing a dynamic characteristics of any physical objects (fields, particles etc. The using of integral functionals gives in considerable thickening of the mathematical tool and change a great deal representations about the physical nature of time and space. Only in the case when the fractal corrections to FD are small, it is possible to present GFD with the help of usual derivatives and integrals as was shown in and used in this paper. In this case the fractionality of dimensions of time (or, similarly, dimensions of space, see ) can be presented by introduction of an ”effective” time and space, in which derivatives (and the integrals, for a case of negative fractional dimensions) are substituted by ”covariant” derivatives and integrals. The consideration of concrete models and Lagrangians, in particular, the Lagrangians of the standard theory of fundamental particles, has been illustrated an opportunity of exposition of gouge invariant fields with covariant derivatives of a new type (15). This covariant derivatives are contains, except for known terms of the gouge invariant standard theory of fundamental particles and the terms including of effective Riemannian space, the terms describing the continuous change of densities of energy of physical fields (their birth or disappearance). Multifractal set of time used in the present paper (also, as well as the mathematical tool GFD, used in a series of other papers \- ) creates peculiar model of the world considered as the open system (see ,), in which there are no invariable states. In viewed model of multifractal time all physical fields are not stationary (note that GFD with respect to time or coordinates of stationary values are not equal to zero), but also their energy continuously changing thow potentials are invariable (”non-potential” change of energy as the result of an exchange of energy with the set of the carrier of a measure, because of the presence at the covariant derivatives the terms with derivatives from densities of Lagrangians. At least, for the case $`|\epsilon |1`$, the theory of fundamental particles defined on sets of multifractal time and space, allows to take into account by the uniform mathematical method the influences of all known physical fields (the fundamental particles, used by the theory, and gravitational field). It is achieved by introduction the new paradigm : the exposition of time and space as multifractal sets with fractional dimensions and introduction, in this connection, mathematical methods of generalized fractional derivatives and integrals. For deriving the modified equations of the standard theory of fundamental particles or various variants of the great unification theory (or any physical theory) it is enough to describe the fields taking into account the multifractal nature of time and space, i.e. replace covariant derivatives such as (1) (or an orderly derivatives in a scalar theories) by generalized covariant derivatives (GFD) (15) or, depending on selection of models, by similar though more composite covariant derivatives for fields with more complicated mathematical nature. The problems of correspondence obtained by such replacement models of the theories of fundamental particles to characteristics of the real world (if use model of multifractal time and spaces presented in ) here are not considered. All equations of theoretical physics wrote down with the help the GFD for the case $`\epsilon =0`$ can be considered as special cases of equations of presented model of multifractal time and space . At last we pay attention on the appearance of the new characteristics in modified thus theories (which were not explored yet) stipulated by the additional factors, and by the new terms in the modified covariant derivatives. |
warning/0002/astro-ph0002188.html | ar5iv | text | # I. Introduction
## I. Introduction
There are many fascinating issues associated with eternal inflation, which will be the main subject of this talk. You have certainly heard other people talk about eternal inflation, but I feel that the topic is important enough so that you should hear about it in some accent other than Russian. I will begin by summarizing the basics of inflation, including a discussion of how inflation works, and why many of us believe that our universe almost certainly evolved through some form of inflation. This material is certainly not new, but I think it is an appropriate introduction to any volume that focuses on inflationary cosmology. Then I will move on to discuss eternal inflation, first explaining how it works. I will then argue the eternal inflation has important implications, and raises important questions, which should not be dismissed as being merely metaphysical.
## II. How Does Inflation Work?
In this section I will review the basics of how inflation works, focusing on the earliest working forms of inflation—new inflation and chaotic inflation . While more complicated possibilities (e.g. hybrid inflation and supernatural inflation ) appear very plausible, the basic scenarios of new and chaotic inflation will be sufficient to illustrate the physical effects that I want to discuss in this article.
The key property of the laws of physics that makes inflation possible is the existence of states of matter that have a high energy density which cannot be rapidly lowered. In the original version of the inflationary theory , the proposed state was a scalar field in a local minimum of its potential energy function.<sup>*</sup><sup>*</sup>*A similar proposal was advanced by Starobinsky , in which the high energy density state was achieved by curved space corrections to the energy-momentum tensor of a scalar field. Such a state is called a false vacuum, since the state temporarily acts as if it were the state of lowest possible energy density. Classically this state would be completely stable, because there would be no energy available to allow the scalar field to cross the potential energy barrier that separates it from states of lower energy. Quantum mechanically, however, the state would decay by tunneling . Initially it was hoped that this tunneling process could successfully end inflation, but it was soon found that the randomness of false vacuum decay would produce catastrophically large inhomogeneities .
This “graceful exit” problem was solved by the invention of the new inflationary universe model , which achieved all the successes that had been hoped for in the context of the original version. In this theory inflation is driven by a scalar field perched on a plateau of the potential energy diagram, as shown in Fig. 1. Such a scalar field is generically called the inflaton. If the plateau is flat enough, such a state can be stable enough for successful inflation. Soon afterwards Linde showed that the inflaton potential need not have either a local minimum or a gentle plateau; in chaotic inflation , the inflaton potential can be as simple as
$$V(\varphi )=\frac{1}{2}m^2\varphi ^2,$$
(2.1)
provided that $`\varphi `$ begins at such a large value that it takes a long time for it to relax. For simplicity of language, I will stretch the meaning of the phrase “false vacuum” to include all of these cases; that is, I will use the phrase to denote any state with a high energy density that cannot be rapidly decreased. While inflation was originally developed in the context of grand unified theories, the only real requirement on the particle physics is the existence of a false vacuum state.
The New Inflationary Scenario:
Suppose that the energy density of a state is approximately equal to a constant value $`\rho _f`$. Then, if a region filled with this state of matter expanded by an amount $`dV`$, its energy would have to increase by
$$dU=\rho _fdV.$$
(2.2)
Something would have to supply that energy. Work would have to be done to cause the region to expand, which implies that the region has a negative pressure, which pulls back against whatever is causing the region to expand. The work done by this negative pressure $`p_f`$ is given by the elementary formula
$$dW=p_fdV.$$
(2.3)
Equating the work with the change in energy, one finds
$$p_f=\rho _f.$$
(2.4)
It is this negative pressure which is the driving force behind inflation. When one puts this negative pressure into Einstein’s equations to find out its gravitational effect, one finds that it leads to a repulsion, causing such a region to undergo exponential expansion. If the region can be approximated as isotropic and homogeneous, this result can be seen from the standard Friedmann-Robertson-Walker (FRW) equations:
$$\frac{d^2a}{dt^2}=\frac{4\pi }{3}G(\rho +3p)a=\frac{8\pi }{3}G\rho _fa.$$
(2.5)
where $`a(t)`$ is the scale factor, $`G`$ is Newton’s constant, and we adopt units for which $`\mathrm{}=c=1`$. For late times the growing solution to this equation has the form
$$a(t)e^{\chi t},\text{ where }\chi =\sqrt{\frac{8\pi }{3}G\rho _f}.$$
(2.6)
Of course inflationary theorists prefer not to assume that the universe began homogeneously and isotropically, but there is considerable evidence for the “cosmological no-hair conjecture” , which implies that a wide class of initial states will approach this exponentially expanding solution.
So the basic scenario of new inflation begins by assuming that at least some patch of the early universe was in this peculiar false vacuum state. In the original papers this initial condition was motivated by the fact that, in many quantum field theories, the false vacuum resulted naturally from the supercooling of an initially hot state in thermal equilibrium. It was soon found, however, that quantum fluctuations in the rolling inflaton field give rise to density perturbations in the universe , and that these density perturbations would be much larger than observed unless the inflaton field is very weakly coupled. For such weak coupling there would be no time for an initially nonthermal state to reach thermal equilibrium. Nonetheless, since thermal equilibrium describes a probability distribution in which all states of a given energy are weighted equally, the fact that thermal equilibrium leads to a false vacuum implies that false vacuum-like states are not uncommon. Thus, even in the absence of thermal equilibrium, even if the universe started in a highly chaotic initial state, it seems reasonable to simply assume that some small patches of the early universe settled into the false vacuum state, as was suggested for example in . The idea that one should consider small patches of the early universe with arbitrary initial configurations of scalar fields was later emphasized by Linde in the context of chaotic inflation. Linde pointed out that even highly improbable initial patches could be important if they inflated, since the exponential expansion could still cause such patches to dominate the volume of the universe. One might hope that eventually a full theory of quantum origins would allow us to calculate the probability of regions settling into the false vacuum, but I will argue in Sec. V that, in the context of eternal inflation, this probability is quite irrelevant.
Once a region of false vacuum materializes, the physics of the subsequent evolution seems rather clear-cut. The gravitational repulsion caused by the negative pressure will drive the region into a period of exponential expansion. If the energy density of the false vacuum is at the grand unified theory scale ($`\rho _f(2\times 10^{16}\text{GeV})^4)`$, Eq. (2.6) shows that the time constant $`\chi ^1`$ of the exponential expansion would be about $`10^{38}`$ sec. For inflation to achieve its goals, this patch has to expand exponentially for at least 60 e-foldings. Then, because this state is only metastable—the inflaton field is perched on top of the hill of the potential energy diagram of Fig. 1—eventually this state will decay. The inflaton field will roll off the hill, ending inflation. And when it does, the energy density that has been locked in the inflaton field is released. Because of the coupling of the inflaton to other fields, that energy becomes thermalized to produce a hot soup of particles, which is exactly what had always been taken as the starting point of the standard big bang theory before inflation was introduced. From here on the scenario joins onto the standard big bang description. The role of inflation is to replace the postulates of the standard big bang theory with dynamically generated initial conditions.
The inflationary mechanism produces an entire universe starting from essentially nothing, so one needs to answer the question of where the energy of the universe came from. The answer is that it came from the gravitational field. I am not saying that the colossal energy of the universe was stored from the beginning in the gravitational field. Rather, the crucial point is that the energy density of the gravitational field is literally negative—a statement which is true both in Newtonian gravity and in general relativity. So, as more and more positive energy materialized in the form of an ever-growing region filled with a high-energy scalar field, more and more negative energy materialized in the form of an expanding region filled with a gravitational field. So the total energy remained very small, and could in fact be exactly zero. There is nothing known that places any limit on the amount of inflation that can occur while the total energy remains exactly zero.In Newtonian mechanics the energy density of a gravitational field is unambiguously negative; it can be derived by the same methods used for the Coulomb field, but the force law has the opposite sign. In general relativity there is no coordinate-invariant way of expressing the energy in a space that is not asymptotically flat, so many experts prefer to say that the total energy is undefined. Either way, there is agreement that inflation is consistent with the general relativistic description of energy conservation.
Chaotic Inflation:
Chaotic inflation can occur in the context of a much more general class of potential energy functions. In particular, even a potential energy function as simple as Eq. (2.1), describing a scalar field with a mass and no interaction, is sufficient to describe chaotic inflation. Chaotic inflation is illustrated in Fig. 2. In this case there is no state that bears any obvious resemblance to the false vacuum of new inflation. Instead the scenario works by supposing that chaotic conditions in the early universe produced one or more patches in which the inflaton field $`\varphi `$ was at some high value $`\varphi =\varphi _0`$ on the potential energy curve. Inflation occurs as the inflaton field rolls down the hill. As long as the initial value $`\varphi _0`$ is sufficiently high on the curve, there will be sufficient inflation to solve all the problems that inflation is intended to solve.
The equations describing chaotic inflation can be written simply, provided that we assume that the universe is already flat enough so that we do not need to include a curvature term. The field equation for the inflaton field in the expanding universe is
$$\ddot{\varphi }+3H\dot{\varphi }=\frac{V}{\varphi },$$
(2.7)
where the overdot denotes a derivative with respect to time $`t`$, and $`H`$ is the time-dependent Hubble parameter given by
$$H^2=\frac{8\pi }{3}GV.$$
(2.8)
For the toy-model potential energy of Eq. (2.1), these equations have a very simple solution:
$$\varphi =\varphi _0\frac{m}{\sqrt{12\pi G}}t.$$
(2.9)
One can then calculate the number $`N`$ of inflationary e-foldings, which is given by
$$N=_{\varphi =\varphi _0}^{\varphi =0}H(t)𝑑t=2\pi G\varphi _0^2.$$
(2.10)
In this free-field model $`N`$ depends only on $`\varphi _0`$ and not on the inflaton mass $`m`$. Thus the number of e-foldings will exceed 60 provided that
$$\varphi _0>\sqrt{\frac{60}{2\pi }}M_\mathrm{P}3.1M_\mathrm{P},$$
(2.11)
where $`M_\mathrm{P}1/\sqrt{G}=1.22\times 10^{19}`$ GeV is the Planck mass. Although the value of the scalar field is larger than $`M_\mathrm{P}`$, the energy density can be low compared to the Planck scale:
$$\rho _0=\frac{1}{2}m^2\varphi _0^2>\frac{60}{4\pi }M_\mathrm{P}^2m^2.$$
(2.12)
For example, if $`m=10^{16}`$ GeV, then the potential energy density is only $`3\times 10^6M_\mathrm{P}^4`$. Since it is presumably the energy density that is relevant to gravity, one does not expect this situation to lead to strong quantum gravity effects.
## III. Evidence for Inflation
The arguments in favor of inflation are pretty much the same no matter which form of inflation we are discussing. In my opinion, the evidence that our universe is the result of some form of inflation is very solid. Since the term inflation encompasses a wide range of detailed theories, it is hard to imagine any alternative. Let me review the basic arguments.
1) The universe is big
First of all, we know that the universe is incredibly large. The visible part of it contains about $`10^{90}`$ particles. It is easy, however, to take this fact for granted: of course the universe is big, it’s the whole universe! In “standard” Friedmann-Robertson-Walker cosmology, without inflation, one simply postulates that about $`10^{90}`$ or more particles were here from the start. If, however, we try to imagine a theory describing the origin of the universe, it would have to somehow output this number of $`10^{90}`$ or more. That is a very big number, and it is hard to imagine it ever coming out of a calculation in which the input consists only geometrical quantities, quantities associated with simple dynamics, and factors of 2 and $`\pi `$. In the inflationary model, the huge number of particles is explained naturally by the exponential expansion, which reduces the problem to explaining 60 or 70 e-foldings of inflation. In fact, it is easy to construct underlying particle theories that will give far more than 70 e-foldings, suggesting that the observed universe is only a tiny speck within the universe as a whole.
2) The Hubble expansion
The Hubble expansion is also easy to take for granted, since it is so familiar. In standard FRW cosmology, the Hubble expansion is part of the list of postulates that define the initial conditions. But inflation offers an explanation of how the Hubble expansion began. The repulsive gravity associated with the false vacuum is exactly the kind of force needed to propel the universe into a pattern of motion in which any two particles are moving apart with a velocity proportional to their separation.
3) Homogeneity and isotropy
The degree of uniformity in the universe is startling. Through careful measurements of the cosmic background radiation, we know that the intensity of this radiation is the same in all directions to an accuracy of 1 part in 100,000. To get some feeling for how high this precision is, we can imagine a marble that is spherical to this accuracy. The surface of the marble would have to be shaped with a tolerance of about 1,000 angstroms, a quarter of the wavelength of light.
Although precision lenses can be ground to quarter-wavelength accuracy, we would nonetheless be shocked if we ever dug up a stone from the ground that was round to this extraordinary accuracy. If such a stone were somehow found, I am confident that we would not accept an explanation of its origin which simply proposed that the stone started out perfectly round. Similarly, in the current era, I do not think it makes sense to consider any theory of cosmogenesis that cannot offer some explanation of how the universe became so incredibly isotropic.
The uniformity of the cosmic background radiation implies that the observed universe had become uniform in temperature by about 300,000 years after the big bang, when the universe cooled enough so that the opaque plasma neutralized into a transparent gas. In standard FRW cosmology, the uniformity could be established by this time only if signals could propagate 100 times faster than light, which is not possible. In inflationary cosmology, however, the uniformity can be created initially on microscopic scales, by normal thermal-equilibrium processes. Then inflation takes over and stretches the regions of uniformity to become large enough to encompass the observed universe.
4) The flatness problem
I find the flatness problem particularly impressive, because the numbers that it leads to are so extraordinary. The problem concerns the value of the ratio
$$\mathrm{\Omega }_{\mathrm{tot}}\frac{\rho _{\mathrm{tot}}}{\rho _c},$$
(3.1)
where $`\rho _{\mathrm{tot}}`$ is the average total mass density of the universe and $`\rho _c=3H^2/8\pi G`$ is the critical density, the density that would make the universe spatially flat. ($`\rho _{\mathrm{tot}}`$ includes any vacuum energy $`\rho _{\mathrm{vac}}=\mathrm{\Lambda }/8\pi G`$ associated with the cosmological constant $`\mathrm{\Lambda }`$, if it is nonzero.)
The present value of $`\mathrm{\Omega }_{\mathrm{tot}}`$ satisfies
0.1
<
Ωtot,0
<
2,
<
0.1subscriptΩtot0
<
20.1\mathrel{\vbox{\vbox{\offinterlineskip\hbox{$<$}
\vskip 0.2pt\hbox{$\sim$}}}}\Omega_{{\rm tot},0}\mathrel{\vbox{\vbox{\offinterlineskip\hbox{$<$}
\vskip 0.2pt\hbox{$\sim$}}}}2\ , (3.2)
but the precise value is not known. Despite the breadth of this range, the value of $`\mathrm{\Omega }_{\mathrm{tot}}`$ at early times is highly constrained, since $`\mathrm{\Omega }_{\mathrm{tot}}=1`$ is an unstable equilibrium point of the standard model evolution. If $`\mathrm{\Omega }_{\mathrm{tot}}`$ was ever exactly equal to one, it would remain so forever. However, if $`\mathrm{\Omega }_{\mathrm{tot}}`$ differed slightly from 1 in the early universe, that difference—whether positive or negative—would be amplified with time. In particular, the FRW equations imply that $`\mathrm{\Omega }_{\mathrm{tot}}1`$ grows as
$$\mathrm{\Omega }_{\mathrm{tot}}1\{\begin{array}{cc}t\hfill & \text{(during the radiation-dominated era)}\hfill \\ t^{2/3}\hfill & \text{(during the matter-dominated era) .}\hfill \end{array}$$
(3.3)
At $`t=1`$ sec, for example, Dicke and Peebles pointed out that $`\mathrm{\Omega }_{\mathrm{tot}}`$ must have equaled one to an accuracy of one part in $`10^{15}`$. Classical cosmology provides no explanation for this fact—it is simply assumed as part of the initial conditions. In the context of modern particle theory, where we try to push things all the way back to the Planck time, $`10^{43}`$ sec, the problem becomes even more extreme. At this time $`\mathrm{\Omega }_{\mathrm{tot}}`$ must have equaled one to 58 decimal places!
While this extraordinary flatness of the early universe has no explanation in classical FRW cosmology, it is a natural prediction for inflationary cosmology. During the inflationary period, instead of $`\mathrm{\Omega }_{\mathrm{tot}}`$ being driven away from 1 as described by Eq. (3.3), $`\mathrm{\Omega }_{\mathrm{tot}}`$ is driven towards 1, with exponential swiftness:
$$\mathrm{\Omega }_{\mathrm{tot}}1e^{2H_{\mathrm{inf}}t},$$
(3.4)
where $`H_{\mathrm{inf}}`$ is the Hubble parameter during inflation. Thus, as long as there is enough inflation, $`\mathrm{\Omega }_{\mathrm{tot}}`$ can start at almost any value, and it will be driven to unity by the exponential expansion.
5) Absence of magnetic monopoles
All grand unified theories predict that there should be, in the spectrum of possible particles, extremely massive magnetic monopoles. By combining grand unified theories with classical cosmology without inflation, Preskill found that magnetic monopoles would be produced so copiously that they would outweigh everything else in the universe by a factor of about $`10^{12}`$. A mass density this large would cause the inferred age of the universe to drop to about 30,000 years! In inflationary models, the monopoles can be eliminated simply by arranging the parameters so that inflation takes place after (or during) monopole production, so the monopole density is diluted to a completely negligible level.
6) Anisotropy of the cosmic background radiation
The process of inflation smooths the universe essentially completely, but quantum fluctuations of the inflaton field can generate density fluctuations as inflation ends. Generically these are adiabatic Gaussian fluctuations with a nearly scale-invariant spectrum . New data is arriving quickly, but so far the observations are in excellent agreement with the predictions of the simplest inflationary models. For a review, see for example Bond and Jaffe , who find that the combined data give a slope of the primordial power spectrum within 5% of the preferred scale-invariant value.
## IV. Eternal Inflation: Mechanisms
Having discussed the mechanisms and the motivation for inflation itself, I now wish to move on the main issue that I want to stress in this article—eternal inflation, the questions that it can answer, and the questions that it raises. In this section I will discuss the mechanisms that make eternal inflation possible, leaving the other issues for the following sections. I will discuss eternal inflation first in the context of new inflation, and then in the context of chaotic inflation, where it is more subtle.
Eternal New Inflation:
The eternal nature of new inflation was first discovered by Steinhardt and Vilenkin in 1983. Although the false vacuum is a metastable state, the decay of the false vacuum is an exponential process, very much like the decay of any radioactive or unstable substance. The probability of finding the inflaton field at the top of the plateau in its potential energy diagram does not fall sharply to zero, but instead trails off exponentially with time . However, unlike a normal radioactive substance such as radium, the false vacuum exponentially expands at the same time that it decays. In fact, in any successful inflationary model the rate of exponential expansion is always much faster than the rate of exponential decay. Therefore, even though the false vacuum is decaying, it never disappears, and in fact the total volume of the false vacuum, once inflation starts, continues to grow exponentially with time, ad infinitum.
Fig. 3 shows a schematic diagram of an eternally inflating universe. The top bar indicates a region of false vacuum. The evolution of this region is shown by the successive bars moving downward, except that I could not show the expansion and still fit all the bars on the page. So the region is shown as having a fixed size in comoving coordinates, while the scale factor, which is not shown, increases from each bar to the next. As a concrete example, suppose that the scale factor for each bar is three times larger than for the previous bar. If we follow the region of false vacuum indicated by the top bar as it evolves into the second bar, in about one third of the region the scalar field rolls down the hill of the potential energy diagram, precipitating a local big bang that will evolve into something that will eventually appear to its inhabitants as a universe. This local big bang region is shown in gray and labeled “Universe.” Meanwhile, however, the space has expanded so much that each of the two remaining regions of false vacuum is the same size as the starting region. Thus, if we follow the region for another time interval of the same duration, each of these regions of false vacuum will break up, with about one third of each evolving into a local universe, as shown on the third bar from the top. Now there are four remaining regions of false vacuum, and again each is as large as the starting region. This process will repeat itself literally forever, producing a kind of a fractal structure to the universe, resulting in an infinite number of the local universes shown in gray. These local universes are often called bubble universes, but that terminology conveys the unfortunate connotation that the local universes are spherical. While bubbles formed in first-order phase transitions are round , the local universes formed in eternal new inflation are generally very irregular, as can be seen for example in the two-dimensional simulation by Vanchurin, Vilenkin, and Winitzki in Fig. 2 of Ref. . I therefore prefer to call them pocket universes, to try to avoid the suggestion that they are round.
The diagram in Fig. 3 is of course an idealization. The real universe is three dimensional, while the diagram illustrates a schematic one-dimensional universe. It is also important that the decay of the false vacuum is really a random process, while I constructed the diagram to show a very systematic decay, because it is easier to draw and to think about. When these inaccuracies are corrected, we are still left with a scenario in which inflation leads asymptotically to a fractal structure in which the universe as a whole is populated by pocket universes on arbitrarily small comoving scales. Of course this fractal structure is entirely on distance scales much too large to be observed, so we cannot expect astronomers to actually find it. Nonetheless, one does have to think about the fractal structure if one wants to understand the very large scale structure of the spacetime produced by inflation.
Most important of all is the simple statement that once inflation begins, it produces not just one universe, but an infinite number of universes.
Eternal Chaotic Inflation:
The eternal nature of new inflation depends crucially on the scalar field lingering at the top of the plateau of Fig. 1. Since the potential function for chaotic inflation, Fig. 2, has no plateau, it does not seem likely that eternal inflation can happen in this context. Nonetheless, Andrei Linde showed in 1986 that chaotic inflation can also be eternal.
The key to eternal chaotic inflation is the role of quantum fluctuations, which is very significant in all inflationary models. Quantum fluctuations are invariably important on very small scales, and with inflation these very small scales are rapidly stretched to become macroscopic and even astronomical. Thus the quantum fluctuations of the inflaton field can have very noticeable effects.
When the mass of the scalar field is small compared to the Hubble parameter $`H`$, these quantum evolution of the scalar field is accurately described as a random walk. It is useful to divide space into regions of physical size $`H^1`$, and to discuss the average value of the scalar field $`\varphi `$ within a given region. In a time $`H^1`$, the quantum fluctuations cause the scalar field to undergo a random Gaussian jump of zero mean and a root-mean-squared magnitude given by
$$\mathrm{\Delta }\varphi _{\mathrm{qu}}=\frac{H}{2\pi }.$$
(4.1)
This random quantum jump is superimposed on the classical motion, as indicated in Fig. (4).
To illustrate how eternal inflation happens in the simplest context, let us consider again the free scalar field described by the potential function of Eq. (2.1). We consider a region of physical radius $`H^1`$, in which the field has an average value $`\varphi `$. Using Eq. (2.9) along with Eqs. (2.8) and (2.1), one finds that the magnitude of the classical change that the field will undergo in a time $`H^1`$ is given in by
$$\mathrm{\Delta }\varphi _{\mathrm{cl}}=\frac{M_\mathrm{P}m}{\sqrt{12\pi }}H^1=\frac{1}{4\pi }\frac{M_\mathrm{P}^2}{\varphi }.$$
(4.2)
Let $`\varphi ^{}`$ denote the value of $`\varphi `$ which is large enough so that
$$\mathrm{\Delta }\varphi _{\mathrm{qu}}(\varphi ^{})=\mathrm{\Delta }\varphi _{\mathrm{cl}}(\varphi ^{}),$$
(4.3)
which implies that
$$\varphi ^{}=\left(\frac{3}{16\pi }\right)^{1/4}\frac{M_\mathrm{P}^{3/2}}{m^{1/2}}.$$
(4.4)
Now consider what happens to a region for which the initial average value of $`\varphi `$ is equal to $`\varphi ^{}`$. In a time interval $`H^1`$, the volume of the region will increase by $`e^320`$. At the end of the time interval we can divide the original region into 20 regions of the same volume as the original, and in each region the average scalar field can be written as
$$\varphi =\varphi ^{}+\mathrm{\Delta }\varphi _{\mathrm{cl}}+\delta \varphi ,$$
(4.5)
where $`\delta \varphi `$ denotes the random quantum jump, which is drawn from a Gaussian probability distribution with standard deviation $`\mathrm{\Delta }\varphi _{\mathrm{qu}}=\mathrm{\Delta }\varphi _{\mathrm{cl}}`$. Gaussian statistics imply that there is a 15.9% chance that a Gaussian random variable will exceed its mean by more than one standard deviation, and therefore there is a 15.9% chance that the net change in $`\varphi `$ will be positive. Since there are now 20 regions of the original volume, on average the value of $`\varphi `$ will exceed the original value in 3.2 of these regions. Thus the volume for which $`\varphi \varphi ^{}`$ does not (on average) decrease, but instead increases by more than a factor of 3. Since this argument can be repeated, the expectation value of the volume for which $`\varphi \varphi ^{}`$ increases exponentially with time. Typically, therefore, inflation never ends, but instead the volume of the inflating region grows exponentially without bound. The minimum field value for eternal inflation is a little below $`\varphi ^{}`$, since a volume increase by a factor of 3.2 is more than necessary—any factor greater than one would be sufficient. A short calculation shows that the minimal value for eternal inflation is $`0.78\varphi ^{}`$.
While the value of $`\varphi ^{}`$ is larger than Planck scale, again we find that this is not true of the energy density:
$$V(\varphi ^{})=\frac{1}{2}m^2\varphi ^2=\sqrt{\frac{3}{64\pi }}mM_\mathrm{P}^3,$$
(4.6)
which for $`m=10^{16}`$ GeV gives an energy density of $`1\times 10^4M_\mathrm{P}^4`$.
If one carries out the same analysis with a potential function
$$V(\varphi )=\frac{1}{4}\lambda \varphi ^4,$$
(4.7)
one finds that
$$\varphi ^{}=\left(\frac{3}{2\pi \lambda }\right)^{1/6}M_\mathrm{P},$$
(4.8)
and
$$V(\varphi ^{})=\left(\frac{3}{16\pi }\right)^{2/3}\lambda ^{1/3}M_\mathrm{P}^4.$$
(4.9)
Since $`\lambda `$ must be very small in any case so that density perturbations are not too large, one finds again that eternal inflation is predicted to happen at an energy density well below the Planck scale.
## V. Eternal Inflation: Implications
When I told Rocky Kolb that I was going to be talking about eternal inflation, he said, “That’s OK, we can talk about physics later.” So that’s the point I’d like to address here. In spite of the fact that the other universes created by eternal inflation are too remote to imagine observing directly, I still believe that eternal inflation has real consequences in terms of the way we extract predictions from theoretical models. Specifically, there are four consequences of eternal inflation that I will highlight.
1) Unobservability of initial conditions
First, eternal inflation implies that all hypotheses about the ultimate initial conditions for the universe—such as the Hartle-Hawking no boundary proposal, the tunneling proposals by Vilenkin or Linde , or the more recent Hawking-Turok instanton —become totally divorced from observation. That is, one would expect that if inflation is to continue arbitrarily far into the future with the production of an infinite number of pocket universes, then the statistical properties of the inflating region should approach a steady state which is independent of the initial conditions. Unfortunately, attempts to quantitatively study this steady state are severely limited by several factors. First, there are ambiguities in defining probabilities, which will be discussed later. In addition, the steady state properties seem to depend strongly on super-Planckian physics which we do not understand. That is, the same quantum fluctuations that make eternal chaotic inflation possible tend to drive the scalar field further and further up the potential energy curve, so attempts to quantify the steady state probability distribution require the imposition of some kind of a boundary condition at large $`\varphi `$. Although these problems remain unsolved, I still believe that it is reasonable to assume that in the course of its perpetual evolution, an eternally inflating universe would lose all memory of the state in which it started.
Although the ultimate origin of the universe would become unobservable, I would not expect that the question of how the universe began would lose its interest. While eternally inflating universes continue forever once they start, they are presumably not eternal into the past. (The word eternal is therefore not technically correct—it would be more precise to call this scenario semi-eternal or future-eternal.) While the issue is not completely settled, it appears likely that eternally inflating universes must necessarily have a beginning. Borde and Vilenkin have shown, provided that certain conditions are met, that spacetimes which are future-eternal must have an initial singularity, in the sense that they cannot be past null geodesically complete. The proof, however, requires the weak energy condition, which is classically valid but quantum-mechanically violated . In any case, I am not aware of any viable model without a beginning, and certainly nothing that we know can rule out the possibility of a beginning. The possibility of a quantum origin of the universe is very attractive, and will no doubt be a subject of interest for some time. Eternal inflation, however, seems to imply that the entire study will have to be conducted with literally no input from observation.
2) Irrelevance of initial probability
A second consequence of eternal inflation is that the probability of the onset of inflation becomes totally irrelevant, provided that the probability is not identically zero. Various authors in the past have argued that one type of inflation is more plausible than another, because the initial conditions that it requires appear more likely to have occurred. In the context of eternal inflation, however, such arguments have no significance.
To illustrate the insignificance of the probability of the onset of inflation, I will use a numerical example. We will imagine comparing two different versions of inflation, which I will call Type A and Type B. They are both eternally inflating—but Type A will have a higher probability of starting, while Type B will be a little faster in its exponential expansion rate. Since I am trying to show that the higher starting probability of Type A is irrelevant, I will choose my numbers to be extremely generous to Type A. First, we must choose a number for how much more probable it is for Type A inflation to begin, relative to type B. A googol, $`10^{100}`$, is usually considered a large number—it is some 20 orders of magnitude larger than the total number of baryons in the visible universe. But I will be more generous: I will assume that Type A inflation is more likely to start than type B inflation by a factor of $`10^{1,000,000}`$. Type B inflation, however, expands just a little bit faster, say by 0.001%. We need to choose a time constant for the exponential expansion, which I will take to be a typical grand unified theory scale, $`\tau =10^{37}`$ sec. ($`\tau `$ represents the time constant for the overall expansion factor, which takes into account both the inflationary expansion and the exponential decay of the false vacuum.) Finally, we need to choose a length of time to let the system evolve. In principle this time interval is infinite (the inflation is eternal into the future), but to be conservative we will watch the system for only one second.
We imagine setting up a statistical ensemble of universes at $`t=0`$, with an expectation value for the volume of Type A inflation exceeding that of Type B inflation by $`10^{1,000,000}`$. For brevity, let the term “weight” to refer to the ensemble expectation value of the volume. Thus, the weights of Type A inflation and Type B inflation will begin with the ratio
$$\frac{W_B}{W_A}|_{t=0}=10^{1,000,000}.$$
(5.1)
After one second of evolution, the expansion factors for Type A and Type B inflation will be
$`Z_A`$ $`=`$ $`e^{t/\tau }=e^{10^{37}}`$ (5.2)
$`Z_B`$ $`=`$ $`e^{1.00001t/\tau }=e^{0.00001t/\tau }Z_A`$ (5.3)
$`=`$ $`e^{10^{32}}Z_A10^{4.3\times 10^{31}}Z_A.`$
The weights at the end of one second are proportional to these expansion factors, so
$$\frac{W_B}{W_A}|_{t=1\mathrm{sec}}=10^{(4.3\times 10^{31}1,000,000)}.$$
(5.4)
Thus, the initial ratio of $`10^{1,000,000}`$ is vastly superseded by the difference in exponential expansion factors. In fact, we would have to calculate the exponent of Eq. (5.4) to an accuracy of 25 significant figures to be able to barely detect the effect of the initial factor of $`10^{1,000,000}`$.
One might criticize the above argument for being naive, as the concept of time was invoked without any discussion of how the equal-time hypersurfaces are to be chosen. I do not know a decisive answer to this objection; as I will discuss later, there are unresolved questions concerning the calculation of probabilities in eternally inflating spacetimes. Nonetheless, given that there is actually an infinity of time available, it is seems reasonable to believe that the form of inflation that expands the fastest will always dominate over the slower forms by an infinite factor.
A corollary to this argument is that new inflation is not dead. While the initial conditions necessary for new inflation cannot be justified on the basis of thermal equilibrium, as proposed in the original papers , in the context of eternal inflation it is sufficient to conclude that the probability for the required initial conditions is nonzero. Since the resulting scenario does not depend on the words that are used to justify the initial state, the standard treatment of new inflation remains valid.
3) Inevitability of eternal inflation
Third, I’d like to claim that, since it appears that a universe is in principle capable of eternally reproducing, it is hard to believe that any other description can make sense at all. To clarify this point, let me raise the analogy of rabbits. We all know that rabbits can reproduce—in fact, they reproduce like rabbits. Suppose that you went out into the woods and found a rabbit that had characteristics indicating that it did not belong to any known rabbit species. Then you would have to theorize about how the rabbit originated. You might entertain the notion that the rabbit was created by some unique, mysterious, cosmic event that you hope to someday understand better. Or you could assume that the rabbit was created by the process of rabbit reproduction that we all know so well. I think that we would all consider the latter possibility to be far more plausible. So, I claim that once we become convinced that universes can reproduce like rabbits, then the situations are similar. When we notice that there is a universe and ask how it originated, the same inferences that we made for the rabbit question should apply to this one.
4) Possibility of restoring the uniqueness of theoretical predictions
A fourth consequence of eternal inflation is the possibility that it offers to rescue the predictive power of theoretical physics. All the indications suggest that string theory or M theory describes an elegantly unique theoretical structure, but nonetheless it seem unlikely that the theory possesses a unique vacuum. Since predictions will ultimately depend on the properties of the vacuum, the predictive power of string/M theory may be limited. Eternal inflation, however, provides a hope that this problem can be remedied. Even if many types of vacua are equally stable, it may turn out that a unique state produces the maximum possible rate of inflation. If so, then this state will dominate the universe, even if its expansion rate is only infinitesimally larger than the other possibilities. Thus, eternal inflation might allow physicists to extract unique predictions, in spite of the multiplicity of stable vacua.
## VI. Difficulties in Calculating Probabilities
In an eternally inflating universe, anything that can happen will happen; in fact, it will happen an infinite number of times. Thus, the question of what is possible becomes trivial—anything is possible, unless it violates some absolute conservation law. To extract predictions from the theory, we must therefore learn to distinguish the probable from the improbable.
However, as soon as one attempts to define probabilities in an eternally inflating spacetime, one discovers ambiguities. Since an eternally inflating universe produces an infinite number of pocket universes, the sample space is infinite. The fraction of universes with any particular property is given by the meaningless ratio of infinity divided by infinity. To obtain a well-defined answer, one needs to invoke some method of regularization. The most straightforward form of regularization consists of truncating the space to a finite subspace, and then taking a limit in which the subspace becomes larger and larger.
To understand the nature of the problem, it is useful to think about the integers as a model system with an infinite number of entities. We can ask, for example, what fraction of the integers are odd. Most people would presumably say that the answer is $`1/2`$, since the integers alternate between odd and even. That is, if the string of integers is truncated after the $`N`$th, then the fraction of odd integers in the string is exactly $`1/2`$ if $`N`$ is even, and is $`(N+1)/2N`$ if $`N`$ is odd. In any case, the fraction approaches $`1/2`$ as $`N`$ approaches infinity.
However, the ambiguity of the answer can be seen if one imagines other orderings for the integers. One could, if one wished, order the integers as
$$1,3,2,5,7,4,9,11,6,\mathrm{},$$
(6.1)
always writing two odd integers followed by one even integer. This series includes each integer exactly once, just like the usual sequence ($`1,2,3,4,\mathrm{}`$). The integers are just arranged in an unusual order. However, if we truncate the sequence shown in Eq. (6.1) after the $`N`$th entry, and then take the limit $`N\mathrm{}`$, we would conclude that 2/3 of the integers are odd. Thus, we see that probabilities can depend nontrivially on the method of regularization that is used.
In the case of eternally inflating spacetimes, one might consider a regularization defined by ordering the pocket universes in the sequence in which they form, and then truncating after the $`N`$th. However, each pocket universe fills its own future light cone, so no pocket universe forms in the future light cone of another. Any two pocket universes are spacelike separated from each other, so different observers can disagree about which formed first. One can arbitrarily choose equal-time surfaces that foliate the spacetime, and then truncate at some value of $`t`$, but this recipe is far from unique. In practice, different ways of choosing equal-time surfaces give different results.
## VII. The Youngness Paradox
If one chooses a regularization in the most naive way, one is led to a set of very peculiar results which I call the youngness paradox.
Specifically, suppose that one constructs a Robertson-Walker coordinate system while the model universe is still in the false vacuum (de Sitter) phase, before any pocket universes have formed. One can then propagate this coordinate system forward with a synchronous gauge condition,By a synchronous gauge condition, I mean that each equal-time hypersurface is obtained by propagating every point on the previous hypersurface by a fixed infinitesimal time interval $`\mathrm{\Delta }t`$ in the direction normal to the hypersurface. and one can define probabilities by truncating at a fixed value $`t_f`$ of the synchronous time coordinate $`t`$. That is, the probability of any particular property can be taken to be proportional to the volume on the $`t=t_f`$ hypersurface which has that property. This method of defining probabilities was studied in detail by Linde, Linde, and Mezhlumian, in a paper with the memorable title “Do we live in the center of the world?” . I will refer to probabilities defined in this way as synchronous gauge probabilities.
The youngness paradox is caused by the fact that the volume of false vacuum is growing exponentially with time with an extraordinarily short time constant, in the vicinity of $`10^{37}`$ sec. Since the rate at which pocket universes form is proportional to the volume of false vacuum, this rate is increasing exponentially with the same time constant. That means that in each second the number of pocket universes that exist is multiplied by a factor of $`\mathrm{exp}\left\{10^{37}\right\}`$. At any given time, therefore, almost all of the pocket universes that exist are universes that formed very very recently, within the last several time constants. The population of pocket universes is therefore an incredibly youth-dominated society, in which the mature universes are vastly outnumbered by universes that have just barely begun to evolve. Although a mature universe has a larger volume then a young one, this multiplicative factor is of little importance, since in synchronous coordinates the volume no longer grows exponentially once the pocket universe forms.
Probability calculations in this youth-dominated ensemble lead to peculiar results, as discussed in Ref. . These authors considered the expected behavior of the mass density in our vicinity, concluding that we should find ourselves very near the center of a spherical low-density region. Here I would like to discuss a less physical but simpler question, just to illustrate the paradoxes associated with synchronous gauge probabilities. Specifically, I will consider the question: “Are there any other civilizations in the visible universe that are more advanced than ours?”. Intuitively I would not expect inflation to make any predictions about this question, but I will argue that the synchronous gauge probability distribution strongly implies that there is no civilization in the visible universe more advanced than we are.
Suppose that we have reached some level of advancement, and suppose that $`t_{\mathrm{min}}`$ represents the minimum amount of time needed for a civilization as advanced as we are to evolve, starting from the moment of the decay of the false vacuum—the start of the big bang. The reader might object on the grounds that there are many possible measures of advancement, but I would respond by inviting the reader to pick any measure she chooses; the argument that I am about to give should apply to all of them. The reader might alternatively claim that there is no sharp minimum $`t_{\mathrm{min}}`$, but instead we should describe the problem in terms of a function which gives the probability that, for any given region within a pocket universe of the size of our visible universe, a civilization as advanced as we are would develop by time $`t`$. I believe, however, that the introduction of such a probability distribution would merely complicate the argument, without changing the result. So, for simplicity of discussion, I will assume that there is some sharply defined minimum time $`t_{\mathrm{min}}`$ required for a civilization as advanced as ours to develop.
Since we exist, our pocket universe must have an age $`t_0`$ satisfying
$$t_0t_{\mathrm{min}}.$$
(7.1)
Suppose, however, that there is some civilization in our visible universe that is more advanced than we are, let us say by 1 second. In that case Eq. (7.1) is not sufficient, but instead the age of our pocket universe would have to satisfy
$$t_0t_{\mathrm{min}}+1\text{ second}.$$
(7.2)
However, in the synchronous gauge probability distribution, universes that satisfy Eq. (7.2) are outnumbered by universes that satisfy Eq. (7.1) by a factor of approximately $`\mathrm{exp}\left\{10^{37}\right\}`$. Thus, if we know only that we are living in a pocket universe that satisfies Eq. (7.1), the probability that it also satisfies Eq. (7.2) is approximately $`\mathrm{exp}\left\{10^{37}\right\}`$. We would conclude, therefore, that it is extraordinarily improbable that there is a civilization in our visible universe that is at least 1 second more advanced than we are.
Perhaps this argument explains why SETI has not found any signals from alien civilizations, but I find it more plausible that it is merely a symptom that the synchronous gauge probability distribution is not the right one.
## VIII. Toy Model of Eternal Inflation
The conceptual issue involved in the youngness paradox can perhaps be clarified by considering a toy model of a highly simplified eternally inflating universe. Suppose that the universe as a whole can be labeled with a global time variable $`t`$, and that it consists of a countably infinite set of pocket universes, each of which is labeled by an index $`i`$. For simplicity, we let each pocket universe have zero spatial dimensions, so a spacetime point is fully specified by the time $`t`$ and the index $`i`$ which indicates the pocket universe in which it is located. We assume that each pocket universe $`i`$ forms at some time $`t_i=n_i\tau `$, where $`n_i`$ is an integer and $`\tau `$ is a fixed time constant characterizing the entire universe. Let the number of pocket universes that form at time $`t=n\tau `$ be equal to $`2^n`$, for each nonnegative integer $`n`$. Assume that each pocket universe exists for a time $`T\tau `$, and then disappears, and that within each pocket universe the interval from the time of formation to disappearance is uniformly populated with “sentient beings.” Within each pocket universe we can define a relative time, $`t_{\mathrm{rel}}tt_i`$, which measures the amount of time since the formation of the pocket universe.
The difficult question, then, is the following: At what relative time $`t_{\mathrm{rel}}`$ does a typical sentient being live? If one answers this question by truncating the spacetime by the criterion
$$tt_c,$$
(8.1)
for some cut-off time $`t_c=n_c\tau `$, then one finds that most of the pocket universes in the truncated spacetime formed within the past few time constants. As $`t_c\mathrm{}`$, the mean value of $`t_{\mathrm{rel}}`$ approaches $`\tau `$. This method is analogous to the synchronous gauge cut-off discussed above. If, however, one truncates the spacetime by including all pocket universes for which the time of formation
$$t_it_c,$$
(8.2)
then the mean value of $`t_{\mathrm{rel}}`$ is equal to $`T/2`$ for any $`t_c`$. The truncation method of Eq. (8.1) leads to the youngness paradox, in which the probability sample is strongly dominated by universes that are extremely young, while the truncation method of Eq. (8.2) does not.
At this point, I have to admit that I do not understand how to resolve the ambiguities associated with this toy model. It is conceivable that there is no meaningful method of regularization, and that $`t_{\mathrm{rel}}`$ is somehow not susceptible to probabilistic predictions. It is also conceivable that there is something wrong with either the truncation (8.1) or (8.2) or both, and that a correct analysis would lead to a unique probability calculation. It is also conceivable that the regularization has to be specified as part of the theory, so that the truncations (8.1) and (8.2) represent two distinct theories, each of which is logically consistent.
## IX. An Alternative Probability Prescription
Since the probability measure depends on the method used to regulate the infinite spacetime of eternal inflation, we are not forced to accept the consequences of the synchronous gauge probabilities. A very attractive alternative has been proposed by Vilenkin , and developed further by Vanchurin, Vilenkin, and Winitzki . This procedure is, roughly speaking, analogous to the truncation of Eq. (8.2).
The key idea of the Vilenkin proposal is to define probabilities within a single pocket universe (which he describes more precisely as a connected, thermalized domain). Thus, unlike the synchronous gauge method, there is no comparison between old pocket universes and young ones. To justify this approach it is crucial to recognize that each pocket universe is infinite, even if one starts the model with a finite region of de Sitter space. The infinite volume arises in the same way as it does for the special case of Coleman-de Luccia bubbles , the interior of which are open Robertson-Walker universes. From the outside one often describes such bubbles in a coordinate system in which they are finite at any fixed time, but in which they grow without bound. On the inside, however, the natural coordinate system is the one that reflects the intrinsic homogeneity, in which the space is infinite at any given time. The infinity of time, as seen from the outside, becomes an infinity of spatial extent as seen on the inside. Thus, at least for continuously variable parameters, a single pocket universe provides an infinite sample space which can be used to define probabilities. The second key idea of Vilenkin’s method is to use the inflaton field itself as the time variable, rather than the synchronous time variable discussed in the previous section.
This approach can be used, for example, to discuss the probability distribution for $`\mathrm{\Omega }`$ in open inflationary models, or to discuss the probability distribution for some arbitrary field that has a flat potential energy function. If, however, the vacuum has a discrete parameter which is homogeneous within each pocket universe, but which takes on different values in different pocket universes, then this method does not apply.
The proposal can be described in terms of Fig. 5. We suppose that the theory includes an inflaton field $`\varphi `$ of the new inflation type, and some set of fields $`\chi _i`$ which have flat potentials. The goal is to find the probability distribution for the fields $`\chi _i`$. We assume that the evolution of the inflaton $`\varphi `$ can be divided into three regimes, as shown on the figure. $`\varphi <\varphi _1`$ describes the eternally inflating regime, in which the evolution is governed by quantum diffusion. For $`\varphi _1<\varphi <\varphi _{\mathrm{end}}`$, the evolution is described classically in a slow-roll approximation, so that $`\dot{\varphi }\mathrm{d}\varphi /\mathrm{d}t`$ can be expressed as a function of $`\varphi `$. For $`\varphi >\varphi _{\mathrm{end}}`$ inflation is over, and the $`\varphi `$ field no longer plays an important role in the evolution. The $`\chi _i`$ fields are assumed to have a finite range of values, such as angular variables, so that a flat probability distribution is normalizable. They are assumed to have a flat potential energy function for $`\varphi >\varphi _{\mathrm{end}}`$, so that they could settle at any value. They are also assumed to have a flat potential energy function for $`\varphi <\varphi _1`$, although they might interact with $`\varphi `$ during the slow-roll regime, however, so that they can affect the rate of inflation.
Since the potential for the $`\chi _i`$ is flat for $`\varphi <\varphi _1`$, we can assume that they begin with a flat probability distribution $`P_0(\chi _i)P(\chi _i,\varphi _1)`$ on the $`\varphi =\varphi _1`$ hypersurface. If the kinetic energy function for the $`\chi _i`$ is of the standard form, we take $`P_0(\chi _i)=const`$. If, however, the kinetic energy is nonstandard,
$$_{\mathrm{kinetic}}=g^{ij}(\chi )_\mu \chi _i^\mu \chi _j,$$
(9.1)
as is plausible for a field described in angular variables, then the initial probability distribution is assumed to take the reparameterization-invariant form
$$P_0(\chi _i)\sqrt{detg}.$$
(9.2)
During the slow-roll era, it is assumed that the $`\chi _i`$ fields evolve classically, so one can calculate the number of e-folds of inflation $`N(\chi _i)`$ as a function of the final value of the $`\chi _i`$ (i.e., the value of $`\chi _i`$ on the $`\varphi =\varphi _{\mathrm{end}}`$ hypersurface). One can also calculate the final values $`\chi _i`$ in terms of the initial values $`\chi _i^0`$ (i.e., the value of $`\chi _i`$ on the $`\varphi =\varphi _1`$ hypersurface). One then assumes that the probability density is enhanced by the volume inflation factor $`e^{3N(\chi _i)}`$. The evolution from $`\chi _i^0`$ to $`\chi _i`$ results in a Jacobian factor. The (unnormalized) final probability distribution is thus given by
$$P(\chi _i,\varphi _{\mathrm{end}})=P_0(\chi _i^0)e^{3N(\chi _i)}det\frac{\chi _j^0}{\chi _k}.$$
(9.3)
Alternatively, if the evolution of the $`\chi _i`$ during the slow-roll era is subject to quantum fluctuations, Ref. shows how to write a Fokker-Planck equation which is equivalent to averaging the result of Eq. (9.3) over a collection of paths that result from interactions with a noise term.
The Vilenkin proposal sidesteps the youngness paradox by defining probabilities by the comparison of volumes within one pocket universe. The youngness paradox, in contrast, arose when one considered a probability ensemble of all pocket universes at a fixed value of the synchronous gauge time coordinate—an ensemble that is overwhelmingly dominated by very young pocket universes.
The proposal has the drawback, however, that it cannot be used to compare the probabilities of discretely different alternatives. Furthermore, although the results of this method seem reasonable, I do not at this point find them compelling. That is, it is not clear what principles of physics or probability theory ensure that this particular method of regularizing the spacetime is the one that leads to correct predictions. Perhaps there is no way to answer this question, so we may be forced to accept this proposal, or something similar to it, as a postulate.
## X. Probabilities with only one universe?
In discussing a probabilistic approach to cosmology, we need to know whether it makes sense to talk about a probability distribution for a cosmic parameter such as $`\mathrm{\Omega }`$, for which we have only one example to measure. I have certainly heard more than one physicist say that he or she doesn’t think that one can meaningfully talk about probabilities for an experiment that can be done only once. The notion that probability requires repetition is very widespread, and I am sure that it is incorporated into many books about probability theory. Nonetheless, I would like to argue that repetition is not at all necessary to make use of probability theory. Instead, I will argue that probability is meaningful whenever one has a strong probabilistic prediction, by which I mean a prediction that the probability for some discernible event is either very close to zero or very close to one.
Thus, if a cosmological theory predicts a probability distribution for $`\mathrm{\Omega }`$ which is reasonably flat, then there is no strong prediction, and the implications of the theory for $`\mathrm{\Omega }`$ do not provide a way of testing the theory. However, if the theory predicts that the probability of $`\mathrm{\Omega }`$ lying outside the range of 0.99 to 1.01 is $`10^6`$, then I would claim that the prediction is meaningful and can be used to test the theory.
My point of view can be explained most easily by considering coin flips. If a flip of an unbiased coin is repeated 20 times, the probability of getting 20 successive heads is a very small number, about $`10^6`$. This is an example of what I call a strong prediction. Many common examples of strong predictions involve repetition. However, if 20 unbiased coins are flipped simultaneously in a single experiment, the probability that all will come up heads is identical, about $`10^6`$. Since the probability of these two results—20 successive heads or 20 simultaneous heads—are both equally small, I would draw the obvious conclusion that we should be equally surprised if either result occurred. It does not matter that the first result involved repetition, while the second did not. Some might argue that the 20 simultaneous coin flips involved the replication of identical experiments, even if they were not performed in succession, so I will take the analogy one step further. Suppose we constructed a roulette wheel that was so finely ruled that the probability of the ball landing on 0 was only $`10^6`$. Again, we should be just as surprised if this result occurred as we would be if 20 successive coins landed heads.
Similarly, if our cosmological theory predicted that the probability of $`\mathrm{\Omega }`$ lying outside the range of 0.99 to 1.01 is $`10^6`$, we should be just as surprised if this outcome occurred as we would be if 20 consecutive coins came up heads. In both cases, we would have good cause to question the assumptions that went into calculating the prediction.
## XI. Conclusion
In this paper I have summarized the workings of inflation, and the arguments that strongly suggest that our universe is the product of inflation. I argued that inflation can explain the size, the Hubble expansion, the homogeneity, the isotropy, and the flatness of our universe, as well as the absence of magnetic monopoles, and even the characteristics of the nonuniformities. The detailed observations of the cosmic background radiation anisotropies continue to fall in line with inflationary expectations, and the evidence for an accelerating universe fits well with the inflationary preference for a flat universe.
Next I turned to the question of eternal inflation, claiming that essentially all inflationary models are eternal. In my opinion this makes inflation very robust: if it starts anywhere, at any time in all of eternity, it produces an infinite number of pocket universes. Eternal inflation has the very attractive feature, from my point of view, that it offers the possibility of allowing unique predictions even if the underlying string theory does not have a unique vacuum. I have also emphasized, however, that there are important problems in understanding the implications of eternal inflation. First, there is the problem that we do not know how to treat the situation in which the scalar field climbs upward to the Planck energy scale. Second, the definition of probabilities in an eternally inflating spacetime is not yet a closed issue, although important progress has been made. And third, I might add that the entire present approach is at best semiclassical. A better treatment may not be possible until we have a much better handle on quantum gravity, but eventually this issue will have to be faced.
## Acknowledgments
The author particularly thanks Andrei Linde, Alexander Vilenkin, Neil Turok, and other participants in the Isaac Newton Institute programme Structure Formation in the Universe for very helpful conversations. This work is supported in part by funds provided by the U.S. Department of Energy (D.O.E.) under cooperative research agreement #DF-FC02-94ER40818, and in part by funds provided by NM Rothschild & Sons Ltd and by the EPSRC. |
warning/0002/math0002205.html | ar5iv | text | # On the existence of absolutely simple abelian varieties of a given dimension over an arbitrary field
## 1. Introduction
An abelian variety over a field $`k`$ is called simple if it has no proper nonzero sub-abelian varieties over $`k`$; it is called absolutely simple (or geometrically simple) if it is simple over the algebraic closure of $`k`$. In this paper we will prove the following theorem:
###### Theorem 1.
Let $`k`$ be a field and let $`n`$ be a positive integer. Then there exists an absolutely simple $`n`$-dimensional abelian variety over $`k`$.
An easy reduction argument, similar in spirit to the one in Section 4 of , shows that an absolutely simple abelian variety over a field $`k`$ remains simple over every extension field of $`k`$, even the non-algebraic ones; thus, it suffices to prove Theorem 1 in the special case where $`k`$ is a prime field. Mori (see also Zarhin ) provides examples of absolutely simple abelian varieties of arbitrary dimension over $`𝐐`$, so we need only prove Theorem 1 for finite prime fields $`k`$. We will in fact prove that over such fields there exist absolutely simple ordinary abelian varieties of every dimension, and in addition we will prove an asymptotic result concerning arbitrary finite fields:
###### Theorem 2.
For every integer $`n0`$ and finite field $`k`$ let $`S(k,n)`$ denote the fraction of the isogeny classes of $`n`$-dimensional abelian varieties over $`k`$ that consist of absolutely simple ordinary abelian varieties. Then for every $`n`$ we have $`S(𝐅_q,n)1`$ as $`q\mathrm{}`$ over the prime powers.
In fact, for every $`n`$ and $`ϵ`$ we provide an explicit value of $`M`$ such that if $`q>M`$ then $`0<1S(𝐅_q,n)<ϵ`$; see Theorems 13 and 14 in Sections 6 and 7.
Suppose $`A`$ is an abelian variety over a finite field $`k`$. One can ask whether there exists an absolutely simple abelian variety over $`k`$ with the same formal isogeny type (see ) as $`A`$. Theorem 1 shows that the answer to this question is yes when $`A`$ is ordinary; we have not considered the question for other formal isogeny types. Lenstra and Oort considered the analogous question when $`k`$ is the algebraic closure of a finite field, and showed that the answer is yes when $`A`$ is not supersingular.
Theorem 1 leads to the question of whether there exist absolutely simple Jacobians of every dimension over a given field $`k`$. Chai and Oort show that the answer is yes if $`k`$ is the algebraic closure of a finite field, and Mori and Zarhin show that the answer is also yes if $`k`$ has characteristic zero. If $`k`$ has positive characteristic $`p`$ but is not algebraic over $`𝐅_p`$, then results of Katz and Sarnak (see Sections 10.1 and 10.2 of ) can be used to show that once again the answer is yes — see also Mori for some partial results for such fields. In fact, the examples provided by Katz and Sarnak, Mori, and Zarhin are Jacobians of explicitly-given hyperelliptic curves. However, the question seems to be open when $`k`$ is a finite field. The techniques we use in this paper do not help settle the general question for finite fields, but our results do at least show that over every finite field $`k`$ there are curves of genus $`2`$ and $`3`$ with absolutely simple Jacobians, as the following argument shows:
As we mentioned above, we show that for every $`n`$ and for every finite field $`k`$ there is an absolutely simple $`n`$-dimensional ordinary abelian variety over $`k`$, and in particular this is true for $`n=2`$ and $`n=3`$. But every absolutely simple ordinary abelian variety of dimension $`2`$ or $`3`$ over a finite field is isogenous to a principally polarized variety (see Corollary 12.6 and Theorem 1.2 of ). The main result of shows that each such principally polarized variety is isomorphic (over the algebraic closure of $`k`$) to a Jacobian of a possibly reducible curve $`C`$, but since the Jacobian of $`C`$ is absolutely simple $`C`$ must be geometrically irreducible. Finally, a simple descent argument shows that $`C`$ has a model defined over $`k`$. Thus, for every finite field $`k`$ there are curves of genus $`2`$ and $`3`$ over $`k`$ with absolutely simple Jacobians.
Our paper is organized as follows. In Section 2, we briefly review the properties of Weil numbers and Weil polynomials that we will use in the proofs of Theorems 1 and 2. In Section 3 we give an easy-to-verify sufficient condition for an abelian variety over a finite field to be absolutely simple. We use this condition in Section 4 to prove Theorem 6, which shows how the characteristic polynomial of Frobenius of a simple ordinary abelian surface over a finite field can be used to quickly determine the splitting behavior of the surface over the algebraic closure. Theorem 6 allows us to give a very short proof of Theorem 1 in the case $`n=2`$; we provide a proof for the case $`n>2`$ in Section 5. In Sections 6 and 7 we prove Theorems 13 and 14, which are effective versions of Theorem 2 in the cases $`n=2`$ and $`n>2`$, respectively. Finally, in Section 8 we prove a lemma about polynomials with prescribed reduction modulo certain primes that is essential for our proof of Theorem 14.
###### Conventions .
Suppose $`A`$ is an abelian variety over a field $`k`$ and suppose $`\mathrm{}`$ is an extension field of $`k`$. We will denote by $`A_{\mathrm{}}`$ the $`\mathrm{}`$-scheme $`A\times _{\mathrm{Spec}k}\mathrm{Spec}\mathrm{}`$. If $`B`$ is another abelian variety over $`k`$, then when we speak of a morphism from $`A`$ to $`B`$ we always mean a $`k`$-morphism; thus, we write $`\mathrm{End}A`$ for what some authors would call $`\mathrm{End}_kA`$.
###### Acknowledgments .
The authors thank David Cantor, Robert Coleman, Daniel Goldstein, Hendrik Lenstra, Bjorn Poonen, and Joel Rosenberg for helpful conversations and correspondence, and Mike Zieve for asking the questions that led to this research and for commenting on an early version of this paper. The authors are grateful to Hendrik Lenstra for suggesting Lemma 5. The authors used the computer packages PARI/GP and MAGMA for some of the computations they performed in the course of writing this paper.
## 2. Weil numbers and Weil polynomials
Suppose $`q`$ is a power of a prime number $`p`$. A Weil $`q`$-number, or simply a Weil number if $`q`$ is clear from context, is an algebraic integer $`\pi `$ such that $`|\phi (\pi )|=q^{1/2}`$ for every embedding $`\phi `$ of $`𝐐(\pi )`$ into the complex numbers. Suppose $`k`$ is a field with $`q`$ elements. To every abelian variety $`A`$ over $`k`$ we associate the characteristic polynomial $`f_A𝐙[x]`$ of its Frobenius endomorphism (acting on the $`\mathrm{}`$-adic Tate modules of $`A`$); the polynomial $`f_A`$ is monic of degree twice the dimension of $`A`$. We call a polynomial $`f`$ a Weil $`q`$-polynomial, or simply a Weil polynomial, if there is an abelian variety $`A`$ over $`k`$ with $`f=f_A`$. Weil proved that all of the roots of a Weil polynomial are Weil numbers, and Honda showed that every Weil number is a root of some Weil polynomial. Furthermore, Tate showed that two abelian varieties over $`k`$ are isogenous if and only if their associated Weil polynomials are equal. If $`A`$ is a simple abelian variety over $`k`$ then $`f_A`$ is a power of an irreducible polynomial, and in fact the Honda-Tate theorem (see \[13, Théorème 1\]) says the map that sends $`A`$ to the set of roots (in $`\overline{𝐐}`$) of $`f_A`$ induces a bijection between the set of isogeny classes of simple abelian varieties over $`k`$ and the set of Galois conjugacy classes of Weil numbers in $`\overline{𝐐}`$. The Honda-Tate theorem also provides a simple number-theoretic criterion for determining whether a polynomial, all of whose roots are Weil numbers, is a Weil polynomial. In addition, the theorem shows how the Weil polynomial of an abelian variety $`A`$ over $`k`$ determines the algebra $`(\mathrm{End}A)𝐐`$.
An abelian variety $`A`$ over $`k`$ is ordinary if the rank of its group of $`p`$-torsion points over the algebraic closure of $`k`$ is equal to the dimension of $`A`$; a Weil polynomial is ordinary if it is the characteristic polynomial of Frobenius of an ordinary abelian variety; and a Weil number is ordinary if its minimal polynomial is an ordinary Weil polynomial. The Honda-Tate theorem simplifies considerably if one considers only ordinary varieties and ordinary Weil polynomials — see Section 3 of . For example, a monic polynomial in $`𝐙[x]`$ is an ordinary Weil $`q`$-polynomial if and only if it is of even degree $`2n`$, all of its roots are Weil numbers, and its middle coefficient (that is, the coefficient of $`x^n`$) is coprime to $`q`$. Furthermore, an ordinary abelian variety $`A`$ over $`k`$ is simple if and only if its Weil polynomial $`f`$ is irreducible. If $`A`$ is simple and ordinary then the algebra $`(\mathrm{End}A)𝐐`$ is generated by the Frobenius endomorphism of $`A`$, and so is isomorphic to the number field defined by $`f`$. Since the characteristic polynomial of Frobenius of $`A`$ has degree equal to twice the dimension of $`A`$, we see that the degree of the number field $`K=(\mathrm{End}A)𝐐`$ over $`𝐐`$ is twice the dimension of $`A`$. In fact, the number field $`K`$ is a CM-field, which means that $`K`$ is a totally imaginary quadratic extension of a totally real field $`K^+`$. (A number field $`L`$ is totally imaginary if it cannot be embedded into $`𝐑`$, and it is totally real if every embedding of $`L`$ into $`𝐂`$ comes from an embedding of $`L`$ into $`𝐑`$.)
## 3. An easy test for absolute simplicity
In this section we will present an easy-to-verify sufficient condition for a simple abelian variety over a finite field to be absolutely simple. For ordinary varieties, the sufficient condition is also necessary. Throughout this section, $`k`$ will be a finite field, $`\overline{k}`$ its algebraic closure, $`A`$ a simple abelian variety over $`k`$, and $`\pi `$ its Frobenius endomorphism. We let $`\mathrm{End}^0A`$ denote the algebra $`(\mathrm{End}A)𝐐`$. Note that the simplicity of $`A`$ implies that the subalgebra $`𝐐(\pi )`$ of $`\mathrm{End}^0A`$ is a field.
###### Proposition 3.
Let $`D`$ be the set of integers $`d>1`$ such that either
1. the minimal polynomial of $`\pi `$ lies in $`𝐙[x^d]`$ or
2. the field $`𝐐(\pi ^d)`$ is a proper subfield of $`𝐐(\pi )`$ and there is a primitive $`d`$th root of unity $`\zeta `$ in $`𝐐(\pi )`$ such that $`𝐐(\pi )=𝐐(\pi ^d,\zeta )`$.
Then:
1. The set $`D`$ is empty if and only if $`𝐐(\pi ^d)=𝐐(\pi )`$ for all $`d>0`$.
2. If $`𝐐(\pi ^d)=𝐐(\pi )`$ for all $`d>0`$ then $`A`$ is absolutely simple. If $`A`$ is ordinary, then the converse is also true.
To prove this proposition we will need two elementary lemmas.
###### Lemma 4.
Let $`\mathrm{}`$ be a finite extension of $`k`$. If $`𝐐(\pi ^{[\mathrm{}:k]})=𝐐(\pi )`$ then $`A_{\mathrm{}}`$ is simple. If $`A`$ is ordinary, then the converse is also true.
###### Proof.
An abelian variety is simple if and only if its endomorphism ring contains no zero-divisors. Thus, if $`A`$ is simple and $`A_{\mathrm{}}`$ is not, there must exist an element of $`\mathrm{End}^0A_{\mathrm{}}`$ that does not come from $`\mathrm{End}^0A`$. But it follows from the Honda-Tate theorem that $`\mathrm{End}^0A_{\mathrm{}}=\mathrm{End}^0A`$ if $`𝐐(\pi ^{[\mathrm{}:k]})=𝐐(\pi )`$. This proves the first statement of the lemma.
If $`A`$ is ordinary and $`𝐐(\pi ^{[\mathrm{}:k]})`$ is a proper subfield of $`𝐐(\pi )`$, then it follows from the Honda-Tate theorem that $`\mathrm{End}^0A_{\mathrm{}}`$ is a matrix algebra over $`𝐐(\pi ^{[\mathrm{}:k]})`$. In particular, $`\mathrm{End}^0A_{\mathrm{}}`$ contains a zero-divisor, so that $`A_{\mathrm{}}`$ is not simple. ∎
###### Lemma 5.
Let $`\alpha `$ be an algebraic number with minimal polynomial $`g𝐐[x]`$, and suppose that $`d`$ is a positive integer such that the field $`L=𝐐(\alpha ^d)`$ is a proper subfield of $`K=𝐐(\alpha )`$ and such that $`𝐐(\alpha ^r)=K`$ for all positive $`r<d`$. Then either $`g𝐐[x^d]`$ or there is a primitive $`d`$th root of unity $`\zeta `$ in $`K`$ such that $`K=L(\zeta )`$.
###### Proof.
Let $`\zeta `$ be a primitive $`d`$th root of unity in an algebraic closure of $`K`$ and let $`M=L(\zeta )K`$. Note that $`M`$ contains $`L`$. Since $`L(\zeta )`$ is a Galois extension of $`L`$ it is also a Galois extension of $`M`$, and it follows that $`L(\zeta )`$ and $`K`$ are linearly disjoint over $`M`$, so that $`[K(\zeta ):L(\zeta )]=[K:M]`$. Let $`m=[K(\zeta ):L(\zeta )]=[K:M]`$. Since $`K(\zeta )=𝐐(\alpha ,\zeta )`$ is a Kummer extension of $`L(\zeta )=𝐐(\alpha ^d,\zeta )`$, we see that $`\alpha ^m`$ lies in $`L(\zeta )`$, and hence also in $`M`$.
Suppose $`m>1`$. Then since $`𝐐(\alpha ^m)`$ is a subfield of the proper subfield $`M`$ of $`K`$, the lemma’s hypothesis shows we must have $`m=d`$. If we let $`h`$ be the minimal polynomial of $`\alpha ^d`$ over $`𝐐`$, then $`g(x)=h(x^d)`$.
Suppose $`m=1`$. Then $`K(\zeta )=L(\zeta )`$, so $`K/L`$ is a subextension of the abelian extension $`K(\zeta )/L`$, and is therefore Galois. Let $`G`$ be its Galois group, and suppose $`\sigma `$ is a non-identity element of $`G`$. Let $`\xi =\sigma (\alpha )/\alpha `$, so that $`\xi `$ lies in the multiplicative group generated by $`\zeta `$. Suppose $`r`$ is a positive integer less than $`d`$. Then the hypothesis of the lemma shows that $`K=𝐐(\alpha ^r)`$, so we must have $`\alpha ^r\sigma (\alpha ^r)=\xi ^r\alpha ^r`$. Thus $`\xi `$ must in fact be a primitive $`d`$th root of unity, which shows that $`\zeta K`$. It follows that $`K=K(\zeta )`$, and this last field is $`L(\zeta )`$ because $`m=1`$. ∎
###### Proof of Proposition 3.
If $`d`$ is an integer in $`D`$ then clearly $`𝐐(\pi ^d)`$ is a proper subfield of $`𝐐(\pi )`$. On the other hand, if there exists some $`d>0`$ such that $`𝐐(\pi ^d)𝐐(\pi )`$ then there exists a smallest such $`d`$, and by Lemma 5 this $`d`$ lies in $`D`$. This proves the first statement of the proposition.
It is clear that $`A`$ is absolutely simple if and only if $`A_{\mathrm{}}`$ is simple for every finite extension $`\mathrm{}`$ of $`k`$. The second statement of the proposition follows from this fact and from Lemma 4. ∎
###### Remark .
A theorem of Silverberg shows that if $`A`$ is an abelian variety over an arbitrary field $`k`$, then to check that $`\mathrm{End}^0A=\mathrm{End}^0A_{\overline{k}}`$ it suffices to check that $`\mathrm{End}^0A=\mathrm{End}^0A_{\mathrm{}}`$ for a certain finite extension $`\mathrm{}`$ of $`k`$; in particular, if one chooses an integer $`m>2`$ not divisible by the characteristic of $`k`$, Silverberg shows that one may take $`\mathrm{}`$ to be the smallest field over which every $`m`$-torsion point of $`A`$ is defined. The degrees of such $`\mathrm{}`$ over $`k`$ may be quite large, even when $`k`$ is a finite field. Lemmas 5 and the proof of Lemma 4 show that Silverberg’s general result can be improved in the special case where $`k`$ is finite.
## 4. Absolutely simple abelian surfaces
In this section we will prove a theorem that shows that, given the characteristic polynomial of Frobenius of a simple ordinary abelian surface over a finite field, it is quite easy to determine whether the surface is absolutely simple. At the end of the section we will use this theorem to prove the special case $`n=2`$ of Theorem 1.
Suppose $`k`$ is a finite field with $`q`$ elements and $`A`$ is an abelian surface over $`k`$. Let $`f`$ be the characteristic polynomial of Frobenius for $`A`$. Then Weil’s “Riemann Hypothesis” shows that $`f`$ is of the form $`x^4+ax^3+bx^2+aqx+q^2`$ for some integers $`a`$ and $`b`$. If neither $`a`$ nor $`b`$ is coprime to $`q`$ then one can use the Honda-Tate theorem to show that $`A`$ becomes isogenous to the square of a supersingular elliptic curve over a finite extension of $`k`$. If $`a`$ is coprime to $`q`$ but $`b`$ is not, then one can again use Honda-Tate to show that $`A`$ is absolutely simple if and only if it is simple, and that $`A`$ is simple if and only if $`f`$ is irreducible. The most interesting situation arises when $`b`$ is coprime to $`q`$, which is the case exactly when $`A`$ is an ordinary abelian variety. In this case, $`A`$ is simple if and only if $`f`$ is irreducible.
###### Theorem 6.
Suppose $`f=x^4+ax^3+bx^2+aqx+q^2`$ is the Weil polynomial of a simple ordinary abelian surface $`A`$ over $`k`$. Then exactly one of the following conditions holds:
1. The variety $`A`$ is absolutely simple.
2. We have $`a=0`$.
3. We have $`a^2=q+b`$.
4. We have $`a^2=2b`$.
5. We have $`a^2=3b3q`$.
In cases (b), (c), (d), and (e), the smallest extension of $`k`$ over which $`A`$ splits is quadratic, cubic, quartic, and sextic, respectively.
###### Proof.
Let $`\pi `$ be the Frobenius endomorphism of $`A`$ and let $`K`$ be the field $`𝐐(\pi )`$. Because $`A`$ is ordinary, the field $`K`$ is a CM-field of degree $`4`$ over $`𝐐`$. The ordinariness of $`A`$ also implies that $`𝐐(\pi ^d)`$ is a CM-field for every positive integer $`d`$, and that $`A`$ splits over the degree-$`d`$ extension of $`k`$ if and only if $`𝐐(\pi ^d)`$ is a proper subfield of $`K`$.
Suppose $`A`$ is not absolutely simple. Then there is a positive integer $`d`$ such that $`𝐐(\pi ^d)`$ is a proper subfield of $`K`$; let us take $`d`$ to be the smallest such integer, and let $`L`$ be the imaginary quadratic field $`𝐐(\pi ^d)`$. By Lemma 5, either $`d=2`$ and $`f𝐙[x^2]`$, or $`d=4`$ and $`f𝐙[x^4]`$, or there is a primitive $`d`$th root of unity $`\zeta `$ in $`K`$ such that $`K=L(\zeta )`$. Let us first show that if the third possibility is the case and if $`d>4`$ then $`d`$ must equal $`6`$.
Suppose, to obtain a contradiction, that we are in the third case and that $`d`$ is greater than $`4`$ but not equal to $`6`$. Then the degree of $`𝐐(\zeta )`$ over $`𝐐`$ is greater than $`2`$, so $`K`$ must be $`𝐐(\zeta )`$. Let $`\sigma \mathrm{Gal}(K/𝐐)`$ be the nontrivial automorphism of $`K`$ that fixes $`L`$. The proof of Lemma 5 shows that we may choose our primitive root of unity $`\zeta `$ so that $`\pi ^\sigma =\zeta \pi `$. Applying $`\sigma `$ to this equality, we find $`\pi =\zeta ^\sigma \pi ^\sigma =\zeta ^\sigma \zeta \pi `$, so that $`\zeta ^\sigma \zeta =1`$. The only element of the Galois group of the cyclotomic field with this property is complex conjugation. But then the fixed field $`L`$ of $`\sigma `$ must be totally real, and we have reached a contradiction.
So we must find, for $`d=2,3,4,`$ and $`6`$, the conditions on the coefficients $`a`$ and $`b`$ in the minimal polynomial of $`\pi `$ that are equivalent to $`\pi ^d`$ lying in a quadratic subfield. Note that the characteristic polynomial of $`\pi ^d`$ is of the form $`x^4+\alpha x^3+\beta x^2+\alpha q^dx+q^{2d}`$, and that such a quartic polynomial is the square of a quadratic polynomial if and only if $`\alpha ^24\beta +8q^d=0`$. It is not difficult to explicitly calculate the characteristic polynomial of $`\pi ^d`$ for each $`d`$ we are considering, and we find that
$$\alpha ^24\beta +8q^d=\{\begin{array}{cc}a^2(a^24b+8q)\hfill & \text{if }d=2\text{;}\hfill \\ (a^2bq)^2(a^24b+8q)\hfill & \text{if }d=3\text{;}\hfill \\ a^2(a^22b)^2(a^24b+8q)\hfill & \text{if }d=4\text{;}\hfill \\ a^2(a^2bq)^2(a^23b+3q)^2(a^24b+8q)\hfill & \text{if }d=6\text{.}\hfill \end{array}$$
We have assumed that $`A`$ is simple over $`k`$, so the characteristic polynomial for $`\pi `$ is irreducible; this means in particular that the quantity $`a^24b+8q`$ is nonzero. Thus, if $`A`$ is not absolutely simple then one of the cases (b), (c), (d), or (e) must hold. Note that if two of these cases were to hold simultaneously, then $`b`$ would equal a multiple of $`q`$, contradicting our assumption that $`b`$ is coprime to $`q`$. Thus exactly one of the cases (a) through (e) must hold.
Finally, the formulas for $`\alpha ^24\beta +8q^d`$ given above make it easy to verify the theorem’s statement about the degree of the minimal splitting field of $`A`$. ∎
Using Theorem 6, it is easy to show that there exist absolutely simple ordinary abelian surfaces over every finite field. If $`q`$ is an arbitrary prime power, then Theorem 1.1 of shows that the polynomial $`x^4+x^3+x^2+qx+q^2`$ is an ordinary Weil polynomial. It is easy to check that this polynomial is irreducible, so it corresponds to an isogeny class of simple abelian varieties over the field $`𝐅_q`$. Then Theorem 6 shows that the varieties in the isogeny class are absolutely simple.
## 5. The existence of absolutely simple abelian varieties of higher dimension
In this section we will prove Theorem 1 in the case where $`n>2`$. As we noted in the Introduction, it suffices to prove the theorem for finite prime fields $`k`$, but we will assume only that $`k`$ is finite. In fact, for such fields we will prove a result that is slightly stronger than Theorem 1.
###### Theorem 7.
Let $`k`$ be a finite field and let $`n>2`$ be an integer. Then there is an absolutely simple $`n`$-dimensional ordinary abelian variety over $`k`$.
The proof of Theorem 7 depends on three lemmas, whose proofs we will postpone until after the proof of the theorem. The first lemma gives sufficient conditions for an ordinary Weil number to correspond to an isogeny class of absolutely simple varieties.
###### Lemma 8.
Let $`q`$ be a prime power and let $`n>2`$ be an integer. Suppose $`\pi `$ is an ordinary Weil $`q`$-number, let $`K=𝐐(\pi )`$, let $`K^+`$ be the maximal real subfield of $`K`$, and let $`n=[K^+:𝐐]`$. Suppose that
1. the minimal polynomial of $`\pi `$ is not of the form $`x^{2n}+ax^n+q^n`$,
2. the field $`K^+`$ has no proper subfields other than $`𝐐`$, and
3. the field $`K^+`$ is not the maximal real subfield of a cyclotomic field.
Then the isogeny class corresponding to $`\pi `$ consists of absolutely simple varieties.
The second lemma shows that any polynomial satisfying a certain set of local conditions also satisfies the hypotheses of Lemma 8. We will use this lemma again in Section 7.
###### Lemma 9.
Let $`q`$ be a prime power and let $`n>2`$ be an integer. Let $`g𝐙[x]`$ be a monic polynomial of degree $`n`$, and let $`f`$ be the polynomial given by $`f(x)=x^ng(x+q/x)`$. Suppose that the following five conditions hold:
1. the polynomial $`f`$ is not of the form $`x^{2n}+ax^n+q^n`$,
2. all of the complex roots of $`g`$ are real numbers of absolute value less than $`2\sqrt{q}`$,
3. the constant term of $`g`$ is coprime to $`q`$,
4. there exists a prime $`p_1`$ such that the reduction of $`g`$ modulo $`p_1`$ is irreducible, and
5. there exists a prime $`p_2`$ such that the reduction of $`g`$ modulo $`p_2`$ is a linear times an irreducible.
Then $`f`$ is an irreducible ordinary Weil polynomial of degree $`2n`$, and its roots $`\pi `$ satisfy the hypotheses of Lemma 8.
The third lemma gives us a way of producing polynomials that meet the hypotheses of Lemma 9.
###### Lemma 10.
Let $`q`$ be a prime power and let $`n>2`$ be an integer. Then there is a monic polynomial $`g`$ in $`𝐙[x]`$ that satisfies the following five conditions:
1. the polynomial $`g`$ can be written
$$g=x^n+cx^{n2}+\text{lower-order terms},$$
where either $`c`$ is equal to $`2n`$ or $`c`$ is not divisible by $`n`$,
2. all of the complex roots of $`g`$ are real numbers of absolute value less than $`2\sqrt{2}`$,
3. the constant term of $`g`$ is coprime to $`q`$,
4. the reduction of $`g`$ modulo $`2`$ is irreducible, and
5. the reduction of $`g`$ modulo $`3`$ is a linear times an irreducible.
###### Proof of Theorem 7.
Let $`g`$ be the polynomial whose existence is guaranteed by Lemma 10. Then $`g`$ satisfies the last four of the five hypotheses of Lemma 9; we will show that it satisfies the first hypothesis as well.
First we will consider the case in which $`q>2`$. Since
$$g=x^n+cx^{n2}+\text{lower-order terms},$$
we find that the polynomial $`f`$ defined in Lemma 9 may be written in the form
$$f=x^{2n}+(qn+c)x^{2n2}+\text{lower-order terms}.$$
Now, $`c`$ is either $`2n`$ or is not a multiple of $`n`$, so the coefficient of $`x^{2n2}`$ in $`f`$ is nonzero. In particular, $`f`$ is not of the form $`x^{2n}+ax^n+q^n`$.
For the case in which $`q=2`$ we use the easily-proven fact that the reduction of $`f`$ modulo $`2`$ is equal to $`x^n`$ times the reduction of $`g`$ modulo $`2`$. Since $`g`$ modulo $`2`$ is irreducible, and since $`x^n+1`$ is not irreducible over $`𝐅_2`$, the polynomial $`f`$ must have an odd coefficient somewhere between $`x^{2n}`$ and $`x^n`$. Again we see that $`f`$ is not of the form $`x^{2n}+ax^n+q^n`$.
Thus $`g`$ satisfies all the hypotheses of Lemma 9, so by Lemma 8 the roots of $`f`$ are Weil numbers that correspond to an isogeny class of absolutely simple ordinary varieties over $`𝐅_q`$. ∎
###### Proof of Lemma 8.
Suppose, to obtain a contradiction, that $`\pi `$ corresponds to an isogeny class that is not absolutely simple. Then by Proposition 3 there is a positive integer $`d`$ such that $`𝐐(\pi ^d)`$ is a proper subfield of $`K`$. Let $`d`$ be the smallest positive integer with this property. Since $`\pi `$ is ordinary, the field $`L=𝐐(\pi ^d)`$ is a CM-field, and its maximal real subfield $`L^+`$ is a proper subfield of $`K^+`$. Hypothesis (2) shows that $`L^+`$ must be $`𝐐`$, so $`L`$ is an imaginary quadratic field.
Lemma 5 shows that either the minimal polynomial $`f`$ of $`\pi `$ lies in $`𝐙[x^d]`$ or $`K=L(\zeta )`$ for some primitive $`d`$th root of unity. The first possibility cannot occur, because it would imply that $`d=n`$, contradicting hypothesis (1). Therefore the second possibility must be the case. We find that the maximal real subfield of $`𝐐(\zeta )`$ is a subfield of $`K^+`$, and since $`K^+`$ is not itself the maximal real subfield of a cyclotomic field (by assumption), we find that the maximal real subfield of $`𝐐(\zeta )`$ must be $`𝐐`$, so that $`𝐐(\zeta )`$ is either a quadratic field or $`𝐐`$ itself. But $`K`$ is the compositum of $`L`$ and $`𝐐(\zeta )`$, so the degree of $`K`$ over $`𝐐`$ is at most $`4`$. This contradicts our assumption that the degree of $`K`$ over $`𝐐`$ is $`2n`$, where $`n>2.`$
###### Proof of Lemma 9.
Since $`g`$ modulo $`p_1`$ is irreducible, $`g`$ itself is irreducible in $`𝐙[x]`$, and since all of its complex roots are real, $`g`$ defines a totally real number field $`K^+`$. Let $`\alpha `$ be a root of $`g`$ in $`K^+`$. The discriminant of the polynomial $`h=x^2\alpha x+q`$ is totally negative because the roots of $`g`$ all have magnitude less than $`2\sqrt{q}`$, so $`h`$ defines a totally imaginary quadratic extension $`K`$ of $`K^+`$. If $`\pi `$ is a root of $`h`$ in $`K`$, then $`K=𝐐(\pi )`$ contains $`K^+`$ because $`\alpha =\pi +q/\pi `$. Thus $`\pi `$ is an algebraic number of degree $`2n`$. Furthermore, if $`\phi `$ is an embedding of $`K`$ into $`𝐂`$, then $`\phi (\pi )`$ is a root of $`x^2\phi (\alpha )x+q`$, and the quadratic formula shows that $`|\phi (\pi )|=\sqrt{q}`$. Thus $`\pi `$ is in fact a Weil number of degree $`2n`$. Since $`\pi `$ is a root of $`f`$, the polynomial $`f`$ must be the minimal polynomial of $`\pi `$. This shows that $`f`$ is an irreducible polynomial whose roots are Weil numbers, and to show that $`f`$ is an ordinary Weil polynomial we need merely check that its middle coefficient is coprime to $`q`$. But this follows from hypothesis (3), because the middle coefficient of $`f`$ differs from the constant term of $`g`$ by a multiple of $`q`$.
Now we must check that a root $`\pi `$ of $`f`$ satisfies the hypotheses of Lemma 8. The first of these hypotheses is identical to the first hypothesis of the lemma we are proving, and is therefore satisfied.
We will show that $`K^+`$ is not the maximal real subfield of a cyclotomic field. It will suffice to show that $`K^+`$ is not Galois over $`𝐐`$. The defining polynomial $`g`$ of $`K^+`$ reduces modulo $`p_2`$ as a linear times an irreducible, so the prime $`p_2`$ splits in $`K^+`$ into two primes with different residue class degrees, so $`K^+/𝐐`$ cannot be Galois.
Finally, we prove that $`K^+`$ has no proper subfields other than $`𝐐`$. For suppose $`K^+`$ had a proper subfield $`L`$ other than $`𝐐`$. Let $`𝔓`$ be the prime of $`K^+`$ over $`p_2`$ whose residue class degree is $`n1`$. Let $`𝔭`$ be the prime of $`L`$ lying under $`𝔓`$. Let $`f_1`$ be the residue class degree of $`𝔓`$ over $`𝔭`$ and let $`f_2`$ be the residue class degree of $`𝔭`$ over $`p_2`$. Then we have the three statements:
1. $`f_1[K^+:L]`$,
2. $`f_2[L:𝐐]`$, and
3. $`f_1f_2=n1=1+[K^+:L][L:𝐐].`$
Statement (c) shows that strict inequality must hold in one of statements (a) and (b); but then, since both of the field extensions $`K^+/L`$ and $`L/𝐐`$ are assumed non-trivial, we find that $`f_1f_2`$ must be less than $`1+[K^+:L][L:𝐐]`$. This contradiction shows that $`K^+`$ has no proper subfields other than $`𝐐`$. ∎
Our proof of Lemma 10 depends on a result of Robinson concerning certain modified Chebyshev polynomials. Before starting on the proof of the lemma we will define these polynomials and present Robinson’s result.
For every positive integer $`i`$ let $`t_i`$ be the $`i`$th Chebyshev polynomial, so that $`t_i(x)=\mathrm{cos}(i\mathrm{arccos}(x))`$. For every positive integer $`i`$ let $`T_i`$ be the polynomial given by $`T_i(x)=22^{i/2}t_i(x/2^{3/2})`$. It is not hard to show that $`T_i`$ is a monic polynomial in $`𝐙[x]`$ and that $`T_ix^imod2`$. Let $`T_0=1`$.
###### Lemma 11.
Suppose $`a_1,\mathrm{},a_n`$ are real numbers such that
$$\left(\underset{i=1}{\overset{n1}{}}\left|\frac{a_i}{2^{i/2}}\right|\right)+\frac{1}{2}\left|\frac{a_n}{2^{n/2}}\right|<1.$$
Then every complex root of the polynomial
$$T_n+a_1T_{n1}+\mathrm{}+a_{n1}T_1+a_nT_0$$
is real and lies in the open interval $`(2\sqrt{2},2\sqrt{2})`$.
###### Proof.
This follows from the techniques of Robinson . ∎
###### Proof of Lemma 10.
If $`n9`$ we can simply choose the appropriate value of $`g`$ from Table 1, so let us assume that $`n>9`$.
Lemma 12 (below) shows that there exist monic degree-$`n`$ polynomials $`g_2`$ in $`𝐅_2[x]`$ and $`g_3`$ in $`𝐅_3[x]`$ such that $`g_2`$ is irreducible, such that $`g_3`$ is a linear times an irreducible and has nonzero constant term, and such that the coefficients of $`x^{n1},\mathrm{},x^{n6}`$ in $`g_2`$ and $`g_3`$ are equal to the reductions (modulo $`2`$ and $`3`$) of the corresponding coefficients of the modified Chebyshev polynomial $`T_n`$. Once we have fixed $`g_2`$ and $`g_3`$, we can choose values of $`a_7,a_8,\mathrm{},a_n`$ in the set $`\{2,1,0,1,2,3\}`$ such that the polynomial
$$g=T_n+a_7T_{n7}+a_8T_{n8}+\mathrm{}+a_{n1}T_1+a_nT_0$$
reduces to $`g_2`$ modulo $`2`$ and to $`g_3`$ modulo $`3`$.
Note that the constant term of $`g`$ is coprime to $`6`$ because $`g_2`$ and $`g_3`$ have nonzero constant terms. Thus, if $`q`$ is a power of $`2`$ or $`3`$ then the constant term of $`g`$ is coprime to $`q`$. If $`q`$ is not a power of $`2`$ or $`3`$, then the constant term of $`g`$ may have a factor in common with $`q`$. If this is the case, replace $`a_n`$ with either $`a_n6`$ or $`a_n+6`$, whichever one lies in the interval $`[6,6]`$; this changes the constant term of $`g`$ by $`6`$, so that the constant term is now coprime to $`q`$ but so that $`g`$ still reduces to $`g_2`$ modulo $`2`$ and to $`g_3`$ modulo $`3`$.
One can calculate that $`T_n=x^n2nx^{n2}+\text{lower-order terms},`$ and since $`g`$ differs from $`T_n`$ by a polynomial of degree at most $`n7`$, we see that $`g`$ may also be written $`g=x^n2nx^{n2}+\text{lower-order terms}.`$ In particular, $`g`$ satisfies the first condition of Lemma 10.
Thus $`g`$ satisfies four of the five conditions listed in the statement of Lemma 10. We are left to show that all of its roots are real, and that they have absolute value less than $`2\sqrt{2}`$. But this follows from Lemma 11; to apply the lemma we must verify that the quantity
$$\left|\frac{a_7}{2^{7/2}}\right|+\left|\frac{a_8}{2^{8/2}}\right|+\mathrm{}+\left|\frac{a_{n1}}{2^{(n1)/2}}\right|+\frac{1}{2}\left|\frac{a_n}{2^{n/2}}\right|$$
is less than $`1`$, and this follows from the fact that $`|a_i|`$ is at most $`3`$ for $`i<n`$, and that $`|a_n|`$ is at most $`6`$.
Thus, the $`g`$ we have written down satisfies all the conditions of the lemma. ∎
###### Lemma 12.
Suppose $`n10`$. Then there exist monic degree-$`n`$ polynomials $`g_2`$ in $`𝐅_2[x]`$ and $`g_3`$ in $`𝐅_3[x]`$ such that $`g_2`$ is irreducible, such that $`g_3`$ is a linear times an irreducible and has nonzero constant term, and such that the coefficients of $`x^{n1}`$ through $`x^{n6}`$ of $`g_2`$ and $`g_3`$ are equal to the reductions modulo $`2`$ and $`3`$ of the corresponding coefficients of the modified Chebyshev polynomial $`T_n`$ defined above.
###### Proof.
For $`n18`$ we choose $`g_2`$ and $`g_3`$ from Table 2. For $`n>18`$ we argue as follows:
Corollary 3.2 (p. 94) of shows that there exists a monic irreducible polynomial in $`𝐅_2[x]`$ of degree $`n`$ with zeroes for the first six coefficients after the leading $`x^n`$. We take this polynomial for our $`g_2`$. The same corollary shows that there is a monic irreducible polynomial $`h`$ in $`𝐅_3[x]`$ such that the first six coefficients of $`(x1)h`$ are equal to those of the reduction of $`T_n`$ modulo $`3`$; we take $`g_3`$ to be $`(x1)h`$. ∎
## 6. Asymptotic results for abelian surfaces
In this section we will prove Theorem 2 in the case $`n=2`$. In fact, we will prove a more precise statement.
###### Theorem 13.
Let $`ϵ`$ be a positive real number. If $`q`$ is a prime power with $`q>(659/ϵ)^2`$ then $`S(𝐅_q,2)>1ϵ`$.
###### Proof.
Let $`r`$ be the arithmetic function defined by $`r(x)=\phi (x)/x`$, where $`\phi `$ is Euler’s $`\phi `$-function, let $`I`$ be the number of isogeny classes of abelian surfaces over $`𝐅_q`$, let $`O_{\mathrm{simple}}`$ be the number of isogeny classes of simple ordinary abelian surfaces, and let $`O_{\mathrm{abs}.\mathrm{simple}}`$ be the number of isogeny classes of absolutely simple ordinary abelian surfaces. Theorem 1.2 of shows that
$$I<\frac{32}{3}r(q)q^{3/2}+3473q+8359q^{1/2};$$
this upper bound is obtained by combining the estimates that Theorem 1.2 gives for the number of ordinary and non-ordinary isogeny classes of abelian surfaces.
The same theorem shows that the number of isogeny classes of ordinary abelian surfaces over $`𝐅_q`$ is at least
$$\frac{32}{3}r(q)q^{3/2}8359q^{1/2}.$$
The isogeny classes of ordinary elliptic curves over $`𝐅_q`$ correspond to the integers $`t`$ such that $`|t|<2q^{1/2}`$ and $`(t,q)=1`$, so there are at most $`4q^{1/2}`$ such isogeny classes. A non-simple isogeny class of ordinary abelian surfaces is determined by its two factors, so there are at most $`4q^{1/2}(4q^{1/2}+1)/2=8q+2q^{1/2}`$ such reducible isogeny classes. Thus we have
$$O_{\mathrm{simple}}>\frac{32}{3}r(q)q^{3/2}8q8361q^{1/2}.$$
Now we must estimate the number of simple ordinary isogeny classes that are not absolutely simple. For this we use Theorem 6. First note that if $`x^4+ax^3+bx^2+aqx+q^2`$ is the Weil polynomial for an ordinary abelian surface over $`𝐅_q`$ then $`|a|<4q^{1/2}`$, and if $`a=0`$ then $`0<|b|<2q.`$ Thus, the number of Weil polynomials of ordinary abelian surfaces that satisfy case (b) of Theorem 6 is at most $`4q`$. Also, for every nonzero integer $`d`$ in the interval $`(4q^{1/2},4q^{1/2})`$ there is at most one Weil polynomial with $`a=d`$ that satisfies case (c) of the theorem; for every nonzero integer $`d`$ in the interval $`(2q^{1/2},2q^{1/2})`$ there is at most one Weil polynomial with $`a=2d`$ that satisfies case (d) of the theorem; and for every nonzero integer $`d`$ in the interval $`((4/3)q^{1/2},(4/3)q^{1/2})`$ there is at most one Weil polynomial with $`a=3d`$ that satisfies case (e) of the theorem. We find that there are at most $`15q^{1/2}`$ simple Weil polynomials $`x^4+ax^3+bx^2+aqx+q^2`$ with $`a0`$ that are not absolutely simple.
Combining these estimates with the lower bound for $`O_{\mathrm{simple}}`$ given above, we find that
$$O_{\mathrm{abs}.\mathrm{simple}}>\frac{32}{3}r(q)q^{3/2}12q8376q^{1/2}.$$
Now suppose $`ϵ`$ is given. If $`ϵ1`$ then the conclusion of the theorem is clearly true for all $`q`$, so we may assume that $`ϵ<1`$ and that $`q>659^2`$. With this lower bound for $`q`$, our bounds for $`I`$ and $`O_{\mathrm{abs}.\mathrm{simple}}`$ show that
$$I<\frac{32}{3}r(q)q^{3/2}+3486q$$
and
$$O_{\mathrm{abs}.\mathrm{simple}}>\frac{32}{3}r(q)q^{3/2}25q.$$
Thus
$$\frac{O_{\mathrm{abs}.\mathrm{simple}}}{I}>\left(1\frac{75}{32r(q)q^{1/2}}\right)/\left(1+\frac{10458}{32r(q)q^{1/2}}\right).$$
The denominator is less than $`2`$, so we have
$`{\displaystyle \frac{O_{\mathrm{abs}.\mathrm{simple}}}{I}}`$ $`>\left(1{\displaystyle \frac{75}{32r(q)q^{1/2}}}\right)\left(1{\displaystyle \frac{10458}{32r(q)q^{1/2}}}\right)`$
$`>1{\displaystyle \frac{10533}{32r(q)q^{1/2}}}`$
$`>1{\displaystyle \frac{10533}{16q^{1/2}}}>1{\displaystyle \frac{659}{q^{1/2}}}>1ϵ,`$
as was to be shown. ∎
## 7. Asymptotic results for abelian varieties of higher dimension
In this section we will prove Theorem 2 in the case $`n>2`$ by proving a more precise result, whose statement requires us to introduce some notation. First we define constants $`c_1,c_2,`$ and $`c_3`$ by setting
$`c_1`$ $`=\sqrt{3}/60.288675,`$
$`c_2`$ $`=\mathrm{exp}(3/2)2(1+\sqrt{2})\sqrt{3}(1+\sqrt{3}/162)^3/312.898608,`$
and
$`c_3`$ $`=c_2/(1+\sqrt{2})5.342778.`$
Next, for every positive integer $`n`$ we let
$$v_n=\frac{2^n}{n!}\underset{j=1}{\overset{n}{}}\left(\frac{2j}{2j1}\right)^{n+1j}$$
and we let
$$G_n=\frac{1}{v_n}6^{n^2}c_1^nc_3\frac{n(n+1)}{(n1)!}.$$
Finally, if $`n>1`$ is an integer and if $`ϵ`$ is a positive real, we let $`k_{n,ϵ}`$ denote the smallest positive integer $`k`$ such that
$$\left(1\frac{1}{2n}\right)^k\frac{ϵ}{8},$$
we let $`m_{n,ϵ}`$ be the product of the first $`k_{n,ϵ}`$ prime numbers, and we let
$$M_{n,ϵ}=\left(\frac{8G_nm_{n,ϵ}}{ϵ}\right)^2.$$
Recall that $`S(𝐅_q,n)`$ denotes the fraction of isogeny classes of $`n`$-dimensional abelian varieties over $`𝐅_q`$ that are ordinary and absolutely simple.
###### Theorem 14.
Let $`n>2`$ be an integer and let $`ϵ`$ be a positive real number. If $`q>M_{n,ϵ}`$ then $`S(𝐅_q,n)>1ϵ.`$
For every prime power $`q`$ and non-negative integer $`n`$ we let $`(q,n)`$ denote the set of isogeny classes of $`n`$-dimensional abelian varieties over $`𝐅_q`$ and we let $`𝒪(q,n)`$ and $`𝒩(q,n)`$ denote the sets of ordinary and non-ordinary isogeny classes in $`(q,n)`$, respectively. Also, we let $`𝒪_{\mathrm{simple}}(q,n)`$ and $`𝒪_{\mathrm{abs}.\mathrm{simple}}(q,n)`$ denote the sets of simple and absolutely simple isogeny classes in $`𝒪(q,n)`$, respectively. As in Section 6 we let $`r`$ be the arithmetic function defined by $`r(x)=\phi (x)/x`$, where $`\phi `$ is Euler’s $`\phi `$-function. Our proof of Theorem 14 breaks into two parts. First we will give an upper bound for $`\mathrm{\#}(q,n)`$.
###### Proposition 15.
Let $`n>2`$ be an integer and let $`ϵ`$ be a positive real number with $`ϵ1`$. If $`q>M_{n,ϵ}`$ then $`\mathrm{\#}(q,n)<(1+ϵ/8)v_nr(q)q^{n(n+1)/4}.`$
Then we will give a lower bound for $`\mathrm{\#}𝒪_{\mathrm{abs}.\mathrm{simple}}(q,n)`$.
###### Proposition 16.
Let $`n>2`$ be an integer and let $`ϵ`$ be a positive real number with $`ϵ1`$. If $`q>M_{n,ϵ}`$ then $`\mathrm{\#}𝒪_{\mathrm{abs}.\mathrm{simple}}(q,n)(17ϵ/8)v_nr(q)q^{n(n+1)/4}.`$
Clearly these two propositions provide a proof of Theorem 14.
###### Proof of Proposition 15.
Combining the estimates for $`\mathrm{\#}𝒪(q,n)`$ and $`\mathrm{\#}𝒩(q,n)`$ given in Theorem 1.2 of , we find that the quantity $`\mathrm{\#}(q,n)v_nr(q)q^{n(n+1)/4}`$ is less than or equal to
$$6^{n^2}c_1^nc_2\frac{n(n+1)}{(n1)!}q^{n(n1)/4}+\left(v_n+6^{n^2}c_1^nc_3\frac{n(n+1)}{(n1)!}\right)q^{(n+2)(n1)/4},$$
so
$$\frac{\mathrm{\#}(q,n)}{v_nr(q)q^{n(n+1)/4}}1+\frac{c_2G_n}{c_3r(q)q^{n/2}}+\frac{1}{r(q)q^{1/2}}+\frac{G_n}{r(q)q^{1/2}}.$$
An easy induction shows that $`G_n>2`$, and certainly $`c_2/c_3<2.5`$, so we have
$`{\displaystyle \frac{\mathrm{\#}(q,n)}{v_nr(q)q^{n(n+1)/4}}}`$ $`<1+{\displaystyle \frac{G_n}{r(q)q^{1/2}}}\left({\displaystyle \frac{c_2}{c_3}}+{\displaystyle \frac{1}{2}}+1\right)`$
$`<1+{\displaystyle \frac{4G_n}{r(q)q^{1/2}}}.`$
Since $`q>M_{n,ϵ}`$ we have $`q^{1/2}>8G_nm_{n,ϵ}/ϵ`$, and combining this with the fact that $`r(q)1/2`$ we find that
$$\frac{\mathrm{\#}(q,n)}{v_nr(q)q^{n(n+1)/4}}1+\frac{ϵ}{m_{n,ϵ}}.$$
But $`m_{n,ϵ}`$ is greater than $`8`$ for $`ϵ1`$, so the right-hand side is at most $`1+ϵ/8`$. This proves the inequality of the proposition. ∎
Our proof of Proposition 16 is based upon Lemmas 8 and 9. We will compute
1. a lower bound on the number of degree-$`n`$ polynomials satisfying hypotheses (2), (4), and (5) of Lemma 9,
2. an upper bound on the number of degree-$`n`$ polynomials satisfying hypothesis (2) but failing hypothesis (3) of Lemma 9, and
3. an upper bound on the number of degree-$`n`$ polynomials satisfying hypothesis (2) but failing hypothesis (1) of Lemma 9.
Subtracting the sum of the latter two estimates from the first estimate will give us a lower bound on the number of degree-$`n`$ polynomials satisfying all the hypotheses of Lemma 9. By Lemma 8, this lower bound will also be a lower bound on $`\mathrm{\#}𝒪_{\mathrm{abs}.\mathrm{simple}}(q,n)`$. The computation of the lower bound on the number of polynomials satisfying hypotheses (2), (4), and (5) of Lemma 9 will depend on the following lemma, whose proof we will postpone until the next section.
###### Lemma 17.
Let $`n>2`$ be a positive integer, let $`ϵ`$ be a real number between $`0`$ and $`1`$, and let $`m=m_{n,ϵ}`$ be as defined at the beginning of this section. Then there are at least $`m^n(1ϵ/4)`$ monic degree-$`n`$ polynomials in $`(𝐙/m𝐙)[x]`$ such that
1. there exists a prime divisor $`p_1`$ of $`m`$ such that the reduction of $`g`$ modulo $`p_1`$ is irreducible, and
2. there exists a prime divisor $`p_2`$ of $`m`$ such that the reduction of $`g`$ modulo $`p_2`$ is a linear times an irreducible.
Before we proceed to the proof of Proposition 16 we should mention a basic correspondence that we will use repeatedly in our argument. Fix our prime power $`q`$. Suppose $`g`$ is a monic polynomial of degree $`n`$ with integer coefficients, say
$$g=x^n+b_1x^{n1}+\mathrm{}+b_n,$$
and let $`f`$ be the polynomial defined by $`f(x)=x^ng(x+q/x)`$, so that
$$f=(x^{2n}+q^n)+a_1(x^{2n1}+q^{n1}x)+\mathrm{}+a_{n1}(x^{n+1}+qx^{n1})+a_nx^n$$
for some integers $`a_i`$. Then the linear map $`\mathrm{\Omega }`$ from $`𝐙^n`$ to $`𝐙^n`$ that sends a vector $`𝐛=(b_1,\mathrm{},b_n)`$ to the vector $`𝐚=(a_1,\mathrm{},a_n)`$ is invertible — in fact, it is represented by a matrix with integer entries that is lower-triangular with $`1`$’s on the diagonal. Thus, if we let $`𝐛`$ range over a set of vectors that reduces modulo some integer $`m`$ to the entire set $`(𝐙/m𝐙)^n`$, then $`\mathrm{\Omega }(𝐛)`$ will also range over such a set, and conversely, if $`𝐚`$ ranges over such a set, then so will $`\mathrm{\Omega }^1(𝐚)`$.
Note that if $`g`$ and $`f`$ are related as above, then $`g`$ satisfies hypothesis (2) of Lemma 9 if and only if the roots of $`f`$ in the complex numbers all have magnitude $`q^{1/2}`$ and the roots of $`f`$ in the real numbers all have even multiplicity. Also, the roots of $`f`$ meet this last condition if and only if the vector $`(a_1q^{1/2},a_2q^1,\mathrm{},a_nq^{n/2})`$ lies in the region $`V_n`$ of $`𝐑^n`$ defined in . Thus we will be interested in estimating the sizes of the intersections of certain lattices with $`V_n`$.
Let $`𝐞_1`$, …, $`𝐞_n`$ denote the standard basis vectors of $`𝐑^n`$. Our arguments will involve two lattices in $`𝐑^n`$: The first lattice, denoted $`\mathrm{\Lambda }`$, is generated by the vectors $`q^{i/2}𝐞_i`$, and the second, denoted $`\mathrm{\Lambda }^{}`$, is generated by the same set of vectors, except with $`q^{n/2}𝐞_n`$ replaced with $`pq^{n/2}𝐞_n`$. Thus $`\mathrm{\Lambda }\mathrm{\Lambda }^{}`$.
###### Proof of Proposition 16.
Let $`m=m_{n,ϵ}`$ and let $`\mathrm{\Lambda }^{\prime \prime }`$ denote the lattice $`m\mathrm{\Lambda }`$. If $`\mathrm{}`$ is a point in $`\mathrm{\Lambda }^{\prime \prime }`$ let $`B_{\mathrm{}}`$ denote the “brick”
$$\mathrm{}+\{(x_1,\mathrm{},x_n)𝐑^ni:0x_i<mq^{i/2}𝐞_i\}.$$
Let $`S`$ denote the set of all $`\mathrm{}\mathrm{\Lambda }^{\prime \prime }`$ such that $`B_{\mathrm{}}V_n`$. The proof of Proposition 2.3.1 of (see especially p. 435) shows that
$$\mathrm{\#}S\frac{\text{volume }V_n}{\text{covolume }\mathrm{\Lambda }^{\prime \prime }}6^{n^2}c_1^nc_3\frac{n(n+1)}{(n1)!}\frac{d}{\text{covolume }\mathrm{\Lambda }^{\prime \prime }}$$
where $`d`$ is the mesh of $`\mathrm{\Lambda }^{\prime \prime }`$ (see p. 434 of ), which is $`mq^{1/2}`$. Since the covolume of $`\mathrm{\Lambda }^{\prime \prime }`$ is $`m^nq^{n(n+1)/4}`$, we find that
$$\mathrm{\#}Sm^nv_nq^{n(n+1)/4}6^{n^2}c_1^nc_3\frac{n(n+1)}{(n1)!}m^{n+1}q^{(n^2+n2)/4}.$$
Thus
$$m^n\mathrm{\#}Sv_nq^{n(n+1)/4}(1mG_nq^{1/2}),$$
and using the fact that $`q>M_{n,ϵ}`$ we find that
$$m^n\mathrm{\#}Sv_nq^{n(n+1)/4}(1ϵ/8).$$
Now suppose $`\mathrm{}`$ is a lattice point in $`S`$, and consider a typical element $`𝐱=(a_1q^{1/2},a_2q^1,\mathrm{},a_nq^{n/2})`$ of $`\mathrm{\Lambda }B_{\mathrm{}}`$. As $`𝐱`$ ranges over all of $`\mathrm{\Lambda }B_{\mathrm{}}`$, the vector $`𝐚=(a_1,\mathrm{},a_n)`$ ranges over a set of $`m^n`$ elements of $`𝐙^n`$ that reduces modulo $`m`$ to all of $`(𝐙/m𝐙)^n`$. Lemma 17 above shows that of the $`m^n`$ polynomials $`g`$ we obtain from the vectors $`\mathrm{\Omega }(𝐚)`$ arising from elements of $`\mathrm{\Lambda }B_{\mathrm{}}`$, at least $`m^n(1ϵ/4)`$ satisfy hypotheses (4) and (5) of Lemma 9. So for each element of $`S`$ we obtain at least $`m^n(1ϵ/4)`$ polynomials satisfying hypotheses (2), (4), and (5) of Lemma 9. Thus the total number of such polynomials is at least $`m^n\mathrm{\#}S(1ϵ/4)`$, and by the results of the preceding paragraph this number is at least
$$v_nq^{n(n+1)/4}(1ϵ/4)(1ϵ/8),$$
which is greater than $`v_nq^{n(n+1)/4}(13ϵ/8).`$
Next we estimate the number of polynomials $`g`$ that satisfy hypothesis (2) of Lemma 9 but that fail to satisfy hypothesis (3). There is a bijection between the set of such polynomials and the set $`\mathrm{\Lambda }^{}V_n`$, and Proposition 2.3.1 of gives upper and lower bounds for the size of the latter set; in particular, we find that the number of such polynomials differs from $`(1/p)v_nq^{n(n+1)/4}`$ by at most
$$\frac{q^{n(n+1)/4}}{pq^{1/2}}6^{n^2}c_1^nc_3\frac{n(n+1)}{(n1)!},$$
which is $`v_nq^{n(n+1)/4}G_n/pq^{1/2}`$. Since $`q`$ is at least $`M_{n,ϵ}`$, this last quantity is at most $`v_nq^{n(n+1)/4}ϵ/(4pm)`$, which is less than $`v_nq^{n(n+1)/4}ϵ/32`$, because $`m>4`$ when $`ϵ1`$. Thus the number of polynomials that satisfy hypothesis (2) but not hypothesis (3) is at most
$$v_nq^{n(n+1)/4}\left(\frac{1}{p}+\frac{ϵ}{32}\right).$$
Finally we estimate the number of polynomials $`g`$ that satisfy hypothesis (2) of Lemma 9 but that fail to satisfy hypothesis (1). Now, a polynomial $`x^{2n}+ax^n+q^n`$ has all of its roots on the circle $`|z|=q^{1/2}`$ if and only if $`|a|2q^{n/2}`$, so there are at most $`4q^{n/2}+1`$ polynomials meeting hypothesis (2) but failing hypothesis (1). It is very easy to show that $`4q^{n/2}+1<v_nq^{n(n+1)/4}ϵ/32`$ when $`q>M_{n,ϵ}`$.
Now, the number of polynomials meeting all five hypotheses of Lemma 9 is at least as large as the number that satisfy hypotheses (2), (4), and (5), less the number that satisfy hypothesis (2) but that fail either hypothesis (1) or hypothesis (3). We find that the number of polynomials meeting all five hypotheses is at least
$$\begin{array}{c}v_nq^{n(n+1)/4}\left(1\frac{3ϵ}{8}\right)v_nq^{n(n+1)/4}\left(\frac{1}{p}+\frac{ϵ}{32}\right)v_nq^{n(n+1)/4}\frac{ϵ}{32}\hfill \\ \hfill =v_nq^{n(n+1)/4}\left(r(q)\frac{7ϵ}{16}\right)v_nr(q)q^{n(n+1)/4}\left(1\frac{7ϵ}{8}\right)\end{array}$$
and this is the statement of Proposition 16. ∎
## 8. Proof of Lemma 17
In this section we will prove Lemma 17. We continue to use the notation set at the beginning of Section 7.
For the moment, let us write $`A_{n,p}`$ for the set of monic degree-$`n`$ irreducible polynomials in $`𝐅_p[x]`$ and $`B_{n,p}`$ for the set of monic degree-$`n`$ polynomials in $`𝐅_p[x]`$ that factor as a linear polynomial times an irreducible.
###### Lemma 18.
Let $`p`$ be a prime. For all $`n>0`$ we have $`\mathrm{\#}A_{n,p}p^n/(2n)`$, and for all $`n>1`$ we have $`\mathrm{\#}B_{n,p}p^n/(2n2)`$.
###### Proof.
The lemma follows easily from the well-known exact formula
$$\mathrm{\#}A_{n,p}=\frac{1}{n}\underset{dn}{}p^d\mu \left(\frac{n}{d}\right),$$
where $`\mu `$ is the Möbius function. ∎
Suppose $`n>1`$. We see from Lemma 18 that if we choose a monic degree-$`n`$ polynomial $`f`$ at random from $`𝐅_p[x]`$ (with the uniform distribution), then
$`\mathrm{Prob}(fA_{n,p})`$ $`1{\displaystyle \frac{1}{2n}}`$
and
$`\mathrm{Prob}(fB_{n,p})`$ $`1{\displaystyle \frac{1}{2n2}}.`$
Suppose that $`ϵ`$ is given, with $`0<ϵ<1`$. Let $`k=k_{n,ϵ}`$ and $`m=m_{n,ϵ}`$ be as at the beginning of Section 7, so that $`m`$ is the product of the first $`k`$ prime numbers. Now suppose we choose a monic degree-$`n`$ polynomial $`f`$ at random from $`(𝐙/m𝐙)[x]`$. By the Chinese remainder theorem, making such a choice is equivalent to choosing a monic degree-$`n`$ polynomial $`f`$ at random from $`𝐅_p[x]`$ for each of the first $`k`$ primes $`p`$. Thus we see that
$`\mathrm{Prob}(pm:(fmodp)A_{n,p})`$ $`\left(1{\displaystyle \frac{1}{2n}}\right)^k`$
$`\mathrm{Prob}(pm:(fmodp)B_{n,p})`$ $`\left(1{\displaystyle \frac{1}{2n2}}\right)^k,`$
and it follows that
$$\begin{array}{c}\mathrm{Prob}(p_1,p_2m:(fmodp_1)A_{n,p_1}\text{ and }(fmodp_2)B_{n,p_2})\hfill \\ \hfill >1\left(1\frac{1}{2n}\right)^k\left(1\frac{1}{2n2}\right)^k>12\left(1\frac{1}{2n}\right)^k.\end{array}$$
But the definition of $`k_{n,ϵ}`$ shows that
$$\left(1\frac{1}{2n}\right)^k\frac{ϵ}{8},$$
so
$$\mathrm{Prob}(p_1,p_2m:(fmodp_1)A_{n,p_1}\text{ and }(fmodp_2)B_{n,p_2})>1ϵ/4.$$
Thus the number of monic degree-$`n`$ polynomials in $`(𝐙/m𝐙)[x]`$ that satisfy the two conditions of Lemma 17 is at least $`m^n(1ϵ/4)`$, as was to be shown. |
warning/0002/hep-ex0002039.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Symmetries are one of the most fundamental concepts for understanding the laws of nature leading to conserving quantities. Unexpected violations of symmetries indicate some dynamical mechanism beyond the current understanding of physics.
Parity violation was discovered in 1957 in nuclear $`\beta `$ decays and pion and muon decays . In the charged current interaction of the standard electroweak theory, parity and charge conjugation symmetries are maximally violated due to the $`VA`$ structure . All the experimental results up to now are in full agreement with the theory.
A surprising discovery of the CP violating $`\mathrm{K}_\mathrm{L}\pi ^+\pi ^{}`$ decays was made in 1964. The neutral kaon system still remains to be the only place CP violation has been seen. The Standard Model with three Fermion families can accommodate all the observed CP violation phenomena through the complex quark mixing matrix, Cabibbo-Kobayashi-Maskawa (CKM) matrix . However, no real precision test has been made due to the large uncertainties in evaluating the effect of hadronic interactions.
Interest in CP violation is not limited to elementary particle physics. It is one of the three necessary ingredients to generate observed excess of matter over antimatter in the universe . The amount of CP violation which can be generated by the Standard Model appears to be insufficient for explaining the observed matter-antimatter asymmetry in the universe , giving a strong motivation to search for new physics.
For CP violation in some B meson decay channels, the Standard Model can make precise predictions with little influence from the strong interactions. Those channels can be used to test the predictions quantitatively to look for a sign of new physics. Also in the B meson system, CP violation is expected in many decay modes. The pattern of CP violation allows us to make a systematic qualitative comparison with the Standard Model predictions. Therefore, it is now widely accepted that the B-meson system provides in future an ideal place for testing the Standard Model for CP violation .
In this article, we first derive the formalism describing the particle antiparticle system, with and without CP violation. Three different mechanisms which can generate CP violation are clearly classified, together with experimental observables which identify contributions from the different mechanisms. Then, CP violation in the neutral kaon system is analysed in this formalism. After a brief discussion on the Standard Model description for CP violation in the neutral kaon system, we proceed to the neutral B meson system. Following the discussion on some Standard Model predictions, some thoughts are made how the situation could change if there exists new physics contributing in the B meson system.
## 2 Description of Particle Antiparticle System
### 2.1 Basic Formalism
Let $`|\mathrm{P}^0`$ and $`|\overline{\mathrm{P}}{}_{}{}^{0}`$ be the states of a neutral pseudoscalar particle $`\mathrm{P}^0`$-meson and its antiparticle $`\overline{\mathrm{P}}^0`$-meson at rest, respectively. They have definite flavour quantum numbers with opposite signs: $`F=+1`$ for $`\mathrm{P}^0`$ and $`F=1`$ for $`\overline{\mathrm{P}}^0`$. Both states are eigenstates of the strong and electromagnetic interaction Hamiltonian, i.e.
$$(H_{\mathrm{st}}+H_{\mathrm{em}})|\mathrm{P}^0=m_0|\mathrm{P}^0\mathrm{and}(H_{\mathrm{st}}+H_{\mathrm{em}})|\overline{\mathrm{P}}{}_{}{}^{0}=\overline{m}{}_{0}{}^{}|\overline{\mathrm{P}}{}_{}{}^{0}$$
where $`m_0`$ and $`\overline{m}_0`$ are the rest masses of $`\mathrm{P}^0`$ and $`\overline{\mathrm{P}}^0`$, respectively. The $`\mathrm{P}^0`$ and $`\overline{\mathrm{P}}^0`$ states are related through CP transformations. For stationary states, the T transformation does not alter them, with the exception of an arbitrary phase. While CP is a unitary operation, T is an antiunitary operation.
In summary, we obtain
$$\begin{array}{c}CP|\mathrm{P}^0=e^{i\theta _{\mathrm{CP}}}|\overline{\mathrm{P}}{}_{}{}^{0}\mathrm{and}CP|\overline{\mathrm{P}}{}_{}{}^{0}=e^{i\theta _{\mathrm{CP}}}|\mathrm{P}^0\\ T|\mathrm{P}^0=e^{i\theta _\mathrm{T}}|\mathrm{P}^0\mathrm{and}T|\overline{\mathrm{P}}{}_{}{}^{0}=e^{i\overline{\theta }_\mathrm{T}}|\overline{\mathrm{P}}{}_{}{}^{0}\end{array}$$
(1)
where the $`\theta `$’s are arbitrary phases, and by assuming $`CPT|\mathrm{P}^0=TCP|\mathrm{P}^0`$ it follows that
$$2\theta _{\mathrm{CP}}=\overline{\theta }{}_{\mathrm{T}}{}^{}\theta _\mathrm{T}.$$
Since T is antiunitary, it follows that
$$Tc=c^{}T$$
where $`c`$ is any complex number. If we define
$$T|\alpha =|\stackrel{~}{\alpha },T|\beta =|\stackrel{~}{\beta }$$
antiunitary operation has to give
$$\alpha |\beta =\left[\stackrel{~}{\alpha }|\stackrel{~}{\beta }\right]^{}.$$
On the other hand,
$$\alpha |\beta =\alpha |\left(T^1T|\beta \right)=\alpha |\left(T^1|\stackrel{~}{\beta }\right),$$
hence
$$\alpha |\left(T^1|\stackrel{~}{\beta }\right)=\left[\stackrel{~}{\alpha }|\stackrel{~}{\beta }\right]^{}.$$
We can then conclude
$$\alpha |\left(T^1|\stackrel{~}{\beta }\right)=\left[\left(\alpha |T^1\right)|\stackrel{~}{\beta }\right]^{}$$
i.e. when the T operator changes the direction of the operation, it must be complex conjugated.
If strong and electromagnetic interactions are invariant under the CPT transformation, which is assumed throughout this paper, it follows that $`m_0=\overline{m}_0`$.
Now we switch on an interaction, $`V`$, and the P can decay into final states f with different flavours ($`|\mathrm{\Delta }F|=1`$ process) and $`\mathrm{P}^0`$ and $`\overline{\mathrm{P}}^0`$ can oscillate to each other ($`|\mathrm{\Delta }F|=2`$ process). Thus, a general state $`|\psi (t)`$ which is a solution of the Schrödinger equation
$$i\frac{}{t}|\psi (t)=\left(H_{\mathrm{st}}+H_{\mathrm{em}}+V\right)|\psi (t)$$
(2)
can be written as
$$|\psi (t)=a(t)|\mathrm{P}^0+b(t)|\overline{\mathrm{P}}{}_{}{}^{0}+\underset{\mathrm{f}}{}c_\mathrm{f}(t)|\mathrm{f}$$
where the sum is taken over all the possible final states f and $`a(t)`$, $`b(t)`$ and $`c_\mathrm{f}(t)`$ are time dependent functions; $`|a(t)|^2`$, $`|b(t)|^2`$ and $`|c_\mathrm{f}(t)|^2`$ give the fractions of $`\mathrm{P}^0`$, $`\overline{\mathrm{P}}^0`$ and f at time $`t`$ respectively. Since the weak interaction is much weaker than strong and electromagnetic interactions, perturbation theory can be applied in order to solve equation 2. Also with the help of the Wigner-Weisskopf method, which neglects the weak interactions between the final states , and we obtain
$$i\frac{}{t}\left(\begin{array}{c}a(t)\\ b(t)\end{array}\right)=𝚲\left(\begin{array}{c}a(t)\\ b(t)\end{array}\right)=\left(𝑴i\frac{𝚪}{2}\right)\left(\begin{array}{c}a(t)\\ b(t)\end{array}\right)$$
(3)
where the $`2\times 2`$ matrices $`𝑴`$ and $`𝚪`$ are often referred to as the mass and decay matrices.
The elements of the mass matrix are given as
$$M_{ij}=m_0\delta _{ij}+i|V|j+\underset{\mathrm{f}}{}𝒫\left(\frac{i|V|\mathrm{f}\mathrm{f}|V|j}{m_0E_\mathrm{f}}\right)$$
(4)
where $`𝒫`$ stands for the principal part and the index $`i=1`$(2) denotes $`\mathrm{P}^0`$($`\overline{\mathrm{P}}^0`$). Note that the sum is taken over all possible intermediate states common to $`\mathrm{P}^0`$ and $`\overline{\mathrm{P}}^0`$ for $`ij`$.
The elements of the decay matrix are given by
$$\Gamma _{ij}=2\pi \underset{\mathrm{f}}{}i|V|\mathrm{f}\mathrm{f}|V|j\delta (m_0E_\mathrm{f})$$
(5)
The sum is taken over only real final states common to $`\mathrm{P}^0`$ and $`\overline{\mathrm{P}}^0`$ for $`ij`$.
If the Hamiltonians are not Hermitian, transition probabilities are not conserved in decays or oscillations, i.e. the number of initial states is not identical to the number of final states. This is also referred to as the break down of unitarity. We assume from now on that all the Hamiltonians are Hermitian, i.e.
$$|a(t)|^2+|b(t)|^2+\underset{\mathrm{f}}{}|c_\mathrm{f}|^2=1,$$
and also
$$M_{ij}=M_{ji}^{},\Gamma _{ij}=\Gamma _{ji}^{}.$$
Clearly $`|a(t)|^2+|b(t)|^2`$ decreases as a function of time, hence $`𝚲`$ is not Hermitian.
Since the CP operator changes a particle state into an antiparticle state, the following relation can be obtained if $`V`$ is invariant under the CP transformation, i.e. $`(CP)^1VCP=V`$:
$$\mathrm{CP}:\left|\Lambda _{12}\right|=\left|\Lambda _{21}\right|,\Lambda _{11}=\Lambda _{22}.$$
Since the T operator induces complex conjugation, which is equivalent to interchanging a bra-state and a ket-state, the following relation can be obtained if $`V`$ is invariant under the T transformation:
$$\mathrm{T}:\left|\Lambda _{12}\right|=\left|\Lambda _{21}\right|.$$
By combining the two, we obtain for the CPT invariant case:
$$\mathrm{CPT}:\Lambda _{11}=\Lambda _{22}.$$
For a rigorous proof, equations 1, 4 and 5 are used.
It follows that
$$\begin{array}{c}\mathrm{if}\Lambda _{11}\Lambda _{22},\mathrm{i}.\mathrm{e}.M_{11}M_{22}\mathrm{or}\Gamma _{11}\Gamma _{22}:\hfill \\ \mathrm{𝐂𝐏𝐓}\mathrm{𝐚𝐧𝐝}\mathrm{𝐂𝐏}\mathrm{are}\mathrm{violated}\hfill \\ \mathrm{if}|\Lambda _{12}||\Lambda _{21}|,\mathrm{i}.\mathrm{e}.\mathrm{sin}(\phi _\Gamma \phi _M)0:\hfill \\ 𝐓\mathrm{𝐚𝐧𝐝}\mathrm{𝐂𝐏}\mathrm{are}\mathrm{violated}.\hfill \end{array}$$
(6)
where $`\phi _M=\mathrm{arg}\left(M_{12}\right)\mathrm{and}\phi _\Gamma =\mathrm{arg}\left(\Gamma _{12}\right)`$. Note that CP violation in the mass and decay matrices cannot be separated from CPT violation or T violation.
While there is no fundamental reason to respect CP and T symmetries, it can be shown based on only few basic assumptions that no self consistent quantum field theory can be constructed that does not conserve CPT symmetry . Therefore, we restrict our further discussion to the case where CPT symmetry is conserved: i.e.
$$M_{11}=M_{22}M,\Gamma _{11}=\Gamma _{22}\Gamma $$
thus
$$\Lambda _{11}=\Lambda _{22}\Lambda .$$
Differential equation 3 can be reduced to
$$\frac{d^2a(t)}{dt^2}+2i\Lambda \frac{da(t)}{dt}+\left(\Lambda _{12}\Lambda _{21}\Lambda ^2\right)a(t)=0$$
(7)
for $`a(t)`$, and a general solution is given by
$$a(t)=C_+e^{i\lambda _+}+C_{}e^{i\lambda _{}}$$
where $`C_\pm `$ are arbitrary constants which can only be defined by the initial condition. For $`b(t)`$, we obtain
$$b(t)=\frac{1}{\Lambda _{12}}\left[i\frac{da(t)}{dt}\Lambda a(t)\right]$$
which can be used once $`a(t)`$ becomes known.
Insertion of $`a(t)`$ into equation 7 leads to
$$\lambda _\pm ^22\Lambda \lambda _\pm +\left(\Lambda _{12}\Lambda _{21}\Lambda ^2\right)=0$$
from which the eigen-frequencies are obtained as
$$\lambda _\pm =\Lambda \pm \sqrt{\Lambda _{12}\Lambda _{21}}m_\pm \frac{i}{2}\Gamma _\pm $$
by solving where
$$m_\pm =\mathrm{}\lambda _\pm =M\pm \mathrm{}\left(\Lambda _{12}\Lambda _{21}\right)^{1/2}$$
(8)
and
$$\Gamma _\pm =2\mathrm{}\lambda _\pm =\Gamma \mathrm{\hspace{0.17em}2}\mathrm{}\left(\Lambda _{12}\Lambda _{21}\right)^{1/2}.$$
(9)
For an initially pure $`\mathrm{P}^0`$ state, we have $`a(t)=1`$ and $`b(t)=0`$ at $`t=0`$, i.e. $`C_+=C_{}=1/2`$, and the solution is given by
$`|\mathrm{P}{}_{}{}^{0}(t)`$ $`=`$ $`a(t)|\mathrm{P}^0+b(t)|\overline{\mathrm{P}}{}_{}{}^{0}`$ (10)
$`=`$ $`f_+(t)|\mathrm{P}^0+\zeta f_{}(t)|\overline{\mathrm{P}}{}_{}{}^{0}`$
$`=`$ $`{\displaystyle \frac{\sqrt{\mathrm{\hspace{0.17em}1}+|\zeta |^2}}{2}}\left(|\mathrm{P}_+e^{i\lambda _+t}+|\mathrm{P}_{}e^{i\lambda _{}t}\right)`$ (11)
where
$$f_\pm (t)=\frac{1}{\mathrm{\hspace{0.17em}2}}\left(e^{i\lambda _+t}\pm e^{i\lambda _{}t}\right)$$
and $`\zeta `$ is
$$\zeta =\sqrt{\frac{\Lambda _{21}}{\Lambda _{12}}}.$$
(12)
The two states $`|\mathrm{P}_+`$ and $`|\mathrm{P}_{}`$ are the eigenstates of $`\lambda _\pm `$ and are given by
$$|\mathrm{P}_\pm =\frac{1}{\sqrt{\mathrm{\hspace{0.17em}1}+|\zeta |^2}}(|\mathrm{P}^0\pm \zeta |\overline{\mathrm{P}}{}_{}{}^{0}).$$
(13)
For an initially pure $`\overline{\mathrm{P}}^0`$ state, we have
$`|\overline{\mathrm{P}}{}_{}{}^{0}(t)`$ $`=`$ $`{\displaystyle \frac{1}{\zeta }}f_{}(t)|\mathrm{P}^0+f_+(t)|\overline{\mathrm{P}}{}_{}{}^{0}`$ (14)
$`=`$ $`{\displaystyle \frac{\sqrt{\mathrm{\hspace{0.17em}1}+|\zeta |^2}}{2\zeta }}\left(|\mathrm{P}_+e^{i\lambda _+t}|\mathrm{P}_{}e^{i\lambda _{}t}\right).`$ (15)
While $`\mathrm{P}^\pm `$ have definite masses and decay widths (as seen from equations 11 and 15), $`\mathrm{P}^0`$ and $`\overline{\mathrm{P}}^0`$ do not and they oscillate to each other (see equations 10 and 14).
### 2.2 CP Conserving Case
If $`V`$ remains invariant under the CP transformation, from equations 1, 4 and 5 it follows that
$$M_{12}=M_{21}e^{i\mathrm{\hspace{0.17em}2}\theta _{\mathrm{CP}}}=M_{12}^{}e^{i\mathrm{\hspace{0.17em}2}\theta _{\mathrm{CP}}}$$
thus
$$\mathrm{arg}M_{12}=\theta _{\mathrm{CP}}+n\pi ,$$
and
$$\Gamma _{12}=\Gamma _{21}e^{i\mathrm{\hspace{0.17em}2}\theta _{\mathrm{CP}}}=\Gamma _{12}^{}e^{i\mathrm{\hspace{0.17em}2}\theta _{\mathrm{CP}}}$$
thus
$$\mathrm{arg}\Gamma _{12}=\theta _{\mathrm{CP}}+n^{}\pi ,$$
where $`n`$ and $`n^{}`$ are arbitrary integer numbers.
For $`\zeta `$, we have
$$\zeta =\sqrt{\frac{\Lambda _{21}}{\Lambda _{12}}}=e^{i(\theta _{\mathrm{CP}}+n^{\prime \prime }\pi )}$$
where $`n^{\prime \prime }`$ is an arbitrary integer number. The two mass eigenstates $`|\mathrm{P}_+`$ and $`|\mathrm{P}_{}`$ become CP eigenstates
$$CP|\mathrm{P}_\pm =\pm (1)^{n^{\prime \prime }}|\mathrm{P}_\pm .$$
The mass and decay width eigenvalues, equations 8 and 9, become
$$m_\pm =M\pm (1)^{n+n^{\prime \prime }}|M_{12}|$$
and
$$\Gamma _\pm =\Gamma \pm (1)^{n^{}+n^{\prime \prime }}|\Gamma _{12}|$$
By examining various combinations of $`n`$, $`n^{}`$ and $`n^{\prime \prime }`$, we can show that the following four possibilities exist:
1. $`n`$=even, $`n^{}`$=even: $`CP=+1`$ state is heavier and decays faster,
2. $`n`$=even, $`n^{}`$=odd: $`CP=+1`$ state is heavier and decays slower,
3. $`n`$=odd, $`n^{}`$=even: $`CP=+1`$ state is lighter and decays faster,
4. $`n`$=odd, $`n^{}`$=odd: $`CP=+1`$ state is lighter and decays slower.
Figure 1 illustrates the phase relations in a pictorial way. The choice of $`n^{\prime \prime }`$ does not alter the conclusion and $`n^{\prime \prime }=0`$ can be adopted without any loss of generality. In this case, $`|\mathrm{P}_+`$ is the $`CP=+1`$ state.
### 2.3 CP Violating Case
Let us consider the time dependent decay rate for the initial $`\mathrm{P}^0`$ decaying into a CP eigenstate f, given by $`|f|V|\mathrm{P}^0(t)|^2`$, and that for the initial $`\overline{\mathrm{P}}^0`$ decaying into f, given by $`|f|V|\overline{\mathrm{P}}{}_{}{}^{0}(t)|^2`$:
$`R_\mathrm{f}(t)`$ $``$ $`|f_+(t)|^2+\left|\zeta {\displaystyle \frac{\overline{A}_\mathrm{f}}{A_\mathrm{f}}}\right|^2|f_{}(t)|^2+2\mathrm{}\left[\zeta {\displaystyle \frac{\overline{A}_\mathrm{f}}{A_\mathrm{f}}}f_+^{}(t)f_{}(t)\right]`$ (16)
$`\overline{R}{}_{\mathrm{f}}{}^{}(t)`$ $``$ $`\left|{\displaystyle \frac{\overline{A}_\mathrm{f}}{A_\mathrm{f}}}\right|^2|f_+(t)|^2+\left|{\displaystyle \frac{1}{\zeta }}\right|^2|f_{}(t)|^2+{\displaystyle \frac{2}{|\zeta |^2}}\mathrm{}\left[\zeta ^{}{\displaystyle \frac{\overline{A}_\mathrm{f}^{}}{A_\mathrm{f}^{}}}f_+^{}(t)f_{}(t)\right]`$ (17)
where the instantaneous decay amplitudes are denoted by $`A_\mathrm{f}\mathrm{f}|V|\mathrm{P}^0`$ etc. and equations 10 and 14 are used.
Since $`R_f(t)`$ and $`\overline{R}{}_{\mathrm{f}}{}^{}(t)`$ describe the CP conjugated processes to each other, any difference between the two is an clear proof of CP violation. As seen from the first terms of equations 16 and 17, CP violation is generated if $`|A_\mathrm{f}||\overline{A}{}_{\mathrm{f}}{}^{}|`$. This is called CP violation in the decay amplitudes.
From the second terms of $`R_\mathrm{f}(t)`$ and $`\overline{R}{}_{\mathrm{f}}{}^{}(t)`$, it can be seen that CP violation is generated if $`|\zeta |1`$ even if there is no CP violation in the decay amplitudes. From equations 11 and 15, it is clear that the oscillation rate for $`\mathrm{P}^0\overline{\mathrm{P}}^0`$ is different from that for $`\overline{\mathrm{P}}{}_{}{}^{0}\mathrm{P}^0`$ if $`|\zeta |1`$, thus this is called CP violation in the oscillation.
The third term can be expanded into
$$2\mathrm{}\left(\zeta \frac{\overline{A}_\mathrm{f}}{A_\mathrm{f}}\right)\mathrm{}\left[f_+^{}(t)f_{}(t)\right]2\mathrm{}\left(\zeta \frac{\overline{A}_\mathrm{f}}{A_\mathrm{f}}\right)\mathrm{}\left[f_+^{}(t)f_{}(t)\right]$$
for $`R_\mathrm{f}(t)`$ and
$$\frac{2}{|\zeta |^2}\mathrm{}\left(\zeta \frac{\overline{A}_\mathrm{f}}{A_\mathrm{f}}\right)\mathrm{}\left[f_+^{}(t)f_{}(t)\right]+\frac{2}{|\zeta |^2}\mathrm{}\left(\zeta \frac{\overline{A}_\mathrm{f}}{A_\mathrm{f}}\right)\mathrm{}\left[f_+^{}(t)f_{}(t)\right]$$
for $`\overline{R}{}_{\mathrm{f}}{}^{}(t)`$. If CP violation in $`\mathrm{P}^0`$-$`\overline{\mathrm{P}}^0`$ oscillation is absent, the first terms are identical. Even in that case, if
$$\mathrm{}\left(\zeta \frac{\overline{A}_\mathrm{f}}{A_\mathrm{f}}\right)0$$
CP violation is still present. Since the process involves the decays of $`\mathrm{P}^0`$ ($`\overline{\mathrm{P}}^0`$) from the initial $`\mathrm{P}^0`$ ($`\overline{\mathrm{P}}^0`$) and decays of the $`\overline{\mathrm{P}}^0`$ ($`\mathrm{P}^0`$) oscillated from the initial $`\mathrm{P}^0`$ ($`\overline{\mathrm{P}}^0`$) into a common final state, it is referred as CP violation due to the interplay between the decays and oscillations.
If CP violation in $`\mathrm{P}^0`$-$`\overline{\mathrm{P}}^0`$ oscillation is small, i.e. $`\left(|\zeta |1\right)^2<<1`$, we can derive
$$|\mathrm{sin}(\phi _\Gamma \phi _M)|<<1$$
from equation 12, where $`\phi _\Gamma =\mathrm{arg}\Gamma _{12}`$ and $`\phi _M=\mathrm{arg}M_{12}`$ as already defined. By introducing $`|\mathrm{\Delta }_{\Gamma /M}|<<1`$ as
$$\phi _\Gamma \phi _M=n\pi \mathrm{\Delta }_{\Gamma /M}$$
(18)
where $`n`$ is an integer number, the following two approximations are possible:
a) $`\phi _\Gamma =\mathrm{arg}\Gamma _{12}`$ base
$`\zeta `$ $``$ $`\left\{1{\displaystyle \frac{2|M_{12}||\Gamma _{12}|\mathrm{\Delta }_{\Gamma /M}}{4|M_{12}|^2+|\Gamma _{12}|^2}}\left[(1)^{n+1}+i{\displaystyle \frac{2|M_{12}|}{|\Gamma _{12}|}}\right]\right\}e^{i\phi _\Gamma }`$
$`m_\pm `$ $`=`$ $`M\pm (1)^n|M_{12}|`$
$`\Gamma _\pm `$ $`=`$ $`\Gamma \pm |\Gamma _{12}|`$
b) $`\phi _M=\mathrm{arg}M_{12}`$ base
$`\zeta `$ $``$ $`\left\{1+{\displaystyle \frac{2|M_{12}||\Gamma _{12}|\mathrm{\Delta }_{\Gamma /M}}{4|M_{12}|^2+|\Gamma _{12}|^2}}\left[(1)^n+i{\displaystyle \frac{|\Gamma _{12}|}{2|M_{12}|}}\right]\right\}e^{i\phi _M}`$ (19)
$`m_\pm `$ $`=`$ $`M\pm |M_{12}|`$
$`\Gamma _\pm `$ $`=`$ $`\Gamma \pm (1)^n|\Gamma _{12}|.`$
## 3 Neutral Kaon System
### 3.1 Adaptation of Formalism
Now we adapt the above developed formalism to the neutral kaon system. As described later, observed CP violation in the $`\mathrm{K}^0`$-$`\overline{\mathrm{K}}^0`$ oscillation is very small. The two mass eigenstates are called $`\mathrm{K}_\mathrm{S}`$ and $`\mathrm{K}_\mathrm{L}`$ with corresponding masses and decay widths referred to as $`m_\mathrm{S}`$, $`m_\mathrm{L}`$, $`\Gamma _\mathrm{S}`$ and $`\Gamma _\mathrm{L}`$ respectively and they are known to be $`m_\mathrm{S}<m_\mathrm{L}`$ and $`\Gamma _\mathrm{S}>\Gamma _\mathrm{L}`$. Therefore, $`M_{12}`$ and $`\Gamma _{12}`$ is almost antiparallel to each other, thus $`n=1`$ in equation 18.
Since the kaon decay properties are experimentally well measured, enough information is available to calculate $`\Gamma _{12}`$ from the data, as described in Section 3.5. We therefore adopt the $`\phi _\Gamma `$ base given in the previous section.
It follows that
$`\zeta =(12ϵ)e^{i\phi _\Gamma }`$ (20)
where the small parameter $`ϵ`$ is given by
$$ϵ=\frac{|M_{12}||\Gamma _{12}|\mathrm{sin}\left(\phi _\Gamma \phi _M\right)}{4|M_{12}|^2+|\Gamma _{12}|^2}\left(1+i\frac{\mathrm{\hspace{0.17em}2}|M_{12}|}{|\Gamma _{12}|}\right).$$
and
$`|\mathrm{K}_\mathrm{S}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{1+|ϵ|^2}}}[|\mathrm{K}^0+(12ϵ)e^{i\phi _\Gamma }|\overline{\mathrm{K}}{}_{}{}^{0}]`$ (21)
$`|\mathrm{K}_\mathrm{L}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{1+|ϵ|^2}}}[|\mathrm{K}^0(12ϵ)e^{i\phi _\Gamma }|\overline{\mathrm{K}}{}_{}{}^{0}].`$ (22)
From the measured lifetimes ,
$$\tau _\mathrm{s}\frac{1}{\Gamma _\mathrm{S}}=(0.8934\pm 0.0008)\times 10^{10}\mathrm{s}$$
and
$$\tau _\mathrm{L}\frac{1}{\Gamma _\mathrm{L}}=(5.17\pm 0.04)\times 10^8\mathrm{s}$$
i.e.
$$\mathrm{\Delta }\Gamma =\Gamma _\mathrm{S}\Gamma _\mathrm{L}=(1.1174\pm 0.0010)\times 10^{10}\mathrm{s}^1$$
and the mass difference,
$$\mathrm{\Delta }mm_\mathrm{L}m_\mathrm{S}=(0.5301\pm 0.0014)\times 10^{10}\mathrm{}\mathrm{s}^1$$
we obtain,
$$\frac{|M_{12}||\Gamma _{12}|}{4|M_{12}|^2+|\Gamma _{12}|^2}=0.24966\pm 0.00004$$
and
$$\frac{\mathrm{\hspace{0.17em}2}|M_{12}|}{|\Gamma _{12}|}=0.9488\pm 0.0026.$$
Since the lifetime of $`\mathrm{K}_\mathrm{L}`$ is much longer than that of $`\mathrm{K}_\mathrm{S}`$, it is possible to produce a $`\mathrm{K}_\mathrm{L}`$ beam. Therefore, many kaon experiments have been done using $`\mathrm{K}_\mathrm{L}`$ beams.
### 3.2 CP Violation in Oscillations
The CPLEAR experiment observed CP violation in the $`\mathrm{K}^0`$-$`\overline{\mathrm{K}}^0`$ oscillation by measuring the difference in the oscillation rates between $`\overline{\mathrm{K}}{}_{}{}^{0}\mathrm{K}^0`$ and $`\mathrm{K}^0\overline{\mathrm{K}}^0`$. The initial neutral kaons were produced by $`\mathrm{p}\overline{\mathrm{p}}`$ annihilations: $`\mathrm{p}\overline{\mathrm{p}}\mathrm{K}^0\mathrm{K}^{}\pi ^+`$ and $`\overline{\mathrm{K}}{}_{}{}^{0}\mathrm{K}_{}^{+}\pi ^{}`$, where the initial flavour can be defined by the charge sign of the accompanying kaon. Semileptonic decays were used in order to determine the flavour at the moment of the decay. Since the $`\mathrm{K}^0`$ contains an $`\overline{\mathrm{s}}`$-quark (and $`\overline{\mathrm{K}}^0`$ an s-quark), $`\mathrm{K}^0`$ ($`\overline{\mathrm{K}}^0`$) can decay only into $`\mathrm{e}^+\pi ^{}\nu `$ ($`\mathrm{e}^{}\pi ^+\overline{\nu }`$) instantaneously. Therefore, the initial $`\mathrm{K}^0`$ ($`\overline{\mathrm{K}}^0`$) can produce the final state $`\mathrm{e}^{}\pi ^+\overline{\nu }`$ ($`\mathrm{e}^+\pi ^{}\nu `$) only through the $`\mathrm{K}^0\overline{\mathrm{K}}^0`$ ($`\overline{\mathrm{K}}{}_{}{}^{0}\mathrm{K}^0`$) oscillation. From the two measured time dependent decay rates, $`R_\mathrm{e}^{}(t)`$ and $`\overline{R}_{\mathrm{e}^+}(t)`$, an asymmetry
$$A_\mathrm{T}(t)=\frac{\overline{R}_{\mathrm{e}^+}(t)R_\mathrm{e}^{}(t)}{\overline{R}_{\mathrm{e}^+}(t)+R_\mathrm{e}^{}(t)}$$
is constructed as shown in Figure 2. Using equations 10, 14 and 20, it follows that
$$A_\mathrm{T}(t)=\frac{1|\zeta |^4}{\mathrm{\hspace{0.17em}1}+|\zeta |^4}=4\mathrm{}ϵ$$
and from the measured $`A_\mathrm{T}(t)=(6.6\pm 1.6)\times 10^3`$ ,
$$|\zeta |=0.9967\pm 0.00081$$
is obtained exhibiting a clear sign of CP violation and T violation in the $`\mathrm{K}^0`$-$`\overline{\mathrm{K}}^0`$ oscillation.
The parameter $`|\zeta |`$ can also be measured from the semileptonic branching fractions of $`\mathrm{K}_\mathrm{L}`$ by the lepton sign asymmetry: using equations 22 and 20, we obtain
$`\delta _{\mathrm{}}`$ $``$ $`{\displaystyle \frac{B(\mathrm{K}_\mathrm{L}\mathrm{}^+\pi ^{}\nu )B(\mathrm{K}_\mathrm{L}\mathrm{}^{}\pi ^+\overline{\nu })}{B(\mathrm{K}_\mathrm{L}\mathrm{}^+\pi ^{}\nu )+B(\mathrm{K}_\mathrm{L}\mathrm{}^{}\pi ^+\overline{\nu })}}`$
$`=`$ $`{\displaystyle \frac{1|\zeta |^2}{1+|\zeta |^2}}=2\mathrm{}ϵ`$
$`=`$ $`(3.27\pm 0.12)\times 10^3`$
where $`\mathrm{}`$ can be e or $`\mu `$ and $`B`$ stands for a branching fraction.
Using all the measurements, we obtain
$$\mathrm{}ϵ=(1.64\pm 0.06)\times 10^3$$
and
$$\mathrm{arg}ϵ=(43.50\pm 0.08)^{}.$$
### 3.3 CP Violation due to Decays and Oscillations
Since the two-pion final state is a CP eigenstate with $`CP=+1`$, $`\mathrm{K}_\mathrm{L}`$ decaying into $`\pi ^+\pi ^{}`$ is a CP violating decay. This was indeed the first observed sign of CP violation. A commonly used CP violation parameter $`\eta _+`$ is defined as
$$\eta _+\frac{\pi ^+\pi ^{}|V|\mathrm{K}_\mathrm{L}}{\pi ^+\pi ^{}|V|\mathrm{K}_\mathrm{S}}=\frac{1\zeta \frac{\overline{A}_+}{A_+}}{\mathrm{\hspace{0.17em}1}+\zeta \frac{\overline{A}_+}{A_+}}$$
(23)
where equations 21, 22 are used and $`A_+`$ and $`\overline{A}_+`$ denote the $`\mathrm{K}^0`$ and $`\overline{\mathrm{K}}{}_{}{}^{0}\pi ^+\pi ^{}`$ decay amplitudes respectively.
The parameter $`\eta _+`$ can be measured from the time dependent decay rates for the initial $`\mathrm{K}^0`$ and $`\overline{\mathrm{K}}^0`$ into $`\pi ^+\pi ^{}`$. From equations 11 and 15, the two rates are given by
$$R_+(t)\frac{1}{2}e^{\Gamma _\mathrm{S}t}+\left|\eta _+\right|^2e^{\Gamma _\mathrm{L}t}+2|\eta _+|e^{\widehat{\Gamma }t}\mathrm{cos}(\mathrm{\Delta }mt\varphi _+)$$
and
$$\overline{R}_+(t)\frac{\mathrm{\hspace{0.17em}1}+4\mathrm{}ϵ}{2}\left[e^{\Gamma _\mathrm{S}t}+\left|\eta _+\right|^2e^{\Gamma _\mathrm{L}t}2|\eta _+|e^{\widehat{\Gamma }t}\mathrm{cos}(\mathrm{\Delta }mt\varphi _+)\right]$$
where $`\varphi _+`$ is the phase of $`\eta _+`$ and $`\widehat{\Gamma }`$ is the $`\mathrm{K}_\mathrm{S}`$-$`\mathrm{K}_\mathrm{L}`$ average decay width. The second term is CP violating $`\mathrm{K}_\mathrm{L}`$ decays and the third term is due to the interference between the $`\mathrm{K}_\mathrm{S}`$ decay and CP violating $`\mathrm{K}_\mathrm{L}`$ decay amplitudes. Figure 3 shows the measured $`R_+(t)`$ and $`R_+(t)`$ together with the CP asymmetry defined as
$$A_+(t)=\frac{\overline{R}_+(t)R_+(t)}{\overline{R}_+(t)+R_+(t)}$$
where the interference term is well isolated. At around $`t=10\tau _\mathrm{S}`$, the $`\mathrm{K}_\mathrm{S}`$ decay rate is reduced to the level of the CP violating $`\mathrm{K}_\mathrm{L}`$ decay rate, thus the asymmetry becomes very large.
This direct comparison between the two CP conjugated processes illustrates another straightforward demonstration of CP violation in the neutral kaon system. From the asymmetry, the value of $`\eta _+`$ is measured to be
$$|\eta _+|=(2.264\pm 0.035)\times 10^3,\varphi _+=(43.19\pm 0.60)^{}$$
which leads to
$$\mathrm{}\left(\zeta \frac{\overline{A}_+}{A_+}\right)=(3.099\pm 0.048)\times 10^3$$
exhibiting that CP violation due to the interference between the decay and oscillation is present.
### 3.4 CP Violation in Decays
The two-pion final state can be in a total isospin state of $`I=0`$ or $`I=2`$. The $`I=1`$ state is not allowed due to Bose statistics. Using the isospin decomposition, we can derive the $`\mathrm{K}^0`$ and $`\overline{\mathrm{K}}^0`$ decay amplitudes to $`\pi ^+\pi ^{}`$ to be
$$A_+=\sqrt{\frac{2}{3}}2\pi (I=0)|V|\mathrm{K}^0+\sqrt{\frac{1}{3}}2\pi (I=2)|V|\mathrm{K}^0$$
and
$$\overline{A}{}_{+}{}^{}=\sqrt{\frac{2}{3}}2\pi (I=0)|V|\overline{\mathrm{K}}{}_{}{}^{0}+\sqrt{\frac{1}{3}}2\pi (I=2)|V|\overline{\mathrm{K}}{}_{}{}^{0}.$$
Using CPT symmetry and the S-matrix, the $`\mathrm{K}^0`$ and $`\overline{\mathrm{K}}^0`$ decay amplitudes can be related and it follows that
$`A_+`$ $`=`$ $`\sqrt{{\displaystyle \frac{2}{3}}}a_0e^{i\delta _0}+\sqrt{{\displaystyle \frac{1}{3}}}a_2e^{i\delta _2}`$
$`\overline{A}_+`$ $`=`$ $`\sqrt{{\displaystyle \frac{2}{3}}}a_0^{}e^{i(\delta _0+\theta _{\mathrm{CP}}\overline{\theta }{}_{\mathrm{T}}{}^{})}+\sqrt{{\displaystyle \frac{1}{3}}}a_2^{}e^{i(\delta _2+\theta _{\mathrm{CP}}\overline{\theta }{}_{\mathrm{T}}{}^{})}`$
where $`a_0`$ and $`a_2`$ are the $`\mathrm{K}^0`$ decay amplitudes into $`2\pi (I=0)`$ and $`2\pi (I=2)`$ states due to the short-range weak interactions and $`\delta _0`$ and $`\delta _2`$ are the $`\pi `$-$`\pi `$ scattering phase shifts for the $`I=0`$ and $`I=2`$ two-pion configuration at $`\sqrt{s}=m_\mathrm{K}`$ respectively. It is important to note that the two-pion scattering is totally dominated by the elastic scattering at the energy scale of the kon mass. Similarly for the $`\pi ^0\pi ^0`$ final state, we have
$`A_{00}`$ $`=`$ $`\sqrt{{\displaystyle \frac{1}{3}}}a_0e^{i\delta _0}+\sqrt{{\displaystyle \frac{2}{3}}}a_2e^{i\delta _2}`$
$`\overline{A}_{00}`$ $`=`$ $`\sqrt{{\displaystyle \frac{1}{3}}}a_0^{}e^{i(\delta _0+\theta _{\mathrm{CP}}\overline{\theta }{}_{\mathrm{T}}{}^{})}+\sqrt{{\displaystyle \frac{2}{3}}}a_2^{}e^{i(\delta _2+\theta _{\mathrm{CP}}\overline{\theta }{}_{\mathrm{T}}{}^{})}.`$
As seen from the amplitudes, $`B(\mathrm{K}_\mathrm{S}\pi ^0\pi ^0)/B(\mathrm{K}_\mathrm{S}\pi ^+\pi ^{})`$ would be $`0.5`$ if $`a_2=0`$. Since the measured ratio is $`0.46`$ , we can conclude that $`|a_2/a_0|<<1`$. It follows that
$$\frac{\overline{A}_+}{A_+}=\left(12ϵ^{}\right)e^{i(2\phi _0+\overline{\theta }{}_{\mathrm{T}}{}^{}\theta _{\mathrm{CP}})}$$
(24)
where the parameter $`ϵ^{}`$ is given by
$$ϵ^{}=\frac{1}{\sqrt{2}}\left|\frac{a_2}{a_0}\right|\mathrm{sin}(\phi _2\phi _0)e^{i(\pi /2+\delta _2\delta _0)}$$
(25)
and $`\phi _{0,2}=\mathrm{arg}a_{0,2}`$.
As seen from equation 24, CP violation in the decay amplitude, $`|A_+||\overline{A}{}_{+}{}^{}|`$, is present if $`\mathrm{}ϵ^{}0`$. From equation 25, this is possible only if
$$\mathrm{sin}(\phi _2\phi _0)0\mathrm{and}\mathrm{sin}(\delta _2\delta _0)0.$$
i.e. both the weak and strong phases have to be different for the $`I=0`$ and $`I=2`$ decay amplitudes. More generally, there must be two processes leading to the identical final state and both the strong and the weak phases must be different between the two processes in order to generate CP violation in the decay amplitudes. It should be noted that from the measured $`\pi `$-$`\pi `$ scattering phase shift values, we have
$$\mathrm{arg}ϵ^{}=(43\pm 6)^{}$$
Using equations 20 and 24, it follows that
$`\zeta {\displaystyle \frac{\overline{A}_+}{A_+}}`$ $`=`$ $`(12ϵ2ϵ^{})e^{i(\phi _\Gamma +2\phi _0+\overline{\theta }{}_{\mathrm{T}}{}^{}\theta _{\mathrm{CP}})}`$
$``$ $`12(ϵ+ϵ^{})i(\phi _\Gamma +2\phi _0+\overline{\theta }{}_{\mathrm{T}}{}^{}\theta _{\mathrm{CP}})`$
where the approximation is made assuming that the phase difference between $`\Gamma _{12}`$ and $`A_0\overline{A}_0`$ is small, which will be justified later. From equation 23, $`\eta _+`$ can be derived to be
$$\eta _+=ϵ+i(\phi _\Gamma +2\phi _0+\overline{\theta }{}_{\mathrm{T}}{}^{}\theta _{\mathrm{CP}})+ϵ^{}.$$
Similarly the CP violation parameter for the $`\pi ^0\pi ^0`$ decay channel, $`\eta _{00}`$, is given by
$$\eta _{00}=ϵ+i(\phi _\Gamma +2\phi _0+\overline{\theta }{}_{\mathrm{T}}{}^{}\theta _{\mathrm{CP}})2ϵ^{}.$$
Thus, we expect CP violation parameters to be different between the $`\pi ^+\pi ^{}`$ and $`\pi ^0\pi ^0`$ decay modes if $`ϵ^{}0`$. It has been shown by four recent experiments, NA31 , E731 , KTeV and NA48 ,
$$\left|\frac{\eta _+}{\eta _{00}}\right|^2=1.0127\pm 0.0028$$
i.e. CP violation in the decay amplitude is present in the neutral kaon system. If we neglect $`(\phi _\Gamma +2\phi _0+\overline{\theta }{}_{\mathrm{T}}{}^{}\theta _{\mathrm{CP}})`$, it follows that
$$\mathrm{}\left(\frac{ϵ^{}}{ϵ}\right)=\frac{1}{6}\left(\left|\frac{\eta _+}{\eta _{00}}\right|^21\right).$$
### 3.5 Phase of Decay Matrix
As seen from equation 5, evaluation of $`\Gamma _{12}`$ involves the decay final states which are common to $`\mathrm{K}^0`$ and $`\overline{\mathrm{K}}^0`$, which are $`2\pi (I=0)`$, $`2\pi (I=2)`$, $`3\pi (I=1)`$, $`3\pi (I=2)`$ and $`3\pi (I=3)`$ states:
$$\Gamma _{12}\underset{I=0,2}{}A_{2\pi (I)}^{}\overline{A}{}_{2\pi (I)}{}^{}+\underset{I=1,2,3}{}A_{3\pi (I)}^{}\overline{A}{}_{3\pi (I)}{}^{}.$$
The contribution from the decay amplitude to the $`2\pi (I=2)`$ state is suppressed by the $`\mathrm{\Delta }I=1/2`$ rule and the small measured value of $`ϵ^{}`$. The contribution from the three-pion decay amplitudes are suppressed by $`\Gamma _\mathrm{L}/\Gamma _\mathrm{S}`$ and the measured upper limits for the CP violation parameter for the $`\pi ^+\pi ^{}\pi ^0`$ and $`\pi ^0\pi ^0\pi ^0`$ final states. In conclusion, the phase of $`\Gamma _{12}`$ is essentially given by the phase of the $`A_0`$ amplitude, and it can be expressed as
$$\phi _\Gamma \mathrm{arg}A_0^{}\overline{A}{}_{0}{}^{}=2\phi _0\overline{\theta }{}_{\mathrm{T}}{}^{}+\theta _{\mathrm{CP}}$$
so that
$$|\phi _\Gamma +2\phi _0+\overline{\theta }{}_{\mathrm{T}}{}^{}\theta _{\mathrm{CP}}|<O(10^5).$$
Thus $`|\phi _\Gamma +2\phi _0+\overline{\theta }{}_{\mathrm{T}}{}^{}\theta _{\mathrm{CP}}|<<|ϵ|`$, justifying the approximations made before.
### 3.6 The Standard Model Description
In the framework of the Standard Model , the short range contribution to $`\mathrm{K}^0`$-$`\overline{\mathrm{K}}^0`$ oscillation is obtained from the box diagrams (Figure 4) to be
$$M_{12}^{\mathrm{box}}=\frac{G_\mathrm{F}^2f_\mathrm{K}^2B_\mathrm{K}m_\mathrm{K}m_\mathrm{W}^2}{12\pi ^2}\left[\eta _1\sigma _\mathrm{c}^2S(x_\mathrm{c})+2\eta _2\sigma _\mathrm{c}\sigma _\mathrm{t}E(x_\mathrm{c},x_\mathrm{t})+\eta _3\sigma _\mathrm{t}^2S(x_\mathrm{t})\right]$$
where $`G_\mathrm{F}`$ is the Fermi constant, $`f_\mathrm{K}`$, $`B_\mathrm{K}`$ and $`m_\mathrm{K}`$ are the decay constant, $`B`$ parameter and mass for the K-meson respectively and $`m_\mathrm{W}`$ is the mass of the W-boson. The QCD correction factors are denoted by $`\eta _1=1.38\pm 0.20`$, $`\eta _2=0.57\pm 0.01`$ and $`\eta _3=0.47\pm 0.04`$ and $`S`$ and $`E`$ are known functions of the mass ratios, $`x_\mathrm{i}=m_\mathrm{i}^2/m_\mathrm{W}^2`$ for top (i=t) and charm (i=c). Note that
$$S(x_\mathrm{c})2.4\times 10^4,S(x_\mathrm{t})2.6,E(x_\mathrm{c},x_\mathrm{c})2.2\times 10^3$$
(26)
for $`m_\mathrm{c}=1.25\mathrm{GeV}/c`$, $`m_\mathrm{t}=174\mathrm{GeV}/c^2`$ and $`m_\mathrm{W}=80\mathrm{GeV}/c`$ . The parameters $`\sigma _\mathrm{c}`$ and $`\sigma _\mathrm{t}`$ are the combination of the elements of the Cabibbo-Kobayashi-Maskawa quark mixing matrix (CKM-matrix),
$$V_{\mathrm{CKM}}=\left(\begin{array}{ccc}V_{\mathrm{ud}}& V_{\mathrm{us}}& V_{\mathrm{ub}}\\ V_{\mathrm{cd}}& V_{\mathrm{cs}}& V_{\mathrm{cb}}\\ V_{\mathrm{td}}& V_{\mathrm{ts}}& V_{\mathrm{tb}}\end{array}\right)$$
$`\sigma _\mathrm{c}=V_{\mathrm{cs}}V_{\mathrm{cd}}^{}`$ and $`\sigma _\mathrm{t}=V_{\mathrm{ts}}V_{\mathrm{td}}^{}`$. We adopt the following approximation of the CKM matrix using the parameters introduced by Wolfenstein :
$`V_{\mathrm{CKM}}\left(\begin{array}{ccc}1\lambda ^2/2& \lambda & A\lambda ^3\left(\rho i\eta \right)\\ \lambda iA^2\lambda ^5\eta & 1\lambda ^2/2& A\lambda ^2\\ A\lambda ^3\left(1\stackrel{~}{\rho }i\stackrel{~}{\eta }\right)& A\lambda ^2iA\lambda ^4\eta & 1\end{array}\right)`$ (30)
where where $`\stackrel{~}{\rho }=\rho (1\lambda ^2/2)`$ and $`\stackrel{~}{\eta }=\eta (1\lambda ^2/2)`$. The parameter $`\lambda `$ is known from the light hadron decays to be $`0.221\pm 0.002`$. From the B-meson decays, $`|V_{\mathrm{cb}}|=0.0402\pm 0.0019`$ and $`|V_{\mathrm{ub}}/V_{\mathrm{cb}}|=0.090\pm 0.025`$ are measured , giving $`A=0.823\pm 0.042`$ and $`\sqrt{\rho ^2+\eta ^2}=0.41\pm 0.11`$. The $`B`$-parameter takes in account the difference between $`0|H_\mathrm{W}\mathrm{K}^\pm `$ and $`f|H_\mathrm{W}|\mathrm{K}^0`$ where $`0|`$ is the hadronic vacuum state and $`f|`$ is the common quark states between $`\mathrm{K}^0`$ and $`\overline{\mathrm{K}}^0`$. The theoretical evaluations for this value vary between 0.5 and 1.
In addition to $`M_{12}^{\mathrm{box}}`$, there are large contributions from long range interactions $`M_{12}^{\mathrm{LR}}`$, which are difficult to evaluate. Therefore, theoretical predication for $`M_{12}=M_{12}^{\mathrm{box}}+M_{12}^{\mathrm{LR}}`$ cannot be given. The long range interaction involves only the light flavours and its contribution to $`M_{12}`$ is real in the CKM phase convention; the imaginary part of $`M_{12}`$ is generated only by the box diagram. Therefore we can derive
$$\mathrm{sin}(\phi _M)=\frac{\mathrm{}M_{12}}{|M_{12}|}=\frac{\mathrm{\hspace{0.17em}2}\mathrm{}M_{12}^{\mathrm{box}}}{\mathrm{\Delta }m}.$$
In the CKM phase convention, $`\Gamma _{12}`$ can be approximated as real. Therefore, it follows that
$$\mathrm{}ϵ=\frac{\mathrm{}M_{12}^{\mathrm{box}}}{2\mathrm{\Delta }m}.$$
Although there are considerable uncertainties to evaluate numerically this expression, the currently allowed range of the Wolfenstein parameters, $`\lambda `$, $`A`$, $`\rho `$ and $`\eta `$ gives a consistent value of $`\mathrm{}ϵ`$ with the experimentally measured value.
Prediction of $`ϵ^{}`$ requires an accurate evaluation of the phase difference between $`a_0`$ and $`a_2`$. For the $`a_0`$ amplitudes, the tree, the gluonic penguin and the electroweak penguin diagrams contribute. Only the tree and electroweak penguin diagrams make contributions to the $`a_2`$ decay amplitude. All the penguin diagrams are shown in Figure 5. Not only the short range interactions, but also the hadronic matrix elements with long range interactions have to be evaluated in the calculations. This makes the numerical determination of $`ϵ^{}`$ very difficult. Within the theoretical uncertainties, values of $`ϵ^{}`$ calculated with the currently allowed range of $`\lambda `$, $`A`$, $`\rho `$ and $`\eta `$ are consistent with the data.
### 3.7 CP Violation in Rare Decays
Experimental detection of $`\mathrm{K}_\mathrm{L}\pi ^0\nu \overline{\nu }`$ is clearly very challenging. The final state is a CP eigenstate with $`CP=+1`$. Therefore, observation of this decay is a sign of CP violation. In the Standard Model, the decay is generated by penguin diagrams or box diagrams as shown in Figure 6.
Since the final state consists with only one hadron, long range strong interactions do not play a role and the decay amplitudes can be denoted as
$`\pi ^0\nu \overline{\nu }|H_\mathrm{W}|\mathrm{K}^0`$ $`=`$ $`a_{\pi ^0\nu \overline{\nu }}`$
$`\pi ^0\nu \overline{\nu }|H_\mathrm{W}|\overline{\mathrm{K}}{}_{}{}^{0}`$ $`=`$ $`a_{\pi ^0\nu \overline{\nu }}^{}e^{i(\theta _{\mathrm{CP}}\overline{\theta }{}_{\mathrm{T}}{}^{})}.`$
Unlike for the $`\mathrm{K}^02\pi `$ decays, $`\varphi _{\pi ^0\nu \overline{\nu }}=\mathrm{arg}a_{\pi ^0\nu \overline{\nu }}`$ could be very different from $`\varphi _0`$, so that we could have a situation
$`\left|\mathrm{sin}(\varphi _\Gamma +2\varphi _{\pi ^0\nu \overline{\nu }}+\theta _{\mathrm{CP}}\overline{\theta }{}_{\mathrm{T}}{}^{})\right|`$ $`=`$ $`\left|\mathrm{sin}\left(2\varphi _{\pi ^0\nu \overline{\nu }}2\varphi _0\right)\right|`$
$`>>`$ $`|ϵ|.`$
The $`\mathrm{K}_\mathrm{L}`$ decay amplitude then becomes
$`\pi ^0\nu \overline{\nu }|H_\mathrm{W}|\mathrm{K}_\mathrm{L}`$ $`=`$ $`{\displaystyle \frac{a_{\pi ^0\nu \overline{\nu }}}{\sqrt{\mathrm{\hspace{0.17em}2}}}}\left[1(12ϵ)e^{i(2\varphi _{\pi ^0\nu \overline{\nu }}2\varphi _0)}\right]`$
$``$ $`\sqrt{\mathrm{\hspace{0.17em}2}}i|a_{\pi ^0\nu \overline{\nu }}|\mathrm{sin}(2\varphi _{\pi ^0\nu \overline{\nu }}2\varphi _0).`$
Using isospin symmetry, the hadronic matrix element of the $`\mathrm{K}^0\pi ^0\nu \overline{\nu }`$ decay amplitude and that of the $`\mathrm{K}^+\pi ^+\mathrm{e}^+\nu `$ decay amplitudes can be related as
$$\pi ^0|H_\mathrm{W}|\mathrm{K}^0=\pi ^0|H_\mathrm{W}|\mathrm{K}^+.$$
This allow us to express the branching fractions for $`\mathrm{K}_\mathrm{L}\pi ^0\nu \overline{\nu }`$ using the branching fractions for $`\mathrm{K}^+\pi ^0\mathrm{e}^+\nu `$ as
$`B(\mathrm{K}_\mathrm{L}\pi ^0\nu \overline{\nu })`$ $`=`$ $`{\displaystyle \frac{|\pi ^0\nu \overline{\nu }|H_\mathrm{W}|\mathrm{K}_\mathrm{L}|^2}{\mathrm{\Gamma }_\mathrm{L}}}`$
$`=`$ $`B(\mathrm{K}^+\pi ^0\mathrm{e}^+\nu ){\displaystyle \frac{\tau _\mathrm{L}}{\tau _+}}{\displaystyle \frac{3\alpha ^2\left[\mathrm{}(V_{\mathrm{ts}}^{}V_{\mathrm{td}})X(m_t)\right]^2}{|V_{\mathrm{us}}|^22\pi ^2\mathrm{sin}^4\mathrm{\Theta }_\mathrm{W}}}`$
$`=`$ $`B(\mathrm{K}^+\pi ^0\mathrm{e}^+\nu ){\displaystyle \frac{\tau _\mathrm{L}}{\tau _+}}{\displaystyle \frac{3\alpha ^2\left[X(m_t)\right]^2}{2\pi ^2\mathrm{sin}^4\mathrm{\Theta }_\mathrm{W}}}A^4\lambda ^8(1\lambda ^2/2)^2\eta ^2`$
$``$ $`3\times 10^{11}`$
where $`X`$ is a known function and $`\mathrm{\Theta }_\mathrm{W}`$ is the weak mixing angle. Since the hadronic matrix element is taken from the data, the theoretical uncertainties in this determination is very small. Also the imaginary part of the amplitude is dominated by the short range interactions which can be reliably calculated. Therefore, the theoretical prediction can be considered to be clean.
It is interesting to note that the CP violation parameter
$$\eta _{\pi ^0\nu \overline{\nu }}=\frac{\pi ^0\nu \overline{\nu }|V|\mathrm{K}_\mathrm{L}}{\pi ^0\nu \overline{\nu }|V|\mathrm{K}_\mathrm{S}}$$
as defined in the $`2\pi `$ case has $`|\eta _{\pi ^0\nu \overline{\nu }}|>>|ϵ|`$, although the both final states have $`CP=+1`$.
The current experimental measurement for this branching fraction is $`<5.9\times 10^7`$ with $`90\%`$ confidence by the KTeV experiment , which is still far from the expected number. However, there are several proposals to observe the decays in the near future.
## 4 B-meson System
### 4.1 The Standard Model Description
#### 4.1.1 Some Elements of The CKM Matrix
Among the nine elements of the CKM matrix, five of them related to the third generation play important roles in the B meson system: $`V_{\mathrm{td}}`$, $`V_{\mathrm{ub}}`$, $`V_{\mathrm{ts}}`$, $`V_{\mathrm{cb}}`$ and $`V_{\mathrm{tb}}`$. In the approximation given in equation 30, the phases of the five elements are given by
$$\mathrm{arg}V_{\mathrm{td}}=\varphi _1,\mathrm{arg}V_{\mathrm{ub}}=\varphi _3,\mathrm{arg}V_{\mathrm{ts}}=\delta \varphi _3+\pi ,\mathrm{arg}V_{\mathrm{cb}}=\mathrm{arg}V_{\mathrm{tb}}=0$$
where
$$\varphi _1=\mathrm{tan}^1\frac{\eta }{1\rho },\varphi _3=\mathrm{tan}^1\frac{\eta }{\rho },\delta \varphi _3=\mathrm{tan}^1\lambda ^2\eta .$$
Figure 7 shows the angles in $`\rho `$ and $`\eta `$ planes. Note that $`\varphi _1`$ and $`\varphi _3`$ are often referred to as $`\beta `$ and $`\gamma `$. Clearly $`\delta \varphi _3`$ is very small, $`0.02`$.
#### 4.1.2 Oscillation Amplitude
In the Standard Model, $`\mathrm{B}`$-$`\overline{\mathrm{B}}`$ oscillation is totally governed by the short range interactions, i.e. the box diagrams. Furthermore, only the top quark plays a role in the box diagram due to the large top quark mass (see equation 26) and the structure of the CKM matrix;
$$\frac{\mathrm{}(V_{\mathrm{td}}^{}V_{\mathrm{tb}})}{\mathrm{}(V_{\mathrm{cd}}^{}V_{\mathrm{cb}})}=(\stackrel{~}{\rho }1)1,\frac{\mathrm{}(V_{\mathrm{td}}^{}V_{\mathrm{tb}})}{\mathrm{}(V_{\mathrm{cd}}^{}V_{\mathrm{cb}})}\frac{1}{\lambda ^2}>>1$$
as seen from equation 30.
Therefore, the off diagonal element of the mass matrix, $`M_{12}`$ is given by
$$M_{12}=\frac{G_\mathrm{F}^2f_{\mathrm{B}_\mathrm{d}}^2B_{\mathrm{B}_\mathrm{d}}m_{\mathrm{B}_\mathrm{d}}m_\mathrm{W}^2}{12\pi ^2}\eta _{\mathrm{B}_\mathrm{d}}S(x_\mathrm{t})(V_{\mathrm{td}}^{}V_{\mathrm{tb}})^2\text{for }\mathrm{B}_\mathrm{d}$$
(31)
where $`f_{\mathrm{B}_\mathrm{d}}`$, $`B_{\mathrm{B}_\mathrm{d}}`$ and $`m_{\mathrm{B}_\mathrm{d}}`$ are the decay constant, B-parameter and the mass of the $`\mathrm{B}_\mathrm{d}`$ meson.
Similarly for the $`\mathrm{B}_\mathrm{s}`$ meson, we obtain
$$M_{12}=\frac{G_\mathrm{F}^2f_{\mathrm{B}_\mathrm{s}}^2B_{\mathrm{B}_\mathrm{s}}m_{\mathrm{B}_\mathrm{s}}m_\mathrm{W}^2}{12\pi ^2}\eta _{\mathrm{B}_\mathrm{s}}S(x_\mathrm{t})(V_{\mathrm{ts}}^{}V_{\mathrm{tb}})^2\text{for }\mathrm{B}_\mathrm{s}$$
where $`f_{\mathrm{B}_\mathrm{s}}`$, $`B_{\mathrm{B}_\mathrm{s}}`$ and $`m_{\mathrm{B}_\mathrm{s}}`$ are the decay constant, B-parameter and the mass of the $`\mathrm{B}_\mathrm{s}`$ meson.
The phase of $`M_{12}`$ is then given by
$$\mathrm{arg}M_{12}=\{\begin{array}{cc}\mathrm{arg}(V_{\mathrm{td}}^{}V_{\mathrm{tb}})^2+\pi =2\varphi _1+\pi \hfill & \text{for }\mathrm{B}_\mathrm{d}\hfill \\ \mathrm{arg}(V_{\mathrm{ts}}^{}V_{\mathrm{tb}})^2+\pi =2\delta \varphi _3+\pi \hfill & \text{for }\mathrm{B}_\mathrm{s}\text{.}\text{}\hfill \end{array}$$
The parameter $`\Gamma _{12}`$ can also be determined by taking the absorptive part of the box diagrams with charm and up quarks in the loops. The phase difference between $`M_{12}`$ and $`\Gamma _{12}`$ is given by
$$\mathrm{arg}M_{12}\mathrm{arg}\Gamma _{12}=\pi +\frac{\mathrm{\hspace{0.17em}8}}{3}\left(\frac{m_\mathrm{c}}{m_\mathrm{b}}\right)^2\eta \times \{\begin{array}{cc}\frac{\text{}1}{\text{}\left(1\rho \right)^2+\eta ^2}& :\mathrm{B}_\mathrm{d}\hfill \\ \lambda ^2& :\mathrm{B}_\mathrm{s}\hfill \end{array}$$
(32)
i.e. $`\mathrm{sin}(\mathrm{arg}M_{12}\mathrm{arg}\Gamma _{12})`$ is small for $`\mathrm{B}_\mathrm{d}`$ and very small for $`\mathrm{B}_\mathrm{s}`$. Note that $`M_{12}`$ and $`\Gamma _{12}`$ are antiparallel. Therefore, the approximations for $`\zeta `$, $`m_\pm `$ and $`\Gamma _\pm `$ given on page 2.3 are valid with $`n=1`$. Since we will rely on the Standard Model description of $`M_{12}`$, and our experimental knowledge of the decay amplitudes is still limited, we adopt b) $`\varphi _M`$ base. We refer the mass eigenstate with larger mass as $`\mathrm{B}_\mathrm{h}`$ (B-heavy) and the other $`\mathrm{B}_\mathrm{l}`$ (B-light) with their masses and decay width are given by:
$$m_\mathrm{h}=M+|M_{12}|,\Gamma _\mathrm{h}=\Gamma |\Gamma _{12}|$$
and
$$m_\mathrm{l}=M|M_{12}|,\Gamma _\mathrm{l}=\Gamma +|\Gamma _{12}|$$
respectively, and $`\mathrm{B}_\mathrm{h}`$ ($`\mathrm{B}_\mathrm{l}`$) corresponds to $`\mathrm{P}_+`$ ($`\mathrm{P}_{}`$) defined in equation 13.
For both $`\mathrm{B}_\mathrm{d}`$ and $`\mathrm{B}_\mathrm{s}`$, we can now derive
$$\frac{\mathrm{\Delta }\Gamma }{\mathrm{\Delta }m}=\left|\frac{\Gamma _{12}}{M_{12}}\right|\frac{3\pi m_\mathrm{b}^2}{2m_\mathrm{W}^2S(x_\mathrm{t})}5\times 10^3\text{for }\mathrm{B}_\mathrm{d}\text{ and }\mathrm{B}_\mathrm{s}$$
(33)
for $`m_\mathrm{b}=4.25`$ GeV, $`m_\mathrm{W}=80`$ GeV and $`m_\mathrm{t}=174`$ GeV, where $`\mathrm{\Delta }m`$ and $`\mathrm{\Delta }\Gamma `$ are defined as positive:
$$\mathrm{\Delta }m=m_\mathrm{h}m_\mathrm{l},\mathrm{\Delta }\Gamma =\Gamma _\mathrm{l}\Gamma _\mathrm{h}.$$
Using the measured values of $`\mathrm{\Delta }m=(0.464\pm 0.018)\times 10^{12}\mathrm{}\mathrm{s}^1`$ and the average lifetime $`\tau =1/\widehat{\Gamma }=(1.54\pm 0.03)\times 10^{12}`$ s for the $`\mathrm{B}_\mathrm{d}`$ mesons, where $`\widehat{\Gamma }`$ is the averaged decay width, it follows that
$$\frac{\mathrm{\Delta }\Gamma }{\widehat{\Gamma }}4\times 10^3\text{for }\mathrm{B}_\mathrm{d}$$
and $`\mathrm{\Delta }\Gamma `$ can be neglected in the decay time distribution for the $`\mathrm{B}_\mathrm{d}`$ system. For the $`\mathrm{B}_\mathrm{s}`$ mesons, using the measured lifetime $`(1.54\pm 0.07)\times 10^{12}`$ s, it follows that
$$\frac{\mathrm{\Delta }\Gamma }{\widehat{\Gamma }}0.1\text{for }\mathrm{B}_\mathrm{s}.$$
The effect of $`\mathrm{\Delta }\Gamma `$ is still not large, but can no longer be neglected in the decay time distributions.
The small decay width differences of the $`\mathrm{B}_\mathrm{d}`$ and $`\mathrm{B}_\mathrm{s}`$ systems do not allow to separate one mass-eigenstate from the other, which can be done for the kaon system by creating a $`\mathrm{K}_\mathrm{L}`$ beam. Therefore, CP violation cannot be established by just observing the decays as in the case of $`\mathrm{K}_\mathrm{L}2\pi `$. We either have to compare the decay rates of the initial $`\mathrm{B}^0`$ and initial $`\overline{\mathrm{B}}^0`$ states or measure the time dependent decay rates of at least one of the two cases, i.e. either initial $`\mathrm{B}^0`$ or $`\overline{\mathrm{B}}^0`$.
Since $`\mathrm{\Delta }m=2|M_{12}|`$, one can extract
$$|V_{\mathrm{td}}|^2=A^4\lambda ^6\left[(1\stackrel{~}{\rho })^2+\stackrel{~}{\eta }^2\right]$$
i.e. $`\rho `$ and $`\eta `$, from the measured $`\mathrm{B}^0`$-$`\overline{\mathrm{B}}^0`$ oscillation frequency $`\Delta m_\mathrm{d}`$ using equation 31. However, theoretical uncertainties in calculating the decay constant and B-parameter are considerable and limit the accuracy on the extracted value of $`|V_{\mathrm{td}}|^2`$. If the $`\mathrm{B}_\mathrm{s}^0`$-$`\overline{\mathrm{B}}_\mathrm{s}^0`$ oscillation frequency $`\Delta m_\mathrm{s}=2|M_{12}^\mathrm{s}|`$ is measured, $`|V_{\mathrm{td}}|^2`$ can be determined with much small uncertainty by using the ratio $`\Delta m_\mathrm{d}/\Delta m_\mathrm{s}`$, due to better controlled theoretical errors in $`f_{\mathrm{B}_\mathrm{d}}/f_{\mathrm{B}_\mathrm{s}}`$ and $`B_{\mathrm{B}_\mathrm{d}}/B_{\mathrm{B}_\mathrm{s}}`$. However, the frequency of the $`\mathrm{B}_\mathrm{s}^0`$-$`\overline{\mathrm{B}}_\mathrm{s}^0`$ oscillation is expected to be $`>1/\lambda ^2=20`$ times larger than that of the $`\mathrm{B}^0`$-$`\overline{\mathrm{B}}^0`$ oscillation and we still have to wait for sometime before it is measured.
Since $`|M_{12}/\Gamma _{12}|<<1`$, $`\zeta `$ given by equation 19 can be further approximated as
$`\zeta \left[1{\displaystyle \frac{1}{\mathrm{\hspace{0.17em}2}}}\mathrm{}\left({\displaystyle \frac{\Gamma _{12}}{M_{12}}}\right)\right]e^{i\phi _M}`$ (34)
where $`\phi _M=\mathrm{arg}M_{12}`$ as before. Seen from equation 32 and 33, the approximation $`|\zeta |1`$ is accurate to $`10^3`$ or better.
Similar to the kaon system, CP violation (and T violation) in the oscillation can be measured from the time-dependent rate asymmetry between the initial $`\overline{\mathrm{B}}^0`$ decaying into semileptonic final states with $`\mathrm{e}^+`$ or $`\mu ^+`$, $`\overline{R}{}_{+}{}^{}(t)`$ and the initial $`\mathrm{B}^0`$ decaying into semileptonic final states with $`\mathrm{e}^{}`$ or $`\mu ^{}`$, $`R_{}(t)`$. The asymmetry is given by
$$\frac{\overline{R}{}_{+}{}^{}(t)R_{}(t)}{\overline{R}{}_{+}{}^{}(t)+R_{}(t)}=\frac{1|\zeta |^4}{1+|\zeta |^4}O(10^3)\text{for }\mathrm{B}_\mathrm{d}\text{ and }<<O(10^3)\text{for }\mathrm{B}_\mathrm{d}$$
which is a very small signal.
From now on, we assume
$$\zeta =e^{i\phi _M}$$
for both $`\mathrm{B}_\mathrm{d}`$ and $`\mathrm{B}_\mathrm{s}`$ and $`\mathrm{\Delta }\Gamma =0`$ for $`\mathrm{B}_\mathrm{d}`$.
In summary, the two mass eigenstates are given by
$`|\mathrm{B}_\mathrm{h}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left[|\mathrm{B}+e^{i\phi _M}|\overline{\mathrm{B}}\right]`$
$`|\mathrm{B}_\mathrm{l}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left[|\mathrm{B}e^{i\phi _M}|\overline{\mathrm{B}}\right]`$
and
$$m_\mathrm{h}=m_0+|M_{12}|,m_\mathrm{l}=m_0|M_{12}|,\mathrm{\Delta }m=m_\mathrm{h}m_\mathrm{l}$$
for $`\mathrm{B}_\mathrm{d}`$ and $`\mathrm{B}_\mathrm{s}`$. For the decay width, we have
$$\begin{array}{cc}\Gamma _\mathrm{l}=\Gamma _\mathrm{h}\hfill & \text{for }\mathrm{B}_\mathrm{d}\hfill \\ \Gamma _\mathrm{l}=\Gamma _0+|\Gamma _{12}|,\Gamma _\mathrm{h}=\Gamma _0|\Gamma _{12}|,\mathrm{\Delta }\Gamma =\Gamma _\mathrm{l}\Gamma _\mathrm{h}\hfill & \text{for }\mathrm{B}_\mathrm{s}\hfill \end{array}$$
#### 4.1.3 Time Dependent Decay Rates
Since $`\mathrm{\Delta }\Gamma `$ is small in the B meson system, it is more convenient to derive the time dependent decay rate from the particle-antiparticle base rather than the mass eigenstate base. Using, equations 10 and 14 the time dependent decay rates for the final state f can be derived as
$`R_\mathrm{f}(t)`$ $``$ $`{\displaystyle \frac{|A_\mathrm{f}|^2}{2}}e^{\widehat{\Gamma }t}\left[I_+(t)+I_{}(t)\right]`$ (35)
$`\overline{R}_\mathrm{f}(t)`$ $``$ $`{\displaystyle \frac{|A_\mathrm{f}|^2}{\mathrm{\hspace{0.17em}2}|\zeta |^2}}e^{\widehat{\Gamma }t}\left[I_+(t)I_{}(t)\right]`$ (36)
where $`\widehat{\Gamma }`$ is the averaged decay time, $`\widehat{\Gamma }=(\Gamma _++\Gamma _{})/2`$, and $`A_\mathrm{f}`$ is the instantaneous decay amplitude for the $`\mathrm{P}^0\mathrm{f}`$ decays. The two time dependent functions, $`I_+(t)`$ and $`I_{}(t)`$, are given by
$`I_+(t)`$ $`=`$ $`(1+|L_\mathrm{f}|^2)\mathrm{cosh}{\displaystyle \frac{\mathrm{\Delta }\Gamma }{2}}t+2\mathrm{}L_\mathrm{f}\mathrm{sinh}{\displaystyle \frac{\mathrm{\Delta }\Gamma }{2}}t`$
$`I_{}(t)`$ $`=`$ $`(1|L_\mathrm{f}|^2)\mathrm{cos}\mathrm{\Delta }mt+2\mathrm{}L_\mathrm{f}\mathrm{sin}\mathrm{\Delta }mt.`$
The parameter $`L_\mathrm{f}`$ is given by
$$L_\mathrm{f}=\zeta \frac{\overline{A}_\mathrm{f}}{A_\mathrm{f}}$$
where $`\overline{A}_\mathrm{f}`$ is the instantaneous decay amplitude for the $`\overline{\mathrm{P}}{}_{}{}^{0}\mathrm{f}`$ decays.
The time dependent decay rate for the CP conjugated final states $`\mathrm{f}^{\mathrm{CP}}`$ are derived to be
$`\overline{R}_{\mathrm{f}^{\mathrm{CP}}}(t)`$ $``$ $`{\displaystyle \frac{|\overline{A}_{\mathrm{f}^{\mathrm{CP}}}|^2}{2}}e^{\overline{\Gamma }t}\left[\overline{I}_+^{\mathrm{CP}}(t)+\overline{I}_{}^{\mathrm{CP}}(t)\right]`$ (37)
$`R_{\mathrm{f}^{\mathrm{CP}}}(t)`$ $``$ $`{\displaystyle \frac{|\overline{A}_{\mathrm{f}^{\mathrm{CP}}}|^2|\zeta |^2}{2}}e^{\overline{\Gamma }t}\left[\overline{I}_+^{\mathrm{CP}}(t)\overline{I}_{}^{\mathrm{CP}}(t)\right]`$ (38)
where $`\overline{A}_{\mathrm{f}^{\mathrm{CP}}}`$ is the instantaneous decay amplitude for the $`\overline{\mathrm{P}}{}_{}{}^{0}\mathrm{f}^{\mathrm{CP}}`$ decays. Two time dependent decay rates, $`\overline{I}_+^{\mathrm{CP}}(t)`$ and $`\overline{I}_{}^{\mathrm{CP}}(t)`$ are given by
$`\overline{I}_+^{\mathrm{CP}}(t)`$ $`=`$ $`(1+|L_\mathrm{f}^{\mathrm{CP}}|^2)\mathrm{cosh}{\displaystyle \frac{\mathrm{\Delta }\Gamma }{2}}t+2\mathrm{}L_\mathrm{f}^{\mathrm{CP}}\mathrm{sinh}{\displaystyle \frac{\mathrm{\Delta }\Gamma }{2}}t`$
$`\overline{I}_{}^{\mathrm{CP}}(t)`$ $`=`$ $`(1|L_\mathrm{f}^{\mathrm{CP}}|^2)\mathrm{cos}\mathrm{\Delta }mt+2\mathrm{}L_\mathrm{f}^{\mathrm{CP}}\mathrm{sin}\mathrm{\Delta }mt`$
where the parameter, $`L_\mathrm{f}^{\mathrm{CP}}`$, is given by
$$L_\mathrm{f}^{\mathrm{CP}}=\frac{1}{\zeta }\frac{A_{\mathrm{f}^{\mathrm{CP}}}}{\overline{A}_{\mathrm{f}^{\mathrm{CP}}}}$$
and $`A_{\mathrm{f}^{\mathrm{CP}}}`$ is the instantaneous decay amplitude for the $`\mathrm{P}^0\mathrm{f}^{\mathrm{CP}}`$ decays.
The decay rates $`R_f(t)`$ and $`\overline{R}_{\mathrm{f}^{\mathrm{CP}}}(t)`$ are CP conjugate to each other and so are $`\overline{R}_f(t)`$ and $`R_{\mathrm{f}^{\mathrm{CP}}}(t)`$. If there exists any difference between the CP conjugated processes, this is a clear sign of CP violation.
The final state f can be classified into the following four different cases:
1. Flavour specific final state ($`A_\mathrm{f}=\overline{A}_{\mathrm{f}^{\mathrm{CP}}}=0`$ or $`A_{\mathrm{f}^{\mathrm{CP}}}=\overline{A}_\mathrm{f}=0`$)
2. Flavour non specific final state
1. CP eigenstate ($`A_\mathrm{f}=A_{\mathrm{f}^{\mathrm{CP}}}`$ and $`\overline{A}_\mathrm{f}=\overline{A}_{\mathrm{f}^{\mathrm{CP}}}`$)
2. mixed CP eigenstate ($`A_\mathrm{f}=A_{\mathrm{f}^{\mathrm{CP}}}`$ and $`\overline{A}_\mathrm{f}=\overline{A}_{\mathrm{f}^{\mathrm{CP}}}`$)
3. CP non eigenstate
#### 4.1.4 CP Violation: Clean Case
The contribution to the $`\mathrm{B}^0`$ decaying into $`\mathrm{J}/\psi \mathrm{K}_\mathrm{S}`$ is dominated by the tree diagram with $`V_{\mathrm{cb}}^{}V_{\mathrm{cs}}`$. Although there exist some contribution from the penguin diagrams, the dominant penguin diagram contribution has the CKM phase $`V_{\mathrm{tb}}^{}V_{\mathrm{ts}}`$ which is close to that of the tree diagram (Figure 8). Thus, we can safely assume that there is no CP violation in the decay amplitude and the ratio of the $`\overline{\mathrm{B}}^0`$ and $`\mathrm{B}^0`$ decay amplitudes is given only by the CKM part. By noting that $`CP(\mathrm{J}/\psi \mathrm{K}_\mathrm{S})=1`$ we obtain
$$\frac{A(\overline{\mathrm{B}}{}_{}{}^{0}\mathrm{J}/\psi \mathrm{K}_\mathrm{S})}{A(\mathrm{B}^0\mathrm{J}/\psi \mathrm{K}_\mathrm{S})}=\frac{(V_{\mathrm{cb}}^{}V_{\mathrm{cs}}V_{\mathrm{us}}^{}V_{\mathrm{ud}})^2}{|V_{\mathrm{cb}}^{}V_{\mathrm{cs}}V_{\mathrm{us}}^{}V_{\mathrm{ud}}|^2}.$$
Using the formulae developed in the previous section, the time dependent rates for the initial $`\mathrm{B}^0`$ decaying into $`\mathrm{J}/\psi \mathrm{K}_\mathrm{S}`$, $`R_{\mathrm{J}/\psi \mathrm{K}_\mathrm{S}}(t)`$, and that for $`\overline{\mathrm{B}}^0`$ decaying into $`\mathrm{J}/\psi \mathrm{K}_\mathrm{S}`$, $`\overline{R}{}_{\mathrm{J}/\psi \mathrm{K}_\mathrm{S}}{}^{}(t)`$ are given by
$`R_{\mathrm{J}/\psi \mathrm{K}_\mathrm{S}}(t)e^{\widehat{\Gamma }t}\left(1+\mathrm{}L_{\mathrm{J}/\psi \mathrm{K}_\mathrm{S}}\mathrm{sin}\mathrm{\Delta }mt\right)`$
$`\overline{R}{}_{\mathrm{J}/\psi \mathrm{K}_\mathrm{S}}{}^{}(t)e^{\widehat{\Gamma }t}(1\mathrm{}L_{\mathrm{J}/\psi \mathrm{K}_\mathrm{S}}\mathrm{sin}\mathrm{\Delta }mt)`$
which allow to extract
$$\mathrm{}L_{\mathrm{J}/\psi \mathrm{K}_\mathrm{S}}=\mathrm{}\left(\zeta \times \frac{A(\overline{\mathrm{B}}{}_{}{}^{0}\mathrm{J}/\psi \mathrm{K}_\mathrm{S})}{A(\mathrm{B}^0\mathrm{J}/\psi \mathrm{K}_\mathrm{S})}\right)=\mathrm{}\left[\frac{(V_{\mathrm{td}}^{}V_{\mathrm{tb}}V_{\mathrm{cb}}^{}V_{\mathrm{cs}}V_{\mathrm{us}}^{}V_{\mathrm{ud}})^2}{|V_{\mathrm{td}}^{}V_{\mathrm{tb}}V_{\mathrm{cb}}^{}V_{\mathrm{cs}}V_{\mathrm{us}}^{}V_{\mathrm{ud}}|^2}\right]$$
With the Wolfenstein parameterization, it follows that
$$\mathrm{}L_{\mathrm{J}/\psi \mathrm{K}_\mathrm{S}}=\mathrm{sin}2\varphi _1.$$
The same argument holds for the $`\mathrm{B}_\mathrm{s}\mathrm{J}/\psi \varphi `$ decays and from the time dependent decay rates
$`R_{\mathrm{J}/\psi \varphi }(t)e^{\widehat{\Gamma }t}\left(\mathrm{cosh}{\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }}{2}}t+2\mathrm{}L_{\mathrm{J}/\psi \varphi }\mathrm{sinh}{\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }}{2}}t+\mathrm{}L_{\mathrm{J}/\psi \varphi }\mathrm{sin}\mathrm{\Delta }mt\right)`$
$`\overline{R}{}_{\mathrm{J}/\psi \varphi }{}^{}(t)e^{\widehat{\Gamma }t}(\mathrm{cosh}{\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }}{2}}t+2\mathrm{}L_{\mathrm{J}/\psi \varphi }\mathrm{sinh}{\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }}{2}}t\mathrm{},L_{\mathrm{J}/\psi \varphi }\mathrm{sin}\mathrm{\Delta }mt)`$
one can extract
$$\mathrm{}L_{\mathrm{J}/\psi \varphi }=\mathrm{}\left[\zeta \times \frac{A(\overline{\mathrm{B}}{}_{\mathrm{s}}{}^{0}\mathrm{J}/\psi \varphi )}{A(\mathrm{B}_\mathrm{s}^0\mathrm{J}/\psi \varphi )}\right]=\mathrm{sin}2\delta \varphi _3$$
Note that we assumed in the calculation above that $`CP(\mathrm{J}/\psi \varphi )=+1`$, i.e. the $`\mathrm{J}/\psi \varphi `$ state is in the lowest orbital angular momentum state of $`l=0`$. If there exists the $`l=1`$ state with $`CP(\mathrm{J}/\psi \varphi )=1`$, the measured $`\mathrm{}L_{\mathrm{J}/\psi \varphi }`$ will be diluted and the fraction of the $`CP=1`$ state must be experimentally measured. If there is the same amount of $`CP=+1`$ state and $`CP=1`$ state, $`\mathrm{}L_{\mathrm{J}/\psi \varphi }`$ will vanish.
An even cleaner decay channel is $`\mathrm{B}^0\mathrm{D}^{}\pi ^\pm `$. There is only one tree diagram, $`\overline{\mathrm{b}}\overline{\mathrm{c}}+\mathrm{W}^+`$ followed by $`\mathrm{W}^+\mathrm{u}+\overline{\mathrm{d}}`$, which contributes to the $`\mathrm{B}^0\mathrm{D}^{}\pi ^+`$ decays. The same final state can be produced from the $`\overline{\mathrm{B}}^0`$ decays with another tree diagram, $`\mathrm{b}\mathrm{u}+\mathrm{W}^{}`$ followed by $`\mathrm{W}^{}\overline{\mathrm{c}}+\mathrm{d}`$ (Figure 9). Therefore, the time dependent rate for the initial $`\mathrm{B}^0`$ decaying into $`\mathrm{D}^{}\pi ^+`$ is given by
$$R_\mathrm{D}^{}(t)e^{\widehat{\Gamma }t}\left[1+\frac{(1|L_{\mathrm{D}^{}\pi ^+}|^2)}{(1+|L_{\mathrm{D}^{}\pi ^+}|^2)}\mathrm{cos}\mathrm{\Delta }mt+\frac{2\mathrm{}L_{\mathrm{D}^{}\pi ^+}}{(1+|L_{\mathrm{D}^{}\pi ^+}|^2)}\mathrm{sin}\mathrm{\Delta }mt\right]$$
where
$$L_{\mathrm{D}^{}\pi ^+}=\zeta \times \frac{A(\overline{\mathrm{B}}{}_{}{}^{0}\mathrm{D}^{}\pi ^+)}{A(\mathrm{B}^0\mathrm{D}^{}\pi ^+)}$$
The weak phase of $`A(\overline{\mathrm{B}}{}_{}{}^{0}\mathrm{D}^{}\pi ^+)`$ is given by $`V_{\mathrm{ub}}V_{\mathrm{cd}}^{}`$ and that of $`A(\mathrm{B}^0\mathrm{D}^{}\pi ^+)`$ by $`V_{\mathrm{cb}}^{}V_{\mathrm{ud}}`$. The phase of $`L_{\mathrm{D}^{}\pi ^+}`$ is then derived to be
$`\mathrm{arg}L_{\mathrm{D}^{}\pi ^+}`$ $`=`$ $`\mathrm{arg}V_{\mathrm{ub}}\mathrm{arg}M_{12}+\phi _\mathrm{S}`$
$`=`$ $`\varphi _3+2\varphi _1+\phi _\mathrm{S}`$
where $`\phi _\mathrm{S}`$ is a possible strong phase difference between the $`\mathrm{b}\mathrm{u}+\mathrm{W}^{}`$ and $`\overline{\mathrm{b}}\overline{\mathrm{c}}+\mathrm{W}^+`$ tree diagrams.
CP conjugated decay amplitudes of $`A(\overline{\mathrm{B}}{}_{}{}^{0}\mathrm{D}^{}\pi ^+)`$ and $`A(\mathrm{B}^0\mathrm{D}^{}\pi ^+)`$, i.e. $`A(\mathrm{B}^0\mathrm{D}^+\pi ^{})`$ and $`A(\overline{\mathrm{B}}{}_{}{}^{0}\mathrm{D}^+\pi ^{})`$ respectively, are obtained by taking the complex conjugate of the weak amplitudes while the strong phase remains unchanged. Thus for $`\mathrm{D}^+\pi ^{}`$ we obtain
$$R_{\mathrm{D}^+}(t)e^{\widehat{\Gamma }t}\left[1\frac{(1|L_{\mathrm{D}^{}\pi ^+}^{\mathrm{CP}}|^2)}{(1+|L_{\mathrm{D}^{}\pi ^+}^{\mathrm{CP}}|^2)}\mathrm{cos}\mathrm{\Delta }mt\frac{2\mathrm{}L_{\mathrm{D}^{}\pi ^+}^{\mathrm{CP}}}{(1+|L_{\mathrm{D}^{}\pi ^+}^{\mathrm{CP}}|^2)}\mathrm{sin}\mathrm{\Delta }mt\right]$$
where
$$L_{\mathrm{D}^{}\pi ^+}^{\mathrm{CP}}=\frac{1}{\zeta }\times \frac{A(\mathrm{B}^0\mathrm{D}^+\pi ^{})}{A(\overline{\mathrm{B}}{}_{}{}^{0}\mathrm{D}^+\pi ^{})}$$
and the phase of $`L_{\mathrm{D}^{}\pi ^+}^{\mathrm{CP}}`$ is given by
$`\mathrm{arg}L_{\mathrm{D}^{}\pi ^+}^{\mathrm{CP}}`$ $`=`$ $`\mathrm{arg}V_{\mathrm{ub}}+\mathrm{arg}M_{12}+\phi _\mathrm{S}`$
$`=`$ $`\varphi _32\varphi _1+\phi _\mathrm{S}`$
From the two time dependent decay rates, we can extract $`\varphi _32\varphi _1`$.
Note that
$$\left|L_{\mathrm{D}^{}\pi ^+}\right|=\left|L_{\mathrm{D}^{}\pi ^+}^{\mathrm{CP}}\right|\left|\frac{V_{\mathrm{ub}}V_{\mathrm{cd}}^{}}{V_{\mathrm{cb}}^{}V_{\mathrm{ud}}}\right|=\lambda ^2\sqrt{\rho ^2\eta ^2}<<1$$
i.e. the effect we have to measure is small.
The CP conjugated time dependent decay rate distributions are given by
$$\overline{R}_{\mathrm{D}^+}(t)e^{\widehat{\Gamma }t}\left[1+\frac{(1|L_{\mathrm{D}^{}\pi ^+}^{\mathrm{CP}}|^2)}{(1+|L_{\mathrm{D}^{}\pi ^+}^{\mathrm{CP}}|^2)}\mathrm{cos}\mathrm{\Delta }mt+\frac{2\mathrm{}L_{\mathrm{D}^{}\pi ^+}^{\mathrm{CP}}}{(1+|L_{\mathrm{D}^{}\pi ^+}^{\mathrm{CP}}|^2)}\mathrm{sin}\mathrm{\Delta }mt\right]$$
and
$$\overline{R}_\mathrm{D}^{}(t)e^{\widehat{\Gamma }t}\left[1\frac{(1|L_{\mathrm{D}^{}\pi ^+}|^2)}{(1+|L_{\mathrm{D}^{}\pi ^+}|^2)}\mathrm{cos}\mathrm{\Delta }mt\frac{2\mathrm{}L_{\mathrm{D}^{}\pi ^+}}{(1+|L_{\mathrm{D}^{}\pi ^+}|^2)}\mathrm{sin}\mathrm{\Delta }mt\right]$$
which can be used to obtain the same information.
A similar method can be used for the $`\mathrm{B}_\mathrm{s}^0\mathrm{D}_\mathrm{s}^{}\mathrm{K}^\pm `$ decays to extract $`\varphi _32\delta \varphi _3`$. The effect is larger since
$$|L_{\mathrm{D}_\mathrm{s}^{}\mathrm{K}^+}|\left|\frac{V_{\mathrm{ub}}V_{\mathrm{cs}}^{}}{V_{\mathrm{cb}}^{}V_{\mathrm{us}}}\right|=\sqrt{\rho ^2+\eta ^2}=O(1).$$
#### 4.1.5 CP Violation: Not So Clean Case
The penguin contribution to the $`\mathrm{B}_\mathrm{d}\pi ^+\pi ^{}`$ decay was originally thought to be small and the decay would be dominated by the $`\mathrm{b}\mathrm{u}+\mathrm{W}`$ tree diagram. However, the discovery of $`B(\mathrm{B}_\mathrm{d}\mathrm{K}^\pm \pi ^{})>B(\mathrm{B}_\mathrm{d}\pi ^+\pi ^{})`$ indicates that the contribution of the penguin diagrams to the $`\mathrm{B}_\mathrm{d}\pi ^+\pi ^{}`$ amplitude should be $`20\%`$ or more.
Due to the penguin contribution, the phase of the $`\mathrm{B}^0\pi ^+\pi ^{}`$ decay amplitude deviates from that of $`V_{\mathrm{ub}}^{}`$. Furthermore, CP violation in the decay amplitude could be present. Evaluation of those effects involves calculating contributions from different diagrams accurately. Strong interactions may play an important role as well. Therefore, this decay mode may not be ideal to make precise determinations of $`\rho `$ and $`\eta `$ from CP violation.
### 4.2 Case with New Physics
Decay processes where only the tree diagrams contribute should be unaffected by the presence of physics beyond the Standard Model. Therefore, $`|V_{\mathrm{cb}}|`$ and $`|V_{\mathrm{ub}}|`$ obtained from the semileptonic decays of B mesons would not be affected by the new physics and $`A`$ and $`\rho ^2+\eta ^2`$ can be obtained even if physics beyond the Standard Model is present.
New physics could generate $`\mathrm{B}^0`$-$`\overline{\mathrm{B}}^0`$ and $`\mathrm{B}_\mathrm{s}^0`$-$`\overline{\mathrm{B}}_\mathrm{s}^0`$ oscillations by new particles generating new box diagrams. They could also generate a tree level flavour changing neutral current contributing to the oscillation. Since these contributions are through “virtual” states, they contribute to $`M_{12}`$ with little effect on $`\Gamma _{12}`$, i.e.
$$M_{12}=M_{12}^{\mathrm{SM}}+M_{12}^{\mathrm{NP}},\Gamma _{12}=\Gamma _{12}^{\mathrm{SM}}$$
where $`M_{12}^{\mathrm{SM}}`$ and $`\Gamma _{12}^{\mathrm{SM}}`$ are due to the Standard Model and $`M_{12}^{\mathrm{NP}}`$ is the contribution from the new physics. The measured $`\mathrm{\Delta }m`$ is given by $`2|M_{12}|`$ and can no longer used to extract $`|V_{\mathrm{td}}|^2`$ due to $`M_{12}^{\mathrm{NP}}`$.
Since
$$\left|\frac{\Gamma _{12}}{M_{12}}\right|=\frac{2\left|\Gamma _{12}^{\mathrm{SM}}\right|}{\mathrm{\Delta }m}$$
remains small, CP violation in the oscillation remains small as seen from equation 34. Therefore,
$$\zeta =e^{i\phi _M}$$
is still valid. However, note that
$$\phi _M\mathrm{arg}M_{12}\mathrm{arg}M_{12}^{\mathrm{SM}}.$$
Decay amplitudes from the penguin diagrams can be affected by physics beyond the Standard Model since new particles can contribute virtually in the loop. Therefore, the modes such as $`\mathrm{B}_\mathrm{d}`$ decaying into $`\pi ^+\pi ^{}`$, $`\mathrm{K}^\pm \pi ^{}`$ may have some contribution from the new physics.
Since the decays $`\mathrm{B}_\mathrm{d}\mathrm{J}/\psi \mathrm{K}_\mathrm{S}`$ and $`\mathrm{B}_\mathrm{s}\mathrm{J}/\psi \varphi `$ are tree dominated, they are little affected by new physics. Therefore we have
$$\frac{A(\overline{\mathrm{B}}{}_{}{}^{0}\mathrm{J}/\psi \mathrm{K}_\mathrm{S})}{A(\mathrm{B}^0\mathrm{J}/\psi \mathrm{K}_\mathrm{S})}=\frac{A(\mathrm{B}_\mathrm{s}^0\mathrm{J}/\psi \varphi )}{A(\overline{\mathrm{B}}{}_{\mathrm{s}}{}^{0}\mathrm{J}/\psi \varphi )}=1$$
with the phase convention due to the Wolfenstein parameterization and
$$L_{\mathrm{J}/\psi \mathrm{K}_\mathrm{S},\mathrm{J}/\psi \varphi }=e^{i\phi _M}:\text{ for }\mathrm{B}_\mathrm{d}\mathrm{J}/\psi \mathrm{K}_\mathrm{S}\text{ and }+\text{ for }\mathrm{B}_\mathrm{s}\mathrm{J}/\psi \varphi $$
and studies of the time dependent decay rates give $`\mathrm{arg}M_{12}`$.
The $`\mathrm{B}_\mathrm{d}\mathrm{D}^{}\pi `$ and $`\mathrm{B}_\mathrm{s}\mathrm{D}_\mathrm{s}\mathrm{K}`$ decays are generated by only the tree diagrams and not affected by new physics. Therefore we have
$$\mathrm{arg}L_{\mathrm{D}^{}\pi ^+}=\varphi _3\mathrm{arg}M_{12}+\phi _\mathrm{S}$$
and
$$\mathrm{arg}L_{\mathrm{D}^+\pi ^{}}=\varphi _3+\mathrm{arg}M_{12}+\phi _\mathrm{S}$$
and studies of the time dependent decay rates provide $`\mathrm{arg}M_{12}+\varphi _3`$. Similarly studies can be done for $`\mathrm{B}_\mathrm{s}\mathrm{D}_\mathrm{s}\mathrm{K}`$.
By combining the measurements of $`\mathrm{B}_\mathrm{d}\mathrm{J}/\psi \mathrm{K}_\mathrm{S}`$ and $`\mathrm{D}^{}\pi `$ or $`\mathrm{B}_\mathrm{s}\mathrm{J}/\psi \varphi `$ and $`\mathrm{D}_\mathrm{s}\mathrm{K}`$, the angle $`\varphi _3`$ can be determined even with presence of physics beyond the Standard Model. By comparing the result from $`\mathrm{B}_\mathrm{d}`$ and that from $`\mathrm{B}_\mathrm{s}`$, consistency of the method can be tested. Since the phase of $`V_{\mathrm{ub}}`$ is given by $`\varphi _3`$ and its modulus is measured from the semileptonic decay, $`\rho `$ and $`\eta `$ can be extracted. Once $`\lambda `$, $`A`$, $`\rho `$ and $`\eta `$ are known, $`M_{12}^{\mathrm{SM}}`$ can be calculated and from the measured $`\mathrm{\Delta }m`$ and $`\mathrm{arg}M_{12}`$, the new physics contribution $`M_{12}^{\mathrm{NP}}`$ is obtained. This can be used to identify the nature of the new physics contributing to the oscillation.
### 4.3 Experimental Prospects
A possible experimental programme for the study of CP violation in the B meson system and search for physics beyond the Standard Model can be summarised in the following steps:
1. Determination of $`|V_{\mathrm{cb}}|`$ and $`|V_{\mathrm{ub}}|`$ from semileptonic (and some hadronic) decays.
2. Measurement of $`\mathrm{\Delta }m`$ for $`\mathrm{B}_\mathrm{d}`$ and $`\mathrm{B}_\mathrm{s}`$
3. Measurement of $`\mathrm{}L_{\mathrm{J}/\psi \mathrm{K}_\mathrm{S}}`$
4. Measurement of $`L_{\mathrm{J}/\psi \varphi }`$, $`L_{\mathrm{D}^{}\pi ^\pm }`$ and $`L_{\mathrm{D}_\mathrm{s}^{}\mathrm{K}^\pm }`$
The first step has been made by ARGUS and CLEO at $`\mathrm{{\rm Y}}(4S)`$ machines and the four LEP experiments. BABAR and BELLE at the new asymmetric $`\mathrm{{\rm Y}}(4S)`$ machines and CLEO will improve the precisions on those determinations. Future improvement of theory is also an important factor. Half of the second step, $`\mathrm{\Delta }m(\mathrm{B}_\mathrm{d})`$ was done by ARGUS, CLEO, UA1 at the SPS Collider, the four LEP experiments, SLD at SLC and CDF at the Tevatron. For $`\mathrm{\Delta }m(\mathrm{B}_\mathrm{s})`$, we may have to wait for the next data taking by CDF, D0 and HERA-B. The third step will be made by BABAR, BELLE, CDF, D0 and possibly HERA-B by the year 2005.
After the second step, four parameters of the CKM matrix are all defined within the framework of the Standard Model, e.g. $`A`$, $`\lambda `$, $`\rho `$ and $`\eta `$. The third step provides an additional information $`\mathrm{tan}^1\eta /(1\rho )`$ within the framework of the Standard Model and consistency of the CKM picture can now be tested.
As demonstrated in the previous chapter, if physics beyond the Standard Model exists, the fourth step is needed to clearly establish the evidence of new physics and separate the effect due to the Standard Model and that from new physics. After the third step, only $`\rho ^2+\eta ^2`$ will be known from $`|V_{\mathrm{ub}}|`$ and the information on $`\mathrm{tan}^1\eta /(1\rho )`$ is spoiled by new physics. Only after the fourth step, $`\rho `$ and $`\eta `$ can be determined, together with isolating the new physics contribution.
For the last step, new generation of experiments with statistics much higher than $`10^{10}`$ B mesons are needed. The $`\mathrm{B}_\mathrm{s}`$ meson is an essential ingredient. After 2005, LHC will be the most powerful source of B mesons. Experiments must be equipped with a trigger efficient for hadronic decay modes to gain high statistics for the necessary final states. Particle identification is also crucial in order to reduce background. LHCb is a detector at the LHC optimised for CP violation studies with B mesons. The two general purpose LHC detectors, ATLAS and CMS can contribute only to a limited aspect of the fourth step. A proposed experiment at Tevatron, BTeV, can also make the last two steps.
Clearly CP violation is expected in many other decay channels. For many of them, there are some theoretical problems for making accurate predictions. However, they can be used to make a systematic study which will provide a global picture whether CP violation can fit into the CKM picture. With all those experiments, we continue to improve our understanding of CP violation and hope to discover physics beyond the Standard Model.
#### Acknowledgement
The author is very grateful to the organizers of this school for their extended hospitality and efforts to prepare such a stimulating environment. R. Forty is acknowledged for reading this manuscript and giving many useful comments. The author appreciates many useful comments by O. Schneider. |
warning/0002/hep-ph0002063.html | ar5iv | text | # MSW Effects in Vacuum Oscillations
## Abstract
We point out that for solar neutrino oscillations with the mass–squared difference of $`\mathrm{\Delta }m^210^{10}10^9`$ eV<sup>2</sup>, traditionally known as “vacuum oscillation” range, the solar matter effects are non-negligible, particularly for the low energy pp neutrinos. One consequence of this is that the values of the mixing angle $`\theta `$ and $`\pi /2\theta `$ are not equivalent, leading to the need to consider the entire physical range of the mixing angle $`0\theta \pi /2`$ when determining the allowed values of the neutrino oscillation parameters.
preprint: UCB-PTH-00/04, LBNL-45037
1. The field of solar neutrino physics is currently undergoing a remarkable change. For 30 years the goal was simply to confirm the deficit of solar neutrinos. The latest experiments, however, such as Super-Kamiokande, SNO, Borexino, KamLAND, etc, aim to accomplish more than that. By collecting high statistics real–time data sets on different components of the solar neutrino spectrum, they hope to obtain unequivocal proof of neutrino oscillations and measure the oscillation parameters. With the physics of solar neutrinos quickly becoming a precision science, it is more important then ever to ensure that all relevant physical effects are taken into account and the right parameter set is used.
It has been a long–standing tradition in solar neutrino physics to present experimental results in the $`\mathrm{\Delta }m^2\mathrm{sin}^22\theta `$ space and to treat separately the “vacuum oscillation” ($`\mathrm{\Delta }m^210^{11}10^9`$ eV<sup>2</sup>) and the MSW ($`\mathrm{\Delta }m^210^810^3`$ eV<sup>2</sup>) regions. In the vacuum oscillation region the neutrino survival probability (i.e. the probability to be detected as $`\nu _e`$) was always computed according to the canonical formula,
$$P=1\mathrm{sin}^22\theta \mathrm{sin}^2\left(1.27\frac{\mathrm{\Delta }m^2L}{E}\right),$$
(1)
where the neutrino energy $`E`$ is in GeV, the distance $`L`$ in km, and the mass–squared splitting $`\mathrm{\Delta }m^2`$ in eV<sup>2</sup>. Eq. (1) makes $`\mathrm{sin}^22\theta `$ seem like a natural parameter choice. As $`\mathrm{sin}^22\theta `$ runs from 0 to 1, the corresponding range of the mixing angle is $`0\theta \pi /4`$. There is no need to treat separately the case of $`\mathrm{\Delta }m^2<0`$ (or equivalently $`\pi /4\theta \pi /2`$), since Eq. (1) is invariant with respect to $`\mathrm{\Delta }m^2\mathrm{\Delta }m^2`$ ($`\theta \pi /2\theta `$).
The situation is different in the MSW region, since neutrino interactions with matter are manifestly flavor-dependent. It is well known that for $`|\mathrm{\Delta }m^2|10^8`$ eV<sup>2</sup> matter effects in the Sun and Earth can be quite large. In this case, if one still chooses to limit the range of the mixing angle to $`0\theta \pi /4`$, one must consider both signs of $`\mathrm{\Delta }m^2`$ to describe all physically inequivalent situations. As was argued in , to exhibit the continuity of physics around the maximal mixing, it is more natural to keep the same sign of $`\mathrm{\Delta }m^2`$ and to vary the mixing angle in the range $`0\theta \pi /2`$.
Historically, a possible argument in favor of not considering $`\theta >\pi /4`$ in the MSW region might have been that this half of the parameter space is “uninteresting”, since for $`\theta >\pi /4`$ there is no level-crossing in the Sun and the neutrino survival probability is always greater than 1/2. However, a detailed analysis reveals that allowed MSW regions can extend to maximal mixing and beyond, as was explored in (see also and for a treatment of 3- and 4- neutrino mixing schemes).
In this letter we point out that for solar neutrinos with low energies, particularly the pp neutrinos, the solar matter effects can be relevant even for neutrino oscillations with $`\mathrm{\Delta }m^210^{10}10^9`$ eV<sup>2</sup>. These effects break the symmetry between $`\theta `$ and $`\pi /2\theta `$ making it necessary to consider the full physical range of the mixing angle $`0\theta \pi /2`$ even in the “vacuum oscillation” case.
2. For simplicity, we will only consider here the two-generation mixing. If neutrino masses are nonzero then, in general, the mass eigenstates $`|\nu _{1,2}`$ are different from the flavor eigenstates $`|\nu _{e,\mu }`$. The relationship between the two bases is given in terms of the mixing angle $`\theta `$:
$`|\nu _1=\mathrm{cos}\theta |\nu _e\mathrm{sin}\theta |\nu _\mu ,`$ (2)
$`|\nu _2=\mathrm{sin}\theta |\nu _e+\mathrm{cos}\theta |\nu _\mu .`$ (3)
In our convention $`|\nu _2`$ is always the heavier of the two eigenstates, i.e. $`\mathrm{\Delta }m^2m_2^2m_1^20`$. Then, as already mentioned, $`0\theta \pi /2`$ encompasses all physically different situations.
Neutrinos are created in the Sun’s core and exit the Sun in the superposition of $`|\nu _1`$ and $`|\nu _2`$. For $`\mathrm{\Delta }m^2`$ in the vacuum oscillation region, the neutrino is produced almost completely in the heavy Hamiltonian eigenstate $`|\nu _+`$. In this case, if the evolution inside the Sun is *adiabatic*, the exit state is purely $`|\nu _2`$. In the case of a *nonadiabatic* transition there is also a nonzero probability $`P_c`$ to find the neutrino in the $`|\nu _1`$ state (a “level crossing” probability). For a given value of $`P_c`$, the survival probability for neutrinos arriving at the Earth is determined by simple 2-state quantum mechanics :
$`P`$ $`=`$ $`P_c\mathrm{cos}^2\theta +(1P_c)\mathrm{sin}^2\theta `$ (4)
$`+`$ $`2\sqrt{P_c(1P_c)}\mathrm{sin}\theta \mathrm{cos}\theta \mathrm{cos}\left(2.54{\displaystyle \frac{\mathrm{\Delta }m^2L}{E_\nu }}+\delta \right).`$ (5)
Here $`\delta `$ is a phase acquired when neutrinos traverse the Sun. In our analysis it is determined numerically . Units are the same as in Eq. (1).
In the adiabatic limit $`P_c=0`$ and Eq. (4) yields $`P=\mathrm{sin}^2\theta `$. Neutrinos exit the Sun in the heavy mass eigenstate and do not oscillate in vacuum. In the opposite limit of small $`\mathrm{\Delta }m^2`$, when the neutrino evolution in the Sun is “extremely nonadiabatic”, $`P_c\mathrm{cos}^2\theta `$. It is trivial to verify that Eq. (4) in this limit reduces to Eq. (1). It has been assumed that in the vacuum oscillation region this limit is reached. Remarkably, however, this is not always the case for the low energy solar neutrinos, especially the pp neutrinos ($`E_\nu 0.42`$ MeV).
The most reliable way to compute $`P_c`$ is by numerically solving the Schrödinger equation in the Sun for different values of $`\mathrm{\Delta }m^2`$ and $`\theta `$. We do this using the latest available BP2000 solar profile . Fig. 1 shows contours of constant $`P_c`$ for the energy of <sup>7</sup>Be neutrino (solid lines). Note that the variable on the horizontal axis is $`\mathrm{tan}^2\theta `$. With this choice, points $`\theta `$ and $`\pi /2\theta `$ are located symmetrically on the logarithmic scale about $`\mathrm{tan}^2\theta =1`$ (see ) . The figure demonstrates that the contours are not symmetric with respect to the $`\mathrm{tan}^2\theta =1`$ line, except in the region of $`\mathrm{\Delta }m^2/E_\nu 10^{10}`$ eV$`{}_{}{}^{2}/`$MeV, where the extreme nonadiabatic limit is reached. This simple observation is the crucial point of this letter.
In the MSW region the value of $`P_c`$ is often computed using the analytical result
$$P_c=\frac{e^{\gamma \mathrm{cos}^2\theta }1}{e^\gamma 1},$$
(6)
where
$$\gamma =2\pi r_0\frac{\mathrm{\Delta }m^2}{2E_\nu },$$
(7)
valid for the exponential solar profile $`n_e\mathrm{exp}(r/r_0)`$, with $`r_0=R_{}/10.54=6.60\times 10^4`$ km . Although originally derived for $`\theta \pi /4`$, it also applies when $`\theta >\pi /4`$, as was demonstrated in . In the region relevant for vacuum oscillation, however, $`0.9R_{}RR_{}`$, the profile falls off faster than the exponential with $`r_0=R_{}/10.54`$. Nevertheless, Eq. (6) can still be used with the appropriately chosen value of $`r_0`$. The dashed lines in Fig. 1 show the contours of $`P_c`$ computed using Eq. (6) with $`r_0=R_{}/18.4=3.77\times 10^4`$ km. As can be seen from the figure, the agreement between the two sets of contours for $`\mathrm{\Delta }m^24\times 10^9`$ eV<sup>2</sup> is very good. Note that a similar result was arrived at in for $`\theta \pi /4`$, where the value of $`r_0=R_{}\times 0.065=6.5\times 10^4`$ km was obtained.
Fig. 1 can also be used to read off the values of $`P_c`$ for different neutrino energies, since $`P_c`$ depends on $`E_\nu `$ through the combination $`\mathrm{\Delta }m^2/E_\nu `$. It is obvious that for neutrinos of lower energies $`P_c`$ starts deviating from its “extreme nonadiabatic” value at even smaller values of $`\mathrm{\Delta }m^2`$, and vice versa. Consequently, as will be seen later, the solar matter effects on vacuum oscillations are most important at the gallium experiments, which are sensitive to the pp neutrinos, while the Super-Kamiokande experiment is practically unaffected.
Using Eqs. (6,4), it is possible to derive a corrected form of Eq. (1), by retaining in the expansion terms linear in $`\gamma `$:
$`P`$ $`=`$ $`1\left(1+{\displaystyle \frac{\gamma }{4}}\mathrm{cos}2\theta \right)\mathrm{sin}^22\theta \mathrm{sin}^2\left(1.27{\displaystyle \frac{\mathrm{\Delta }m^2L}{E}}\right)+`$ (8)
$`+`$ $`O(\gamma ^2)`$ (9)
Notice that the first order correction contains $`\mathrm{cos}2\theta `$ and hence is manifestly not invariant under the transformation $`\theta \pi /2\theta `$. Using Eq. (7) with $`r_0=3.8\times 10^4`$ km, we see that for the pp neutrinos ($`E_\nu 0.42`$ MeV) this correction is indeed non-negligible already for $`\mathrm{\Delta }m^210^{10}10^9`$ eV<sup>2</sup>.
With matter effects being relevant already at $`\mathrm{\Delta }m^210^{10}`$ eV<sup>2</sup> one might wonder if the separation between vacuum oscillation solutions and MSW solutions is somewhat artificial. To fix the terminology, we will adopt a definition of vacuum oscillations as the situation when the value of neutrino survival probability depends on the distance $`L`$ from the Sun, regardless of whether matter effects are negligible or not. The transition between the vacuum and the MSW regions will be discussed shortly.
3. To illustrate the role of matter effects in vacuum oscillations, we present fits to the total rates of the Homestake , GALLEX and SAGE , and Super-Kamiokande experiments. We combine experimental rates and uncertainties for the two gallium experiments and use the latest available 825-day Super-Kamiokande data set. The experimental results are conveniently collected and tabulated in .
We fit the data to the theoretical predictions of the BP98 standard solar model . Predicted fluxes and uncertainties for various solar reactions were kindly made available by J. N. Bahcall at . To compute the rate suppression caused by neutrino oscillations, we numerically integrate the neutrino survival probability, Eq. (4), over the energy spectra of the pp, <sup>7</sup>Be, <sup>8</sup>B, pep, <sup>13</sup>N, and <sup>15</sup>O neutrinos. In addition, to account for the fact that the Earth–Sun distance $`L`$ varies throughout the year as a consequence of the eccentricity of the Earth’s orbit
$$L=L_0(1ϵ\mathrm{cos}(2\pi t/\mathrm{year}))$$
(10)
we also integrate over time to find an average event rate. In Eq. (10) $`t`$ is time measured in years from the perihelion, $`L_0=1.5\times 10^8`$ km is one astronomical unit, and $`ϵ=1.7\%`$.
In Fig. 2 (A) we show the vacuum oscillation regions allowed by the total rates of GALLEX and SAGE. For comparison, we also show the regions one would obtain by neglecting the neutrino interactions with the solar matter (dark outlines), i.e. by setting $`P_c=\mathrm{cos}^2\theta `$ (black contours). The allowed regions were defined as the sets of points where the theoretically predicted and experimentally observed rates are consistent with each other at the 2$`\sigma `$ C.L. for 1 d.o.f. ($`\chi ^2=4.0`$) . The plot demonstrates that the matter effects at the gallium experiments are quite important, with their contribution being significant for $`\mathrm{\Delta }m^22\times 10^{10}`$ eV<sup>2</sup>.
The remaining two plots in Fig. 2 show the vacuum regions allowed at $`3\sigma `$ C.L. by the rates of GALLEX, SAGE, and Super-Kamiokande (B) (2 d.o.f., $`\chi ^2=11.83`$, in the same convention as before), and all four experiments combined (C) (3 d.o.f., $`\chi ^2=14.15`$). In order to properly account for the correlation between the theoretical errors of the different experiments, we followed the technique developed in and . The matter effects are noticeable for $`\mathrm{\Delta }m^2>6\times 10^{10}`$ eV<sup>2</sup>.
4. An important question is how well future experiments will be able to cover vacuum oscillation solutions with $`\theta >\pi /4`$. In Fig. 3 we show the sensitivity of the Borexino experiment to anomalous seasonal variations for the entire physical range of the mixing angle $`0\theta \pi /2`$. This is an extension of the analysis performed in , where the details of the procedure are described. The sensitivity region shows a clear asymmetry as a result of the solar matter effects.
Fig. 3 gives us an opportunity to discuss the extent of the vacuum oscillation region. There are two primary physical reasons why the neutrino event rate becomes independent of $`L`$ (and anomalous seasonal variations disappear) for sufficiently large $`\mathrm{\Delta }m^2`$:
* *Adiabatic evolution in the Sun.* As $`P_c0`$ the last term in Eq. (4) vanishes.
* *Integration over neutrino energy spectrum.* To compute the event rate one has to integrate Eq. (4) over neutrino energies. For sufficiently large $`\mathrm{\Delta }m^2`$ the last term averages out to zero, leading effectively to the loss of coherence between the two mass eigenstates.
As $`\mathrm{\Delta }m^2`$ increases, coherence is first lost for reactions with broad energy spectra, such as pp and <sup>8</sup>B, and persist the longest for neutrinos produced in two-body final states. The most important such reaction is <sup>7</sup>Be$`+e^{}^7`$Li $`+\nu _e`$, which produces the <sup>7</sup>Be neutrinos. The <sup>7</sup>Be neutrinos have an energy spread of only a few keV, arising from the Doppler shift due to the motion of the <sup>7</sup>Be nucleus and the thermal kinetic energy of the electron. A detailed discussion of this phenomenon can be found in .
In order to properly take these effects into account, in our codes we numerically integrate over the exact <sup>7</sup>Be line profile, computed in . As Fig. 3 shows, the neutrino survival probability becomes independent of $`L`$ for $`\mathrm{\Delta }m^26\times 10^9`$ eV<sup>2</sup>. For this reason, we present our fits for $`\mathrm{\Delta }m^2`$ ranging from $`10^{11}`$ eV<sup>2</sup> to $`10^8`$ eV<sup>2</sup>. Unfortunately, in the literature vacuum oscillations are usually studied in the range from $`10^{11}`$ eV<sup>2</sup> to $`10^9`$ eV<sup>2</sup> , although the allowed regions in all these papers seem to extend above $`10^9`$ eV<sup>2</sup>.
5. In summary, the preceding examples clearly illustrate the importance of including the solar matter effects when studying vacuum oscillation of solar neutrinos with $`\mathrm{\Delta }m^210^{10}`$ eV<sup>2</sup>. Because to describe such effects one has to use the full range of the mixing angle $`0\theta \pi /2`$, future fits to the data should be extended to $`\theta >\pi /4`$. This seems especially important in light of the latest analyses , , which in addition to the total rates also use the information on the neutrino spectrum and time variations at Super-Kamiokande. In this case the allowed vacuum oscillation regions are mostly located in the $`\mathrm{\Delta }m^24\times 10^{10}`$ eV<sup>2</sup> region , precisely where the matter effects are relevant. (The best fit to the Super-Kamiokande electron recoil spectrum is achieved for $`\mathrm{\Delta }m^2=6.3\times 10^{10}`$ eV<sup>2</sup>, $`\mathrm{sin}^22\theta =1`$ .) It would be very desirable to repeat these analyses with the solar matter effects included.
Additionally, since the <sup>7</sup>Be neutrinos remain (partially) coherent for $`\mathrm{\Delta }m^2>10^9\text{ eV}^2`$, it is desirable to present the results of the fits in the range $`10^{11}\text{ eV}^2<\mathrm{\Delta }m^2<10^8\text{ eV}^2`$, as was done in .
###### Acknowledgements.
I am very grateful to Hitoshi Murayama and John Bahcall for their support. I would like to thank John Bahcall for including in the BP2000 solar model the data for the outer regions of the Sun. I would also like to thank James Pantaleone, Plamen Krastev, M.C. Gonzalez-Garcia, and Yosef Nir for their valuable input. This work was in part supported by the U.S. Department of Energy under Contract DE-AC03-76SF00098. |
warning/0002/nucl-th0002012.html | ar5iv | text | # Proton recoil polarization in exclusive (𝒆,𝒆'"pp") reactions
## I Introduction
It has always been a great challenge of nuclear physics to develop experiments and theoretical models able to explore short-range correlations (SRC). A detailed investigation of these correlations, which are induced by the repulsive components of the nucleon-nucleon ($`NN`$) interaction, should give insight into the structure of the interaction of two nucleons in the nuclear medium.
Since a long time electromagnetically induced two-nucleon knockout reactions have been devised as a powerful tool for such investigation, since the probability that a real or virtual photon is absorbed by a pair of nucleons should be a direct measure of the correlations between these nucleons . In particular, the ($`e,e^{}pp`$) reaction has been envisaged as the preferential tool for studying SRC in nuclei, since in such a reaction the competing contributions of two-body currents are highly suppressed. Such triple coincidence experiments have been made possible only recently by the progress in accelerator and detector technology. First measurements of the exclusive <sup>16</sup>O($`e,e^{}pp`$)<sup>14</sup>C reaction have been performed at NIKHEF in Amsterdam and MAMI in Mainz . Investigations on these data indicate that resolution of discrete final states provides an interesting tool to disentangle and thus separately investigate contributions of one-body currents, due to SRC, and two-body isobar currents. In particular, direct and clear evidence for SRC has been obtained for the transition to the ground state of <sup>14</sup>C. This result opens up good perspectives that further theoretical and experimental efforts on two-nucleon knockout reactions will be able to determine SRC.
Good opportunities to increase the richness of information available from two-nucleon knockout reactions is offered by polarization experiments. In fact, only the use of polarization degrees of freedom allows one to obtain complete information on all possible reaction matrix elements. Reactions with polarized particles depend on a larger number of observables, which are hidden in the unpolarized case, where the sum and/or average over spin states is performed. These observables are represented by new response functions , whose determination can impose more severe constraints on theoretical models. In fact, some of these observables are expected to be sensitive to the small components of the transition amplitudes, which without polarization are generally masked by the dominant ones. These small amplitudes often contain interesting information on subtle effects and may thus represent a stringent test of theoretical models. This is the place where polarization observables enter, because in general they contain interference terms of the various matrix elements in different ways. Thus, a small amplitude may be considerably amplified by the interference with a dominant one.
A complete experimental determination of the transition amplitudes requires a long-term program involving the development of new experimental techniques and many different polarization measurements able to separate various sets of structure functions. Unfortunately, some of these measurements are extremely difficult in general and in particular in the case of two-nucleon knockout. Therefore, a complete determination represents at present a too ambitious programme.
In order to exploit, at least partially, the potentiality of polarization measurements, it is possible, however, to select and investigate specific situations corresponding to simpler experiments that appear feasible in the near future. They should allow one to gain a better control on the reaction mechanism of two-nucleon knockout and hopefully study SRC.
As a first step in the study of spin degrees of freeedom in electromagnetic two-nucleon knockout reactions, we consider in this work nucleon recoil polarization. First theoretical predictions have been already presented in Refs. within the theoretical model developed by the Gent group. Here we give the most general formalism of ($`\stackrel{}{e},e^{}\stackrel{}{N}N`$) reactions in terms of structure functions and polarization observables. Numerical results are shown for the specific case of the exclusive <sup>16</sup>O($`\stackrel{}{e},e^{}\stackrel{}{p}p`$)<sup>14</sup>C reaction for transitions to the lowest-lying states of the residual nucleus. Calculations have been performed within the theoretical framework of Ref. . This model, that has been successfully applied in the analysis of the first experimental cross sections , gives a reasonably realistic base to our numerical predictions. Results for the components of the outgoing proton polarization and of the polarization transfer coefficient are presented and discussed in different kinematics. Measurements of the proton recoil polarization require an efficient proton polarimeter and a double scattering. In addition, measurements of the polarization transfer coefficient require also a polarized electron beam. Such experiments are therefore not easy, but seem reasonably within reach of available facilities.
The general formalism is given Sec. II and in Appendix. The numerical results are presented and discussed in Sec. III. Some conclusions are drawn in Sec. IV.
## II Nucleon recoil polarization in ($`\stackrel{}{𝒆},𝒆^{}\stackrel{}{𝑵}𝑵`$) reactions
In general, the polarization of an outgoing nucleon, which is the expectation value of the spin, can be calculated as
$$𝑷=\frac{\mathrm{Tr}(MM^{}𝝈)}{\mathrm{T}r(MM^{})},$$
(1)
where $`M`$ is the scattering amplitude of the reaction.
The coincidence cross section of an ($`\stackrel{}{e},e^{}\stackrel{}{N}N`$) reaction, where two nucleons, with momenta $`𝒑_1^{}`$ and $`𝒑_2^{}`$ and energies $`E_1^{}`$ and $`E_2^{}`$, are emitted and the polarization of only one nucleon, with spin directed along $`\widehat{𝒔}`$, is detected can be written
$$\frac{\mathrm{d}^5\sigma }{\mathrm{d}E_0^{}\mathrm{d}\mathrm{\Omega }_0^{}\mathrm{d}\mathrm{\Omega }_1^{}\mathrm{d}\mathrm{\Omega }_2^{}\mathrm{d}E_1^{}}=\sigma _0\frac{1}{2}\left[1+𝑷\widehat{𝒔}+h(A+𝑷^{}\widehat{𝒔})\right],$$
(2)
where $`E_0^{}`$ is the energy of the outgoing electron, $`\sigma _0`$ the unpolarized differential cross section, $`𝑷`$ the outgoing nucleon polarization, $`h`$ the electron helicity, $`A`$ the electron analyzing power and $`𝑷^{}`$ the polarization transfer coefficient.
The cross section and the polarization can be expressed in terms of the components of the hadron tensor
$$W_{\alpha \alpha ^{}}^{\mu \nu }=\overline{\underset{\mathrm{i}}{}}\underset{\mathrm{f}}{}J_\alpha ^\mu (𝒒)J_\alpha ^{}^\nu (𝒒)\delta (E_\mathrm{i}E_\mathrm{f}),$$
(3)
where $`\alpha `$ and $`\alpha ^{}`$ are the eigenvalues of the spin of the considered particle in the given reference frame and $`𝒒`$ is the momentum transfer. The quantities $`J_\alpha ^\mu (𝒒)`$ are the Fourier transforms of the transition matrix elements of the nuclear charge-current density operator between initial and final nuclear states
$$J^\mu (𝒒)=\mathrm{\Psi }_\mathrm{f}|\widehat{J}^\mu (𝒓)|\mathrm{\Psi }_\mathrm{i}\mathrm{e}^{\mathrm{i}𝒒𝒓}d𝒓.$$
(4)
The symmetrical and antisymmetrical combinations of the components of the hadron tensor produce nine spin dependent structure functions. The spin dependence can be explicited in a spherical basis (see Appendix) as
$$f_{\lambda \lambda ^{}}=h_{\lambda \lambda ^{}}^\mathrm{u}+\widehat{𝒔}𝒉_{\lambda \lambda ^{}},$$
(5)
where $`\widehat{𝒔}`$ is the unit vector in the spin space. The structure functions represent the response of the nucleus to the longitudinal and transverse components of the electromagnetic interaction and only depend on the energy and momentum transfer $`\omega `$ and $`q`$, the momenta of the two outgoing nucleons $`p_1^{}`$ and $`p_2^{}`$ and the angles $`\gamma _1`$, between $`𝒑_1^{}`$ and $`𝒒`$, $`\gamma _2`$, between $`𝒑_2^{}`$ and $`𝒒`$, and $`\gamma _{12}`$, between $`𝒑_1^{}`$ and $`𝒑_2^{}`$. When the outgoing nucleon polarization is not detected and the cross section is summed over the spin quantum numbers of the outgoing nucleon, the spin independent structure functions $`h_{\lambda \lambda ^{}}^\mathrm{u}`$ go over to the structure functions $`f_{\lambda \lambda ^{}}`$ of the unpolarized case . Thus, in this case only nine structure functions are obtained. New structure functions are produced in the polarized case by the components of $`𝒉_{\lambda \lambda ^{}}`$
The components of the polarization and of the polarization transfer coefficient are obtained through the structure functions and, involving also the matrix elements of the lepton tensor $`\rho _{\lambda \lambda ^{}}`$ and $`\rho _{\lambda \lambda ^{}}^{}`$ , are given by
$$P^i=\frac{_{\lambda \lambda ^{}}\rho _{\lambda \lambda ^{}}h_{\lambda \lambda ^{}}^i}{_{\lambda \lambda ^{}}\rho _{\lambda \lambda ^{}}h_{\lambda \lambda ^{}}^\mathrm{u}}$$
(6)
and
$$P^i=\frac{_{\lambda \lambda ^{}}\rho _{\lambda \lambda ^{}}^{}h_{\lambda \lambda ^{}}^i}{_{\lambda \lambda ^{}}\rho _{\lambda \lambda ^{}}h_{\lambda \lambda ^{}}^\mathrm{u}}.$$
(7)
The vectors $`𝑷`$ and $`𝑷^{}`$ are usually projected onto the basis of unit vectors given by $`\widehat{𝑳}`$ (parallel to the momentum $`𝒑^{}`$ of the outgoing particle), $`\widehat{𝑵}`$ (in the direction of $`𝒒\times 𝒑^{}`$) and $`\widehat{𝑺}=\widehat{𝑵}\times \widehat{𝑳}`$, which define the cm helicity frame of the particle.
The explicit expressions of the cross section in terms of the structure functions and of the structure functions in terms of the components of the hadron tensor can be found in Appendix.
For the ($`\stackrel{}{e},e^{}\stackrel{}{N}N`$) reaction in an unrestricted kinematics all the components $`P^N,P^L,P^S`$ and $`P^N,P^L,P^S`$ are allowed, as parity conservation does not impose in this case any restrictions. In this general situation 36 structure functions are present: nine spin independent structure functions in $`\sigma _0`$ ($`h^\mathrm{u}`$) and $`A`$ ($`h^\mathrm{u}`$), and nine spin dependent structure functions in each one of the components $`i=N,L,S`$ of $`𝑷`$ ($`h^i`$) and $`𝑷^{}`$ ($`h^i`$).
This number is reduced in particular situations (see Appendix for more details). When the angle $`\alpha `$ between the $`𝒑_1^{}`$ $`𝒒`$ plane and the electron scattering plane is equal to zero, all the 36 structure functions are in general non-vanishing, but those of them which are multiplied by $`\mathrm{sin}\alpha `$ or $`\mathrm{sin}2\alpha `$ in Eq. (16) do not contribute to the cross section and to the components of $`𝑷`$ and $`𝑷^{}`$. As a consequence, the condition $`\alpha =0`$ reduces to 24 the number of structure functions.
When the vectors $`𝒒`$, $`𝒑_1^{}`$, $`𝒑_2^{}`$ lie in the same plane, parity conservation combined with the general properties of the hadron tensor reduces the number of non-vanishing structure functions to 18: five in $`\sigma _0`$ and $`A`$ and 13 in the components of $`𝑷`$ and $`𝑷^{}`$. This result is similar to that obtained for the ($`\stackrel{}{e},e^{}\stackrel{}{N}`$) reaction in an unrestricted kinematics .
For a coplanar kinematics, i.e. when the initial and final electrons and the outgoing nucleons lie in the same plane, those of the 18 non-vanishing structure functions which are multiplied by $`\mathrm{sin}\alpha `$ or $`\mathrm{sin}2\alpha `$ do not contribute and the number is reduced to 12, 4 spin independent and 8 spin dependent, and only the components $`P^N,P^L`$ and $`P^S`$ survive. This result is similar to that obtained for the ($`\stackrel{}{e},e^{}\stackrel{}{N}`$) reaction in a coplanar kinematics .
A particular situation occurs in the interesting case of the super-parallel kinematics, where the two outgoing nucleons are ejected parallel and antiparallel to the momentum transfer. In this case only two structure functions, $`h_{00}^\mathrm{u}`$ and $`h_{11}^\mathrm{u}`$, do not vanish in the unpolarized cross section and three, $`h_{01}^N,\overline{h}_{01}^S,h_{11}^L`$, when polarization is considered. As a consequence, $`P^N,P^L`$ and $`P^S`$ are each directly proportional to only one structure function. This result is similar to the one obtained for the ($`\stackrel{}{e},e^{}\stackrel{}{N}`$) reaction in parallel kinematics .
It is worthwhile to investigate what happens when final-state interactions (FSI) are neglected and the plane-wave (PW) approximation is used for the outgoing nucleons wave functions.
The behaviour of the hadron tensor under time reversal and parity transformation has the property
$$W^{\mu \nu }(𝒔,())=W^{\nu \mu }(𝒔,(+)),$$
(8)
where $`𝒔`$ is the spin vector in the ejectile rest frame, and the dependence on the final state boundary condition for incoming $`()`$ and outgoing $`(+)`$ scattered waves is shown. For nucleon knockout, the $`()`$ condition is appropriate. When the boundary conditions can be ignored, as in the PW approximation, Eq. (8) states that the symmetric part of $`W^{\mu \nu }`$ is independent of $`𝒔`$ and the antisymmetric part is proportional to $`𝒔`$. This is because, owing to the spin-$`\frac{1}{2}`$ nature of the nucleon, the dependence of $`W^{\mu \nu }`$ on $`𝒔`$ is at most linear. Therefore, in the PW approximation $`P^N=P^L=P^S=0`$, while $`𝑷^{}`$ does not vanish.
## III Proton recoil polarization in the <sup>16</sup>O($`\stackrel{}{𝒆},𝒆^{}\stackrel{}{𝒑}𝒑`$)<sup>14</sup>C reaction
An experimental separation of the various structure functions would be of great interest, but appears extremely difficult. A measurement of the nucleon recoil polarization would be simpler and less affected by experimental errors, as it is obtained through the determination of asymmetries. Therefore, in this section we present numerical predictions for the outgoing proton polarization $`𝑷`$ and the polarization transfer coefficient $`𝑷^{}`$ of the <sup>16</sup>O($`\stackrel{}{e},e^{}\stackrel{}{p}p`$)<sup>14</sup>C reaction leading to the lowest-lying discrete states in the residual nucleus. This reaction is of particular interest for our investigation, due to the presence of discrete states in the experimental spectrum of the residual nucleus <sup>14</sup>C, well separated in energy and that can be separated with high-resolution experiments. Cross sections calculations pointed out that transitions to different states can be differently affected by the two reaction processes due to SRC and two-body currents . Thus, the experimental separation of different final states can act as a filter for the study of the two processes. Recent experiments at NIKHEF and MAMI have been able to resolve the lowest-lying states of <sup>14</sup>C and have confirmed, in comparison with the theoretical results, the predicted selectivity of the exclusive reaction involving different transitions. In particular, clear evidence of the dominant contribution of SRC has been obtained for the transition to the 0<sup>+</sup> ground state, while the transition to the 1<sup>+</sup> state at 11.31 MeV appears better dominated by the $`\mathrm{\Delta }`$ isobar current.
New and complementary information is in principle available from polarization observables. The aim of our investigation is to clarify the sensitivity of $`𝑷`$ and $`𝑷^{}`$ to the two competing reaction processes and in particular to SRC.
Calculations have been performed within the same theoretical model used for the analysis of the available cross section data. A detailed description of the theoretical framework can be found in Ref. and in a series of previous papers where the different aspects of the model have been developed . Here we summarize only the main features.
The transition matrix elements $`J^\mu (𝒒)`$ in Eq. (4), for an exclusive reaction and under the assumption of a direct knockout mechanism, can be reduced to a form which contains three main ingredients: the two-nucleon overlap integral, the nuclear current and the final-state wave function of the two outgoing nucleons.
In the calculations the scattering state is given by the product of two uncoupled single-particle distorted wave functions, eigenfunctions of a complex phenomenological optical potential which contains a central, a Coulomb and a spin-orbit term.
The nuclear current operator is the sum of a one-body and a two-body part. These two parts correspond to the two reaction processes. In fact, while two nucleons are naturally ejected by a two-body current, even if correlations are not explicitly included in the two-nucleon wave function, they cannot be ejected by a one-body current without correlations. Thus, the contribution of the one-body current is entirely due to correlations. In the calculations the one-body part contains a Coulomb, a convective and a spin term. For $`pp`$ knockout the two-body current, which is completely transverse, contains only the contributions of non charge-exchange processes with intermediate $`\mathrm{\Delta }`$-isobar configurations in the intermediate state .
The two-nucleon overlap integrals are taken, for the different final states, from the calculation of the two-proton spectral function of <sup>16</sup>, where long-range and short-range correlations are consistently included. They are expressed in terms of a sum of products of relative and cm wave functions. Different components of relative and cm motion contribute to different transitions. They are weighed in the sum by two-proton removal amplitudes calculated within a large shell-model basis. SRC are included in the radial wave functions of relative motion through defect functions, defined by the difference between correlated and uncorrelated relative wave functions, and which are different for different relative states and for different $`NN`$ potentials.
Therefore, SRC play a different role in different relative states. They are quite strong for the $`{}_{}{}^{1}S_{0}^{}`$ state and much weaker for $`{}_{}{}^{3}P_{j}^{}`$ states (the notation $`{}_{}{}^{2S+1}l_{j}^{}`$, for $`l=S,P\mathrm{}`$, is here used for the relative states). An opposite effect is given by the two-body $`\mathrm{\Delta }`$ current, whose contribution is strongly reduced for $`{}_{}{}^{1}S_{0}^{}`$ $`pp`$ knockout, since there the generally dominant contribution of that current, due to the magnetic dipole $`NNN\mathrm{\Delta }`$ transition, is suppressed . Thus, $`{}_{}{}^{1}S_{0}^{}`$ $`pp`$ knockout is generally dominated by SRC, while the $`\mathrm{\Delta }`$ current is more important in $`{}_{}{}^{3}P_{j}^{}`$ $`pp`$ knockout.
This result explains the above mentioned selectivity of the exclusive cross sections involving the transitions to the 0<sup>+</sup> ground state and the 1<sup>+</sup> state. In fact, only $`{}_{}{}^{3}P`$ relative states, $`{}_{}{}^{3}P_{0}^{}`$ $`{}_{}{}^{3}P_{1}^{}`$ $`{}_{}{}^{3}P_{2}^{}`$, contribute for the 1<sup>+</sup> state, whose cross section is generally dominated by the $`\mathrm{\Delta }`$ current. The two relative waves $`{}_{}{}^{1}S_{0}^{}`$ and $`{}_{}{}^{3}P_{1}^{}`$ contribute for the 0<sup>+</sup> ground state and it is possible to envisage suitable kinematics where the role of $`{}_{}{}^{1}S_{0}^{}`$ $`pp`$ knockout and thus of SRC becomes dominant.
In our investigation of polarization observables in the <sup>16</sup>O($`\stackrel{}{e},e^{}\stackrel{}{p}p`$)<sup>14</sup>C reaction calculations have been performed in different kinematics. In order to obtain a more complete information it is interesting to consider both coplanar and out-of-plane kinematics. Experiments in coplanar kinematics are certainly simpler, but give access only to the components $`P^N`$, $`P^L`$ and $`P^S`$, while all the components of $`𝑷`$ and $`𝑷^{}`$ are present in an out-of-plane kinematics.
A special and interesting case of coplanar kinematics is represented by the so-called super-parallel kinematics , where the knocked-out nucleons are detected parallel and anti-parallel to the transferred momentum $`𝒒`$. This kinematics, that has been realized in the experiments at MAMI , is favored by the fact that only two structure functions, $`f_{00}`$ ($`h_{00}^\mathrm{u}`$) and $`f_{11}`$ ($`h_{11}^\mathrm{u}`$), contribute to the unpolarized cross section $`\sigma _0`$ and $`P^N`$, $`P^L`$, $`P^S`$ are each directly proportional to one structure function: $`h_{01}^N`$, $`h_{11}^L`$, $`\overline{h}_{01}^S`$, respectively.
The unpolarized differential cross section as well as the two unpolarized structure functions in the super-parallel kinematics of the MAMI experiment have been already shown and discussed in Ref. . The polarization observables $`P^N`$, $`P^L`$, $`P^S`$ are displayed in Figs. 1 and 2 for the transitions to the $`0^+`$ ground state and the $`1^+`$ state of <sup>14</sup>C, respectively, as a function of the recoil ($`p_\mathrm{B}`$) or missing momentum ($`p_{2\mathrm{m}}`$), defined by
$$𝒑_{2\mathrm{m}}=𝒑_\mathrm{B}=𝒒𝒑_1^{}𝒑_2^{}.$$
(9)
In Figs. 1 and 2 the result given by the sum of the one-body and the two-body current is compared with the separate contribution of the one-body current. In Fig. 1, for the transition to the ground state, where effects of SRC are expected to be more relevant, results given by the defect functions for the Bonn-A and Reid $`NN`$ potentials are compared.
In the calculations each state is characterized by a particular value of the missing energy, given by
$$E_{2\mathrm{m}}=\omega T_1^{}T_2^{}T_\mathrm{B}=E_\mathrm{s}+E_\mathrm{x},$$
(10)
where $`T_1^{}`$, $`T_2^{}`$ and $`T_\mathrm{B}`$ are the kinetic energies of the two outgoing nucleons and of the residual nucleus, respectively, $`E_\mathrm{s}`$ is the separation energy at threshold for two-nucleon emission and $`E_\mathrm{x}`$ is the excitation energy of the residual nucleus.
In the super-parallel kinematics all possible values of $`p_\mathrm{B}`$ are explored, for a fixed value of the energy and momentum transfer and for a particular final state, changing the values of the kinetic energies of the outgoing nucleons.
An analysis of the results as a function of the recoil momentum appears of particular interest. The shape of the unpolarized differential cross section is determined by the value of the cm orbital angular momentum $`L`$ of the knocked out pair . Different components of relative and cm motion contribute to the two-nucleon overlap functions for each final state. The shape of the calculated cross section is therefore driven by the component which gives the major contribution. For the $`0^+`$ state, the $`{}_{}{}^{1}S_{0}^{}`$ and $`{}_{}{}^{3}P_{1}^{}`$ relative waves are combined with $`L=0`$ and $`L=1`$, respectively; for the $`1^+`$ state the $`{}_{}{}^{3}P_{0}^{}`$, $`{}_{}{}^{3}P_{1}^{}`$ and $`{}_{}{}^{3}P_{2}^{}`$ relative waves are all combined with $`L=1`$ . Cross sections calculations in the considered kinematics indicate that the transition to the $`0^+`$ ground state has an $`s`$ wave shape, which is due to the major role played by $`{}_{}{}^{1}S_{0}^{}`$ $`pp`$ knockout, dominated by SRC, for low values of the recoil momentum. The $`p`$ wave component, which is due to the $`{}_{}{}^{3}P_{1}^{}`$ relative state and is dominated by the $`\mathrm{\Delta }`$ current, becomes meaningful only at large values of $`p_\mathrm{B}`$, beyond $`150200`$ MeV/$`c`$, where the contribution of the $`s`$ wave becomes much smaller . For the transition to the $`1^+`$ state the cross section has a $`p`$ wave shape and is almost entirely due to the $`\mathrm{\Delta }`$ current . These theoretical findings have been confirmed in comparison with data .
The polarization observables displayed in Fig. 1 for the $`0^+`$ state confirm the dominant role of the one-body current up to about $`150200`$ MeV/$`c`$. For larger values of $`p_\mathrm{B}`$, where the cross section is smaller, the contribution of the $`\mathrm{\Delta }`$ current becomes more relevant. This is due to its interference with the one-body current, which is meaningful in the interference longitudinal-transverse structure functions $`h_{01}^N`$ and $`\overline{h}_{01}^S`$ in this region of recoil momenta. For the component $`P^L`$, which is proportional to the transverse polarized structure function $`h_{11}^L`$, the separate contribution of the $`\mathrm{\Delta }`$ current is large and of about the same size as that of the one-body current for all the considered values of $`p_\mathrm{B}`$. Thus, in this case, effects of interference between the one-body current and the two-body current are important over all the momentum distribution.
The results given by the two sets of defect functions from Bonn-A and Reid potentials on the components of $`P^N`$, $`P^L`$ and $`P^S`$ are qualitatively similar. The numerical differences, however, are appreciable and even large at large values of $`p_\mathrm{B}`$, especially for the component $`P^N`$. This effect is predominantly due to the different interference between the one-body current and the $`\mathrm{\Delta }`$ current.
The results in Fig. 1 for the transitions to the $`0^+`$ state indicate that in the super-parallel kinematics the considered polarization observables are sizable and sensitive to effects of SRC. The component of the proton recoil polarization $`P^N`$ appears very well suited to study correlation effects. A measurement of this component would be simpler than a measurement of $`P^L`$ and $`P^S`$, which would require also a polarized electron beam. On the other hand, since all the components of $`𝑷`$ vanish in the PW approximation, $`P^N`$ could be also sensitive to the treatment of FSI.
The polarization observables displayed in Fig. 2 for the transition to the $`1^+`$ state are smaller and, as expected, more sensitive to the contribution of the $`\mathrm{\Delta }`$ current. The results in Figs. 1 and 2 confirm that different final states in the exclusive <sup>16</sup>O($`e,e^{}pp`$)<sup>14</sup>C reaction may act as a filter for the study of the two reaction processes due to SRC and two-body currents.
The super-parallel kinematics represents a special case of coplanar kinematics where only 5 structure functions contribute to the unpolarized and polarized cross sections. In general, in a coplanar kinematics 12 structure functions are present (see Sec. II and Appendix): 4 are contained in the unpolarized cross section $`\sigma _0`$ and 8 in the surviving components of the of the proton recoil polarization $`P^N`$ and of the polarization transfer coefficient $`P^L`$ and $`P^S`$.
As an example, we have here considered a specific kinematical setting included in the first experiment on <sup>16</sup>O carried out at NIKHEF , with an incident electron energy of 584 MeV, $`\omega =212`$ MeV and $`q=300`$ MeV/$`c`$. The kinetic energy of the first outgoing proton $`T_1^{}`$ is 137 MeV and the angle $`\gamma _1=30^\mathrm{o}`$, on the opposite side of the outgoing electron with respect to the momentum transfer. Changing the angle $`\gamma _2`$ on the other side, different values of the recoil momentum $`𝒑_\mathrm{B}`$ are explored in the range between $`250`$ and $`300`$ MeV/$`c`$, including the zero values at $`\gamma _2120^\mathrm{o}`$ .
The unpolarized cross section has been already discussed in Ref. and is shown again, for the transition to the $`0^+`$ ground state of <sup>14</sup>C, in the top panels of Fig. 4. The qualitative features are similar to those obtained in the super-parallel kinematics of the MAMI experiment. For the $`0^+`$ state the shape of the recoil momentum or angular distribution is driven by the component with $`L=0`$ , that is by $`{}_{}{}^{1}S_{0}^{}`$ $`pp`$ knockout, and is thus dominated by SRC at low values of $`p_\mathrm{B}`$. The contribution with $`L=1`$, combined with the $`{}_{}{}^{3}P_{1}^{}`$ relative state, which is better driven by the $`\mathrm{\Delta }`$ current, is negligible when $`p_\mathrm{B}`$ is low, but becomes meaningful at large values of the recoil momentum, where the contribution of the $`s`$ wave component is strongly suppressed. In contrast, the transition to the $`1^+`$ state, which contains only $`{}_{}{}^{3}P`$ relative waves and $`L=1`$ cm components, has a $`p`$ wave shape and is dominated by the $`\mathrm{\Delta }`$ current.
The polarization observables $`P^N`$, $`P^L`$, and $`P^S`$ are displayed, as a function of the angle $`\gamma _2`$, in Fig. 3 for the transitions to the $`0^+`$ ground state and to the $`1^+`$ state of <sup>14</sup>C. For the $`0^+`$ state $`P^N`$ is large and dominated over all the angular distribution by the one-body current and thus by SRC. The components of the polarization transfer coefficient $`P^L`$ and $`P^S`$ are also large, but appear better driven by the $`\mathrm{\Delta }`$ current. Thus, also in this kinematics the component $`P^N`$ turns out to be very well suited to study SRC, while in $`P^L`$ and $`P^S`$ the contribution of the $`\mathrm{\Delta }`$ current is relevant and intertwined with that of the one-body current. Therefore, for these two components, whose measurement requires more complicated experiments, the separation of either contribution of the two reaction processes appears difficult.
The comparison between the results with the two kinematical settings in Figs. 1 and 3 for the $`0^+`$ state indicates that in Fig. 3, where a larger number of structure functions contribute, the polarization observables are generally larger and the role of SRC in $`P^N`$ is dominant over all the distribution. From this point of view, this kinematics seems better suited for the study of SRC with polarization measurements. On the other hand, the role of the $`\mathrm{\Delta }`$ current in $`P^L`$ and $`P^S`$ is more relevant than in the super-parallel kinematics of Fig. 1.
Also in the kinematics of Fig. 3 the polarization observables for the transition to the $`1^+`$ state turn out to be generally smaller than for the $`0^+`$ state. For this transition $`P^N`$, $`P^L`$ and $`P^S`$ are sizable only in the region of low vales of $`p_\mathrm{B}`$, where the cross section has the minimum, and are driven over all the distribution by the $`\mathrm{\Delta }`$ current. This result confirms that the transition to the $`1^+`$ state is dominated by two-body currents and indicates that polarization measurements for this final state appear more difficult.
The results of Fig. 3 have been obtained with the defect functions from the Bonn-A potential. Calculations performed with the set of defect functions from the Reid potential give appreciable differences, especially for the $`0^+`$ state and for the observable $`P^N`$, which is most affected by correlations. The main qualitative features of the results remain however unchanged with respect to those presented in Fig. 3.
All the components $`N,L,S`$ of $`𝑷`$ and $`𝑷^{}`$ can be explored with an out-of-plane kinematics. Such kinematics can be realized when the angle $`\alpha `$, between the $`𝒑_1^{}`$ $`𝒒`$ plane and the electron scattering plane, or the azimuthal angle $`\varphi `$ of the outogoing proton whose polarization is not considered, or even both $`\alpha `$ and $`\varphi `$ are different from zero. In the most general case, where $`\alpha 0`$ and $`\varphi 0`$, 36 structure functions are obtained (see Sec. II and Appendix). When $`\varphi 0`$, all these structure functions are in general non-vanishing, but if $`\alpha =0`$ those of them which are multiplied by $`\mathrm{sin}\alpha `$ or $`\mathrm{sin}2\alpha `$ do not give any contributions. Thus only 24 structure functions are active. If $`\varphi =0`$, the vectors $`𝒒`$, $`𝒑_1^{}`$ and $`𝒑_2^{}`$ lie in the same plane and the structure functions are the same as in coplanar kinematics, but if $`\alpha 0`$ also those of them which are multiplied by $`\mathrm{sin}\alpha `$ or $`\mathrm{sin}2\alpha `$ contribute. This gives 18 structure functions.
In order to explore all the components of $`𝑷`$ and $`𝑷^{}`$, two examples of out-of-plane kinematics are here considered: $`\alpha =45^\mathrm{o}`$ $`\varphi =0^\mathrm{o}`$ and $`\alpha =0^\mathrm{o}`$ $`\varphi =30^\mathrm{o}`$. The other kinematical variables are taken as in the coplanar kinematics of Fig. 3. Calculations have been performed only for the transition to the $`0^+`$ ground state, where effects of SRC are known to be most relevant.
The differential cross sections in the two considered out-of-plane kinematics are displayed in Fig. 4 as a function of the angle $`\gamma _2`$ and compared with the cross section of the corresponding coplanar kinematics. The final results are compared in the three cases with the separate contributions of the one-body and the two-body currents and of the $`{}_{}{}^{1}S_{0}^{}`$ and $`{}_{}{}^{3}P_{1}^{}`$ components of relative motion.
In the kinematics with $`\alpha =45^\mathrm{o}`$ the size of the cross section is only slightly lower than in the coplanar kinematics. Also the shapes of the curves given by the separate contributions are similar to those obtained in the coplanar situation. The contribution of the $`{}_{}{}^{1}S_{0}^{}`$ component is, however, reduced and that of $`{}_{}{}^{3}P_{1}^{}`$ enhanced. As a consequence, the role of $`{}_{}{}^{3}P_{1}^{}`$ $`pp`$ knockout, dominated by the $`\mathrm{\Delta }`$ current, is enhanced in the final cross section. This effect slightly changes the shape of the final distribution. On the other hand, when $`\alpha =45^\mathrm{o}`$ the role of the $`\mathrm{\Delta }`$ current becomes much more important also on the $`{}_{}{}^{1}S_{0}^{}`$ component, as can be seen from the comparison between the results in the left and right panels. Therefore, in the final cross section the contribution of the $`\mathrm{\Delta }`$ current becomes competitive with that of the one-body current even for the transition to the $`0^+`$ state. This is the main difference with respect to the result of the coplanar kinematics and is due to the different role played by the structure functions which are multiplied by factors including the angle $`\alpha `$. In particular, a significant role is played by the structure functions multiplied by $`\mathrm{sin}\alpha `$ and $`\mathrm{sin}2\alpha `$, which do not contribute in the coplanar kinematics, and where the effect of the $`\mathrm{\Delta }`$ current is important.
In the kinematics with $`\varphi =30^\mathrm{o}`$ the cross section is about an order of magnitude lower than in the peak region of the coplanar kinematics and has a different shape. In practice, the peak has disappeared and there is no more evidence of an $`s`$ shape distribution. In fact, the two contributions of $`{}_{}{}^{1}S_{0}^{}`$ and $`{}_{}{}^{3}P_{1}^{}`$ turn out to be of comparable size in this case. In spite of that, the major role is still played by the one-body current, whose contribution is relevant, even though not dominant, also on the $`{}_{}{}^{3}P_{1}^{}`$ relative wave. A part of the differences with respect to the results of the coplanar kinematics is due to the different structure functions and to their dependence on the angle $`\varphi `$. The main reason of the differences, however, can be attributed to a kinematic effect. In fact, when the outgoing nucleon is taken out of the plane, different values of the recoil momentum $`p_\mathrm{B}`$ are obtained in the considered range of $`\gamma _2`$. In particular, low values in the range between $`160`$ and $`160`$ MeV/$`c`$ are forbidden. Thus, the main source of difference with respect to the coplanar situation is that the region of momenta between $`160`$ and $`160`$ MeV/$`c`$, where the $`s`$ wave has the maximum and the $`p`$ wave a minimum, has been cut. This explains the different shape of the angular distributions in the top and bottom panels of Fig. 4.
The components of $`𝑷`$ and $`𝑷^{}`$ in the two kinematics with $`\alpha =45^\mathrm{o}`$ and $`\varphi =30^\mathrm{o}`$ are displayed in Figs. 5 and 6, respectively.
In Fig. 5, with $`\alpha =45^\mathrm{o}`$, the components $`P^N`$, $`P^L`$ and $`P^S`$, already present in coplanar kinematics, are significantly different from those displayed in Fig. 3 for the same transition. The differences are larger for $`P^N`$ and $`P^L`$. The $`\mathrm{\Delta }`$ current gives the major contribution to $`P^L`$ and $`P^S`$. The one-body current gives the main contribution to $`P^N`$, but in this kinematics the effect of the $`\mathrm{\Delta }`$ current is meaningful also on this polarization component and larger than in coplanar kinematics. Of the three observables $`P^N`$, $`P^L`$ and $`P^S`$, which are present only in an out-of-plane kinematics, $`P^N`$ is small, while $`P^L`$ and $`P^S`$ are sizable. However, both contributions of SRC and two-body currents are large and intertwined in these polarization components. Therefore, the kinematics with $`\alpha =45^\mathrm{o}`$ does not seem particularly well suited to disentangle the two reaction processes and study SRC.
Correlation effects are much more important in Fig. 6, for the kinematics with $`\varphi =30^\mathrm{o}`$. In this case all the components of the outgoing proton polarization $`𝑷`$ are sizable and driven by the one-body current, which gives the main contribution also to $`P^N`$. This component of the polarization transfer coefficient, however, is small also in this kinematics. The other components $`P^L`$ and $`P^S`$ are large, but also in this case strongly affected by both contributions of the two reaction processes due to SRC and two-body currents.
## IV Summary and conclusions
In this paper we have discussed the general properties of the nucleon recoil polarization in the ($`\stackrel{}{e},e^{}\stackrel{}{N}N`$) reaction. In the most general situation no restrictions are due to parity conservation and 36 structure functions are active, more than in the ($`\stackrel{}{e},e^{}\stackrel{}{N}`$) reaction, where parity conservation reduces the number of structure functions available in an unconstrained kinematics to 18. The same formal situation is obtained in ($`\stackrel{}{e},e^{}\stackrel{}{N}N`$) only when the momenta of the two outgoing nucleons and the momentum transfer lie in the same plane. A minor number of structure functions can be obtained in more restricted kinematics: 12 in a coplanar kinematics, as in the coplanar kinematics of ($`\stackrel{}{e},e^{}\stackrel{}{N}`$), and 5 in the special case of the super-parallel kinematics, as in the parallel kinematics of ($`\stackrel{}{e},e^{}\stackrel{}{N}`$).
An experimental determination of the structure functions would be of great interest, but appears at present extremely difficult. In order to exploit the potentiality of polarization measurements to increase the richness of information available in two-nucleon knockout reactions, it is anyhow possible to envisage other interesting polarization experiments which appear more feasible. As an example, a measurement of the nucleon recoil polarization, which is obtained through the determination of asymmetries, should be reasonably within reach of available facilities.
A complete investigation of all the components of the outgoing nucleon polarization $`𝑷`$ and of the polarization transfer coefficient $`𝑷^{}`$, which implies also a polarized electron beam, requires out-of-plane experiments. In the simpler case of a coplanar kinematics only the components $`P^N`$, $`P^L`$ and $`P^S`$ are available, as in the coplanar kinematics of the ($`\stackrel{}{e},e^{}\stackrel{}{N}`$) reaction. Moreover, as in ($`\stackrel{}{e},e^{}\stackrel{}{N}`$), all the components of the outgoing nucleon polarization $`𝑷`$ vanish in the PW approximation.
The aim of this work was to explore the capability of polarization measurements to disentangle and separately investigate the two reaction processes due to correlations and two-body currents. With this aim and within our theoretical framework, we have checked the sensitivity of $`𝑷`$ and $`𝑷^{}`$ to the two competing processes.
Numerical predictions have been presented for the specific case of the exclusive <sup>16</sup>O($`\stackrel{}{e},e^{}\stackrel{}{p}p`$)<sup>14</sup>C reaction leading to the lowest-lying discrete states of the residual nucleus. Our results confirm that the contribution of SRC as compared with that of the $`\mathrm{\Delta }`$ current depends on the final state of the residual nucleus and that the transition to the ground state of <sup>14</sup>C is particularly sensitive to correlation effects. Calculations performed in the coplanar kinematical settings realized for the cross section measurements at NIKHEF and MAMI indicate that for the transition to the ground state the major contribution to the only surviving component of the outgoing proton polarization, perpendicular to the plane, is given by correlations, while effects of two-body currents are almost negligible. In contrast, these effects are much more important or even dominant on the components of the polarization transfer coefficient.
In more complicated non coplanar kinematics the number of available observables increases, but their sensitivity to nuclear correlations is reduced. It is possible to envisage particular conditions where correlation effects are dominant, but they correspond to situations where measurements are more difficult and cross sections generally smaller.
In conclusion, we believe that a combined measurement of cross sections and polarizations in ($`\stackrel{}{e},e^{}\stackrel{}{p}p`$) reactions would provide a unique tool to determine short-range nuclear correlations and could positively contribute to clarify the behaviour of the short-range interaction of nucleons in the nuclear medium.
The analysis of polarization observables can be extended to other electromagnetic two-nucleon knockout reactions, such as ($`e,e^{}pn`$), ($`\gamma ,pn`$) and ($`\gamma ,pp`$).
Recent calculations have shown that the cross section of the exclusive <sup>16</sup>O($`\stackrel{}{e},e^{}\stackrel{}{p}n`$)<sup>14</sup>N reaction are sensitive to details of the nuclear correlations considered and in particular to the presence of the tensor component. Therefore, a study of polarization observables in the ($`e,e^{}pn`$) reaction would be of particular interest for the study of tensor correlations.
Cross sections and photon asymmetries calculated for both ($`\gamma ,pp`$) and ($`\gamma ,pn`$) knockout reactions are dominated by two-body currents and only slightly affected by correlation effects . Thus, reactions induced by real photons do not seem particularly well suited to study correlations, but are anyhow interesting to give complementary information on other theoretical ingredients, for instance on the nuclear currents. Moreover, polarization observables could amplify the role of small amplitudes and of subtle effects hidden in the unpolarized case. Therefore, also ($`\gamma ,pp`$) and ($`\gamma ,pn`$) reactions deserve a careful investigation.
## Appendix
The components of the hadron tensor $`W^{\mu \nu }`$ are to be restricted by the conditions of current conservation
$$W^{\mu \nu }q_\mu =W^{\mu \nu }q_\nu =0.$$
(11)
Therefore, the sixteen components of the tensor are reduced to nine independent quantities. By separating the symmetrical and antisymmetrical terms, one obtains six symmetrical and three antysimmetrical independent components.
A spherical basis can be used and defined by the four-vectors
$$ϵ_{\pm 1}^\mu =\frac{1}{\sqrt{2}}(0,1,\pm \mathrm{i},0),ϵ_0^\mu =(\frac{|𝒒|}{Q},0,0,\frac{\omega }{Q}),$$
(12)
where $`Q^2=|𝒒|^2\omega ^2`$. (Note that this spherical basis is consistent with Ref. , and is different from the one used in other papers, where $`ϵ_0^\mu =(1,0,0,0)`$.)
Then, nine structure functions are obtained as a function of the hadron tensor components . Their expressions are given in Table 1, in the reference frame where the $`𝒛`$ axis is taken parallel to $`𝒒`$ and the momentum $`𝒑_1^{}`$ lies in the $`𝒙`$$`𝒛`$ plane.
The triple coincidence cross section for the electron induced reaction where two nucleons are emitted is obtained from the contraction between the lepton tensor and the hadron tensor as a linear combination of the nine structure functions
$`{\displaystyle \frac{\mathrm{d}^5\sigma }{\mathrm{d}E_0^{}\mathrm{d}\mathrm{\Omega }_0^{}\mathrm{d}\mathrm{\Omega }_1^{}\mathrm{d}\mathrm{\Omega }_2^{}\mathrm{d}E_1^{}}}=\sigma _\mathrm{M}`$ $`\mathrm{\Omega }_\mathrm{f}`$ $`f_{\mathrm{rec}}{\displaystyle \underset{\lambda \lambda ^{}}{}}\{\rho _{\lambda \lambda ^{}}f_{\lambda \lambda ^{}}+h\rho _{\lambda \lambda ^{}}^{}f_{\lambda \lambda ^{}}^{}\}`$ (13)
$`=K`$ $`\mathrm{\Omega }_\mathrm{f}`$ $`f_{\mathrm{rec}}\{ϵ_\mathrm{L}f_{00}+f_{11}+\sqrt{ϵ_\mathrm{L}(1+ϵ)}(f_{01}\mathrm{cos}\alpha +\overline{f}_{01}\mathrm{sin}\alpha )`$ (14)
$``$ $`ϵ`$ $`(f_{11}\mathrm{cos}2\alpha +\overline{f}_{11}\mathrm{sin}2\alpha )`$ (15)
$`+`$ $`h`$ $`[\sqrt{ϵ_\mathrm{L}(1ϵ)}(f_{01}^{}\mathrm{sin}\alpha +\overline{f}_{01}^{}\mathrm{cos}\alpha )+\sqrt{1ϵ^2}f_{11}^{}]\},`$ (16)
where $`\sigma _\mathrm{M}`$ is the Mott scattering cross section,
$$\mathrm{\Omega }_\mathrm{f}=|𝒑_1^{}|E_1^{}|𝒑_2^{}|E_2^{},$$
(17)
$$f_{\mathrm{rec}}^1=1\frac{E_2^{}}{E_\mathrm{r}}\frac{𝒑_2^{}𝒑_\mathrm{r}}{|𝒑_2^{}|^2},$$
(18)
where $`E_\mathrm{r}`$ and $`𝒑_\mathrm{r}`$ are the relativistic energy and momentum of the residual nucleus,
$$K=\frac{e^4}{16\pi ^2}\frac{E_0^{}}{Q^2E_0(1ϵ)},$$
(19)
$$ϵ=\left(1+2\frac{|𝒒|^2}{Q^2}\mathrm{tan}^2\frac{\theta }{2}\right)^1,$$
(20)
$$ϵ_\mathrm{L}=\frac{Q^2}{|𝒒|^2}ϵ$$
(21)
and $`E_0`$ and $`E_0^{}`$ are the energies of the incident and outgoing electrons. The components of the lepton tensor $`\rho _{\lambda \lambda ^{}}`$ and $`\rho _{\lambda \lambda ^{}}^{}`$ can be deduced from Eq. (16).
The structure functions depend on the kinematical variables of the particular process under investigation. Even if a large number of variables are involved, at most four independent four-momenta are available in the Lorentz frame. Therefore, the hadron tensor can be expanded on the basis of four independent variables . In the case of two-nucleon emission different choices are possible. One can refer to the target four-momentum $`P^\mu `$, the momentum transfer $`q^\mu `$, the ejectile four-momentum $`p_1^\mu `$ and the four vector
$$\xi ^\mu =ϵ^{\alpha \beta \gamma \mu }q_\alpha p_{1\beta }^{}P_\gamma .$$
(22)
The expansion coefficients can depend on the invariants $`Q^2,Pq,qp_1^{},qp_2^{},Pp_1^{},Pp_2^{},p_1^{}p_2^{}`$ and $`\xi s_1`$, where $`s_1`$ is the spin of the outgoing nucleon with momentum $`𝒑_1^{}`$, which are scalars and $`Ps_1,qs_1,p_2^{}s_1`$ and $`\xi p_2^{}`$, which are pseudoscalars. It is worthwhile to notice that both scalars and pseudoscalars, both dependent on and independent of $`s_1`$, are available.
The four-vector $`s_1`$, in the reference frame where $`𝒑_1^{}=0`$, is given by $`s_1`$ = (0, $`𝒔_1`$) and can be boosted in the laboratory frame as
$$s_1^\mu =(\frac{𝒔_1𝒑_1^{}}{m},𝒔_1+\frac{𝒔_1𝒑_1^{}}{m(E_1^{}+m)}𝒑_1^{}).$$
(23)
Therefore, we obtain the invariants
$$\xi s_1=M_T𝒒\times 𝒑_1^{}𝒔_1,$$
(24)
$$Ps_1=\frac{M_T}{m}𝒑_1^{}𝒔_1,$$
(25)
where $`m`$ is the nucleon mass and $`M_T`$ the target mass, and
$$qs_1=\left(\frac{\omega }{m}\frac{𝒒𝒑_1^{}}{m(E_1^{}+m)}\right)𝒑_1^{}𝒔_1𝒒𝒔_1.$$
(26)
It is clear from the above equations that $`\xi s_1`$ is proportional to $`\widehat{𝑵}𝒔_1`$, $`Ps_1`$ to $`\widehat{𝑳}𝒔_1`$ and $`qs_1`$ has a component proportional to $`\widehat{𝑺}𝒔_1`$. Since we can obtain from these invariants both scalars and pseudoscalars which linearly contain the spin, il follows that no restrictions due to parity conservation can in general be derived concerning the dependence of the different structure functions on the components of the spin, in contrast with what happens in the ($`\stackrel{}{e},e^{}\stackrel{}{N}`$) reaction.
These restrictions are recovered in the particular case of a reaction where the outgoing nucleon momenta $`𝒑_1^{}`$ and $`𝒑_2^{}`$ and the momentum transfer $`𝒒`$ lie all in the same plane. In this case the invariant $`\xi p_2^{}`$ = 0, and neither scalars proportional to $`\widehat{𝑳}𝒔_1`$ and $`\widehat{𝑺}𝒔_1`$, nor pseudoscalars proportional to $`\widehat{𝑵}𝒔_1`$ are obtained. This result is independent of the electron plane and corresponds to the typical situation, in an unrestricted kinematics, of the ($`\stackrel{}{e},e^{}\stackrel{}{N}`$) reaction.
In general, for the ($`\stackrel{}{e},e^{}\stackrel{}{N}N`$) reaction, the dependence of the structure functions on the spin can be explicited remembering that for a particle with spin $`\frac{1}{2}`$ only a linear dependence is allowed, as
$$f_{\lambda \lambda ^{}}=h_{\lambda \lambda ^{}}^\mathrm{u}+\widehat{𝒔}𝒉_{\lambda \lambda ^{}},$$
(27)
where $`\widehat{𝒔}`$ is the unit vector in the spin space. Thus, when the polarization of the outgoing nucleon is considered, 36 structure functions $`h_{\lambda \lambda ^{}}^\mathrm{u}`$ and $`h_{\lambda \lambda ^{}}^i`$ are obtained. The explicit expressions of these structure functions in terms of the components of the hadron tensor $`W_{\alpha \alpha ^{}}^{\mu \nu }`$ can be easily obtained from Table I, by simply substituting the quantities $`W^{\mu \nu }`$, wherever they appear, with the following expressions:
$`W_{++}^{\mu \nu }+W_{}^{\mu \nu }`$ $`\mathrm{for}`$ $`\mathrm{i}=\mathrm{u},`$ (28)
$`W_+^{\mu \nu }+W_+^{\mu \nu }`$ $`\mathrm{for}`$ $`\mathrm{i}=\mathrm{x},`$ (29)
$`\mathrm{i}(W_{++}^{\mu \nu }W_{}^{\mu \nu })`$ $`\mathrm{for}`$ $`\mathrm{i}=\mathrm{y},`$ (30)
$`W_{++}^{\mu \nu }W_{}^{\mu \nu }`$ $`\mathrm{for}`$ $`\mathrm{i}=\mathrm{z}.`$ (31)
Usually, the quantities $`𝒉_{\lambda \lambda ^{}}`$ and $`𝒉_{\lambda \lambda ^{}}^{}`$ are projected onto the basis of unit vectors $`\widehat{𝑵}`$, $`\widehat{𝑳}`$ and $`\widehat{𝑺}`$ defined in Sec. II and the structure functions are thus given for the components $`i=N,L,S`$.
In the most general situation all the 36 structure functions do not vanish. When the electrons are on the $`𝒙`$$`𝒛`$ plane, i.e. when $`\alpha `$ = 0, the structure functions which are multiplied by $`\mathrm{sin}\alpha `$ or $`\mathrm{sin}2\alpha `$ in Eq. (16) do not contribute. In this case only 24 structure functions contribute to the cross section and polarizations. When the two outgoing nucleons and the momentum transfer lie in the same plane, we recover formally the typical situation of the ($`\stackrel{}{e},e^{}\stackrel{}{N}`$) reaction and only 18 structure functions do not vanish. In all these cases all the components of the polarization and of the polarization transfer coefficient are in principle different from zero. In a coplanar kinematics, i.e. when the momenta of the electrons and of the outgoing nucleons lie in the same plane, we recover the typical situation of the ($`\stackrel{}{e},e^{}\stackrel{}{N}`$) reaction in coplanar kinematics, where only 12 structure functions contribute and only $`P^N,P^L`$ and $`P^S`$ do not vanish. Finally, in the super-parallel kinematics only 5 structure functions survive, as in the parallel kinematics of the ($`\stackrel{}{e},e^{}\stackrel{}{N}`$) reaction. |
warning/0002/math0002148.html | ar5iv | text | # Higher order scattering on asymptotically Euclidean Manifolds
## 1. introduction
In this paper, we develop a scattering theory for powers of the Laplacian on a class of manifolds which includes perturbations of Euclidean space and apply this theory to obtain new inverse scattering results on Euclidean space. This theory is a natural extension of the work of Melrose , who developed a theory of scattering for the Laplacian on asymptotically Euclidean manifolds. We show that the higher order scattering matrix has very similar properties in this case. In particular, we show that the Melrose-Zworski calculus of Legendrian distributions, , can be applied to construct the Poisson operator for the scattering problem and thus deduce that the scattering matrix is a Fourier integral operator associated to geodesic flow at time $`\pi .`$ For higher order operators, this result appears to be new for the class of perturbations we consider even for $`^n`$. This theory is then also applied to extend the inverse results of Joshi and Sá Barreto on recovering the asymptotics of perturbations, , , , , , to this higher order case.
An asymptotically Euclidean manifold is a smooth manifold with boundary $`(X,X)`$ which is equipped with a scattering metric. A scattering metric is a smooth Riemannian metric, $`g,`$ on the interior of $`X`$ which blows up in a prescribed way at the boundary: there exists a product decomposition close to the boundary, $`p(x,y)[0,ϵ)\times X,`$ such that $`g`$ takes the form
(1.1)
$$g=\frac{dx^2}{x^4}+\frac{h(x,y,dy)}{x^2},$$
with $`h`$ smooth on the closed space and $`h_{|x=0}`$ a non-degenerate metric on $`X.`$ This is slightly different from Melrose’s definition in but was shown to be equivalent in . The co-tensor, $`h_{|x=0},`$ is independent of the decomposition chosen and thus we have a natural metric on $`X.`$ We shall assume throughout that a product decomposition close to the boundary has been chosen and fixed.
It is important to realize that $`^n`$ with the Euclidean metric is a special case of such a manifold. To see this, put $`x=|z|^1`$ and $`\omega =z|z|^1;`$ we then have
(1.2)
$$dz^2=\frac{dx^2}{x^4}+\frac{d\omega ^2}{x^2},$$
which also shows that the induced metric on the sphere at infinity is the Euclidean metric on the unit sphere.
Melrose showed that given a smooth function, $`f,`$ on $`X,`$ and $`\lambda \{0\}`$ that there is a unique function $`u,`$ smooth on the interior of $`X`$, of the form
(1.3)
$$e^{i\lambda /x}x^{\frac{n1}{2}}f_++e^{i\lambda /x}x^{\frac{n1}{2}}f_{},$$
with $`f_\pm `$ smooth functions on $`(X,X)`$ and $`f_{}`$ restricted to the boundary equal to $`f`$ such that $`(\mathrm{\Delta }\lambda ^2)u=0.`$ The scattering matrix is then defined to be the map, $`S(\lambda ),`$ on $`C^{\mathrm{}}(X),`$ defined by
(1.4)
$$S(\lambda ):ff_{+}^{}{}_{|X}{}^{}.$$
It was shown in that $`S(\lambda )`$ extends to a unitary operator on $`L^2(X)`$ with the density induced by $`h.`$ Melrose and Zworski () studied the micro-local structure of this operator and showed that it is a zeroth order, classical Fourier integral operator associated to geodesic flow at time $`\pi .`$ In the special case that $`h(x,y,dy)=h(0,y,dy)+𝒪(x^{\mathrm{}}),`$ Christiansen () and Parnovski () showed that for $`\lambda <0`$, modulo smoothing,
(1.5)
$$S(\lambda )=ie^{i\pi \sqrt{\mathrm{\Delta }_X+\frac{(n2)^2}{4}}}.$$
Here we develop analogous results for perturbations of powers of the Laplacian. For $`k`$, we define a short range perturbation of $`\mathrm{\Delta }^k`$ to be a symmetric differential operator of order $`2k1`$ which close to the boundary can be written in local coordinates in the form $`x^2P(x,y,x^2D_x,xD_y)`$ where $`P(x,y,\tau ,\eta )`$ is smooth in $`(x,y)`$ and is a polynomial of order $`2k1`$ in $`(\tau ,\eta ).`$ In particular, a real-valued, smooth function vanishing to second order at $`X`$ defines a short range perturbation. In Section 2 only, we allow certain pseudodifferential perturbations as well.
We prove
###### Theorem 1.1.
Let $`V`$ be a short range perturbation of $`\mathrm{\Delta }^k`$ and let $`\lambda \{0\}`$. Then, given $`fC^{\mathrm{}}(X)`$, there exists a smooth function, $`u,`$ on $`X^0`$ such that $`(\mathrm{\Delta }^k+V\lambda ^{2k})u=0`$ and $`u`$ is of the form
(1.6)
$$e^{i\lambda /x}x^{\frac{n1}{2}}f_++e^{i\lambda /x}x^{\frac{n1}{2}}f_{},$$
with $`f_\pm `$ smooth functions on $`X`$ such that $`f_{}`$ restricted to the boundary is equal to $`f.`$ The function $`u`$ is unique modulo smooth functions vanishing to infinite order at $`X.`$
The non-uniqueness here corresponds to the possibility of embedded discrete spectrum. The scattering matrix $`S(\lambda )`$ can then be defined precisely as before as the indeterminacy will not affect the lead term at the boundary, and we show that it has a unitary extension.
We also prove
###### Theorem 1.2.
Let $`V`$ be a short range perturbation of $`\mathrm{\Delta }^k`$ and let $`\lambda \{0\}`$. Then $`S(\lambda )`$ is a zeroth order classical Fourier integral operator associated to geodesic flow at time $`\pi `$ on $`X.`$
It is interesting to compare the scattering matrix for different powers and the following is an immediate corollary to our construction.
###### Corollary 1.1.
Let $`0<k_1k_2,`$ and suppose $`V_j`$ is a short range perturbation of $`\mathrm{\Delta }^{k_j},`$ of the form $`x^lp_j(x,y,x^2D_x,xD_y),`$ with $`p_j(x,y,\xi ,\eta )`$ a polynomial in $`(\xi ,\eta ).`$ Let $`S_j(\lambda )`$ be the scattering matrix associated to $`\mathrm{\Delta }^{k_j}+V_j\lambda ^{2k_j}.`$ Then $`S_1(\lambda )S_2(\lambda )`$ is a Fourier integral operator of order $`1l.`$
It is therefore immediate that (1.5) holds also for the higher order scattering matrix when $`V=0.`$
In the special case of $`^n`$ we study the problem of recovering asymptotics of a perturbation from scattering data. As in , we need an aradiality condition to recover the perturbation. In fact, in it was observed that recovery is not possible without it. We shall say a perturbation is aradial modulo Schwartz functions if it is asymptotically equal to a sum
(1.7)
$$\underset{l=2}{\overset{\mathrm{}}{}}\underset{\alpha }{}f_{\alpha ,l}D_z^\alpha ,$$
with each term $`\underset{\alpha }{}f_{\alpha ,l}D_z^\alpha `$ of the form $`|z|^l`$ times a composition of vector fields tangent to the sphere. The aradiality conditions allows any zeroth order perturbation.
###### Theorem 1.3.
Let $`V_1,V_2`$ be short range perturbations of $`\mathrm{\Delta }^k`$ on $`(^n,dz^2)`$ such that $`V_1V_2`$ is an aradial differential operator of order $`l.`$ Let $`S_j(\lambda )`$ be the scattering matrix associated to $`\mathrm{\Delta }^k+V_j\lambda ^{2k},`$ and suppose that for $`l+1`$ values of $`\lambda >0,`$ $`S_1(\lambda )S_2(\lambda )`$ is smoothing. Then the coefficients of $`V_1`$ and $`V_2`$ agree modulo Schwartz functions.
Note that we assume neither that $`V_1,V_2`$ are of order $`l`$ nor that they are aradial.
Our approach to this higher order scattering problem is highly influenced by that of Melrose, , and that of Melrose-Zworski, . In particular, we use the scattering calculus developed by Melrose and used for the case $`k=1`$ to study the general case and use techniques similar to those of to establish the existence of the scattering matrix. To establish the micro-local structure of the scattering matrix we proceed as in to construct the Poisson operator for the scattering problem as a Legendrian distribution associated to a pair of intersecting Legendrian submanifolds using the calculus developed there.
To prove the inverse result, we choose to follow the approach of rather than in order to increase the readability of the paper for non-experts. In particular, we establish our results for the case of $`^n`$ without explicitly using the Melrose-Zworski Legendrian calculus. The reader may regard these proofs as a warm-up for the construction of the Poisson operator in the general case. As in , , , , , the proof proceeds by establishing that the principal symbol of the difference of the scattering matrices determines and is determined by a weighted integral of the lead term of the difference of the perturbations over geodesics of length $`\pi .`$ The injectivity of this transformation is then deduced by using some elementary calculus and some deep results of Bailey and Eastwood, , on the integral geometry of tensor fields on projective space.
The main reference for higher order scattering on $`^n`$ is Chapter 14 of where a scattering theory for perturbations of a much more general class of constant coefficient operators on $`^n`$ is developed - this extends ideas developed by Agmon and Hörmander in and Agmon in . The complementary problem of studying the recovery of compactly supported perturbations of higher order operators on $`^n`$ was studied by Liu in .
The results of Bailey and Eastwood () which we use in section 4 to prove the inverse results are restricted to projective spaces and spheres. Since in addition, we believe the principal interest in this problem is in Euclidean space, we restrict ourselves to studying that important special case as the analysis is much more accessible. The problem of recovering the asymptotics of metrics for the $`k=1`$ case has been studied in and it is likely that a similar result could be proven in the higher order case.
The question of recovering the entire perturbation from the scattering matrix at fixed energy is an interesting one but is still open even in the case with $`k=1.`$
We are grateful to Mike Eastwood for helpful conversations and to the Royal Society for the travel grant which made those conversations possible. We also thank the London Mathematical Society for supporting this collaborative research through its small grants scheme. The first author is grateful for partial support from a University of Missouri S.R.F.
## 2. Basic scattering theory of $`\mathrm{\Delta }^k+V`$
In this section, we construct solutions of $`\mathrm{\Delta }^k+V\lambda ^{2k}`$ having specified behaviour at the boundary, leading to the definition of the scattering matrix.
Throughout, $`\mathrm{\Delta }`$ denotes the Laplacian associated to a scattering metric on a compact manifold $`(X,X)`$, $`k`$ is a positive integer, and $`\lambda \{0\}`$. We assume a product decomposition with boundary defining function $`x`$ has been fixed. We call the smooth functions vanishing to infinite order at the boundary the Schwartz functions on $`X.`$ We use the product decomposition to extend functions on the boundary smoothly into $`X`$ by making them constant in the normal direction and cutting off.
The results of this section are very closely related to the results of . We have tried to state the results in a manner accessible to those unfamiliar with that paper, but in order to avoid repetition we omit proofs which follow essentially as in , only giving some indication of how they need to be modified to work in this setting. We recall some of the definitions and results of but refer the reader to the original for full details.
In this section only, we will allow a somewhat larger class of perturbations of $`\mathrm{\Delta }^k`$. In order to describe them, we recall some notation.
Let $`𝒱_{\mathrm{sc}}(𝒳)`$ be the space of all smooth vector fields of finite length with respect to a scattering metric. Near a point on the boundary, $`x^2\frac{}{x}`$ and $`x\frac{}{y_i}`$, $`i=1,\mathrm{},n1,`$ form a basis for $`𝒱_{\mathrm{sc}}(𝒳)`$, where $`y_i`$ are coordinates on $`X`$. The set of scattering differential operators of order $`m`$ is
$$\mathrm{Diff}_{\mathrm{sc}}^m(X)=\mathrm{span}_{0jm}(𝒱_{\mathrm{sc}}(𝒳))^𝒿.$$
The short range perturbations which we are allowing in the remainder of this paper are elements of $`x^2\mathrm{Diff}_{\mathrm{sc}}^{2k1}(X)`$.
Most of the results of this section hold for a wider class of perturbations: elements of $`\mathrm{\Psi }_{\mathrm{sc}}^{2k1,2}(X)`$, part of the (small) calculus of scattering pseudodifferential operators. We refer the reader to \[14, Sections 4 and 5\] for the full definition and some properties of $`\mathrm{\Psi }_{\mathrm{sc}}^{m,l}(X)`$. However, we remark that $`x^j\mathrm{Diff}_{\mathrm{sc}}^m(X)\mathrm{\Psi }_{\mathrm{sc}}^{m,j}(X)`$ and $`(\mathrm{\Delta }z)^1\mathrm{\Psi }_{\mathrm{sc}}^{2,0}(X)`$ when $`z[0,\mathrm{})`$ (\[14, Theorem 1\]). Moreover, if we consider the manifold $`(^n,dz)`$ (which by an earlier discussion under radial compactification becomes a compact manifold $`(^n)_{rc}`$ with scattering metric), the operators corresponding to $`\mathrm{\Psi }_{\mathrm{sc}}^{m,l}((^n)_{rc})`$ have Schwartz kernels of the form
$$A(z,z^{})=\frac{1}{(2\pi )^n}e^{i(zz^{})\zeta }a_L(z,\zeta )𝑑\zeta $$
with $`a_L`$ satisfying
$$|D_z^\alpha D_\zeta ^\beta a_L(z,\zeta )|C_{\alpha \beta }(1+|z|)^{l|\alpha |}(1+|\zeta |)^{m|\beta |}$$
(\[14, (4.1),(4.2)\]).
###### Definition 2.1.
We shall say a differential operator $`V`$ is a short range perturbation of $`\mathrm{\Delta }^k`$ if it is symmetric and if $`Vx^2\mathrm{Diff}_{\mathrm{sc}}^{2k1}(X)`$. We call a pseudodifferential operator a pseudodifferential short range perturbation of $`\mathrm{\Delta }^k`$ if it is symmetric and it is an element of $`\mathrm{\Psi }_{\mathrm{sc}}^{2k1,2}(X)`$.
Throughout this section $`V`$ is a pseudodifferential short range perturbation of $`\mathrm{\Delta }^k`$.
We prove
###### Proposition 2.1.
Let $`V`$ be a pseudodifferential short range perturbation of $`\mathrm{\Delta }^k`$ and $`\lambda \{0\}.`$ Given $`fC^{\mathrm{}}(X),`$ there exists $`uC^{\mathrm{}}(X),`$ unique modulo Schwartz functions, such that $`u_{|X}=f`$ and $`(\mathrm{\Delta }^k+V\lambda ^{2k})(e^{i\lambda /x}x^{\frac{n1}{2}}u)`$ is Schwartz.
Note we make no assumptions on the sign of $`\lambda .`$
This will follow easily once we have proven
###### Lemma 2.1.
Let $`V`$ be a pseudodifferential short range perturbation of $`\mathrm{\Delta }^k`$ and $`\lambda \{0\}.`$ If $`fC^{\mathrm{}}(X),`$ then
$$(\mathrm{\Delta }^k+V\lambda ^{2k})(e^{i\lambda /x}x^{\frac{n1}{2}+\alpha }f)=e^{i\lambda /x}(k\lambda ^{2k1}C_\alpha x^{\frac{n+1}{2}+\alpha }f+x^{\frac{n+3}{2}+\alpha }g)$$
with $`gC^{\mathrm{}}(X),`$ where $`C_0=0`$ and $`C_\alpha 0`$ for $`\alpha 0.`$
###### Proof.
For $`k=1,`$ this follows as in , using \[15, Lemma 8\] for the mapping properties of pseudodifferential short range perturbations. For general $`k,`$ one simply iterates. ∎
Now to prove Proposition 2.1, we simply choose the lead term to be $`f`$ and then repeatedly iterate to compute all the terms of the Taylor series. The result follows from Borel’s lemma. The uniqueness follows from the fact that $`C_\alpha 0`$ for $`\alpha 0.`$
We will repeatedly use the operator
(2.1)
$$Q=Q(\lambda )=\underset{j=0}{\overset{k1}{}}\lambda ^{2j}\mathrm{\Delta }^{kj1}$$
since $`\mathrm{\Delta }^k\lambda ^{2k}=Q(\mathrm{\Delta }\lambda ^2)=(\mathrm{\Delta }\lambda ^2)Q`$. Using the same techniques as Theorem 1 of , one can show that for $`\lambda \{0\}`$, $`Q(\lambda )^1\mathrm{\Psi }_{\mathrm{sc}}^{2k+2,0}(X)`$.
A key idea of many of the proofs is that
(2.2)
$$\mathrm{\Delta }^k\lambda ^{2k}+V=Q(\mathrm{\Delta }\lambda ^2+V^{}),$$
where $`V^{}=Q^1V\mathrm{\Psi }_{\mathrm{sc}}^{1,2}(X)`$. The importance of this is that the symbol of $`V^{}`$ vanishes to second order at the boundary. As noted in Remark 3 of , many of the results there hold if $`\mathrm{\Delta }`$ is replaced by $`\mathrm{\Delta }+W`$, when $`Wx^2C^{\mathrm{}}(X)`$, because they depend on the properties of the principal symbol and the boundary symbol, which are unchanged by the addition of such a $`W`$. What we are doing here is allowing a somewhat more general perturbation, but with the same kind of decay at the boundary. Thus, the results of Melrose’s Propositions 9-11 hold if $`\mathrm{\Delta }`$ is replaced by $`\mathrm{\Delta }+W`$, with $`W\mathrm{\Psi }_{\mathrm{sc}}^{1,2}(X)`$.
The following proposition is closely related to Proposition 11 of and follows from the modifications of Propositions 9-11 indicated above and from (2.2).
###### Proposition 2.2.
If $`\lambda \{0\}`$, $`uL_{sc}^2(X)`$, $`V`$ is a pseudodifferential short range perturbation of $`\mathrm{\Delta }^k`$ and $`(\mathrm{\Delta }^k+V\lambda ^{2k})u=0`$, then $`u`$ is Schwartz.
If $`k=1`$ and $`V`$ is a short range (differential) perturbation, then $`\mathrm{\Delta }+V`$ has no positive eigenvalues. However, for (differential) short range perturbations of $`\mathrm{\Delta }^k`$, $`k>1`$, the situation can be different.
###### Proposition 2.3.
There are (differential) short-range perturbations $`V`$ of $`\mathrm{\Delta }^k`$, $`k>1`$, even, such that $`\mathrm{\Delta }^k+V`$ has positive eigenvalues.
###### Proof.
We give two ways of constructing such examples. Let $`j`$.
Using the min-max principle, for any asymptotically Euclidean manifold $`X`$ one can find $`\stackrel{~}{V}x^2C^{\mathrm{}}(X)`$ so that $`\mathrm{\Delta }+\stackrel{~}{V}`$ has a negative eigenvalue $`\sigma `$. Then $`(\mathrm{\Delta }+\stackrel{~}{V})^{2j}`$ has a positive eigenvalue $`\sigma ^{2j}`$. It is possible to choose $`\stackrel{~}{V}`$ to be compactly supported and such that $`\mathrm{\Delta }+\stackrel{~}{V}`$ has many negative eigenvalues.
To construct a Schwartz potential $`V`$ such that $`\mathrm{\Delta }^{2j}+V`$ has an embedded eigenvalue, choose $`\tau >0`$ and let $`x`$ be a globally defined boundary-defining function, $`xC^{\mathrm{}}(X)`$. Then
$$(\mathrm{\Delta }^{2j}\tau ^{4j})e^{\tau /x}x^{(n1)/2}=𝒪(𝓍^{(𝓃+\mathcal{3})/\mathcal{2}}^{\tau /𝓍}).$$
Just as in Proposition 2.1 we can use Lemma 2.1 to successively solve away the errors, resulting in a $`u`$ which satisfies
$$(\mathrm{\Delta }^{2j}\tau ^{4j})u=g=𝒪(𝓍^{\mathrm{}}^{\tau /𝓍})$$
with $`u=e^{\tau /x}x^{(n1)/2}(1+𝒪(𝓍))`$. In fact, by doing the asymptotic summation judiciously, we can ensure that $`x^{(n1)/2}e^{\tau /x}u`$ is nowhere vanishing. Set $`V=u^1g`$. Then $`\mathrm{\Delta }^{2j}+V`$ has eigenvalue $`\tau ^{4j}`$ with eigenfunction $`u`$. ∎
We continue with the results needed to define the scattering matrix.
###### Theorem 2.1.
If $`\lambda \{0\}`$, $`V\mathrm{\Psi }_{\mathrm{sc}}^{2k1,2}(X)`$ is symmetric, and $`f`$ is a Schwartz function on $`X`$, orthogonal to the $`L^2`$ null space of $`\mathrm{\Delta }^k+V\lambda ^{2k}`$ if there is any, then there exists $`u`$ with $`x^{(n1)/2}e^{i\lambda /x}uC^{\mathrm{}}(X)`$ such that
$$(\mathrm{\Delta }^k+V\lambda ^{2k})u=f.$$
This is an analogue of \[14, Propositions 12 and 14\].
###### Proof.
As in Proposition 14 of ,
$$u=\underset{t0}{lim}(\mathrm{\Delta }^k+V(\lambda +it)^{2k})^1f$$
where the limit exists in $`x^{1/2\delta }H_{\mathrm{sc}}^{\mathrm{}}(X)`$ for any $`\delta >0`$. To show that the limit exists, let $`u_t=(\mathrm{\Delta }^k+V(\lambda +it)^{2k})^1f`$. Then, since $`f`$ is orthogonal to the null space of $`\mathrm{\Delta }^k+V\lambda ^{2k}`$, so is $`u_t`$ for $`t>0`$. Having made this observation, the proof follows as in \[14, Proposition 14\], giving us the same additional microlocal regularity as well.
To finish the proof, one uses the analog of Proposition 12 of . That is, if $`\stackrel{~}{u}C^{\mathrm{}}(X)`$, $`WF_{\mathrm{sc}}^{,1/2}(\stackrel{~}{u})R_+(\lambda )=\mathrm{}`$, and $`(\mathrm{\Delta }+Q^1V\lambda ^2)\stackrel{~}{u}\dot{C}^{\mathrm{}}(X),`$ then $`e^{i\lambda /x}x^{(n1)/2}\stackrel{~}{u}C^{\mathrm{}}(X)`$. The proof follows just as the proof of \[14, Proposition 12\], first noting that $`[\stackrel{~}{\mathrm{\Delta }}_0,Q^1V]\mathrm{\Psi }_{\mathrm{sc}}^{2,1}(X)`$ and thus
$$[\stackrel{~}{\mathrm{\Delta }}_0,Q^1V]:x^sH_{\mathrm{sc}}^{\mathrm{}}(X)x^{s+1}H_{\mathrm{sc}}^{\mathrm{}}(X).$$
The second observation that allows the proof to proceed as in \[14, Prop. 12\] is that if $`\mathrm{Diff}_c^l(X)uH_{\mathrm{sc}}^{\mathrm{},m}(X)`$, then $`\mathrm{Diff}_c^{ll^{}}(X)A\mathrm{Diff}_c^l^{}(X)uH_{\mathrm{sc}}^{\mathrm{},m+s}(X)`$ when $`l^{}l`$ is a nonnegative integer and $`A\mathrm{\Psi }_{\mathrm{sc}}^{,s}(X)`$. This then allows the inductive step in the proof of \[14, Proposition 12\] to proceed just as it does there. ∎
We shall also use the following boundary pairing result, the analogue of \[14, Proposition 13\].
###### Proposition 2.4.
Suppose $`V\mathrm{\Psi }_{\mathrm{sc}}^{2k1,2}(X)`$ is symmetric, $`\lambda \{0\}`$, $`u_i`$, $`i=1,2`$ satisfies $`(\mathrm{\Delta }^k+V\lambda ^{2k})u_i=f_i`$, $`f_i`$ is Schwartz, and $`u_i=x^{(n1)/2}(e^{i\lambda /x}a_i^++e^{i\lambda /x}a_i^{})`$, with $`a_i^\pm C^{\mathrm{}}(X).`$ Let $`b_i^\pm =(a_i^\pm )_{|X}`$. Then
$$2ik\lambda ^{2k1}_X(b_1^+\overline{b_2^+}b_1^{}\overline{b_2^{}})𝑑h=_Xu_1\overline{f_2}f_1\overline{u_2}dg.$$
###### Proof.
Using $`Q=Q(\lambda )=_{j=0}^{k1}\lambda ^{2j}\mathrm{\Delta }^{kj1}`$, we have
$`{\displaystyle _X}u_1\overline{f_2}f_1\overline{u_2}dg`$
$`={\displaystyle _X}u_1\overline{(\mathrm{\Delta }\lambda ^2+VQ^1)Qu_2}Q(\mathrm{\Delta }\lambda ^2+Q^1V)u_1\overline{u_2}dg`$
(2.3) $`={\displaystyle _X}u_1\overline{(\mathrm{\Delta }\lambda ^2+VQ^1)Qu_2}(\mathrm{\Delta }\lambda ^2+Q^1V)u_1\overline{Qu_2}dg`$
as $`Q`$ is self-adjoint and $`(\mathrm{\Delta }\lambda ^2+Q^1V)u_1\dot{C}^{\mathrm{}}(X)`$. Recall that $`V`$ is self-adjoint as well. Choose a function $`\varphi C^{\mathrm{}}()`$ such that $`\varphi (t)=0`$ if $`t<1`$ and $`\varphi (t)=2`$ if $`t>2`$. Then (2) is equal to
$$\begin{array}{c}\underset{ϵ0}{lim}_X\varphi (x/ϵ)\left(u_1\overline{(\mathrm{\Delta }\lambda ^2+VQ^1)Qu_2}(\mathrm{\Delta }\lambda ^2+Q^1V)u_1\overline{Qu_2}\right)𝑑g\hfill \\ \hfill =\underset{ϵ0}{lim}_X[\mathrm{\Delta }+Q^1V,\varphi (x/ϵ)]u_1\overline{Qu_2}𝑑g.\end{array}$$
Since $`V`$ is short-range, $`lim_{ϵ0}_X[Q^1V,\varphi (x/ϵ)]u_1\overline{Qu_2}𝑑g=0.`$ To compute the remainder, we just use the corresponding results from \[14, Proposition 13\] and the fact that at the boundary of $`X`$,
$$Qu_2=k\lambda ^{2k2}x^{(n1)/2}e^{i\lambda /x}b_2^++k\lambda ^{2k2}x^{(n1)/2}e^{i\lambda /x}b_2^{}+𝒪(𝓍^{(𝓃+\mathcal{1})/\mathcal{2}}).$$
We will need
###### Theorem 2.2.
If $`\lambda \{0\}`$, $`V`$ is a pseudodifferential short range perturbation of $`\mathrm{\Delta }^k`$, and $`fC^{\mathrm{}}(X)`$, then there exists $`f_\pm C^{\mathrm{}}(X)`$ such that $`u=x^{\frac{n1}{2}}\left(e^{i\lambda /x}f_++e^{i\lambda /x}f_{}\right)`$ satisfies
$$(\mathrm{\Delta }^k+V\lambda ^{2k})u=0$$
and the restriction of $`f_{}`$ to $`X`$ is $`f.`$ The function $`u`$ is unique if $`\lambda ^{2k}`$ is not an eigenvalue of $`\mathrm{\Delta }^k+V`$ and is unique up to the addition of a Schwartz function if $`\lambda ^{2k}`$ is an eigenvalue.
###### Proof.
Most of the proof follows from Proposition 2.1 and Theorem 2.1. We use Theorem 2.1 to solve away the error obtained by constructing the formal expansion as in Proposition 2.1. We note that since the error we wish to solve away is of the type $`(\mathrm{\Delta }^k+V\lambda ^2)g`$, where $`gx^{1/2ϵ}L_{sc}^2(X)`$, and since eigenfunctions are Schwartz we may still integrate by parts to obtain that the error is orthogonal to the eigenspace. Finally, suppose that there are two such $`u`$ satisfying the conditions of the theorem. Then their difference $`v`$ satisfies $`(\mathrm{\Delta }^k+V\lambda ^{2k})v=0`$ and $`v=x^{(n1)/2}e^{i\lambda /x}g_++x^{(n+1)/2}e^{i\lambda /x}g_{}`$ with $`g_\pm C^{\mathrm{}}(X).`$ Then the boundary pairing result of Proposition 2.4 gives us that $`(g_+)_{|X}=0`$, and thus $`v`$ is in $`L_{sc}^2(X)`$. Consequently, $`v`$ is an eigenfunction and is thus Schwartz. ∎
We can therefore define the scattering matrix as in the usual case:
###### Definition 2.2.
The scattering matrix is a map on $`C^{\mathrm{}}(X)`$ taking a function $`f`$ to $`f_+`$ restricted to $`X`$ where $`f,f_+`$ are as in Theorem 2.2.
Using Proposition 2.4 and the uniqueness it implies, it is easy to check that the scattering matrix can be extended to a unitary operator on $`L^2(X)`$ which satisfies $`S(\lambda )^1=S(\lambda )`$.
We remark that if $`V\mathrm{\Psi }_{\mathrm{sc}}^{2k1,2}(X)`$ is symmetric, then $`\mathrm{\Delta }^k+V`$, initially considered as an operator on the Schwartz functions, has a unique self-adjoint extension to a domain in $`L_{\mathrm{sc}}^2(X)`$. We end this section by describing the continuous part of its spectral measure. Let $`P`$ be the Poisson operator whose existence is given by Theorem 2.2:
$$P(\lambda ):C^{\mathrm{}}(X)fux^{(n1)/2}e^{i\lambda /x}C^{\mathrm{}}(X)+x^{(n1)/2}e^{i\lambda /x}C^{\mathrm{}}(X)$$
where $`u`$ is the function in Theorem 2.2. Then the continuous part of the spectral measure of $`\mathrm{\Delta }^k+V`$ is given by
$$dE_c(\lambda )=\frac{1}{2\pi }P(\lambda )P^{}(\lambda )d\lambda $$
where $`\lambda ^{2k}`$ is the spectral variable and $`\lambda [0,\mathrm{})`$. This can be seen by first noting that if $`V=0`$ this follows from \[4, Lemma 2.2\] (see also \[5, Lemma 5.2\]) and the general case then follows as in \[17, Appendix to XI.6\]. We note that $`\mathrm{\Delta }^k+V`$ has no singular continuous spectrum, though, as previously noted, there may be discrete spectrum.
## 3. The Poisson Operator for Euclidean space
As an introduction to the general case, we use the techniques of to construct the Poisson operator in the Euclidean case, where the construction can be done more explicitly. Following we look for a Poisson operator in the form $`e^{i\lambda z.\omega }a(z,\omega )`$ with $`a`$ a polyhomogeneous symbol in $`z`$ and $`\omega `$ is a smooth parameter.
Let $`\stackrel{~}{a}(z,\omega )`$ be any polyhomogeneous symbol in $`z`$, with $`\omega `$ a parameter. Then
(3.1)
$$\mathrm{\Delta }(e^{i\lambda z.\omega }\stackrel{~}{a})=(\lambda ^2\stackrel{~}{a}2i\lambda \omega .\frac{\stackrel{~}{a}}{z}+\mathrm{\Delta }\stackrel{~}{a})e^{i\lambda z.\omega },$$
so we conclude that
(3.2)
$$\mathrm{\Delta }^k(e^{i\lambda z.\omega }\stackrel{~}{a})=(\lambda ^{2k}\stackrel{~}{a}2ki\lambda ^{2k1}\omega .\frac{\stackrel{~}{a}}{x}+b)e^{i\lambda z.\omega },$$
with $`b`$ a symbol two orders lower than $`\stackrel{~}{a}.`$
Now $`V`$ is a short range perturbation of $`\mathrm{\Delta }^k`$, which for $`^n`$ means that
(3.3)
$$V=\underset{|\alpha |2k1}{}f_\alpha D_x^\alpha ,$$
with $`f_\alpha S_{cl}^2(^n).`$ We conclude that it maps $`e^{i\lambda z.\omega }S^m`$ to $`e^{i\lambda z.\omega }S^{m2}`$ and hence that if $`\stackrel{~}{a}S_{phg}^m`$ we have,
(3.4)
$$(\mathrm{\Delta }^k+V\lambda ^{2k})(e^{i\lambda z.\omega }\stackrel{~}{a})=(2ki\lambda ^{2k1}\omega .\frac{\stackrel{~}{a}}{z}+c)e^{i\lambda z.\omega },$$
with $`cS_{phg}^{m2}.`$
Returning to the Poisson operator, by taking the lead term of $`a`$ to be the constant function $`1`$ we have an error in $`S_{phg}^2.`$ To solve away this and subsequent error terms, we proceed precisely as in \[9, Section 2\], using a term in $`S_{phg}^{j+1}`$ to solve away an error term in $`S_{phg}^j`$ (modulo terms in $`S_{phg}^{j1}`$). We solve the transport equations,
$$2ki\lambda ^{2k1}\omega .\frac{b}{z}=d,$$
where $`d`$ is the error, along the geodesics on the unit sphere from $`\omega `$ to $`\omega .`$ The only difference from the $`k=1`$ case of is a factor of $`\lambda ^{2(k1)}.`$ As before, the transport equation degenerates on approach to $`\omega `$ and in particular blows up like $`(\pi s)^r`$ as $`s,`$ the geodesic distance from $`\omega ,`$ tends to $`\pi ,`$ when solving for the term of homogeneity of order $`r.`$
So as before we need a different ansatz close to $`\omega .`$ We can locally in $`\omega `$ rotate our coordinate system so that $`\omega `$ is the north pole. So close to the south pole, we look for the Poisson operator in the form
(3.5)
$$\underset{0}{\overset{\mathrm{}}{}}\left(\frac{1}{S|z|}\right)^\gamma S^\alpha e^{i\lambda (Sz^{}.\mu \sqrt{1+S^2}|z|)}a(\frac{1}{S|z|},S,\mu )𝑑S𝑑\mu ,$$
with $`a(t,S,\mu )`$ a smooth function compactly supported on $`[0,ϵ)\times [0,ϵ)\times S^{n2}.`$ We denote the class of functions that can be written in this form plus a Schwartz error by $`I^{\gamma ,\alpha }.`$ We recall from that this class is asymptotically complete in $`\gamma .`$ It follows from a stationary phase computation carried out in that away from the south pole this is equivalent to the previous ansatz with a symbol of order $`\gamma +\frac{n1}{2}.`$
We recall from ,
###### Proposition 3.1.
If $`u(z,\omega )I^{\gamma ,\alpha }`$ and $`fC^{\mathrm{}}(X\times X)`$ then
$$e^{i\lambda |z|}u(|z|\theta ,\omega )f(\theta ,\omega )𝑑\theta 𝑑\omega $$
is a smooth symbolic function in $`|z|`$ of order $`1\alpha `$ and its lead coefficient is $`|z|^{\alpha 1}K,f`$ where $`K`$ is the pull-back of the Schwartz kernel of a pseudo-differential operator of order $`\alpha \gamma (n2)`$ by the map $`\theta \theta .`$ The principal symbol of $`K`$ determines and is determined by the lead term of the symbol, $`a(t,S,\mu ),`$ of $`u`$ as $`S0+.`$
We also have as a special case of Proposition 15 from ,
###### Proposition 3.2.
If $`uI^{\gamma ,\alpha }`$, then
$$(\mathrm{\Delta }^k\lambda ^{2k})uI^{\gamma +1,\alpha +1},$$
and
$$VuI^{\gamma +2,\alpha +2}.$$
It is also important to note that,
###### Lemma 3.1.
If $`uI^{\gamma ,\frac{n3}{2}}`$ and $`(\mathrm{\Delta }^k+V\lambda ^{2k})uI^{\gamma +j,\frac{n1}{2}}`$, then $`(\mathrm{\Delta }^k+V\lambda ^{2k})uI^{\gamma +j,\frac{n+1}{2}}.`$
This follows from a slight extension of the argument in the proof of Lemma 3.2 in .
We also recall from that
###### Proposition 3.3.
If $`uI^{\mathrm{},\alpha }=\underset{\gamma }{}I^{\gamma ,\alpha }`$, then $`u=e^{i\lambda |z|}f(z),`$ with $`f`$ a classical symbol of order $`\alpha 1.`$
To carry out our construction we first use the original ansatz, obtaining a symbol which blows up on approach to the south pole. Near the south pole, we use the second ansatz, (3.5), taking $`\gamma =\frac{n1}{2},\alpha =\frac{n3}{2}.`$ The argument is then identical to the one in . At the end of the construction, we obtain an approximate Poisson operator, $`\stackrel{~}{P},`$ such that $`(\mathrm{\Delta }^k+V\lambda ^{2k})\stackrel{~}{P}(\lambda )`$ is Schwartz. This error will in fact be orthogonal to the eigenspace at energy $`\lambda `$ (if there is one) as the eigenfunctions are necessarily Schwartz and the orthogonality follows from self-adjointness and integration by parts. We can therefore apply the resolvent (Theorem 2.1) and gain a term of the from $`e^{i\lambda |z|}f`$ with $`f`$ smooth, solving away the error.
## 4. The inverse problem
In this section we prove Theorem 1.3. Suppose we have two short range perturbations, $`V_1,V_2,`$ of $`\mathrm{\Delta }^k`$ and we want to compare their scattering matrices. In particular, suppose
(4.1)
$$V_1V_2=\underset{|\alpha |2k1}{}a_\alpha (z)D_z^\alpha ,$$
with $`a_\alpha S_{phg}^{1r}(^n).`$ If we carry out our construction for each $`V_j,`$ then the first $`r`$ terms of the construction with the first ansatz will be the same but the term of homogeneity $`r`$ will be different. In particular, the forcing terms in the transport equations will differ by
(4.2)
$$W_r=e^{i\lambda z.\omega }(V_1V_2)e^{i\lambda z.\omega }=\underset{|\alpha |2k1}{}\omega ^\alpha a_{\alpha ,r1}\lambda ^{|\alpha |},$$
where $`a_{\alpha ,r1}`$ is the lead term of $`a_\alpha .`$ We therefore conclude that the terms of homogeneity $`r`$ will differ at $`(\gamma (s),\omega )`$ by
(4.3)
$$\frac{i|z|^r}{2k\lambda ^{2k1}(\mathrm{sin}s)^r}\underset{0}{\overset{s}{}}W_r(\gamma (s^{}),\omega )(\mathrm{sin}s^{})^{r1}𝑑s^{},$$
where $`\gamma `$ is a geodesic on the sphere running from $`\omega `$ to $`\omega .`$ As we have shown that the lead singularity of the first ansatz as $`s\pi `$ is essentially the principal symbol of the scattering matrix, we conclude by the same argument that the difference of the scattering matrices will be of order $`r`$ and that the principal symbol of the difference will determine and be determined by,
(4.4)
$$\frac{1}{\lambda ^{2k1}}\underset{0}{\overset{\pi }{}}W_r(\gamma (s),\omega )(\mathrm{sin}s^{})^{r1}𝑑s,$$
for all geodesics $`\gamma `$ with $`\omega `$ equal to $`\gamma (0).`$
Note that as $`W`$ depends polynomially on $`\lambda `$, so does the principal symbol. In particular, if we know the scattering matrix for $`2k`$ different values of $`\lambda >0`$, then we can separate out the parts coming from the differing orders of $`|\alpha |`$ in the forcing term. More generally, if we have a perturbation of order $`l2k1`$, then we determine the asymptotics from $`l+1`$ values. So in the sequel, assuming the perturbation is of order $`l2k1`$, we only consider forcing terms of the form
(4.5)
$$\underset{|\alpha |=l}{}a_\alpha (z)D_z^\alpha ,$$
where $`l2k1.`$ To solve the inverse problem, then, we must show that $`a_\alpha (z)`$, $`|\alpha |=l`$, can be recovered from knowledge of the transform (4.4), where we replace the sums in (4.1) and (4.2) by sums only over the terms with $`|\alpha |=l`$. That is, we need to show that the transform (4.4), a weighted integral along geodesics of length $`\pi `$, is invertible. The injectivity of the transform was proven in for the case $`l=0.`$ For higher order perturbations ($`l>0`$) the question is more subtle and requires the notion of aradiality. We say a perturbation is aradial if it has no radial component, that is, if the perturbation is a span of vector fields tangent to the sphere. Clearly, all zeroth order perturbations are aradial. The injectivity of the transform (4.4) for aradial first order perturbations was shown in and for second order aradial perturbations in . The impossibility of recovering first order perturbations that are not aradial was also shown in . We now look at the general case.
First we re-express the forcing term more invariantly. We regard the perturbation as an $`l`$-form, $`\mu =\underset{|\alpha |=l}{}a_{\alpha ,r1}(z)dz^\alpha .`$ Since $`dz^\alpha `$ is symmetric, $`\mu `$ is symmetric. The aradiality means that $`\mu `$ makes sense as a form on the sphere and is determined by its values as a map from the tangent space of the sphere to $`.`$
To re-express the forcing term we rotate so that $`\omega =(0,\mathrm{},0,1)`$ and $`\gamma (s)=(0,\mathrm{},0,\mathrm{sin}(s),\mathrm{cos}(s)).`$ We then have $`\frac{d\gamma }{ds}(s)=(0,\mathrm{},0,\mathrm{cos}(s),\mathrm{sin}(s)).`$ The transform is then just $`\underset{0}{\overset{\pi }{}}a_{(0,\mathrm{},0,l)}(\gamma (s))(\mathrm{sin}s)^{r1}𝑑s.`$ We show that this is equal to
$$(1)^l\underset{0}{\overset{\pi }{}}(\mathrm{sin}s)^{l+r1}\mu ,\frac{d\gamma }{ds}(s)𝑑s$$
where $`\mu ,\frac{d\gamma }{ds}(s)`$ is the pairing of $`\mu `$ with $`\frac{d\gamma }{ds}\mathrm{}\frac{d\gamma }{ds}`$ ($`l`$ copies). To prove this, observe that as $`\gamma `$ lies in a plane this is really a two-dimensional question, so without loss of generality its enough to take $`n=2.`$ Now $`\mu `$ is a symmetric $`l`$form so
$$\mu ,v=\underset{j=0}{\overset{l}{}}a_{(j,lj)}v^{(j,lj)}.$$
Thus
$$\mu ,\frac{d\gamma }{ds}(s)=\underset{j=0}{\overset{l}{}}a_{(j,lj)}(\gamma (s))(\mathrm{sin}(s))^{lj}(\mathrm{cos}(s))^j.$$
It follows from aradiality and the fact that the perturbation is of order $`l`$ that
$$a_{(j,lj)}=\left(\genfrac{}{}{0pt}{}{l}{j}\right)\left(\frac{\mathrm{cos}(s)}{\mathrm{sin}(s)}\right)^ja_{(0,l)},$$
and thus we have
$$(\mathrm{sin}(s))^l\mu ,\frac{d\gamma }{ds}(s)=a_{(0,l)}\underset{j=0}{\overset{l}{}}\left(\genfrac{}{}{0pt}{}{l}{j}\right)(\mathrm{cos}(s))^{2j}(\mathrm{sin}(s))^{2l2j},$$
which is of course equal to $`a_{(0,l)}.`$ This establishes the equality.
We now want to show that if $`\mu `$ is a symmetric $`l`$form on the sphere and
(4.6)
$$\underset{0}{\overset{\pi }{}}\mu (\gamma (s)),\frac{d\gamma }{ds}(s)(\mathrm{sin}(s))^{l+r1}𝑑s=0,$$
for all geodesics $`\gamma `$ on the sphere, then $`\mu =0.`$ Since we know this is true for any geodesic $`\gamma ,`$ we have
$$I_{l+r1,\gamma ,\alpha }=\underset{0}{\overset{\pi }{}}\mu (\gamma (s+\alpha )),\frac{d\gamma }{ds}(s+\alpha )(\mathrm{sin}(s))^{l+r1}𝑑s=0$$
for any $`\alpha .`$ Differentiating with respect to $`\alpha `$ (see ), we deduce that $`I_{j,\gamma ,\alpha }=0`$ implies that $`I_{j2,\gamma ,\alpha }=0`$ and thus, differentiating repeatedly, if $`l+r1`$ is even and $`I_{l+r1,\gamma ,\alpha }=0`$ then $`\mu `$ is even. Also, if $`l+r1`$ is odd and $`I_{l+r1,\gamma ,\alpha }=0`$, then $`\mu `$ is odd.
Since we can always reduce by two, it is enough to consider the cases where $`r`$ is $`1`$ or $`2.`$ Considering $`r=1,`$ we have for any geodesic $`\gamma `$ that
$$\underset{0}{\overset{\pi }{}}(\mathrm{sin}(s))^l\mu (\gamma (s)),\frac{d\gamma }{ds}(s)𝑑s=0.$$
If we take a geodesic starting at $`z_n=0`$, then we deduce that
$$\underset{0}{\overset{\pi }{}}(z_n^l\mu )(\gamma (s)),\frac{d\gamma }{ds}(s)𝑑s=0.$$
By rotational invariance we have that
$$\underset{0}{\overset{\pi }{}}(p\mu )(\gamma (s)),\frac{d\gamma }{ds}(s)𝑑s=0,$$
for all homogeneous polynomials, $`p,`$ of order $`l.`$ As $`p\mu `$ is even, we can regard it as a symmetric tensor on projective space which is in the kernel of the generalized x-ray transform on symmetric $`l`$tensors. Fortunately, the kernel of this operator has been identified by Bailey and Eastwood, . They showed that the kernel is precisely the symmetrized covariant derivatives of symmetric $`(l1)`$-tensors. Therefore, to show that $`\mu =0`$ we need to show that if $`p\mu `$ is a symmetrized covariant derivative for all homogeneous polynomials $`p`$ of order $`l`$, then $`\mu =0.`$
We show this by showing that $`\mu `$ has to vanish at the north pole. By rotational invariance it will then follow that $`\mu `$ vanishes everywhere. For the case $`r=2,`$ we apply the same arguments to $`z_n\mu `$ and the result will follow from the case $`r=1.`$
We first prove
###### Lemma 4.1.
Let $`\mu `$ be a symmetric co-tensor of order $`l`$ on $`^{n1}`$ such that $`p\mu `$ is a symmetrized covariant derivative for every homogeneous polynomial of order $`l`$. Then $`\mu `$ vanishes at the origin.
###### Proof.
To prove this lemma, we introduce some new notation. If $`\alpha =(\alpha _1,\mathrm{},\alpha _r)\{1,\mathrm{},n1\}^r`$ then
(4.7)
$$_\alpha =\frac{}{x_{\alpha _1}}\mathrm{}\frac{}{x_{\alpha _r}}.$$
Note that $`_\alpha ^\alpha `$ even when both sides make sense. We also put
(4.8)
$$dx_\alpha =dx_{\alpha _1}\mathrm{}dx_{\alpha _r}.$$
Let $`\stackrel{~}{\alpha }_t=(\alpha _1,\mathrm{},\alpha _{t1},\alpha _{t+1},\mathrm{}\alpha _t).`$
The symmetrized covariant derivative, $`_s\eta ,`$ of a tensor $`\eta `$ is obtained by taking the usual covariant derivative and then averaging over the symmetric group to make it symmetric. On $`^{n1}`$, if the symmetric $`l1`$ tensor $`\eta =\varphi _\gamma dx_\gamma ,`$ and $`_s\eta =\psi _\alpha dx_\alpha ,`$ then
(4.9)
$$\psi _\alpha =\frac{1}{l}\underset{j=1}{\overset{l}{}}_{\alpha _j}\varphi _{\stackrel{~}{\alpha }_j}.$$
Our proof of the lemma follows from the observation that the symmetric $`l1`$ tensors must satisfy certain PDEs. If $`\alpha ,\beta \{1,\mathrm{},n1\}^r,`$ we define an exchange to be a map swapping certain places in $`\alpha `$ with ones in $`\beta .`$ We include the case where no swaps take place. The order of the exchange is the number of swaps. The sign of the exchange will be $`1`$ to the power of the order. If the exchange is $`e,`$ we denote the sign of $`e`$ by $`\mathrm{sgn}(e)`$ and the new value of $`\alpha `$ by $`e(\alpha ,\beta )^{(1)}`$ and of $`\beta `$ by $`e(\alpha ,\beta )^{(2)}.`$ For example, if $`e`$ exchanges $`\alpha _1`$ with $`\beta _1`$ then
$`e(\alpha ,\beta )^{(1)}`$ $`=(\beta _1,\alpha _2,\mathrm{},\alpha _r)`$
$`e(\alpha ,\beta )^{(2)}`$ $`=(\alpha _1,\beta _2,\mathrm{},\beta _r)`$
$`\mathrm{sgn}(e)`$ $`=1.`$
If $`\alpha ,\beta \{1,\mathrm{},n1\}^l,`$ and $`\mu =\psi _\gamma dx_\gamma `$ is a covariant symmetrized derivative then
(4.10)
$$\underset{e}{}\mathrm{sgn}(e)_{e(\alpha ,\beta )^{(2)}}\psi _{e(\alpha ,\beta )^{(1)}}=0,$$
where the sum is taken over all exchanges. To prove this one substitutes the expression for a covariant derivative and observes that the terms with no swaps will cancel with some of the terms from swaps of order one but that the swaps of order one will then have remainder terms which are canceled by swaps of order two and so on.
Suppose we have that every homogeneous polynomial of order $`l`$ times $`\mu `$ is a symmetrized covariant derivative. We show that each term $`\psi _\alpha `$ vanishes at the origin. We consider two cases. The first case is that $`\alpha `$ does not contain all possible values from $`\{1,\mathrm{},n1\},`$ (which will always happen when $`l<n1`$). Let $`r`$ be the value not taken. Let $`\beta =(r,\mathrm{},r).`$ We have that
(4.11)
$$\underset{e}{}\mathrm{sgn}(e)_{e(\alpha ,\beta )^{(2)}}(x_\beta \psi _{e(\alpha ,\beta )^{(1)}})=0$$
where $`x_\beta =x_{\beta _1}x_{\beta _2}\mathrm{}x_{\beta _l}`$. Upon evaluating at $`x=0`$, all terms will vanish except $`_\beta (x_\beta \psi _\alpha )`$, which will equal $`\psi _\alpha `$, which must therefore be zero.
In the second case, $`\alpha `$ takes all values from $`\{1,\mathrm{},n1\}.`$ We let $`\beta _j`$ equal $`2`$ if $`\alpha _j`$ equals $`1`$ and $`1`$ otherwise. As before, (4.11) is satisfied. Upon evaluation at zero, the only terms that will not vanish are those for which $`e(\alpha ,\beta )=(\alpha ,\beta )`$ (up to the order of $`\alpha _j`$ and the order of $`\beta _j`$) and such $`e`$ will have positive sign as an even number of swaps will be involved. We therefore have that a positive multiple of $`\psi _\alpha `$ is zero and thus that $`\psi _\alpha `$ is zero. ∎
This completes the proof for $`^{n1}`$ with the Euclidean metric. We, however, wish to obtain a similar result for the sphere. Projecting the upper hemi-sphere onto the plane, we obtain coordinates on the sphere and the metric agrees with the Euclidean one at the north pole (which corresponds to the origin). The symmetrized covariant derivative will then agree with the Euclidean one up to terms vanishing there. As we are applying the result to polynomials of order $`l`$ these additional terms do not affect the argument and so we conclude that the result will also hold for the sphere.
So we have shown that if the principal symbol vanishes then the lead term of an aradial perturbation vanishes also and thus by induction the asymptotics of aradial perturbations are recoverable from fixed energy scattering data.
## 5. Review of Legendrian Distributions
In this section we review and rephrase the material we need from and . Here $`X`$ is a compact manifold with boundary $`X`$ and $`g`$ is a scattering metric on $`X`$ and we have chosen a product decomposition of the form (1.1). Our account is necessarily brief and we refer the reader to and for more details.
There is a natural bundle over $`X`$ called the scattering cotangent bundle which is denoted $`^{sc}T^{}(X).`$ This is the dual to the bundle of smooth vector fields of bounded length with respect to some (and hence all) scattering metrics on $`X.`$ The restriction of $`^{sc}T^{}(X)`$ to $`X`$ is denoted $`^{sc}T^{}(X)_{|X}`$ and carries a natural contact structure. If $`y`$ are local coordinates on $`X`$ and $`\mu `$ are the corresponding dual coordinates, then $`(y,\mu ,\tau )`$ form local coordinates on $`^{sc}T^{}(X)_{|X},`$ where $`\tau `$ is the coefficient of $`\frac{dx}{x^2}.`$
We recall from Section 2 that a differential operator $`P(x,y,xD_y,x^2D_x)`$ will be in $`\mathrm{\Psi }_{sc}^{m,k}(X,^{sc}\mathrm{\Omega }^{1/2})`$ if it is of order $`m`$ and the total symbol as an operator in $`xD_y,x^2D_x`$ vanishes to $`k^{th}`$ order at the boundary. It then has a well-defined symbol at the boundary
$$j(P)=x^kp_k+x^{k+1}p_{k+1},p_k,p_{k+1}C^{\mathrm{}}(\times T^{}X).$$
###### Definition 5.1.
An intersecting pair with conic points is a subset, $`\stackrel{~}{W},`$ of $`^{sc}T^{}(X)_{|X}`$ which is a union of the closure of a smooth Legendrian submanifold, $`W`$, and $`W^\mathrm{\#},`$ which is a finite union of global sections of the form $`W^\mathrm{\#}(\lambda _j)=\{(y,0,\lambda _j)\},`$ and contains $`\overline{W}W.`$ We also require $`\overline{W}`$ to have an at most conic singularity at $`\mu =0`$; that is, it is smooth if polar coordinates are introduced along $`\overline{W}W.`$
The process of introducing polar coordinates along $`\overline{W}W`$ can be given an invariant meaning and is then called blow-up. We denote the blown-up manifold by $`\widehat{W}.`$
The metric $`g`$ induces a metric $`h`$ on the boundary as nearby it is of the form
$$\tau ^2+h^{}(y,\mu )+xg^{}$$
as a function on $`^{sc}T^{}X,`$ we obtain $`h^{}`$ from $`h`$ via the isomorphism
$$\mu .\frac{dy}{x}\mu .dy.$$
###### Example 5.1.
For each $`y^{}X`$ and $`0\lambda ,`$ let $`G_y^{}(\lambda )`$ be equal to the set of $`(\tau ,y,\mu ),`$ such that $`\tau ^2+|\mu |^2=\lambda ^2,\mu 0,`$ and putting $`\mu =|\mu |\widehat{\mu },`$ such that
(5.4)
$$\begin{array}{c}\tau =|\lambda |\mathrm{cos}(s)\\ |\mu |=|\lambda |\mathrm{sin}(s)\\ (y,\widehat{\mu })=\mathrm{exp}(sH_{\frac{1}{2}h})(y^{},\widehat{\mu }^{})\end{array}$$
where $`s(0,\pi ),`$ $`(y^{},\widehat{\mu }^{})T^{}X,`$ and $`h(y^{},\widehat{\mu }^{})=1.`$ Then $`G_y^{}(\lambda )\{(\lambda ,y,0)\}`$ is an intersecting pair with conic points. We denote this pair $`\stackrel{~}{G}(\lambda ).`$ This is the pair in which we are interested here. The set $`G^{\mathrm{}}(\lambda )=\{(\lambda ,y,0,y^{},0)\}`$ is also important in our construction. Note that $`G^{\mathrm{}}(\lambda )`$ is the initial or incoming surface and that $`G^{\mathrm{}}(\lambda )`$ is the outgoing surface. Note that in the coordinates defined by (5.4), $`G^{\mathrm{}}(\lambda )`$ is $`s=\pi `$ and $`G^{\mathrm{}}(\lambda )`$ is $`s=0,`$ when $`\lambda `$ is positive.
Associated with these intersecting pairs at each conic point is a unique homogeneous Lagrangian submanifold $`\mathrm{\Lambda }(\stackrel{~}{W},\lambda _i)`$ of $`T^{}(X).`$ For the pair $`\stackrel{~}{G}(\lambda )`$ we are interested in, this is precisely the relation of being $`\pi `$ apart along a lifted geodesic. (See Proposition 4 of .) For simplicity, we shall henceforth take $`\lambda `$ to be positive. The $`\lambda `$ negative case is similar, or could be deduced from the positive case.
Melrose and Zworski associated to any such intersecting pair a class of smooth functions whose asymptotics on approach to the boundary are determined by symbols on the Legendrians. A symbol bundle over the smooth Legendrian $`W(\lambda )`$ in the pair $`\stackrel{~}{W}`$ can be defined and is denoted $`\widehat{E}^{m,p}.`$ The sections of this bundle are of the form
$$aS^{pm}|dx|^{mn/4}$$
where $`a`$ is a smooth section of $`C^{\mathrm{}}(\widehat{W};\mathrm{\Omega }_b^{\frac{1}{2}}M_{\widehat{H}})`$, $`S`$ is a defining function of the boundary of $`W,`$ $`M`$ is the Maslov bundle, and $`\mathrm{\Omega }_b^{\frac{1}{2}}`$ is the $`b`$half density bundle. For $`G`$ above, one could take $`S=\mathrm{sin}s.`$ Melrose and Zworski remove this singularity at the endpoints by rescaling but for us it will be easier not to do so.
###### Proposition 5.1.
If $`\stackrel{~}{W}(\lambda )`$ is an intersecting pair with conic points then there is a class of smooth half-densities on $`X^o`$, denoted $`I_{sc}^{m,p}(X,\stackrel{~}{W}),`$ such that $`\underset{m,p}{}I_{sc}^{m,p}(X,\stackrel{~}{W})`$ is equal to the class of half-densities vanishing to infinite order at the boundary. There exists a symbol map
$$\widehat{\sigma }_{sc,m,p}:I_{sc}^{m,p}(X,\stackrel{~}{W},^{sc}\mathrm{\Omega }^{1/2})C^{\mathrm{}}(\widehat{W};\widehat{E}^{m,p})$$
which gives a short exact sequence
$$0I_{sc}^{m+1,p}(X,\stackrel{~}{W},^{sc}\mathrm{\Omega }^{1/2})I_{sc}^{m,p}(X,\stackrel{~}{W},^{sc}\mathrm{\Omega }^{1/2})C^{\mathrm{}}(\widehat{W};\widehat{E}^{m,p})0.$$
This is Proposition 12 from . The Legendrian half-densities of order $`m`$ are given locally away from the conic points by oscillatory integrals
(5.5)
$$u=(2\pi )^{\frac{n}{4}\frac{k}{2}}e^{i\varphi (y,u)}a(x,y,u)x^{mk/2+n/4}𝑑u,$$
with $`a`$ smooth on $`[0,ϵ)\times U\times U^{}`$ with $`U,U^{}`$ open and $`\varphi `$ parameterizes $`W,`$ that is
(5.6)
$$W=\{(y,\varphi (y,u),d_y\varphi ):d_u\varphi =0\}.$$
Near the conic singularity a more general form is required and we refer the reader to . The order $`m`$ here is adjusted by $`k/2+n/4`$ and so is inconsistent with section 3. We keep this inconsistency as the second definition provides for better invariance properties but less clarity.
An important related fact we need to know is how Legendrian distributions map under scattering pseudo-differential operators. We recall Proposition 13 from .
###### Proposition 5.2.
Suppose $`P\mathrm{\Psi }_{sc}^{l,k}(X,^{sc}\mathrm{\Omega }^{1/2})`$ has symbol $`x^kp_k+x^{k+1}p_{k+1}`$ with respect to a product decomposition of $`X`$ near $`X,`$ and suppose that
$$W^{sc}T_X^{}(X)$$
is a smooth Legendre submanifold. Then for any $`m,`$
(5.7)
$$P:I_{sc}^m(X,W;^{sc}\mathrm{\Omega }^{1/2})I_{sc}^{m+k}(X,W;^{sc}\mathrm{\Omega }^{1/2})$$
(5.8)
$$\sigma _{sc,m+k}(Pu)=(p_{k}^{}{}_{|G}{}^{})\sigma _{sc,m}(u)|dx|^k.$$
Furthermore, if $`p_k`$ vanishes identically on $`W`$ then
$$P:I_{sc}^m(X,W;^{sc}\mathrm{\Omega }^{1/2})I_{sc}^{m+k+1}(X,W;^{sc}\mathrm{\Omega }^{1/2})$$
and
$$\sigma _{sc,m+k+1}(Pu)=$$
$$\left(\frac{1}{i}\left(L_V+\left(\frac{1}{2}(k+1)+m\frac{n}{4}\right)\frac{p_k}{\tau }\right)+p_{k+1}^{}{}_{|W}{}^{}\right)a|dx|^{m+k+1\frac{n}{4}}$$
where $`\sigma _{sc,m}(u)=a|dx|^{m\frac{n}{4}}`$ and $`V`$ is the rescaled Hamiltonian vector field associated to $`p_k.`$
We omit the definition of the rescaled Hamiltonian vector field but recall that for $`\mathrm{\Delta }`$ on the pair $`G`$ we are studying it is equal to
$$2\lambda \mathrm{sin}s\frac{}{s}$$
in the semi-global coordinates given by (5.4).
We also need two push-forward theorems, Propositions 16 and 17, from . They relate the singularities of the scattering matrix to the asymptotics in small $`x`$ of the Poisson operator. Given a product decomposition near the boundary, there is a natural pairing
(5.9)
$$B:C^{\mathrm{}}(X,^{sc}\mathrm{\Omega }^{1/2})\times C^{\mathrm{}}(X;^{sc}\mathrm{\Omega }^{1/2})C^{\mathrm{}}([0,ϵ),^{sc}\mathrm{\Omega }^{1/2})$$
(5.10)
$$B(u,f)=x^{\frac{n1}{2}}_Xu(x,y)f(y).$$
###### Proposition 5.3.
For any intersecting pair of Legendre submanifolds with conic points, $`W,`$ the partial pairing (5.10) gives a map
$$B:I_{sc}^{m,p}(X,\stackrel{~}{W};^{sc}\mathrm{\Omega }^{1/2})\times C^{\mathrm{}}(X;^{sc}\mathrm{\Omega }^{1/2})\underset{j}{}I^{p+\frac{n1}{4}}([0,ϵ),W^{}(\overline{\tau }_j;^{sc}\mathrm{\Omega }^{1/2}))$$
where the $`W^{}(\overline{\tau }_j)=\{(0,\tau _jdx/x^2)\}`$ are the Legendre submanifolds corresponding to the components of $`W^\mathrm{\#}`$ and
$$B(u,f)=\underset{j}{}e^{i\overline{\tau }/x}x^{p+n/4}Q_{\overline{\tau }_j}^0(u,f)\left|\frac{dx}{x^2}\right|^{\frac{1}{2}}+O(x^{p+n/4+1})$$
with
$$Q_{\overline{\tau }}^0(u)I_{phg}^{pm\frac{n1}{4}}(X,\mathrm{\Lambda }(\stackrel{~}{W},\overline{\tau })),$$
and the principal symbol of $`Q_{\overline{\tau }}^0(u)`$ determines and is determined by the lead singularity of the principal symbol of $`u`$ on $`W`$ on approach to $`W^{}(\overline{\tau }_j).`$
When the Legendrian distribution is actually associated to a smooth Legendrian submanifold the push-forward becomes much simpler and this simplifies the construction of the Poisson operator.
###### Proposition 5.4.
If $`G`$ is a smooth Legendre variety and $`uI_{sc}^m(X,G^{},^{sc}\mathrm{\Omega }^{1/2})`$ near $`\tau =\overline{\tau },`$ then the distribution $`Q_{\overline{\tau }}^0`$ is a Dirac delta distribution.
## 6. The Poisson operator in the general case
In this section, we apply the calculus reviewed in Section 5 to construct the Poisson operator and prove that the higher order scattering matrix is indeed a zeroth order, classical, Fourier integral operator, proving Theorem 1.2 and Corollary 1.1. We shall refer heavily to , Section 15 as our construction is a modification of the one there.
We assume a product decomposition of $`X`$ close to the boundary of the form (1.1) has been chosen and is fixed throughout this section. We then have as in that $`\mathrm{\Delta }`$, the intrinsic Laplacian acting on scattering half-densities on $`X`$, induces an operator
$$\mathrm{\Delta }_X\mathrm{Diff}_{\mathrm{sc}}^2(X\times X,^{sc}\mathrm{\Omega }^{1/2}(X\times X))$$
by
$$\mathrm{\Delta }_X\left(u\right|\frac{dx}{x^2}|^{\frac{1}{2}}\left|\frac{dy}{x^{n1}}|^{\frac{1}{2}}\right|\frac{dy^{}}{x^{n1}}|^{\frac{1}{2}})=\mathrm{\Delta }\left(u(.,y^{})\right|\frac{dx}{x^2}|^{\frac{1}{2}}\left|\frac{dy}{x^{n1}}|^{\frac{1}{2}}\right)|\frac{dy^{}}{x^{n1}}|^{\frac{1}{2}},$$
where $`(x,y,y^{})`$ is a point in $`X\times X.`$ Throughout this section $`V`$ will be a short range higher order perturbation, that is a scattering differential operator of order $`2k1`$ as a differential operator and order $`2`$ at the boundary.
We recall from , using the notation of (1.1) that
###### Lemma 6.1.
The symbol at the boundary of $`\mathrm{\Delta }_X`$ is $`p=p_0+xp_1`$ with $`p_0=\tau ^2+h(0,y,\mu )`$ and with $`p_1`$ equal to $`i(n1)\tau +c`$ where $`c=\frac{}{x}h(x,y,\mu )_{|x=0}`$ is quadratic in $`\mu `$.
Our choice of product decomposition ensures that there is no $`\tau `$ term in $`c.`$ We need to compute the symbols at the boundary of $`\mathrm{\Delta }_X^k\lambda ^{2k}+V.`$ The short range assumption on $`V`$ ensures that it does not contribute - indeed this is one reason why this definition of short range is appropriate. Now we can decompose
$$(\mathrm{\Delta }^k\lambda ^{2k})=Q(\lambda )(\mathrm{\Delta }_X\lambda ^2),$$
with
$$Q(\lambda )=\underset{j=0}{\overset{k1}{}}\lambda ^{2j}\mathrm{\Delta }_X^{k1j}.$$
We deduce that the lead symbol of $`(\mathrm{\Delta }_X^k\lambda ^{2k})`$ is $`p_0^k\lambda ^{2k}`$ and the second symbol (subprincipal term) is equal to $`k\lambda ^{2(k1)}`$ times $`p_1`$ on the zero set of the principal symbol. It also follows that the Hamiltonian on the boundary of the lead term is just $`k\lambda ^{2(k1)}`$ times that of $`p_0`$ on the zero set of the principal symbol. The fact that both these terms have been changed simply by multiplication by $`k\lambda ^{2(k1)}`$ ensures the simplicity of extending results from the $`k=1`$ case to the higher order case.
We also note the following lemma which is important in our construction to show that the transport equations are solvable.
###### Lemma 6.2.
If $`LI^{m,\frac{1}{4}}(X\times X,\stackrel{~}{G}(\lambda ),^{sc}\mathrm{\Omega }^{1/2})`$ is such that
$$(\mathrm{\Delta }_X^k\lambda ^{2k}+V)LI^{m+j,\frac{3}{4}}(X\times X,\stackrel{~}{G}(\lambda ),^{sc}\mathrm{\Omega }^{1/2}),$$
then
$$(\mathrm{\Delta }_X^k\lambda ^{2k}+V)LI^{m+j,\frac{7}{4}}(X\times X,\stackrel{~}{G}(\lambda ),^{sc}\mathrm{\Omega }^{1/2}).$$
This is a modification of Lemma 15 from and in fact, the $`k=1`$ case, though not explicitly stated, is essential to the construction there also.
###### Proof.
The proof is no different from that of the special case in $`^n`$ (Lemma 3.1). ∎
###### Proposition 6.1.
For any $`0\lambda `$ there exists
$$KI^{m,p}(X\times X,\stackrel{~}{G}(\lambda );^{sc}\mathrm{\Omega }^{1/2})$$
such that
$$(\mathrm{\Delta }_X^k\lambda ^{2k}+V)K𝒞^{\mathrm{}}(X\times X;^{sc}\mathrm{\Omega }^{1/2})\text{ and }$$
$$Q_\lambda ^0(K)=\mathrm{Id},$$
with
$$m=\frac{2n1}{4},p=\frac{1}{4},$$
and the principal symbol of $`K`$ on $`G`$ is
$$C\mathrm{sin}(s)^{\frac{n1}{2}}\frac{|ds|^{\frac{1}{2}}|dy|^{\frac{1}{2}}|d\widehat{\mu }|^{\frac{1}{2}}}{(\mathrm{sin}s)^{\frac{1}{2}}}|dx|^{m\frac{2n1}{4}},$$
where $`C(y,\widehat{\mu })`$ is a non-zero smooth function.
###### Proof.
As in , we first construct $`K^bI^{m,p}(X\times X,\stackrel{~}{G}(\lambda );^{sc}\mathrm{\Omega }^{1/2})`$ such that
(6.1)
$$(\mathrm{\Delta }_X^k\lambda ^{2k}+V)K^bI_{sc}^{\frac{3}{4}}(X\times X,G^{\mathrm{}}(\lambda ))\text{ and }$$
(6.2)
$$Q_\lambda ^0(K^b)=\mathrm{Id}.$$
We construct $`K^b`$ as an asymptotic sum of
$$K_jI_{sc}^{\frac{2n1}{4}+j,\frac{1}{4}}(X\times X,\stackrel{~}{G}(\lambda );^{sc}\mathrm{\Omega }^{1/2}).$$
We wish
(6.3)
$$(\mathrm{\Delta }_X^k\lambda ^{2k}+V)K_0I_{sc}^{\frac{2n1}{4}+2,\frac{3}{4}}(X\times X,\stackrel{~}{G}(\lambda );^{sc}\mathrm{\Omega }^{1/2})\text{ and }$$
(6.4)
$$\sigma _0(Q_\lambda ^0(K_0))=\sigma _0(\mathrm{Id})$$
and then it is automatically in
$$I_{sc}^{\frac{2n1}{4}+2,\frac{7}{4}}(X\times X,\stackrel{~}{G}(\lambda );^{sc}\mathrm{\Omega }^{1/2})$$
by Lemma 6.2. We also want
$$(\mathrm{\Delta }_X^k\lambda ^{2k}+V)\left(\underset{l=0}{\overset{j1}{}}K_l\right)I_{sc}^{\frac{2n1}{4}+j+2,\frac{3}{4}}(X\times X,\stackrel{~}{G}(\lambda );^{sc}\mathrm{\Omega }^{1/2})$$
and this, of course, implies that it will also be an element of
$$I_{sc}^{\frac{2n1}{4}+j+2,\frac{7}{4}}(X\times X,\stackrel{~}{G}(\lambda );^{sc}\mathrm{\Omega }^{1/2}).$$
Now near $`GG^{\mathrm{}}(\lambda ),`$ where $`G`$ is smooth, we can as in give an explicit construction and it is then only necessary to have that the principal symbol of $`Q_\lambda ^0(K_0)`$ is equal to $`1`$ to ensure that $`Q_\lambda ^0(K^b)=\mathrm{Id}.`$ We look for $`K_j`$ of the form
$$x^je^{i\lambda \varphi (y,y^{})/x}a_j(x,y,y^{},\lambda )v,a_jC^{\mathrm{}}(X\times X),$$
with $`v`$ a fixed scattering half-density and $`\varphi `$ the cosine of the Riemannian distance from $`y`$ to $`y^{}.`$ Let $`a_j^{^{}}`$ be the restriction of $`a_j`$ to $`x=0.`$ Taking geodesic normal coordinates, $`y,`$ about each $`y^{}`$ the transport equations for $`a_j^{^{}}`$ is of the form
$$(y_y+j)a_j^{^{}}+b_ja_j^{^{}}=c_jC^{\mathrm{}}(X\times X)$$
near $`y=0`$ where $`c_0`$ is identically zero and $`b_j`$ vanishes quadratically at $`y=0.`$ So as in , the terms $`K_j`$ exist sc-microlocally close to $`G^{\mathrm{}}(\lambda ).`$
We now need to continue each $`K_j`$ up to $`G^{\mathrm{}}(\lambda ).`$ We do so by solving transport equations for the principal symbols and iteratively solving away the error.
The principal symbol of $`K_0,`$ $`\sigma _m(K_0),`$ is of the form
$$b\frac{|ds|^{\frac{1}{2}}|dy|^{\frac{1}{2}}|d\widehat{\mu }|^{\frac{1}{2}}}{(\mathrm{sin}s)^{\frac{1}{2}}}|dx|^{m\frac{2n1}{4}}.$$
On the lifted geodesic $`\beta (s)`$ the sub-principal term $`c(\beta (s))=2k\lambda ^{2k1}\mathrm{sin}(s)d(\beta (s))`$ for some smooth $`d.`$ From Proposition 5.2, the transport equation for $`b`$ is
$$\frac{2k\lambda ^{2k1}}{i}\left(\mathrm{sin}(s)\frac{d}{ds}+\left(\frac{1n}{2}\right)\mathrm{cos}(s)+i\mathrm{sin}(s)d(\beta (s))\right)b=0$$
Writing $`\stackrel{~}{b}=(e^{i{\scriptscriptstyle d(\beta (s^{}))𝑑s^{}}}\mathrm{sin}(s)^{\frac{1n}{2}})b,`$ we thus have
$$\frac{d\stackrel{~}{b}}{ds}=0.$$
This means that
(6.5)
$$b=C\mathrm{sin}(s)^{\frac{n1}{2}}e^{i{\scriptscriptstyle d(\beta (s^{}))𝑑s^{}}}.$$
As $`s\pi ,`$ that is near $`G^{\mathrm{}}(\lambda ),`$ this has a singularity of the form $`(\pi s)^{\frac{n1}{2}}`$ and as $`s0+,`$ of the form $`s^{\frac{n1}{2}}.`$
As the order on $`G^{\mathrm{}}(\pm \lambda )`$ is $`\frac{1}{4},`$ the order $`pm`$ is equal to the order of singularity and thus the solution of the transport equation is a legitimate symbol and we can construct $`K_0.`$
Now by Lemma 6.2,
$$(\mathrm{\Delta }_X^k+V\lambda ^{2k})(K_0)I^{\frac{2n1}{4}+2,\frac{1}{4}+2}(X\times X,\stackrel{~}{G};^{sc}\mathrm{\Omega }^{1/2}),$$
and we look for
$$K_1I^{\frac{2n1}{4}+1,\frac{1}{4}}(X\times X,\stackrel{~}{G};^{sc}\mathrm{\Omega }^{1/2}),$$
such that
$$(\mathrm{\Delta }_X^k+V\lambda ^{2k})(K_0+K_1)I^{\frac{2n1}{4}+3,\frac{1}{4}+1}(X\times X,\stackrel{~}{G};^{sc}\mathrm{\Omega }^{1/2}),$$
and thus by Lemma 6.2 is in $`I^{\frac{2n1}{4}+3,\frac{1}{4}+2}(X\times X,\stackrel{~}{G};^{sc}\mathrm{\Omega }^{1/2}).`$ Letting the principal symbol of $`K_1`$ be $`b_1|dx|^{\frac{2n1}{4}+1\frac{2n1}{4}}`$ times the trivializing density above, we obtain a transport equation; arguing as above, it becomes
$$\begin{array}{c}\mathrm{sin}(s)^{\frac{n1}{2}}e^{i{\scriptscriptstyle d(\beta (s^{}))𝑑s^{}}}\frac{d}{ds}e^{i{\scriptscriptstyle d(\beta (s^{}))𝑑s^{}}}\left(\mathrm{sin}(s)^{\frac{1n}{2}+1}b_1\right)\hfill \\ \hfill =g(s)e^{i{\scriptscriptstyle d(\beta (s^{}))𝑑s^{}}}\mathrm{sin}(s)^{\frac{n1}{2}},\end{array}$$
with $`g(s)`$ a smooth function on $`[0,\pi ]`$ (and smoothly depending on the suppressed parameters). Canceling, we obtain that
$$\frac{d}{ds}\left(e^{i{\scriptscriptstyle }d(\beta (s^{})ds^{}}\mathrm{sin}(s)^{\frac{1n}{2}+1}b_1\right)=g(s),$$
which has a solution in the appropriate symbol class. The same argument, after appropriately shifting indices, constructs all the terms $`K_j.`$
Asymptotically summing, we obtain $`K^b`$ such that
$$(\mathrm{\Delta }_X^k+V\lambda ^{2k})K^bI^{7/4}(G^{\mathrm{}}(\lambda )).$$
These errors can now be removed by an iterative construction of their Taylor series, cf Lemma 16 of , and we obtain $`K`$ as desired. ∎
We have therefore constructed the Poisson operator modulo smooth terms vanishing to infinite order at the boundary. We wish to remove this error by applying the resolvent. As before, this is possible even if there is embedded discrete spectrum, as the error is orthogonal to the eigenspace at that energy. This is seen by a simple integration by parts, since the elements of the eigenspace are Schwartz they are orthogonal to the image of smooth functions of tempered growth. We have thus constructed $`P(\lambda )`$, the Poisson operator for the problem, as a paired Legendrian distribution. The remainder of the proof that $`S(\lambda )`$ is a Fourier integral operator now follows as in , Proposition 19.
To prove Corollary 1.1, one observes that if $`P_1(\lambda )`$ is the Poisson operator for $`\mathrm{\Delta }\lambda ^2,`$ then
(6.6)
$$(\mathrm{\Delta }^k+V\lambda ^{2k})P_1(\lambda )=Q(\lambda )(\mathrm{\Delta }\lambda ^2)P_1(\lambda )+VP_1(\lambda )=VP_1(\lambda ).$$
This means that the Poisson operator for $`\mathrm{\Delta }^k+V\lambda ^{2k}`$ can be constructed as a perturbation of $`P_1(\lambda )`$ and as $`VP_1(\lambda )I^{\frac{2n1}{4}+l,\frac{1}{4}+l},`$ the scattering matrices will agree to order $`1l,`$ and the corollary follows. Note, however, that even if one fixes the perturbation, the scattering matrices will not agree to order more than $`1l`$ as the solutions of the transport equations will differ.
Department of Mathematics, University of Missouri, Columbia, Missouri 65211 U.S.A.
E-mail address: tjc@@math.missouri.edu
Darwin College, Cambridge CB3 9EU U.K.
Current Address: NatWest Group Risk, Level 9, 135 Bishopsgate, London EC2M 3UR U.K.
E-mail address: markjoshi@@alum.mit.edu |
warning/0002/gr-qc0002086.html | ar5iv | text | # References
Macroscopic Einstein - Maxwell equations for a system of interacting particles to second-order accuracy in the interaction constan.
A. V. Zakharov
Dept.Gen.Rel.& Grav.
Kazan State University,
Kremljevskaja Str. 18
420008 Kazan, Russia
e-mail: Alexei.Zakharov@ksu.ru
Macroscopic Einstein - Maxwell equations for a system of interacting particles to second-order accuracy in the interaction constant.
A. V. Zakharov
Abstract.
In this paper the macroscopic Einstein and Maxwell equations for system, in which the electromagnetic interactions are dominating (for instance, the cosmological plasma before the moment of recombination), are derived.
Ensemble averaging of the microscopic Einstein - Maxwell equations and the Liouville equations for the random functions leads to a closed system of macroscopic Einstein - Maxwell equations and kinetic equations for one-particle distribution functions. The macroscopic Einstein equations for a relativistic plasma differ from the classical Einstein equations in that their left-hand sides contain additional terms due to particle interaction. The terms are traceless tensors with zero divergence. An explicit covariant expression for these terms is given in the form of momentum-space integrals of expressions depending on one-particles distribution functions of the interacting particles of the medium.
The additional terms in the left-hand side of the macroscopic Einstein equations for a relativistic plasma has murch in common with the additional terms in the left-hand side of the macroscopic Einstein equations for a system of self - gravitating particles (refer to. , ).
The macroscopic Maxwell equations alsow differ from the classical macroscopic Maxwell equations in that their left-hand sides contain additional terms due to particle interaction as well the effects of general relativity.
Macroscopic Einstein - Maxwell equations for a system of interacting particles to second-order accuracy in the interaction constant.
A. V. Zakharov
1. Introduction.
The idea of macroscopic gravity can be considered as an extension of Lorentz’ idea (refer to. ), formulated first for electrodynamics, about the existence of two levels, microscopic and macroscopic, of understanding classical physical phenomena. Lorentz showed that Maxwell’s electrodynamics is a macroscopic theory of electromagnetism, and the Maxwell equations could be derived from a system of microscopic field equations called now the Maxwell - Lorentz ones, by infinitisimal space - time regions averaging (refer to. , ).
As it is known the macroscopic Maxwell equation for continuous media can be alsow obtained from the microscopic Maxwell equations by ensemble averaging the latter (refer to.).
The Einstein equations, whose right-hand side contain the energy-momentum tensor of matter, are phenomenological equations. It is natural to suppose that the Einstein equations (or their generalizations) for continuous media can also be obtained from the microscopic Einstein equations, i.e., Einstein equations whose right-hand sides contain the sum of the energy - momentum tensors of individual particles. However, due to the nonlinearity of the left-hand side of Einstein equations, the averaging of the microscopic Einstein equations is more complicated than one of the microscopic Maxwell equations (refer to. - ).
In (refer to.) a method is developed for the ensemble averaging of the microscopic Einstein equations for interacting particles. ( We use the ensemble averaging procedure introduced by Klimontovich (refer to. ) to derive the relativistic kinetic equation for a plasma. The same procedure was used by the present author in (refer to. ) to derive a relativistic kinetic equation for a system of gravitationally and electromagnetical interacting particles in General Relativity accurate to within the second order for the interaction constant.) This results to macroscopic Einstein equations for continuous media that are accurate to second-order terms in the interaction constant. The macroscopic Einstein equations for a system of interacting particles differ from the classical Einstein equations in that their left-hand sides contain additional components due to particle interaction. The components are expressed in terms of the two-particle correlation function of the particles.
In (refer to. ) we got covariant expressions for additional components for the system of self-gravitating particles. The components are traceless tensors with zero divergence. The expressions were obtained in the form of momentum-space integrals of expressions depending on one-particle distribution function of the gravitationally interacting particles of the medium. The given expressions are proportional to the cube of the Einstein constant and the square of the particle number density. The latter relationship implies that interaction effects manifest themselves in systems of very high density (the Universe in the early stages of its evolution, dense objects close to gravitational collapse, etc.)
The present paper is a direct continuation of earlier papers (refer to., ), devoted to the derivation of the macroscopic Einstein equations for a system of self-gravitating particles to within terms of second order in the interaction constant.
The objective of this paper is to obtain the macroscopic Einstein equations for system, in which electromagnetic interactions (for instance, cosmological plasma before a moment of recombination,) are dominating.
The macroscopic Einstein equations for relativistic plasma differ from the classical Einstein equations in that their left-hand side contains additional terms due to particle interaction. The terms are traceless tensors with zero divergence. An explicit covariant expression for these terms is given in the form of momentum-space integrals of expressions depending on one-particles distribution functions of the interacting particles of the medium.
The additional terms in the left-hand side of the macroscopic Einstein equations for a relativistic plasma has march in common with the additional terms in the left-hand side of the macroscopic Einstein equations for a system of self - gravitating particles (refer to. , ).
The macroscopic Maxwell equations alsow differ from the classical macroscopic Maxwell equations for their left-hand sides contain additional terms due to particle interaction as well the effects of general relativity.
2. Microscopic Einstein - Maxwell equations
The method of deriving the macroscopic Einstein equations is discussed in (refer to.). The notation we use here are the same that in (refer to.).
Briefly, the method we used to obtain the macroscopic Einstein - Maxwell equations is the following.
We start from the microscopic Einstein and Maxwell equations
$$\stackrel{~}{G}^{ij}=\chi \stackrel{~}{T}_{(m)}^{ij}+\chi \stackrel{~}{T}_{(el)}^{ij},$$
(1)
$$\stackrel{~}{}_k\stackrel{~}{F}^{ik}=\frac{4\pi }{c}\stackrel{~}{J}^i.$$
(2)
Here $`\stackrel{~}{G}^{ij}`$ is the Einstein tensor in a Riemannian space with metric $`\stackrel{~}{g}_{ij}`$ , $`\chi =8\pi k/c^4`$ is Einstein’s constant (where $`k`$ is the gravitational constant, $`c`$ is the velocity of light), $`\stackrel{~}{T}^{ij}`$ is the microscopic energy-momentum tensor of particles, $`\stackrel{~}{F}^{ik}`$ is the Maxwell’s tensor, $`\stackrel{~}{J}^i`$ is the microscopic current vector of particles, $`\stackrel{~}{T}_{(el)}^{ij}`$ is the energy - momentum tensor of electromagnetical field. Raising and lowering of indexes is accomplishment with the metric $`\stackrel{~}{g}_{ij}`$, $`\stackrel{~}{}_k`$ is a covariant derivative in a Riemannian space with metric $`\stackrel{~}{g}_{ij}`$.
The tensor $`\stackrel{~}{T}_{(el)}^{ij}`$ have the form
$$\stackrel{~}{T}_{(el)}^{ij}=\frac{1}{4\pi }\left(\stackrel{~}{F}_{.l}^i\stackrel{~}{F}^{lj}+\frac{1}{4}\stackrel{~}{g}^{ij}\stackrel{~}{F}_{lm}\stackrel{~}{F}^{lm}\right).$$
(3)
The tensors $`\stackrel{~}{T}_{(m)}^{ij}`$ and $`\stackrel{~}{J}^i`$ has the form
$$\stackrel{~}{T}_{(m)}^{ij}=\underset{a}{}m_ac^2\frac{d^4\stackrel{~}{p}_a}{\sqrt{\stackrel{~}{g}}}\stackrel{~}{u}_a^i\stackrel{~}{u}_a^j\stackrel{~}{N}_a(q^i,\stackrel{~}{p}_i),$$
(4)
$$\stackrel{~}{J}^i=\underset{b}{}e_bc\frac{d^4\stackrel{~}{p}}{\sqrt{\stackrel{~}{g}}}\stackrel{~}{u}_b^i\stackrel{~}{N}_b(q,\stackrel{~}{p}).$$
(5)
Here $`e_a`$ is the charge of particles of species ”a”, $`m_a`$ is their mass, $`\stackrel{~}{g}`$ is the determinant of $`\stackrel{~}{g}_{ij}`$, $`\stackrel{~}{p}_a^i`$ is the momentum of particles of spesies ”a”, $`\stackrel{~}{u}_a^i=\stackrel{~}{p}_a^i/\sqrt{\stackrel{~}{g}_{kj}\stackrel{~}{p}_a^k\stackrel{~}{p}_a^j}`$,
$$\frac{d^4\stackrel{~}{p}}{\sqrt{\stackrel{~}{g}}}$$
\- is the invariant volume element in momentum spase .
$`N_a(q^i,\stackrel{~}{p}_a^i)`$ \- is the Klimontovich random function :
$$\stackrel{~}{N}_a(q^i,\stackrel{~}{p}_j)=\underset{i=1}{\overset{n_a}{}}𝑑\stackrel{~}{s}\delta ^4(q^iq_{(l)}^i)\delta ^4(\stackrel{~}{p}_j\stackrel{~}{p}_j^{(l)}(\stackrel{~}{s})).$$
(6)
Nere $`n_a`$ is the number of particles belonging to species ”a”, $`\stackrel{~}{s}`$ is the canonical parameter along the particle trajectories: $`d\stackrel{~}{s}=\sqrt{g_{ij}dq^idq^j}`$; $`q_{(l)}^i`$ and $`\stackrel{~}{p}_j^{(l)}`$ are the coordinates and momentum of the l-th particle of spesies ”a”, which are found by solving the equations of motion:
$$\frac{dq_{(l)}^i}{d\stackrel{~}{s}}=\frac{\stackrel{~}{p}_{(l)}^i}{m_ac},\frac{d\stackrel{~}{p}_i^{(l)}}{d\stackrel{~}{s}}=\frac{1}{m_ac}\stackrel{~}{\mathrm{\Gamma }}_{j,ik}\stackrel{~}{p}_{(l)}^j\stackrel{~}{p}_{(l)}^k+\frac{e_a}{c}\stackrel{~}{F}_{ik}\stackrel{~}{p}_{(l)}^k.$$
(7)
Here $`\stackrel{~}{\mathrm{\Gamma }}_{j,ik}`$ is the Christoffel symbol of the first kind given by the metric $`\stackrel{~}{g}_{ij}`$.
In view of Eqs. (7) the random function (6) obeys the equation
$$\stackrel{~}{p}^i\frac{\stackrel{~}{N}_a}{q^i}+\stackrel{~}{\mathrm{\Gamma }}_{j,ik}\stackrel{~}{p}^j\stackrel{~}{p}^k\frac{\stackrel{~}{N}_a}{\stackrel{~}{p}_i}+\frac{e_a}{c}\stackrel{~}{F}_{ik}\stackrel{~}{p}^k\frac{\stackrel{~}{N}_a}{\stackrel{~}{p}_i}=0.$$
(8)
Next we write the metric $`\stackrel{~}{g}_{ij}`$ as
$$\stackrel{~}{g}_{ij}=g_{ij}+h_{ij},$$
(9)
and $`\stackrel{~}{F}_{ik}`$ as
$$\stackrel{~}{F}_{ik}=F_{ik}+\omega _{ik},$$
(10)
Here $`g_{ij}=\stackrel{~}{g}_{ij}`$ is the ensemble averege of the metric $`\stackrel{~}{g}_{ij}`$ , $`\stackrel{~}{F}_{ik}`$ is the ensemble average of $`\stackrel{~}{F}_{ik}`$. Note that $`h_{ij}0`$ and $`\omega _{ik}0`$.
Parallel with the momenta $`\stackrel{~}{p}_{(l)}^i=m_acdq_{(l)}^i/d\stackrel{~}{s}`$ we use the momenta $`p^i`$ measured in the metric $`g_{ij}`$:
$$p_{(l)}^i=\alpha ^1(q,p)\stackrel{~}{p}_{(l)}^i,\alpha (q,p)=ds/d\stackrel{~}{s}=(g_{ij}p^ip^j)^{1/2}(\stackrel{~}{g}_{lk}\stackrel{~}{p}^l\stackrel{~}{p}^k)^{1/2}.$$
(11)
Here $`s`$ is the canonical parametr introduced by $`g_{ij}`$.
The transformation from $`\stackrel{~}{p}_i`$ to $`p_i`$ is given by
$$\stackrel{~}{p}_j=\stackrel{~}{g}_{jk}\stackrel{~}{p}^k=\alpha \stackrel{~}{g}_{jk}g^{ki}p_i.$$
(12)
THe Jacobian of transformation (12), is (see ):
$$\frac{\stackrel{~}{p}_i}{p_j}=\alpha ^4\frac{\stackrel{~}{g}}{g},$$
(13)
where $`g`$ is the determinant of $`g_{ij}`$.
Now we introduce the function $`N_a(q^i,\stackrel{~}{p}_j)`$ defined in the eight - dimensional phase space with coordinates $`(q,p)`$ as
$$N_a(q,p)=\underset{l=1}{\overset{n_a}{}}𝑑s\delta ^4(q^iq_{(l)}^i(s))\delta ^4(p_jp_j^{(l)}(s)),$$
(14)
where $`q_{(l)}^i`$ $`p_j^{(l)}`$ are found by solving equations obtained from (7) with the transformation (12) taken into account ($`p^i=g^{ij}p_j`$).
Note that the functions $`\stackrel{~}{N}_a`$ and $`N_a`$ are related in the folloing manner:
$$\stackrel{~}{N}_a(q,\stackrel{~}{p})=\frac{g}{\stackrel{~}{g}\alpha ^5}N_a(q,p).$$
(15)
Equation for $`N_a(q,p)`$ can be obtained directly from (8) by replaicing the variables (12) and (15):
$$p^i\frac{N_a}{q^i}+\mathrm{\Gamma }_{j,ik}p^jp^k\frac{N_a}{p_i}+\frac{e_a}{c}F_{ik}p^k\frac{N_a}{p_j}=$$
$$=\frac{}{p_i}\left[\left(\mathrm{\Omega }_{jk}^m\mathrm{\Delta }_{mi}p^jp^k\frac{e_a}{c}\psi _{.k}^l\mathrm{\Delta }_{lj}p^k\right)N_a\right].$$
(16)
Here
$$\mathrm{\Delta }_{ki}=g_{ki}u_ku_i;u^k=\frac{p^k}{\sqrt{p^lp_l}};\mathrm{\Omega }_{kj}^m=\stackrel{~}{\mathrm{\Gamma }}_{kj}^m\mathrm{\Gamma }_{kj}^m$$
(17)
\- is the difference of the Christoffel symbols of second kind for the metrics $`\stackrel{~}{g}_{ij}`$ and $`g_{ij}`$,
$$\psi _{.k}^l=\frac{1}{\alpha (q,p)}\stackrel{~}{F}_{.k}^lF_{.k}^l=\frac{1}{\alpha (q,p)}\stackrel{~}{g}^{lm}\stackrel{~}{F}_{mk}g^{lm}F_{mk}.$$
(18)
If in (4) and (5) we turn to the variables $`p_i`$ and $`N_a`$ we get
$$\stackrel{~}{T}^{ij}=\underset{a}{}m_a^2\frac{d^4p_a}{\sqrt{(g)}}\alpha (q,p)\sqrt{\frac{g}{\stackrel{~}{g}}}u_a^iu_a^jN_a(q,p_a),$$
(19)
$$\stackrel{~}{J}^i=\underset{b}{}e_bc\frac{d^4p_b}{\sqrt{g}}\sqrt{\frac{g}{\stackrel{~}{g}}}u_b^iN_b(q,p_b).$$
(20)
where $`d^4p/\sqrt{g}`$ is the invariant volume element in the unperturbed momentum spase.
For subsequent calculation is it convinient to write the Einstein equations as
$$R_{ij}+_m\mathrm{\Omega }_{ij}^m_j\mathrm{\Omega }_{im}^m+\mathrm{\Omega }_{mn}^m\mathrm{\Omega }_{ij}^n\mathrm{\Omega }_{jn}^m\mathrm{\Omega }_{im}^n=$$
$$=\chi \underset{a}{}m_ac^2\frac{d^4p_a}{\sqrt{g}}\alpha \sqrt{\frac{g}{\stackrel{~}{g}}}\left(\stackrel{~}{g}_{ik}\stackrel{~}{g}_{jm}\frac{1}{2}\stackrel{~}{g}_{ij}\stackrel{~}{g}_{km}\right)u_a^ku_a^mN_a(q,p_a)+\chi \stackrel{~}{T}_{ij}^{(el)}.$$
(21)
Here $`R_{ij}`$ is the Richi tensor of the Riemannian spase with the metric $`g_{ij}`$, $`_m`$ is the covariant derivative in this spase.
With the help (9) and (10) we can write the Maxwell equations (2) to within the first-order terms in $`h_{ij}`$:
$$_kF^{ik}+_k\left(h_m^iF^{km}h_m^kF^{im}\right)+\frac{1}{2}F^{ik}_kh+_k\omega ^{ik}+$$
$$+_k\left(h_m^i\omega ^{km}h_m^k\omega ^{im}\right)+\frac{1}{2}\omega ^{ik}_kh=4\pi \underset{a}{}e_a\frac{d^4p_a}{\sqrt{g}}u_{(a)}^i\left(1\frac{1}{2}h\right)N_a(q,p_a).$$
(22)
In (21), 22) and below when raising and lowering the indexes we use the averaged metric $`g_{ij}`$, $`h=h_l^l`$.
Let’s expand the Einstein equations (21) up to the secand -order members in small $`h_{ij}`$ and $`\omega _{ij}`$:
$$R_{ij}+R_{ij}^{(1)}+R_{ij}^{(2)}+\mathrm{}=\underset{a}{}\chi m_ac^2\frac{d^4p_a}{\sqrt{g}}\left(\stackrel{~}{g}_{ik}\stackrel{~}{g}_{jm}\frac{1}{2}\stackrel{~}{g}_{ij}\stackrel{~}{g}_{km}\right)u_a^ku_a^mN_a(q,p_a)+$$
$$+\underset{a}{}\chi m_ac^2\frac{d^4p_a}{\sqrt{g}}L_{ijkm}^{(1)}(h)u_a^ku_a^mN_a+\mathrm{}+\chi T_{ij}^{(el)}+\chi \left(T_{ij}^{(el)}\right)^{(1)}+\chi \left(T_{ij}^{(el)}\right)^{(2)}.$$
(23)
Here $`R_{ij}^{(1)}`$ is the sum of all terms of expansion $`\stackrel{~}{R}_{ij}`$ that are first-order in $`h_{ij}`$, $`R_{ij}^{(2)}`$ is the sum of all second-order terms in $`h_{ij}`$, ets.; $`L_{ijkm}^{(1)}(h)`$ is the sum of all first - order terms in $`h_{ij}`$ of expansion the expression
$$\alpha \sqrt{\frac{g}{\stackrel{~}{g}}}\left(\stackrel{~}{g}_{ik}\stackrel{~}{g}_{jm}\frac{1}{2}\stackrel{~}{g}_{ij}\stackrel{~}{g}_{km}\right),$$
$`T_{ij}^{(el)}`$ is the energy-momentum tensor of averaged electromagnetic field $`F_{ij}`$, $`\left(T_{ij}^{(el)}\right)^{(1)}`$ is the sum of all terms of expansion $`\stackrel{~}{T}_{ij}^{(el)}`$ that are first-order in $`h_{ij}`$ and $`\omega _{ij}`$, $`\left(T_{ij}^{(el)}\right)^{(2)}`$ is the sum of all terms of expansion $`\stackrel{~}{T}_{ij}^{(el)}`$ that are second-order in $`h_{ij}`$ and $`\omega _{ij}`$.
This expressions has the forms:
$$R_{ij}^{(1)}=_m\mathrm{\Omega }_{ij}^{(1)m}_j\mathrm{\Omega }_{im}^{(1)m},$$
(24)
$$R_{ij}^{(2)}=_m\mathrm{\Omega }_{ij}^{(2)m}_j\mathrm{\Omega }_{im}^{(2)m}+\mathrm{\Omega }_{mn}^{(1)m}\mathrm{\Omega }_{ij}^{(1)n}\mathrm{\Omega }_{jn}^{(1)m}\mathrm{\Omega }_{im}^{(1)n},$$
(25)
$$\mathrm{\Omega }_{ij}^{(1)m}=\frac{1}{2}g^{ml}(_lh_{ij}+_ih_{lj}+_jh_{li}),$$
(26)
$$\mathrm{\Omega }_{ij}^{(2)}=\frac{1}{2}h^{ml}(_lh_{ij}+_ih_{lj}+_jh_{li})=\frac{1}{2}h_l^m\mathrm{\Omega }^{(1)l}ij.$$
(27)
$$L_{ijkm}^{(1)}=\frac{1}{2}\left(h_{st}u^su^t+h_{st}g^{st}\right)\left(g_{ik}g_{jm}\frac{1}{2}g_{ij}g_{km}\right)+$$
$$+h_{ik}g_{jm}+g_{ik}h_{jm}\frac{1}{2}h_{ij}g_{km}\frac{1}{2}g_{ij}h_{km},$$
(28)
$$T_{ij}^{(el)}=\frac{1}{4\pi }\left(F_{il}F_{jm}g^{lm}+\frac{1}{4}g_{ij}F_{lm}F^{lm}\right),$$
(29)
$$\left(T_{ij}^{(el)}\right)^{(1)}=\frac{1}{4\pi }[F_{il}F_{im}h^{lm}+\frac{1}{4}h_{ij}F_{lm}F^{lm}\frac{1}{2}g_{ij}F_{lm}F_{kn}g^{lk}h^{mn}$$
$$\omega _{il}F_j^{.l}\omega _{jl}F_i^{.l}+\frac{1}{2}F^{lm}\omega _{lm}],$$
(30)
$$\left(T_{ij}^{(el)}\right)^{(2)}=\frac{1}{4\pi }[F_{il}F_{jm}g_{st}+\frac{1}{2}g_{ij}F_{lk}F_m^{.k}g_{st}+\frac{1}{4}g_{ij}F_{lm}F_{st}$$
$$\frac{1}{2}F_{.t}^kF_{km}g_{il}g_{js}]h^{ls}h^{tm}+$$
$$+\frac{1}{4\pi }\left[F_{it}g_{jl}g_{sm}+F_{jt}g_{il}g_{sm}g_{ij}F_{lm}g_{st}+\frac{1}{2}g_{it}g_{jm}F_{ls}\right]\omega ^{ls}h^{tm}+$$
$$+\frac{1}{4\pi }\left(g_{it}g_{is}+\frac{1}{4}g_{ij}g_{st}\right)\omega _{.l}^t\omega ^{sl}.$$
(31)
2. Averaging of microscopic system of Einstein and Maxwell equations for the relativistic plasma
We average (22) and (23) over the paths (Ref. -) and introduce the one-particle distribution function
$$f_a(q,p)=𝑑s\delta ^4(q^iq_{a(l)}^i(s))\delta ^4(p_jp_j^{a(l)}(s))=\frac{1}{n_a}N_a.$$
(32)
As a result we have the averaged Einstein equations in the form
$$R_{ij}+\mathrm{\Lambda }_{ij}=\chi (T_{ij}^{(m)}(1/2)T^{(m)}g_{ij})+$$
$$+\chi (T_{ij}^{(el)}\frac{1}{2}T^{(el)}g_{ij})++\chi (T_{ij}^{(r)}\frac{1}{2}T^{(r)}g_{ij}).$$
(33)
Here
$$T_{(m)}^{ij}=\underset{a}{}n_am_ac^2\frac{d^4p_a}{\sqrt{g}}u_a^iu_a^jf_a(q,p_a)$$
(34)
is the macroscopic energy-momentum tensor of medium, $`T_{(el)}^{ij}`$ is the energy-momrntum tensor of macroscopic electromagnetic field (see. (29),
$$T_{(r)}^{ij}=\frac{1}{4\pi }\omega _{.l}^i\omega ^{jl}+\frac{1}{4}g^{ij}\omega _{lm}\omega ^{lm}$$
(35)
is the macroscopic energy-momentum tensor of radiation in plasma,
$$\mathrm{\Lambda }_{ij}=R_{ij}^{(2)}\underset{a}{}\chi m_ac^2\frac{d^4p_a}{\sqrt{g}}N_aL_{ijkm}^{(1)}u_a^mu_a^k.$$
(36)
When obtained (33) we taking into account that
$$R_{ij}^{(1)}=0,\left(T_{ij}^{(el)}\right)^{(1)}=0.$$
Next we assum that
$$F_{ik}h_j^k\omega _{ij}ets.$$
(37)
inside the correlation region. Consequantly $`\left(T_{ij}^{(el)}\right)^{(2)}`$ is equal approximately to $`T_{ij}^{(r)}`$.
Taking into account the (25) - (28) we can write $`\mathrm{\Lambda }_{ij}`$ in the form
$$\mathrm{\Lambda }_{ij}=_k\phi _{ij}^k+\mu _{ij},$$
(38)
where
$$\phi _{ij}^k=\frac{1}{2}\left(\delta _n^k\delta _j^s\delta _j^k\delta _n^s\right)P_{is}^n,$$
(39)
$$P_{is}^n=h_l^n\mathrm{\Omega }^{(1)l}is,$$
(40)
$$\mu _{ij}=\left(\delta _n^k\delta _j^s\delta _j^k\delta _n^s\right)Q_{kis}^n+\lambda _{ij},$$
(41)
$$Q_{kis}^n=\mathrm{\Omega }_{kl}^{(1)n}\mathrm{\Omega }_{is}^{(1)l},$$
(42)
$$\lambda _{ij}=\underset{a}{}\chi m_ac^2\frac{d^4p}{\sqrt{(g)}}\{\frac{1}{2}u_iu_ju^ku^m\frac{1}{4}g_{ij}u^ku^m\frac{1}{2}u_iu_jg^{km}+$$
$$+\frac{1}{4}g_{ij}g^{km}+u_iu^k\delta _j^m+u_ju^k\delta _i^m\frac{1}{2}\delta _i^k\delta _j^m\}<N_ah_{km}>.$$
(43)
In (43) we reject the indices ”a” on momentums and velosities of particles of species ”a”.
The macroscopic Maxwell equations, obtained from (22) after averaging, have the form
$$_kF^{ik}+_k\phi ^{ik}+\mu ^i=\frac{4\pi }{c}J^i,$$
(44)
where
$$\phi ^{ik}=h_m^i\omega ^{km}h_m^k\omega ^{im},$$
(45)
$$\mu ^i=\frac{1}{2}\omega ^{ik}_kh+\lambda ^i,$$
(46)
$$\lambda ^i=2\pi \underset{b}{}e_b\frac{d^4p}{\sqrt{g}}u_{(b)}^iN_bh,$$
(47)
$$J^i=\underset{b}{}e_bcn_b\frac{d^4p}{\sqrt{g}}u_b^if_b$$
(48)
is the macroscopic current vector.
To simplify still furthe, we only have to calculate $`h_{ij}`$ and $`\omega _{ij}`$ inside the region determined by the correlations radius and corresponding correlation time. Note, That distant collisions provide the main contribution to caculated macroscopic quantities. To consider this contribution it is enough to find $`h_{ij}`$ and $`\omega _{ij}`$ from the Einstein-Maxwell equations linearised with respect to the metric $`g_{ij}`$ and the macroscopic electromagnetic field $`F_{ij}`$.
We assume the average gravitational field generated by the particles to be constant within the correlation region. In this case we can interpret $`g_{ij}`$ within the correlation region as the Minkowski metric. We assume alsow, that the influence of macroscopic electromagnetical field on the microscopic field in correlation region is small (see (37)) .
As a result we have linearized Einstein and Maxwell equations with respect the Minkowski metric $`g_{ij}`$. By employing the gauge $`_k\gamma ^{ik}=0`$, where $`\gamma _{ij}=h_{ij}(1/2)hg_{ij}`$, we get the folloing form of the linearized Einstein and Maxwell equations
$$\mathrm{}\gamma ^{ij}=\underset{b}{}2\chi m_bc^2d^4p_{}^{}{}_{b}{}^{}\mathrm{\Phi }_b(q,p_{}^{}{}_{b}{}^{})u_{}^{}{}_{b}{}^{i}u_{}^{}{}_{b}{}^{j},$$
(49)
$$_k\omega ^{ik}=4\pi \underset{b}{}e_bd^4p_{}^{}{}_{b}{}^{}\mathrm{\Phi }_b(q.p_{}^{}{}_{b}{}^{})u_{}^{}{}_{b}{}^{i},$$
(50)
Here $`\mathrm{}=g^{ij}_i_j`$, $`\mathrm{\Phi }_b=N_bn_bf_b`$.
Thus subsequent calculatios do not have a covariant form, but they are all done for the purpose of determining the components of the tensors $`\phi _{ij}^k`$, $`\mu _{ij}`$, $`\phi _{ij}`$, $`\mu _i`$ and $`T_{ij}^{(r)}`$ at some (arbitrary) point $`(q)`$ in the locally Lorentzian frame. In this reference frame the interval $`ds^2`$ has the form
$$ds^2=d\eta ^2(dq^1)^2(dq^2)^2(dq^3)^2.$$
(51)
The final result must be written in covariant form.
The expressions for $`h_{ij}`$ and $`\mathrm{\Omega }_{ij}^l`$ we get from the Einstein equations (49)linearized with respect to the Minkowski metric (which we still denote by $`g_{ij}`$)) were found in ((refer to.):
$$h_{ij}(\eta ,𝐪)=\underset{b}{}d^4p^{}d^3q^{}d^3k_{\mathrm{}}^\eta d\eta ^{}e^{i𝐤(𝐪𝐪^{})}\times $$
$$\times h_{ij}^{(b)}(\eta ,\eta ^{},p^{},𝐤)\mathrm{\Phi }_b(\eta ^{},𝐪^{},p^{}),$$
(52)
$$\mathrm{\Omega }_{jk}^i(\eta ,𝐪)=\underset{b}{}d^4p^{}d^3q^{}d^3k_{\mathrm{}}^\eta d\eta ^{}e^{i𝐤(𝐪𝐪^{})}\times $$
$$\times \mathrm{\Omega }_{jk}^{i(b)}(\eta ,\eta ^{},p^{},𝐤)\mathrm{\Phi }_b(\eta ^{},𝐪^{},p^{}).$$
(53)
where ($`𝐪=(q^1,q^2,q^3)`$ is the three - dimensional radius vector in the given reference frame, and $`𝐤=(k_1,k_2,k_3)`$),
$$\mathrm{\Phi }_b=N_bn_bf_b$$
,
$$h_{ij}^{(b)}(\eta ,\eta ^{},p^{},𝐤)=\frac{i\chi m_bc^2}{(2\pi )^3k}(u_i^{}u_j^{}\frac{1}{2}g_{ij})\left\{e^{ik(\eta ^{}\eta )}e^{ik(\eta ^{}\eta )}\right\},$$
(54)
$$\mathrm{\Omega }_{jk}^{i(b)}(\eta ,\eta ^{},p^{},𝐤)=\frac{\chi m_bc^2}{2(2\pi )^3k}\{[(u_j^{}u_k^{}\frac{1}{2}g_{jk})k_+^i(u_j^{}u^i\frac{1}{2}\delta _j^i)k_k^+$$
$$(u_k^{}u^i\frac{1}{2}\delta _k^i)k_j^+]e^{ik(\eta ^{}\eta )}[(u_j^{}u_k^{}\frac{1}{2}g_{jk})k_{}^i(u_j^{}u^i\frac{1}{2}\delta _j^i)k_k^{}$$
$$(u_k^{}u^i\frac{1}{2}\delta _k^i)k_j^{}]e^{ik(\eta ^{}\eta )}\}.$$
(55)
In (54) and (55) the following vectors were introduced:
$$k_i^+=(k,𝐤),k_i^{}=(k,𝐤),$$
where
$$k=\sqrt{[(k_1)^2+(k_2)^2+(k_3)^2]}=|𝐤|.$$
Obviously, $`k_i^{}(𝐤)=k_i^+(𝐤)`$.
To obtain the additional terms $`_k\phi ^{ik}`$ and $`\mu ^i`$ in macroscopic Maxwell equations we have to calculate the $`h=h_l^l`$ and $`_kh`$:
$$h(\eta ,𝐪)=\underset{b}{}d^4p^{}d^3q^{}d^3k_{\mathrm{}}^\eta 𝑑\eta ^{}e^{i𝐤(𝐪𝐪^{})}$$
$$h^{(b)}(\eta ,\eta ^{},p^{},𝐤)\mathrm{\Phi }_b(\eta ^{},𝐪^{},p^{}),$$
(56)
where
$$h^{(b)}(\eta ,\eta ^{},p^{},𝐤)=\frac{i\chi m_bc^2}{(2\pi )^3k}\left\{e^{ik(\eta ^{}\eta )}e^{ik(\eta ^{}\eta )}\right\}.$$
(57)
$$_kh(\eta ,𝐪)=\underset{b}{}d^4p^{}d^3q^{}d^3k_{\mathrm{}}^\eta 𝑑\eta ^{}e^{i𝐤(𝐪𝐪^{})}$$
$$h_{;k}^{(b)}(\eta ,\eta ^{},p^{},𝐤)\mathrm{\Phi }_b(\eta ^{},𝐪^{},p^{}),$$
(58)
where
$$h_{;k}^{(b)}(\eta ,\eta ^{},p^{},𝐤)=\frac{\chi m_bc^2}{(2\pi )^3k}\left\{k_k^+e^{ik(\eta ^{}\eta )}k_k^{}e^{ik(\eta ^{}\eta )}\right\}.$$
(59)
Let’s write the solution of (50) in the form
$$\omega _{ik}=_iA_k_kA_i.$$
where :
$$A_i(\eta ,𝐪)=\underset{b}{}d^4p^{}d^3q^{}d^3k_{\mathrm{}}^\eta 𝑑\eta ^{}e^{i𝐤(𝐪𝐪^{})}$$
$$A_i^{(b)}(\eta ,\eta ^{},p^{},𝐤)\mathrm{\Phi }_b(\eta ^{},𝐪^{},p^{}).$$
(60)
Here
$$A_i^{(b)}(\eta ,\eta ^{},p^{},𝐤)=\frac{ie_b}{(2\pi )^2k}u_i^{}\left\{e^{ik(\eta ^{}\eta )}e^{ik(\eta ^{}\eta )}\right\}.$$
(61)
The $`\omega _{ik}`$ have the form:
$$\omega _{ik}(\eta ,𝐪)=\underset{b}{}d^4p^{}d^3q^{}d^3k_{\mathrm{}}^\eta d\eta ^{}e^{i𝐤(𝐪𝐪^{})}\times $$
$$\times \omega _{ik}^{(b)}(\eta ,\eta ^{},p^{},𝐤)\mathrm{\Phi }_b(\eta ^{},𝐪^{},p^{}),$$
(62)
where
$$\omega _{ik}^{(b)}(\eta ,\eta ^{},p^{},𝐤)=\frac{e_b}{(2\pi )^2k}\left\{(k_i^+u_k^{}k_k^+u_i^{})e^{ik(\eta ^{}\eta )}(k_i^{}u_k^{}k_k^{}u_i^{})e^{ik(\eta ^{}\eta )}\right\}.$$
(63)
If substitute the (52), (53), (56), (58), (62) to (40), (42), (43), (45) — (47) we get the folloing expressions for $`P_{is}^n`$ , $`Q_{kis}^n`$, $`\lambda _{ij}`$, $`h_m^i\omega ^{km}`$, $`\omega ^{ik}_kh`$, $`\lambda _i`$:
$$P_{is}^n=\underset{bc}{}d^4p^{}d^4p^{\prime \prime }d^3q^{}d^3q^{\prime \prime }_{\mathrm{}}^\eta d\eta ^{}_{\mathrm{}}^\eta d\eta ^{\prime \prime }d^3k^{}d^3k^{\prime \prime }\times $$
$$\times e^{i𝐤^{}(𝐪𝐪^{})}e^{i𝐤^{\prime \prime }(𝐪𝐪^{\prime \prime })}h_l^{n(b)}(\eta ,\eta ^{},p^{},𝐤^{})\mathrm{\Omega }_{is}^{l(c)}(\eta ,\eta ^{\prime \prime },p^{\prime \prime },𝐤^{\prime \prime })n_bn_cg_{bc}(x^{},x^{\prime \prime }),$$
(64)
$$Q_{kis}^n=\underset{bc}{}d^4p^{}d^4p^{\prime \prime }d^3q^{}d^3q^{\prime \prime }_{\mathrm{}}^\eta d\eta ^{}_{\mathrm{}}^\eta d\eta ^{\prime \prime }d^3k^{}d^3k^{\prime \prime }\times $$
$$\times e^{i𝐤^{}(𝐪𝐪^{})}e^{i𝐤^{\prime \prime }(𝐪𝐪^{\prime \prime })}\mathrm{\Omega }_{kl}^{n(b)}(\eta ,\eta ^{},p^{},𝐤^{})\mathrm{\Omega }_{is}^{l(c)}(\eta ,\eta ^{\prime \prime },p^{\prime \prime },𝐤^{\prime \prime })n_bn_cg_{bc}(x^{},x^{\prime \prime }),$$
(65)
$$\lambda _{ij}=\underset{bc}{}\chi m_cc^2d^4p^{}d^4p^{\prime \prime }d^3q^{}_{\mathrm{}}^\eta d\eta ^{}d^3k^{}e^{i𝐤^{}(𝐪𝐪^{})}\times $$
$$\times \{\frac{1}{2}u_{}^{\prime \prime }{}_{i}{}^{}u_{}^{\prime \prime }{}_{j}{}^{}u_{}^{\prime \prime }{}_{}{}^{k}u_{}^{\prime \prime }{}_{}{}^{m}\frac{1}{4}g_{ij}u_{}^{\prime \prime }{}_{}{}^{k}u_{}^{\prime \prime }{}_{}{}^{m}\frac{1}{2}u_{}^{\prime \prime }{}_{i}{}^{}u_{}^{\prime \prime }{}_{j}{}^{}g^{km}+\frac{1}{4}g_{ij}g^{km}+u_{}^{\prime \prime }{}_{i}{}^{}u_{}^{\prime \prime }{}_{}{}^{k}\delta _j^m+$$
$$+u_{}^{\prime \prime }{}_{j}{}^{}u_{}^{\prime \prime }{}_{}{}^{k}\delta _i^m\frac{1}{2}\delta _i^k\delta _j^m\}h^{(b)}_{km}(\eta ,\eta ^{},p^{},𝐤^{})n_bn_cg_{bc}(x^{};\eta ,𝐪,p^{\prime \prime }),$$
(66)
$$h_m^i\omega ^{km}=\underset{bc}{}d^4p^{}d^4p^{\prime \prime }d^3𝐪^{}d^3𝐪^{\prime \prime }_{\mathrm{}}^\eta 𝑑\eta ^{}_{\mathrm{}}^\eta 𝑑\eta ^{\prime \prime }d^3𝐤^{}d^3𝐤^{\prime \prime }$$
$$e^{i𝐤^{}(𝐪𝐪^{})}e^{i𝐤^{\prime \prime }(𝐪𝐪^{\prime \prime })}h_m^{i(b)}(\eta ,\eta ^{},p^{},𝐤^{})\omega _{(c)}^{km}(\eta ,\eta ^{\prime \prime },p^{\prime \prime },𝐤^{\prime \prime })n_bn_cg_{bc}(x^{},x^{\prime \prime }),$$
(67)
$$\omega ^{ik}_kh=\underset{bc}{}d^4p^{}d^4p^{\prime \prime }d^3𝐪^{}d^3𝐪^{\prime \prime }_{\mathrm{}}^\eta 𝑑\eta ^{}_{\mathrm{}}^\eta 𝑑\eta ^{\prime \prime }d^3𝐤^{}d^3𝐤^{\prime \prime }$$
$$e^{i𝐤^{}(𝐪𝐪^{})}e^{i𝐤^{\prime \prime }(𝐪𝐪^{\prime \prime })}h_{;k}^{(b)}(\eta ,\eta ^{},p^{},𝐤^{})\omega _{(c)}^{ik}(\eta ,\eta ^{\prime \prime },p^{\prime \prime },𝐤^{\prime \prime })n_bn_cg_{bc}(x^{},x^{\prime \prime }),$$
(68)
$$\lambda ^i=\underset{bc}{}2\pi e_cd^4p^{}d^4p^{\prime \prime }d^3𝐪^{}_{\mathrm{}}^\eta 𝑑\eta ^{}d^3𝐤^{}e^{i𝐤^{}(𝐪𝐪^{})}u_{}^{\prime \prime }{}_{}{}^{i}$$
$$h^{(b)}(\eta ,\eta ^{},p^{},𝐤^{})n_bn_cg_{bc}(x^{};\eta ,𝐪,p^{\prime \prime }).$$
(69)
In this expressions unprimed, primed and double-primed quantities refer to particles belonging to species ”a”, ”b” and spesies ”c” respectively.
In (64) - (69) we introduced the two -particle correlation function $`g_{ab}(x^{},x^{\prime \prime })`$ (see Ref. , , , ):
$$f_{ab}(x,x^{})=f_a(x)f_b(x^{})+g_{ab}(x,x^{}).$$
(70)
Here $`f_a(x)`$, $`f_{ab}(x,x^{})`$, $`f_{abc}(x,x^{},x^{\prime \prime })`$ are the one-particle, two-particle and tree-particle distribution functions respectively:
$$𝑑s\delta (xx_a(s))=f_a(x),$$
$$𝑑s\delta (xx_a(s))𝑑s^{}\delta (x^{}x_b(s^{}))=f_{ab}(x,x^{}),$$
$$𝑑s\delta (xx_a(s))𝑑s^{}\delta (x^{}x_b(s^{}))𝑑s^{\prime \prime }\delta (x^{\prime \prime }x_c(s^{\prime \prime }))=f_{abc}(x,x^{},x^{\prime \prime }),$$
where
$$\delta (xx_a(s))=\delta ^4(q^iq_a^i(s))\delta ^4(p_jp_j^a(s)).$$
We denote the set of all variable $`(\eta ,𝐪,p_i)`$ by $`x`$, the set $`(\eta ^{},𝐪^{},p_i^{})`$ by $`x^{}`$, while the momenta $`p_{(b)}^{}`$ are denoted by $`p^{}`$, and the $`p_{(c)}^{\prime \prime }`$ by $`p^{\prime \prime }`$.
For the moments of random functions we have the formulas (Ref., ):
$$N_a(x)=n_af_a(x),$$
(71)
$$N_a(x)N_b(x^{})=(n_an_bn_a\delta _{ab})f_{ab}(x,x^{})+$$
,
$$+n_a\delta _{ab}f_a(x)𝑑s^{}\delta (x^{}x_a(s^{}/x)),$$
(72)
$$N_a(x)N_b(x^{})N_c(x^{\prime \prime })=(n_an_bn_cn_an_b\delta _{ac}n_an_b\delta _{bc}$$
$$n_an_c\delta _{ab}+2n_a\delta _{ab}\delta _{bc})f_{abc}(x,x^{},x^{\prime \prime })+$$
$$+(n_an_cn_a\delta _{ac})\delta _{ab}f_{ac}(x,x^{\prime \prime })𝑑s^{}\delta (x^{}x_a(s^{}/x))+$$
$$+(n_an_bn_a\delta _{ab})\delta _{ac}f_{ab}(x,x^{})𝑑s^{\prime \prime }\delta (x^{\prime \prime }x_a(s^{\prime \prime }/x))+$$
$$+(n_an_bn_a\delta _{ab})\delta _{bc}f_{ab}(x,x^{})𝑑s^{\prime \prime }\delta (x^{\prime \prime }x_b(s^{\prime \prime }/x^{}))+$$
$$+n_a\delta _{ab}\delta _{bc}f_a(x)𝑑s^{}\delta (x^{}x_a(s^{}/x))𝑑s^{\prime \prime }\delta (x^{\prime \prime }x_a(s^{\prime \prime }/x)).$$
(73)
Here $`x_a(s/x)`$ stands for the particle path through point $`x`$ of the phase space. Bearing in mind that $`\mathrm{\Phi }_a=N_an_af_a`$ and that $`f_a`$ is not a random function, we can easily obtain expressions for the avereges
$$N_a(x)\mathrm{\Phi }_b(x^{}),N_a(x)\mathrm{\Phi }_b(x^{})\mathrm{\Phi }_c(x^{\prime \prime }).$$
In deriving (64) — (69) we assumed that $`n_a1`$ and that $`x^{\prime \prime }=x_b(s^{\prime \prime }/x^{})`$, i.e. point $`x^{\prime \prime }`$ is not on the path of particles of spesies ”b” passing through the point $`x^{}`$ of the phase spase.
In work two-particles correlation functions $`g_{ab}(x^{},x^{\prime \prime })`$ are found for the system gravitationally interacting particles. The two-particles correlation functions the for system of electromagnetically interacting particles were found by the author in Ref. , when getting the relativistic kinetic equation for the plasma (Eq. (18) from Ref. ).
In our case we should find two-particles correlation function $`g_{ab}(x,x^{})`$ caused by electromagnetic and gravitationall interactions simultaneously.
To obtain correlation function inside the correlation region we assume that $`h_{ik}1`$, $`\omega _{ik}1`$. Therefore $`\mathrm{\Omega }_{ij}^k\mathrm{\Omega }_{ij}^{k(1)}`$, $`\psi _{.k}^l\omega _{.k}^l`$.
After substituting (53), (62) into (16), multiplying (16) by $`\mathrm{\Phi }_b(x^{})`$ and averaging we get
$$p^i\frac{}{q^i}N_a(x)\mathrm{\Phi }_b(x^{})+\mathrm{\Gamma }_{j,ik}p^jp^k\frac{}{p_i}N_a(x)\mathrm{\Phi }_b(x^{})+\frac{e_a}{c}F_{jk}p^k\frac{}{p_j}N_a(x)\mathrm{\Phi }_b(x^{})=$$
$$=\frac{}{p_i}\{\underset{c}{}d^4p_b^{\prime \prime }d^3𝐪^{\prime \prime }d^3𝐤_{\mathrm{}}^\eta d\eta ^{\prime \prime }\mathrm{exp}[i𝐤(𝐪𝐪^{})]\times $$
$$\times [\mathrm{\Omega }_{lm}^{j(c)}(\eta ,\eta ^{\prime \prime },p_b^{\prime \prime },𝐤)p^lp^m\mathrm{\Delta }_{ji}\frac{e_a}{c}\omega _{.k}^{l(c)}(\eta ,\eta ^{},p^{\prime \prime },𝐤)p^k\mathrm{\Delta }_{lj}]N_a(x)\mathrm{\Phi }_b(x^{})\mathrm{\Phi }_c(x^{\prime \prime })\}.$$
(74)
Next we assume that
$$f_{abc}(x,x^{},x^{\prime \prime })f_a(x)f_b(x^{}),f_c(x^{\prime \prime })$$
.
In view of Eqs.(72), (73) the two partial correlation function $`g_{ab}(x,x^{})`$ obeys the equation
$$p^i\frac{}{q^i}g_{ab}(x,x^{})+\mathrm{\Gamma }_{i,jk}p^ip^k\frac{}{p_j}g_{ab}(x,x^{})+\frac{e_a}{c}F_{jk}p^k\frac{}{p_j}g_{ab}(x,x^{})=$$
$$=\frac{}{p_i}\{d^4p_b^{\prime \prime }d^3𝐪^{\prime \prime }d^3𝐤_{\mathrm{}}^\eta d\eta ^{\prime \prime }\mathrm{exp}[i𝐤(𝐪𝐪^{\prime \prime })]\times $$
$$\times [\mathrm{\Omega }_{lm}^{j(b)}(\eta ,\eta ^{\prime \prime },p^{\prime \prime },𝐤)p^lp^m\mathrm{\Delta }_{ji}\frac{e_a}{c}\omega _{jk}^{(c)}(\eta ,\eta ^{\prime \prime },p^{\prime \prime },𝐤)p^k]\times $$
$$\times f_a(x)f_b(x^{})ds^{\prime \prime }\delta (x^{\prime \prime }x_b(s^{\prime \prime }/x^{}))\}.$$
(75)
In this equation for $`g_{ab}(x,x^{})`$ we should put $`\mathrm{\Gamma }_{i,jk}=0`$, $`F_{jk}=0`$, since we assume that within the correlation region the metric $`g_{ij}`$ are constant and that the influence of macroscopic electromagnetic field is small.
So we have the first order linear equation (75), whose right-hand side contains the sum of two terms. The first term caused by gravitationall interactions, the second one caused by electromagnetic interaction.
Consequently, we can write the solution of Eqs. (75) in the form
$$g_{ab}(x,x^{})=g_{ab}^{(gr)}(x,x^{})+g_{ab}^{(el)}(x,x^{}).$$
(76)
Here $`g_{ab}^{(gr)}(x,x^{})`$ caused by gravitational interaction, and $`g_{ab}^{(el)}(x,x^{})`$ caused by electromagnetic interaction.
The equation for $`g_{ab}^{(gr)}(x,x^{})`$ coinside with one in Ref. . In we got the $`g_{ab}^{(gr)}(x,x^{})`$ in the form (see Eq. (46) in Ref. ):
$$g_{ab}^{(gr)}(x,x^{})=d^3k_{\mathrm{}}^\eta \frac{d\tau }{p^0}\left[\frac{}{p_i}(p^lp^m\mathrm{\Delta }_{ij}f_a(x))\right]_\tau _{\mathrm{}}^\tau \frac{d\tau ^{}}{u^0}f_b(x^{})\times $$
$$\times \mathrm{\Omega }_{lm}^{j(b)}(\tau ,\tau ^{},p^{},𝐤)exp[i𝐤(𝐪𝐪^{})+\frac{i}{c}(\mathrm{𝐤𝐯})(\eta \tau )+\frac{i}{c}(\mathrm{𝐤𝐯}^{})(\tau ^{}\eta ^{})]+$$
$$+d^3k_{\mathrm{}}^\eta ^{}\frac{d\tau ^{}}{p_{}^{}{}_{}{}^{0}}\left[\frac{}{p_{}^{}{}_{i}{}^{}}(p_{}^{}{}_{}{}^{l}p_{}^{}{}_{}{}^{m}\mathrm{\Delta }_{}^{}{}_{ij}{}^{}f_b(x^{}))\right]_\tau ^{}_{\mathrm{}}^\tau \frac{d\tau }{u^0}f_a(x)\times $$
$$\times \mathrm{\Omega }_{lm}^{j(a)}(\tau ^{},\tau ,p,𝐤)exp[i𝐤(𝐪𝐪^{})+\frac{i}{c}(\mathrm{𝐤𝐯}^{})(\eta ^{}\tau ^{})+\frac{i}{c}(\mathrm{𝐤𝐯})(\tau \eta )],$$
(77)
where $`𝐯=c𝐮_a/u_a^0`$, $`𝐯^{}=c𝐮_{}^{}{}_{b}{}^{}/u_{}^{}{}_{b}{}^{0}`$, $`𝐮_a=(u^1,u^2,u^3)`$, $`𝐮_{}^{}{}_{b}{}^{}=(u_{}^{}{}_{}{}^{1},u_{}^{}{}_{}{}^{2},u_{}^{}{}_{}{}^{3})`$. Here the subscript $`\tau `$ indicates that after calculating the derivatives with respect $`p`$ we must replace the arguments $`\eta `$, and $`𝐪`$ by $`\tau `$ and $`𝐪+\frac{𝐯}{c}(\tau \eta )`$, respectively. The subscript $`\tau ^{}`$ indicates that after calculating the derivatives with respect $`p^{}`$ we must replace the arguments $`\eta ^{}`$, and $`𝐪^{}`$ by $`\tau ^{}`$ and $`𝐪^{}+\frac{𝐯^{}}{c}(\tau ^{}\eta ^{})`$, respectively.
The equation for $`g_{ab}^{(el)}(x,x^{})`$ coinside with one in Ref.. In Ref. we got the solution for $`g_{ab}^{(el)}(x,x^{})`$ (See (18) from Ref. ).
With the preceding notation the result (18) from Ref. takes the form:
$$g_{ab}^{(el)}(x,x^{})=\frac{e_be_c}{2\pi ^2c}\frac{d^3k}{k}_{\mathrm{}}^\eta ^{}d\tau ^{}\frac{u^k}{u^0}f_c(x^{\prime \prime })\frac{f_b(x^{})}{p_i^{}}_{\mathrm{}}^\tau ^{}d\tau ^{\prime \prime }\times $$
$$\times \{\frac{u_k^{\prime \prime }}{u^{\prime \prime 0}}\frac{}{q_\tau ^{}^i}[exp(i𝐤(𝐪^{}𝐪^{\prime \prime })sin(k(\tau ^{}\tau ^{\prime \prime }))]$$
$$\frac{u_i^{\prime \prime }}{u^{\prime \prime 0}}\frac{}{q_\tau ^{}^k}[exp(i𝐤(𝐪^{}𝐪^{\prime \prime })sin(k(\tau ^{}\tau ^{\prime \prime }))]\}\times $$
$$\times exp[\frac{i}{c}(\mathrm{𝐤𝐯}^{})(\eta ^{}\tau ^{})+\frac{i}{c}(\mathrm{𝐤𝐯}^{\prime \prime })(\tau ^{\prime \prime }\eta ^{\prime \prime })]$$
$$\frac{e_be_c}{2\pi ^2c}\frac{d^3k}{k}_{\mathrm{}}^{\eta ^{\prime \prime }}d\tau ^{\prime \prime }\frac{u^{\prime \prime k}}{u^{\prime \prime 0}}f_b(x^{})\frac{f_c(x^{\prime \prime })}{p_i^{\prime \prime }}_{\mathrm{}}^{\tau ^{\prime \prime }}d\tau ^{}\times $$
$$\times \{\frac{u_k^{}}{u^0}\frac{}{q_{\tau ^{\prime \prime }}^{\prime \prime i}}[exp(i𝐤(𝐪^{\prime \prime }𝐪^{})sin(k(\tau ^{\prime \prime }\tau ^{}))]$$
$$\frac{u_i^{}}{u^0}\frac{}{q_{\tau ^{\prime \prime }}^{\prime \prime k}}[exp(i𝐤(𝐪^{\prime \prime }𝐪^{})sin(k(\tau ^{\prime \prime }\tau ^{}))]\}\times $$
$$\times exp[\frac{i}{c}(\mathrm{𝐤𝐯}^{\prime \prime })(\eta ^{\prime \prime }\tau ^{\prime \prime })+\frac{i}{c}(\mathrm{𝐤𝐯}^{})(\tau ^{}\eta ^{})].$$
(78)
Here $`q_\tau ^i=(\tau ,q^\alpha )`$.
After performing a differentiation with respect to $`q_\tau ^i`$ in (26) we get a following expression for correlation function:
$$g_{ab}^{(el)}(x,x^{})=\frac{e_be_c}{4\pi ^2c}\frac{d^3k}{k}_{\mathrm{}}^\eta ^{}\frac{d\tau ^{}}{u^0}f_c(x^{\prime \prime })\frac{f_b(x^{})}{p_i^{}}_{\mathrm{}}^\tau ^{}\frac{d\tau ^{\prime \prime }}{u^{\prime \prime 0}}e^{i𝐤(𝐪^{}𝐪^{\prime \prime })}[(u^{}u^{\prime \prime })\delta _i^k$$
$$u_i^{\prime \prime }u^k](k_k^+e^{ik(\tau ^{\prime \prime }\tau ^{})}k_k^{}e^{ik(\tau ^{\prime \prime }\tau ^{})})exp[\frac{i}{c}(\mathrm{𝐤𝐯}^{})(\eta ^{}\tau ^{})+\frac{i}{c}(\mathrm{𝐤𝐯}^{\prime \prime })(\tau ^{\prime \prime }\eta ^{\prime \prime })]$$
$$\frac{e_be_c}{4\pi ^2c}\frac{d^3k}{k}_{\mathrm{}}^{\eta ^{\prime \prime }}\frac{d\tau ^{\prime \prime }}{u^{\prime \prime 0}}f_b(x^{})\frac{f_c(x^{\prime \prime })}{p_i^{\prime \prime }}_{\mathrm{}}^{\tau ^{\prime \prime }}\frac{d\tau ^{}}{u^0}e^{i𝐤(𝐪^{\prime \prime }𝐪^{})}[(u^{}u^{\prime \prime })\delta _i^k$$
$$u_i^{}u^{\prime \prime k}](k_k^+e^{ik(\tau ^{}\tau ^{\prime \prime })}k_k^{}e^{ik(\tau ^{}\tau ^{\prime \prime })})exp[\frac{i}{c}(\mathrm{𝐤𝐯}^{\prime \prime })(\eta ^{\prime \prime }\tau ^{\prime \prime })+\frac{i}{c}(\mathrm{𝐤𝐯}^{})(\tau ^{}\eta ^{})].$$
(79)
It is evident, that the electromagnetic interactios in plasma are dominating. Consequently
$$g_{ab}^{(gr)}(x,x^{})g_{ab}^{(el)}(x,x^{}).$$
That is why one can put the $`g_{ab}^{(el)}(x,x^{})`$ instead of $`g_{ab}(x,x^{})`$ in (64) - (69).
If we now substitude (79) into (64) - (69) and integrate with respect $`𝐪^{}`$, $`𝐪^{\prime \prime }`$, $`𝐤^{}`$, $`𝐤^{\prime \prime }`$ we get the following expressions:
$$P_{is}^n=\underset{bc}{}(\frac{e_be_cn_bn_c(2\pi )^6}{4\pi ^2c})d^4p^{}d^4p^{\prime \prime }_{\mathrm{}}^\eta d\eta ^{}_{\mathrm{}}^\eta d\eta ^{\prime \prime }\frac{d^3k}{k}\times $$
$$\times \{_{\mathrm{}}^\eta ^{}\frac{d\tau ^{}}{u^0}_{\mathrm{}}^\tau ^{}\frac{d\tau ^{\prime \prime }}{u^{\prime \prime 0}}f_c(x^{\prime \prime })\frac{f_b(x^{})}{p_j^{}}[(u^{}u^{\prime \prime })\delta _j^ku_j^{\prime \prime }u^k]\times $$
$$\times \left(k_k^{}e^{ik(\tau ^{\prime \prime }\tau ^{})}+k_k^+e^{ik(\tau ^{\prime \prime }\tau ^{})}\right)exp[\frac{i}{c}(\mathrm{𝐤𝐯}^{\prime \prime })(\eta ^{\prime \prime }\tau ^{\prime \prime })+\frac{i}{c}(\mathrm{𝐤𝐯}^{})(\tau ^{}\eta ^{})]+$$
$$+_{\mathrm{}}^{\eta ^{\prime \prime }}\frac{d\tau ^{\prime \prime }}{u^{\prime \prime 0}}_{\mathrm{}}^{\tau ^{\prime \prime }}\frac{d\tau ^{}}{u^0}f_b(x^{})\frac{f_c(x^{\prime \prime })}{p_j^{\prime \prime }}[(u^{}u^{\prime \prime })\delta _j^ku_j^{}u^{\prime \prime k}]\times $$
$$\times (k_k^+e^{ik(\tau ^{\prime \prime }\tau ^{})}k_k^{}e^{ik(\tau ^{\prime \prime }\tau ^{})})exp[\frac{i}{c}(\mathrm{𝐤𝐯}^{\prime \prime })(\eta ^{\prime \prime }\tau ^{\prime \prime })+\frac{i}{c}(\mathrm{𝐤𝐯}^{})(\tau ^{}\eta ^{})]\}\times $$
$$\times h_l^{n(b)}(\eta ,\eta ^{},p^{},𝐤)\mathrm{\Omega }_{is}^{l(c)}(\eta ,\eta ^{\prime \prime },p^{\prime \prime },𝐤),$$
(80)
$$Q_{kis}^n=\underset{bc}{}(\frac{e_be_cn_bn_c(2\pi )^6}{4\pi ^2c})d^4p^{}d^4p^{\prime \prime }_{\mathrm{}}^\eta d\eta ^{}_{\mathrm{}}^\eta d\eta ^{\prime \prime }\frac{d^3k}{k}\times $$
$$\times \{_{\mathrm{}}^\eta ^{}\frac{d\tau ^{}}{u^0}_{\mathrm{}}^\tau ^{}\frac{d\tau ^{\prime \prime }}{u^{\prime \prime 0}}f_c(x^{\prime \prime })\frac{f_b(x^{})}{p_j^{}}[(u^{}u^{\prime \prime })\delta _j^ku_j^{\prime \prime }u^k]\times $$
$$\times \left(k_k^{}e^{ik(\tau ^{\prime \prime }\tau ^{})}+k_k^+e^{ik(\tau ^{\prime \prime }\tau ^{})}\right)exp[\frac{i}{c}(\mathrm{𝐤𝐯}^{\prime \prime })(\eta ^{\prime \prime }\tau ^{\prime \prime })+\frac{i}{c}(\mathrm{𝐤𝐯}^{})(\tau ^{}\eta ^{})]+$$
$$+_{\mathrm{}}^{\eta ^{\prime \prime }}\frac{d\tau ^{\prime \prime }}{u^{\prime \prime 0}}_{\mathrm{}}^{\tau ^{\prime \prime }}\frac{d\tau ^{}}{u^0}f_b(x^{})\frac{f_c(x^{\prime \prime })}{p_j^{\prime \prime }}[(u^{}u^{\prime \prime })\delta _j^ku_j^{}u^{\prime \prime k}]\times $$
$$\times (k_k^+e^{ik(\tau ^{\prime \prime }\tau ^{})}k_k^{}e^{ik(\tau ^{\prime \prime }\tau ^{})})exp[\frac{i}{c}(\mathrm{𝐤𝐯}^{\prime \prime })(\eta ^{\prime \prime }\tau ^{\prime \prime })+\frac{i}{c}(\mathrm{𝐤𝐯}^{})(\tau ^{}\eta ^{})]\}\times $$
$$\times \mathrm{\Omega }_{kl}^{n(b)}(\eta ,\eta ^{},p^{},𝐤)\mathrm{\Omega }_{is}^{l(c)}(\eta ,\eta ^{\prime \prime },p^{\prime \prime },𝐤),$$
(81)
$$\lambda _{ij}=\underset{bc}{}\chi (2\pi )e_be_cn_bn_cm_ccd^4p^{}d^4p^{\prime \prime }[\frac{1}{2}u_i^{\prime \prime }u_j^{\prime \prime }u^{\prime \prime k}u^{\prime \prime m}\frac{1}{4}g_{ij}u^{\prime \prime k}u^{\prime \prime m}$$
$$\frac{1}{2}u_i^{\prime \prime }u_j^{\prime \prime }g^{km}+\frac{1}{4}g_{ij}g^{km}+u_i^{\prime \prime }u^{\prime \prime k}\delta _j^m+u_j^{\prime \prime }u^{\prime \prime k}\delta _i^m\frac{1}{2}\delta _i^k\delta _j^m]\times $$
$$\times \frac{d^3k}{k}_{\mathrm{}}^\eta d\eta ^{}h_{km}^{(b)}(\eta ,\eta ^{},p^{},𝐤)\{_{\mathrm{}}^\eta ^{}\frac{d\tau ^{}}{u^0}_{\mathrm{}}^\tau ^{}\frac{d\tau ^{\prime \prime }}{u^{\prime \prime 0}}f_c(x^{\prime \prime })\frac{f_b(x^{})}{p_r^{}}\times $$
$$\times ((u^{}u^{\prime \prime })\delta _r^lu_r^{\prime \prime }u^l)(k_l^+e^{ik(\tau ^{\prime \prime }\tau ^{})}k_l^{}e^{ik(\tau ^{\prime \prime }\tau ^{})})\times $$
$$\times exp[\frac{i}{c}(\mathrm{𝐤𝐯}^{\prime \prime })(\eta \tau ^{\prime \prime })+\frac{i}{c}(\mathrm{𝐤𝐯}^{})(\tau ^{}\eta ^{})]+$$
$$+_{\mathrm{}}^\eta \frac{d\tau ^{\prime \prime }}{u^{\prime \prime 0}}_{\mathrm{}}^{\tau ^{\prime \prime }}\frac{d\tau ^{}}{u^0}f_b(x^{})\frac{f_c(x^{\prime \prime })}{p_r^{\prime \prime }}((u^{}u^{\prime \prime })\delta _r^lu_r^{}u^{\prime \prime l})\times $$
$$\times (k_l^+e^{ik(\tau ^{\prime \prime }\tau ^{})}k_l^{}e^{ik(\tau ^{\prime \prime }\tau ^{})})exp[\frac{i}{c}(\mathrm{𝐤𝐯}^{\prime \prime })(\eta \tau ^{\prime \prime })+\frac{i}{c}(\mathrm{𝐤𝐯}^{})(\tau ^{}\eta ^{})]\}.$$
(82)
$$h_m^i\omega ^{km}=\underset{bc}{}(\frac{e_be_cn_bn_c(2\pi )^6}{4\pi ^2c})\left(\frac{i\chi m_bc^2}{(2\pi )^3k}\right)\left(\frac{e_c}{4\pi ^2k}\right)d^4p^{}d^4p^{\prime \prime }\times $$
$$\times _{\mathrm{}}^\eta d\eta ^{}_{\mathrm{}}^\eta d\eta ^{\prime \prime }\frac{d^3k}{k}\{_{\mathrm{}}^\eta ^{}\frac{d\tau ^{}}{u^0}_{\mathrm{}}^\tau ^{}\frac{d\tau ^{\prime \prime }}{u^{\prime \prime 0}}f_c(x^{\prime \prime })\frac{f_b(x^{})}{p_l^{}}[(u^{}u^{\prime \prime })\delta _l^su_l^{\prime \prime }u^s]\times $$
$$\times \left(k_s^{}e^{ik(\tau ^{\prime \prime }\tau ^{})}+k_s^+e^{ik(\tau ^{\prime \prime }\tau ^{})}\right)exp[\frac{i}{c}(\mathrm{𝐤𝐯}^{\prime \prime })(\eta ^{\prime \prime }\tau ^{\prime \prime })+\frac{i}{c}(\mathrm{𝐤𝐯}^{})(\tau ^{}\eta ^{})]+$$
$$+_{\mathrm{}}^{\eta ^{\prime \prime }}\frac{d\tau ^{\prime \prime }}{u^{\prime \prime 0}}_{\mathrm{}}^{\tau ^{\prime \prime }}\frac{d\tau ^{}}{u^0}f_b(x^{})\frac{f_c(x^{\prime \prime })}{p_l^{\prime \prime }}[(u^{}u^{\prime \prime })\delta _l^su_l^{}u^{\prime \prime s}]\times $$
$$\times (k_s^+e^{ik(\tau ^{\prime \prime }\tau ^{})}k_s^{}e^{ik(\tau ^{\prime \prime }\tau ^{})})exp[\frac{i}{c}(\mathrm{𝐤𝐯}^{\prime \prime })(\eta ^{\prime \prime }\tau ^{\prime \prime })+\frac{i}{c}(\mathrm{𝐤𝐯}^{})(\tau ^{}\eta ^{})]\}\times $$
$$\times (u^iu_m^{}\frac{1}{2}\delta _m^i)(e^{ik(\eta ^{}\eta )}e^{ik(\eta ^{}\eta )})[u_{}^{\prime \prime }{}_{}{}^{m}((k_+^ke^{ik(\eta ^{\prime \prime }\eta )}k_{}^ke^{ik(\eta ^{\prime \prime }\eta )})$$
$$u_{}^{\prime \prime }{}_{}{}^{k}((k_+^me^{ik(\eta ^{\prime \prime }\eta )}k_{}^me^{ik(\eta ^{\prime \prime }\eta )})],$$
(83)
$$\omega ^{ik}_kh=\underset{bc}{}(\frac{e_be_cn_bn_c(2\pi )^6}{4\pi ^2c})\left(\frac{\chi m_bc^2}{(2\pi )^3k}\right)\left(\frac{e_c}{4\pi ^2k}\right)d^4p^{}d^4p^{\prime \prime }\times $$
$$\times _{\mathrm{}}^\eta d\eta ^{}_{\mathrm{}}^\eta d\eta ^{\prime \prime }\frac{d^3k}{k}\{_{\mathrm{}}^\eta ^{}\frac{d\tau ^{}}{u^0}_{\mathrm{}}^\tau ^{}\frac{d\tau ^{\prime \prime }}{u^{\prime \prime 0}}f_c(x^{\prime \prime })\frac{f_b(x^{})}{p_l^{}}[(u^{}u^{\prime \prime })\delta _l^su_l^{\prime \prime }u^s]\times $$
$$\times \left(k_s^{}e^{ik(\tau ^{\prime \prime }\tau ^{})}+k_s^+e^{ik(\tau ^{\prime \prime }\tau ^{})}\right)exp[\frac{i}{c}(\mathrm{𝐤𝐯}^{\prime \prime })(\eta ^{\prime \prime }\tau ^{\prime \prime })+\frac{i}{c}(\mathrm{𝐤𝐯}^{})(\tau ^{}\eta ^{})]+$$
$$+_{\mathrm{}}^{\eta ^{\prime \prime }}\frac{d\tau ^{\prime \prime }}{u^{\prime \prime 0}}_{\mathrm{}}^{\tau ^{\prime \prime }}\frac{d\tau ^{}}{u^0}f_b(x^{})\frac{f_c(x^{\prime \prime })}{p_l^{\prime \prime }}[(u^{}u^{\prime \prime })\delta _l^su_l^{}u^{\prime \prime s}]\times $$
$$\times (k_s^+e^{ik(\tau ^{\prime \prime }\tau ^{})}k_s^{}e^{ik(\tau ^{\prime \prime }\tau ^{})})exp[\frac{i}{c}(\mathrm{𝐤𝐯}^{\prime \prime })(\eta ^{\prime \prime }\tau ^{\prime \prime })+\frac{i}{c}(\mathrm{𝐤𝐯}^{})(\tau ^{}\eta ^{})]\}\times $$
$$\times (k_k^{}e^{ik(\eta ^{}\eta )}k_k^+e^{ik(\eta ^{}\eta )})[u_{}^{\prime \prime }{}_{}{}^{k}((k_+^ie^{ik(\eta ^{\prime \prime }\eta )}k_{}^ie^{ik(\eta ^{\prime \prime }\eta )})$$
$$u_{}^{\prime \prime }{}_{}{}^{i}((k_+^ke^{ik(\eta ^{\prime \prime }\eta )}k_{}^ke^{ik(\eta ^{\prime \prime }\eta )})],$$
(84)
$$\lambda ^i=\underset{bc}{}\frac{i\chi e_be_c^2m_bcn_bn_c}{2\pi }d^4p^{}d^4p^{\prime \prime }_{\mathrm{}}^\eta d\eta ^{}\frac{d^3k}{k^2}\times $$
$$\times u_{}^{\prime \prime }{}_{}{}^{i}(e^{ik(\eta ^{}\eta )}e^{ik(\eta ^{}\eta )})\{_{\mathrm{}}^\eta ^{}\frac{d\tau ^{}}{u^0}_{\mathrm{}}^\tau ^{}\frac{d\tau ^{\prime \prime }}{u_{}^{\prime \prime }{}_{}{}^{0}}f_c(x^{\prime \prime })\frac{f_b(x^{})}{p_l^{}}((u^{}u^{\prime \prime })\delta _l^s$$
$$u_{}^{\prime \prime }{}_{l}{}^{}u^s\left)\right(k_s^+e^{ik(\tau ^{\prime \prime }\tau ^{})}k_l^{}e^{ik(\tau ^{\prime \prime }\tau ^{})})exp[\frac{i}{c}(\mathrm{𝐤𝐯}^{\prime \prime })(\eta \tau ^{\prime \prime })+\frac{i}{c}(\mathrm{𝐤𝐯}^{})(\tau ^{}\eta ^{})]+$$
$$+_{\mathrm{}}^\eta \frac{d\tau ^{\prime \prime }}{u^{\prime \prime 0}}_{\mathrm{}}^{\tau ^{\prime \prime }}\frac{d\tau ^{}}{u^0}f_b(x^{})\frac{f_c(x^{\prime \prime })}{p_l^{\prime \prime }}((u^{}u^{\prime \prime })\delta _l^su_l^{}u^{\prime \prime s})(k_s^+e^{ik(\tau ^{\prime \prime }\tau ^{})}$$
$$k_s^{}e^{ik(\tau ^{\prime \prime }\tau ^{})})exp[\frac{i}{c}(\mathrm{𝐤𝐯}^{\prime \prime })(\eta \tau ^{\prime \prime })+\frac{i}{c}(\mathrm{𝐤𝐯}^{})(\tau ^{}\eta ^{})]\}.$$
(85)
To simplify (80) — (85) still further, we proceed as follows. We assume that the distribution function changes little inside the correlation region, so that in calculating the integrals in (80) — (85) we can ignore, in first approximation, the temporal coordinate dependence on $`f`$. We substitute the explicit expressions for $`h_{ij}^{(b)}`$ and $`\mathrm{\Omega }_{kj}^{i(b)}`$ (Eqs. (54) and (55)) into (80) \- (82) and evaluate the integrals with respect to $`\tau ^{},\tau ^{\prime \prime },\eta ^{},\eta ^{\prime \prime }`$ and $`𝐤`$. Then the expression for $`P_{is}^n`$ becomes
$$P_{is}^n=\underset{bc}{}\frac{\chi ^2e_be_cm_bm_cn_bn_cc^3}{8(\pi )^2}d^4p^{}d^4p^{\prime \prime }(u^nu_l^{}\frac{1}{2}\delta _l^n)[(u_i^{\prime \prime }u_s^{\prime \prime }\frac{1}{2}g_{is})g^{lf}$$
$$(u^{\prime \prime l}u_i^{\prime \prime }\frac{1}{2}\delta _i^l)\delta _s^f(u^{\prime \prime l}u_s^{\prime \prime }\frac{1}{2}\delta _s^l)\delta _i^f\left]\right\{[(u^{}u^{\prime \prime })\delta _j^mu_j^{\prime \prime }u^m]f_c(x^{\prime \prime })\frac{f_b(x^{})}{p_j^{}}K_{fm}^{(1)}(u^{},u^{\prime \prime })+$$
$$+[(u^{}u^{\prime \prime })\delta _j^mu_j^{}u^{\prime \prime m}]f_b(x^{})\frac{f_c(x^{\prime \prime })}{p_j^{\prime \prime }}K_{fm}^{(2)}(u^{},u^{\prime \prime })\}.$$
(86)
Here we have introduced the notation $`K_{fm}^{(1)}(u^{},u^{\prime \prime })`$ and $`K_{fm}^{(2)}(u^{},u^{\prime \prime })`$ for tensors that in locally Lorentzian reference frame, in which $`g_{ij}=\eta _{ij}`$ is the Minkowski tensor, have the following form:
$$K_{fm}^{(1)}(u^{},u^{\prime \prime })=\frac{i}{u^0u^{\prime \prime 0}}\frac{d^3k}{k^3}_{\mathrm{}}^\eta d\eta ^{}_{\mathrm{}}^\eta d\eta ^{\prime \prime }_{\mathrm{}}^\eta ^{}d\tau ^{}_{\mathrm{}}^\tau ^{}d\tau ^{\prime \prime }(e^{ik(\eta ^{}\eta )}$$
$$e^{ik(\eta ^{}\eta )}\left)\right(k_f^+e^{ik(\eta ^{\prime \prime }\eta )}k_f^{}e^{ik(\eta ^{\prime \prime }\eta )}\left)\right(k_m^+e^{ik(\tau ^{\prime \prime }\tau ^{})}k_m^{}e^{ik(\tau ^{\prime \prime }\tau ^{})})\times $$
$$\times exp[\frac{i}{c}(\mathrm{𝐤𝐯}^{\prime \prime })(\eta ^{\prime \prime }\tau ^{\prime \prime })+\frac{i}{c}(\mathrm{𝐤𝐯}^{})(\tau ^{}\eta ^{})],$$
$$K_{fm}^{(2)}(u^{},u^{\prime \prime })=\frac{i}{u^0u^{\prime \prime 0}}\frac{d^3k}{k^3}_{\mathrm{}}^\eta d\eta ^{}_{\mathrm{}}^\eta d\eta ^{\prime \prime }_{\mathrm{}}^{\eta ^{\prime \prime }}d\tau ^{\prime \prime }_{\mathrm{}}^{\tau ^{\prime \prime }}d\tau ^{}(e^{ik(\eta ^{}\eta )}$$
$$e^{ik(\eta ^{}\eta )}\left)\right(k_f^+e^{ik(\eta ^{\prime \prime }\eta )}k_f^{}e^{ik(\eta ^{\prime \prime }\eta )}\left)\right(k_m^+e^{ik(\tau ^{\prime \prime }\tau ^{})}k_m^{}e^{ik(\tau ^{\prime \prime }\tau ^{})})\times $$
$$\times exp[\frac{i}{c}(\mathrm{𝐤𝐯}^{\prime \prime })(\eta ^{\prime \prime }\tau ^{\prime \prime })+\frac{i}{c}(\mathrm{𝐤𝐯}^{})(\tau ^{}\eta ^{})],$$
After carrying out the integrals with respect $`\tau ^{}`$, $`\tau ^{\prime \prime }`$, $`\eta ^{}`$, and $`\eta ^{\prime \prime }`$ these expressions take the form:
$$K_{fm}^{(1)}(u^{},u^{\prime \prime })=\frac{2\pi c^5}{u^0u^{\prime \prime 0}}\frac{d^3k}{k^2}\delta (\mathrm{𝐤𝐯}^{\prime \prime }\mathrm{𝐤𝐯}^{})\{\frac{k_f^+k_m^+}{(kc\mathrm{𝐤𝐯}^{\prime \prime })(kc+\mathrm{𝐤𝐯}^{\prime \prime })^3}+$$
$$+\frac{k_f^+k_m^{}+k_f^{}k_m^+}{(kc\mathrm{𝐤𝐯}^{\prime \prime })^2(kc+\mathrm{𝐤𝐯}^{\prime \prime })^2}+\frac{k_f^{}k_m^{}}{(kc\mathrm{𝐤𝐯}^{\prime \prime })^3(kc+\mathrm{𝐤𝐯}^{\prime \prime })}\}=K_{fm}(u^{},u^{\prime \prime }),$$
(87)
$$K_{fm}^{(2)}(u^{},u^{\prime \prime })=K_{fm}^{(1)}(u^{},u^{\prime \prime })=K_{fm}(u^{},u^{\prime \prime })$$
(88)
The above equalities hold only in a locally Lorentzian reference frame. To obtain covariant expressions for the tensors $`K_{fm}^{(1)}(u^{},u^{\prime \prime })`$ and $`K_{fm}^{(2)}(u^{},u^{\prime \prime })`$, we take into account the following fact. The quantities $`K_{fm}^{(1)}(u^{},u^{\prime \prime })`$ and $`K_{fm}^{(2)}(u^{},u^{\prime \prime })`$ appeared in (86) after the correlation function $`g_{ab}(x^{},x^{\prime \prime })`$ was substituted to (64) and result was integrated with respect to $`𝐪^{},𝐪^{\prime \prime },𝐤^{}`$ and $`𝐤^{\prime \prime }`$. But the expression (79) for two - particle correlation function is a sum of two terms, which differ in that primed quantities referring to particles of species $`\mathrm{"}a\mathrm{"}`$ are replaced by double - primed quantities referring to particles of species $`\mathrm{"}b\mathrm{"}`$, and vice versa. It is after these terms were integrated with respect $`𝐪^{},𝐪^{\prime \prime },𝐤^{}`$ and $`𝐤^{\prime \prime }`$ that $`K_{fm}^{(1)}(u^{},u^{\prime \prime })`$ and $`K_{fm}^{(2)}(u^{},u^{\prime \prime })`$ appeared in (86). Obviously, the both must be calculated in the same reference frame, for which it is convenient to take the center - of - mass reference frame, in which
$$𝐯^{}=𝐯,𝐯^{\prime \prime }=𝐯,u^0=u^{\prime \prime 0}=1/\sqrt{1v^2/c^2}=u^0$$
. In this reference frame
$$K_{00}=K_{0\alpha }=0,K_{\alpha \beta }=\frac{2\pi ^2c}{vu_0^2k_{min}^2}\left(\delta _{\alpha \beta }\frac{v_\alpha v_\beta }{v^2}\right)$$
(89)
Here $`v=\sqrt{v_1^2+v_2^2+v_3^2}`$ , where $`v_\alpha =v^\alpha =u^\alpha /u^0`$ are spatial components of the vector $`𝐯`$.
A covariant generalization of (89) has the form
$$K_{ij}(u^{},u^{\prime \prime })=\frac{4\pi ^2}{k_{min}^2[(u^{}u^{\prime \prime })^21]^{3/2}}\{[(u^{}u^{\prime \prime })^21]g_{ij}$$
$$u_i^{}u_j^{}u_i^{\prime \prime }u_j^{\prime \prime }+(u^{}u^{\prime \prime })(u_i^{}u_j^{\prime \prime }+u_i^{\prime \prime }u_j^{})\}$$
(90)
The expressions for $`K_{fm}^{(1)}(u^{},u^{\prime \prime })`$ and $`K_{fm}^{(2)}(u^{},u^{\prime \prime })`$ diverge as $`k0`$, i.e., for large impact parameters. The reason is that we integrate over an infinite region, while actually we should integrate only over the correlation region, where the metric is assumed to vary only weakly. This difficulty is resolved, as well as in the case of kinetic equation deriving, by introducing a cutoff procedure in the divergent integral
$$_0^{\mathrm{}}\frac{dk}{k^3}$$
.
We set the lower integration limit to $`k_{min}=1/r_{max}`$, rather than zero, where $`r_{max}`$ is the size of the correlation region (the correlation radius). Then the above integral assumes the value $`1/2k_{min}^2=(1/2)r_{max}^2`$.
As the experience of deriving the relativistic kinetic equation (refer to. , , , ) shows, more thorough investigations suggest that the integrals become convergent as $`r\mathrm{}`$, with the contribution from the region where $`r>r_{max}`$ being infinitesimal. In Ref. , there are estimates for $`r_{max}`$ in the case where the average metric $`g_{ij}`$ is the metric of isotropic cosmological model and in the case of gravitational interaction of particles.
In the case of electromagnetical interaction of particles the parameter $`k_{min}`$ is equal to $`\frac{1}{r_D}`$, where $`r_D`$ a radius of Debit, since electromagnetic interactions in the plasma are shielded under $`r>r_D`$.
The tensor (90) possesses the following properties:
$$K_{ij}(u^{},u^{\prime \prime })=K_{ij}(u^{\prime \prime },u^{});K_{ij}u^i=K_{ij}u^{\prime \prime i}=0;K_{ij}=K_{ji}.$$
(91)
Because of this the expression for $`P_{is}^n`$ simplifies considerably. The macroscopic Einstein equations incorporate not $`P_{is}^n`$, but the tensor $`\phi _{ij}^k=(1/2)(\delta _n^k\delta _j^s\delta _j^k\delta _n^s)P_{is}^n`$. The expression for this tensor can be written as follows:
$$\phi _{ij}^k=(1/2)(\delta _n^k\delta _j^s\delta _j^k\delta _n^s)P_{is}^n=\underset{bc}{}\frac{\chi ^2e_be_cm_bm_cn_bn_cc^3}{16(\pi )^2}\frac{d^4p^{}}{\sqrt{(g)}}\frac{d^4p^{\prime \prime }}{\sqrt{(g)}}[\frac{1}{2}g^{fk}u_i^{\prime \prime }u_j^{\prime \prime }+$$
$$+u^k(u^{}u^{\prime \prime })(\delta _j^fu_i^{\prime \prime }+\delta _i^fu_j^{\prime \prime })](u^{}u^{\prime \prime })K_{fr}(u^{},u^{\prime \prime })(f_c(x^{\prime \prime })\frac{f_b(x^{})}{p_r^{}}f_b(x^{})\frac{f_c(x^{\prime \prime })}{p_r^{\prime \prime }})$$
(92)
Note that
$$g^{ij}\phi _{ij}^k=0,\phi _{ij}^i=0,\phi _{ij}^k=\phi _{ji}^k.$$
(93)
Reasoning along similar lines, we can simplify the expression for the tensor $`\mu _{ij}`$, $`\phi _{ij}`$, $`\mu _i`$, which assumes the following form:
$$\mu _{ij}=\underset{bc}{}\frac{\chi ^2e_be_cm_bm_cn_bn_cc^3}{16\pi ^3}\frac{d^4p^{}}{\sqrt{(g)}}\frac{d^4p^{\prime \prime }}{\sqrt{(g)}}[(z^2+\frac{1}{2})(u_i^{\prime \prime }u_j^{\prime \prime }+$$
$$+u_i^{}u_j^{})g^{qr}+(z^2\frac{1}{2})g_{ij}g^{qr}2z(u_i^{}u_j^{\prime \prime }+u_i^{\prime \prime }u_j^{})g^{qr}(z^2\frac{1}{2})(\delta _i^q\delta _j^r+\delta _j^q\delta _i^r)]\times $$
$$\times \left(z\delta _f^mu_f^{\prime \prime }u^m\right)J_{rqm}(u^{},u^{\prime \prime })f_c(x^{\prime \prime })\frac{f_b(x^{})}{p_f^{}},$$
(94)
$$\phi ^{ik}=\underset{bc}{}\frac{\chi e_be_c^2n_bn_cc}{2\pi }\frac{d^4p^{}}{\sqrt{(g)}}\frac{d^4p^{\prime \prime }}{\sqrt{(g)}}(u^{}u^{\prime \prime })K_{fl}(u^{},u^{\prime \prime })\times $$
$$\times [(u^{}u^{\prime \prime })(u^ig^{kf}u^kg^{if})(u^{\prime \prime i}g^{kf}u^{\prime \prime k}g^{if})]\left(f_c(x^{\prime \prime })\frac{f_b(x^{})}{p_l^{}}f_b(x^{})\frac{f_c(x^{\prime \prime })}{p_l^{\prime \prime }}\right),$$
(95)
$$\mu ^i=\underset{bc}{}\frac{\chi e_b^2e_cm_cn_bn_cc}{4\pi }\frac{d^4p^{}}{\sqrt{(g)}}\frac{d^4p^{\prime \prime }}{\sqrt{(g)}}[((u^{}u^{\prime \prime })\delta _l^su_l^{\prime \prime }u^s)]\times $$
$$u^{\prime \prime k}J_{ks}^i(u^{},u^{\prime \prime })f_c(x^{\prime \prime })\frac{f_b(x^{})}{p_l^{}}$$
(96)
Here $`z=(u^{}u^{\prime \prime })`$.
In (94), (96) we have introduced the notation $`J_{rqm}(u^{},u^{\prime \prime })`$ for tensor that in locally Lorentzian reference frame have the form
$$J_{lmn}(u^{},u^{\prime \prime })=\frac{1}{u^0u^{\prime \prime 0}}\frac{d^3k}{k^3}_{\mathrm{}}^\eta d\eta ^{}_{\mathrm{}}^\eta d\eta ^{\prime \prime }_{\mathrm{}}^\eta ^{}d\tau ^{}_{\mathrm{}}^\tau ^{}d\tau ^{\prime \prime }(k_l^+e^{ik(\eta ^{}\eta )}$$
$$k_l^{}e^{ik(\eta ^{}\eta )}\left)\right(k_m^+e^{ik(\eta ^{\prime \prime }\eta )}k_m^{}e^{ik(\eta ^{\prime \prime }\eta )}\left)\right(k_n^{}e^{ik(\tau ^{\prime \prime }\tau ^{})}$$
$$k_n^+e^{ik(\tau ^{\prime \prime }\tau ^{})})exp[\frac{i}{c}(\mathrm{𝐤𝐯}^{\prime \prime })(\eta ^{\prime \prime }\tau ^{\prime \prime })+\frac{i}{c}(\mathrm{𝐤𝐯}^{})(\tau ^{}\eta ^{})].$$
After evaluating the integrals with respect $`\eta ^{},\eta ^{\prime \prime },\tau ^{}`$ and $`\tau ^{\prime \prime }`$, we get
$$J_{lmn}(u^{},u^{\prime \prime })=\frac{c^4}{u^0u^{\prime \prime 0}}\frac{d^3k}{k^3}\frac{V.p.}{(\mathrm{𝐤𝐯}^{\prime \prime }\mathrm{𝐤𝐯}^{})}\{\frac{k_l^+k_m^+k_n^+}{(kc+\mathrm{𝐤𝐯}^{\prime \prime })^3}+$$
$$+\frac{k_l^+k_m^+k_n^{}+k_l^+k_m^{}k_n^++k_l^{}k_m^+k_n^+}{(kc+\mathrm{𝐤𝐯}^{\prime \prime })^2(kc\mathrm{𝐤𝐯}^{\prime \prime })}+\frac{k_l^+k_m^{}k_n^{}+k_l^{}k_m^+k_n^{}+k_l^{}k_m^{}k_n^+}{(kc+\mathrm{𝐤𝐯}^{\prime \prime })(kc\mathrm{𝐤𝐯}^{\prime \prime })^2}+$$
$$+\frac{k_l^{}k_m^{}k_n^{}}{(kc\mathrm{𝐤𝐯}^{\prime \prime })^3}\}$$
(97)
The symbol $`V.p.`$ indicates that the integral is calculated as a principal value.
Just as in the previous case, we specify (97) in the center-of-mass reference frame, where
$$𝐯^{}=𝐯,𝐯^{\prime \prime }=𝐯,u^0=u^{\prime \prime 0}=1/\sqrt{1v^2/c^2}=u^0$$
. In this reference frame the components of $`J_{lmn}(u^{},u^{\prime \prime })`$ have the following form (the spatial indexes of three-dimensional velocity $`v^\alpha `$ are lowered by using the tree-dimensional Kronecker symbol $`\delta _{\alpha \beta }`$)
$$J_{000}=\alpha (v)\frac{v^2}{c^2},J_{00\alpha }=\alpha (v)\frac{v_\alpha }{c},J_{0\alpha \beta }=\alpha (v)\delta _{\alpha \beta }+\beta (v)\left(\delta _{\alpha \beta }\frac{v_\alpha v_\beta }{v^2}\right),$$
(98)
$$J_{\alpha \beta \gamma }=\frac{c^2}{v^2}\alpha (v)\left[\delta _{\alpha \beta }\frac{v_\gamma }{c}+\delta _{\alpha \gamma }\frac{v_\beta }{c}+\delta _{\beta \gamma }\frac{v_\alpha }{c}2\frac{v_\alpha v_\beta v_\gamma }{cv^2}\right]+$$
$$+\frac{c^2}{v^2}\beta (v)\left[\left(\delta _{\alpha \beta }\frac{v_\alpha v_\beta }{v^2}\right)\frac{v_\gamma }{c}+\left(\delta _{\alpha \gamma }\frac{v_\alpha v_\gamma }{v^2}\right)\frac{v_\beta }{c}+\left(\delta _{\beta \gamma }\frac{v_\beta v_\gamma }{v^2}\right)\frac{v_\alpha }{c}\right].$$
(99)
The function $`\alpha `$ and $`\beta `$ in (98) and (99) depend on the velocity $`v=\sqrt{v_1^2+v_2^2+v_3^2}`$ only and have the explicit form
$$\alpha =\frac{\pi c^3}{u_0^2v^3k_{min}}\left[\frac{2\frac{v}{c}\left(1+\frac{v^2}{c^2}\right)}{\left(1\frac{v^2}{c^2}\right)^2}+\mathrm{ln}\left(\frac{1\frac{v}{c}}{1+\frac{v}{c}}\right)\right],$$
(100)
$$\beta =\frac{\pi c^3}{2u_0^2v^3k_{min}}\left[\frac{2\frac{v}{c}\left(32\frac{v^2}{c^2}+3\frac{v^4}{c^4}\right)}{\left(1\frac{v^2}{c^2}\right)^2}+3\left(1+\frac{v^2}{c^2}\right)\mathrm{ln}\left(\frac{1\frac{v}{c}}{1+\frac{v}{c}}\right)\right].$$
(101)
Here we introduced the following notation for the integral
$$\frac{1}{k_{min}}=_{k_{min}}^{\mathrm{}}\frac{dk}{k^2}.$$
We set the lower integration limit to $`k_{min}=1/r_D`$.
A covariant generalization of this results, which were obtained in the locally Lorentzian center-of-mass reference frame, to arbitrary reference frames has the form
$$J_{ijk}(u^{},u^{\prime \prime })=A[(g_{ij}u_k^{}+g_{ik}u_j^{}+g_{jk}u_i^{})z(g_{ij}u_k^{\prime \prime }+g_{ik}u_j^{\prime \prime }+g_{jk}u_i^{\prime \prime })$$
$$(u_i^{}u_j^{\prime \prime }u_k^{\prime \prime }+u_i^{\prime \prime }u_j^{}u_k^{\prime \prime }+u_i^{\prime \prime }u_j^{\prime \prime }u_k^{})+3zu_i^{\prime \prime }u_j^{\prime \prime }u_k^{\prime \prime }]+$$
$$+C[u_i^{}u_j^{}u_k^{}z(u_i^{}u_j^{}u_k^{\prime \prime }+u_i^{}u_j^{\prime \prime }u_k^{}+u_i^{\prime \prime }u_j^{}u_k^{})+$$
$$+z^2(u_i^{}u_j^{\prime \prime }u_k^{\prime \prime }+u_i^{\prime \prime }u_j^{}u_k^{\prime \prime }+u_i^{\prime \prime }u_j^{\prime \prime }u_k^{})z^3u_i^{\prime \prime }u_j^{\prime \prime }u_k^{\prime \prime }],$$
(102)
where $`z=(u^{}u^{\prime \prime })=(u^iu_i^{\prime \prime })`$ ,
$$A=\frac{2\pi \sqrt{2}}{k_{min}}\left[\frac{(z2)}{(z1)^2(z+1)^{1/2}}+\frac{(2z1)}{(z+1)(z1)^{5/2}}\mathrm{ln}\left(z+\sqrt{z^21}\right)\right],$$
(103)
$$C=\frac{2\pi \sqrt{2}}{k_{min}}\left[\frac{(z6)}{(z1)^3(z+1)^{3/2}}+\frac{(6z1)}{(z+1)^2(z1)^{7/2}}\mathrm{ln}\left(z+\sqrt{z^21}\right)\right].$$
(104)
The tensor $`J_{ijk}(u^{},u^{\prime \prime })`$ satisfies the identity
$$J_{ijk}(u^{},u^{\prime \prime })u^{\prime \prime k}=0.$$
(105)
Note that the tensor $`\mu _{ij}`$ is traceless:
$$g^{ij}\mu _{ij}=0.$$
(106)
Let us now simplify the tensor $`T_{ij}^{(r)}`$ (see. (35). Substitution of (62) to (35) yields the expression for $`T_{ij}^{(r)}`$. In view (70) - (73) we have (if $`n_a1`$):
$$T_{ij}^{(r)}=\frac{1}{4\pi }(g_{it}g_{js}+\frac{1}{4}g_{ij}g_{ts})\underset{bc}{}d^4p^{}d^4p^{\prime \prime }d^3q^{}d^3q^{\prime \prime }\times $$
$$\times _{\mathrm{}}^\eta d\eta ^{}_{\mathrm{}}^\eta d\eta ^{\prime \prime }d^3k^{}d^3k^{\prime \prime }e^{i𝐤^{}(𝐪𝐪^{})}e^{i𝐤^{\prime \prime }(𝐪𝐪^{\prime \prime })}\times $$
$$\times \omega _{.l}^{t(b)}(\eta ,\eta ^{},p^{},𝐤^{})\omega ^{(c)sl}(\eta ,\eta ^{\prime \prime },p^{\prime \prime },𝐤^{\prime \prime })n_bn_cg_{bc}(x^{},x^{\prime \prime }),$$
(107)
Here we can not neglect by $`g_{bc}^{(gr)}(x^{},x^{\prime \prime })`$ in expression (76) for $`g_{bc}(x^{},x^{\prime \prime })`$.
That is why
$$T_{ij}^{(r)}=\tau _{ij}^{(r)}+\tau _{ij}^{(gr)}.$$
(108)
One can get $`\tau _{ij}^{(r)}`$ if replace $`g_{bc}(x^{},x^{\prime \prime })`$ by $`g_{bc}^{(el)}(x^{},x^{\prime \prime })`$ (see (79)) in expression (107) and replace $`g_{bc}(x^{},x^{\prime \prime })`$ by $`g_{bc}^{(gr)}(x^{},x^{\prime \prime })`$ (see (77) to obtain the expression for $`\tau _{ij}^{(gr)}`$.
Substitution (77) and (79) to (107) yields the folloing expressions for $`\tau _{ij}^{(r)}`$ and $`\tau _{ij}^{(gr)}`$:
$$\tau _{ij}^{(r)}=\underset{bc}{}\frac{e_b^2e_c^2n_bn_c}{2(2\pi )^4c}\frac{d^4p^{}}{\sqrt{(g)}}\frac{d^4p^{\prime \prime }}{\sqrt{(g)}}[2z\delta _i^r\delta _j^f+zg_{ij}g^{rf}+$$
$$+(\delta _i^fu_{}^{\prime \prime }{}_{j}{}^{}+\delta _j^fu_{}^{\prime \prime }{}_{i}{}^{})u_{}^{}{}_{}{}^{r}g^{rf}(u_{}^{}{}_{i}{}^{}u_{}^{\prime \prime }{}_{j}{}^{}+u_{}^{\prime \prime }{}_{i}{}^{}u_{}^{}{}_{j}{}^{})]\times $$
$$\times \left(z\delta _n^su_n^{\prime \prime }u^s\right)J_{rfs}^{(el)}(u^{},u^{\prime \prime })f_c(x^{\prime \prime })\frac{f_b(x^{})}{p_n^{}}.$$
(109)
$$\tau _{ij}^{(gr)}=\underset{bc}{}\frac{\chi e_be_cm_bm_cn_bn_cc^3}{16\pi ^2}\frac{d^4p^{}}{\sqrt{(g)}}\frac{d^4p^{\prime \prime }}{\sqrt{(g)}}[2z\delta _i^p\delta _j^qzg_{ij}g^{pq}$$
$$(\delta _i^qu_{}^{\prime \prime }{}_{j}{}^{}+\delta _j^qu_{}^{\prime \prime }{}_{i}{}^{})u_{}^{}{}_{}{}^{p}+g^{pq}(u_{}^{}{}_{i}{}^{}u_{}^{\prime \prime }{}_{j}{}^{}+u_{}^{\prime \prime }{}_{i}{}^{}u_{}^{}{}_{j}{}^{})]J^{(gr)}_{pqf}(u^{},u^{\prime \prime })f_c(x^{\prime \prime })\times $$
$$\times \frac{}{p_n^{}}.\left\{f_b(x^{})[(z^2\frac{1}{2})\delta _n^f+(z^2+\frac{1}{2})u_{}^{}{}_{n}{}^{}u_{}^{}{}_{}{}^{f}2zu_{}^{\prime \prime }{}_{n}{}^{}u_{}^{}{}_{}{}^{f}]\right\}.$$
(110)
Note that the tensors $`\tau _{ij}^{(gr)}`$ and $`\tau _{ij}^{(r)}`$ are traceless:
$$g^{ij}\tau _{ij}^{(gr)}=0,g^{ij}\tau _{ij}^{(r)}=0$$
(111)
In (109) and (110) the tensors $`J_{rpq}^{(el)}(u^{},u^{\prime \prime })`$ and $`J_{rpq}^{(gr)}(u^{},u^{\prime \prime })`$ have the form (102), where $`A`$ and $`B`$ have the forms (103) and (104) respectively. But in the expression for $`J_{rpq}^{(el)}(u^{},u^{\prime \prime })`$ we must put $`k_{min}=1/r_D`$, where $`r_D`$ is the radius of Debit, since the electromagnetic interaction in plasma are shielded under $`r>r_D`$. In the expression for $`J_{rpq}^{(gr)}(u^{},u^{\prime \prime })`$ we must put $`k_{min}=1/r_g`$, where $`r_g`$ is the radius of correlation for gravitational interaction. As the experience of deriving the relativistic kinetic equation (refer to. , , , ) shows, more thorough investigations suggest than the integrals become convergent as $`r\mathrm{}`$, with the contribution from the region where $`r>r_g`$ being infinitesimal. In Ref. , there are estimates for $`r_g`$ in the case where the average metric $`g_{ij}`$ is the metric of isotropic cosmological model and in the case of gravitational interaction of particles.
4. Macroscopic system of Einstein and Maxwell equations for relativistic plasma
As a result were obtained the macroscopic Einstein and Maxwell equations in relativistic plasma. They have the forms:
$$G_{ij}+_k\phi _{ij}^k+\mu _{ij}\chi \tau _{ij}^{(gr)}=\chi T_{ij},$$
(112)
$$_kF^{ik}+_k\phi ^{ik}+\mu ^i=\frac{4\pi }{c}J^i.$$
(113)
Here $`G_{ij}`$ is the Einstein’s tensor of the Riamannian space with macroscopic metric $`g_{ij}`$, $`F^{ik}`$ is the macroscopic tensor of electromagnetic field (Maxwell’s tensor), $`J^i`$ is macroscopic current vector, $`T_{ij}`$ is the macroscopic energy-momentum tensor. The last is the sum of macroscopic energy-momentum tensor of medium $`T_{ij}^{(m)}`$ (34), energy momentum tensor of macroscopic electromagnetical field $`T_{ij}^{(el)}`$ (29) and macroscopic energy-momentum tensor $`\tau _{ij}^{(r)}`$ (109) of electromagnetical radiation in plasma. (In cosmological plasma in the last case one should say about the energy - momentum thensor of relict radiation.)
The Einstein equations of the gravitational field for continium media, obtained here, differ from the classical Einstein equations by the presence of additional terms $`_k\phi _{ij}^k`$, $`\mu _{ij}`$ and $`\chi \tau _{ij}^{(gr)}`$ in the left-hand side. It caused by particle interaction. The forms of this tensors are (92), (94) and (110). The third term, $`(\chi \tau _{ij}^{(gr)})`$ is the addition to macroscopic energy - momentum tensor of electromagnetic radiation, caused by gravitational interaction which multiplying on $`\chi `$ and moving from the right-hand side of macroscopic equations to the left-hand side.
The macroscopic Maxwell equations differ from the classical Maxwell equations by the presence of additional terms $`_k\phi ^{ki}+\mu ^i`$. The additional terms in Maxwell equations caused by particle interaction and by effects of general relativity.
The tensors $`_k\phi _{ij}^k`$, $`\mu _{ij}`$, $`\tau _{ij}^{(gr)}`$, $`_k\phi ^{ki}`$ and $`\mu ^i`$ are expressed in (92) — (96) and (110) in terms of one-particle distribution function $`f_b`$ specified in the eight-dimensional phase spase in which all four components of momentum are independent. The transition to the seven-dimensional distribution function $`F_a(q^i,p_\alpha )`$ is made according to the formula
$$n_af_a(q^i,p_j)=F_a(q^i,p_\alpha )\delta (\sqrt{g^{lm}p_lp_m}m_ac).$$
(114)
Here the function $`F_a`$ depends on the spatial components of momentum only. Greek indexes are used to denote spartial components.
By integrating (92) — (96) and (110) with respect to $`p_{}^{}{}_{0}{}^{}`$ and $`p_{}^{\prime \prime }{}_{0}{}^{}`$ we can whrite down the tensors $`_k\phi _{ij}^k`$, $`\mu _{ij}`$, $`\tau _{ij}^{(gr)}`$, $`_k\phi ^{ki}`$ and $`\mu ^i`$ as
$$\phi _{ij}^k=\underset{bc}{}\frac{\chi ^2e_be_cm_b^2m_c^2c^5}{16(\pi )^2}\frac{d^3p^{}}{p_{}^{}{}_{}{}^{0}\sqrt{(g)}}\frac{d^3p^{\prime \prime }}{p_{}^{\prime \prime }{}_{}{}^{0}\sqrt{(g)}}[\frac{1}{2}g^{fk}u_i^{\prime \prime }u_j^{\prime \prime }+$$
$$+u^k(u^{}u^{\prime \prime })(\delta _j^fu_i^{\prime \prime }+\delta _i^fu_j^{\prime \prime })](u^{}u^{\prime \prime })K_{f\alpha }(u^{},u^{\prime \prime })(F_c(x^{\prime \prime })\frac{F_b(x^{})}{p_\alpha ^{}}F_b(x^{})\frac{F_c(x^{\prime \prime })}{p_\alpha ^{\prime \prime }})$$
(115)
$$\mu _{ij}=\underset{bc}{}\frac{\chi ^2e_be_cm_b^2m_c^2c^5}{16\pi ^2}\frac{d^3p^{}}{p_{}^{}{}_{}{}^{0}\sqrt{(g)}}\frac{d^3p^{\prime \prime }}{p_{}^{\prime \prime }{}_{}{}^{0}\sqrt{(g)}}[(z^2+\frac{1}{2})(u_i^{\prime \prime }u_j^{\prime \prime }+$$
$$+u_i^{}u_j^{})g^{qr}+(z^2\frac{1}{2})g_{ij}g^{qr}2z(u_i^{}u_j^{\prime \prime }+u_i^{\prime \prime }u_j^{})g^{qr}(z^2\frac{1}{2})(\delta _i^q\delta _j^r+\delta _j^q\delta _i^r)]\times $$
$$\times \left(z\delta _\alpha ^mu_\alpha ^{\prime \prime }u^m\right)J_{rqm}(u^{},u^{\prime \prime })F_c(x^{\prime \prime })\frac{F_b(x^{})}{p_\alpha ^{}}.$$
(116)
$$\phi ^{ik}=\underset{bc}{}\frac{\chi e_be_c^2m_bm_cc^3}{2\pi }\frac{d^3p^{}}{p_{}^{}{}_{}{}^{0}\sqrt{(g)}}\frac{d^3p^{\prime \prime }}{p_{}^{\prime \prime }{}_{}{}^{0}\sqrt{(g)}}(u^{}u^{\prime \prime })K_{f\alpha }(u^{},u^{\prime \prime })\times $$
$$\times [(u^{}u^{\prime \prime })(u^ig^{kf}u^kg^{if})(u^{\prime \prime i}g^{kf}u^{\prime \prime k}g^{if})]\left(F_c(x^{\prime \prime })\frac{F_b(x^{})}{p_\alpha ^{}}F_b(x^{})\frac{F_c(x^{\prime \prime })}{p_\alpha ^{\prime \prime }}\right),$$
(117)
$$\mu ^i=\underset{bc}{}\frac{\chi e_b^2e_cm_bm_c^2c^3}{4\pi }\frac{d^3p^{}}{p_{}^{}{}_{}{}^{0}\sqrt{(g)}}\frac{d^3p^{\prime \prime }}{p_{}^{\prime \prime }{}_{}{}^{0}\sqrt{(g)}}[((u^{}u^{\prime \prime })\delta _\alpha ^su_\alpha ^{\prime \prime }u^s)]\times $$
$$u^{\prime \prime k}J_{ks}^i(u^{},u^{\prime \prime })F_c(x^{\prime \prime })\frac{F_b(x^{})}{p_\alpha ^{}}$$
(118)
$$\tau _{ij}^{(gr)}=\underset{bc}{}\frac{\chi e_be_cm_b^2m_c^2c^5}{16\pi ^2}\frac{d^3p^{}}{p_{}^{}{}_{}{}^{0}\sqrt{(g)}}\frac{d^3p^{\prime \prime }}{p_{}^{\prime \prime }{}_{}{}^{0}\sqrt{(g)}}[2z\delta _i^p\delta _j^qzg_{ij}g^{pq}$$
$$(\delta _i^qu_{}^{\prime \prime }{}_{j}{}^{}+\delta _j^qu_{}^{\prime \prime }{}_{i}{}^{})u_{}^{}{}_{}{}^{p}+g^{pq}(u_{}^{}{}_{i}{}^{}u_{}^{\prime \prime }{}_{j}{}^{}+u_{}^{\prime \prime }{}_{i}{}^{}u_{}^{}{}_{j}{}^{})]J^{(gr)}_{pqf}(u^{},u^{\prime \prime })F_c(x^{\prime \prime })\times $$
$$\times \frac{}{p_n^{}}.\left\{F_b(x^{})[(z^2\frac{1}{2})\delta _n^f+(z^2+\frac{1}{2})u_{}^{}{}_{n}{}^{}u_{}^{}{}_{}{}^{f}2zu_{}^{\prime \prime }{}_{n}{}^{}u_{}^{}{}_{}{}^{f}]\right\}.$$
(119)
Here
$$\frac{d^3p^{}}{p^0\sqrt{(g)}}and\frac{d^3p^{\prime \prime }}{p^{\prime \prime 0}\sqrt{(g)}}$$
are the invariant volume elements in tree - dimensional momentum spase of particles spesies ”b” and ”c” respectively.
The greek index $`\alpha `$ in (115) - (118) takes the values 1,2 and 3 only (the spartial index). The derivative with respect to $`p_n^{}`$ in (119) should by calculated as all four components of momentum are independent. The dependence of $`p_0^{}`$ on $`p_\alpha ^{}`$ is taken into account after differentiation with respect $`p_n^{}`$ is completed only.
In (118) and (119) the tensors $`J_{rpq}^{(el)}(u^{},u^{\prime \prime })`$ and $`J_{rpq}^{(gr)}(u^{},u^{\prime \prime })`$ have the form (102), where $`A`$ and $`B`$ have the forms (103) and (104) respectively. But in expression for $`J_{rpq}^{(el)}(u^{},u^{\prime \prime })`$ we must put $`k_{min}=1/r_D`$, where $`r_D`$ is the radius of Debit, since the electromagnetic interaction in plasma are shielded under $`r>r_D`$. In the expression for $`J_{rpq}^{(gr)}(u^{},u^{\prime \prime })`$ we must put $`k_{min}=1/r_g`$, where $`r_g`$ is the radius of correlation for gravitational interaction. As the experience of deriving the relativistic kinetic equation (refer to. , , , ) shows, more thorough investigations suggest than the integrals become convergent as $`r\mathrm{}`$, with the contribution from the region where $`r>r_g`$ being infinitesimal. In Ref. , there are estimates for $`r_g`$ in the case where the average metric $`g_{ij}`$ is the metric of isotropic cosmological model and in the case of gravitational interaction of particles.
The tensors $`\phi _{ij}^k`$ , $`\mu _{ij}`$, $`\tau _{ij}^{(gr)}`$ $`\mu ^i`$ must obey the additional conditions
$$g^{lj}_l\left(_k\phi _{ij}^k+\mu _{ij}\chi \tau _{ij}^{(gr)}\right)=0,$$
(120)
$$_i\mu ^i=0,$$
(121)
since the divergenses of $`G_{ij}`$, $`T_{ij}`$, $`_kF^{ik}`$, $`_k\phi ^{ik}`$, $`J^i`$ vanish.
Equations (120), (121) impose some restrictions on the parametres $`r_D`$ and $`r_g`$ dependence on the coordinates and the relative velosity of particles. The latter can be expressed via of $`z=(u^{}u^{\prime \prime })`$.
The macroscopic energy-momentum tensor $`T_{ij}^{(m)}`$ of medium and the current vector $`J^i`$ can by alsow written in terms of one-particle distribution function as follows:
$$T_{ij}^{(m)}=\underset{a}{}c\frac{d^3p}{p^0\sqrt{(g)}}p_ip_jF_a(p),$$
(122)
$$J^i=\underset{a}{}e_ac\frac{d^3p}{p^0\sqrt{(g)}}p^iF_a(p).$$
(123)
The system of equations (112). (113) must by augmented by the kinetic equation for $`F_b`$ in relativistic plasma. In the case when the electromagnetical interaction og particles are dominating, the equation for $`F_a`$ was derived in Refs. , .
The covariant form of this kinetic equations is
$$u^i\frac{f_a}{q^i}+\mathrm{\Gamma }_{j,ik}\frac{f_a}{p_i}+\frac{e_a}{c}<F_{ik}>u^k\frac{f_a}{p_i}=$$
$$=\underset{b}{}\frac{}{p_i}\frac{d^4p^{}}{\sqrt{(g)}}E_{ij}(p,p^{})\left(\frac{f_a}{p_j}f_b^{}\frac{f_b^{}}{p_j^{}}f_a\right),$$
(124)
where
$$E_{ij}(p,p^{})=\frac{2\pi e_a^2e_b^2Ln_b}{c^2}[(u,u^{})^21]^{3/2}(u^{},u)^2\times $$
$$\times \left\{g_{ij}[(u,u^{})^21]u_iu_ju_i^{}u_j^{}+(u,u^{})(u_iu_j^{}+u_i^{}u_j)\right\}$$
(125)
with $`(u,u^{})=u_i^{}u^i`$ . Primed and nonprimed quantities refer to particles belonging to species $`a`$, and $`b`$ respectively and $`L`$ is the Coulomb logarithm
$$L=_{k_{min}}^k_{\mathrm{}}\frac{dk}{k}$$
(126)
5. Conclusion
The macroscopic equations of the gravitational field in relativistic plasma differ from the classical Einstein equations by the presence of additional terms
$$Z_{ij}=_k\phi _{ij}^k+\mu _{ij}\chi \tau _{ij}^{(gr)}$$
on the left-hand side due to partial interaction.
These terms are proportional to the square of Einstein constant and to the square of particle namber density.
The macroscopic equations of the electromagnetical field in relativistic plasma differ from classical Maxwell equations by the presence of additional terms
$$Z^i=_k\phi ^{ki}+\mu ^i.$$
on the left-hand side due to particle interaction and due to effects of general relativity.
This terms are proportional to the first power of Einstein constant and to the square of the particle number density.
Hence these terms can play an important role in continuous media of very high density only. Such density are possible in the early stages of the evolution of the Universe and inside objects hear gravitational collapse. Therefore, it is natural to look for applications of the derived equations primarily in the theory of early stages of the Universe evolution and in gravitational collapse theory. |
warning/0002/astro-ph0002196.html | ar5iv | text | # Neutrino Pair Annihilation in the Gravitation of Gamma Ray Burst Sources
## 1 INTRODUCTION
The relativistic fireball (Shemi & Piran 1990; Rees & Mészáros 1992; Mészáros & Rees 1993; Sari & Piran 1995; Sari, Narayan & Piran 1996) is one of the most promising models of gamma ray bursts (GRBs). However, even if the fraction of the baryon rest energy is only $`10^3`$ in the fireball, the relativistic bulk flow, which is indispensable to GRBs, cannot be realized. Notwithstanding the very high energy phenomenon ($`10^{52}`$ ergs), the baryon density in the fireball must be extremely small. This is the famous baryon contamination problem and still remains unsolved. Thus the central engine of GRBs is still beyond deep mist. The source of GRBs may be one super massive (failed) supernovae (Woosley 1993; Paczyński 1998) or may be a merger of two neutron stars or of a neutron star and a black hole (e.g. Eichler et al. 1989; Narayan, Paczyński & Piran 1992; Mészáros & Rees 1992a; Katz 1997; Ruffert & Janka 1998, 1999).
In these compact high energy objects, the neutrino-antineutrino annihilation into electrons and positrons (hereafter neutrino pair annihilation) is a possible and important candidate to explain the energy source of GRBs (Paczyński 1990; Mészáros & Rees 1992b; Janka & Ruffert 1996; Ruffert et al. 1997; Ruffert & Janka 1998, 1999). Motivated by the delayed explosion of Type II supernovae, the energy deposition rate due to the neutrino pair annihilation above the neutrinosphere has been calculated (Goodman, Dar & Nussinov 1987; Cooperstein, Van Den Horn & Baron 1987; Berezinsky & Prilutsky 1987). The energy deposition rate is proportional to $`r^8`$ ($`r`$ is the distance from the center of the neutrinosphere) for a large $`r`$, and almost all deposition occurs near the neutrinosphere. As they themselves noted in their paper, Goodman et al. (1987) neglected the general relativistic effects on the energy deposition rate, which may change their numerical value seriously. In simulations of the neutrino pair annihilation rate, it is very important to confirm whether or not the energy deposition rate is altered or not by the gravitational effects. In the recent study, Salmonson & Wilson (1999) concluded that the energy deposition rate in Type II supernovae is enhanced about 4 times as a result of the gravitational effects. We must check whether or not their results can be applied to the central engine of GRBs.
One of the most probable candidates for the central engine of GRBs is the accretion disk around a black hole (Woosley 1993; Popham, Woosley & Fryer 1999; MacFadyen & Woosley 1999; Ruffert & Janka 1999). The system of an accretion disk and a black hole may be formed by the merging of two neutron stars, the merging of a black hole and a neutron star, or the failed supernovae. In general, the baryon density has the lowest value along the rotation axis just above the black hole (e.g. see Ruffert & Janka 1999). This region might be a key to resolving the baryon contamination problem. The hot accretion disk emits neutrinos and antineutrinos. The energy deposited in the lowest density region is a candidate for the central engine of GRBs. Using hydrodynamic simulations, Ruffert & Janka (1999) showed that the neutrino pair annihilation deposits energy in the vicinity of the torus at a rate of $`(35)\times 10^{50}`$ ergs $`\mathrm{s}^1`$. They concluded that the gravitational effect on the energy deposition rate around the accretion disk is small. We must supplement their results from the analytical side.
In various arguments on the energy deposition in the central engine of GRBs, the order estimation of the deposited energy is sufficient, at least at present. In this article, based on simple models, we study semianalytically the gravitational effects on the energy deposition rate for two cases. In one case neutrinos are emitted spherically symmetrically. In the other case the hot accretion disk emits neutrinos. We have derived the gravitational effects on the former case independently of Salmonson & Wilson (1999). Some differences of our work from Salmonson & Wilson in the formulation, the interpretation of the energy deposition, and the additional factor are mentioned. As for the disk case, we assume that the accretion disk is isothermal and that the gravitational field is dominated solely by the central Schwarzschild black hole. These assumptions enable us to treat the energy deposition around the disk semianalytically. Thus, in both two cases gravitation is described by the Schwarzschild metric, and the essential differences between the two cases come from the shape of the neutrino emitters.
The gravitational effects consist of three factors: they are the bending of neutrino trajectories, the gravitational redshift, and the trapping of deposited energy into the central gravitational source. We show that the energy deposition rate is indeed enhanced rather crucially by the effect of neutrino bending. However, it is also shown that the gravitational redshift and the trapping of the deposited energy reduce this enhancement. As a result, the gravitational effects do not substantially change the energy deposition rate for either the spherical symmetric case or the disk case.
This paper is organized as follows. In section two we investigate neutrino pair annihilation near the neutrinosphere. The same process around the accretion disk is discussed in section three. The last section is devoted to conclusions.
## 2 NEUTRINO PAIR ANNIHILATION NEAR THE NEUTRINOSPHERE
In this section we study the general relativistic effects on neutrino pair annihilation near the neutrinosphere. This study has been already done by Salmonson & Wilson (1999). Using another method, we formulate the same problem independently of the work of Salmonson & Wilson. Some alterations in the interpretation of the energy deposition in Salmonson & Wilson are mentioned.
The number of reaction, $`\nu +\overline{\nu }e^++e^{}`$, per unit volume per unit time (Goodman, Dar & Nussinov 1987) is written as
$$\frac{dN(𝒓)}{dtdV}=f_\nu (𝒑_\nu ,𝒓)f_{\overline{\nu }}(𝒑_{\overline{\nu }},𝒓)\sigma \left|𝒗_\nu 𝒗_{\overline{\nu }}\right|d^3𝒑_\nu d^3𝒑_{\overline{\nu }}.$$
(1)
Here $`f_\nu `$ ($`f_{\overline{\nu }}`$) is the number density of neutrinos (antineutrinos) in phase space, $`𝒗_\nu `$ ($`𝒗_{\overline{\nu }}`$) is the velocity of neutrinos (antineutrinos), and $`\sigma `$ is the rest-frame cross section. The left handside of equation (1) is Lorentz invariant, since both the numerator, $`dN`$, and denominator, $`dtdV=\sqrt{g}d^4x`$, are Lorentz invariant. Since $`f_\nu `$ and $`d^3𝒑_\nu /\epsilon _\nu `$ (where $`\epsilon _\nu `$ is the proper energy of neutrinos) of the right handside are also Lorentz invariant, $`\epsilon _\nu \epsilon _{\overline{\nu }}|𝒗_\nu 𝒗_{\overline{\nu }}|\sigma `$ should be Lorentz invariant. The latter is written in a manifest Lorentz-invariant form as $`\sigma c^3(p_\nu p_{\overline{\nu }})`$, where $`(p_\nu p_{\overline{\nu }})`$ is the inner product of the 4-momenta. The standard model predicts that the cross section is expressed as
$$\sigma =2c^2KG_\mathrm{F}^2(p_\nu p_{\overline{\nu }}),$$
(2)
where the dimensionless parameter $`K`$ is written as
$`K(\nu _\mu \overline{\nu _\mu })`$ $`=`$ $`K(\nu _\tau \overline{\nu _\tau })={\displaystyle \frac{14\mathrm{sin}^2\theta _\mathrm{W}+8\mathrm{sin}^4\theta _\mathrm{W}}{6\pi }},`$
$`K(\nu _\mathrm{e}\overline{\nu _\mathrm{e}})`$ $`=`$ $`{\displaystyle \frac{1+4\mathrm{sin}^2\theta _\mathrm{W}+8\mathrm{sin}^4\theta _\mathrm{W}}{6\pi }}.`$ (3)
Here the Fermi constant $`G_\mathrm{F}^2=5.29\times 10^{44}\mathrm{cm}^2\mathrm{MeV}^2`$ and the Weinberg angle $`\mathrm{sin}^2\theta _\mathrm{W}=0.23`$.
Let us incorporate the effects of gravitational force due to the neutron star or black hole on the neutrino pair annihilation rate. We assume that the gravitational field is described by the Schwarzschild metric:
$$ds^2=g_{ij}dx^idx^j=\left(1\frac{r_g}{r}\right)c^2dt^2\frac{1}{1\frac{r_g}{r}}dr^2r^2\left(d\theta ^2\mathrm{sin}^2\theta d\phi ^2\right),$$
(4)
where $`r_g=2GM/c^2`$ is the Schwarzschild radius. In this field the eikonal for a massless particle (Landau & Lifshitz 1979) is written as
$$\psi =\omega _0t+L\phi +\psi _r(r),$$
(5)
where $`\omega _0`$ and $`L`$ are constants. $`\psi _r(r)`$ satisfies the equation
$$\frac{\psi _r(r)}{r}=\sqrt{\frac{\omega _0^2}{c^2}\left(1\frac{r_g}{r}\right)^2\frac{L^2}{r^2}\frac{1}{1\frac{r_g}{r}}}.$$
(6)
From equation (5), we can obtain the momentum of a neutrino by $`p_i=\mathrm{}\frac{\psi }{x^i}`$.
Let us consider a neutrino and an antineutrino moving on the same surface, $`\theta =\pi /2`$. In this case, the inner product of the two particles is written by
$`(p_\nu p_{\overline{\nu }})`$ $`=`$ $`g^{ij}p_{\nu i}p_{\overline{\nu }j}`$ (8)
$`=`$ $`{\displaystyle \frac{\epsilon _\nu \epsilon _{\overline{\nu }}}{c^2}}(1\sqrt{1\left({\displaystyle \frac{\rho _\nu }{r}}\right)^2\left(1{\displaystyle \frac{r_g}{r}}\right)}\sqrt{1\left({\displaystyle \frac{\rho _{\overline{\nu }}}{r}}\right)^2\left(1{\displaystyle \frac{r_g}{r}}\right)}`$
$`{\displaystyle \frac{\rho _\nu \rho _{\overline{\nu }}}{r^2}}(1{\displaystyle \frac{r_g}{r}})),`$
where
$$\rho _\nu \frac{cL_\nu }{\omega _{0\nu }}.$$
(9)
The proper energy of the neutrino has been written as
$$\epsilon _\nu =\frac{\mathrm{}\omega _{0\nu }}{\sqrt{1\frac{r_g}{r}}}\frac{\epsilon _{0\nu }}{\sqrt{1\frac{r_g}{r}}},$$
(10)
where $`\epsilon _{0\nu }`$ is the energy observed at infinity. Thus the proper energy is redshifted, as is well known. If we define an angle $`\theta _\nu `$ as
$$\mathrm{sin}\theta _\nu =\frac{\rho _\nu }{r}\sqrt{1\frac{r_g}{r}},$$
(11)
equation (8) becomes a simple and natural form,
$$(p_\nu p_{\overline{\nu }})=\frac{\epsilon _\nu \epsilon _{\overline{\nu }}}{c^2}\left(1\mathrm{cos}(\theta \nu \theta \overline{\nu })\right).$$
(12)
The angle $`\theta _\nu `$ ($`\theta _{\overline{\nu }}`$) represents the angle between $`𝒑_\nu `$ ($`𝒑_{\overline{\nu }}`$) and the position vector $`𝒓`$ (see Figure 1). We assume that the neutrinosphere emits neutrinos and antineutrinos isotropically. Then we can write the number densities as $`f_\nu (𝒑_\nu ,𝒓)d^3𝒑_\nu =n(\epsilon _\nu )\epsilon _\nu ^2d\epsilon _\nu d\mathrm{\Omega }`$. Because $`\rho _\nu `$ is constant along a neutrino ray, the maximum angle, $`\theta _\mathrm{M}`$, is obtained by substituting $`\pi /2`$ for $`\theta _\nu `$ at the radius of the neutrinosphere, $`R_\nu `$, in equation (11). Thus we obtain
$$\mathrm{sin}\theta _\mathrm{M}=\frac{R_\nu }{r}\sqrt{\frac{1\frac{r_g}{r}}{1\frac{r_g}{R_\nu }}}.$$
(13)
The effect of the orbital bending is apparent in this equation. Until now we have discussed the maximum angle on the surface of $`\theta =\pi /2`$. In general cases, the angles between $`𝒑_\nu `$ and $`𝒓`$ or the inner product, $`(p_\nu p_{\overline{\nu }})`$, are expressed by the two angles, $`\theta _\nu `$ and $`\phi _\nu `$. From the symmetry, the behaviour of $`\theta _\mathrm{M}`$ is obviously the same as described in equation (13), and $`\phi _\nu `$ varies from $`0`$ to $`2\pi `$.
Using the effective temperature of the neutrinosphere, $`T_{\mathrm{eff}}=T_0/\sqrt{g_{00}}`$ with a constant $`T_0`$, we can write the density
$$n(\epsilon _\nu )=\frac{g_\nu }{(hc)^3}\frac{1}{\mathrm{exp}\left(\frac{\epsilon _\nu }{kT_{\mathrm{eff}}}\right)+1},$$
(14)
where $`g_\nu `$ is a statistical factor ($`g_\nu =1`$ for a neutrino). $`\epsilon _\nu /(kT_{\mathrm{eff}})`$ is constant along a neutrino ray, since the redshift is cancelled out. Thus $`n(\epsilon _\nu )`$ is conserved along a neutrino ray in accordance with Liouville’s theorem in curved spacetime (Misner, Thorne & Wheeler 1975). From the above formulation, one can find
$$\frac{dN(𝒓)}{dtdV}=2cKG_\mathrm{F}^2F(r)𝑑\epsilon _\nu 𝑑\epsilon _{\overline{\nu }}n(\epsilon _\nu )n(\epsilon _{\overline{\nu }})\epsilon _\nu ^3\epsilon _{\overline{\nu }}^3,$$
(15)
where the dimensionless factor $`F(r)`$ is written by
$`F(r)`$ $`=`$ $`{\displaystyle _0^{\theta _\mathrm{M}}}𝑑\theta _\nu \mathrm{sin}\theta _\nu {\displaystyle _0^{\theta _\mathrm{M}}}𝑑\theta _{\overline{\nu }}\mathrm{sin}\theta _{\overline{\nu }}{\displaystyle _0^{2\pi }}𝑑\phi _\nu {\displaystyle _0^{2\pi }}𝑑\phi _{\overline{\nu }}`$ (16)
$`\times \left(1\mathrm{sin}\theta _\nu \mathrm{sin}\theta _{\overline{\nu }}\mathrm{cos}(\phi _\nu \phi _{\overline{\nu }})\mathrm{cos}\theta _\nu \mathrm{cos}\theta _{\overline{\nu }}\right)^2`$
$`=`$ $`{\displaystyle \frac{2\pi ^2(1X)^4}{3}}(X^2+4X+5),`$ (17)
where
$$X=\sqrt{1\left(\frac{R_\nu }{r}\right)^2\frac{1\frac{r_g}{r}}{1\frac{r_g}{R_\nu }}}.$$
(18)
In our assumption, the energy deposited by the neutrino pair annihilation is propagated outward as a fireball or a shock wave, and observed as a GRB by a distant observer. Thus the energy we need to calculate is $`\epsilon _{0\nu }`$, not the proper energy $`\epsilon _\nu `$. In this case the energy deposition rate is obtained by putting a factor $`(\epsilon _{0\nu }+\epsilon _{0\overline{\nu }})`$ in the integrand in equation (15);
$`{\displaystyle \frac{dE_0(𝒓)}{dtdV}}`$ $`=`$ $`{\displaystyle \frac{2cKG_\mathrm{F}^2}{\left(1\frac{r_g}{r}\right)^4}}{\displaystyle _0^{\mathrm{}}}{\displaystyle _0^{\mathrm{}}}𝑑\epsilon _{0\nu }𝑑\epsilon _{0\overline{\nu }}`$ (19)
$`\times n(\epsilon _{0\nu })n(\epsilon _{0\overline{\nu }})\epsilon _{0\nu }^3\epsilon _{0\overline{\nu }}^3(\epsilon _{0\nu }+\epsilon _{0\overline{\nu }})F(r)`$
$`=`$ $`{\displaystyle \frac{21\pi ^4}{4}}\zeta (5){\displaystyle \frac{KG_\mathrm{F}^2g_\nu ^2}{h^6c^5}}{\displaystyle \frac{\left(1\frac{r_g}{R_\nu }\right)^{\frac{9}{2}}}{\left(1\frac{r_g}{r}\right)^4}}(kT_{\mathrm{eff}})^9F(r).`$ (20)
The integrals for $`\epsilon _{0\nu }`$ and $`\epsilon _{0\overline{\nu }}`$ should be defined in the range in which the total energy produced by the pair annihilation is larger than the mass of created electrons, and smaller than the masses of weak bosons. Here we have approximated the integrals as expressed in equation (19) in the same manner as Salmonson & Wilson (1999) did. This is because the cross section decreases with the energy of neutrinos, and the number of neutrinos whose energy is larger than the masses of weak bosons is also very small in our assumption ($`kT_{\mathrm{eff}}`$ is of the order of several MeV). The factor $`(1r_g/R_\nu )^{9/2}/(1r_g/r)^4`$ represents the effect of the gravitational redshift, and $`F(r)`$ includes the effect of the orbital bending.
As is understood from the Lorentz invariant, $`dtdV=\sqrt{g}d^4x`$, if we integrate equation (20) over proper volume, $`dV^{}=\sqrt{g_{rr}g_{\theta \theta }g_{\phi \phi }}drd\theta d\phi `$, we can obtain the total energy deposition per unit proper time, $`d\tau =\sqrt{g_{00}}dt`$. It is natural to evaluate the energy deposition rate by the world time $`dt`$ for a distant observer. We integrate over the volume, $`dV=\sqrt{g_{\theta \theta }g_{\phi \phi }}drd\theta d\phi `$. Thus the energy deposition per unit world time is expressed as
$$\frac{dE_0}{dt}=\frac{21\pi ^4}{4}\zeta (5)\frac{KG_\mathrm{F}^2g_\nu ^2}{h^6c^5}(kT_{\mathrm{eff}})^9\left(1\frac{r_g}{R_\nu }\right)^{\frac{9}{2}}_{R_\nu }^{\mathrm{}}𝑑r4\pi r^2\frac{F(r)}{\left(1\frac{r_g}{r}\right)^4}C(r),$$
(21)
where we have put a factor,
$$C(r)=\frac{1}{2}\left(1+\sqrt{1\frac{27}{4}\left(\frac{r_g}{r}\right)^2\left(1\frac{r_g}{r}\right)}\right),$$
(22)
in the integrand. This is the escape probability of the deposited energy at $`r`$ from the gravitational attraction (Chandrasekhar 1983; Shapiro & Teukolsky 1983; Ruffert & Janka 1999). The electrons, positrons and photons which are captured by the gravitational attraction cannot contribute to the energy source of GRBs. Apart from $`C(r)`$, the radial profile in the integrand of equation (21) is different from those of Salmonson & Wilson (1999), since Salmonson & Wilson calculated the proper energy deposition per unit proper time. Of course, the results of Salmonson & Wilson are not mistakes for the estimate of the energy deposition rate in supernovae. For the source of GRBs, however, equation (21) is adequate.
Let us investigate the effects of the redshift, orbital bending and gravitational capture. We integrate equation (21) and obtain the energy deposition rate for $`\nu _e`$ as
$$\frac{dE_0}{dt}=1.27\times 10^{42}\left(\frac{kT_{\mathrm{eff}}}{1\mathrm{M}\mathrm{e}\mathrm{V}}\right)^9\left(\frac{R_\nu }{10\mathrm{k}\mathrm{m}}\right)^3f\text{ergs }\mathrm{s}^1,$$
(23)
where the dimensionless factor $`f`$ expresses the effects of the general relativity ($`f=1`$ when we neglect the gravitation). The energy deposition rates for $`\nu _\mu `$ and $`\nu _\tau `$ are 0.64 times equation (23). We numerically estimate $`f`$ including the effects of the redshift only, the orbital bending only, or both redshift and orbital bending. Last, the total effects of the redshift, orbital bending, and gravitational capture are calculated. The results are listed in Table 1. As Table 1 or equation (21) indicates, the effect of the redshift reduces the energy deposition rate, and the effect of the orbital bending increases it. Although each effect, that of the redshift and that of orbital bending, is substantial, the effects partly cancel each other. As a result, the order of the energy deposition rate for the most probable case, $`R_\nu /r_g=2.5`$, is not altered. When we neglect the general relativistic effects, the energy deposition rate increases by 1.3 times as $`R_\nu `$ becomes 10% larger, and also increases by 2.4 times as the temperature becomes 10% higher. Therefore, the gravitational effects are not so large in comparison with the errors due to the uncertainties of $`R_\nu `$ or $`T_{\mathrm{eff}}`$, and are overwhelmed by them. The effect of the gravitational capture becomes important as $`R_\nu /r_g`$ decreases. As is plotted in Figure 2, in the cases of both the presence and absence of gravitation, the energy deposition mainly occurs near the neutrinosphere.
Salmonson & Wilson (1999) concluded that the effects of gravity enhance the energy deposition rate up to a factor of more than 4 for $`R_\nu 2.5r_g`$. However, our results show that the gravitational effects reduce the energy deposition rate. This discrepancy survives even if we omit the escape factor $`C(r)`$. The proper energy deposition per unit proper time is enhanced by both the effects of the redshift and that of orbital bending. Additionally, Salmonson & Wilson expressed the general relativistic effects with the fixed neutrino luminosity at infinity $`L_{\mathrm{}}`$, whereas we have done so with the local physical quantity $`T_{\mathrm{eff}}`$. Therefore, an additional factor coming from the redshift of the luminosity ($`L(R_\nu )T_{\mathrm{eff}}^4`$, $`L_{\mathrm{}}=(1r_g/R_\nu )L(R_\nu )`$) enhances the energy deposition rate in the work of Salmonson & Wilson. However, the quantity $`L_{\mathrm{}}`$ of GRBs is not directly observable at present. It is more natural to study the effects for the given local parameters, $`T_{\mathrm{eff}}`$ or $`L(R_\nu )`$, which is restricted or provided by models of the central engine.
## 3 NEUTRINO PAIR ANNIHILATION AROUND THE ACCRETION DISK
In this section we investigate the energy deposition rate around the accretion disk. In order to simplify our formulation, we assume that the accretion disk is isothermal and that the gravitational field is dominated by the central Schwarzschild black hole. We neglect the rotation of the black hole. The accretion disk is assumed to be thin, and its self-gravitational effects are neglected. Of course, these idealizations may be far from the case of the realistic accretion disk. However, we consider that this simple method is sufficient for qualitatively studying the gravitational effects on the energy deposition rate. In this case the equation of the energy deposition rate is the same as equation (20) provided that $`F(r)`$ is replaced by $`F(r,\theta )`$ (it will be given below). The effect of the gravitational redshift can be easily incorporated, whereas the formulation of the neutrino bending is difficult to do because the accretion disk emits neutrinos anisotropically.
First, we calculate the dimensionless factor $`F(r,\theta )`$ without the effect of gravity. The accretion disk is placed on the equatorial plane, $`\theta =\pi /2`$. The black hole is at the origin, and we consider a point $`P=(r,\theta ,0)`$ where pair annihilations occur (see Figure 3). A neutrino is emitted from an arbitrary point on the disk $`S=(R,\pi /2,\phi )`$, where $`R`$ is limited in the range from $`R_{\mathrm{in}}`$ to $`R_{\mathrm{out}}`$. The neutrino emitted from $`S`$ travels straight and arrives at the point $`P`$. Let us denote the angle components of the vector joining $`S`$ and $`P`$ by $`(\theta _\nu ,\phi _\nu )`$. They are given by
$$\mathrm{cos}\theta _\nu =\frac{r\mathrm{cos}\theta }{\sqrt{r^2+R^22rR\mathrm{sin}\theta \mathrm{cos}\phi }},$$
(24)
$$\mathrm{sin}\phi _\nu =\frac{R\mathrm{sin}\phi }{\sqrt{r^2\mathrm{sin}^2\theta +R^22rR\mathrm{sin}\theta \mathrm{cos}\phi }}.$$
(25)
Thus $`\theta _\nu `$ and $`\phi _\nu `$ are functions of $`R`$ and $`\phi `$ for fixed $`r`$ and $`\theta `$. The Jacobian $`J(\theta _\nu ,\phi _\nu )/(R,\phi )`$ is
$$J=\frac{rR\mathrm{cos}\theta }{\sqrt{r^2\mathrm{sin}^2\theta +R^22rR\mathrm{sin}\theta \mathrm{cos}\phi }(r^2+R^22rR\mathrm{sin}\theta \mathrm{cos}\phi )}.$$
(26)
Consequently, we obtain $`F(r,\theta )`$ as
$`F(r,\theta )`$ $`=`$ $`{\displaystyle _{R_{\mathrm{in}}}^{R_{\mathrm{out}}}}𝑑R{\displaystyle _{R_{\mathrm{in}}}^{R_{\mathrm{out}}}}𝑑R^{}{\displaystyle _0^{2\pi }}𝑑\phi {\displaystyle _0^{2\pi }}𝑑\phi ^{}JJ^{}`$ (27)
$`\times `$ $`\mathrm{sin}\theta _\nu \mathrm{sin}\theta _{\overline{\nu }}\left(1\mathrm{sin}\theta _\nu \mathrm{sin}\theta _{\overline{\nu }}\mathrm{cos}(\phi _\nu \phi _{\overline{\nu }})\mathrm{cos}\theta _\nu \mathrm{cos}\theta _{\overline{\nu }}\right)^2.`$
In equation (27) we adopt $`R_{\mathrm{in}}=3r_g`$, the innermost stable orbit, and $`R_{\mathrm{out}}=10r_g`$ as Woosely (1993) assumed. $`F(r,\theta )`$ derived from the numerical integral of equation (27) is plotted in Figure 4(a) and (b). As is shown in these figures, the energy deposition rate is maximized in the vicinity of the accretion disk, where $`F(r,\theta )3033`$. The simulation of a neutron star merger by Ruffert & Janka (1999) showed that the Paczyński-Wiita potential (Paczyński & Wiita 1980), which mimics the effects of the general relativity, gives a relatively more transparent disk for neutrinos than that given by the Newtonian potential. The profile of the energy deposition rate in the Paczyński-Wiita potential is similar to our analytical one depicted in Figure 4(a), which shows that the rate takes its maximum value on the surface of the disk. On the other hand, the simulated deposition rate in the Newtonian potential is maximized near the rotation axis. Let us calculate the energy deposition rate near the rotation axis, where is the lowest baryon density region. Thus we calculate in the region $`\theta \pi /4`$ and obtain
$$\frac{dE_0}{dt}=5.22\times 10^{43}\left(\frac{kT_{\mathrm{eff}}}{1\mathrm{M}\mathrm{e}\mathrm{V}}\right)^9\left(\frac{r_g}{10\mathrm{k}\mathrm{m}}\right)^3Gf\text{ergs }\mathrm{s}^1,$$
(28)
where the dimensionless quantity $`G`$ shows the relative contributions from various regions in the absence of gravitation. $`G`$ is normalized to unity when we integrate over the volume for $`\theta \pi /4`$ and $`r=2r_g10r_g`$. The values of $`G`$ in the other regions are summarized in Table 2, from which we can obtain the energy deposition rate in the respective region. We neglect the energy deposited inside $`r=2r_g`$, since the baryon density in this region is very high and the energy contribution is small for the small volume and deposition rate. In the case of the spherical emitter in section 2, the deposition rate decreases as $`r^8`$. On the other hand, around the accretion disk, as is seen from Table 2, there remains a marginal deposition rate even at regions relatively distant from the center. Of course the deposition rate per unit volume at distant positions is small. However, large volume results in a non-negligible contribution at the regions distant from the center.
Untill now we have neglected the gravitational effects. It is easy to incorporate the effects of the redshift and trapping by the central gravitational source in the preceding arguments of this section. However, the bending effect is difficult to treat, unlike the case of the neutrinosphere, since the accretion disk emits neutrinos anisotropically. Thus we are forced to make some approximation. As is shown in Figure 4(b), the $`\theta `$ dependence of $`F(r,\theta )`$ is weak for small $`\theta `$. We may set $`F(r,\theta )F(r,0)`$ for $`\theta \pi /4`$. In the absence of gravitation if we adopt this approximation in the region $`\theta \pi /4`$ and $`r=2r_g10r_g`$, we obtain $`G=0.81`$. The exact value of $`G`$ is unity, and this approximation is not necessarily satisfactory. However, this approximation may be sufficient for the order estimate of the gravitational effects.
We can obtain $`F(r,0)`$ including the effect of orbital bending with comparative ease, since the geometry of this case maintain the symmetry. A neutrino is emitted from the disk at $`R`$ and $`\theta =\pi /2`$, and it arrives at a point at $`r`$ and $`\theta =0`$. The nearest distance, $`r_0`$, from the origin to the orbit of neutrinos (Landau & Lifshitz 1979) is numerically obtained from
$$\pi /2=_\mathrm{C}\frac{dr^{}}{r^{}\sqrt{\left(\frac{r^{}}{r_0}\right)^2\left(1\frac{r_g}{r_0}\right)\left(1\frac{r_g}{r^{}}\right)}}.$$
(29)
Here, in the case in which a neutrino passes through $`r_0`$ until it arrives at a point at $`\theta =0`$, the integration for $`r^{}`$ is performed from $`r_0`$ to $`R`$ and $`r`$. When the distance from the origin to the neutrino varies monotonically, the integration is performed from the smaller to the larger of $`r`$ and $`R`$. We can get $`\theta _\nu `$ at $`\theta =0`$ numerically from $`r_0`$ and the following equation;
$$\mathrm{sin}\theta _\nu =\frac{r_0}{r}\sqrt{\frac{1\frac{r_g}{r}}{1\frac{r_g}{r_0}}}.$$
(30)
The constant, $`r_0`$, or $`\theta _\nu `$, is a function of $`r`$ and $`R`$ in this case. As is easily understood, a neutrino coming from $`R_{\mathrm{in}}`$ forms $`\theta _\mathrm{m}`$, the minimum value of $`\theta _\nu `$, at $`\theta =0`$ and that from $`R_{\mathrm{out}}`$ forms the maximum value of $`\theta `$, $`\theta _\mathrm{M}`$. Integrating equation (16) from $`\theta _\mathrm{m}`$ to $`\theta _\mathrm{M}`$, we obtain $`F(r,0)`$ involving gravitational effects. In Figure 5 we plot $`F(r,0)`$ for both the case when the bending is taken into consideration and the case when it is not. In comparison with the spherical case in Figure 2, the deposited energy at distant regions in the disk case is marginally substantial. In the presence of bending, the peak of the energy deposition rate is shifted to a little bit larger $`r`$ and the value of the rate at the peak is about twice as larger as values obtained in the absence of bending.
Although the $`\theta `$ dependence of the gravitational effects may not necessarily be small, unlike $`F(r,\theta )`$ in the absence of gravitation, we assume it is small here. Using $`F(r,0)`$ with the bending effect, we calculate the energy deposition rate in the range $`\theta \pi /4`$ and $`r=2r_g10r_g`$. Table 3 lists the values of the factor $`f`$ that shows the gravitational effects. The $`\theta `$-dependence is neglected in our calculation except for the case involving the redshift only. Thus $`f`$ is normalized to unity when $`F(r,\theta )`$ (in the case involving the effect of the redshift only) or $`F(r,0)`$ (in the other cases) is integrated over $`r`$ and $`\theta `$ in the absence of gravitation. As is easily seen, the gravitational effects cancel one another out. This is analogous to the neutrino sphere case in the previous section. This result strongly supports that of Ruffert & Janka (1999). They treated the system of the accretion torus and a black hole unlike our system of the disk and a black hole. Using an approximation similar to ours, they analytically calculated the energy deposition rate due to the neutrino pair annihilation. Their result is that the gravitational effects reduce the deposition rate by a factor of $`1030\%`$. It agrees well with our result.
In order to circumvent the baryon contamination problem, the energy fraction of baryonic matter in the fireball must be less than about $`10^5`$ (Shemi & Piran 1990). If we adopt the duration time of the neutrino radiation to be $`t_{\mathrm{dur}}=0.1`$s and $`T_{\mathrm{eff}}=10`$MeV, the highest mean mass densities $`\overline{\rho }`$ inside $`\theta =0\pi /3`$ to resolve the above problem are $`10^6`$g/$`\mathrm{cm}^3`$ for $`r=2r_g5r_g`$, $`10^5`$g/$`\mathrm{cm}^3`$ for $`r=5r_g10r_g`$ and $`10^4`$g/$`\mathrm{cm}^3`$ for $`r=10r_g20r_g`$. Since some fraction of energy really escapes from the considered regions during the finite duration time, the above restrictions may become more stringent.
## 4 CONCLUSIONS
In this article we have investigated semianalytically the neutrino pair annihilation near the neutrinosphere and around the thin accretion disk assuming that the gravitational sources in both cases are described by the Schwarzschild metric. The accretion disk has been assumed to be a blackbody and isothermal. These assumptions enable us to treat these two cases based in an almost unified fashion, which also clarifies the physical differences between these two cases. We have studied the general relativistic effects only near the rotation axis, because that region is especially of interest to the source of GRBs and estimating the effect of orbital bending for large $`\theta `$ is difficult.
The general relativistic effects as a whole do not enhance the neutrino energy deposition rate in either case. The energy deposition rate is enhanced by the effect of orbital bending toward the center. However, the enhancement is cancelled out by the effects of the redshift and capture by the gravitational attraction. Consequently, numerical simulations of the neutrino energy deposition rate in various models can correctly estimate the order of the rate without considering the gravitational effects, since it is supposed that the thickness, shape, or temperature distribution of the disk or sphere does not greatly affect the gravitational effects themselves. Taking into account also the results of Ruffert & Janka (1999), the conclusions mentioned above are strongly suggested to be valid in the following geometrical forms of the neutrino source: sphere, thin disk and torus. We have also shown that the neutrinos emitted from the disk can deposit energy at more distant regions than the neutrinos emitted from the sphere. The importance in this article resides in the qualitative properties of the general relativistic effects. The quantitative calculations in this paper are not so important, and should be investigated on the basis of more sophisticated models and simulations.
We appreciate the helpful advice of M. Ruffert. This work was partly supported by a Research Fellowship of the Japan Society for the Promotion of Science.
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Figure captions
| $`R_\nu /r_g`$ | $`f`$ | | | |
| --- | --- | --- | --- | --- |
| | Redshift Only | Bending Only | Redshift and Bending | Whole |
| 1.5 | 0.32 | 4.7 | 0.97 | 0.57 |
| 2.5 | 0.60 | 1.6 | 0.87 | 0.73 |
| 5 | 0.81 | 1.2 | 0.93 | 0.89 |
Table 1. The dimensionless factor $`f`$ represents the general relativistic effects (see eq. ) when neutrinos are emitted isotropically from the neutrinosphere; $`f`$ is normalized to unity in the absence of gravitation. The column headings ”Redshift Only” and so on indicate the incorporated effects of gravitation; ”Whole” over the last column means that we incorporate all gravitational effects, redshift, bending, and trapping. The energy deposition rate is enhanced by orbital bending and reduced by the redshift and trapping.
| $`\theta `$ | $`2r_g5r_g`$ | $`5r_g10r_g`$ | $`10r_g20r_g`$ |
| --- | --- | --- | --- |
| $`0\pi /4`$ | 0.35 | 0.65 | 0.22 |
| $`\pi /4\pi /3`$ | 0.32 | 0.77 | 0.17 |
Table 2. The dimensionless factor $`G`$. It represents the fraction of the energy deposition rate for each region (see eq. ) when neutrinos are emitted from the disk. $`G`$ is normalized to unity for the region surrounded by $`r=2r_g10r_g`$ and $`\theta =0\pi /4`$. Here we neglect the effects of gravitation.
| range | $`f`$ | | | |
| --- | --- | --- | --- | --- |
| | redshift only | bending only | redshift and bending | whole |
| $`2r_g5r_g`$ | 0.66 | 2.6 | 1.6 | 1.4 |
| $`5r_g10r_g`$ | 0.31 | 2.5 | 0.78 | 0.75 |
Table 3. The dimensionless factor $`f`$ for $`\theta =0\pi /4`$ (see equation ). Neutrinos are emitted from the disk; $`f`$ is normalized to unity when $`F(r,\theta )`$ (in the case involving the effect of the redshift only) or $`F(r,0)`$ (in the other cases) is integrated over $`r`$ and $`\theta `$ in the absence of gravitation. |
warning/0002/astro-ph0002351.html | ar5iv | text | # Transport Phenomena and Light Element Abundances in the Sun and Solar Type Stars
## 1. Introduction
Element diffusion and mixing processes in stellar interiors are now widely constrained, first by detailed observations of abundances, second by helio and asteroseismology. In most cases however, pure microscopic diffusion in stars would lead to abundance variations much larger than those observed : mild macroscopic motions in stellar radiative zones are definitely needed to account for the observations. This gives strong constraints on the kind of mixing processes allowed. Other constraints come from the consequences of the nuclear reactions occuring in stellar interiors : in some cases stellar mixing from the atmosphere down to the regions of nuclear processing is needed to explain the observed element abundances. This is the case, for example, to account for the depletion of lithium in the Sun and solar type stars.
Lithium observations in main-sequence population I field stars and galactic clusters show a large abundance dispersion which has been extensively studied in the literature (see reviews by Deliyannis 2000, Charbonnel 2000, Michaud 2000 and Pinsonneault 2000). The lithium abundance decreases for decreasing effective temperature below 5500K and the depletion increases with increasing age. This is generally attributed to the deepening of the convective zone, associated with some mild mixing process connecting the bottom of the convective zone with the nuclear destruction region.
Lithium is also depleted in F-type stars (the so-called “Boesgaard dip”). Several possible reasons have been invoked to explain this feature, most related to mixing and nuclear destruction. Element segregation has been proved negligible here as it would lead to unobserved variations of metal abundances (Turcotte et al 1998) and beryllium (Boesgaard 2000).
On the other hand, observations of lithium in main-sequence population II field stars show remarkably constant abundances, with a very small dispersion (e.g. Bonifacio and Molaro 1997) Why is lithium destroyed in Pop I stars while it does not seem destroyed in Pop II stars?
For the same effective temperatures, the convective zone is smaller in Pop II stars than in Pop I stars because of their smaller metallicity. Meanwhile they have a smaller rotation velocity on the average. This could explain why the lithium destruction induced by nuclear reactions is smaller in these stars than in Pop I stars. However the element segregation is more important for smaller densities and smaller rotation, so that this process should lead to a visible lithium depletion, which is not observed (Vauclair and Charbonnel 1995 and 1998). This represents the so-called “lithium paradox”. Here we suggest that the influence of $`\mu `$-gradients on the rotation-induced mixing may help solving this paradox.
## 2. Competition between rotation induced mixing and element diffusion
In rotating stars, the equipotentials of “effective gravity” (including the centrifugal acceleration) have ellipsoidal shapes while the energy transport still occurs in a spherically symetrical way. The resulting thermal imbalance must be compensated by macroscopic motions: the so-called “meridional circulation” (Von Zeipel 1924). The stellar regions outside the convective zones cannot be in complete radiative equilibrium. They are subject to entropy variations given by :
$`\rho T\left({\displaystyle \frac{S}{t}}+𝐮S\right)`$ $`=`$ $`𝐅+\rho \epsilon _n`$ (1)
$`=`$ $`\rho \epsilon _\mathrm{\Omega }(0)`$
where $`𝐅`$ represents the heat flux, $`\epsilon _n`$ the nuclear energy production and $`\epsilon _\mathrm{\Omega }`$ an energy generation rate which results from sources and sinks of energy along the equipotentials.
The vertical component of the meridional velocity $`u_r`$ is computed as a function of $`\epsilon _\mathrm{\Omega }`$ in the stationary regime (from eq. 1):
$$u_r=\left(\frac{P}{C_p\rho T}\right)\frac{\epsilon _\mathrm{\Omega }}{g}$$
(2)
which, for a perfect gas, reduces to:
$$u_r=\frac{\epsilon _\mathrm{\Omega }}{g}\frac{_{\mathrm{ad}}}{_{\mathrm{ad}}+_\mu }$$
(3)
where $`g`$ represents the local gravity, $`_{\mathrm{ad}}`$ and $``$ the usual adiabatic and real ratios $`\left({\displaystyle \frac{d\mathrm{ln}T}{d\mathrm{ln}P}}\right)`$ and $`_\mu `$ the mean molecular weight contribution $`\left({\displaystyle \frac{d\mathrm{ln}\mu }{d\mathrm{ln}P}}\right)`$.
The expression of $`\epsilon _\mathrm{\Omega }`$ is computed by expanding the right-hand-side of eq. (1) on a level surface and writing that its mean value vanishes.
Mestel (1953, 1957 and 1965) pointed out that, in the presence of vertical $`\mu `$-gradients, $`\epsilon _\mathrm{\Omega }`$ contains two kinds of terms : those related to the resulting horizontal variations of $`\mu `$: the so-called “$`\mu `$-induced currents” $`E_\mu `$ and those independent of $`\mu `$, the so-called “$`\mathrm{\Omega }`$-induced currents” $`E_\mathrm{\Omega }`$ . The expression of $`\epsilon _\mathrm{\Omega }`$ obtained in this case has been derived in detail by Maeder and Zahn (1998), who took into account several effects which were not included in the previous computations: more general equations of state instead of perfect gas law, presence of a thermal flux induced by horizontal turbulence, non-stationary cases.
Vauclair (1999) discussed more simple expressions, valid only for negligible differential rotation. In this case $`\mu `$-currents are opposite to $`\mathrm{\Omega }`$-currents in most of the star and $`\epsilon _\mathrm{\Omega }`$ may be written :
$$\epsilon _\mathrm{\Omega }=\left(\frac{L}{M}\right)\left(E_\mathrm{\Omega }+E_\mu \right)P_2(\mathrm{cos}\theta )$$
(4)
with:
$`E_\mathrm{\Omega }`$ $`=`$ $`{\displaystyle \frac{8}{3}}\left({\displaystyle \frac{\mathrm{\Omega }^2r^3}{GM}}\right)\left(1{\displaystyle \frac{\mathrm{\Omega }^2}{2\pi G\overline{\rho }}}\right)`$ (5)
$`E_\mu `$ $`=`$ $`{\displaystyle \frac{\rho _m}{\overline{\rho }}}\left\{{\displaystyle \frac{r}{3}}{\displaystyle \frac{d}{dr}}\left[\left(H_T{\displaystyle \frac{d\mathrm{\Lambda }}{dr}}\right)(\chi _\mu +\chi _T+1)\mathrm{\Lambda }\right]{\displaystyle \frac{2H_T\mathrm{\Lambda }}{r}}\right\}`$ (6)
Here $`\overline{\rho }`$ represents the density average on the level surface $`(\rho )`$ while $`\rho _m`$ is the mean density inside the sphere of radius $`r`$; $`H_T`$ is the temperature scale height; $`\mathrm{\Lambda }`$ represents the horizontal $`\mu `$ fluctuations $`{\displaystyle \frac{\stackrel{~}{\mu }}{\overline{\mu }}}`$; $`\chi _\mu `$ and $`\chi _T`$ represent the derivatives:
$$\chi _\mu =\left(\frac{\mathrm{ln}\chi }{\mathrm{ln}\mu }\right)_{P,T};\chi _T=\left(\frac{\mathrm{ln}\chi }{\mathrm{ln}T}\right)_{P,\mu }$$
(7)
Vertical $`\mu `$-gradients may occur in stars due to two different processes : first the nuclear reactions which occur in the stellar cores, second the helium settling which occurs in the outer layers. The importance of the first process in reducing or even suppressing the meridional motions has been demonstrated several times in the literature (e.g. Huppert and Spiegel 1977). The second process on the other hand has not been extensively studied. We claim here that it may play a crucial role for understanding the lithium problem in Pop I and Pop II stars.
## 3. Application to Pop II stars
Computations of $`\mu `$-currents induced by the helium settling in halo stars have been performed by Vauclair 1999 and Théado and Vauclair 2000 a and b. We found that, for slow rotation, $`\mu `$-currents cancel $`\mathrm{\Omega }`$-currents for very small concentration gradients, corresponding to $`\mu `$-gradients of order $`10^{15}`$ cm<sup>-1</sup>.
Let us summarize the situation of a slowly rotating star in which element settling leads to an increase of the $`\mu `$-gradient below the outer convection zone. At the beginning, the star is homogeneous and meridional circulation can occur, leading to upward flows in the polar regions and downward flows in the equatorial parts (except in the very outer layers where the Gratton-Öpik term becomes important, which we do not discuss here). The $`\mu `$-currents, opposite to the classical $`\mathrm{\Omega }`$-currents, are first negligible. The $`\mu `$-gradients increasing with time because of helium settling, the order of magnitude of the $`\mu `$-currents also increases until it reaches the value for which the circulation vanishes.
This does not occur all at once: as the $`\mu `$-gradient decreases with depth below the convective zone, we expect that the meridional circulation freezes out step by step (see figure 1 of Théado and Vauclair 2000a). An equilibrium situation may be reached, in which the temperature and mean molecular weight gradients along the level surfaces are such that $`\mathrm{\Omega }`$-currents and $`\mu `$-currents cancel each other.
Once it is reached, this equilibrium situation is quite robust. Suppose that some mechanism leads to a decrease of the $`\mu `$-gradient: then $`|E_\mu |`$ becomes smaller than $`|E_\mathrm{\Omega }|`$ and the circulation tends to be restablished in the $`|E_\mathrm{\Omega }|`$ direction, thereby restoring the original $`\mu `$ gradient. Suppose now that the $`\mu `$-gradient is increased. Then $`|E_\mu |`$ becomes larger than $`|E_\mathrm{\Omega }|`$ and the circulation begins in the $`E_\mu `$ direction. Here again the original gradient is restored.
When the meridional circulation is frozen below the convective zone, helium settling could proceed further; however, due to the increase of the diffusion time scale with depth, this would modify the $`\mu `$-gradient. We may thus expect that $`\mu `$-currents would take place and restore the original equilibrium gradient, thereby strongly reducing the microscopic diffusion (Théado and Vauclair 2000b). This self-regulating process could be the reason for the low dispersion of the lithium abundance in the lithium plateau of halo stars.
## 4. Discussion : Pop I versus Pop II stars
There are many observations in stars which give evidences of mixing processes occuring below the outer convective zones as, for example, the lithium depletion observed in the Sun and in galactic clusters. The process we have described above should not apply in all these stars. The reason could be related to the rapid rotation of young stars on the ZAMS and to their subsequent rotational braking.
The abundance determinations in the solar photosphere show that lithium has been depleted by a factor of about 140 compared to the protosolar value while beryllium has not been depleted by more than a factor 2, and maybe much less, as discussed by Balachandran and Bell (1997). These values represent strong constraints on the mixing processes in the solar interior.
Observations of the <sup>3</sup>He/<sup>4</sup>He ratio in the solar wind and in the lunar rocks (Geiss 1993, Geiss and Gloecker 1998) show that this ratio may not have increased by more than $`10\%`$ since 3 Gyr in the Sun. While the occurence of some mild mixing below the solar convective zone is needed to explain the lithium depletion , the <sup>3</sup>He/<sup>4</sup>He observations put a strict constraint on its efficiency. The only way to obtain such a result is to postulate a mild mixing, which would be efficient down to the lithium nuclear burning region but not too far below, to preserve the original <sup>3</sup>He abundance. The efficiency of this mixing should also decrease with time, as the <sup>3</sup>He peak itself builts up during the solar life.
It is interesting to compute the minimum enhancement of the <sup>3</sup>He/<sup>4</sup>He ratio implied by the lithium observed depletion. Vauclair and Richard 1998 showed that it is possible to deplete lithium by a factor larger than $`100`$ as observed and not increase <sup>3</sup>He/<sup>4</sup>He by more than 5 percent since the solar origin. In this case beryllium is only depleted by about 10 percent.
Such a confined mixing zone is also needed from helioseismology : although the introduction of pure element settling in the solar models considerably improves the consistency with the seismic Sun, some discrepancies do remain, particularly below the convective zone where a ”spike” appears in the sound velocity (Richard et al 1996, Turck-Chièze et al. 1998). It has been shown that this behavior may be due to the helium gradient which would be too strong in case of pure settling. Mild macroscopic motions below the convective zone slightly decrease this gradient and helps reducing the discrepancy (Richard et al 1996, Corbard et al 1998, Brun et al 1998). The helium profiles directly obtained from helioseismology (Basu 1998, Antia and Chitre 1998) show indeed a helium gradient smoother than the gradient obtained with pure settling.
The constraints implied by both the helioseismic inversions and abundance determinations in the Sun converge towards the existence of a small mild mixing region below the convective zone, which would extend down to a depth of the order of one scale height. The implied mixing region must be very mild, with diffusion coefficients of $`10^3`$ \- $`10^4`$ only. It must also be completely deconnected from the solar core. No mixing can indeed be allowed down to the nuclear energy production region as it would lead to a sound velocity incompatible with helioseismology. In particular the mixing processes invoked by Morel and Schatzman 1996 to decrease the neutrino fluxes are excluded by helioseismology (Richard and Vauclair 1997).
Mixing processes localized at the boundary between convective and radiative regions include overshooting and regions of large differential rotation like the “tachocline” below the solar convective zone. Up to now, overshooting was generally treated in the models simply as a continuation of the convective zone on a fraction of a pressure scale height. Recent parametrisations use a diffusion coefficient which decreases exponentially with decreasing radius (Freytag et al 1996). The tachocline, which represents in the present Sun the small boundary between the region of large differential rotation (in the convective zone) and the region of solid rotation (in the radiative zone below) is also treated as a mixed layer with an exponentially decreasing diffusion coefficient (Brun et al 1998, Richard 1999). Results are encouraging, although more sophisticated numerical simulation including 2-D abundance variations would be needed to go further.
In any case, the self-regulating process that we have discussed for halo stars in section 3 would not apply below the convective zone in the Sun and solar type stars because of the differential rotation which takes place there. Such a differential rotation would not be expected in halo stars if we suppose that they always rotated slowly and thus did not suffer large transport of angular momentum. The different behavior for the lithium abundance in Pop I and Pop II stars could thus be directly related to their rotation history.
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warning/0002/hep-lat0002001.html | ar5iv | text | # Green’s Function Monte Carlo study of correlation functions in the (2+1)D 𝑈(1) lattice gauge theory
## I Introduction
The two major variants of lattice gauge theory (LGT) are the “Euclidean” formulation of Wilson , and the “Hamiltonian” version of Kogut and Susskind . In the Euclidean régime, classical Monte Carlo simulations have proved to be extremely powerful in extracting quantitative predictions from the theory, as first shown by Creutz . This approach is preferred by an overwhelming majority of lattice gauge theorists at the present time.
The Hamiltonian formulation is still worthy of study, however. It can provide a valuable check of universality, for instance. Lattice gauge theory relies on the fundamental assumption that quantities such as mass ratios calculated in the continuum limit (a critical point of the lattice model) must be ‘universal’, i.e. independent of the microscopic lattice structure or space-time formulation. There is not much real doubt that this is correct, but it is important to provide checks where possible . Another reason is that many techniques imported from quantum many-body theory and condensed matter physics can be employed in this arena, which may give useful results. Examples include the strong-coupling series approach , the $`t`$-expansion , the coupled-cluster method , and others. Nevertheless, it seems likely that Monte Carlo simulations will provide the most robust and accurate numerical techniques in this area also. Our aim in this paper is to discuss some further applications of these Quantum Monte Carlo methods .
The use of quantum Monte Carlo methods in Hamiltonian LGT has a long and somewhat chequered history, and lags a good ten years behind the Euclidean developments. The first calculations used a strong-coupling basis involving discrete “electric field” link variables, and a “Projector Monte Carlo” approach , which used the Hamiltonian itself to project out the ground state. A later version of this was the “stochastic truncation” approach of Allton et al. . Using this approach one can successfully compute string tensions and mass gaps for Abelian models . For non-Abelian models, however, some problems arose . Using an electric field representation for the link variables and a Robson-Webber recoupling scheme at the vertices requires the use of Clebsch-Gordan coefficients or 6$`j`$-symbols, which are not known to high order for $`SU(3)`$; and furthermore, the ‘minus sign’ problem rears its head, in that destructive interference occurs between different paths to the same final state. It may well be that a better choice of strong-coupling basis, such as the ‘loop representation’, might avoid these problems; but this has not yet been demonstrated.
Heys and Stump and Chin et al. pioneered the use of “Greens Function Monte Carlo” (GFMC) or “Diffusion Monte Carlo” techniques in Hamiltonian LGT, in conjunction with a weak-coupling representation involving continuous gauge field link variables. This was successfully adapted to non-Abelian Yang-Mills theories , with no minus sign problem arising. In this representation, however, one is simulating the wave function in gauge field configuration space by a discrete ensemble or density of random walkers: it is not possible to determine the derivatives of the gauge fields for each configuration, or to enforce Gauss’s law explicitly, and the ensemble always relaxes back to the ground state sector. Hence one cannot compute the string tensions and mass gaps directly as Hamiltonian eigenvalues corresponding to ground states in different sectors, as one does in the strong-coupling representation. Chin, Long and Robson thus resorted to “variational Monte Carlo” (VMC) techniques to compute mass gaps. They obtained some reasonable results; but this approach always suffers from the major drawback that there is an unknown systematic error in the results, due to their dependence on the form of the variational wave function which is chosen.
It appears, therefore, that to make unbiased measurements of mass gaps in the weak-coupling representation one is forced back to the more laborious approach used in Euclidean calculations: namely, to measure an appropriate correlation function, and estimate the mass gap as the inverse of the correlation length. In ref. the GFMC method was tested on the (2+1)D U(1) model using a “secondary amplitude” technique to compute expectation values: but this proves to be expensive and prone to bias . In this paper we will show how the standard ‘forward-walking’ technique used in many-body theory can be used for this purpose. The forward-walking method has already been applied to lattice spin models by Runge and Samaras and Hamer . Here we apply it to the compact $`U(1)`$ Yang-Mills theory in (2+1) dimensions, which has been a standard test-bed for Hamiltonian LGT.
A brief discussion of the $`U(1)`$ model is given in Section II. Our Monte Carlo methods are outlined in Section III. The GFMC method is briefly summarized, and then the forward-walking method for estimating expectation values is discussed, together with a technique for measuring timelike correlations. In Section IV the results are presented. Our conclusions are summarized in Section V.
## II Model Hamiltonian
In a weak coupling basis the Hamiltonian for the compact $`U(1)`$ LGT in (2+1)D is given by :
$$H=\underset{l}{}\frac{^2}{A_l^2}2x\underset{P}{}\mathrm{cos}\theta _P$$
(1)
where $`A_l`$ is the gauge field variable on link $`l`$ and
$$\theta _PA_{l_1}+A_{l_2}A_{l_3}A_{l_4}$$
(2)
is the plaquette variable for a lattice plaquette $`P`$, formed by the four links $`l_1,\mathrm{},l_4`$, as illustrated in Fig. 1(a). We consider a periodic, square lattice of linear dimension $`L`$ and lattice spacing $`a`$. The ‘strong coupling’ parameter $`x=1/e^4a^2=1/g^4`$ approaches infinity in the continuum limit $`a0`$.
This is an interesting model, which possesses some important similarities with QCD (for a more extensive review, see for example ref. ). If one takes the ‘naive’ continuum limit at a fixed energy scale, one regains the simple continuum theory of non-interacting photons ; but if one renormalizes or rescales in the standard way so as to maintain the mass gap constant, then one obtains a confining theory of free massive bosons, as discussed by Polyakov , and proven by Göpfert and Mack . The Hamiltonian version of the model has been well studied by a variety of methods: some of the more recent include series expansions , finite-lattice techniques , the $`t`$-expansion , and coupled-cluster techniques , as well as QMC . Quite accurate estimates have been obtained for the string tension and mass gaps, which can be used as comparisons for our Monte Carlo results. The finite-size scaling properties of the model can be predicted using an effective Lagrangian approach combined with a weak-coupling expansion , and the predictions agree very well with finite-lattice data .
## III Monte Carlo Methods
### A Greens Function Monte Carlo
We use the Green’s Function Monte Carlo \[GFMC\] method , which was adapted to the $`U(1)`$ model by Heys and Stump , and Chin et al. . A brief summary of the method can be given as follows.
In a weak-coupling representation, the basis states are taken to be eigenstates of the plaquette angles $`\theta _P`$, which can take continuous values. The Hamiltonian (1) can be written as
$$H=\underset{l}{}\frac{^2}{A_l^2}+V(\mathrm{\Theta }),$$
(3)
where
$$V(\mathrm{\Theta })=2x\underset{P}{}\mathrm{cos}\theta _P,$$
(4)
and the plaquette angles $`\theta _P`$ and link angles $`A_l`$ are related by equation (2). The imaginary time Schrödinger equation for the system is
$$\frac{}{\tau }\mathrm{\Phi }(\mathrm{\Theta },\tau )=[\underset{l}{}\frac{^2}{A_l^2}+V(\mathrm{\Theta })E_\text{T}]\mathrm{\Phi }(\mathrm{\Theta },\tau ),$$
(5)
where $`E_\text{T}`$ is a trial energy, representing a constant shift in the zero of energy, which will prove useful. The imaginary time evolution operator $`\mathrm{exp}[(HE_\text{T})\tau ]`$ acts as a projector onto the ground state $`|\mathrm{\Phi }_0`$:
$$|\mathrm{\Phi }_0=\underset{\tau \mathrm{}}{lim}e^{\tau (HE_\text{T}}|\mathrm{\Psi }_\text{T}$$
(6)
for any trial state $`|\mathrm{\Psi }_\text{T}`$, provided that $`|\mathrm{\Psi }_\text{T}`$ is not orthogonal to $`|\mathrm{\Phi }_0`$.
Equation (5) is a diffusion equation in configuration space, and is easily simulated by the Green’s Function Monte Carlo method. It is assumed that the ground-state wave function can be chosen positive everywhere, and it is simulated by the density distribution of an ensemble of random walkers $`\{\mathrm{\Theta }_i\}`$ in configuration space, with weights {$`w_i`$}. The first term on the right of Eq. (5) produces diffusion, and is simulated by a Gaussian random walk of the members of the ensemble as time proceeds, while the term $`[V(\mathrm{\Theta })E_\text{T}]`$ produces a growth or decay in the density which is simulated by a branching process.
### B Variational Guidance
The efficiency and accuracy of the simulation are greatly enhanced by the use of variational guidance or importance sampling . Let $`\mathrm{\Psi }_\text{T}(\mathrm{\Theta })`$ be a variational approximation to the true ground-state wave function, and define a new probability distribution
$$f(\mathrm{\Theta },\tau )=\mathrm{\Phi }(\mathrm{\Theta },\tau )\mathrm{\Psi }_\text{T}(\mathrm{\Theta }),$$
(7)
Then the modified imaginary time Schrödinger equation for $`f(\mathrm{\Theta },\tau )`$ reads
$$\frac{f}{\tau }=\underset{l}{}\frac{^2f}{A_l^2}+[E_\text{L}(\mathrm{\Theta })E_\text{T}]f+\underset{l}{}\frac{}{A_l}(fF_{Ql}(\mathrm{\Theta })),$$
(8)
where
$$E_\text{L}(\mathrm{\Theta })=\frac{1}{\mathrm{\Psi }_\text{T}(\mathrm{\Theta })}H\mathrm{\Psi }_\text{T}(\mathrm{\Theta })$$
(9)
is the local energy obtained from the trial function, and
$$F_{Ql}(\mathrm{\Theta })\frac{2}{\mathrm{\Psi }_\text{T}(\mathrm{\Theta })}\frac{\mathrm{\Psi }_\text{T}(\mathrm{\Theta })}{A_l}$$
(10)
is a “quantum force” term, which produces a directed drift in the ensemble towards the configurations favoured by the trial wave function. By a good choice of $`\mathrm{\Psi }_\text{T}`$ and $`E_\text{T}`$ the “excess local energy” term $`[E_\text{L}(\mathrm{\Theta })E_\text{T}]`$ can be made very small, which reduces the amount of branching necessary, and reduces the statistical fluctuations in the results.
For small time steps $`\mathrm{\Delta }\tau `$, an approximate Green’s function solution to Eq. (8) is
$`G(\mathrm{\Theta }\mathrm{\Theta }^{},\mathrm{\Delta }\tau )`$ $``$ $`\mathrm{exp}\{[E_\text{L}(\mathrm{\Theta })E_\text{T}]\mathrm{\Delta }\tau \}`$ (12)
$`\times {\displaystyle \underset{l}{}}\left({\displaystyle \frac{1}{\sqrt{\overline{4\pi \mathrm{\Delta }\tau }}}}\mathrm{exp}\{[A_l^{}A_l\mathrm{\Delta }\tau F_{Ql}(\mathrm{\Theta })]^2/4\mathrm{\Delta }\tau \}\right).`$
In the Monte Carlo simulation, each iteration corresponds to a time step $`\mathrm{\Delta }\tau `$. At each iteration, we sweep through each link in turn, and simulate the corresponding exponential factor in the sum on the right of (12) by a random displacement of the link variable for each walker:
$$\mathrm{\Delta }A_l=\mathrm{\Delta }\tau F_{Ql}(\mathrm{\Theta })+\chi ,$$
(13)
where $`\chi `$ is randomly chosen from a Gaussian distribution with standard deviation $`\sqrt{\overline{(2\mathrm{\Delta }\tau )}}`$. The first term in (13) is the “drift” term, and the second is the “diffusion” term. The first exponential on the right of (12) is simulated by simply multiplying the “weight” of each walker $`w_i`$ by an equivalent amount.
At the end of each iteration, the trial energy $`E_\text{T}`$ is adjusted to compensate for any change in the total weight of all walkers in the ensemble; and a “branching” process is carried out, so that walkers with weight greater than (say) 2 are split into two new walkers, while any two walkers with weight less than (say) 1/2 are combined into one, chosen randomly according to weight from the originals. This procedure of “Runge smoothing” maximizes statistical accuracy by keeping the weights of all the walkers within fixed bounds, while mimimizing any fluctuations in the total weight due to the branching process.
When equilibrium is reached after many sweeps through the lattice, the average value of the trial energy $`E_\text{T}`$ will give an estimate of the ground-state energy $`E_0`$, and the weight density of the ensemble in configuration space will be proportional to $`\mathrm{\Phi }_0\mathrm{\Psi }_\text{T}`$. Various corrections due to the finite time interval $`\mathrm{\Delta }\tau `$ have been ignored in this discussion, and the limit $`\mathrm{\Delta }\tau \mathrm{}`$ must be taken in some fashion to eliminate such corrections.
In the simulations presented here, a trial function for the ground state was chosen as
$$\psi _\text{T}(\mathrm{\Theta })=\mathrm{exp}\left[c\underset{P}{}\mathrm{cos}\theta _P+d\underset{<PP^{}>}{}\mathrm{cos}(\theta _P+\theta _P^{})\right]$$
(14)
with two variational parameters $`c`$ and $`d`$, where the sum over $`PP^{}`$ denotes a sum over nearest-neighbour pairs of plaquettes, forming rectangles ($`1\times 2`$ Wilson loops) on the lattice. That is, the uncorrelated single plaquette trial function used in has been extended to include a correlated term. In the single plaquette case ($`d=0`$) it is straightforward to determine the optimal value for $`c`$ . For the correlated case ($`d0`$), the optimum values for $`c`$ and $`d`$ must be found by means of a separate Variational Monte Carlo (VMC) calculation, where the variational energy
$$E_{\text{VMC}}(c,d)=\frac{(\mathrm{\Psi }_\text{T}(\mathrm{\Theta }))^2E_\text{L}(\mathrm{\Theta })D\mathrm{\Theta }}{(\mathrm{\Psi }_\text{T}(\mathrm{\Theta }))^2D\mathrm{\Theta }}$$
(15)
is minimised. The procedure is described in where a six- parameter variational wave function was used to estimate the mass gaps in the $`U(1)`$ model at the VMC level. In short, one samples the local trial energy over configurations $`\{\mathrm{\Theta }\}`$ randomly generated by a Metropolis procedure, distributed with respect to the weight function $`(\mathrm{\Psi }_\text{T}(\mathrm{\Theta }))^2`$. A downhill simplex method was used to perform the minimisation. To this end $`E_{\text{VMC}}(c,d)`$ was evaluated in a region centered around a guessed minimum $`(c_0,d_0)`$ using the reweighting procedure described in , whereby uncorrelated configurations generated with respect to the weight function $`(\mathrm{\Psi }(\mathrm{\Theta })_{c=c_0,d=d_0})^2`$ were reweighted to give a distribution corresponding to the required weight function for a pair $`(c,d)`$ in the neighbourhood of $`(c_0,d_0)`$. This procedure was iterated until the minimum $`(c_0,d_0)`$ converged. We used 100,000 independent configurations on an $`8\times 8`$ lattice to determine $`E_{\text{VMC}}(c,d)`$. This allowed the optimal $`c`$ and $`d`$ to be fixed to within around 1%. Though it is possible to optimise $`c`$ and $`d`$ for each lattice size, for a given value of the coupling $`x`$, we used the same values of $`c`$ and $`d`$ for all lattice sizes $`L`$. The values of $`c`$ and $`d`$ used are listed in Table I.
The use of a correlated trial wave function is expected to markedly decrease the variance of estimators and the amount of branching in the population smoothing process, particularly in the physically interesting weak coupling regime ($`x`$ large). Moreover, because the trial function is closer to the true ground state, fewer forward walking iterations (see section III C) are required in order to derive converged estimators for correlation functions. These advantages come at a price in that updating of GFMC configurations becomes more complicated: The local trial energy is
$$E_\text{L}(\mathrm{\Theta })=(4c2x)\underset{P}{}\mathrm{cos}\theta _P+6d\underset{PP^{}}{}\mathrm{cos}\left(\theta _P+\theta _P^{}\right)\frac{1}{4}\underset{l}{}F_{Ql}^2,$$
(16)
and the quantum force term is given by
$`F_{Ql}`$ $`=`$ $`c(\mathrm{sin}\theta _{P_1(l)}\mathrm{sin}\theta _{Q_1(l)})+d{\displaystyle \underset{j=2}{\overset{4}{}}}\{\mathrm{sin}(\theta _{P_1(l)}+\theta _{P_j(l)})`$ (18)
$`\mathrm{sin}(\theta _{Q_1(l)}+\theta _{Q_j(l)})\}.`$
Here, $`P_1(l)`$,…,$`P_4(l)`$ and $`Q_1(l)`$,…,$`Q_4(l)`$ denote the 4 closest plaquettes on either side of the link $`l`$, as illustrated in Fig. 1(b). These expressions are more difficult to compute than in the single parameter ($`d=0`$) case. However, the increased complexity is easily offset by the gains from variance reduction.
### C Forward-walking method
The quickest method of estimating an expectation value $`Q_0`$ is simply to form the weighted average of $`Q`$ over the ensemble of random walkers (we assume Q is diagonal in the chosen basis, for simplicity). This produces an estimate according to the distribution $`\mathrm{\Psi }_\text{T}\mathrm{\Phi }_0`$, rather than $`\mathrm{\Phi }_0^2`$. The estimate can be perturbatively improved , but there remains an unknown systematic error due to the dependence on the trial function $`\mathrm{\Psi }_\text{T}`$. For this reason, we have preferred to use the so-called “forward-walking” method for estimating expectation values.
The forward-walking method is a robust technique for estimating expectation values , based on the following equation for an operator $`Q`$:
$`<Q>_0`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Phi }_0|Q|\mathrm{\Phi }_0}{\mathrm{\Phi }_0|\mathrm{\Phi }_0}}`$ (19)
$`\stackrel{}{J\mathrm{}}`$ $`{\displaystyle \frac{\mathrm{\Psi }_\text{T}|K^JQ|\mathrm{\Phi }_0}{\mathrm{\Psi }_\text{T}|K^J|\mathrm{\Phi }_0}}`$ (20)
$`=`$ $`{\displaystyle \frac{\stackrel{~}{K}(\mathrm{\Theta }_J,\mathrm{\Theta }_{J1})\mathrm{}\stackrel{~}{K}(\mathrm{\Theta }_2,\mathrm{\Theta }_1)Q(\mathrm{\Theta }_1)\stackrel{~}{\mathrm{\Phi }}_0(\mathrm{\Theta }_1)}{\stackrel{~}{K}(\mathrm{\Theta }_J,\mathrm{\Theta }_{J1})\mathrm{}\stackrel{~}{K}(\mathrm{\Theta }_2,\mathrm{\Theta }_1)\stackrel{~}{\mathrm{\Phi }}_0(\mathrm{\Theta }_1)}}`$ (21)
where $`K(\mathrm{\Theta }_J,\mathrm{\Theta }_{J1})`$ is the evolution operator for time $`\mathrm{\Delta }\tau `$, and $`\stackrel{~}{K}(\mathrm{\Theta }_J,\mathrm{\Theta }_{J1})`$ is the same operator in the similarity transformed basis. Again we have assumed that the operator $`Q`$ is diagonal in the basis of plaquette variables $`\mathrm{\Theta }`$.
This equation is implemented by the following procedure:
* Starting from the trial state, iterate until equilibrium is achieved, then begin a measurement;
* Record the value $`Q(\mathrm{\Theta }_i)`$ for each walker (“ancestor”) at the beginning of the measurement;
* Propagate the ensemble as normal for $`J`$ iterations, keeping a record of the “ancestor” of each walker in the current population;
* Take the weighted average of the $`Q(\mathrm{\Theta }_i)`$ with respect to the weights of the descendants of $`\mathrm{\Theta }_i`$ after the $`J`$ iterations, using sufficient iterations $`J`$ that the estimate reaches a ‘plateau’.
This procedure has been tested for the magnetization in the 2D lattice Heisenberg model by Runge , and for correlation functions in the 1D transverse Ising model by Samaras and Hamer , and works very well. The drawback to the procedure is that after a large number of iterations $`J`$ many of the “ancestors” will die out, leaving no descendants, which leads to a progressive loss of statistical accuracy. Thus it is even more crucial in this connection to use a good guiding wavefunction $`\mathrm{\Psi }_\text{T}`$ so as to minimize “branching”.
### D Timelike Correlations
In order to estimate mass gaps in this model, we have again used the forward-walking technique to measure correlations between operators at different times. The mass gaps can then be estimated from the decay constants for these correlation functions.
A timelike correlator $`Q_1(\tau )Q_2(0)_0`$ can be found from the following equations:
$$<Q_1(\tau )Q_2(0)>_0=\frac{<\mathrm{\Phi }_0|Q_1e^{H\tau }Q_2|\mathrm{\Phi }_0>}{<\mathrm{\Phi }_0|\mathrm{\Phi }_0>}$$
(22)
$$\stackrel{}{J\mathrm{}}\frac{<\mathrm{\Psi }_\text{T}|K^JQ_1K^NQ_2|\mathrm{\Phi }_0>}{<\mathrm{\Psi }_\text{T}|K^{J+N}|\mathrm{\Phi }_0>}$$
(23)
$$=\frac{\stackrel{~}{K}^JQ_1\stackrel{~}{K}^NQ_2\stackrel{~}{\mathrm{\Phi }}_0}{\stackrel{~}{K}^{J+N}\stackrel{~}{\mathrm{\Phi }}_0}$$
(24)
where $`N\mathrm{\Delta }\tau =\tau `$. This can be implemented in much the same way as an expectation value (assuming both $`Q_1`$ and $`Q_2`$ are diagonal in the weak-coupling representation). At the beginning of the measurement, record the ‘ancestor’ configurations as in Sec. (III C). Then allow the ensemble to propagate for time $`\tau `$, with the branching process turned off so that each state retains its identity. At the end of time $`\tau `$, record the initial value $`Q_1`$ and the final value $`Q_2`$. Propagate each state for a further $`J`$ iterations as before, and then average, weighting each ‘ancestor’ state according to the forward-walking prescription of Sec. (III C).
## IV Results
Simulations were carried out for $`L\times L`$ lattices up to $`L=16`$ sites, using runs of 10,000 walkers over 50,000 iterations. At each iteration several sweeps of the lattice were performed, after which Runge smoothing (the branching (combining) of high (low) weight walkers) was imposed on the walker population. The average percentage of walkers branched/combined per iteration depends on a number of factors: clearly, the time step $`\mathrm{\Delta }\tau `$ and number of lattice sweeps performed per iteration can be adjusted in order to control the extent of the branching. We wish to choose values of $`\mathrm{\Delta }\tau `$ which are sufficiently small that time discretization errors can be made negligible. Furthermore, the number of sweeps performed per iteration was chosen so that essentially uncorrelated measurements could be made roughly every 150 iterations. We found that, for the lattice sizes $`L`$ and couplings $`x`$ considered, this could be achieved in such a way that on average no more than 10% of the walkers were branched/combined per iteration. However, the extent of the smoothing per iteration depends on the quality of the variational guiding function. Generally, the two-parameter guiding function (14) works better for smaller couplings $`x`$ and smaller lattice sizes $`L`$. Under the most strained conditions considered ($`x=4`$, $`L=16`$) around 10% of the walkers were branched/combined per iteration for the values of $`\mathrm{\Delta }\tau `$ chosen and the number of sweeps per iteration required to get independent measurements, whereas for $`x1`$ the ratio was a fraction of a percent. The first 2000 iterations were discarded to allow for equilibration. Two values of the time step $`\mathrm{\Delta }\tau `$ were used, differing by a factor of 5 (e.g. $`\mathrm{\Delta }\tau =0.02`$ and $`\mathrm{\Delta }\tau =0.005`$). The number of sweeps performed for the smaller $`\mathrm{\Delta }\tau `$ value was 5 times as large as that for the larger value and as a result measurements were equally uncorrelated in the two cases. The typical number of sweeps used per iteration, e.g. for $`\mathrm{\Delta }\tau =0.02`$ and $`\mathrm{\Delta }\tau =0.005`$ was 1 and 5 respectively. Linear extrapolations were made to the limit $`\mathrm{\Delta }\tau =0`$. The linear dependence on $`\mathrm{\Delta }\tau `$ has been checked previously .
### A Ground-state Energy
Results for the ground-state energy per site $`E_0/L^2`$ at coupling $`x=4`$ are graphed against $`1/L^3`$ in Fig. 2. They agree within errors with those of the earlier study up to $`L=10`$, and extend them to $`L=16`$. It can be seen that the data are fitted quite well by the form
$$E_0/L^2=4.414(1)\frac{2.53}{L^3},$$
(25)
which compares very well with effective Lagrangian theory and the weak-coupling series prediction
$$E_0/L^2=4.413(2)\frac{2.48(1)}{L^3},$$
(26)
The estimated bulk value $`4.414(1)`$ compares well with estimates of $`4.43(2)`$ from strong-coupling series , $`4.415(6)`$ from a $`t`$-expansion , and $`4.412`$ from the coupled cluster method (CCM) .
### B Wilson Loops
The expectation values of Wilson loops were computed using the forward-walking method outlined in section III C. Measurements of the observables were made in cycles starting around every 150 iterations, with typically 10 forward-walking weighted averages over the ensemble taken at time steps of 12–15 iterations. For each time step the weighted averages were then block averaged over successive measurement cycles. Fig. 3(a) shows the forward-walking convergence of the $`3\times 3`$ Wilson loop $`W(3,3)`$ measured on the $`16\times 16`$ lattice at coupling $`x=4`$ for a trial run involving 1200 walkers, 20,000 iterations, a time step $`\mathrm{\Delta }\tau =0.005`$, 3 sweeps of the lattice per iteration, with 12 forward-walking weighted averages being taken every 20 iterations in measurement cycles made 250 iterations apart. Results are shown for two different guiding functions, using the 1-parameter and 2-parameter forms respectively. It can be seen that in both cases the data relax exponentially towards a common equilibrium value, which can be estimated by making an exponential fit to the data. The equilibrium value is then taken as the final result for the Wilson loop. It can also be seen that the 2-parameter form reduces the variance and produces much more rapid convergence to the asymptotic value.
Typical results from a production run (using the 2-parameter guiding function) are shown in Fig. 3(b). In this case $`L=12`$, $`x=2`$ and 50,000 iterations are performed with 10,000 walkers. Measurement cycles involving 10 forward-walking weighted averages 16 iterations apart were started every 165 iterations. Results for two values of $`\mathrm{\Delta }\tau `$ (0.05 and 0.01), using 1 and 5 lattice sweeps per iterations respectively, are shown.
Fig. 4 displays the values for the ‘mean plaquette’ $`P=W(1,1)`$ at coupling $`x=4`$, graphed against $`1/L^3`$. There is evidently a discrepancy here between our present results and those of , obtained using the secondary amplitude technique. Either the errors in have been underestimated for the larger lattices and couplings, or the discrepancy could have arisen from bias in the secondary amplitude technique . Our present results show a consistent finite-size scaling behaviour for $`P(L)`$:
$$P(L)=0.7584(4)+\frac{0.22}{L^3},$$
(27)
to be compared with the weak-coupling series prediction
$`P(L)`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{d}{dx}}[E_0/L^2]`$ (28)
$`=`$ $`0.7593(3)+{\displaystyle \frac{0.183(3)}{L^3}},`$ (29)
obtained using the Hellmann-Feynman theorem. The agreement is once again quite good. Extrapolation to the axis gives a bulk limit of $`P=0.7584(4)`$, to be compared with strong-coupling series 0.80(3), the $`t`$-expansion 0.757(4), and the CCM 0.7585. The results for other Wilson loops scale similarly with lattice size. For example, the finite-size scaling of the $`4\times 4`$ Wilson loop is shown in Fig. 5.
Of course, because the model has a gap for any finite value of $`x`$, one should expect to see a crossover from this algebraic (free-photon theory) scaling to exponential convergence for large enough $`L`$. Indeed, for smaller values of the coupling, exponential scaling sets in rapidly enough that essentially bulk results can be obtained on relatively small lattices. This is illustrated in Fig. 6 in the case $`x=2`$ where we see a definite crossover from algebraic to exponential scaling for the mean plaquette and the $`2\times 2`$ Wilson loop when $`L`$ reaches around 10.
Table II lists the final estimates we have obtained for the bulk ground-state energy per site and Wilson loop values at some selected couplings.
### C String Tension
In a confining model, the Wilson loops are expected to behave asymptotically as
$$W(m,n)\mathrm{exp}[KA]$$
(30)
where $`A=mn`$ is the area of the loop and $`K`$ is the string tension. In Fig. 7 we illustrate this behaviour for the $`x=2`$ case on the $`12\times 12`$ lattice by plotting the Wilson loops $`W(n,n)`$ and $`W(n,n1)`$ as a function of area. The standard estimator of the string tension is the Creutz ratio
$$KR_n=\mathrm{ln}\left[\frac{W(n,n)W(n1,n1)}{W(n,n1)^2}\right]$$
(31)
Fig. 8(a) shows the Creutz ratios $`R_n`$ graphed against $`1/\stackrel{~}{A}`$ (where $`\stackrel{~}{A}n(n1)`$ is the average area of the loops used to form the ratio) at $`x=4`$, from which we obtain the rough estimate
$$K=0.05(2)$$
(32)
Because the Creutz ratios decrease very rapidly with $`n`$, and the relative error correspondingly increases, it is very difficult to determine a reliable value for the string tension. Given the limited number of data points, it is probably more fruitful to use a two-point estimator
$$KR_n^{}=\frac{1}{n}\mathrm{ln}\left[\frac{W(n,n)}{W(n,n1)}\right],$$
(33)
formed from successive pairs of Wilson loops. Included in Fig. 8(a) are plots of these estimators against $`1/\stackrel{~}{A}`$. Unfortunately, the two-point estimators display oscillations (which the Creutz ratios were designed to eliminate) so it is again difficult to perform a linear extrapolation to obtain an estimate of the string tension. Nevertheless, a linear fit through the points gives an estimate of the string tension of $`K0.03`$, somewhat lower than the result above.
We have performed a similar analysis for $`x=2`$, working with the $`12\times 12`$ lattice which is sufficiently close to the bulk limit. The results for the string tension estimators are shown in Fig. 8(b). In this case a linear extrapolation of the two-point estimators yields an estimate of $`K0.08`$ for the string tension.
The string tension has previously been calculated as an energy per unit length, $`\sigma `$. The best available values at $`x=2`$ are $`\sigma =0.28(1)`$ from stochastic truncation , and $`\sigma =0.282(2)`$ from an ‘exact linked cluster expansion’ . The quantities are related by
$$\sigma =vK$$
(34)
where $`v`$ is the speed of light, estimated as $`v=2.27`$ at $`x=2`$ from the weak-coupling expansion . Hence $`\sigma =0.28(1)`$ corresponds to $`K=0.12(1)`$. This is in rough agreement with the result above.
We also attempted the analysis for smaller values of $`x`$. However, because the Wilson loops decay extremely rapidly in these cases, it is not possible to obtain enough data points to make sensible extrapolations.
### D Mass Gaps
We have attempted to estimate mass gaps from the exponential decay of correlation functions, as is done in the Euclidean approach. The lowest-lying excitations in the strong-coupling limit are the single-plaquette excitations, antisymmetric or symmetric under reflections respectively. At first we attempted to use the spatial correlations between plaquette operators for this purpose; but the signal from the connected spatial correlations turns out to be very small, and can even change sign, giving no good exponential decay signal. We can understand this problem as follows: : we are restricted to a plane in the spatial directions, and can therefore only measure correlations between plaquette operators “edge on” to each other, as illustrated in Fig. 9(a), as opposed to the “face on” configuration in Fig. 9(b). This means there is very poor overlap with the configurations of interest in the intermediate state.
We are thus compelled to fall back on the correlations in imaginary time $`\tau `$, which do correspond to “face on” plaquettes as in Fig. 9(b). Using the method of III D, we have measured correlation functions
$$f(\tau )=(\underset{P}{}P(0))(\underset{P^{}}{}P^{}(\tau ))_0$$
(35)
where the plaquette operators have been summed over all positions to project out the zero-momentum intermediate states, and the sums are either symmetric or antisymmetric under reflections.
As with the Wilson loop calculations, typically 50,000 iterations were performed with 10,000 walkers and forward-walking measurment cycles were started every 150 iterations. As mentioned, at the start of each measurement cycle, Runge smoothing was switched off so that no walkers were created or destroyed while timelike correlation functions were calculated. Between 10 and 20 timelike correlations (between the plaquette operators at the start of the measurement cycle $`\tau =0`$ and later times $`\tau `$)were measured at intervals of one or two iterations. For each of the timelike measurements the “physical time” $`\tau `$ is given by:
$$\tau =\mathrm{\Delta }\tau \times N_{\text{SWEEP}}\times N_{\text{INTERVAL}}$$
where $`\mathrm{\Delta }\tau `$ is the basic timestep, $`N_{\text{SWEEP}}`$ is the number of sweeps of the lattice and $`N_{\text{INTERVAL}}`$ denotes the number of iterations between the calculation of initial and final plaquette operators.
At this point Runge smoothing was switched on again and forward-walking weighted averages of the timelike correlations were calculated over most of the measurement cycle at intervals of 10–20 iterations. Again the weighted averages were block averaged over all measurement cycles, and exponential fitting was used to find the forward-walking relaxed estimates of the timelike correlation functions. As with the Wilson loops, two values of the time step $`\mathrm{\Delta }\tau `$ were used and linear extrapolation was used to obtain $`\mathrm{\Delta }\tau =0`$ estimates of the correlators. Again, the two values of $`\mathrm{\Delta }\tau `$ used differed by a factor of 5, and 5 times as many lattice sweeps were performed per iteration for the smaller $`\mathrm{\Delta }\tau `$ so that the “physical” times $`\tau `$ over which the correlations $`f(\tau )`$ were measured were the same (i.e. $`\mathrm{\Delta }\tau \times N_{\text{SWEEP}}`$ was kept fixed).
A typical plot of the resulting estimates for the timelike correlation function $`f(\tau )`$, as a function of $`\tau `$, is shown in Fig. 10 for the antisymmetric correlator with $`x=0.5`$ and $`L=4`$. It can be seen that $`f(\tau )`$ shows a smooth exponential decay with $`\tau `$. Defining an ‘effective mass’ for each pair of points by
$$m(\tau )=\mathrm{ln}[f(\tau )/f(\tau 1)],$$
(36)
one finds that the effective mass decreases somewhat with $`\tau `$, as might be expected. An ad hoc fit in $`1/\tau `$, illustrated in Fig. 11, gives a final estimate for the mass of the corresponding intermediate state. These values are listed as a function of lattice size in Table III for couplings $`x=0.5`$, 1 and 2. For the higher couplings the exponential decay is so slow that we have been unable to obtain useful results for the mass. To do so would require measurements to be made over substantially larger time scales with significantly greater accuracy. Unfortunately, our results were not sufficiently accurate to sensibly resolve the form of the lattice size dependence for the couplings studied. In Table III we list results of estimates of the bulk mass from other sources. It can be seen that for the couplings considered, the agreement between methods is reasonable.
## V Summary and Conclusions
In this paper we have shown how the standard ‘forward-walking’ methods of quantum many-body theory can be used to calculate expectation values and correlation functions in Hamiltonian lattice gauge theory. Accurate values for the Wilson loops have been obtained, and their finite-size scaling behaviour has been demonstrated. The string tension can then be estimated using two-point and Creutz ratios. The Wilson loop values dropped away too quickly with lattice size to give accurate string tensions at small couplings $`x`$; but this problem could presumably be rectified by looking at timelike correlators—either timelike Wilson loops, or timelike correlations between Polyakov loops.
We have also shown how mass gaps can be estimated from the exponential decay of timelike correlators. This approach is very reminiscent of the techniques used in the Euclidean regime: our freedom of choice of the timelike measurement interval is akin to the choice of an anisotropic lattice in the ‘timelike’ direction in the Euclidean case. Reasonable results have been obtained for the masses at small couplings $`x`$, but not at large $`x`$.
The accuracy of the final results for the string tension and mass gaps is not remarkable; but it must be recalled that this is the first time that unbiased estimates have been obtained for these quantities using a ‘weak-coupling’ representation and GFMC techniques. There are no doubt many ways by which the results could be improved: for example, by the use of improved guiding wavefunctions, or improved (“fuzzed”) correlation functions, or improved Hamiltonian operators. Our aim here has been to ‘prove the concept’, that forward-walking techniques can give reliable estimates for the correlation functions. In future work we plan to apply similar methods to the (3+1)D $`SU(3)`$ Yang-Mills theory.
## Acknowledgments
This work is supported by the Australian Research Council. Calculations were performed on the SGI Power Challenge Facility at the New South Wales Center for Parellel Computing and the Fujitsu VPP300 vector machine at the Australian National Universtiy Supercomputing Facility. |
warning/0002/math0002048.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Yangians $`Y(g)`$ were introduced by V.G.Drinfeld as the quantizations of rational solutions of the classical Yang-Baxter equation (CYBE). It was demonstrated by A.Stolin that rational solutions with values in the tensor square of the Lie algebra $`g`$ space can be written in the form
$$\frac{C^2}{uv}+a_0+b_1u+b_2v+cuv$$
with $`a_0,b_i,cg`$ and $`C^2`$ – the second Casimir of the algebra $`g`$. Notice that this does not necessarily signify the simplicity of $`g`$. The existence of Yangians for nonsemisimple Lie algebras was predicted in and explicitly demonstrated in .
The important class of rational solutions is of the form
$$\frac{C^2}{uv}+r_0,$$
where the additional summand is a constant solution of the CYBE. As it was proved in the quantization of such solutions leads to the twist deformations $`Y_{}(g)`$ of the Yangians $`Y(g)`$ (the latter refer to the canonical rational solutions $`\frac{C^2}{uv}`$). The deformed Yangian $`Y_{}(g)`$ has the same multiplication as in $`Y(g)`$, the twisted coproduct
$$\mathrm{\Delta }_{}\left(y\right)=\mathrm{\Delta }\left(y\right)^1,$$
and the transformed $``$-matrix
$$_{}=\left(\right)_{21}^1.$$
The twisting element $``$ has to satisfy the equations :
$$\begin{array}{c}_{12}(\mathrm{\Delta }id)()=_{23}(id\mathrm{\Delta })(),\hfill \\ (ϵid)()=(idϵ)()=1.\hfill \end{array}$$
(1)
When $`𝒜`$ and $``$ are the universal enveloping algebras: $`𝒜=U(\text{})=U(\text{})`$ with $`\text{}\text{}`$and $`U(\text{})`$ is the minimal subalgebra on which $``$ is completely defined as $`U(\text{})U(\text{l})`$ then l is called the carrier algebra for $``$.
The existence of a twist can be formulated in terms of a nondegenerate bilinear form on the carrier algebra. The generators dual to the PBW basic elements of $`U(g)`$ are very important. They provide the explicit presentation of the twisting element .
The two known examples of twists that were written explicitly correspond to the two-dimensional carrier subalgebra $`B(2)`$ with the generators $`H`$ and $`E`$,
$$[H,E]=E,$$
and the four-dimensional carrier subalgebra $`𝐋`$:
$$\begin{array}{c}[H,E]=E,[H,A]=\alpha A,[H,B]=\beta B,\hfill \\ [A,B]=E,[E,A]=[E,B]=0,\alpha +\beta =1.\hfill \end{array}$$
(2)
The first one is called the Jordanian twist and has the twisting element
$$\mathrm{\Phi }_𝒥=e^{H\sigma },\sigma =\mathrm{ln}(1+E).$$
(3)
The second one is the extended Jordanian twist suggested in . The corresponding twisting element
$$_{(\alpha ,\beta )}=\mathrm{\Phi }_{(\alpha ,\beta )}\mathrm{\Phi }_𝒥$$
(4)
contains the Jordanian factor $`\mathrm{\Phi }_𝒥`$ and the extension
$$\mathrm{\Phi }_{(\alpha ,\beta )}=\mathrm{exp}\{ABe^{\beta \sigma }\}.$$
(5)
In general the composition of two twists is not a twist. But there are some important examples of the opposite behaviour. When $`𝐋`$ is a subalgebra $`𝐋\text{}`$ there may exist several pairs of generators of the type $`(A,B)`$ arranged so that the Jordanian twist can acquire several similar extensions . This demonstrates that some twistings can be applied successively to the initial Hopf algebra even in the case when their carrier subalgebras are nontrivially linked. In the universal enveloping algebras for classical Lie algebras there exists the possibility to construct systematically the special sequences of twists called chains :
$$_{_{p0}}__p_{_{p1}}\mathrm{}__0.$$
(6)
The factors $`__k=\mathrm{\Phi }__k\mathrm{\Phi }_{𝒥_k}`$ of the chain are the twisting elements of the extended Jordanian twists for the initial Hopf algebra $`𝒜_0`$. Here the extensions $`\left\{\mathrm{\Phi }__k,k=0,\mathrm{}p1\right\}`$ contain the fixed set of normalized factors $`\mathrm{\Phi }_{(\alpha ,\beta )}=\mathrm{exp}\{ABe^{\beta \sigma }\}`$ , the full set. It was proved that in the classical Lie algebras that conserve symmetric invariant forms such chains can be made maximal and proper. This means that for the algebras $`U(A_N)`$, $`U(B_N)`$ and $`U(D_N)`$ of the three classical series there exist chains $`_{_{p0}}`$ that cannot be reduced to a chain for a simple subalgebra and their full sets of extensions are the maximal sets in the sense described below.
To construct a maximal proper chain for $`𝒜=U(g)`$ (where $`g`$ is a classical Lie algebra with the root system $`\mathrm{\Lambda }_𝒜`$) the sequences $`𝒜𝒜_0𝒜_1\mathrm{}𝒜_{p1}𝒜_p`$ of Hopf subalgebras are to be fixed in $`𝒜`$. For each element $`𝒜_k`$ of the sequence there must exist the so called initial root $`\lambda _0^k`$ and the set $`\pi _k`$ of its constituent roots,
$$\pi _k=\left\{\lambda ^{},\lambda ^{\prime \prime }\right|\lambda ^{}+\lambda ^{\prime \prime }=\lambda _0^k;\lambda ^{}+\lambda _0^k,\lambda ^{\prime \prime }+\lambda _0^k\neg \mathrm{\Lambda }_𝒜\}$$
(7)
For any $`\lambda ^{}\pi _k`$ there must be an element $`\lambda ^{\prime \prime }\pi _k`$ that $`\lambda ^{}+\lambda ^{\prime \prime }=\lambda _0^k`$. So, $`\pi _k`$ is naturally decomposed as
$$\pi _k=\pi _k^{}\pi _k^{\prime \prime },\pi _k^{}=\{\lambda ^{}\},\pi _k^{\prime \prime }=\{\lambda ^{\prime \prime }\}.$$
(8)
The triples $`(𝒜_k,\lambda _0^k,\pi _k)`$ are subject to the following conditions:
1. $`\lambda _0^k`$ must be orthogonal to the roots of the subalgebra $`𝒜_{k1}`$,
2. the subsets $`\pi _k^{}`$ and $`\pi _k^{\prime \prime }`$ must form the diagrams of conjugate representations for $`𝒜_{k1}`$.
In these terms the factors $`__k`$ of the chain (6) are fixed as follows:
$$__k=\mathrm{\Phi }__k\mathrm{\Phi }_{𝒥_k}$$
(9)
with
$$\mathrm{\Phi }_{𝒥_k}=\mathrm{exp}\{H_{\lambda _0^k}\sigma _0^k\},\sigma _0^k=\mathrm{ln}(1+L_{\lambda _0^k});$$
(10)
$$\mathrm{\Phi }__k=\underset{\lambda ^{}\pi _k^{}}{}\mathrm{\Phi }__\lambda ^{}=\underset{\lambda ^{}\pi _k^{}}{}\mathrm{exp}\{L_\lambda ^{}L_{\lambda _0^k\lambda ^{}}e^{\frac{1}{2}\sigma _0^k}\}$$
(11)
(here $`L_\lambda `$ is the generator associated with the root $`\lambda `$).
Quantizations $`𝒜_{_{_{p0}}}`$ of classical universal enveloping algebras produce the chains of $`_{}`$-matrices:
$$_{_{p0}}=(__p)_{21}(_{_{p1}})_{21}\mathrm{}(__0)_{21}__0^1\mathrm{}_{_{p1}}^1__p^1.$$
(12)
The deformation parameters can be introduced in chains by rescaling the generators in the subalgebras $`_k`$. Each $`_k`$ can be rescaled separately with an independent variable $`\xi _k`$. When all these scaling factors are proportional to the deformation parameter $`\xi `$, i.e. $`\xi _k=\xi \eta _k`$, then in the classical limit the parameters $`\eta _k`$ appear as the multipliers in the classical $`r`$-matrix:
$$r_{_{p0}}=\underset{k=0,1,\mathrm{},p}{}\eta _k\left(H_{\lambda _0^k}L_{\lambda _0^k}+\underset{\lambda ^{}\pi _k}{}L_\lambda ^{}L_{\lambda _0^k\lambda ^{}}\right).$$
(13)
To obtain the necessary background for the integrable models with deformed Yangian symmetry the simplest way is to construct the defining representation of the universal $`_𝒴`$-matrix, $`R=d\left(_𝒴\right)`$. In some special situations the quantum $`R_{}`$-matrix corresponding to the twisted algebra can be obtained directly from its classical counterpart $`r_{}`$ (in particular this happens when $`R=\mathrm{exp}\left(r\right)`$ ). Such simplifications are unavailable when the $`r_{}`$-matrices of the type (13) are considered for the orthogonal classical Lie algebras $`g`$. In such cases more information about the $`_{}`$-matrix is necessary.
In this report it will be shown how the chains of extended twists (6) lead to the series of deformed Yangians. To illustrate the situation the explicit formulas will be presented for the case of $`g=so(2N+1)`$.
## 2 Chains and twisted $``$-matrices
Let $`g`$ be a classical Lie algebra of the type $`B_N`$ or $`D_N`$ and $`_{_{p0}}`$ $``$ $`__p_{_{p1}}\mathrm{}`$ $`\mathrm{}__0`$ – its full proper chain of extended twists.
To construct a maximal proper chain for $`𝒜=U(g)`$ the following sequences $`𝒜𝒜_0𝒜_1\mathrm{}𝒜_{p1}𝒜_p`$ of Hopf subalgebras are to be fixed:
$$\begin{array}{c}U(so(2N))U(so(2(N2))\mathrm{}U(so(2(N2k))\mathrm{}\hfill \end{array}\mathrm{for}D_N$$
(14)
$$U(so(2N+1))U(so(2(N2)+1)\mathrm{}U(so(2(N2k)+1)\mathrm{}\mathrm{for}B_N$$
(15)
In both cases the initial roots for $`𝒜_k`$ can be chosen to be $`\lambda _0^k=e_1+e_2`$ (here the root subsystems are considered for $`𝒜_k`$ separately and all the roots are written in the standard $`e`$-basis).
The twisting element for a chain can be written explicitly as
$$\begin{array}{ccc}_{_{p0}}\hfill & =\hfill & _{\lambda ^{}\pi _p^{}}\left(\mathrm{exp}\{L_\lambda ^{}L_{\lambda _0^p\lambda ^{}}e^{\frac{1}{2}\sigma _0^p}\}\right)\mathrm{exp}\{H_{\lambda _0^p}\sigma _0^p\}\hfill \\ & & _{\lambda ^{}\pi _{p1}^{}}\left(\mathrm{exp}\{L_\lambda ^{}L_{\lambda _0^{p1}\lambda ^{}}e^{\frac{1}{2}\sigma _0^{p1}}\}\right)\mathrm{exp}\{H_{\lambda _0^{p1}}\sigma _0^{p1}\}\hfill \\ & & \mathrm{}\hfill \\ & & _{\lambda ^{}\pi _0^{}}\left(\mathrm{exp}\{L_\lambda ^{}L_{\lambda _0^0\lambda ^{}}e^{\frac{1}{2}\sigma _0^0}\}\right)\mathrm{exp}\{H_{\lambda _0^0}\sigma _0^0\}\hfill \end{array}$$
(16)
Let us introduce the deformation parameters $`\xi _k=\xi \eta _k`$ and rescale the generators $`\{L_{\lambda _0^k},L_{\lambda _0^k\lambda ^{}}\}`$ by $`\xi _k`$ in each subalgebra $`𝒜_k`$. Using the expressions (10) and (11) one can get the first terms of the expansion for the twisting element
$$_{_{p0}}\left(\xi \right)=II+\xi \rho _{}+𝒪\left(\xi ^2\right).$$
Here
$$\rho _{}=\underset{k=0,1,\mathrm{},p}{}\eta _k\left(H_{\lambda _0^k}L_{\lambda _0^k}+\underset{\lambda ^{}\pi _k}{}L_\lambda ^{}L_{\lambda _0^k\lambda ^{}}\right)$$
(17)
is ”a half” of the $`r_{_{p0}}`$-matrix:
$$r_{_{p0}}=\rho _{}\tau \rho _{}.$$
The carrier subalgebras for chains of extended twists are solvable. As it was already mentioned in Section 1 the generators of the carrier form dual sets. In the case of extended twists one of these sets $`=\left\{L_{\lambda _0^k},L_{\lambda _0^k\lambda ^{}}\right|\lambda ^{}\pi _k^{}\}`$ forms a nilpotent subalgebra. In the defining representation $`d(g)`$ the corresponding matrix ring is nilpotent with the index $`\kappa =3`$. As it was demonstrated in the twisting element can always be presented in the form $`=\mathrm{exp}\left(_i^{}\psi ^i\left(\right)\right)`$and the expansions for the elements $`\psi ^i\left(\right)`$ starts with the term $`\xi ^i`$. Applying these results to the case of chains we get the following property.
Lemma. Let $`g`$ be a Lie algebra of the type $`A_N`$, $`B_N`$ or $`D_N`$ and $`d(g)`$ – the defining representation of $`g`$. Then for the full proper chain of extended twists $`_{_{p0}}\left(\xi \right)=II+\xi \rho _{}+𝒪\left(\xi ^2\right)`$ the following relations are true:
1. $`d\left(\rho _{}^3\right)=0,`$
2. $`d\left(__k\left(\xi \right)\right)=\mathrm{exp}\left(\xi d\left(\rho __k\right)\right),`$
3. $`R_{_{p0}}\left(\xi \right)=d\left(_{_{p0}}\left(\xi \right)\right)=II\xi d\left(\rho _{}\right)+\frac{\xi ^2}{2}\left(d^2\left(\rho _{}\right)+d^2\left(\tau \rho _{}\right)\right)\xi ^2d\left(\tau \rho _{}\right)d\left(\rho _{}\right).\mathrm{}`$
The construction of the $`R`$-matrix differs for linear and orthogonal algebras. The most interesting is the case $`g=so(M)`$, i.e. the series $`B_N`$ or $`D_N`$. Let us fix $`g`$ to be of $`B_N`$ type, $`g=so(2N+1)`$. Thus $`𝖢^{\mathrm{𝟤}𝖭+\mathrm{𝟣}}`$ is the space of the defining representation. Let $`P`$ be the permutation matrix acting in $`𝖢^{\mathrm{𝟤}𝖭+\mathrm{𝟣}}𝖢^{\mathrm{𝟤}𝖭+\mathrm{𝟣}}`$
$$P=M^{}M^{\prime \prime }\mathrm{Mat}(2N+1,𝖢)^2$$
Define also the matrix
$$K=\left(M^{}\right)^𝖳M^{\prime \prime }.$$
It was shown in that if $``$ is the twisting element for $`U(so(2N+1))`$ with the twisted $`R`$-matrix $`R_{}`$ the corresponding rational solution of the quantum Yang-Baxter equation in the defining representation has the form
$$R=d\left(_𝒴\right)=uR_{}+P\frac{u}{u+N1/2}d\left(_{21}\right)Kd\left(^1\right).$$
(18)
Applying the Lemma proved above we get the following final expression for $`R`$.
$$\begin{array}{c}R=u\left(II\xi d\left(\rho _{}\right)+\frac{\xi ^2}{2}d\left(\rho _{}^2+\left(\tau \rho _{}\right)^22\left(\tau \rho _{}\right)\rho _{}\right)\right)+P\hfill \\ \frac{u}{u+N1/2}d\left(II+\tau \rho _{}+\frac{1}{2}\left(\tau \rho _{}\right)^2\right)Kd\left(II\rho _{}\frac{1}{2}\left(\rho _{}\right)^2\right).\hfill \end{array}$$
(19)
This $`R`$-matrix describes the Yangians $`Y_{_{_{p0}}}(so(2N+1))`$ deformed by the full chains of extended twists. To get the final answer the only term that must be calculated is the defining representation for the $`\rho _{}`$-matrix (17).
## 3 Example. QYBE solutions for so(2N+1)
Let $`g=so(2N+1)`$. To make the illustration maximally visual we shall consider the simplest nontrivial chain, i. e. put $`p=0`$ and consider a chain with a single factor $`__0`$ in (6). The algorithm for other factors $`__k`$ is similar to that of the first one and the case of a full chain can be reconstructed using the formulas (9-11) and the expressions presented below.
The matrices of the defining representation $`d(g)`$ are written in terms of basic antisymmetric Okubo matrices $`M_{ij}`$:
$$d\left(H_{e_i+e_j}\right)=i\left(M_{2i1,2i}+M_{2j1,2j}\right)H_{i+j},$$
$$d\left(L_{e_k}\right)=M_{2k,2N+1}iM_{2k1,2N+1}E_k,k=1,\mathrm{},N$$
$$\begin{array}{ccc}d\left(L_{e_i\pm e_j}\right)\hfill & =& \frac{1}{2}\left(M_{2i,2j}\pm iM_{2i,2j1}+iM_{2i1,2j}\pm M_{2i1,2j1}\right)\hfill \\ & & E_{i+j},i<j\hfill \end{array}$$
In this representation the factors of the sequence
$$d\left(__0\left(\xi \right)\right)=d\left(\mathrm{\Phi }__k\left(\xi \right)\right)d\left(\mathrm{\Phi }_{𝒥_0}\left(\xi \right)\right)$$
have the following form:
$$d\left(\mathrm{\Phi }_{𝒥_0}\left(\xi \right)\right)=\mathrm{exp}\{\xi H_{1+2}E_{1+2}\},$$
$$d\left(\mathrm{\Phi }__k\left(\xi \right)\right)=\mathrm{exp}\left\{\xi \left(E_1E_2+2\underset{j>2}{\overset{N}{}}E_{1\pm j}E_{2j}\right)\right\}.$$
The $`R`$ -matrix can be easily obtained using the general formula (12),
$$\begin{array}{c}R__0=II\xi \left(H_{1+2}E_{1+2}+E_1E_2+2_{j>2}^NE_{1\pm j}E_{2j}\right)+\hfill \\ +\frac{1}{2}\xi ^2\left(E_1^2E_2^2+E_2^2E_1^2+2E_{1+2}E_{1+2}2E_2E_1E_1E_2\right)+\hfill \\ +2\xi ^2_{j>2}^N(E_{1\pm j}E_{1j}E_{2j}E_{2\pm j}+E_{2j}E_{2\pm j}E_{1\pm j}E_{1j}\hfill \\ 2E_{2j}E_{1\pm j}E_{1\pm j}E_{2j})\hfill \end{array}$$
The corresponding matrix $`d(\rho )`$ is
$$d\left(\rho __0\right)=H_{1+2}E_{1+2}+E_1E_2+2\underset{j>2}{\overset{N}{}}E_{1\pm j}E_{2j}.$$
This expression together with the general formula (19) defines the final result – the set of $`R_𝒴`$-matrices for the deformed Yangians $`Y__0(so(2N+1))`$.
## 4 Conclusions
We have demonstrated that for some types of simple Lie algebras the Yangians deformed by chains of extended twists are completely determined by the matrix $`d(\rho )`$ which is the logarithm of the twisting element in the defining representation of an algebra. This gives rise to a set $`R_𝒴`$ of solutions to the quantum Yang-Baxter equation and, therefore, to a series of integrable models with the corresponding local Hamiltonians.
When the chain has the index $`p>1`$ the $`R_𝒴`$-matrices naturally acquire the parameters $`\eta _k`$. Moreover, in chains of extended twists one can cut any number of factors from the left, i.e. use the discrete parametrization of the last nontrivial sequence $`__p`$.
## 5 Acknowledgements
I am grateful to professors P.P.Kulish and A.Stolin for valuable discussions on the subject. It is a pleasure for me to thank the organizers of the International Seminar ”Supersymmetries and Quantum Symmetries” (SQS’99, 27-31 July, 1999) for warm hospitality.
This work has been partially supported by the Russian Foundation for Basic Research under the grant N 97-01-01152. |
warning/0002/math0002093.html | ar5iv | text | # 0 Introduction
## 0 Introduction
An $`n`$-dimensional submanifold $`X`$ of a projective space $`P^N(𝐂)`$ is called tangentially degenerate if the rank of its Gauss mapping $`\gamma :XG(n,N)`$ is less than $`n,r=\text{rank}\gamma <n`$. Here $`xX,\gamma (x)=T_x(X)`$, and $`T_x(X)`$ is the tangent subspace to $`X`$ at $`x`$ considered as an $`n`$-dimensional projective space $`P^n`$. The number $`r`$ is also called the rank of $`X,r=\text{rank}X`$. The case $`r=0`$ is trivial one: it gives just an $`n`$-plane. A submanifold $`X`$ is called tangentially degenerate if $`0r<n`$, and it is denoted by $`V_r^n,X=V_r^n`$. The submanifolds of rank $`r<n`$ have been the object of numerous investigations because of their analogy to developable surfaces in a three-dimensional space and because of their significance in the theory of the curvature of submanifolds.
The tangentially degenerate submanifolds $`X`$ of rank $`r<n`$ were first considered by É. Cartan \[C 16\] in connection with his study of metric deformations of hypersurfaces, and in \[C 19\] in connection with his study of manifolds of constant curvature. In particular, Cartan proved that if $`V_r^n`$ is a tangentially degenerate submanifold of dimension $`n`$ and rank $`r`$ and the dimension $`\rho `$ of the osculating subspace is $`\rho =n+\frac{1}{2}r(r+1)`$, then $`V_r^n`$ is a cone with an $`(nr1)`$-dimensional vertex.
It appeared that tangentially degenerate submanifolds are less rigid under the affine and projective deformations. Yanenko \[Ya 53\] studied them in connection with his study of metric deformations of submanifolds. Akivis \[A 57, 62\] studied them in multidimensional projective space, considered their focal images (the locus of singular points and the locus of singular hyperplanes) and applied the latter to clarify the structure of the tangentially degenerate submanifolds. Savelyev \[Sa 57, 60\] found a classification of tangentially degenerate submanifolds and described in detail the tangentially degenerate submanifolds of rank 2. Ryzhkov \[R 58\] showed that a tangentially degenerate submanifolds $`X`$ of rank $`r`$ can be constructed by using the Peterson transformation of $`r`$-dimensional submanifolds, and in \[R 60\] he proved that such a construction is quite general. That is, by means of it, an arbitrary tangentially degenerate submanifold $`X`$ of rank $`r`$ can be obtained (see also the survey paper \[AR 64\]). In particular, Ryzhkov \[R 60\] generalized the above mentioned Cartan’s result to the case $`n+1+\frac{1}{2}r(r1)<\rho <n+\frac{1}{2}r(r+1)`$. Brauner \[Br 38\], Wu \[Wu 95\], and Fischer and Wu \[FW 95\] studied such submanifolds in an Euclidean $`N`$-space.
For a submanifold $`V^n`$ of a Riemannian space $`V^N,n<N`$, Chern and Kuiper \[CK 52\] introduced the notion of the index of relative nullity $`\mu (x)`$, where $`xV^n`$ (see also \[KN 69\], p. 348). The submanifolds $`V^n`$, for which $`\mu (x)`$ is constant and greater than 0 for all points $`xV^n`$, are called strongly parabolic. Akivis \[A 87\] proved that if a space $`V^N`$ admits a projective realization (this is always the case for the simply connected Riemannian spaces of constant curvature, see \[W 72\], Ch. 2) and if the index $`\mu (x)`$ is constant on $`V^nV^N`$, then the index $`\mu (x)`$ is connected with the rank of a submanifold $`V^n`$ by the relation $`\mu (x)=nr`$. This implies that the results of the papers \[A 57, 62\] as well as the results of the current paper can be applied to the study of strongly parabolic submanifolds of the Euclidean and non-Euclidean spaces. In particular, in \[A 87\] Akivis proved the existence of tangentially degenerate submanifolds in such spaces without singularities and constructed examples of such submanifolds. The main results of papers indicated above can be found in Chapter 4 of the book \[AG 93\]. In the same paper \[A 87\], Akivis also proved that a hypersurface in a four-dimensional Euclidean space $`𝐑^4`$ considered by Sacksteder \[S 60\] is of rank 2 and without singularities, and that this hypersurface is a particular case of a series of examples presented in \[A 87\]. Later Akivis and Goldberg \[AG 00\] proved that a similar example constructed (but not published) by Bourgain and published in \[Wu 95\], \[I 98, 99a, 99b\], and \[WZ 99\] coincides with Sacksteder’s example up to a coordinate transformation.
Griffiths and Harris \[GH 79\] (Section 2, pp. 383–393) considered tangentially degenerate submanifolds from the point of view algebraic geometry. They used the term “submanifolds with degenerate Gauss mappings” instead of the term “tangentially degenerate varieties”. They used this term to avoid a confusion with submanifolds with degenerate tangential varieties considered in Section 5 of the same paper.
Following \[GH 79\], Landsberg \[L 96\] considered tangentially degenerate submanifolds. His recently published notes \[L 99\] are in some sense an update to the paper \[GH 79\]. Section 5 (pp. 47–50) of these notes is devoted to tangentially degenerate submanifolds.
Griffiths and Harris \[GH 79\] asserted a structure theorem for submanifolds with degenerate Gauss mappings, that is, for the varieties $`X=V_r^n`$ such that $`dim\gamma (X)<dimX`$. They asserted that such varieties are “built up from cones and developable varieties” (see \[GH 79\], p. 392). They gave a proof of this assertion in the case $`n=2`$. However, their assertion is not completely correct. In a recent note \[AGL\], Akivis, Goldberg, and Landsberg present counter-examples to Griffiths–Harris’ assertion when $`n>2`$, and in particular, they prove that this assertion is false even for hypersurfaces with one-dimensional fibers.
Recently four papers \[I 98, 99a, 99b\] and \[IM 97\] on tangentially degenerate submanifolds (called “developable” in these papers) were published. In \[IM 97\], the authors found the connection between such submanifolds and solutions of Monge-Ampère equations, with the foliation of plane generators $`L`$ of $`X`$ called the Monge-Ampère foliation. In \[IM 97\], the authors proved that the rank of the Gauss map of a compact tangentially degenerate $`C^{\mathrm{}}`$-hypersurface $`X𝐑P^N`$ is an even integer satisfying the inequality $`\frac{r(r+3)}{2}<N,r0`$, and that if $`r1`$, then $`X`$ is necessarily a projective hyperplane of $`𝐑P^N`$. If $`N=3`$ or $`N=5`$, then a compact tangentially degenerate $`C^{\mathrm{}}`$-hypersurface is a projective hyperplane.
In \[I 98, 99b\], Ishikawa found real algebraic cubic nonsingular tangentially degenerate hypersurface in $`𝐑P^N`$ for $`N=4,7,13,25`$, and in \[I 99a\] he studied singularities of tangentially degenerate $`C^{\mathrm{}}`$-hypersurfaces.
In 1997 Borisenko published the survey paper \[B 97\] in which he discussed results on strongly parabolic submanifolds and related questions in Riemannian and pseudo-Riemannian spaces of constant curvature and in particular, in an Euclidean space $`E^N`$. Among other results, he gives a description of certain classes of submanifolds of arbitrary codimension that are analogous to the class of parabolic surfaces in an Euclidean space $`E^3`$. Borisenko also investigates the local and global metric and topological properties; indicates conditions which imply that a submanifold of an Euclidean space $`E^N`$ is cylindrical; presents results on strongly parabolic submanifolds in pseudo-Riemannian spaces of constant curvature, and finds the relationship with minimal surfaces.
In the current paper we study systematically the differential geometry of tangentially degenerate submanifolds of a projective space $`P^N(𝐂)`$. By means of the focal images, three basic types of submanifolds are discovered: cones, tangentially degenerate hypersurfaces, and torsal submanifolds. Moreover, for tangentially degenerate submanifolds, a structural theorem is proven. By this theorem, tangentially degenerate submanifolds that do not belong to the basic types are foliated into submanifolds of basic types. In the proof we introduce irreducible, reducible, and completely reducible tangentially degenerate submanifolds. It is found that cones and tangentially degenerate hypersurfaces are irreducible and torsal submanifolds are completely reducible while all other tangentially degenerate submanifolds not belonging to basic types are reducible. Particular examples of tangentially degenerate submanifolds as well as tangentially degenerate submanifolds of low dimensions is considered in \[AGL\]. Some examples of tangentially degenerate submanifolds can be also found in \[A 87\].
In this paper we apply the method of exterior forms and moving frames of Cartan \[C 45\] which was often successfully used in differential geometry.
## 1 The main results
Let $`XP^N(𝐂)`$ be an $`n`$-dimensional smooth submanifold with a degenerate Gauss map $`\gamma :XG(n,N),\gamma (x)=T_x(X),xX`$. Suppose that $`\text{rank}\gamma =r<n`$. Denote by $`L`$ a leave of this map, $`L=\gamma ^1(T_x)X`$.
###### Theorem 1
A leave $`L`$ of the Gauss map $`\gamma `$ is a subspace of $`P^N,dimL=nr=l`$ or its open part.
For the proof of this theorem see \[AG 93\] (p. 115, Theorem 4.1).
The foliation on $`X`$ with leaves $`L`$ is called the Monge-Ampère foliation (see, for example, \[D 89\] or \[I 98, 99b\]). In this paper, we extend the leaves of the Monge-Ampère foliation to a projective space $`P^l`$ assuming that $`LP^l`$ is a plane generator of the submanifold $`X`$. As a result, we have $`X=f(P^l\times M^r)`$, where $`M^r`$ is a parametric variety, and $`f`$ is a differentiable map $`f:P^l\times M^rP^N`$.
However, unlike a traditional definition of the foliation (see for example, \[DFN 85\], §29), the leaves of the Monge-Ampère foliation have singularities. This is a reason that in general its leaves are not diffeomorphic to a standard leaf.
The tangent subspace $`T_x(X)`$ is fixed when a point $`x`$ moves along $`L`$. This is the reason that we denote it by $`T_L,LT_L`$. A pair $`(L,T_L)`$ on $`X`$ depends on $`r`$ parameters.
With a second-order neighborhood of a pair $`(L,T_L)`$, two systems of square $`r\times r`$ matrices $`B^\alpha =(b_{pq}^\alpha )`$ and $`C_i=(c_{pi}^q),i=0,1,\mathrm{},l;p,q=l+1,\mathrm{},n;\alpha =n+1,\mathrm{},N,`$ are associated. The equation $`det(C_ix^i)=0`$ defines the set of singular points $`x=(x^i)L`$ of the map $`\gamma `$. This set is an algebraic hypersurface $`F_LL`$ of degree $`r`$ which is called the focus hypersurface. The equation $`det(\xi _\alpha B^\alpha )=0`$ defines the set of singular tangent hyperplanes $`\xi =(\xi _\alpha )T_L`$ of the map $`\gamma `$. This set is an algebraic hypercone $`\mathrm{\Phi }_L`$ with vertex $`T_L`$ which is called the focus hypercone.
It appears that the products $`H_i^\alpha =B^\alpha C_i=(h_{ipq}^\alpha )`$ are symmetric. They define on $`X`$ the second fundamental forms $`h_i^\alpha =h_{ipq}^\alpha \theta ^p\theta ^q`$, where $`\theta ^p`$ are basis forms of the manifold $`M`$. We assume that for all values of parameters $`u=(u^p)M`$, the system of forms $`h_i^\alpha `$ is regular, i.e., among them there is at least one nondegenerate quadratic form.
Denote by $`S_L`$ the osculating subspace of $`X`$ which is constant at all points $`xL`$ of its generator $`L`$. Its dimension is $`dimS_L=n+m`$, where $`m`$ is the number of linearly independent second fundamental forms of $`X`$.
We assume that the conditions of all theorems in this paper are satisfied for all values of parameters $`uM`$.
###### Theorem 2
Suppose that $`l1`$ and $`m2`$, and the focus hypersurfaces $`F_L`$ and the focus hypercones $`\mathrm{\Phi }_L`$ do not have multiple components. Then the submanifolds $`X`$ is foliated into $`r`$ families of torses with $`l`$-dimensional plane generators $`L`$. Each of these families depends on $`r1`$ parameters.
For the definition of torse see Example 2 in Section 2. A manifold $`X`$ described in Theorem 2 is called torsal.
###### Theorem 3
Suppose that $`l2`$, and the focus hypersurfaces $`F_L`$ do not have multiple components and are indecomposable. Then the submanifold $`X`$ is a hypersurface of rank $`r`$ in a subspace $`P^{n+1}P^N`$.
###### Theorem 4
Suppose that $`m2`$, and the focus hypercones $`\mathrm{\Phi }_L`$ do not have multiple components and are indecomposable. Then the submanifold $`X`$ is a cone with an $`(l1)`$-dimensional vertex and $`l`$-dimensional plane generators.
The system of matrices $`B^\alpha `$ and $`C_i`$ associated with a submanifold $`X`$ is said to be reducible if these matrices can be simultaneously reduced to a block diagonal form:
$$C_i=\text{diag}(C_{i1},\mathrm{},C_{is}),B^\alpha =\text{diag}(B_1^\alpha ,\mathrm{},B_s^\alpha ),$$
(1)
where $`C_{it}`$ and $`B_t^\alpha ,t=1,\mathrm{},s,`$ are square matrices of orders $`r_t`$, and $`r_1+r_2+\mathrm{}+r_s=r`$. If such a decomposition of matrices is not possible, the system of matrices $`B^\alpha `$ and $`C_i`$ is called irreducible. If $`r_1=r_2=\mathrm{}=r_s=1`$, then the system of matrices $`B^\alpha `$ and $`C_i`$ is called completely reducible.
A manifold $`X`$ with a degenerate Gauss mapping is said to be reducible, irreducible or completely reducible if for any values of parameters $`uM`$ the matrices $`B^\alpha `$ and $`C_i`$ are reducible, irreducible or completely reducible, respectively.
###### Theorem 5
Suppose that a manifold $`X`$ is reducible, and its matrices $`B_t^\alpha `$ and $`C_{it}`$ defined in $`(1)`$ are of order $`r_t,t=1,\mathrm{},s`$. Then $`X`$ is foliated into $`s`$ families of $`(l+r_t)`$-dimensional submanifolds of rank $`r_t`$ with $`l`$-dimensional plane generators. For $`r_t=1`$, these submanifolds are torses, and for $`r_t2`$, they are irreducible submanifolds described in Theorems $`3`$ and $`4`$.
## 2 Examples of submanifolds with degenerate <br>Gauss maps
Consider a few examples of submanifolds with a degenerate Gauss map.
###### Example 1
For $`r=0`$, a submanifold $`X`$ is an $`n`$-dimensional subspace $`P^n,n<N`$. This submanifold is the only tangentially degenerate submanifold without singularities in $`P^N`$.
###### Example 2
Let $`Y`$ be a curve of class $`C^p`$ in the space $`P^N`$. Suppose that $`P^N`$ is a space of minimal dimension containing the curve $`Y`$. Denote by $`L_y`$ the osculating subspace of order $`l,lp,lN1,`$ of the curve $`Y`$ at a point $`yY`$. Since $`dimL_y=l`$, it follows that when a point $`y`$ moves along the curve $`Y`$, the subspace $`L_y`$ sweeps a submanifold $`X`$ of dimension $`n=l+1`$ in the space $`P^N`$. At any point $`xL_y`$, the tangent subspace $`T_x(X)`$ coincides with the osculating subspace $`L_y^{}`$ of order $`l+1`$ of the curve $`Y`$, $`dimL_y^{}=l+1`$, and the focus hypersurface in $`L_y`$ is the osculating subspace $`{}_{}{}^{}L_{y}^{}`$ of order $`l1`$ and dimension $`l1`$ of the curve $`Y`$. Thus $`dimX=l+1`$, and the manifold $`X`$ is tangentially degenerate of rank $`r=1`$. Such a manifold $`X`$ is called a torse. Conversely, a submanifold of dimension $`n`$ and rank 1 is a torse formed by a family of osculating subspaces of order $`n1`$ of a curve of class $`C^p,pn1`$, in the space $`P^N`$.
In what follows, unless otherwise stated, we always assume that $`r>1`$.
###### Example 3
Suppose that $`S`$ is a subspace of the space $`P^N,dimS=l1`$, and $`T`$ is its complementary subspace, $`dimT=Nl,TS=\mathrm{}`$. Let $`Y`$ be a smooth tangentially nondegenerate submanifold of the subspace $`T`$, $`dimY=\text{rank}Y=r<Nl`$. Consider an $`r`$-parameter family of $`l`$-dimensional subspaces $`L_y=Sy,yY`$. This manifold is a cone $`X`$ with vertex $`S`$ and the director manifold $`Y`$. The subspace $`T_x(X)`$ tangent to the cone $`X`$ at a point $`xL_y`$ is defined by its vertex $`S`$ and the subspace $`T_y(Y),T_x(X)=ST_y(X)`$, and $`T_x(X)`$ remains fixed when a point $`x`$ moves in the subspace $`L_y`$. As a result, the cone $`X`$ is a tangentially degenerate submanifold of dimension $`n=l+r`$ and rank $`r`$, with plane generators $`L_y`$ of dimension $`l`$.
###### Example 4
Let $`T`$ be a subspace of dimension $`n+1`$ in $`P^N`$, and let $`Y`$ be an $`r`$-parameter family of hyperplanes $`\xi `$ in general position in $`T,r<n`$. Such a family has an $`n`$-dimensional envelope $`X`$ that is a tangentially degenerate submanifold of dimension $`n`$ and rank $`r`$ in the subspace $`T`$. It is foliated into an $`r`$-parameter family of plane generators $`L`$ of dimension $`l=nr`$ along which the tangent subspace $`T_x(X),xL,`$ is fixed and coincides with a hyperplane $`\xi `$ of the family in question. Thus $`X`$ is a tangentially degenerate hypersurface of rank $`r`$ with $`(nr)`$-dimensional flat generators $`L`$ in the space $`T`$.
###### Example 5
If $`\text{rank}X=dimX=n`$, then $`X`$ is a submanifold of complete rank. $`X`$ is also called a tangentially nondegenerate in the space $`P^N`$.
## 3 Application of the duality principle
By the duality principle, to a point $`x`$ of a projective space $`P^N`$, there corresponds a hyperplane $`\xi `$. A set of hyperplanes of space $`P^N`$ forms the dual projective space $`(P^N)^{}`$ of the same dimension $`N`$. Under this correspondence, to a subspace $`P`$ of dimension $`p`$, there corresponds a subspace $`P^{}(P^N)^{}`$ of dimension $`Np1`$. Under the dual map, the incidence of subspaces is preserved, that is, if $`P_1P_2`$, then $`P_1^{}P_2^{}`$.
In the space $`P^N`$, if a point $`x`$ describes a tangentially nondegenerate manifold $`X,dimX=r`$, then in general, the hyperplane $`\xi `$ corresponding to $`x`$ envelopes a hypersurface $`X^{}`$ of rank $`r`$ with $`(Nr1)`$-dimensional generators. The dual map sends a smooth manifold $`XP^N`$ of dimension $`n`$ and rank $`r`$ with plane generators $`L`$ of dimension $`l=nr`$ to a manifold $`X^{}`$ of dimension $`n^{}=Nl1`$ and the same rank $`r`$ with plane generators $`L^{}`$ of dimension $`l^{}=Nn1`$. Under this map, to a tangent subspace $`T_x(X)`$ of the submanifold $`X`$ there corresponds the plane generator $`L^{}`$, and to a plane generator $`L`$ there corresponds the tangent subspace of the manifold $`X^{}`$.
Note that Ein \[E 85, 86\] applied the duality principle for studying tangentially degenerate varieties with small dual varieties.
Let us determine which manifolds correspond to tangentially degenerate manifolds considered in Examples 2–4. To a cone $`X`$ of rank $`r`$ with vertex $`S`$ of dimension $`l1`$ (see Example 3), there corresponds a manifold $`X^{}`$ lying in the subspace $`T=S^{},dimT=Nl`$. Since $`dimX^{}=n^{}=Nl1`$, the manifold $`X^{}`$ is a hypersurface of rank $`r`$ in the subspace $`T`$. Such a hypersurface was considered in Example 4. It follows that Examples 3 and 4 are mutually dual one to another.
Suppose that the subspace $`T`$ containing a hypersurface $`X`$ of rank $`r`$ coincides with the space $`P^N`$. Then under the dual map in $`T`$, to the hypersurface $`X`$ there corresponds an $`r`$-dimensional tangentially nondegenerate manifold $`X^{}`$.
Under the dual map, in $`P^N`$ to torses (see Example 2) there correspond torses enveloping a one-parameter family of hyperplanes $`\eta `$—the images of points of the curve $`Y`$. These torses are of dimension $`N1`$ and rank 1, and their plane generators are of dimension $`l=N2`$.
Note that the dual map allows us to construct tangentially degenerate manifolds in the space $`P^{n+1}`$ from the general tangentially nondegenerate manifolds.
###### Example 6
Consider a set of conics in a projective plane $`P^2`$. They are defined by the equation
$$a_{ij}x^ix^j=0,i,j=0,1,2,$$
(2)
where the coefficients $`a_{ij}`$ are defined up to a constant factor, and $`a_{ij}=a_{ji}`$. Thus to the curve (2) there corresponds a point of a five-dimensional projective space $`P^5`$ whose coordinates coincide with the coefficients of equation (2).
Among the conics defined by equation (2) there are singular curves which decompose into two straight lines. Such conics are distinguished by the condition
$$det(a_{ij})=0.$$
(3)
Equation (3) is of degree three and defines a hypersurface in $`P^5`$ which is called the cubic symmetroid.
We prove now that the hypersurface defined by equation $`(3)`$ is tangentially degenerate of rank two and bears a two-parameter family of two-dimensional plane generators. To this end, we write equation (3) in the form
$$F=det\left(\begin{array}{ccc}a_{00}& a_{01}& a_{02}\\ a_{10}& a_{11}& a_{12}\\ a_{20}& a_{21}& a_{22}\end{array}\right)=0.$$
(4)
The equation of tangent subspaces of the hypersurface (4) has the form
$$\xi ^{ij}a_{ij}=0,$$
where $`\xi ^{ij}={\displaystyle \frac{F}{a_{ij}}},\xi ^{ij}=\xi ^{ji},`$ are the cofactors of the entries $`a_{ij}`$ in the determinant in equation (4). The quantities $`\xi ^{ij}`$ are coordinates in the projective space $`P^5`$, which is dual to the space $`P^5`$ where we defined the cubic symmetroid (4).
Consider the symmetric matrix $`\xi =(\xi ^{ij})`$. A straightforward computation shows that the rank of this matrix is equal to one. Thus the entries of the matrix $`\xi `$ can be represented in the form
$$\xi ^{ij}=\xi ^i\xi ^j,$$
(5)
where $`\xi ^i`$ are projective coordinates of a point in the projective plane $`P^2`$. But equations (5) define the Veronese surface in the space $`P^5`$. Since the Veronese surface is of dimension two, its points depend on two affine parameters $`u=\frac{\xi ^1}{\xi ^0}`$ and $`v=\frac{\xi ^2}{\xi ^0}`$. As a result, the tangent hyperplanes of the cubic symmetroid (4) depend on two parameters, and this hypersurface is of rank two.
To the plane generators of the hypersurface (4), there correspond two-parameter families of conics in $`P^2`$ decomposing into pairs of straight lines with a common intersection point.
Thus, the cubic symmetroid $`(4)`$ and the Veronese surface $`(5)`$ are mutually dual submanifolds of a five-dimensional projective space.
## 4 Basic equations and focal images
We study tangentially degenerate submanifolds applying the method of moving frames in a projective space $`P^N`$. In $`P^N`$, we consider a manifold of projective frames $`\{A_0,A_1,\mathrm{},A_N\}`$. On this manifold
$$dA_u=\omega _u^vA_v,u,v,=0,1,\mathrm{},N,$$
(6)
where the sum $`\omega _u^u=0`$. The 1-forms $`\omega _u^v`$ are linearly expressed in terms of the differentials of parameters of the group of projective transformations of the space $`P^N`$. The total number of these parameters is $`N^21`$. These 1-forms satisfy the structure equations
$$d\omega _u^v=\omega _u^w\omega _w^v$$
(7)
of the space $`P^N`$ (see, for example, \[AG 93\], p. 19). Equations (7) are the conditions of complete integrability of equations (6).
Consider a tangentially degenerate submanifold $`XP^N,dimX=n,\text{rank}X=rn`$. In addition, as above, let $`L`$ be a rectilinear generator of the manifold $`X,dimL=l`$; let $`T_L,dimT_L=n,`$ be the tangent subspace to $`X`$ along the generator $`L`$, and let $`M`$ be a base manifold for $`X,dimM=r`$. Denote by $`\theta ^p,p=l+1,\mathrm{},n`$, basis forms on the variety $`M`$. These forms satisfy the structure equations
$$d\theta ^p=\theta ^q\theta _q^p,p,q=l+1,\mathrm{},n,$$
(8)
of the variety $`M`$. Here $`\theta _q^p`$ are 1-forms defining transformations of first-order frames on $`M`$.
For a point $`xL`$, we have $`dxT_L`$. With $`X`$, we associate a bundle of projective frames $`\{A_i,A_p,A_\alpha \}`$ such that $`A_iL,i=0,1,\mathrm{},l;A_pT_L,p=l+1,\mathrm{},n`$. Then
$$\{\begin{array}{cc}dA_i=\omega _i^jA_j+\omega _i^pA_p,\hfill & \\ dA_p=\omega _p^iA_i+\omega _p^qA_q+\omega _p^\alpha A_\alpha ,\alpha =n+1,\mathrm{},N.\hfill & \end{array}$$
(9)
It follows from the first equation of (9) that
$$\omega _i^\alpha =0.$$
(10)
Since for $`\theta ^p=0`$ the subspaces $`L`$ and $`T_L`$ are fixed, we have
$$\omega _i^p=c_{qi}^p\theta ^q,\omega _p^\alpha =b_{qp}^\alpha \theta ^q.$$
(11)
Since the manifold of leaves $`LX`$ depends on $`r`$ essential parameters, the rank of the system of 1-forms $`\omega _i^p`$ is equal to $`r`$, $`\text{rank}(\omega _i^p)=r`$. Similarly, we have $`\text{rank}(\omega _p^\alpha )=r`$.
Denote by $`C_i`$ and $`B^\alpha `$ the $`r\times r`$ matrices occurring in equations (11):
$$C_i=(c_{qi}^p),B^\alpha =(b_{qp}^\alpha ).$$
These matrices are defined in a second-order neighborhood of the submanifold $`X`$.
Exterior differentiation of equations (10) by means of structure equations (7) leads to the exterior quadratic equations
$$\omega _p^\alpha \omega _i^p=0.$$
Substituting expansions (11) into the last equations, we find that
$$b_{qs}^\alpha c_{pi}^s=b_{ps}^\alpha c_{qi}^s.$$
(12)
Equations (10), (11), and (12) are called the basic equations in the theory of tangentially degenerate submanifolds.
Relations (12) can be written in the matrix form
$$(B^\alpha C_i)^T=(B^\alpha C_i),$$
i.e., the matrices
$$H_i^\alpha =B^\alpha C_i=(b_{qs}^\alpha c_{pi}^s)$$
are symmetric.
Let $`x=x^iA_i`$ be an arbitrary point of a leaf $`L`$. For such a point we have
$$dx=(dx^i+x^j\omega _j^i)A_i+x^i\omega _i^pA_p.$$
It follows that
$$dx(A_pc_{qi}^px^i)\theta ^q(modL).$$
The tangent subspace $`T_x`$ to the manifold $`X`$ at a point $`x`$ is defined by the points $`A_i`$ and
$$\stackrel{~}{A}_q(x)=A_pc_{qi}^px^i,$$
and therefore $`T_xT_L`$.
A point $`x`$ is a regular point of a leaf $`L`$ if $`T_x=T_L`$. Regular points are determined by the condition
$$J(x)=det(c_{pi}^qx^i)0.$$
(13)
If $`J(x)=0`$ at a point $`x`$, then $`T_x`$ is a proper subspace of $`T_L`$, and a point $`x`$ is said to be a singular point of a leaf $`L`$.
The determinant (13) is the Jacobian of the map $`f:P^l\times M^rP^N`$. Singular points of a leaf $`L`$ are determined by the condition $`J(x)=0`$. In a leaf $`L`$, they form an algebraic submanifold of dimension $`l1`$ and degree $`r`$. This hypersurface (in $`L`$) is called the focus hypersurface and is denoted by $`F_L`$. By (13), the equations $`J(x)=0`$ of the focus hypersurface on the plane generator $`L`$ of the manifold $`X`$ can be written as
$$det(c_{pi}^qx^i)=0.$$
(14)
We calculate now the second differential of a point $`xL`$:
$$d^2xA_\alpha \omega _s^\alpha \omega _i^sx^i(modT_x).$$
This expression is the second fundamental quadratic form of the manifold $`X`$:
$$II_x=A_\alpha \omega _s^\alpha \omega _i^sx^i=A_\alpha b_{ps}^\alpha c_{qi}^sx^i\theta ^p\theta ^q.$$
(15)
Suppose that $`\xi =\xi _\alpha x^\alpha =0`$ is the tangent hyperplane to $`X`$ at $`xL,\xi T_L`$. Then
$$(\xi ,II_x)=h_{pq}(\xi ,x)\theta ^p\theta ^q,$$
where
$$h_{pq}(\xi ,x)=\xi _\alpha b_{ps}^\alpha c_{qi}^sx^i,h_{pq}=h_{qp},$$
is the second fundamental quadratic form of the manifold $`X`$ at $`x`$ with respect to the hyperplane $`\xi `$. Since at regular points $`xL`$ condition (13) holds, the rank of this matrix is the same as the rank of the matrix
$$B(\xi )=(\xi _\alpha b_{pq}^\alpha )=\xi _\alpha B^\alpha ,$$
(16)
and this rank is the same at all regular points $`xL`$.
We call a tangent hyperplane $`\xi `$ singular if
$$det(\xi _\alpha b_{pq}^\alpha )=0,$$
(17)
i.e., if the rank of the matrix (16) is reduced. Condition (17) is an equation of degree $`r`$ with respect to the tangential coordinates $`\xi _\alpha `$ of the hyperplane $`\xi `$. This condition defines an algebraic hypercone whose vertex is the tangent subspace $`T_L`$. This hypercone is called the focus hypercone and denoted by $`\mathrm{\Phi }_L`$ (see \[AG 93\], p. 119).
The determinant $`det(\xi _\alpha b_{pq}^\alpha )`$ in the left-hand side of equation (17) is the Jacobian of the dual map $`f^{}:L^{}\times M^r(P^N)^{},f^{}(L^{}\times M^r)=X^{}(P^N)^{}`$, where $`X^{}`$ is a manifold dual to $`X`$ and $`L^{}`$ is a bundle of hyperplanes of the space $`P^N`$ passing through the tangent subspace $`T_L`$ of the manifold $`X,dimL^{}=Nn1`$.
The focus hypersurface $`F_LL`$ and the focus hypercone $`\mathrm{\Phi }_L`$ with vertex $`T_L`$ are called focal images of the manifold $`X`$ with a degenerate Gauss map.
Note that under the passage from the manifold $`XP^N`$ to its dual manifold $`X^{}(P^N)^{}`$, the systems of square matrices $`C_i`$ and $`B^\alpha `$ as well as the focus hypersurfaces $`F_L`$ and the focus cones $`\mathrm{\Phi }_L`$ exchange their roles.
Since
$$d^2xA_\alpha b_{qs}^\alpha c_{pi}^sx^i\theta ^p\theta ^q(modT_L,xL),$$
the points
$$A_{pq}=A_\alpha b_{qs}^\alpha c_{pi}^sx^i,A_{pq}=A_{qp},$$
(18)
together with the points $`A_i`$ and $`A_p`$ define the osculating subspace $`T_L^2(X)`$. Its dimension is
$$dimT_L^2(X)=n+m,$$
where $`m`$ is the number of linearly independent points among the points $`A_{pq},m\frac{r(r+1)}{2}`$. But since at a regular point $`xL`$ condition (13) holds, the number $`m`$ is the number of linearly independent points among the points
$$\stackrel{~}{A}_{pq}=A_\alpha b_{pq}^\alpha .$$
We also use the notation $`S_L`$ for the osculating space $`T_L^2(X)`$.
On a generator $`L`$ of the manifold $`X`$, consider the system of equations
$$c_{pi}^qx^i=0.$$
(19)
Its matrix $`C=(c_{pi}^q)`$ has $`r^2`$ rows and $`l+1`$ columns. Denote the rank of this matrix by $`m^{}`$. If $`m^{}<l+1`$, then system (19) defines a subspace $`K_L`$ of dimension $`k=lm^{}`$ in $`L`$. This subspace belongs to the focus hypersurface $`F_L`$ defined by equation (14). If $`l>m^{}`$, then the hypersurface $`F_L`$ becomes a cone with vertex $`K_L`$. We call the subspace $`K_L`$ the characteristic subspace of the generator $`L`$.
Note also that by the duality principle in $`P^N`$, the osculating subspace $`S_L`$ and the characteristic subspace $`K_L`$ constructed for a pair $`(L,T_L)`$ correspond one to another.
## 5 Proof of Theorem 2
###### Lemma 6
Suppose that $`l1`$, and the focus hypersurface $`F_LL`$ does not have multiple components. Then all matrices $`B^\alpha `$ can be simultaneously diagonalized, $`B^\alpha =(\text{diag}b_{pp}^\alpha )`$, and the focus hypercone $`\mathrm{\Phi }_L`$ decomposes into $`r`$ bundles of hyperplanes $`\mathrm{\Phi }_p`$ in $`P^N`$ whose axes are $`(n+1)`$-planes $`T_LB_p`$, where $`B_p=b_{pp}^\alpha A_\alpha `$ are points located outside of the tangent subspace $`T_L`$. The dimension $`n+m`$ of the osculating subspace $`S_L`$ of the manifold $`X`$ along a generator $`L`$ does not exceed $`n+r`$.
Proof. Since the hypersurface $`F_LL`$ does not have multiple components, a general straight line $`\lambda `$ lying in $`L`$ intersects $`F_L`$ at $`p`$ distinct points. We place the vertices $`A_0`$ and $`A_1`$ of our moving frame onto the line $`\lambda `$. By (14), the coordinates of the common points of $`\lambda `$ and $`F_L`$ are defined by the equation
$$det(c_{p0}^qx^0+c_{p1}^qx^1)=0.$$
Suppose that $`A_0F_L`$. Then $`det(c_{p0}^q)0`$, and it is easy to prove that the matrices $`C_0`$ and $`C_1`$ can be simultaneously diagonalized, $`C_0=(\delta _q^p)`$ and $`C_1=(\text{diag}c_{p1}^p)`$. Since the common points of $`\lambda `$ and $`F`$ are not multiple points, we have $`c_{p1}^pc_{q1}^q`$ for $`pq`$.
Next we write equations (12) for $`i=0,1`$:
$$b_{pq}^\alpha =b_{qp}^\alpha ,b_{qp}^\alpha c_{p1}^p=b_{pq}^\alpha c_{q1}^q.$$
Since $`c_{p1}^pc_{q1}^q`$ , it follows that $`b_{pq}^\alpha =0`$ for $`pq`$. As a result, all matrices $`B^\alpha `$ can be simultaneously diagonalized, $`B^\alpha =(\text{diag}b_{pp}^\alpha )`$. Equation (17) takes the form
$$\underset{p}{}(\xi _\alpha b_{pp}^\alpha )=0,$$
and the focus hypercone $`\mathrm{\Phi }_L`$ decomposes into $`r`$ bundles of hyperplanes $`\mathrm{\Phi }_p`$ in $`P^N`$ whose axes are $`(n+1)`$-planes $`T_LB_p`$, where $`B_p=b_{pp}^\alpha A_\alpha `$ are points located outside of the tangent subspace $`T_L`$. The osculating subspace $`S_L`$ of the manifold $`X`$ along a generator $`L`$ is the span of the tangent subspace $`T_L`$ and the points $`B_{l+1},\mathrm{},B_n`$. Thus, the dimension of this subspace does not exceed $`n+r`$.
###### Lemma 7
Suppose that $`l1,m2`$, the focus hypersurfaces $`F_LL`$ do not have multiple components, and all the bundles $`\mathrm{\Phi }_p`$ into which the hypercone $`\mathrm{\Phi }_L`$ decomposes are of multiplicity one. Then the focus hypersurface $`F_L`$ decomposes into $`r`$ hyperplanes $`F_pL`$.
Proof. Consider the matrix
$$B=(b_{pp}^\alpha )$$
(20)
composed from the eigenvalues of the matrices $`B^\alpha `$. Matrix (20) has $`r`$ columns and $`m`$ independent rows, $`mr`$. Write equations (12) for the diagonal matrices $`B^\alpha `$:
$$b_{qq}^\alpha c_{pi}^q=b_{pp}^\alpha c_{qi}^p.$$
(21)
Since the matrix $`B`$ has $`m`$ linearly independent columns and $`m2`$, it follows from equation (21) that $`c_{pi}^q=0`$ for $`pq`$, and all the matrices $`C_i`$ and $`B^\alpha `$ can be simultaneously diagonalized. Now the equation of the focus hypersurface $`F_L`$ takes the form
$$\underset{p=l+1}{\overset{n}{}}c_{pi}^px^i=0,$$
and the hypersurface $`F_L`$ decomposes into $`r`$ hyperplanes $`F_p`$ defined in $`L`$ by the equation $`c_{pi}^px^i=0`$. Since by the conditions of Theorem 2, the focus hypersurface $`F_L`$ does not have multiple components, all hyperplanes $`F_p`$ are distinct.
Consider a rectangular $`r\times (l+1)`$ matrix
$$C=(c_{pi}^p)$$
(22)
formed by the eigenvalues of the matrix $`C_i`$.
Proof of Theorem 2. From the conditions of Theorem 2 and Lemmas 6 and 7, it follows that the matrices $`C_i`$ and $`B^\alpha `$ can be simultaneously diagonalized:
$$C_i=\text{diag}(c_{l+1,i}^{l+1},\mathrm{},c_{ni}^n),B^\alpha =\text{diag}(b_{l+1,l+1}^\alpha ,\mathrm{},b_{nn}^\alpha ).$$
This implies that formulas (11) take the form:
$$\omega _i^p=c_{pi}^p\theta ^p,\omega _p^\alpha =b_{pp}^\alpha \theta ^p,$$
(23)
where $`\theta ^p`$ are basis forms on the parametric variety $`M`$, and there is no summation over the index $`p`$. Exterior differentiation of equations (23) leads to the following exterior quadratic equations:
$$c_{pi}^p\theta ^p+c_{pi}^pd\theta ^p+\underset{qp}{}c_{qi}^q\omega _q^p\theta ^q=0,$$
(24)
$$b_{pp}^\alpha \theta ^p+b_{pp}^\alpha d\theta ^p\underset{qp}{}b_{qq}^\alpha \omega _p^q\theta ^q=0,$$
(25)
where
$$c_{pi}^p=dc_{pi}^pc_{pj}^p\omega _i^j+c_{pi}^p\omega _p^p,b_{pp}^\alpha =db_{pp}^\alpha +b_{pp}^\beta \omega _\beta ^\alpha b_{pp}^\alpha \omega _p^p.$$
Suppose now that at least one of the matrices $`B`$ and $`C`$ defined by equations (20) and (22) has mutually linearly independent columns. Then it follows from equations (24) and (25) that
$$\{\begin{array}{cc}d\theta ^p0(mod\theta ^p),\hfill & \\ \omega _q^p\theta ^q0(mod\theta ^p),\hfill & \\ \omega _p^q\theta ^q0(mod\theta ^p).\hfill & \end{array}$$
This implies that
$$d\theta ^p=\theta ^p\theta _p^p,\omega _p^q=l_{pp}^q\theta ^p+l_{pq}^q\theta ^q.$$
(26)
We write the expression of $`dA_p`$ from (9) in more detail:
$$dA_p=\omega _p^iA_i+\omega _p^pA_p+\underset{qp}{}\omega _p^qA_q+\omega _p^\alpha A_\alpha .$$
(27)
Let the index $`p`$ be fixed and $`qp;p,q=l+1,\mathrm{},n`$. Consider a submanifold $`X_pX`$ defined by the equations
$$\theta ^q=0,qp.$$
(28)
The first equation of (26) implies that $`dimX_p=1`$. By (28), (9) and (27), we find that
$$dA_i=\omega _i^jA_j+c_{pi}^p\theta ^pA_p,$$
(29)
$$dA_p=\omega _p^iA_i+\omega _p^pA_p+B_p\theta ^p,$$
(30)
where $`B_p=_{qp}l_{pp}^qA_q+b_{pp}^\alpha A_\alpha `$. Since the points $`A_i`$ are basis points of a generator $`L`$ of the manifold $`X`$, equations (29) and (30) prove that the subspace $`LA_p`$ is tangent to the submanifold $`X_p`$ at all regular points of the generator $`L`$, and the subspace $`LA_pB_p`$ is the osculating subspace of $`X_p`$. Singular points of a generator $`L`$ of the manifold $`X_p`$ are determined by the equations $`c_{pi}^px^i=0`$ and form a hyperplane in $`L`$ (see Example 2 in Section 2). Thus, the submanifold $`X_p`$ is a torse with $`l`$-dimensional plane generators. Hence the manifold $`X`$ is foliated into $`r`$ families of torses, each of these families depends on $`r1`$ parameters $`u^q`$, and the forms $`\theta ^q`$ are expressed in terms of the differentials of these parameters, $`\theta ^q=l^qdu^q`$.
## 6 Proof of Theorems 3 and 4
###### Lemma 8
Suppose that $`l2`$, the focus hypersurfaces $`F_LL`$ do not have multiple components and are indecomposable. Then the hypercone $`\mathrm{\Phi }_L`$ is an $`r`$-multiple bundle of hyperplanes with an $`(n+1)`$-dimensional vertex in $`P^N`$.
Proof. From the conditions of the lemma and equation (21) it follows that the columns of the matrix $`B`$ are linearly dependent, and thus all points $`B_p=b_{pp}^\alpha A_\alpha `$ defined by the columns of this matrix lie in a subspace of dimension $`n+1`$ defined by the tangent subspace $`T_L`$ and one of the points $`B_p`$.
Therefore all bundles of hyperplanes $`\mathrm{\Phi }_p`$ into which the hypercone $`\mathrm{\Phi }_L`$ decomposes coincide.
We formulate the lemma which is dual to Lemma 8.
###### Lemma 9
Suppose that $`m2`$ and the focus hypercones $`\mathrm{\Phi }_L`$ with their vertices $`T_L`$ do not have multiple components and are indecomposable. Then the focus hypersurface $`F_LL`$ is an $`r`$-multiple hyperplane in $`L`$.
Proof of Theorem 3. From the conditions of Theorem 3 and Lemma 8, it follows that all pairs of columns of the matrix $`B`$ are linearly dependent. Thus all matrices $`B^\alpha `$ associated with $`X`$ are mutually linearly dependent. Hence, we have
$$b_{pq}^\alpha =b^\alpha b_{pq}.$$
(31)
Now conditions (12) take the form
$$b_{ps}c_{qi}^s=b_{qs}c_{pi}^s.$$
(32)
Although the matrix $`B=(b_{pq})`$ is still can be diagonalized, in general, the matrices $`C_i`$ do not possess this property. Thus we write equations (11) in the form
$$\omega _i^p=c_{qi}^p\theta ^q,\omega _p^\alpha =b^\alpha b_{pq}\theta ^q.$$
(33)
Exterior differentiation of equations (33) and applying structure equations (7) of the parametric variety $`M`$, we obtain the following exterior quadratic equations:
$$c_{qi}^p\theta ^q=0,(b_{pq}b^\alpha +b^\alpha b_{pq})\theta ^q=0,$$
(34)
where
$$\{\begin{array}{cc}c_{qi}^p=dc_{qi}^pc_{qj}^p\omega _i^j+c_{qi}^s\omega _s^pc_{si}^p\theta _p^s,\hfill & \\ b^\alpha =db^\alpha +b^\beta \omega _\beta ^\alpha ,\hfill & \\ b_{pq}=db_{pq}b_{sq}\omega _p^sb_{ps}\theta _q^s.\hfill & \end{array}$$
As we noted already in the proof of Theorem 2, if the point $`A_0`$ does not belong to the focus hypersurface $`F_L`$ of a generator $`L`$, then the matrix $`(c_{q0}^p)`$ is nonsingular, and by means of a frame transformation in the tangent space $`T_u(M)`$ to the variety $`M`$ we can reduce this matrix to the form $`(c_{q0}^p)=(\delta _q^p)`$. As a result, equation (32) corresponding to $`i=0`$ takes the form
$$b_{pq}=b_{qp}.$$
(35)
Hence the matrix $`B=(b_{pq})`$ becomes symmetric. This matrix is the matrix of the second fundamental form $`II`$ of the manifold $`X`$ at the point $`A_0`$. Since $`A_0`$ is a regular point of $`X`$, the matrix $`B`$ is nonsingular, $`det(b_{pq})0`$.
Next, we find from the second equation of (34) that
$$b^\alpha =b_p^\alpha \theta ^p,b_{pq}=b_{pqs}\theta ^s.$$
Substituting these expansions into equation (34) and equating to 0 the coefficients in independent products $`\theta ^s\theta ^q`$, we find that
$$b_{pq}b_s^\alpha b_{ps}b_q^\alpha +b^\alpha (b_{pqs}b_{psq})=0.$$
Contracting these equations with the matrix $`B^1=(b^{pq})`$ which is the inverse of $`B`$, we find that
$$(r1)b_s^\alpha +b^\alpha (b_{pqs}b_{psq})b^{pq}=0.$$
Since by the theorem hypothesis $`r1`$, it follows that
$$b_s^\alpha =b^\alpha b_s,$$
(36)
where
$$b_s=\frac{1}{r1}(b_{pqs}b_{psq})b^{pq}.$$
Next consider equations (17) of the focus hypercone $`\mathrm{\Phi }_L`$ of the manifold $`X`$. By (31), this equation becomes
$$(b^\alpha \xi _\alpha )^rdet(b_{pq})=0,$$
and since $`det(b_{pq})0`$, the hypercone $`\mathrm{\Phi }_L`$ becomes an $`r`$-multiple bundle of hyperplanes. The axis of this bundle is a subspace of dimension $`n+1`$ which is the span of the tangent subspace $`T`$ and the point $`b^\alpha A_\alpha `$.
Let us prove that if the parameters $`u`$ on the variety $`M`$ vary, this subspace is fixed. In fact, the basis points of the subspace $`T`$ are the points $`A_i`$ and $`A_p`$, and by (9) and (36), we have
$$\{\begin{array}{cc}dA_i=\omega _i^jA_j+\omega _i^pA_p,\hfill & \\ dA_p=\omega _p^iA_i+\omega _p^qA_q+b_{pq}\theta ^qb^\alpha A_\alpha .\hfill & \end{array}$$
If we differentiate the point $`b^\alpha A_\alpha `$ and apply equation (36), we find that
$$d(b^\alpha A_\alpha )b_s\theta ^sb^\alpha A_\alpha (modA_i,A_p).$$
This proves our last assertion.
Thus the subspace $`P^{n+1}=Tb^\alpha A_\alpha `$ is fixed when the tangent subspace $`T`$ moves along $`X`$, and $`X`$ is a hypersurface in the subspace $`P^{n+1}`$.
Proof of Theorem 4. Theorem 4 is dual to Theorem 3 and can be proved by applying Lemma 9 in the same way as we used Lemma 8 to prove Theorem 3.
## 7 Proof of Theorem 5
Equations (14) and (17) of the focal images imply the following lemma.
###### Lemma 10
If the matrices $`C_i`$ and $`B^\alpha `$ of a manifold $`X`$ can be reduced to the form $`(1)`$, then each of its focus hypersurfaces $`F_LL`$ decomposes into $`s`$ components $`F_t`$ of dimension $`l1`$ and degree $`r_1,r_2,\mathrm{},r_s`$, and each of its focus hypercones $`\mathrm{\Phi }_L`$ decomposes into $`s`$ hypercones $`\mathrm{\Phi }_t`$ of the same degrees $`r_1,r_2,\mathrm{},r_s;r_1+r_2+\mathrm{}+r_s=r,`$ and with the same vertex $`T`$. In particular, if $`r_1=r_2=\mathrm{}=r_s=1`$, then a focus hypersurface $`F_L`$ decomposes into $`r`$ hyperplanes, and a focus hypercone $`\mathrm{\Phi }_L`$ decomposes into $`r`$ bundles of hyperplanes with $`(n+1)`$-dimensional axes.
Proof of Theorem 5.
We prove Theorem 5 assuming that the index $`t`$ takes only two values, $`t=1,2,r=r_1+r_2`$, and the indices $`p`$ and $`q`$ have the following values:
$$p_1,q_1=l+1,\mathrm{},l+r_1,p_2,q_2=l+r_1+1,\mathrm{},n.$$
Then equations (11) become
$$\{\begin{array}{cc}\omega _i^{p_1}=c_{q_1i}^{p_1}\theta ^{q_1},\hfill & \omega _{p_1}^\alpha =b_{p_1q_1}^\alpha \theta ^{q_1},\hfill \\ \omega _i^{p_2}=c_{q_2i}^{p_2}\theta ^{q_2},\hfill & \omega _{p_2}^\alpha =b_{p_2q_2}^\alpha \theta ^{q_2}.\hfill \end{array}$$
(37)
Exterior differentiation of equations (37) gives
$$c_{q_1i}^{p_1}\theta ^{q_1}+(c_{q_2i}^{s_2}\omega _{s_2}^{p_1}c_{s_1i}^{p_1}\theta _{q_2}^{s_1})\theta ^{q_2}=0,$$
(38)
$$b_{p_1q_1}^\alpha \theta ^{q_1}(b_{s_2q_2}^\alpha \omega _{p_1}^{s_2}+b_{p_1s_1}^\alpha \theta _{q_2}^{s_1}\theta ^{q_2}=0,$$
(39)
$$c_{q_2i}^{p_2}\theta ^{q_2}+(c_{q_1i}^{s_1}\omega _{s_1}^{p_2}c_{s_2i}^{p_2}\theta _{q_1}^{s_2})\theta ^{q_1}=0,$$
(40)
$$b_{p_2q_2}^\alpha \theta ^{q_2}(b_{s_1q_1}^\alpha \omega _{p_2}^{s_1}+b_{p_2s_2}^\alpha \theta _{q_1}^{s_2}\theta ^{q_1}=0,$$
(41)
where
$$\{\begin{array}{cc}c_{q_1i}^{p_1}=dc_{q_1i}^{p_1}c_{q_1j}^{p_1}\omega _i^j+c_{q_1i}^{s_1}\omega _{s_1}^{p_1}c_{s_1i}^{p_1}\theta _{q_1}^{s_1},\hfill & \\ b_{p_1q_1}^\alpha =db_{p_1q_1}^\alpha +b_{p_1q_1}^\beta \omega _\beta ^\alpha b_{s_1q_1}^\alpha \omega _{p_1}^{s_1}b_{p_1s_1}^\alpha \theta _{s_1}^{q_1},\hfill & \\ c_{q_2i}^{p_2}=dc_{q_2i}^{p_2}c_{q_2j}^{p_2}\omega _i^j+c_{q_2i}^{s_2}\omega _{s_2}^{p_2}c_{s_2i}^{p_2}\theta _{q_2}^{s_2},\hfill & \\ b_{p_2q_2}^\alpha =db_{p_2q_2}^\alpha +b_{s_2q_2}^\beta \omega _\beta ^\alpha b_{s_2q_2}^\alpha \omega _{p_2}^{s_2}b_{p_2s_2}^\alpha \theta _{q_2}^{s_2}.\hfill & \end{array}$$
Consider the system of equations
$$\theta ^{q_1}=0$$
(42)
on the manifold $`X`$. By (8), its exterior differentiation gives
$$\theta ^{q_2}\theta _{q_2}^{q_1}=0.$$
(43)
It follows from (43) that the conditions of complete integrability of equations (42) have the form
$$\theta _{q_2}^{q_1}=l_{q_2s_2}^{q_1}\theta ^{s_2},$$
(44)
where $`l_{q_2s_2}^{q_1}=l_{s_2q_2}^{q_1}.`$
By equations (42), the system of equations (38) takes the form
$$(c_{q_2i}^{s_2}\omega _{s_2}^{p_1}c_{s_1i}^{p_1}\theta _{q_2}^{s_1})\theta ^{q_2}=0.$$
(45)
Suppose that the component $`F_1`$ of the focus hypersurface $`F_L`$ does not have multiple components. Assuming that $`l1`$, we write equations (45) for two different values of the index $`i`$, for example, for $`i=0,1`$. Since the matrices $`(c_{s_1i}^{p_1})`$ and $`(c_{s_2i}^{p_2})`$ are not proportional, then it follows from (45) that two terms occurring in (45) vanish separately. In particular, this means that
$$c_{s_1i}^{p_1}\theta _{q_2}^{s_1}\theta ^{q_2}=0.$$
(46)
Since the number of linearly independent forms among the 1-forms $`\omega _i^{p_1}`$ connected with the basis forms by relations (47) is equal to the number of linearly independent forms $`\theta ^{q_1}`$ (i.e., it is equal $`r_1`$), then it follows from (45) that
$$\theta _{q_2}^{s_1}\theta ^{q_2}=0.$$
But the last equations coincide with equations (43) and are conditions of complete integrability of (42). Thus the manifold $`X`$ is foliated into an $`r_1`$-parameter family of submanifolds of dimension $`l+r_2`$ and of rank $`r_2`$, and these submanifolds belong to the types described in Theorems 2 or 3.
In a similar way, one can prove the complete integrability of equations $`\theta ^{q_2}=0`$ on the manifold $`X`$. Thus the manifold $`X`$ is foliated also into an $`r_2`$-parameter family of submanifolds of dimension $`l+r_1`$ and of rank $`r_1`$.
By induction over $`s`$, we can prove the result, which we have proved for $`s=2`$ components, for the case of any number $`s`$ of components.
Thus, Theorem 5 describes the structure of tangentially degenerate manifolds of general types. As a result, this theorem is a structure theorem for such manifolds.
Note that the torsal manifolds described in Theorem 2 are completely reducible, and the manifolds $`X`$ described in Theorems 3 and 4 are irreducible manifolds.
Note that Theorem 5 does not cover tangentially degenerate submanifolds with multiple nonlinear components of their focal images. This gives rise to the following problem.
Problem. Construct an example of a submanifold $`XP^N(𝐂)`$ with a degenerate Gauss map whose focal images have multiple nonlinear components or prove that such submanifolds do not exist.
## 8 Additional results
In conclusion we prove two additional theorems.
###### Theorem 11
Let $`XP^N`$ be a tangentially degenerate submanifold of dimension $`n`$ and rank $`r<n`$. Suppose that all matrices $`B^\alpha `$ can be simultaneously diagonalized, $`B^\alpha =\text{diag}(b_{l+1,l+1}^\alpha ,\mathrm{},b_{nn}^\alpha )`$. Suppose also that the rectangular matrix $`B`$ $`(`$defined by $`(20))`$ composed from the eigenvalues of the matrices $`B^\alpha `$ has a rank $`r_1r1`$, and this rank does not reduce when we delete any column of this matrix. Then the submanifold $`X`$ belongs to a subspace $`P^{n+r_1}`$ of the space $`P^N`$.
Proof. Under the conditions of Theorem 11, the second group of equations (11) takes the form
$$\omega _p^\alpha =b_{pp}^\alpha \theta ^p,p=l+1,\mathrm{},n,\alpha =n+1,\mathrm{},N.$$
(47)
The matrix $`B`$ has only $`r_1`$ linearly independent rows. Thus by means of transformations of moving frame’s vertices located outside of the tangent subspace $`T_L`$, equations (47) can be reduced to the form
$$\omega _p^\lambda =b_{pp}^\lambda \theta ^p,\omega _p^\sigma =0,$$
(48)
where $`\lambda =n+1,\mathrm{},n+r_1,\sigma =n+r_1+1,\mathrm{},N.`$ The second group of equations (9) takes the form
$$dA_p=\omega _p^iA_i+\omega _p^qA_q+\omega _p^\lambda A_\lambda ,$$
and the points $`A_\lambda `$ together with the points $`A_i`$ and $`A_q`$ define the osculating subspace $`S_L`$ of the submanifold $`X`$ for all points $`xL`$. The dimension of $`S_L`$ is $`n+r_1,dimS_L=n+r_1`$.
Differentiation of the points $`A_\lambda `$ gives
$$dA_\lambda =\omega _\lambda ^iA_i+\omega _\lambda ^pA_p+\omega _\lambda ^\mu A_\mu +\omega _\lambda ^\sigma A_\sigma ,$$
(49)
where $`\lambda ,\mu =n+1,\mathrm{},n+r_1;\sigma =n+r_1+1,\mathrm{},N`$. If $`\theta ^p=0`$, then the osculating subspace $`S_L`$ of $`X`$ remains fixed. It follows from equations (49) that the 1-forms $`\omega _\lambda ^\rho `$ are expressed in terms of the basis forms $`\theta ^p`$ of $`X`$, that is,
$$\omega _\lambda ^\rho =l_{\lambda p}^\rho \theta ^p.$$
(50)
Taking exterior derivatives of the second group of equations (48), we find that
$$\omega _p^\lambda \omega _\lambda ^\rho =0.$$
(51)
Substituting the values of the 1-forms $`\omega _p^\lambda `$ and $`\omega _\lambda ^\rho `$ from equations (48) and (50) into equation (51), we find that
$$b_{pp}^\lambda \theta ^pl_{\lambda q}^\rho \theta ^q=0.$$
In this equation the summation is carried over the indices $`\lambda `$ and $`q`$, but there is no summation over the index $`p`$. It follows from these equations that
$$b_{pp}^\lambda l_{\lambda q}^\rho =0,pq.$$
(52)
System (52) is a system of linear homogeneous system with respect to the unknown variables $`l_{\lambda q}^\rho `$. For each pair of the values $`\rho `$ and $`q`$, system (52) has the rank $`r1`$ and $`r_1`$ unknowns. Since $`r_1r1`$, under the conditions of Theorem 11, the rank of the matrix of coefficients of this system is equal $`r_1`$. As a result, the system has only the trivial solution $`l_{\lambda q}^\rho =0`$. Thus equations (50) take the form
$$\omega _p^\lambda =0.$$
(53)
It follows from (49) and (53) that the osculating subspace $`S_L`$ of $`X`$ remains fixed when $`L`$ moves in $`X`$. Thus $`XP^{n+r_1}`$.
Remark. If $`r_1=r`$ and $`N>n+r`$, then the osculating subspace $`S_L`$ of $`X`$ can move in $`P^N`$ when $`L`$ moves in $`X`$. In this case the submanifold $`X`$ is torsal.
Theorem 11 is similar to Theorem 3.10 from the book \[AG 93\] and was proved in \[AG 93\] for submanifolds of a space $`P^N`$ bearing a net of conjugate lines. Note that Theorem 3.10 from \[AG 93\] generalizes a similar theorem of C. Segre (see \[Se 07\], p. 571) proved for submanifolds $`X`$ of dimension $`n`$ of the space $`P^N`$ which has at each point $`xX`$ the osculating subspace $`S_x`$ of dimension $`n+1`$. By this theorem, a submanifold $`X`$ either belongs to a subspace $`P^{n+1}`$ or is a torse.
The theorem dual to Theorem 11 is also valid.
###### Theorem 12
Let $`XP^N`$ be a tangentially degenerate submanifold of dimension $`n`$ and rank $`r<n`$. Suppose that all matrices $`C_i`$ can be simultaneously diagonalized, $`C_i=\text{diag}(c_{l+1,i}^{l+1},\mathrm{},c_{ni}^n)`$. Suppose also that the rectangular matrix $`C=(c_{pi}^p)`$ composed from the eigenvalues of the matrices $`C_i`$ has a rank $`r_2r1`$, and this rank is not reduced when we delete any column of this matrix. Then the submanifold $`X`$ is a cone with an $`(lr_2)`$-dimensional vertex $`K_L`$.
Proof. The proof of this theorem is similar to the proof of Theorem 11.
Authors’ addresses:
| M. A. Akivis | V. V. Goldberg |
| --- | --- |
| Department of Mathematics | Department of Mathematical Sciences |
| Jerusalem College of Technology—Mahon Lev | New Jersey Institute of Technology |
| Havaad Haleumi St., P. O. B. 16031 | University Heights |
| Jerusalem 91160, Israel | Newark, N.J. 07102, U.S.A. |
| E-mail address: akivis@avoda.jct.ac.il | E-mail address: vlgold@m.njit.edu | |
warning/0002/astro-ph0002306.html | ar5iv | text | # Flash-Heating of Circumstellar Clouds by Gamma Ray Bursts
## 1 Introduction
The identification of flaring and fading X-ray, optical and radio counterparts to gamma-ray burst (GRB) sources (e.g., Costa et al. 1997; van Paradijs et al. 1997; Djorgovski et al. 1997; Frail et al. 1997), and the large energy releases implied by redshift measurements, find a consistent explanation in an expanding relativistic blast-wave model (Paczyński & Rhoads 1993; Mészáros & Rees 1997). As a result of Beppo-SAX and optical follow-on observations, the redshifts of about one dozen GRBs with durations greater than $`1`$ s have been measured. The distribution of redshifts is broad and centered near $`z1`$, corresponding to the cosmological epoch of active star formation (Hogg & Fruchter 1999). GRBs are extremely luminous and energetic at hard X-ray and $`\gamma `$-ray energies. The degree of GRB collimation is unknown, but peak directional $`\gamma `$-ray luminosities and energy releases as large as $`L/\mathrm{\Omega }3\times 10^{51}`$ ergs (s-sr)<sup>-1</sup> and $`E/\mathrm{\Omega }3\times 10^{53}`$ ergs sr<sup>-1</sup>, respectively, have been measured (Kulkarni et al. 1999). Less powerful GRBs and less luminous episodes during the GRB produce smaller $`\gamma `$-ray powers, but the apparent isotropic $`\gamma `$-ray luminosities from typical GRBs could regularly reach values exceeding $`10^{50}L_{50}`$ ergs s<sup>-1</sup> with $`L_{50}1`$, with some GRBs reaching $`L_{50}>10^2`$. Because the energy radiated in $`\gamma `$ rays is less than the total energy released by a GRB, the apparent isotropic energy release of GRB sources could often reach values of $`10^{54}E_{54}`$ ergs, with $`E_{54}`$1.
Cosmological gamma-ray burst and afterglow observations are best explained through the fireball/blast-wave model, where the deposition of large quantities of energy into a small region yields a fireball that expands until it reaches a relativistic speed determined by the amount of baryons mixed into the fireball (see, e.g., Piran 1999 for a review). Nonthermal synchrotron radiation from energetic electrons in the relativistic blast wave is thought to account for the origin of the prompt $`\gamma `$-ray emission and afterglow radiation. This paradigm has been called into question, however, by observations of very hard X-ray emission during the prompt $`\gamma `$-ray luminous phase of a significant number of GRBs (Crider et al. 1997; Preece et al. 1998). Photon fluxes $`\varphi (ϵ)ϵ^{\alpha _X}`$ with $`\alpha _X`$ 0, where $`ϵ=h\nu /m_ec^2`$ is the dimensionless photon energy, have been observed in 5-10% of GRBs that are bright enough to permit spectral analysis. This strongly contradicts the optically-thin synchrotron shock model, which predicts that only radiation spectra with $`\alpha _X2/3`$ can emerge from the blast wave. In view of the severity of this challenge to the model, these observations have been termed the “line-of-death” to the synchrotron shock model. Possible explanations for this phenomeonon involve photoelectric absorption by optically thick cold matter (Liang & Kargatis 1994; Brainerd 1994; Böttcher et al. 1999), synchrotron self-absorption (Crider & Liang 1999, Granot, Piran, & Sari 2000, Lloyd & Petrosian 1999), Compton scattering (Liang 1997; Liang et al. 1999), or the existence of a pair-photosphere (Mészáros & Rees 2000) within the blast wave. Except for the last model cited, these explanations are inconsistent with the standard synchrotron shock model. Here we offer a solution to this problem that is consistent with the standard model and recent observations pointing to a massive star origin of GRBs.
## 2 Massive Star Origin of GRBs
Considerable evidence linking the sources of GRBs with star-forming regions has recently been obtained (e.g., Lamb 1999). For example, the associated host galaxies have blue colors, consistent with galaxy types that are undergoing active star formation. GRB counterparts are found within the optical radii and central regions of the host galaxies (e.g. Bloom et al. 1999a), rather than far outside the galaxies’ disks, as might be expected in a scenario of merging neutron stars and black holes (Narayan, Paczyński, & Piran 1992). Lack of optical counterparts in some GRBs could be due to extreme reddening from large quantities of gas and dust in the host galaxy. This, together with the appearance of supernova-like emissions in the late time optical decay curves of a few GRBs (e.g., Bloom et al. 1999b) and weak X-ray evidence for Fe K<sub>α</sub>-line signatures (Piro et al. 1999), supports a massive star hypernova/collapsar (Woosley 1993; Paczyński 1998) origin for the long duration gamma-ray bursters.
The observations thus favor a model for GRBs involving the collapse of the core of a $`\stackrel{>}{}30M_{}`$ star to a black hole, with the collapse events producing fireballs and relativistic outflows with large directed energy releases. Earlier treatments of the blast-wave model considered systems where the density of the surrounding medium is either uniform or monotonically decreasing as a result of a circumstellar medium formed by a hot stellar wind (Mészáros, Rees, & Wijers 1998). Until recently (Chevalier & Li 1999; Li & Chevalier 1999), less attention has been paid to the actual environment found in the vicinity of massive stars. For this we consider $`\eta `$ Carinae (Davidson & Humphreys 1997), the best-studied massive star that might correspond to a GRB progenitor. It is an evolved star with present-day mass $`90M_{}`$, distance of $`2300\pm 200`$ pc, and lifetime of about 3 million years. It anisotropically ejects mass at a current rate of $`\stackrel{<}{}0.003M_{}`$ yr<sup>-1</sup> to form its unusual “homonculus nebula.” Several Solar masses of material surround $`\eta `$ Carinae. In the immediate vicinity of the central star, dense clouds of slow-moving gas with radii $`r10^{15}`$ cm and densities between $`10^7`$ and $`10^{10}`$ cm<sup>-3</sup> were discovered with speckle techniques (Hofman & Weigelt 1988) and confirmed with high-resolution HST observations (Davidson et al. 1995). This material, moving with speeds of $`50`$ km s<sup>-1</sup>, is apparently ejected nonuniformly from the equatorial zone, but may remain trapped by the gravitational field of the star. Inferences (Davidson & Humphreys 1997) from \[FeII\] observations suggest that $`\stackrel{>}{}0.02M_{}`$ of gas are contained within $`2\times 10^{16}`$ cm, implying a volume-averaged gas density $`\stackrel{>}{}7\times 10^5`$ cm<sup>-3</sup>. Model results (Böttcher 1999) imply that a dense ($`n10^{12}`$ cm<sup>-3</sup>) torus of gas at mean distance $`d2\times 10^{15}`$ cm and with a 10-fold enhancement of Fe relative to Solar abundance is required to explain the Fe K<sub>α</sub> emission weakly detected from GRB 970508 with Beppo-SAX (Piro et al. 1999).
Guided by the optical observations of $`\eta `$ Carinae, we assume that the volume-averaged density of gas at $`d\stackrel{<}{}10^{16}`$ cm of a GRB source is $`n=10^6n_6`$ cm<sup>-3</sup>. Dense clouds of radius $`10^{15}r_{15}`$ cm and radial Thomson depths $`\tau _T=n_c\sigma _Tr`$ are assumed to be embedded within this region, so that the mean density of particles in a cloud is $`n_c=1.5\times 10^9\tau _T/r_{15}`$ cm<sup>-3</sup>. Thus $`\tau _T1`$ clouds located very close to a GRB source are consistent with the observations of dense blobs near $`\eta `$ Carinae. The deceleration length scale of a blast wave with initial Lorentz factor $`\mathrm{\Gamma }_0=100\mathrm{\Gamma }_2`$ in a uniform medium is $`x_d=(3E_0/4\pi \mathrm{\Gamma }_0^2nm_pc^2)^{1/3}=2.5\times 10^{15}(E_{54}/\mathrm{\Gamma }_2^2n_6)^{1/3}`$ cm; hence the blast wave would emit a significant fraction of its energy before reaching distances of $`10^{16}`$ cm. The deceleration time, which corresponds to the duration of the prompt $`\gamma `$-ray luminous phase of a GRB in the external shock model (Rees & Mészáros 1992), is $`t_d=(1+z)x_d/(c\mathrm{\Gamma }_0^2)8(1+z)(E_{54}/\mathrm{\Gamma }_2^8n_6)^{1/3}`$ s. These parameters are not unique, and we expect that GRBs display a wide range of energies, Lorentz factors and surrounding mean densities that could accommodate the diverse range of GRB observations.
## 3 Blast-Wave/Cloud Interaction
A wave of photons impinging on a cloud located $`10^{16}d_{16}`$ cm from the explosion center will photoionize and Compton-scatter the ambient electrons to energies characteristic of the incident $`\gamma `$ rays (Madau & Thompson 1999). The $`\gamma `$-ray photon front has a width of $``$10-100 lt-s, corresponding to the duration of the GRB, whereas the plasma cloud has a width of $`3\times 10^4r_{15}`$ lt-s, so that the radiation effects must be treated locally. The radiation force driving the electrons outward is balanced by strong electrostatic forces from the more massive protons and ions that anchor the system until the net impulse is sufficient to drive the entire plasma cloud outward. The Compton back-scattered photons provide targets for successive waves of incident GRB photons through $`\gamma \gamma `$ pair-production interactions (Thompson & Madau 1999). Higher-energy photons are preferentially attenuated, forming an additional injection source of $`\stackrel{>}{}1`$ MeV electron-positron pairs. The nonthermal electrons and pairs will Compton scatter successive waves of photons, thereby modifying the incident spectrum. The pairs, no longer bound by electrostatic attraction with the ions, will be driven outward by both radiation forces and restoring electrostatic fields to form a mildly relativistic pair wind passing through the more slowly moving normal plasma. Shortly after the $`\gamma `$-ray photon front has passed, the decelerating blast wave from the GRB will plow into the cloud, shock-heating the relativistic plasma.
Nonthermal synchrotron photons with energy $`ϵ`$ impinge on the atoms in the cloud with a flux which can be parametrized as
$$\mathrm{\Phi }(ϵ)=(4\pi d^2)^1\frac{L}{m_ec^2ϵ_0^2\zeta _1}\left[\frac{1}{(ϵ/ϵ_0)^{2/3}+(ϵ/ϵ_0)^{\alpha _\gamma }}\right]$$
(1)
(Dermer, Chiang, & Böttcher 1999), where $`ϵ_01`$ is the photon energy of the peak of the $`\nu F_\nu `$ spectrum, $`\alpha _\gamma 2`$-3 is the photon spectral index at energies $`ϵϵ_0`$, and $`\zeta _1[3/4+(\alpha _\gamma 2)^1]`$. Hydrogen, the most abundant species in the cloud, will be ionized on a time scale of $`5\times 10^5(1+z)d_{16}^2\zeta _1ϵ_0^{4/3}/L_{50}`$ s. Fe features might persist briefly during the early periods of very weak GRBs on a time scale of $`4\times 10^3(1+z)d_{16}^2\zeta _1ϵ_0^{4/3}/L_{50}`$ s, and would be identified by a rapidly evolving Fe absorption feature at $`9.1/(1+z)`$ keV. After the H and Fe are ionized, the coupling between the GRB photons and gas is dominated by Compton scattering interactions. Pair production through photon-particle processes are negligible by comparison with Compton interactions except for photons with $`ϵ\stackrel{>}{}200`$. A lower limit to the time scale for an electron to be scattered by a photon is $`t_T(s)15(1+z)d_{16}^2\zeta _1ϵ_0/\zeta _2L_{50}`$, where $`\zeta _2=[3+(\alpha _\gamma 1)^1]`$, assuming that all Compton scattering events occur in the Thomson limit. The Klein-Nishina decline in the Compton cross section will increase this estimate by a factor of $`1`$-3, depending on the incident spectrum. Most of the electrons in the cloud will therefore be scattered to high energies during a very luminous ($`L_{50}1`$) GRB, or when the cloud is located at $`d_{16}1`$.
The average energy transferred to an electron at rest when Compton-scattered by a photon with energy $`ϵ`$ is $`\mathrm{\Delta }ϵϵ^2/(1+1.5ϵ)`$ (this expression is accurate to better than 18% for $`ϵ<10^3`$). Defining $`\eta =\gamma 1`$ as the dimensionless electron kinetic energy, we can easily estimate the production rate $`f(\eta )`$ of electrons scattered to energy $`\eta `$ in the nonrelativistic ($`\eta 1`$) and extreme relativistic ($`\eta 1`$) limits, noting that $`f(\eta )d\eta \mathrm{\Phi }(ϵ)\sigma _C(ϵ)dϵ`$ and letting $`\eta \mathrm{\Delta }ϵ`$. In the former limit, the Compton cross section $`\sigma _C(ϵ)\sigma _T`$ and $`\mathrm{\Phi }(ϵ)ϵ^{2/3}`$, so that $`f(\eta )\eta ^{5/6}`$ when $`\eta `$ min(1, $`ϵ_0^2`$). In the high energy limit, $`\sigma _C(ϵ)\mathrm{ln}(3.3ϵ)/ϵ`$ and $`\mathrm{\Phi }(ϵ)ϵ^{\alpha _\gamma }`$, so that $`f(\eta )\eta ^{(\alpha _\gamma +1)}\mathrm{ln}(2.2\eta )`$ when $`\eta `$ max$`(1,ϵ_0^2`$). Thus electrons are Compton-scattered on the time scale derived above to form a hard spectrum that turns over at kinetic energies of $`\stackrel{>}{}500\times `$min($`1,ϵ_0^2`$) keV. For a GRB with $`ϵ_01`$, most of the kinetic energy is therefore carried by nonthermal electrons with energies of $`500`$ keV. Successive waves of photons that pass through this plasma will continue to Compton-scatter the nonthermal electrons. Only the lowest energy photons, however, will be strongly affected by the radiative transfer because both the Compton scattering cross section and energy change per scattering is largest for the lowest energy photons.
Following the initial wave of photons, successive photon fronts also encounter the back-scattered radiation (Madau & Thompson 1999; Thompson & Madau 1999). The kinematics of the Compton process dictate that the energy $`ϵ_s`$ of a photon back-scattered through $`180^{}`$ by an electron at rest is $`ϵ_s=ϵ/(1+2ϵ)`$; thus $`ϵ_s`$ cannot exceed 1/2 the electron rest-mass energy. Head-on collisions of the back-scattered photons by primary GRB photons with $`ϵ_1>2/ϵ_s=2+2\sqrt{3}`$ can thus produce nonthermal e<sup>+</sup>-$`e^{}`$ pairs. The cross section for $`\gamma \gamma `$ pair production peaks near threshold with a value of $`\sigma _T/3`$. The $`\gamma \gamma `$ pair-production optical depth $`\tau _{\gamma \gamma }(ϵ_1)`$ of a photon that trails the onset of the GRB by $`\mathrm{\Delta }t`$ seconds can be estimated by noting that the photon traverses a distance $`r`$ through a backscattered radiation field of spectral density $`n_s(ϵ_s)n_e\sigma _T\mathrm{\Delta }t\mathrm{\Phi }[ϵ_s/(12ϵ_s)]`$ – a more accurate calculation would replace the term $`\mathrm{\Delta }t\mathrm{\Phi }[ϵ_s/(12ϵ_s)]`$ by an integral over the time-varying flux. Approximating $`\tau _{\gamma \gamma }(ϵ_1)r(\sigma _T/3)\mathrm{\Delta }ϵ_sn_s(2/ϵ_1)`$, where $`\mathrm{\Delta }ϵ_s2/ϵ_1`$ is the bandwidth that is effective for producing pairs, we obtain
$$\tau _{\gamma \gamma }(ϵ_1)0.02\frac{\tau _T\mathrm{\Delta }t[s]L_{50}k(ϵ_1)}{d_{16}^2ϵ_1ϵ_0^2\zeta _1}\left[\frac{1}{(ϵ^{}/ϵ_0)^{2/3}+(ϵ^{}/ϵ_0)^{\alpha _\gamma }}\right],$$
(2)
where $`ϵ^{}=2/(ϵ_14)`$. The coefficient results from a more detailed derivation, and the term $`k(ϵ_1)=14ϵ_1^1+ϵ_1/(ϵ_14)`$ is a Klein-Nishina correction. Eq. (2) shows that photons with energies above several MeV will be severely attenuated in Thomson thick clouds if $`L_{50}1`$ or $`d_{16}1`$. Photons with MeV energies are most severely attenuated, and $`\tau _{\gamma \gamma }(ϵ_1)ϵ_1^{1/3}`$ at energies $`ϵ_1\mathrm{max}(1,ϵ_0)`$. The $`\gamma \gamma `$ pair injection process provides another source of nonthermal leptons with $`\eta 1`$. The pairs will not, however, be electrostatically bound but will be accelerated by the photon pressure and electrostatic field.
Fig. 1 shows Monte Carlo simulations of radiation spectra described by eq. (1) that pass through a hot electron scattering medium. For simplicity, we approximate the hard nonthermal electron spectrum by a thermal distribution with temperatures of 100 and 300 keV. These calculations show that the lowest energy photons of the primary synchrotron spectrum are most strongly scattered, and that the “line-of-death” problem of the synchrotron shock model of GRBs (Preece et al. 1998) can be solved by radiation transfer effects through a hot scattering cloud with $`\tau _T\stackrel{>}{}`$ 1-2.
According to this interpretation, GRBs displaying very hard spectra could display one break from the intrinsic synchrotron shock spectrum and a second break from the scattering process. In the examples shown in Fig. 1, these two breaks are so close to each other that they appear as one smooth turnover. Two breaks are observed from GRB 970111 (Crider & Liang 1999), a GRB that strongly violates the “line of death.” A prediction of this model is that GRBs showing such flat X-ray spectra should also display softer MeV spectra than typical GRBs due to $`\gamma \gamma `$ attenuation processes in the hot scattering cloud.
## 4 Observational Signatures of the Flash-Heated Cloud
The Compton-scattered electrons transfer momentum to the $`N_p=4\times 10^{54}r_{15}^3n_9`$ protons of the cloud through their electrostatic coupling. If the radiation efficiency is $`\xi _r`$, then the Compton impulse gives each proton in the cloud $`\xi _r[1\mathrm{exp}(\tau _T)]E\pi r^2/(4\pi d^2N_p)40[1\mathrm{exp}(\tau _T)](\xi _r/0.1)E_{54}/(d_{16}^2r_{15}n_9)`$ MeV of directed energy. Pairs, by contrast, will be accelerated to mildly relativistic speeds until Compton drag or streaming instabilities limit further acceleration. In the simplification that the medium interior to the cloud is uniform, and neglecting pair-loading of the swept-up material (Thompson & Madau 1999), the decelerating blast wave follows the dynamical equation $`\mathrm{\Gamma }(x)=\mathrm{\Gamma }_0(x/x_d)^g`$, where $`g=3/2`$ and 3 for adiabatic and radiative blast waves, respectively. Using the standard parameters adopted here, the blast wave slows to between $`\xi =0.010.1`$ of its initial Lorentz factor before reaching a cloud at $`d=10^{16}`$ cm. Even considering the radiative acceleration of the cloud, the blast wave reaches the cloud at time $`t_{bw}=t_d(d/x_d)^{(2g+1)}/(2g+1)\stackrel{<}{}t_{\mathrm{dyn}}`$, where the dynamical time scale of the cloud is $`t_{\mathrm{dyn}}=r/c=3\times 10^4r_{15}`$ s. Because the cloud is so dense, a large fraction of the residual energy of the blast wave is deposited into the $`N_p`$ particles of the cloud. Thus each proton in the cloud receives an additional $`m_p\beta _p^2c^2\xi E\pi r^2/(4\pi d^2N_p)40(\xi /0.1)E_{54}/(d_{16}^2r_{15}n_9`$) MeV of kinetic energy, divided roughly equally into directed outflow and random thermal energy.
If the circumstellar medium at $`d_{16}1`$ is much more dilute than the interior region, as suggested by observations of $`\eta `$ Carinae, then we can neglect further interactions of the cloud/blast wave system with their surroundings. The observational signatures and fate of the cloud at late times can be outlined by comparing time scales. The cloud expands on a time scale $`t_{\mathrm{ex}}=t_{\mathrm{dyn}}/\beta =1.5\times 10^5r_{15}/(\beta /0.2)`$ s. The basic time scale governing radiative processes in the cloud is the Thomson time $`t_\mathrm{T}=(n_e\sigma _Tc)^1=5\times 10^4/n_9`$ s $`=t_{\mathrm{dyn}}/\tau _\mathrm{T}`$. The electrons thermalize on a time scale $`t_\mathrm{T}/\mathrm{ln}\mathrm{\Lambda }t_{\mathrm{dyn}}`$, where the Coulomb logarithm $`\mathrm{ln}\mathrm{\Lambda }20`$. The protons transfer their energy to the electrons on the time scale $`t_{ep}\theta ^{3/2}(m_p/m_e)t_\mathrm{T}/\mathrm{ln}\mathrm{\Lambda }`$, where $`\theta =kT/m_ec^2`$ is an effective dimensionless electron temperature, and we assume collective plasma processes for energy exchange to be negligible.
The flash-heated cloud can evolve in two limiting regimes. When the external soft-photon energy density is small, the system emits by a hard bremsstrahlung spectrum with $`\theta `$ 1 and luminosity $`L_{ff}N_e\alpha _fm_ec^2\theta ^{1/2}/t_\mathrm{T}5\times 10^{41}r_{15}^3n_9^2\theta ^{1/2}`$ ergs s<sup>-1</sup>. In the more likely case when abundant soft photons are present, for example, from the Compton echo (Madau, Blandford, & Rees 1999), then Compton cooling will balance ion heating to produce a luminous Comptonized soft-photon spectrum with effective temperature $`\theta 0.1`$ and luminosity $`L_\mathrm{C}\xi E\pi r^2/(4\pi d^2t_{ep})2\times 10^{45}(\xi /0.1)E_{54}r_{15}^2n_9d_{16}^2(\theta /0.1)^{3/2}`$ ergs s<sup>-1</sup>. In either case, the spectra persist until the plasma expands and adiabatically cools, that is, for a period $`t_{ex}`$ day. The hot bremsstrahlung plasma will be too dim to be detectable with current instrumentation, but a 50-100 keV Comptonized plasma at redshift $`z1`$ and luminosity distance of $`10^{28}D_{28}`$ cm would have a flux of $`1.6\times 10^{12}(\xi /0.1)E_{54}r_{15}^2n_9d_{16}^2(\theta /0.1)^{3/2}D_{28}^2`$ ergs cm<sup>-2</sup> s<sup>-1</sup>. The hot plasma formed by a nearby GRB at $`z0.1`$ would be easily detectable with the INTEGRAL and Swift missions at hard X-ray and soft $`\gamma `$-ray energies. In either case, e<sup>+</sup>-e<sup>-</sup> pairs would be formed with moderate efficiency, and the cooling, expanding plasma would produce a broad pair annihilation feature (Guilbert & Stepney 1985). The residual pairs formed in the relativistic plasma and the pair wind would diffuse into the dilute interstellar medium with density $`n_{\mathrm{ISM}}`$ to annihilate on a time scale $`(n_{\mathrm{ISM}}\sigma _\mathrm{T}c)^12\times 10^6/n_{\mathrm{ISM}}`$ yr. If the energy intercepted by a single cloud is converted to pairs with a conservative 1% pair yield, past GRBs in the Milky Way would be revealed by localized regions of annihilation radiation with flux $`2\times 10^5E_{54}n_{\mathrm{ISM}}(d/10\mathrm{k}\mathrm{p}\mathrm{c})^2`$ 0.511 MeV ph cm<sup>-2</sup> s<sup>-1</sup>. The high-latitude annihilation feature discovered with OSSE on the Compton Gamma Ray Observatory (Purcell et al. 1997), or other localized hot spots of annihilation radiation that will be mapped in detail with INTEGRAL, could reveal sites of past GRB explosions.
CD thanks B. Paczyński for stressing the importance of massive star observations in developing blast-wave models of GRBs. The work of CD is supported by the Office of Naval Research and NASA Astrophysical Theory Program (DPR S-13756G). The work of MB is supported by NASA through Chandra Postdoctoral Fellowship grant PF 9-10007, awarded by the Chandra X-ray Center, which is operated by the Smithsonian Astrophysical Observatory for NASA under contract NAS 8-39073. |
warning/0002/math0002129.html | ar5iv | text | # Hereditary indecomposability and the Intermediate Value Theorem
## Introduction
The classical Intermediate Value Theorem (IVT for short) states that if $`f`$ is a continuous function from the interval $`[a,b]`$ to $`\mathrm{}`$ with $`f(a)f(b)<0`$ then there is $`c`$ in $`(a,b)`$ such that $`f(c)=0`$. In Henriksen, Larson and Martinez investigated forms of this theorem in lattice-ordered rings, where, because of the absence of any natural topology, they restricted their attention to polynomials. We mention some of their results for the ring $`C^{}(X)`$ of bounded real-valued continuous functions on the topological space $`X`$; let us call $`X`$ an IVT-space if the ring $`C^{}(X)`$ satisfies the Intermediate Value Theorem (the precise formulation of the IVT in this context follows below). The results are:
1. every IVT-space is an $`F`$-space;
2. every compact and zero-dimensional $`F`$-space is an IVT-space;
3. every compact IVT-space is hereditarily indecomposable.
In this note we establish a partial converse to this last result in that we show that every compact hereditarily indecomposable space satisfies the IVT for a restricted class of polynomials.
## 1. Preliminaries
We shall only deal with rings of the form $`C^{}(X)`$, so we can, for the time being restrict our attention to compact Hausdorff spaces.
### 1.1. The intermediate value theorem
In the ring $`C^{}(X)`$ the IVT takes on the following form: Let $`p`$ be a polynomial with coefficients in $`C^{}(X)`$ and let $`u`$ and $`v`$ be elements of $`C^{}(X)`$ such that $`p(u)\mathrm{𝟎}p(v)`$, where $`\mathrm{𝟎}`$ denotes the zero function. Then there is $`wC^{}(X)`$ such that $`uvwuv`$ and $`p(w)=\mathrm{𝟎}`$. The reason for working with $`uv`$ and $`uv`$ is of course that it is usually not the case that $`u(x)v(x)`$ for all $`x`$ (or $`v(x)u(x)`$ for all $`x`$).
To get some feeling for what the IVT says in this context let $`pC^{}(X)[t]`$, so $`p(t)=_{i=0}^nf_it^i`$ for some elements $`f_0`$, …, $`f_n`$ of $`C^{}(X)`$, and let $`u,vC^{}(X)`$ be such that $`p(u)\mathrm{𝟎}p(v)`$. For every separate $`xX`$ we get an ordinary polynomial $`p_x(t)=_{i=0}^nf_i(x)t^i`$; and the assumptions on $`u`$ and $`v`$ imply that $`p\left(u(x)\right)0p\left(v(x)\right)`$. The classical IVT therefore guarantees that there is a function $`w:X\mathrm{}`$ such that $`uvwuv`$ and $`p(w)=\mathrm{𝟎}`$; the IVT for $`C^{}(X)`$ demands that this $`w`$ be continuous.
That this puts severe restrictions on the space $`X`$ may be seen as follows: let $`fC^{}(X)`$ and consider the polynomial $`p(t)=|f|tf`$. Now $`p(\mathrm{𝟏})=|f|f0`$ and $`p(\mathrm{𝟏})=|f|f\mathrm{𝟎}`$, so if $`X`$ is an IVT-space there must be a continuous function $`w`$ such that $`\mathrm{𝟏}w\mathrm{𝟏}`$ and $`f=w|f|`$. This however is one of the characterizations of $`F`$-spaces — see Gillman and Jerison .
### 1.2. Hereditarily indecomposable spaces
Much of what follows is taken from Oversteegen and Tymchatyn , which is a convenient survey on hereditarily indecomposable spaces.
To begin we recall that a continuum is said to be *indecomposable* if it cannot be written as the union of two proper subcontinua; it is *hereditarily indecomposable* if every subcontinuum is indecomposable.
We use the following characterization of hereditarily indecomposable continua.
###### Theorem 1.1.
A continuum $`X`$ is hereditarily indecomposable if and only if whenever two disjoint closed sets $`A`$ and $`B`$ and open neighbourhoods $`U`$ and $`V`$ respectively are given we can write $`X`$ as the union of three closed sets $`X_0`$, $`X_1`$ and $`X_2`$ such that $`AX_0`$, $`BX_2`$, $`X_0X_1V`$, $`X_0X_2=`$, and $`X_1X_2U`$.
The property in this theorem can also be used to characterize those compact spaces (connected or not) for which every closed connected subspace is indecomposable; we shall call these compact spaces hereditarily indecomposable as well. Observe that with this definition compact zero-dimensional spaces are hereditarily indecomposable as well.
## 2. The IVT implies hereditary indecomposability
In this section we reprove Theorem 3.2 from Henriksen, Larson and Martinez , which states that compact IVT-spaces are hereditarily indecomposable. In their proof these authors used a polynomial of degree $`7`$ with two potentially irreducible quadratic factors. We use a completely factored polynomial of degree $`3`$.
###### Theorem 2.1.
Compact IVT-spaces are hereditarily indecomposable.
###### Proof.
Let $`X`$ be a compact IVT-space. To show that $`X`$ is hereditarily indecomposable we take disjoint closed sets $`A`$ and $`B`$ and open sets $`U`$ and $`V`$ such that $`AU`$ and $`BV`$. We must exhibit three closed sets $`X_0`$, $`X_1`$ and $`X_2`$ such that $`AX_0`$, $`BX_2`$, $`X_0X_1V`$, $`X_0X_2=`$, $`X_1X_2U`$ and $`X_0X_1X_2=X`$.
Choose a continuous function $`f:X[0,1]`$ such that $`fA0`$, $`fB1`$, $`f^1\left[[0,\frac{1}{2})\right]U`$ and $`f^1\left[(\frac{1}{2},1]\right]V`$. Using $`f`$ we define three continuous functions, $`f_1`$, $`f_2`$ and $`f_3`$, as follows: first
$$f_1(x)=\{\begin{array}{cc}f(x)\frac{1}{4}\hfill & \text{if }f(x)\frac{1}{4}\hfill \\ 0\hfill & \text{if }\frac{1}{4}f(x)\frac{3}{4}\hfill \\ f(x)\frac{3}{4}\hfill & \text{if }\frac{3}{4}f(x)\text{;}\hfill \end{array}$$
second
$$f_2(x)=\{\begin{array}{cc}\frac{1}{2}(f(x)\frac{1}{4})\hfill & \text{if }f(x)\frac{1}{4}\hfill \\ 2(f(x)\frac{1}{4})\hfill & \text{if }\frac{1}{4}f(x)\frac{3}{4}\hfill \\ \frac{1}{2}(f(x)+\frac{5}{4})\hfill & \text{if }\frac{3}{4}f(x)\text{;}\hfill \end{array}$$
and third
$$f_3(x)=\{\begin{array}{cc}f(x)+\frac{3}{4}\hfill & \text{if }f(x)\frac{1}{4}\hfill \\ 1\hfill & \text{if }\frac{1}{4}f(x)\frac{3}{4}\hfill \\ f(x)+\frac{1}{4}\hfill & \text{if }\frac{3}{4}f(x)\text{;}\hfill \end{array}$$
(At this point the reader may find it instructive to draw the graphs of $`f_1`$, $`f_2`$ and $`f_3`$ in case $`X=[0,1]`$ and $`f(x)=x`$. The zig-zag that appears when one follows the graph of $`f_3`$ left-to-right until it meets the graph of $`f_2`$ then follows the graph of $`f_2`$ right-to-left until it meets the graph of $`f_1`$ and finally the graph of $`f_1`$ left-to-right until the end is characteristic of hereditarily indecomposable spaces.) Note that $`f_1f_2f_3`$.
Consider the polynomial $`p`$ defined by $`p(t)=(tf_1)(tf_2)(tf_3)`$. Then $`p(\mathrm{𝟎})\mathrm{𝟎}p(\mathrm{𝟏})`$, for one readily checks that
* $`f_2(x)<0<f_3(x)<1`$ if $`f(x)<\frac{1}{4}`$;
* $`f_1(x)=0`$ and $`f_3(x)=1`$ if $`\frac{1}{4}f(x)\frac{3}{4}`$ and
* $`0<f_1(x)<1<f_2(x)`$ if $`f(x)>\frac{3}{4}`$.
An application of the Intermediate Value Theorem gives us a continuous function $`w:X[0,1]`$ such that $`p(w)=\mathrm{𝟎}`$.
Let $`X_0=\{x:w(x)=f_3(x)\}`$, $`X_1=\{x:w(x)=f_2(x)\}`$ and $`X_2=\{x:w(x)=f_1(x)\}`$. We check that these sets have all the required properties.
* $`X_0X_1X_2=X`$ because $`p(w)=\mathrm{𝟎}`$;
* $`AX_0`$ because if $`xA`$ then $`f(x)=0`$, hence $`w(x)=f_3(x)`$;
* $`BX_2`$ because if $`xB`$ then $`f(x)=1`$, hence $`w(x)=f_1(x)`$;
* $`X_0X_1V`$ because if $`xX_0X_1`$ then $`f_3(x)=w(x)=f_2(x)`$ hence $`f(x)=\frac{3}{4}`$ and $`xV`$;
* $`X_1X_2U`$ because if $`xX_1X_2`$ then $`f_1(x)=w(x)=f_2(x)`$ hence $`f(x)=\frac{1}{4}`$ and $`xU`$ and
* $`X_0X_2=`$ because $`f_3f_1=\mathrm{𝟏}`$.
We conclude that $`X`$ is indeed hereditarily indecomposable. ∎
As announced before, in the next section we shall see that hereditary indecomposability is in fact characterized by the particular instance of the Intermediate Value Theorem that was actually employed.
## 3. Hereditary indecomposability implies part of the IVT
In this section we show that every compact hereditarily indecomposable $`F`$-space $`X`$ satisfies the Intermediate Value Theorem for *completely factored polynomials*, that is, polynomials that can be written as $`_{i=1}^n(tf_i)`$, where the $`f_i`$ are elements of $`C(X)`$.
This is a rather limited class of polynomials of course but, as we saw in Section 2, the case $`n=3`$ is already strong enough to imply hereditary indecomposability. The Intermediate Value Theorem for this class of polynomials therefore characterizes hereditary indecomposability for $`F`$-spaces.
So let $`X`$ be a hereditarily indecomposable $`F`$-space and let $`p`$, defined by $`p(t)=_{i=1}^n(tf_i)`$, be a completely factored polynomial in $`C(X)`$. Assume furthermore that $`u,vC(X)`$ are such that $`p(u)\mathrm{𝟎}p(v)`$. Through a series of reductions we show that there is $`wC(X)`$ such that $`p(w)=\mathrm{𝟎}`$ and $`uvwuv`$.
###### Lemma 3.1.
We may assume that $`f_1f_2\mathrm{}f_n`$.
###### Proof.
For each $`in`$ define $`g_i`$ by
$$g_i=\underset{|F|=i}{}\underset{jF}{}f_j.$$
Observe that $`g_1g_2\mathrm{}g_n`$ and that, for each individual $`x`$, the sets of values $`\{g_1(x),g_2(x),\mathrm{},g_n(x)\}`$ and $`\{f_1(x),f_2(x),\mathrm{},f_n(x)\}`$ are equal. It follows from this that the coefficients of $`t^0`$, $`t^1`$, …, $`t^{n1}`$ in $`_{i=1}^n(tf_i)`$ and $`_{i=1}^n(tg_i)`$ are the same and hence that the polynomials are the same. ∎
The case $`n=1`$ should offer no problems and the case $`n=2`$ is dealt with in the following proposition, which is a special case of Theorem 2.3 (b) of Henriksen, Larson and Martinez . In fact, the polynomial $`p`$ need not even be factored; it can always be factored by completing the square.
###### Proposition 3.2.
Every space satisfies the Intermediate Value Theorem for monic quadratic polynomials.
###### Proof.
Let $`p(t)=t^2+2ft+g`$ be such a polynomial and assume that there are $`u`$ and $`v`$ such that $`p(u)\mathrm{𝟎}p(v)`$. Completing the square gives us $`q(t)=(t+f)^2+gf^2`$. Now because $`p(u)\mathrm{𝟎}p(v)`$ we know that $`f^2g\mathrm{𝟎}`$ so that we can write $`f^2g=h^2`$ for some nonnegative $`hC(X)`$. We find that $`p(t)=(t+fh)(t+f+h)`$; write $`fh=f_1`$ and $`f+h=f_2`$.
Observe that for each $`x`$ either $`v(x)f_1(x)`$ or $`v(x)f_2(x)`$ and that $`f_1(x)u(x)f_2(x)`$. We cover our space by three closed sets: $`P=\mathrm{cl}\left\{x:u(x)<v(x)\right\}`$, $`Q=\{x:u(x)=v(x)\}`$ and $`R=\mathrm{cl}\left\{x:u(x)>v(x)\right\}`$. We now note that $`uf_2v`$ on $`P`$ (because $`u(x)f_2(x)v(x)`$ whenever $`u(x)<v(x)`$) and that $`vf_1u`$ on $`R`$. We define $`w`$ as the combination
$$(f_2P)(uQ)(f_1R).$$
Note that $`w`$ is well-defined because, by continuity, $`uvf_2`$ on $`PQ`$ and $`uvf_1`$ on $`QR`$. Also $`p(w)(x)=0`$ for all $`x`$; this is clear on $`PR`$ and on $`Q`$ it holds because $`p(u)(x)0p(v)(x)`$ and $`p(u)(x)=p(v)(x)`$. Finally, $`w`$ is continuous because it is the combination of continuous functions defined on closed subsets. ∎
We have given such an extensive proof of Proposition 3.2 because it contains elements that we will use quite often in what follows, to wit breaking the space into closed pieces according to the position of the $`f_i(x)`$ with respect to $`u(x)`$ and $`v(x)`$, and defining $`w`$ by cases. From now on we assume that $`n3`$.
To begin, for every $`x`$ we have $`f_n(x)u(x)f_{n1}(x)`$ or $`f_{n2}(x)u(x)f_{n3}(x)`$ etc., *because* $`p(u)(x)0`$; if, for example, $`f_{n1}(x)>u(x)>f_{n2}(x)`$ then clearly $`p(u)(x)>0`$. This sequence ends with $`f_2(x)u(x)f_1(x)`$ if $`n`$ is even and with $`f_1(x)u(x)`$ if $`n`$ is odd.
Likewise, for all $`x`$ we have $`v(x)f_n(x)`$ or $`f_{n1}(x)v(x)f_{n2}(x)`$ or … or $`f_1(x)v(x)`$ if $`n`$ is even and $`f_2(x)v(x)f_1(x)`$ if $`n`$ is odd.
We shall also employ the cover of $`X`$ by the sets $`P=\mathrm{cl}\left\{x:u(x)<v(x)\right\}`$, $`Q=\{x:u(x)=v(x)\}`$ and $`R=\mathrm{cl}\left\{x:u(x)>v(x)\right\}`$. On $`Q`$ there is no choice: the only admissible solution is $`w_Q=uQ=vQ`$. However, once we have found solutions $`w_P`$ on $`P`$ and $`w_R`$ on $`R`$ then $`w=w_Pw_Qw_R`$ is the desired solution. On $`P`$ we have $`uw_Pv`$ so by continuity we know that $`u(x)=w_P(x)=v(x)`$ for all $`xPQ`$. Likewise $`u(x)=w_R(x)=v(x)`$ for all $`xQR`$. Thus, $`w`$ is well-defined and as a combination of continuous functions defined on closed subsets it is continuous.
Because hereditary indecomposability is a closed hereditary property we can work inside $`P`$ and $`R`$ respectively without worrying about the rest of $`X`$.
### 3.1. Reduction to odd $`n`$
Assume $`n`$ is even and recall that in this case $`f_1uf_n`$.
We show that on $`P`$ we have $`q(u)0q(v)`$, where $`q(t)=_{i=2}^n(tf_i)`$. Indeed the possible positions of $`u(x)`$ ensure that $`q(u)(x)0`$ for all $`x`$. Also, for all $`x`$ with $`u(x)<v(x)`$ we have $`v(x)f_2(x)`$ because $`f_1(x)u(x)<v(x)<f_2(x)`$ would imply $`p(v)(x)<0`$. Hence, by continuity, $`vf_2`$ on $`P`$, so that $`q(v)\mathrm{𝟎}`$ on $`P`$.
On the set $`R`$ we can show in a similar fashion that $`vf_{n1}`$ and hence that $`r(u)\mathrm{𝟎}r(v)`$, where $`r(t)=_{i=1}^{n1}(tf_i)`$.
Both $`q`$ and $`r`$ are of degree $`n1`$.
From now on we assume $`n3`$ and $`n`$ odd.
### 3.2. Reduction to $`uv`$
If $`u(x)>v(x)`$ then, because $`f_nu`$, we must have $`v(x)f_{n1}(x)`$ and because $`vf_1`$ we must have $`u(x)f_2(x)`$. So on $`R`$ we get, by continuity, $`vf_{n1}`$ and $`uf_2`$. Consider now the polynomial $`q(t)=_{i=2}^{n1}(tf_i)`$. Because of the possible positions for $`u(x)`$ and $`v(x)`$ listed above we conclude that $`q(u)\mathrm{𝟎}q(v)`$.
### 3.3. The final case
We now show how to produce $`w`$, given that 1) $`X`$ is the closure of $`\{x:u(x)<v(x)\}`$, 2) $`p(u)p(v)`$ and 3) $`n`$ is odd.
Let $`k`$ be such that $`n=2k+1`$. For each $`ik`$ consider the closed sets $`A_i=\mathrm{cl}\left\{x:v(x)<f_{2i+1}(x)\right\}`$ and $`B_i=\mathrm{cl}\left\{x:u(x)>f_{2i1}(x)\right\}`$.
Note that, because of the positioning of the values $`u(x)`$ and $`v(x)`$ we have $`A_iC_i=\{x:v(x)f_{2i}(x)\}`$ and $`B_iD_i=\mathrm{cl}\left\{x:u(x)f_{2i}(x)\right\}`$. Now note that $`C_iD_i\{x:u(x)=f_{2i}(x)=v(x)\}`$; as the set on the right-hand side is nowhere dense it follows that $`\mathrm{int}C_i`$ and $`\mathrm{int}D_i`$ are disjoint.
Also, because $`X`$ is an $`F`$-space, we know that $`A_i\mathrm{int}C_i`$ and $`B_i\mathrm{int}D_i`$.
Now apply hereditary indecomposability to find three closed sets $`X_i`$, $`Y_i`$ and $`Z_i`$ that cover $`X`$ and with the following properties: $`A_iX_i`$, $`B_iZ_i`$, $`X_iY_i\mathrm{int}D_i`$, $`Y_iZ_i\mathrm{int}C_i`$ and $`X_iZ_i=`$. We note the following facts:
1. $`uf_{2i1}`$ on $`X_iY_i`$ because this set is disjoint from $`B_i`$;
2. $`vf_{2i+1}`$ on $`Y_iZ_i`$ because this set is disjoint from $`A_i`$;
3. $`u=f_{2i1}=f_{2i}`$ on $`X_iY_i`$ because this set is contained in $`D_i`$ and because of (1);
4. $`v=f_{2i+1}=f_{2i}`$ on $`Y_iZ_i`$ because the set is contained in $`C_i`$ and because of (2); and
5. $`uf_{2i1}f_{2i}f_{2i+1}v`$ on $`Y_i`$ because of (1) and (2).
Now we are ready to define $`w`$. We start by letting $`w=f_1`$ on $`X_1`$ and $`w=f_2`$ on $`Y_1`$. We continue by letting, for $`i>1`$, $`w=f_{2i1}`$ on $`X_i_{j<i}Z_j`$ and $`w=f_{2i}`$ on $`Y_i_{j<i}Z_j`$. Finally, on $`_{ik}Z_i`$ we let $`w=f_n`$.
#### We check that $`w`$ is well-defined.
By (3) we have $`f_{2i1}=f_{2i}`$ on $`X_iY_i`$ for every $`i`$. If $`j<i`$ then $`X_jX_i_{l<i}Z_l=`$; on $`Y_jX_i_{l<i}Z_lY_jZ_jZ_{i1}`$ we have $`v=f_{2j+1}=f_{2j}`$ and $`vf_{2i1}`$ and so $`f_{2i1}=f_{2j}`$, and on $`Y_jY_i_{l<i}Z_lY_jZ_jZ_{i1}`$ we have $`v=f_{2j=1}=f_{2j}`$ and $`vf_{2i+1}f_{2i}`$ and so $`f_{2j}=f_{2i}`$. Finally, on $`Y_j_{ik}Z_iY_jZ_j`$ we have $`vf_n`$ and $`v=f_{2j}`$ so $`f_{2j}=f_n`$.
#### We check that $`uwv`$.
On $`X_1`$ we surely have $`uf_1v`$ and if $`i>1`$ then on $`X_i_{j<i}Z_j`$ we have $`uf_{2i1}`$ because of (1) and $`f_{2i1}v`$ because of (1) for $`i1`$. On each $`Y_i`$ we have $`uf_{2i}v`$ by (5). Finally, on $`_{ik}Z_i`$ we have $`vf_nu`$ by (2).
We see that $`w`$ is a well-defined continuous function on $`X`$ such that $`uwv`$ and, because for all $`x`$ there is an $`i`$ with $`w(x)=f_i(x)`$, such that $`p(w)=\mathrm{𝟎}`$.
## 4. Questions and Conjectures
The basic question as to what actually characterizes IVT-spaces remains. On the basis of the evidence from Section 3 we conjecture that hereditarily indecomposable spaces also satisfy the IVT for monic polynomials.
The general case seems more complicated in that the leading coefficient (and others) may vanish at certain places. It may very well be that that the full IVT characterizes zero-dimensionality. |
warning/0002/gr-qc0002064.html | ar5iv | text | # Symmetries of asymptotically flat axisymmetric spacetimes with null dust
## I Introduction and Summary
Recently a unique role of boost-rotation symmetric electrovacuum spacetimes describing “uniformly accelerated particles” of various kinds was exhibited by a theorem which states that in axially symmetric, asymptotically flat spacetimes the only additional symmetry that does not exclude radiation is the boost symmetry (in Ref. for vacuum spacetimes with hypersurface orthogonal Killing vectors and in Ref. for electrovacuum spacetimes with Killing vectors which are in general not hypersurface orthogonal). Our effort in this paper is to prove a similar theorem for asymptotically flat spacetimes with null dust fields. We also specialize the spacetime to be axially symmetric (with the axial Killing vector which is in general not hypersurface orthogonal) - this assumption simplifies lengthy calculations.
If one is interested in gravitational radiation from a general bounded matter source, i.e., in the behaviour of gravitational field far from the source, one has to turn to approximation methods. The Bondi-Sachs formalism is a powerful instrument for the treatment not only of asymptotically flat vacuum and electrovacuum fields, but also for the investigation of asymptotically flat null dust fields . The energy-momentum tensor of null dust or pure radiation
$$T_{\alpha \beta }=\rho n_\alpha n_\beta ,n_\alpha n^\alpha =0,\rho >0,$$
(1)
with $`\rho `$ being the energy density of the radiation field, describes a field of massless radiation propagating along a null congruence with the tangent vector $`n^\alpha `$. This field is the incoherent superposition of waves with random phases and different polarizations where the radiation arises from electromagnetic null field, massless scalar field, neutrino field or gravitational field itself. The field equations of these originating fields are not considered. Exact solutions of this class are the Vaidya solution which can model the exterior of a spherically symmetric shining star and the Kinnersley photon rocket , a particle emitting photons and accelerating because of the recoil. Also null dust fields with rotation are known .
In Sec. II we start out from the general form of an axially symmetric metric in Bondi-Sachs coordinates $`\{u,r,\theta ,\varphi \}`$, where the null coordinate $`u`$ and the spherical angles $`\theta `$, $`\varphi `$ are constant along null rays while the luminosity distance $`r`$ varies. We consider their asymptotic series expansions at $`r\mathrm{}`$ assuming the Einstein equations for the null dust field to be satisfied and the spacetime to be asymptotically flat.
Then we assume that, in addition, another symmetry exists, i.e., an additional Killing vector field which forms with the axial one a two-dimensional Lie algebra. By first decomposing the additional Killing vector field $`\eta ^\alpha `$ in the null tetrad and then solving the Killing equations asymptotically in the leading order for this new Killing vector, we find that the additional Killing vector asymptotically generates either a supertranslation or a boost along the symmetry axis. However, developing and solving Killing equations in higher orders and considering Lie derivatives of the energy momentum tensor for null dust to vanish for both the supertranslational Killing vector in Sec. III and the boost Killing vector in Sec. IV we conclude that in fact the only allowable additional Killing vector of axially symmetric spacetimes with null dust is a supertranslational Killing vector and then the gravitational field is non-radiative (the Weyl tensor has a non-radiative character).
Our conventions for the Riemann and Ricci tensors follow those of Ref. but our signature is $`2`$.
## II Axisymmetric null dust spacetimes with another symmetry
Consider an axially symmetric spacetime with the corresponding Killing vector field denoted by $`/\varphi `$. Assume that at least the ”piece of $`𝒥^+`$” exists in the sense of Ref. . Then one can introduce the Bondi-Sachs coordinate system { $`u`$$`r`$$`\theta `$$`\varphi `$ } $``$ { $`x^0`$$`x^1`$$`x^2`$$`x^3`$} in which the metric has the form
$`\mathrm{d}s^2`$ $`=`$ $`\left({\displaystyle \frac{V}{r}}\mathrm{e}^{2\beta }r^2\mathrm{e}^{2\gamma }U^2\mathrm{cosh}2\delta r^2\mathrm{e}^{2\gamma }W^2\mathrm{cosh}2\delta 2r^2UW\mathrm{sinh}2\delta \right)\mathrm{d}u^2`$ (4)
$`+2\mathrm{e}^{2\beta }\mathrm{d}u\mathrm{d}r+2r^2(\mathrm{e}^{2\gamma }U\mathrm{cosh}2\delta +W\mathrm{sinh}2\delta )\mathrm{d}u\mathrm{d}\theta +2r^2(\mathrm{e}^{2\gamma }W\mathrm{cosh}2\delta +U\mathrm{sinh}2\delta )\mathrm{sin}\theta \mathrm{d}u\mathrm{d}\varphi `$
$`r^2\left[\mathrm{cosh}2\delta (\mathrm{e}^{2\gamma }\mathrm{d}\theta ^2+\mathrm{e}^{2\gamma }\mathrm{sin}^2\theta \mathrm{d}\varphi ^2)+2\mathrm{sinh}2\delta \mathrm{sin}\theta \mathrm{d}\theta \mathrm{d}\varphi \right],`$
where the six metric functions $`U`$, $`V`$, $`W`$, $`\beta `$, $`\gamma `$, $`\delta `$ do not depend on $`\varphi `$ because of axial symmetry. The spacetime (4) is filled with null dust described by the energy-momentum tensor (1) which is assumed to be axially symmetric, too.
In Ref. the expansions of the metric functions of (4) for asymptotically flat spacetimes with null dust in Bondi-Sachs coordinates are derived. For that, the null expansion vector $`n^\alpha `$ is chosen to be identified with the null vector $`k^\alpha `$ of the Bondi-Sachs tetrad (24) at $`𝒥^+`$, i.e., $`𝒥^+`$ has to exist in the direction determined by the null vector $`n^\alpha `$ of the null dust. Then its contravariant components are
$`n^u`$ $`=`$ $`{\displaystyle \frac{𝒰(u,\theta )}{r^2}}+𝒪(r^3),`$ (5)
$`n^r`$ $`=`$ $`1+{\displaystyle \frac{(u,\theta )}{r}}+𝒪(r^2),`$ (6)
$`n^\theta `$ $`=`$ $`{\displaystyle \frac{𝒯(u,\theta )}{r^2}}+𝒪(r^3),`$ (7)
$`n^\varphi `$ $`=`$ $`{\displaystyle \frac{(u,\theta )}{r^2}}+𝒪(r^3),`$ (8)
and the covariant components read
$`n_u`$ $`=`$ $`1+{\displaystyle \frac{}{r}}+𝒪(r^2),`$ (9)
$`n_r`$ $`=`$ $`{\displaystyle \frac{𝒰}{r^2}}+𝒪(r^3),`$ (10)
$`n_\theta `$ $`=`$ $`𝒯+𝒪(r^1),`$ (11)
$`n_\varphi `$ $`=`$ $`\mathrm{sin}^2\theta +𝒪(r^1).`$ (12)
As the vector $`n^\alpha `$ is null ($`n^\alpha n_\alpha =0`$), functions entering (7), (11) have to satisfy
$$𝒰={\scriptscriptstyle \frac{1}{2}}(𝒯^2+^2\mathrm{sin}^2\theta ),$$
(13)
and similarly for the higher-order coefficients. Then the equations for the null dust field (1) can be solved. The metric coefficients have in the first order in $`r^k`$ the same form as Eq. (4) in Ref. :
$`\gamma `$ $`=`$ $`{\displaystyle \frac{c}{r}}+𝒪(r^3),`$ (14)
$`\delta `$ $`=`$ $`{\displaystyle \frac{d}{r}}+𝒪(r^3),`$ (15)
$`\beta `$ $`=`$ $`{\displaystyle \frac{1}{4}}(c^2+d^2){\displaystyle \frac{1}{r^2}}+𝒪(r^4),`$ (16)
$`U`$ $`=`$ $`(c,_\theta +2c\mathrm{cot}\theta ){\displaystyle \frac{1}{r^2}}+𝒪(r^3),`$ (17)
$`W`$ $`=`$ $`(d,_\theta +2d\mathrm{cot}\theta ){\displaystyle \frac{1}{r^2}}+𝒪(r^3),`$ (18)
$`V`$ $`=`$ $`r2M+𝒪(r^1).`$ (19)
For the radiation density $`\rho (u,r,\theta )`$ we write
$$\rho (u,r,\theta )=\frac{\rho _2(u,\theta )}{r^2}+𝒪(r^3)$$
(20)
and from the field equations
$$M,_u=(c,_u^2+d,_u^2){\scriptscriptstyle \frac{1}{2}}\kappa _0\rho _2+{\scriptscriptstyle \frac{1}{2}}(c,_{\theta \theta }+3c,_\theta \mathrm{cot}\theta 2c),_u$$
(21)
follows. The energy balance at null infinity (where null infinity admits a regular spherical cross section) then shows that the mass loss $`m,_u`$ results from a linear superposition of the pure and the gravitational radiation parts
$$m,_u={\scriptscriptstyle \frac{1}{2}}\underset{0}{\overset{\pi }{}}(c,_u^2+d,_u^2+{\scriptscriptstyle \frac{1}{2}}\kappa _0\rho _2)\mathrm{sin}\theta \mathrm{d}\theta 0$$
(22)
with the function $`\rho _2(u,\theta )`$ being an analogue to the news functions of electromagnetic field squared $`X^2+Y^2`$ (see (14) in ).
Since we admit spacetimes with only ”local” $`𝒥^+`$, we assume Eqs. (4)–(22) to be satisfied for $`\varphi 0,2\pi )`$, however not necessarily on the whole sphere, i.e., for all $`\theta 0,\pi `$, but only in some open interval of $`\theta `$.
Let us follow a similar procedure to that one used for the electrovacuum case and assume here again the spacetime to have another Killing vector $`\eta `$ which forms a two-dimensional Lie algebra with the axial one, $`\xi =/\varphi `$, i.e., we assume $`[\eta ,\xi ]=0`$ (see the Lemma in Sec. 2 in ). Hence, the components of $`\eta ^\alpha `$ are independent of $`\varphi `$.
We introduce the standard null tetrad $`\{k^\alpha ,m^\alpha ,t^\alpha ,\overline{t}^\alpha \}`$ (for details see (11) and the paragraph above in ), with bar denoting the complex conjugation:
$`k_\alpha `$ $`=`$ $`[1,0,0,0],\text{ }m_\alpha =[{\scriptscriptstyle \frac{1}{2}}Vr^1\mathrm{e}^{2\beta },\mathrm{e}^{2\beta },0,0],`$ (23)
$`t_\alpha `$ $`=`$ $`{\scriptscriptstyle \frac{1}{2}}r(\mathrm{cosh}2\delta )^{\frac{1}{2}}[(1+\mathrm{sinh}2\delta )\mathrm{e}^\gamma U+\mathrm{cosh}2\delta \mathrm{e}^\gamma W+\mathrm{i}[(1\mathrm{sinh}2\delta )\mathrm{e}^\gamma U\mathrm{cosh}2\delta \mathrm{e}^\gamma W],`$ (24)
$`\text{ }0,(1+\mathrm{sinh}2\delta +\mathrm{i}(1\mathrm{sinh}2\delta ))\mathrm{e}^\gamma ,(1\mathrm{i})\mathrm{cosh}2\delta \mathrm{sin}\theta \mathrm{e}^\gamma ],`$ (25)
and decompose the additional Killing vector $`\eta ^\alpha `$ in this null tetrad
$$\eta ^\alpha =Ak^\alpha +Bm^\alpha +\stackrel{~}{f}(t_R^\alpha +t_I^\alpha )+\stackrel{~}{g}(t_R^\alpha t_I^\alpha ),$$
(26)
where $`A`$, $`B`$, $`\stackrel{~}{f}`$, $`\stackrel{~}{g}`$ are general functions of $`u`$, $`r`$, $`\theta `$ and $`t^\alpha =t_R^\alpha +\mathrm{i}t_I^\alpha `$ .
The Killing vector $`\eta ^\alpha `$ has to satisfy the Killing equations (all of them are written down in )
$$_\eta g_{\alpha \beta }=0.$$
(27)
The easiest one among them is the equation
$$_\eta g_{11}=2\mathrm{e}^{2\beta }B,_r=0,$$
(28)
which implies
$$B=B(u,\theta ).$$
(29)
We solve the other Killing equations asymptotically assuming that the coefficients $`A`$, $`\stackrel{~}{f}`$, $`\stackrel{~}{g}`$ can be expanded in powers of $`r^k`$. Then equations $`_\eta g_{22}=0`$, $`_\eta g_{12}=0`$, $`_\eta g_{13}=0`$ imply
$`A`$ $`=`$ $`A^{(1)}r+A^{(0)}+{\displaystyle \frac{A^{(1)}}{r}}+𝒪(r^2),`$ (30)
$`\stackrel{~}{f}`$ $`=`$ $`f^{(1)}r+f^{(0)}+{\displaystyle \frac{f^{(1)}}{r}}+𝒪(r^2),`$ (31)
$`\stackrel{~}{g}`$ $`=`$ $`g^{(1)}r+g^{(0)}+{\displaystyle \frac{g^{(1)}}{r}}+𝒪(r^2),`$ (32)
where the coefficients $`A^{(k)}`$, $`f^{(k)}`$, $`g^{(k)}`$ are functions of $`u`$ and $`\theta `$.
Since the null dust field decays at infinity in the same way as the electromagnetic field in , it does not enter the Killing equations in the leading order in $`r^k`$ as in (see Eqs. (19)–(25) therein) and their solution is thus identical to the solution in the electrovacuum case in and even the solution obtained in the vacuum case examined in :
$`A^{(1)}`$ $`=`$ $`k\mathrm{cos}\theta ,`$ (33)
$`f^{(1)}`$ $`=`$ $`k\mathrm{sin}\theta ,`$ (34)
$`B`$ $`=`$ $`ku\mathrm{cos}\theta +\alpha (\theta ),`$ (35)
where $`k=\text{const}`$ and $`\alpha `$ is an arbitrary function of $`\theta `$ and
$$g^{(1)}=h\mathrm{sin}\theta ,$$
(36)
where $`h=\text{const}`$. One can easily find (using Eqs. (24), (26) and (31)) that the contribution of $`h`$ to the vector field $`\eta ^\alpha `$ is just constant multiple of the axial Killing vector $`/\varphi `$, $`\eta ^\varphi =h+𝒪(r^1)`$, and so we may, without loss of generality, put $`h=0`$. Therefore, in the lowest order of $`r^1`$ the general asymptotic form of the Killing vector $`\eta `$ turns out to be
$$\eta ^\alpha =[ku\mathrm{cos}\theta +\alpha (\theta ),kr\mathrm{cos}\theta +𝒪(r^0),k\mathrm{sin}\theta +𝒪(r^1),𝒪(r^1)],$$
(37)
where $`k`$ is a constant, $`\alpha `$ – an arbitrary function of $`\theta `$. Thus, assuming the presence of a null dust field satisfying the boundary conditions (7)–(13) and Killing vectors which need not be hypersurface orthogonal, we arrive in the leading order of the asymptotic expansion at the same conclusion obtained in Ref. in the vacuum case with hypersurface orthogonal $`/\varphi `$ or in Ref. for the electrovacuum case which Killing vectors are not hypersurface orthogonal. When $`k=0`$, the vector field (37) generates supertranslations.
Assuming $`k0`$, one can find a Bondi-Sachs coordinate system with $`\alpha =0`$ by making a supertranslation, as was shown in . Hence, we put $`\alpha =0`$ in Eq. (37) and without loss of generality we choose $`k=1`$. Then $`B=u\mathrm{cos}\theta `$, $`A^{(1)}=\mathrm{cos}\theta `$, $`f^{(1)}=\mathrm{sin}\theta `$ and $`g^{(1)}=0`$ and the asymptotic form of the Killing vector field $`\eta `$ is
$$\eta ^\alpha =[u\mathrm{cos}\theta ,r\mathrm{cos}\theta +𝒪(r^0),\mathrm{sin}\theta +𝒪(r^1),𝒪(r^1)],$$
(38)
that is the boost Killing vector. It generates the Lorentz transformations along the axis of axial symmetry.
The conclusion of this section is thus following:
Suppose that an axially symmetric spacetime with null dust admits a “piece” of $`𝒥^+`$ in the sense that the Bondi-Sachs coordinates can be introduced in which the metric takes the form (4), (17) and the asymptotic forms of the energy-mass density and the null vector field of the null dust is given by (7)–(13). If this spacetime admits an additional Killing vector forming with the axial Killing vector a two-dimensional Lie algebra, then the additional Killing vector has asymptotically the form (37). For $`k=0`$ it generates a supertranslation; for $`k0`$ it generates a boost along the symmetry axis.
However in the next sections we see that the boost symmetry is in fact not allowable.
## III The supertranslational Killing field
In this section, assuming $`k=0`$ for the Killing field (37), we consider the Killing equations in higher orders of $`r^1`$ and arrive at the same equations as (31)–(48) in with the same solutions (49)–(55) therein:
$`c`$ $`=`$ $`{\scriptscriptstyle \frac{1}{2}}uB^1(B,_{\theta \theta }B,_\theta \mathrm{cot}\theta ),\text{ }d=d(\theta ),\text{ }M=uc,_u^2B^1[A^{(1)}+B,_\theta (c,_\theta +2c\mathrm{cot}\theta )],`$ (39)
$`A^{(0)}`$ $`=`$ $`{\scriptscriptstyle \frac{1}{2}}(B,_{\theta \theta }+B,_\theta \mathrm{cot}\theta +B),\text{ }f^{(0)}=f^{(0)}(\theta )=B,_\theta ,\text{ }g^{(0)}=0,`$ (40)
$`A^{(1)}`$ $`=`$ $`A^{(1)}(\theta ),\text{ }f^{(1)}=B(c,_\theta +2c\mathrm{cot}\theta ),\text{ }g^{(1)}=B(d,_\theta +2d\mathrm{cot}\theta )B,_\theta d,`$ (41)
where $`B`$, $`A^{(1)}`$ and $`d`$ are arbitrary functions of $`\theta `$. Thus, in the supertranslational case for the null dust the Weyl tensor is also non-radiative as it was in electrovacuum spacetimes. Substituting the metric functions (17) into the null tetrad (24) and coefficients (40), the expansion of the Killing vector reads
$`\eta ^\mu `$ $`=`$ $`[B(\theta ),{\scriptscriptstyle \frac{1}{2}}(B,_{\theta \theta }+B,_\theta \mathrm{cot}\theta )+[B,_{\theta \theta }2B,_\theta B,_{\theta \theta \theta }+2B,_{\theta }^{}{}_{}{}^{2}B,_{\theta \theta }B^12B,_{\theta }^{}{}_{}{}^{3}\mathrm{cot}\theta B^1`$ (42)
$`\text{ }+B,_{\theta }^{}{}_{}{}^{2}(3\mathrm{cot}\theta ^22\mathrm{sin}^2\theta )]B^1{\displaystyle \frac{u}{4r}}+𝒪(r^2),B,_\theta {\displaystyle \frac{1}{r}}+B,_\theta {\displaystyle \frac{c}{r^2}}+𝒪(r^3),B,_\theta {\displaystyle \frac{d}{r^2\mathrm{sin}\theta }}+𝒪(r^3)].`$ (43)
Now let us turn to the asymptotic behaviour of the null dust. It is easy to show that if $`\eta `$ is a Killing vector then the Lie derivative of the Riemann tensor with respect to this vector vanishes, and then also the Lie derivative of the Ricci tensor vanishes. And if, in addition, the Ricci scalar is zero and Einstein’s equations are satisfied, then the following equations hold
$$_\eta T_{\mu \nu }=0.$$
(44)
Substituting $`\eta `$ from (43), equations (44) get in the leading order the form
$`L_\eta T_{00}=0\text{ }(r^2):`$ $`\text{ }\rho _2,_uB=0,`$ (45)
$`L_\eta T_{01}=0\text{ }(r^4):`$ $`\text{ }\rho _2,_u𝒰B+\rho _2(𝒰,_uB𝒯B,_\theta )=0,`$ (46)
$`L_\eta T_{02}=0\text{ }(r^2):`$ $`\text{ }\rho _2,_u𝒯B+\rho _2(𝒯,_uBB,_\theta )=0,`$ (47)
$`L_\eta T_{03}=0\text{ }(r^2):`$ $`\text{ }\rho _2,_uB+\rho _2,_uB=0,`$ (48)
$`L_\eta T_{11}=0\text{ }(r^6):`$ $`𝒰[\rho _2,_u𝒰B+2\rho _2(𝒰,_uB𝒯B,_\theta )]=0,`$ (49)
$`L_\eta T_{12}=0\text{ }(r^4):`$ $`\text{ }\rho _2,_u𝒰𝒯B+\rho _2[(𝒰𝒯),_uB(𝒰+𝒯^2)B,_\theta ]=0,`$ (50)
$`L_\eta T_{13}=0\text{ }(r^4):`$ $`\text{ }\rho _2,_u𝒰B+\rho _2[(𝒰),_uB𝒯B,_\theta ]=0,`$ (51)
$`L_\eta T_{22}=0\text{ }(r^2):`$ $`𝒯[\rho _2,_u𝒯B+2\rho _2(𝒯,_uBB,_\theta )]=0,`$ (52)
$`L_\eta T_{23}=0\text{ }(r^2):`$ $`\text{ }\rho _2,_u𝒯B+\rho _2[(𝒯),_uBB,_\theta ]=0,`$ (53)
$`L_\eta T_{33}=0\text{ }(r^2):`$ $`\text{ }\rho _2,_uB+2\rho _2,_uB=0.`$ (54)
If $`\rho _2=0`$, then all equations are trivially satisfied and it can be shown that we deal with a vacuum spacetime with an arbitrary supertranslational symmetry. If we suppose $`B`$ and $`\rho _2`$ to be non-vanishing, then the first equation (45) implies
$$\rho _2=\rho _2(\theta ),$$
(55)
and from Eqs. (48) or (54), (47) or (52), and (46), (49) or (51) it follows
$`,_u`$ $`=`$ $`0,`$ (56)
$`𝒯,_u`$ $`=`$ $`B,_\theta B^1,`$ (57)
$`𝒰,_u`$ $`=`$ $`B,_\theta B^1𝒯,`$ (58)
which yield
$``$ $`=`$ $`_0(\theta ),`$ (59)
$`𝒯`$ $`=`$ $`B,_\theta B^1u+𝒯_0(\theta ),`$ (60)
$`𝒰`$ $`=`$ $`{\scriptscriptstyle \frac{1}{2}}B,_\theta ^2B^2u^2+B,_\theta B^1𝒯_0u+{\scriptscriptstyle \frac{1}{2}}(𝒯_{0}^{}{}_{}{}^{2}+_{0}^{}{}_{}{}^{2}\mathrm{sin}^2\theta ),`$ (61)
where we used condition (13). Comparing (21) with (40) we obtain the following equation determining the function $`B(\theta )`$ for the given “null dust” news function $`\rho _2(\theta )`$:
$$\left\{\frac{B^2}{\mathrm{sin}\theta }\left[\frac{\mathrm{sin}^3\theta }{2B}\left(\frac{B,_\theta }{\mathrm{sin}\theta }\right),_\theta \right],_\theta \right\},_\theta =\kappa _0\rho _2B^2\mathrm{sin}\theta .$$
(62)
If a spacetime admits a global null infinity this case with a non-zero $`\rho _2`$ is not very physical since (22) then implies permanent linear decreasing of the total Bondi mass $`m`$.
## IV The boost Killing field
In this section we investigate the boost case, $`k=1`$ and $`\alpha =0`$, similarly as the previous one by expanding the Killing equations in higher orders of $`r^1`$. We obtain the same conditions for the coefficients $`A^{(k)}`$, $`f^{(k)}`$ and $`g^{(k)}`$ as Eqs. (81)–(89) and (104)–(112) in with the identical solutions
$`A^{(0)}(u,\theta )`$ $`=`$ $`{\scriptscriptstyle \frac{1}{2}}u\mathrm{cos}\theta ,`$ (63)
$`f^{(0)}(u,\theta )`$ $`=`$ $`u^2w+𝒦(w)=u\mathrm{sin}\theta +𝒦(\mathrm{sin}\theta /u),`$ (64)
$`g^{(0)}(u,\theta )`$ $`=`$ $`g^{(0)}(w),`$ (65)
$`A^{(1)}(u,\theta )`$ $`=`$ $`{\displaystyle \frac{\mathrm{cos}\theta }{w}}({\scriptscriptstyle \frac{1}{2}}𝒦,_{ww}w^2+2𝒦,_ww+𝒦){\displaystyle \frac{\mathrm{cos}\theta }{8\mathrm{sin}^2\theta }}(4𝒦^2+g_{}^{(0)}{}_{}{}^{2})+{\displaystyle \frac{\mathrm{cos}\theta }{u^2}}(w),`$ (66)
$`f^{(1)}(u,\theta )`$ $`=`$ $`{\displaystyle \frac{1}{4\mathrm{sin}\theta }}(2𝒦^2+g_{}^{(0)}{}_{}{}^{2})+u\mathrm{cot}^2\theta (𝒦,_ww+𝒦),`$ (67)
$`g^{(1)}(u,\theta )`$ $`=`$ $`{\scriptscriptstyle \frac{1}{2}}ug^{(0)}{\displaystyle \frac{g^{(0)}𝒦}{\mathrm{sin}\theta }}+{\scriptscriptstyle \frac{1}{2}}u\mathrm{cot}^2\theta (g^{(0)},_ww+g^{(0)}),`$ (68)
$`c(u,\theta )`$ $`=`$ $`{\displaystyle \frac{𝒦(w)}{uw}}\text{ }c,_u(u,\theta )={\displaystyle \frac{𝒦(w),_w}{u^2}},`$ (69)
$`d(u,\theta )`$ $`=`$ $`{\displaystyle \frac{g^{(0)}(w)}{2\mathrm{sin}\theta }}\text{ }d,_u(u,\theta )={\displaystyle \frac{g^{(0)}(w),_w}{2u^2}},`$ (70)
$`M(u,\theta )`$ $`=`$ $`{\displaystyle \frac{1}{2u}}(𝒦,_{ww}w+2𝒦,_w)+{\displaystyle \frac{(w)}{u^3}}={\displaystyle \frac{1}{2\mathrm{sin}\theta }}(w^2𝒦,_w),_w+{\displaystyle \frac{}{u^3}},`$ (71)
where $`𝒦(w)`$ and $`g^{(0)}(w)`$ are arbitrary functions of $`w`$
$$w=\frac{\mathrm{sin}\theta }{u},$$
(72)
and the integration function $`(w)`$ entering the expression for the mass aspect has the form
$$(w)=\frac{\lambda (w)}{w^3},$$
(73)
with $`\lambda `$ satisfying the equation
$$\lambda (w),_w=w^2(𝒦,_w^2+{\scriptscriptstyle \frac{1}{4}}g^{(0)},_w^2+{\scriptscriptstyle \frac{1}{2}}\kappa _0\omega )\frac{1}{2w}(w^3𝒦,_{ww}),_w.$$
(74)
Here, $`𝒦`$ and $`g^{(0)}`$ determine the gravitational news functions, $`c,_u`$ and $`d,_u`$, by the relations (69) and (70), $`\omega `$ determines the null dust news function, $`\rho _2`$, given below by the relation (89). Hence, solving the last equation for $`\lambda `$ for given $`𝒦`$, $`g^{(0)}`$, and $`\omega `$, we find $`(w)`$ and thus the mass aspect $`M(u,\theta )`$ in the form of Eq. (71). The total mass $`m`$ at $`𝒥^+`$ is then given by integrating Eq. (71) over the sphere:
$$m(u)={\scriptscriptstyle \frac{1}{2}}\underset{0}{\overset{\pi }{}}M(u,\theta )\mathrm{sin}\theta \mathrm{d}\theta ={\scriptscriptstyle \frac{1}{4}}\underset{0}{\overset{\pi }{}}(w^2𝒦,_w),_w\mathrm{d}\theta +{\scriptscriptstyle \frac{1}{2}}\underset{0}{\overset{\pi }{}}\frac{w}{u^2}\mathrm{d}\theta .$$
(75)
Substituting the expansions of the metric functions, Eq. (17) into the null tetrad, Eq. (24), and coefficients $`A`$, $`B`$, $`\stackrel{~}{f}`$ and $`\stackrel{~}{g}`$, Eqs. (63)–(71), into Eq. (26), we find the asymptotic form of the boost Killing vector to be
$`\eta ^\mu `$ $`=`$ $`[u\mathrm{cos}\theta ,r\mathrm{cos}\theta +u\mathrm{cos}\theta +\mathrm{cos}\theta (𝒦,_w+{\displaystyle \frac{𝒦}{w}}){\displaystyle \frac{1}{r}}+𝒪(r^2),`$ (76)
$`\text{ }\mathrm{sin}\theta u\mathrm{sin}\theta {\displaystyle \frac{1}{r}}+uc\mathrm{sin}\theta {\displaystyle \frac{1}{r^2}}+𝒪(r^3),\text{ }ud{\displaystyle \frac{1}{r^2}}+𝒪(r^3)].`$ (77)
Finally, let us turn our attention to the asymptotic properties of the null dust represented by the energy-momentum tensor $`T_{\mu \nu }`$ which, as was shown in the previous section, has to have the vanishing Lie derivative (44) with respect to the Killing vector $`\eta ^\alpha `$, (77). Regarding (63)–(71), these equations in the first orders look as follows:
$`_\eta T_{00}=0\text{ }(r^2):`$ $`\mathrm{cos}\theta (u\rho _2,_u+\mathrm{tan}\theta \rho _2,_\theta 4\rho _2)=0,`$ (78)
$`_\eta T_{01}=0\text{ }(r^4):`$ $`\mathrm{cos}\theta [(u\rho _2,_u+\mathrm{tan}\theta \rho _2,_\theta )𝒰+\rho _2(u𝒰,_u+\mathrm{tan}\theta 𝒰,_\theta +4𝒰+u\mathrm{tan}\theta 𝒯)]=0,`$ (79)
$`_\eta T_{02}=0\text{ }(r^2):`$ $`\text{ }\mathrm{cos}\theta [(u\rho _2,_u+\mathrm{tan}\theta \rho _2,_\theta )𝒯+\rho _2(u𝒯,_u+\mathrm{tan}\theta 𝒯,_\theta +4𝒯+u\mathrm{tan}\theta )]=0,`$ (80)
$`_\eta T_{03}=0\text{ }(r^2):`$ $`\text{ }\mathrm{cos}\theta [(u\rho _2,_u+\mathrm{tan}\theta \rho _2,_\theta )+\rho _2(u,_u+\mathrm{tan}\theta ,_\theta +5)]=0,`$ (81)
$`_\eta T_{11}=0\text{ }(r^6):`$ $`\mathrm{cos}\theta [(u\rho _2,_u+\mathrm{tan}\theta \rho _2,_\theta )𝒰+2\rho _2(u𝒰,_u+\mathrm{tan}\theta 𝒰,_\theta +2𝒰+u\mathrm{tan}\theta 𝒯)]=0,`$ (82)
$`_\eta T_{12}=0\text{ }(r^4):`$ $`\text{ }\mathrm{cos}\theta \{(u\rho _2,_u+\mathrm{tan}\theta \rho _2,_\theta )𝒰𝒯+\rho _2[u(𝒰𝒯),_u`$ (84)
$`\text{ }+\mathrm{tan}\theta (𝒰𝒯),_\theta +4𝒰𝒯+u\mathrm{tan}\theta (𝒰+𝒯^2)]\}=0,`$
$`_\eta T_{13}=0\text{ }(r^4):`$ $`\text{ }\mathrm{cos}\theta \{(u\rho _2,_u+\mathrm{tan}\theta \rho _2,_\theta )𝒰+\rho _2[u(𝒰),_u+\mathrm{tan}\theta (𝒰),_\theta +5𝒰+u\mathrm{tan}\theta 𝒯]\}=0,`$ (85)
$`_\eta T_{22}=0\text{ }(r^2):`$ $`\mathrm{cos}\theta [(u\rho _2,_u+\mathrm{tan}\theta \rho _2,_\theta )𝒯+2\rho _2(u𝒯,_u+\mathrm{tan}\theta 𝒯,_\theta +2𝒯+u\mathrm{tan}\theta )]=0,`$ (86)
$`_\eta T_{23}=0\text{ }(r^2):`$ $`\mathrm{cos}\theta \{(u\rho _2,_u+\mathrm{tan}\theta \rho _2,_\theta )𝒯+\rho _2[u(𝒯),_u+\mathrm{tan}\theta (𝒯),_\theta +5𝒯+u\mathrm{tan}\theta ]\}=0,`$ (87)
$`_\eta T_{33}=0\text{ }(r^2):`$ $`\mathrm{cos}\theta [(u\rho _2,_u+\mathrm{tan}\theta \rho _2,_\theta )+2\rho _2(u,_u+\mathrm{tan}\theta ,_\theta +3)]=0.`$ (88)
The trivial solution is again $`\rho _2=0`$ which implies a vacuum boost-rotation symmetric spacetime. Let us assume $`\rho _20`$. Using variable $`w`$ given by (72), Eq. (78) can be solved to yield
$$\rho _2=\frac{\omega (w)}{u^2},$$
(89)
with an arbitrary function $`\omega (w)`$. The sum of Eqs. (81) and (88) gives
$$=\frac{_0(w)}{u},$$
(90)
($`_0(w)`$ being an arbitrary function) and their difference is then identically zero. Similarly, summing Eqs. (80) and (86), we obtain the equation for $`𝒯`$,
$$𝒯,_u=\mathrm{tan}\theta =\frac{uw}{\sqrt{1u^2w^2}},$$
(91)
which leads to the solution
$$𝒯=\frac{\sqrt{1u^2w^2}}{w}+𝒯_0(w)=u\mathrm{cot}\theta +𝒯_0(w),$$
(92)
where $`𝒯_0(w)`$ is an arbitrary integration function. Next, the difference of Eqs. (80) and (86) identically vanishes. We repeat the procedure for the sum of Eqs. (79) and (82) and find the solution
$$𝒰={\scriptscriptstyle \frac{1}{2}}u^2+𝒯_0(w)\frac{\sqrt{1u^2w^2}}{w}+𝒰_0(w)={\scriptscriptstyle \frac{1}{2}}u^2+u\mathrm{cot}\theta 𝒯_0(w)+𝒰_0(w),$$
(93)
($`𝒰_0(w)`$ is an arbitrary function of $`w`$). Then their difference is identically zero and also Eqs. (84), (85), (87) are identically satisfied.
The coefficients $`𝒰`$, $`𝒯`$ and $``$ have, in addition, to fulfil the condition for the null vector (13). This, however, is in contradiction with the boost-rotation symmetric solutions (90), (92) and (93). Consequently, there are no asymptotically flat boost-rotation symmetric solutions of the Einstein equations with null dust. And the final conclusion reads:
Theorem
Suppose that an axially symmetric spacetime with null dust admits a “piece” of $`𝒥^+`$ in the sense that the Bondi-Sachs coordinates can be introduced in which the metric takes the form (4), (17) and the asymptotic forms of the radiation density and the null vector field of the null dust are given by (7)–(13). If this spacetime admits an additional Killing vector forming a two-dimensional Lie algebra with the axial Killing vector, then the additional Killing vector which has asymptotically the form (43) generates asymptotically supertranslations and the Weyl tensor is non-radiative (although one of the “gravitational” news functions, $`c,_u`$, is non-vanishing, however, it is a function only of $`\theta `$ as the “null dust” news function, $`\rho _2`$).
ACKNOWLEDGMENTS
We are grateful to prof. Bičák for suggesting the topic and for discussions. |
warning/0002/math0002097.html | ar5iv | text | # Calibrated Fibrations on Complete Manifolds via Torus Action
## 1 Introduction
In this paper we will use structure preserving torus actions on non-compact manifolds with calibrations to construct calibrated submanifolds (both for Calabi-Yau manifolds and for 7-manifolds with a $`G_2`$ structure). We will assume that no element of the torus acts trivially. Throughout the paper we will use the notion of a calibrated fibration:
###### Definition 1
Let $`(M,\phi )`$ be a Riemannian manifold with a calibrating form $`\phi `$. Then we say that $`M`$ has a calibrated fibration on it if there is a surjective map $`\alpha :MV`$ onto a topological space $`V`$ and a subset $`SM`$ s.t.
i) For any point $`mMS`$ the level set $`L_m`$ of $`\alpha `$ through $`m`$ is a smooth submanifold, calibrated by $`\phi `$.
ii) The set $`S`$ is locally contained in a finite union of submanifolds of codimension $`4`$ in $`M`$.
In Section 2 we consider a Kahler manifold $`(M^{2n},\omega )`$ with a non-vanishing holomorphic $`(n,0)`$ form $`\phi `$. We can define, as in , Special Lagrangian (SLag) submanifolds $`L`$ by the conditions:
$$\omega |_L=0,Im\phi |_L=0.$$
If $`g`$ is the Riemannian metric corresponding to $`\omega `$, then we can conformally scale $`g`$ to a metric $`g^{}`$ on $`M`$ so that the form $`\phi `$ will have length $`\sqrt{2}^n`$ with respect to $`g^{}`$. Then SLag submanifolds will be calibrated by $`Re\phi `$ with respect to $`g^{}`$. In particular, they will be minimal submanifolds of $`(M,g^{})`$ and Lagrangian submanifolds of $`\omega `$.
If $`M`$ is compact and simply connected, then for any Kahler form $`\omega `$ on $`M`$ Yau’s celebrated resolution of the Calabi conjecture gives a (unique) Ricci-flat Kahler form $`\omega ^{}`$ in the same cohomology class as $`\omega `$ (see ). Also the SYZ conjecture (see ) states that $`(M,Re\phi )`$ has a calibrated fibration with generic fiber being a SLag torus with respect to a Ricci-flat Kahler metric. We can ask an analogous question for any Kahler metric on $`M`$ and we showed in that this holds for a choice of Kahler metric on a Borcea-Voisin threefold. In this paper we will be interested in non-compact Calabi-Yau manifolds with a structure-preserving torus action. The main results of Section 2 are as follows:
###### Theorem 1
Suppose we have a Hamiltonian structure-preserving $`T^k`$-action on $`M^{2n}`$. Then any smooth symplectic reduction $`M_{red}`$ has a natural holomorphic volume form $`\phi _{red}`$. Moreover SLag submanifolds of $`M_{red},\omega _{red},\phi _{red}`$ lift to SLag submanifolds of $`M`$, invariant under $`T^k`$-action. Vice versa, let $`L`$ be a connected, $`T`$-invariant SLag submanifold of $`M`$ s.t. $`T`$ acts freely on $`L`$. Then $`L`$ lies on a level set $`a`$ of the moment map $`\mu `$ and one can find a SLag submanifold $`L^{}`$ in the smooth part of the symplectic reduction through level set $`a`$ s.t. $`L^{}`$ lifts to $`L`$.
###### Theorem 2
Suppose that $`k=n1`$ and $`H^1(M,)=0`$. Then $`M`$ has a calibrated fibration $`\alpha `$ over an open subset of $`^n`$ with the set $`S`$ of singular points being the non-regular points of the torus action (i.e points there the differential of the action is not injective). Moreover for a generic point $`p`$ (outside of a countable union of $`(n2)`$-dimensional planes in $`^n`$), the fiber $`\alpha ^1(p)`$ is a smooth SLag submanifold. Connected components of each smooth fiber are diffeomorphic to an $`(n1)`$-torus times $``$. Singular fibers have singularities of codimension at least 2, and near singular points they are diffeomorphic to a product of a cone with a Euclidean ball.
If we make certain assumptions on the set of non-regular points of $`T`$-action, then we can replace the countable union by a finite union in Theorem 2 (see Theorem 3 in Section 2).
In Section 3 we consider a 7-manifold $`M`$ with a $`G_2`$-form $`\phi `$ (see ). Let $`H^1(M,)=0`$. If $`M`$ has a 3-torus action, then $`M`$ is covered by a family of non-intersecting coassociative submanifolds. Suppose $`M`$ admits a 2-torus action. Then we will define certain $`G_2`$-reductions $`M_{red}`$, which are symplectic 4-manifolds with a compatible almost complex structure and trivial canonical bundle. We will see that 2-dimensional complex sub-varieties of $`M_{red}`$ lift to $`T`$-invariant, coassociative submanifolds of $`M`$. Finally suppose $`M`$ admits an $`SO(3)`$-action, which is not regular in at least one point. If $`SO(3)`$ acts freely on the set $`M^{}`$ of regular points of the $`SO(3)`$-action, then $`M^{}`$ is covered by a family of non-intersecting coassociative submanifolds, diffeomorphic to $`SO(3)\times `$.
In Section 4 we will consider some applications of the results of the two previous sections. For the Calabi-Yau case, we will show that crepant resolutions of singularities of a finite Abelian subgroup of $`SU(n)`$ acting on $`^n`$ have SLag fibrations. Also for any Kahler-Einstein manifold $`N`$ with positive scalar curvature, the total space $`K(N)`$ of it’s canonical bundle is a Calabi-Yau manifold (see ). We investigate SLag submanifolds on $`K(N)`$. For each orientable minimal Lagrangian submanifold of $`N`$ we associate a 1-parameter family $`(L_\lambda |\lambda )`$ of SLag submanifolds of $`K(N)`$. Also $`L_0`$ is invariant under scaling of $`K(N)`$ by a real number. For any compact Kahler manifold $`N^{2n}`$ with an effective n-torus action we prove that one of the regular orbits of the action is a minimal Lagrangian submanifold of $`N`$. If $`N`$ is Kahler-Einstein with nonzero scalar curvature $`t`$ and toric, then we prove that precisely 1 such orbit $`L`$ is a minimal Lagrangian submanifold of $`N`$. For $`t>0`$ we use Theorem 2 to construct a SLag fibration on $`K(N)`$ and we prove that all fibers are asymptotic at infinity to the fiber $`L_0`$. We conjecture that any K-E manifold $`N`$ with positive scalar curvature has a minimal Lagrangian submanifold $`L`$. Moreover $`K(N)`$ fibers with generic fiber being a SLag submanifold of $`K(N)`$ and all fibers are asymptotic to $`L_0`$ at infinity.
In the $`G_2`$ case, Bryant and Salamon have constructed in some examples of complete metrics with holonomy $`G_2`$. Some metrics are on the total space of the spin bundle over a 3-dimensional space form. Others are on a total space of a bundle $`\mathrm{\Lambda }_{}^2`$ of anti-self-dual 2-forms over a self-dual Einstein 4-manifold. Many examples admit $`T^2`$ and $`SO(3)`$-actions, and we show that in one example the $`G_2`$ manifold $`M`$ can be covered by non-intersecting coassociative $`SO(3)`$-invariant submanifolds.
Acknowledgments : This paper is a part of author’s work towards his Ph.D. at the Massachusetts Institute of Technology. The author wants to express his gratitude to his advisor, Tom Mrowka, for initiating him into this subject and for continuing support.
After writing this paper the author learned that Mark Gross has independently obtained results, which are similar to some results of Sections 2 and 4.3.
## 2 Torus action on Calabi-Yau manifolds
Let $`(M^{2n},\omega ,\phi )`$ be a Calabi-Yau manifold with a structure-preserving Hamiltonian $`T^k`$-action. For any element $`v`$ of the Lie algebra $`𝒢`$ of $`T^k`$ we associate the infinitesimal flow vector field $`X_v`$ on $`M`$, induced by the differential of the action. Then the $`X_v`$ commute and their flows preserve $`\omega `$ and $`\phi `$. Let $`v_1,\mathrm{},v_l`$ be elements of $`𝒢`$ and $`X_1,\mathrm{},X_l`$ be the corresponding vector fields on $`M`$. We claim that the $`(nl,0)`$\- form $`\phi ^{}=i_{X_1}\mathrm{}i_{X_l}\phi `$ (obtained by contraction of $`\phi `$ by the vector fields $`X_1,\mathrm{},X_l`$) is a closed $`(nl,0)`$form. We will prove it by induction on $`l`$. If $`l=1`$ then the $`X_1`$-flow preserves $`\phi `$ and so
$$0=_{X_1}\phi =d(i_{X_1}\phi )=d\phi ^{}$$
Now we use induction on $`l`$. The $`X_1`$ flow preserves $`X_2,\mathrm{},X_l`$ and it preserves $`\phi `$, so it preserves $`\phi ^{\prime \prime }=i_{X_2}\mathrm{}i_{X_l}\phi `$. So
$$0=_{X_1}\phi ^{\prime \prime }=d(i_{X_1}\phi ^{\prime \prime })+i_{X_1}(d\phi ^{\prime \prime })$$
and we are done by induction.
A moment map $`\mu `$ for the $`T^k`$-action is a map $`\mu :M𝒢^{}`$ (the dual Lie Algebra of $`T^k`$), which satisfies
$$d(\mu (v))=i_{X_v}\omega $$
for any $`v𝒢`$. The moment map is $`T^k`$-invariant.
###### Theorem 1
Any smooth symplectic reduction $`M_{red}`$ has a natural holomorphic volume form $`\phi _{red}`$. Moreover SLag submanifolds of $`M_{red}`$ lift to SLag submanifolds of $`M`$, invariant under $`T^k`$action. Vice versa let $`L`$ be a connected, $`T`$-invariant SLag submanifold of $`M`$. Then $`L`$ lies on a level set $`a`$ of the moment map $`\mu `$ and moreover there is a SLag submanifold $`L^{}`$ on a smooth part of symplectic reduction through level set $`a`$ s.t. $`L`$ is the lift of $`L^{}`$.
Proof : Consider a level set $`\mathrm{\Sigma }_a`$ of the moment map $`\mu `$, on which $`T^k`$ acts freely. Since $`\mu `$ is $`T^k`$invariant, the vector fields $`X_v`$ are tangent to $`\mathrm{\Sigma }_a`$. Consider the bundle $`V=span(X_v)`$ over $`\mathrm{\Sigma }_a`$. Since $`d(\mu (v))=i_{X_v}\omega `$, the tangent bundle to $`\mathrm{\Sigma }_a`$ is the $`\omega `$orthogonal complement of $`V`$. Also $`VJV=0`$ (here $`J`$ is the complex structure on $`M`$). Let $`W=(VJV)^{}`$ ( here $``$ is with respect to the metric). Then $`W`$ is a complex vector bundle over $`\mathrm{\Sigma }_a`$, the tangent bundle to $`\mathrm{\Sigma }_a`$ is $`WV`$ and the quotient of $`W`$ by the $`T`$-action can be viewed as a tangent bundle to the symplectic reduction $`M_{red}=\mathrm{\Sigma }_a/T^k`$.
Let $`v_1,\mathrm{},v_k`$ be a basis for the Lie algebra of $`T^k`$ and $`X_1,\mathrm{},X_k`$ corresponding vector fields on $`M`$. Let $`\phi ^{}=i_{X_1}\mathrm{}i_{X_k}\phi `$. Then as we have seen, $`\phi ^{}`$ is a holomorphic $`(nk,0)`$ form on $`M`$. Also $`\phi ^{}|_W`$ is a holomorphic volume form on $`W`$. Since $`\phi ^{}`$ is $`T^k`$-invariant, it is clear that there is a unique $`(nk,0)`$ form $`\phi _{red}`$ on $`M_{red}`$ s.t. $`\pi ^{}(\phi _{red})=\phi ^{}`$ on $`\mathrm{\Sigma }_a`$ (here $`\pi :\mathrm{\Sigma }_aM_{red}`$ is the quotient map). Now $`\phi _{red}`$ is closed (since $`\phi ^{}`$ is), and hence it is a holomorphic volume form on $`M_{red}`$.
Let $`L^{}`$ be a SLag submanifold of $`M_{red}`$ and $`L=\pi ^1(L^{})`$. It is clear that $`L`$ is a SLag submanifold of $`M`$, invariant under $`T^k`$-action. Vice versa, let $`L`$ be a connected SLag submanifold of $`M`$, invariant under $`T^k`$-action and $`T`$ acts freely on $`L`$. Since $`L`$ is invariant under the torus action, the vector fields $`X_v`$ are tangent to $`L`$. Let $`u`$ be some vector, tangent to $`L`$. Since $`L`$ is Lagrangian, we have
$$0=\omega (X_v,u)=d(\mu (v))(u)$$
So the differential of the moment map is zero on $`L`$, hence $`L`$ lies on some level set $`\mathrm{\Sigma }_a`$ of the moment map. Since $`T`$-action is free on $`L`$, it is freely on a neighbourhood $`U`$ of $`L`$ in $`\mathrm{\Sigma }_a`$. Moreover $`U`$ is a smooth submanifold of $`M`$. So $`U_{red}=U/T`$ will be an open set in the smooth part of the symplectic reduction through $`a`$. It is clear that the quotient $`L^{}=L/T^k`$ is a SLag submanifold on $`U_{red}`$. Q.E.D.
In case $`k=n1`$ $`M`$ has a calibrated fibration by SLag submanifolds, invariant under $`T^k`$-action:
###### Theorem 2
Let $`k=n1`$ and $`H^1(M,)=0`$. Then
i) $`M`$ has a calibrated fibration $`\alpha `$ over an open subset of $`^n`$ with the set $`S`$ of singular points being the non-regular points of the $`T`$-action.
ii) For a generic point $`p`$ (outside of a countable union of $`(n2)`$-planes in $`^n`$), the fiber $`\alpha ^1(p)`$ is a smooth SLag submanifold of $`M`$.
iii) Connected components of each smooth fiber are diffeomorphic to a product of an $`(n1)`$-torus with $``$.
iv) Singular fibers have singularities of codimension at least 2, and near a singular point they are diffeomorphic to a product of a cone with a Euclidean ball.
Proof: Define the form $`\phi ^{}`$ as in the proof of Theorem 1. Then $`\phi ^{}`$ is a holomorphic $`(1,0)`$-form, invariant under the torus action. Since $`H^1(M,)=0`$, there is a holomorphic function $`f=\eta +i\xi `$ s.t. $`df=\phi ^{}`$. It is clear that $`f`$ is also invariant under $`T^k`$-action. Let $`L`$ be a connected SLag submanifold, invariant under the torus action. As we have seen, $`L`$ must lie on the level set of the moment map $`\mu `$. Also, since $`L`$ is Special, one easily deduces that $`Im\phi ^{}|_L=0`$, so $`L`$ must lie on a level set of $`\xi =Imf`$, i.e. $`L`$ lies on a level set of n-functions
$$\mu =a,\xi =c$$
The moment map $`\mu `$ goes to $`𝒢^{}`$, which we identify with $`^{n1}`$ by choosing a basis of $`𝒢`$. We define $`\alpha =(\mu ,\xi ):M^n`$.
Let $`S`$ be the set of non-regular points of the torus action. We claim that a level set $`L_m`$ of $`(\mu ,\xi )`$, that passes through a regular point $`mMS`$ is smooth n-dimensional SLag submanifold of $`M`$ near $`m`$. Indeed let $`\mathrm{\Sigma }_a`$ be the level set of the moment map passing through $`m`$ and $`V`$ and $`W`$ be vector bundles on $`\mathrm{\Sigma }_a`$ near $`m`$ as in the proof of Theorem 1. Let $`v_1,\mathrm{},v_k`$ be a basis for the Lie algebra of $`T^k`$. Then $`d\mu (v_1),\mathrm{},d\mu (v_k)`$ is basis of $`(JV)^{}`$. Also those 1-forms vanish on $`W`$. Now $`d\xi =Im\phi ^{}`$ restricted to $`W`$ is non-zero. So the forms $`d\mu (v_1),\mathrm{},d\mu (v_k),d\xi `$ are linearly independent at $`m`$, and so the level set $`L_m`$ is a smooth submanifold of $`M`$ near $`m`$.
Next we prove that $`L_m`$ is SLag near $`m`$. Since $`\xi `$ and $`\mu `$ are $`T^k`$-invariant, then so is $`L_m`$. So the $`X_v`$ are in the tangent space to $`L_m`$ at $`m`$. Since $`L_m`$ is on the level set of $`\mu `$, the tangent space to $`L_m`$ at $`m`$ is $`\omega `$-orthogonal to $`X_v^{}s`$, so it must be Lagrangian. Also $`Im\phi ^{}|_{L_m}=0`$ implies that $`L_m`$ is Special at $`m`$.
Now we prove iii): Let $`L`$ be a level set of $`(\mu ,\xi )`$ s.t. all points on $`L`$ are regular points for $`T`$-action on $`M`$. Let $`L^{}`$ be a connected component of $`L`$. We have $`Re\phi ^{}=d\eta `$. One easily sees that $`Re\phi ^{}|_L0`$ for all points of $`L`$. One also easily shows that $`\eta `$ along $`L`$ is tangent to $`L`$, hence it coincides with $`(\eta |_L)`$.
Consider a $`T^{n1}`$-orbit $`T`$ on $`L^{}`$. $`T`$ lives on some level set of $`\eta `$ on $`L^{}`$, hence it must coincide with a connected component of this level set. The normalized gradient flow of $`T`$ by $`\frac{\eta }{|\eta |^2}`$ on $`L^{}`$ gives a diffeomorphism between $`L^{}`$ and $`T^{(n1)}`$ times $``$. Indeed this flow is defined on an interval $`(a_{},a_+)`$, there $`a_{},a_+`$ are independent on the choice of a point in $`T`$. One easily deduces that the orbit of $`T`$ under this flow is isolated in $`L`$, hence it is equal to $`L^{}`$.
Next we prove ii) and show that $`S`$ is locally contained in a finite union of submanifolds of codimension 4 in $`M`$. To prove this we need to understand the picture near a point in $`S`$. Let $`mS`$. The differential of the action is not injective at $`m`$ and $`m`$ has a stabilizer $`T^{}`$ of positive dimension $`l`$ and an orbit $`O`$. To prove ii) we need to see what is the image of $`(\mu ,\xi )`$ on $`S`$ near $`m`$.
The symplectic form $`\omega `$ restricts trivially to $`O`$, hence we have $`\omega =d\gamma `$ for some 1-form $`\gamma `$ in a neighbourhood of $`O`$. Since $`\omega `$ is $`T^{n1}`$-invariant, we can make $`\gamma `$ invariant as well (by integrating over $`T^{n1}`$). For any $`v𝒢`$ we have $`0=_{X_v}\gamma =d(\gamma (X_v))+i_{X_v}\omega `$. So the map $`\mu ^{}(v)=\gamma (X_v)`$ is a moment map near $`O`$ and $`\mu \mu ^{}`$ is a constant. Obviously $`Im\phi |_O=0`$, so we can write $`Im\phi =d\beta `$ for a $`T`$-invariant $`(n1)`$-form $`\beta `$ in a neighbourhood of $`O`$. Let $`\xi ^{}=\beta (X_1,\mathrm{},X_{n1})`$. Arguing as in the proof of Theorem 1 we get that $`\phi ^{}=d\xi ^{}`$, so $`\xi \xi ^{}`$ is a constant (here $`\xi =Imf`$ as before). We will prove that the image of $`S`$ near $`O`$ by $`(\mu ^{},\xi ^{})`$ is contained in a finite union of $`(n2)`$-planes in $`^n`$. This will prove that one can find a neighbourhood $`U`$ of $`m`$ and a finite union $`H`$ of $`(n2)`$-planes in $`^n`$ s.t $`(\mu ,\xi )(SU)H`$. Since $`M`$ is paracompact, one can find a countable union of $`(n2)`$-panes $`H^{}`$ in $`^n`$ s.t. $`(\mu ,\xi )(S)H^{}`$ and ii) follows.
Obviously $`\xi ^{}=\beta (X_1,\mathrm{},X_{n1})=0`$ on $`S`$. Next we prove that there is a collection $`v_1,\mathrm{},v_l`$ of linearly independent elements of $`𝒢`$ s.t. at any point $`p^{}S`$ the flow field corresponding to $`v_iv_j`$ vanishes for some $`i`$ and $`j`$. This will imply that the image of $`\mu ^{}`$ on $`S`$ will lie on a finite collection of hyper-planes in the dual Lie Algebra $`𝒢^{}`$. Hence the image of $`(\mu ^{},\xi ^{})`$ on $`S`$ near $`O`$ will be contained in a finite collection of $`(n2)`$-planes.
We have the tangent bundle $`TO`$ and the normal bundle $`N(O)`$, which splits as a direct sum $`N(O)=J(TO)W`$, $`W=(TO+J(TO))^{}`$. $`W`$ is a complex vector bundle of dimension $`l+1`$. $`T^{}`$-action preserves $`O`$ and it’s differential preserves $`TO`$ and $`J(TO)`$, hence $`T^{}`$ acts faithfully on $`W`$ ( because no element of $`T^{}`$ acts trivially on $`M`$). We get an injective homomorphism $`\rho `$ from $`T^{}`$ to $`SU(W)`$, whose image is in some maximal torus of $`SU(W)`$. By dimension count this image is the maximal torus of $`SU(W)`$. We identify a small neighbourhood $`V`$ of $`O`$ in $`M`$ with a small ball in $`N(O)`$ by the exponential map. Then the action of $`T^{}`$ under this identification is trivial on $`J(TO)`$ and equal to the action by representation $`\rho `$ on $`W`$. Let $`\pi :VN(O)O`$ be the projection. Take any element $`v`$ of the Lie Algebra of $`T`$. Then it is clear that $`\pi _{}(X_v)=0`$ iff $`v`$ is in the Lie Algebra $`𝒢^{}`$ of $`T^{}`$. So a point $`pV`$ is in $`S`$ iff there is an element $`0v𝒢^{}`$ s.t. $`X_v(p)=0`$.
We identify the fiber $`W(m)`$ at $`m`$ with $`^{l+1}`$. We have a torus $`T^{}T^l`$ acting as a maximal torus of $`SU(l+1)`$ on $`^{l+1}`$ by
$$(e^{i\theta _1},\mathrm{},e^{i\theta _l})(z_1,\mathrm{},z_{l+1})=(e^{i\theta _1}z_1,\mathrm{},e^{\theta _l}z_l,e^{\mathrm{\Sigma }\theta _i}z_{l+1})$$
(1)
It is clear that the the non-regular points in $`\pi ^1(m)`$ are subspaces $`H_{i,j}=(z_i=z_j=0)J(TO)`$ and the vector field $`_{\theta _i}_{\theta _j}`$ is in the kernel of the differential of the action at these points. Let $`U^{}`$ be some submanifold of $`T`$, which passes through $`0T`$ and is transversal to $`T^{}`$ at $`0`$. Then the image of $`m`$ under $`U^{}`$-action is a neighbourhood of $`m`$ in $`O`$. Let $`U_{i,j}`$ be the image of $`H_{i,j}`$ under $`U^{}`$-action. The flow vector field of $`_{\theta _i}_{\theta _j}`$ vanishes along $`U_{i,j}`$. We claim that near $`m`$ $`S`$ is contained in the union of $`U_{i,j}`$. Indeed let $`p^{}S`$ and let $`m^{}=\pi (p^{})`$. There is an element of $`U^{}`$ which sends $`m`$ to $`m^{\prime \prime }`$ and hence it sends $`p`$ to $`p^{}`$ for some $`p\pi ^1(m)`$. One easily deduces that $`pH_{i,j}`$, hence $`p^{}U_{i,j}`$. Thus we can take $`v_i=_{\theta _i}`$ and we are done. Moreover $`U_{i,j}`$ is obviously a submanifold of $`M`$ of codimension 4, thus we also proved that $`S`$ is locally contained in a finite union of submanifolds of codimension 4 in $`M`$.
To complete the proof of i) and iv) we still need to investigate the structure of the singular fiber $`L_m`$ through $`m`$ and to prove that the image of $`(\mu ,\xi )`$ is open. Let $`e_1,\mathrm{},e_l`$ be the basis of Lie algebra of $`T^{}`$. We extend it by $`e_{l+1},\mathrm{},e_{n1}`$ to be the basis of $`𝒢`$. Let $`\mu ^{\prime \prime }=(\mu (e_{l+1}),\mathrm{},\mu (e_{n1}))`$. Then the differential of $`\mu ^{\prime \prime }`$ is surjective along $`O`$ and the level set $`\mathrm{\Sigma }`$ of $`\mu ^{\prime \prime }`$ through $`m`$ is a smooth submanifold of $`M`$ (containing $`O`$). Also obviously $`L_m\mathrm{\Sigma }`$. We can investigate $`L_m`$ by means of a local symplectic reduction. Let $`Q=span(e_{l+1},\mathrm{},e_{n1})`$. Take a small ball $`U`$ containing the origin in $`Q`$. $`U`$ can be identified with a submanifold (still called $`U`$) of $`T^{n1}`$ via the exponential map. Also consider the induced metric on $`\mathrm{\Sigma }`$ and let $`Z`$ be the image of a small ball in the normal bundle to $`O`$ in $`\mathrm{\Sigma }`$ at $`m`$ by the exponential map. So $`Z`$ is $`T^{}`$-invariant, contains $`m`$ and is transversal to $`O`$. We will define an equivalence relation on a small neighbourhood $`V^{}`$ of $`m`$ in $`\mathrm{\Sigma }`$ by making the equivalence classes to be the orbits of $`U`$-action through points of $`Z`$. The quotient $`M^{}`$ can be thought of a local symplectic reduction of $`M`$ by the action of $`U`$. So $`M^{}`$ is a Kahler manifold. By Theorem 1 we have a natural trivialization $`\phi ^{\prime \prime }`$ of the canonical bundle of $`M^{}`$. We have a structure-preserving $`T^{}`$-action on $`M^{}`$. Let $`\mu ^{}`$ be the restriction of $`\mu `$ to the dual Lie Algebra of $`T^{}`$. Then $`\mu ^{}`$ is a moment map for $`T^{}`$-action on $`M^{}`$. Also $`\xi `$ descends to $`M^{}`$ and the level sets of $`(\mu ^{},\xi )`$ are SLag submanifolds of $`M^{}`$, which lift to level sets of $`(\mu ,\xi )`$ on $`M`$.
We will investigate the level sets of $`(\mu ^{},\xi )`$ on $`M^{}`$. Let $`\tau :V^{}M^{}`$ be the quotient map. Then $`\tau :ZM^{}`$ is a diffeomorphism. Let $`L_m^{}`$ be a level set of $`(\mu ^{},\xi )`$ through $`\tau (m)`$. We will prove that $`L_m^{}`$ is diffeomorphic to an $`(l+1)`$-dimensional cone and moreover all points on $`L_m^{}\tau (m)`$ are regular points for the $`T^{}`$-action. We claim that i) and iv) follow from this. Indeed $`L_m`$ is an orbit of the $`U`$-action on $`\tau ^1(L_m^{})Z`$. Thus $`L_m`$ will be locally a product of a cone with a Euclidean ball. Also it’s singular set is of codimension $`l+12`$. Moreover $`L_m`$ will contain regular points for the $`T^{n1}`$-action. The differential of $`(\mu ,\xi )`$ is surjective at those points, and thus we will deduce that the image of $`(\mu ,\xi )`$ is open.
So we have a Kahler manifold $`M^{}`$ with a trivialization $`\phi ^{\prime \prime }`$ of the canonical bundle and a structure-preserving $`T^{}T^l`$-action, which preserves $`m^{}=\tau (m)M^{}`$ and induces an action of a maximal torus of $`SU(l+1)`$ on the tangent space at $`m^{}`$ as in equation (1). We would like to understand the level set $`\mu ^{}(m^{})`$ of the moment map $`\mu ^{}`$ (which contains $`L_m^{}`$). To do this we introduce Equivariant Darboux coordinates in a following way : We identify a small neighbourhood of $`M^{}`$ with a ball $`Y`$ on a tangent bundle $`T_m^{}M^{}`$ via the exponential map. The action of $`T^{}`$ will be linear on $`Y`$, and it will preserve the symplectic form $`\omega ^{}`$ on $`Y`$, induced by the exponential map. We identify $`Y`$ with a ball in $`^{l+1}`$. Then we will also have a standard symplectic form $`\omega _0`$ on $`Y`$. Moser’s proof of the Darboux theorem gives a embedding $`\varphi `$ of a possibly smaller ball $`Y^{}`$ into $`Y`$, s.t $`\varphi ^{}(\omega _0)=\omega ^{}`$. Now both $`\omega ^{}`$ and $`\omega _0`$ are $`T^{}`$-invariant, so $`\varphi `$ will be $`T^{}`$-equivariant.
The moment map of $`\omega ^0`$ is $`\mu ^0=(\mu _1,\mathrm{},\mu _l)`$ with $`\mu _i=|z_i|^2|z_{l+1}|^2`$. The non-regular points of the action are, as we saw, the planes $`(z_i=z_j=0)`$ and the zero set of $`\mu ^0`$ intersects them only at the origin. The zero set of $`\mu ^0`$ is a cone $`|z_i|=|z_j|`$ in $`^{l+1}`$. It’s symplectic reduction thus will be a 2-dimensional cone with a singular point at the origin. Similarly we can take a level set $`P_0`$ of $`\mu ^{}`$ through $`m^{}`$ to get a symplectic reduction $`M^{\prime \prime }`$. Let $`m^{\prime \prime }M^{\prime \prime }`$ be the image of $`m^{}`$ under the quotient map. We take a fixed compact neighbourhood $`K`$ of $`m^{\prime \prime }`$ in $`M^{\prime \prime }`$ and $`K`$ will be a 2-dimensional cone with boundary.
We had a holomorphic function $`\eta +i\xi `$ on $`M`$ as before, and this function descends to a holomorphic function on $`M^{}`$ and on $`M^{\prime \prime }`$. W.l.o.g. we assume that $`\xi (m)=0`$. The zero set of $`\xi `$ on $`M^{\prime \prime }`$ lifts to the fiber $`L_m^{}`$. The gradients of $`\xi `$ and $`\eta `$ on $`M`$ are orthogonal to the orbits of $`T^{n1}`$-action. Thus the gradient flow of, say, $`\eta `$ on $`M`$ projects to the gradient flow of $`\eta `$ on $`M^{}`$ and on $`M^{\prime \prime }`$. Also the gradients of $`\xi `$ and $`\eta `$ are linearly independent over $`M^{\prime \prime }m^{\prime \prime }`$ and $`\eta =J\xi `$. Thus the gradient flow of $`\eta `$ is the Hamiltonian flow of $`\xi `$, and so it preserves $`\xi `$. Since $`d\xi 0`$ on $`M^{\prime \prime }m^{\prime \prime }`$ then near every point on the zero set of $`\xi `$ in $`M^{\prime \prime }m^{\prime \prime }`$, $`\xi ^1(0)`$ is an orbit of the gradient flow of $`\eta `$. The orbits of this gradient flow on $`\xi ^1(0)`$ are obviously isolated. One end of such orbit might flow to the critical point $`m^{\prime \prime }`$, but the other must flow to the boundary of $`K`$. From this we deduce that there are only finitely many of these orbits in $`\xi ^1(0)`$ in $`K`$. We will look even for a smaller $`K^{}K`$ so that one end of each orbit in $`K^{}`$ will flow to $`m^{\prime \prime }`$. So these will be orbits $`d_1,\mathrm{},d_p`$. We will prove that $`p>0`$. We claim that from this it follows that $`L_m^{}`$ diffeomorphic to a cone modeled on a $`p`$ $`(l)`$-tori (i.e $`p`$ $`l`$-tori will be the base of the cone).
Indeed let $`P_0`$ be the level set $`\mu ^{}(m^{})`$ of $`\mu ^{}`$ on $`M^{}`$ as before and let $`\nu :P_0M^{\prime \prime }`$ be the quotient map. Take a point $`q_i`$ on $`\nu ^1(d_i)`$ in $`M^{}`$. Then the gradient flow of $`\eta `$ through $`q_i`$ projects under $`\nu `$ to the gradient flow on $`M^{\prime \prime }`$. Hence the gradient flow of $`\eta `$ through $`q_i`$ must terminate at $`m^{}`$. Let $`d_i^{}`$ be the trajectory of this flow. Then the orbit $`D_i`$ of $`T^{}`$-action on $`d_i^{}`$ is diffeomorphic to a cone modeled on an $`l`$-torus. Moreover $`\nu ^1(d_i)=D_i`$. So $`L_m^{}`$ is diffeomorphic to a cone, modeled at $`p`$ $`l`$-tori.
Finally we prove that $`p>0`$. We have $`\mu ^0`$ and on every non-regular point of $`T^{}`$-action some of the functions $`\mu _i\mu _j`$ vanish. The set of all such points on which some of $`\mu _i\mu _j`$ vanish is a union of hypersurfaces in $`M^{}`$. If we take a point $`s`$ outside of these hypersurfaces then the level set $`P_s`$ of the moment map $`\mu ^{}`$ containing $`s`$ is smooth. Take now a sequence of such points $`s_j`$ converging to $`m^{}`$. Fix a compact neighbourhood $`B`$ of $`m^{}`$ in $`M^{}`$. We consider the positive direction $`\eta `$-gradient flow lines $`a_j`$ through $`s_j`$. There are no critical points of $`\eta `$ on $`P_{s_j}`$. So those lines $`a_j`$ must intersect the boundary $`B`$. We saw that the level set of $`P_0`$ is smooth outside of $`m^{}`$. So $`P_{s_j}`$ converge to $`P_0`$ outside of $`m^{}`$. It is easy to show that (after passing to a subsequence) $`a_j`$ converge to a flow line on $`P_0`$ terminating at $`m^{}`$. We project it to $`M^{\prime \prime }`$ and get a gradient flow line on $`M^{\prime \prime }`$ terminating at $`m^{\prime \prime }`$ and we are done. Q.E.D.
Remark: The proof of Theorem 2 in fact gave a construction of Special Lagrangian submanifolds as level sets
$$\mu =a,\xi =c$$
Thus we effectively got an algebraic construction of Special Lagrangian submanifolds. We will utilize this construction for some examples in Section 4.
The countable union of planes in Theorem 2 stems from the fact that $`M`$ is non-compact. If we make certain assumptions on the set of non-regular points of the $`T`$-action, then we can replace the countable union by a finite union:
###### Theorem 3
Let $`k=n1`$ as in Theorem 2. Suppose that the set of non-regular points of the $`T`$-action on $`M`$ is a finite union $`S=S_i`$ of connected submanifolds s.t. each $`S_i`$ has a positive-dimensional stabilizer $`T_iT`$. Then for all points $`p`$ outside a finite union $`H=H_i`$ of $`(n2)`$-planes in $`^n`$, the fiber $`(\mu ,\xi )^1(p)`$ is a smooth SLag submanifold of $`M`$.
Proof: We have $`d\xi =Im\phi ^{}`$. On each $`S_i`$ the action is non-regular, thus $`\phi ^{}=0`$ on $`S_i`$. In particular $`d\xi |_{S_i}=0`$, i.e $`\xi `$ is a constant $`\xi _i`$ on $`S_i`$.
Let $`0e_i`$ be an element in the Lie algebra of $`T_i`$. Then the flow v.field $`X_i`$ of $`e_i`$ vanishes along $`S_i`$. Thus $`d\mu (e_i)=i_{X_i}\omega =0`$ along $`S_i`$. So in particular $`\mu (e_i)=\mu _i=const`$. So the image of $`\mu `$ on $`S_i`$ lives on a hyperplane in the dual Lie algebra $`𝒢^{}`$ of $`T`$.
From all this we deduce that the image of $`(\mu ,\xi )`$ on $`S`$ lives on a finite union $`H`$ of $`(n2)`$-planes in $`^n`$. Q.E.D.
## 3 Group actions and coassociative submanifolds
Let $`M^7`$ be a 7-manifold with a $`G_2`$-form $`\phi `$. This means that for each point $`mM`$ there is an isomorphism $`\sigma `$ between the tangent space $`T_mM`$ and $`^7`$ s.t. $`\sigma ^{}(\phi _0)=\phi `$, there $`\phi _0`$ is the standard $`G_2`$ form on $`^7`$ (see ). Since the group $`G_2`$ preserves the Cayley product on $`^8=^7`$, then the bundle $`TM`$ over $`M`$ acquires a structure of an algebra, isomorphic to Cayley numbers (see ).
We will assume that $`\phi `$ is closed and co-closed, hence parallel and the Holonomy of $`M`$ is contained in the group $`G_2`$. The 4-form $`\phi `$ is a calibration and a calibrated submanifold $`L`$ is called a coassociative submanifold. This is equivalent to $`\phi |_L=0`$. We will use group actions to construct coassociative submanifolds on $`M`$. We assume that $`b_1(M)=0`$. We will treat 3 cases :
1) A 3-torus action. Let $`v_i`$ be a basis of the Lie Algebra of $`T^3`$ and $`X_i`$ are the corresponding (commuting) flow vector fields on $`M`$. Then we define the 1-forms $`\sigma _1=i_{X_2}i_{X_3}\phi `$, $`\sigma _2=i_{X_3}i_{X_1}\phi `$, $`\sigma _3=i_{X_1}i_{X_2}\phi `$. As in section 2, one can easily show that $`\sigma _i`$ are closed, $`T^3`$-invariant 1-forms. Hence $`\sigma _i=df_i`$ for some $`T^3`$-invariant functions $`f_i`$. Since $`f_1`$ is $`T^3`$-invariant we have
$$0=df_1(X_1)=\sigma _1(X_1)=\phi (X_1,X_2,X_3)$$
Consider now a level set $`L=(f_i=const)`$. Suppose some point $`mL`$ is a regular point of $`T^3`$-action. Since $`\phi (X_1,X_2,X_3)=0`$ one easily sees that $`\sigma _i`$ are linearly independent at $`m`$ and hence $`m`$ is a smooth point of $`L`$. Also $`L`$ is $`T^3`$-invariant, hence $`X_i(m)`$ are in the tangent bundle of $`L`$ at $`m`$. Hence one easily deduces that $`\phi |_L=0`$, i.e. $`L`$ is coassociative.
2) A 2-torus action. We have vector fields $`X_1`$ and $`X_2`$ and a 1-form $`\sigma =i_{X_1}i_{X_2}\phi `$. Once again $`\sigma =df`$ for a ($`T^2`$-invariant) f. Consider a level set $`N=(f=const)`$. Suppose $`T^2`$ acts freely on $`N`$ and consider the quotient $`M_{red}=N/T^2`$, which we call a $`G_2`$-reduction. We have a projection $`\pi :NM_{red}`$. Consider the (closed) 2-forms $`\eta _i=i_{X_i}\phi `$ and $`\omega =i_{X_1}i_{X_2}(\phi )`$. One can easily show that there are unique, closed 2-forms $`\eta _i^{}`$ and $`\omega ^{}`$ on $`M_{red}`$ s.t. $`\pi ^{}(\eta _i^{})=\eta _i`$ and $`\pi ^{}(\omega ^{})=\omega `$. Let $`L^{}`$ be a 2-submanifold of $`M_{red}`$ and $`L=\pi ^1(L^{})`$. Then obviously $`L`$ is coassociative iff $`\eta _i^{}|_L^{}=0`$. We will reformulate this condition as a pseudoholomorphic condition on $`L^{}`$.
Consider a bundle $`V=span(X_1,X_2)`$ over $`N`$. Pick $`nN`$ and let $`e_1,e_2`$ be an orthonormal basis of $`V`$ at $`n`$, compatible with the orientation, given by $`X_1`$ and $`X_2`$. Let $`e_3=e_1\times e_2`$ (here $`\times `$ is the Cayley product). $`e_3`$ doesn’t depend on the choice of $`e_1`$ and $`e_2`$ and thus gives rise to a section of $`TM`$ over $`N`$. Consider the bundle $`W`$ over $`N`$, which is the orthogonal complement of $`V(e_3)`$ in $`TM`$. Then one easily verifies that the tangent bundle of $`N`$ is $`WV`$ and the quotient of $`W`$ by $`T`$-action can be viewed as a tangent bundle to $`M_{red}`$. Let again $`nN`$ and let $`J_i`$ be a right Cayley multiplication by $`e_i`$. Then $`J_i`$ preserve the fiber $`W_n`$ of $`W`$ at $`n`$ and they give complex structures on $`W_n`$, which form a HyperKahler package. Also $`J_3`$ gives rise an almost complex structure on $`M_{red}`$. Let $`\omega _i`$ be the corresponding symplectic forms on $`W_n`$. Then one easily verifies that $`\omega `$ is proportional to $`\omega _3`$ on $`W_n`$ and $`span(\eta _1|_{W_n},\eta _2|_{W_n})=span(\omega _1,\omega _2)`$ in the space of 2-forms on $`W_n`$. From all this linear algebra we get that $`\omega ^{}`$ is a symplectic form on $`M_{red}`$ and $`J_3`$ is a compatible almost complex structure. Also for a 2-submanifold $`L^{}`$ of $`M_{red}`$ the conditions $`\eta _i^{}|_L^{}=0`$ are equivalent to $`L^{}`$ being $`J_3`$-holomorphic.
3) An $`SO(3)`$ action. We will also assume that the action is not regular in at least one point. Let $`e_1,e_2,e_3`$ be the o.n. basis of the Lie Algebra of $`SO(3)`$. We have the following relations :
$$[e_1,e_2]=e_3,[e_2,e_3]=e_1,[e_3,e_1]=e_2$$
Let $`X_i`$ be the corresponding vector fields on $`M`$. Let $`\sigma =i_{X_1}i_{X_2}\phi `$. Then
$$d\sigma =_{X_1}(i_{X_2}\phi )=i_{[X_1,X_2]}\phi =i_{X_3}\phi $$
Let $`f=\phi (X_1,X_2,X_3)`$. Then
$$df=_{X_3}\sigma i_{X_3}d\sigma $$
We easily deduce that both terms are $`0`$, hence $`df=0`$. Also $`f=0`$ in at least 1 point. Hence $`f=0`$. So we might hope to find coassociative submanifolds, invariant under $`SO(3)`$-action.
Let $`\alpha =i_{X_1}i_{X_2}i_{X_3}(\phi )`$. Then using arguments as before we can show that $`\alpha `$ is a closed, $`SO(3)`$-invariant 1-form. So $`\alpha =dg`$ for an $`SO(3)`$-invariant function $`g`$. Let $`v=g`$. Let $`m`$ be a regular point of the action. Then $`v0`$ at $`m`$. Also the scalar product of $`v`$ and $`X_i`$ is $`0`$, so $`v`$ and $`X_i`$ are linearly independent and span a 4-dimensional space $`W`$. Using some Cayley algebra one can easily show that $`W`$ is a coassociative subspace of $`TM`$.
We assume that $`SO(3)`$ acts freely on the space $`M^{}`$ of regular points of the action. Also the complement $`MM^{}`$ corresponds precisely to the critical points of $`g`$. Let $`l`$ be a non-constant trajectory of the gradient flow of $`g`$. Then $`l`$ is contained in $`M^{}`$. Let $`L=SO(3)\times l`$. Then $`L`$ is coassociative. Also trajectories of the gradient flow are embedded 1-submanifolds, $`g`$ is $`SO(3)`$-invariant and increases on the trajectories. From all this we deduce that $`L`$ is an embedded submanifold. Thus $`M^{}`$ is covered by a family of non-intersecting coassociative submanifolds, diffeomorphic to $`SO(3)\times `$.
We can’t in general say anything about the set of non-regular points. We will do this in one example in section 4.
## 4 Examples
In this section we will give a number of examples, there results of the two previous sections are applicable.
### 4.1 $`^n`$
There is a $`T^{n1}`$ action on $`^n`$ given by
$$(e^{i\theta _1},\mathrm{},e^{i\theta _{n1}})(z_1,\mathrm{},z_n)=(e^{i\theta _1}z_1,\mathrm{},e^{i\theta _{n1}}z_{n1},e^{i(\theta _1+\mathrm{}+\theta _{n1})}z_n)$$
The moment map for this action is $`\mu =(\mu _1,\mathrm{},\mu _{n1})`$ with $`\mu _i=|z_i|^2|z_n|^2`$. The 1-form $`\phi ^{}`$, defined in the proof of Theorem 1, is $`\phi ^{}=i^{n1}\mathrm{\Sigma }dz_i(z_1\mathrm{}\stackrel{}{z_i}\mathrm{}z_n)=d(i^{n1}z_1\mathrm{}z_n)`$. So the fibration is given by
$$|z_i|^2|z_n|^2=c_i,Im(i^{n1}z_1\mathrm{}z_n)=c_n$$
This is a classical example of Harvey and Lawson (see ).
### 4.2 $`K(N)`$ and the Calabi construction
Consider $`^n`$ and a $`_n`$-action on it with $`k_n`$ acts by multiplication by $`e^{2\pi ki/n}`$. Then the quotient has a resolution of singularities, which is a total space of $`\gamma ^n`$, there $`\gamma `$ is the universal line bundle over $`P^{n1}`$, i.e. the resolution of singularities is the total space of the canonical bundle over $`P^{n1}`$.
Let $`K(N)`$ be a total space of a canonical bundle of a complex manifold $`N`$ and $`\pi :K(N)N`$ be a projection. There is a canonical $`(n,0)`$-form $`\rho `$ on $`K(N)`$ defined by $`\rho (a)(v_1,\mathrm{},v_n)=a(\pi _{}(v_1),\mathrm{},\pi _{}(v_n))`$, $`aK(N)`$. The form $`\phi =d\rho `$ is a holomorphic volume form on $`K(N)`$. If $`z_1,\mathrm{},z_n`$ are local coordinates on $`N`$ then $`(z_1,\mathrm{},z_n,y=dz_1\mathrm{}dz_n)`$ are coordinates on $`K(N)`$ and $`\phi =dz_1\mathrm{}dz_ndy`$.
$`P^{n1}`$ is a Kahler-Einstein manifold with positive scalar curvature. If $`N`$ is a K-E manifold with positive scalar curvature then $`K(N)`$ has a Ricci-flat Kahler metric on it (see , p.108). The metric is constructed as follows : The connection on $`K(N)`$ induces a horizontal distribution for the projection $`\pi `$, with a corresponding splitting of the tangent bundle of $`K(N)`$ into horizontal and vertical distributions. We can view the horizontal space at each point $`mK(N)`$ as a tangent space to $`N`$ at $`\pi (m)`$. Let $`r^2:K(N)_+`$ be the square of the length of an element in $`K(N)`$ and $`u:_+_+`$ be a positive function with a positive first derivative. We define the metric $`\omega _u`$ on $`K(N)`$ as follows: We put the horizontal and the vertical distributions to be orthogonal. On the horizontal distribution we define the metric to be $`u(r^2)\pi ^{}(\omega )`$ and on the vertical distribution we define it to be $`t^1u^{}(r^2)\omega ^{}`$. Here $`\omega `$ is the Kahler-Einstein metric on $`N`$, $`t`$ is it’s scalar curvature and $`\omega ^{}`$ is the induced metric on linear fibers of $`\pi `$. The Kahler-Einstein condition ensures that the corresponding 2-form $`\omega _u`$ defining this metric on $`K(N)`$ is closed, i.e. the metric is Kahler. If we take $`u(r^2)=(tr^2+l)^{\frac{1}{n+1}}`$ for some constant $`l`$ (see , p.109), then $`\omega _u`$ is Ricci-flat.
For $`P^{n1}`$ the Ricci-flat metric on $`K(P^n)`$ has a Kahler potential $`f`$ outside of the zero section and $`f`$ is a function $`f=h(r^2)`$, there $`r^2=\mathrm{\Sigma }|z_i|^2`$ on $`(^n0)/_n`$. For instance then $`n=2`$ we have the Eguchi-Hanson potential $`h(x)=\sqrt{x^2+1}+logxlog(\sqrt{x^2+1}+1)`$. Also the metric is asymptotic to the Euclidean metric on $`^n/_n`$ at infinity.
We have an $`(n1)`$-torus action on $`^n`$ as in the first example, and this action commutes with the $`_n`$-action, hence it induces an action on $`K(P^{n1})`$. This action preserves the Calabi-Yau structure on $`K(P^{n1})`$, and hence Theorem 2 applies. To compute the moment map, we note that $`\omega _u=i\overline{}f=d(if)`$. Now the $`T^{n1}`$-action preserves $`f`$ and so it preserves $`f`$. Let $`v𝒢`$ and $`X_v`$ be the vector field on $`K(P^{n1})`$ as before. Then
$$0=_{X_v}(if)=i_{X_v}\omega _u+d(if(X_v))$$
Now $`if(X_v)=i(df(X_V)idf(JX_v))`$. Also $`df(X_v)=0`$, so $`if(X_v)=df(JX_v)=h^{}(r^2)dr^2(JX_v)`$. If $`v=_{\theta _i}`$ then $`dr^2(Jx_v)=|z_n|^2|z_i|^2`$.
So the moment map $`\mu =(\mu _1,\mathrm{},\mu _{n1})`$ satisfies $`d\mu _i=i_{X_i}\omega _u=d(h^{}(r^2)(|z_i|^2|z_n|^2))`$. So $`\mu _i=h^{}(r^2)(|z_i|^2|z_n|)^2`$. By similar reasoning the function $`\xi `$ is given, as in the previous example, by $`\xi =Im(i^{n1}z_1\mathrm{}z_n)`$. So SLag fibration on $`K(P^n)`$ is given by
$$h^{}(r^2)(|z_i|^2|z_n|^2)=c_i,Im(i^{n1}z_1\mathrm{}z_n)=c_n$$
(2)
Also $`h^{}(r^2)`$ converges to 1 at infinity, so this fibration is asymptotic at infinity to the fibration on $`^n`$ in the previous example.
We will now make the following general observation : Let $`L`$ be an oriented Lagrangian submanifold of a K-E manifold $`N`$. We endow $`K(N)`$ with a Kahler metric $`\omega _u`$ as above for any choice of a function $`u`$. For any point $`lL`$ there is a unique element $`\kappa _l`$ in the fiber of $`K(N)`$ over $`l`$ which restricts to the volume form on $`L`$. Various $`\kappa _l`$ give rise to a section $`\kappa `$ of $`K(N)`$ over $`L`$. For any $`\lambda `$ we define a submanifold $`L_\lambda K(N)`$ by
$$L_\lambda =((l,\mu )|lL,\mu =(a+i\lambda )\kappa _l,a)$$
We have the following:
###### Lemma 1
$`L`$ is a minimal Lagrangian submanifold of $`N`$ iff any of $`L_\lambda `$ is a Special Lagrangian submanifold of $`N(K)`$
Proof : First we note that $`L_\lambda `$ are Special, i.e. $`Im\phi |_{L_\lambda }=0`$. Indeed one easily verifies that $`Im\rho |_{L_\lambda }=\lambda \pi ^{}(\kappa )`$ and hence $`Im\phi |_{L_\lambda }=\lambda \pi ^{}(d\kappa )=0`$.
We now prove that $`L_\lambda `$ is Lagrangian iff $`L`$ is minimal. Let $`m`$ be a point on $`L_\lambda `$, $`l=\pi (m)`$ and $`m=(a+i\lambda )\kappa _l`$. The tangent space of $`L_\lambda `$ at $`m`$ is spanned by $`\kappa _l`$ (viewed as a vertical vector in $`T_mK(N)`$) and vectors $`(e+(a+i\lambda )_e\kappa )`$. Here $`e`$ is any tangent vector to $`L`$ at $`l`$ (viewed as an element of the horizontal distribution of $`T_mK(N)`$) and $`(a+i\lambda )_e\kappa `$ lives in the vertical distribution of $`T_mK(N)`$. To compute $`_e\kappa `$ take an orthonormal frame $`(v_j)`$ of $`T_lL`$ and extend it to an orthonormal frame of $`L`$ in a neighbourhood $`U`$ of $`l`$ in $`L`$ s.t. $`^Lv_i=0`$ at $`l`$ (here $`^L`$ is the Levi-Civita connection of $`L`$). We get that
$$_e\kappa =\kappa _e\kappa (v_1,\mathrm{},v_n)=\kappa (e(\kappa (v_1,\mathrm{},v_n))\mathrm{\Sigma }\kappa (v_1,\mathrm{},_ev_j,\mathrm{},v_n))$$
Now $`e(\kappa (v_1,\mathrm{},v_n))=0`$. Also clearly
$$\kappa (v_1,\mathrm{},_ev_j,\mathrm{},v_n)=i<_ev_j,Jv_j>=i<_{v_j}e,Jv_j>=i<e,J(_{v_j}v_j)>$$
Here $`J`$ is the complex structure on $`N`$. Thus we get that
$$(a+i\lambda )_e\kappa =i(a+i\lambda )(Jhe)\kappa _l$$
Here $`h=\mathrm{\Sigma }_{v_j}v_j`$ is the trace of the second fundamental form of $`L`$. From this one easily deduces that $`L_\lambda `$ is Lagrangian iff $`h=0`$, i.e. $`L`$ is minimal. Q.E.D.
We will now investigate toric K-E manifolds. For recent results and examples we refer the reader to and . We begin with the following lemma.
###### Lemma 2
Let $`(M^{2n},\omega )`$ be a compact symplectic manifold and $`g`$ some Riemannian metric on $`M`$. Suppose that we have an effective Hamiltonian n-torus action on $`M`$, which preserves $`g`$. Then there is a regular orbit of the action, which is a minimal submanifold with respect to $`g`$.
Here by regular orbits we mean orbits with a finite stabilizer.
Proof: We have a moment map $`\mu `$ and smooth orbits are levels set of the moment map. For an orbit $`L`$ to be a minimal submanifold, it is obviously necessary to be a critical point of the volume functional on the orbits. We note that it is also sufficient. Indeed let $`v`$ be any element of the Lie Algebra $`𝒢`$ of the torus $`T^n`$. Then $`\mu (v)`$ is $`T^n`$-invariant, and so is the gradient $`\mu (v)`$. Also this gradient is orthogonal to the orbits. Consider now this gradient flow. It commutes with $`T^n`$-action, hence it sends orbits to orbits. Since $`L`$ is critical for the volume functional, we get from the first variation formula $`_Lh\mu (v)=0`$. Here $`h`$ is a trace of the second fundamental form of $`L`$. But both $`h`$ and $`\mu (v)`$ are $`T^n`$-invariant, hence we are integrating a constant. So $`h\mu (v)=0`$ pointwise. Now $`v`$ was arbitrary, hence $`h=0`$.
Finally we note that at least one orbit is critical for the volume functional on the orbits. We use the following easy
###### Lemma 3
Let $`L`$ be a orbit with a positive dimensional stabilizer $`T^{}T`$ and $`xL`$. Then for any $`ϵ>0`$ there is a neighbourhood $`U`$ of $`x`$ s.t. any orbit passing through $`U`$ has volume $`<ϵ`$.
Indeed we can take a (unit) vector $`e_1`$ in the Lie Algebra of $`T^{}`$. Then the corresponding flow vector field $`X_1`$ vanishes along $`L`$. Extend $`e_1`$ to an o.n. basis $`e_2,\mathrm{},e_n`$ of the Lie algebra of $`T`$. The vector fields $`X_i`$ will have uniformly bounded lengths. We choose a neighbourhood $`U`$ of $`x`$ in which $`X_1`$ has sufficiently small length and it is clear that volumes of orbits through $`U`$ will be sufficiently small.
So now we try maximize volume among regular orbits. Let $`L_i`$ be a sequence of orbits, whose volume goes to supremum $`s`$ of volumes of regular orbits. Then by the previous lemma it is clear that a limiting orbit (of some subsequence) $`L`$ is regular. Now the differential of the moment map on $`M`$ is surjective along $`L`$ and $`L_i`$ are level sets of $`\mu `$. It is clear that $`L_iL`$ as manifolds and hence the volume of $`L`$ is $`s`$ and we are done. Q.E.D.
On $`P^{n1}`$ we have the following $`T^{n1}`$-invariant minimal Lagrangian torus
$$L=((z_1,\mathrm{},z_n)||z_i|=|z_j|)$$
and corresponding SLag submanifolds $`L_\lambda `$ in $`K(P^{n1})`$. Those submanifolds are invariant under our $`T^{n1}`$-action and they are in the moduli-space $`(2)`$ we constructed (in fact $`L_0`$ is a submanifold, which corresponds to $`c_i=0`$ in equation (2)). Moreover any other element in our moduli-space is asymptotic at infinity to $`L_0`$. By that we mean the following : Let $`B`$ be the unit ball of $`K(P^{n1})`$ with respect to $`r^2`$. $`L_0`$ is invariant under scaling of $`K(P^{n1})`$ by a real number. If $`L`$ is another element in our moduli-space then scaling of $`L`$ by $`k`$ as $`k`$ goes to infinity converges in $`C^{\mathrm{}}`$ to $`L_0`$ on compact subsets of $`BP^{n1}`$. It turns out that analogous situation holds for any toric K-E manifold N with positive scalar curvature.
Suppose we have an effective, structure-preserving $`T^n`$-action on $`N`$. We make the following definition: We have a $`T^n`$ action on $`N`$ and this action induces a $`T^n`$-action on $`K(N)`$. Let $`𝒢`$ be the Lie algebra of $`T`$, $`v𝒢`$, $`X_v`$ be the flow vector field on $`N`$ and $`X_v^{}`$ be the flow vector field on $`K(N)`$. So $`\pi _{}(X_v^{})=X_v`$. Let $`lN`$ and $`mK_l=\pi ^1(l)`$. Let $`R(m)`$ be the vertical part of $`X_v^{}`$ at $`m`$. Since $`R(m)`$ is vertical, it can be viewed as an element of $`K_l`$. The correspondence $`mR(m)`$ is a linear correspondence on $`K_l`$. Hence there is a complex number $`\sigma _l(v)`$ s.t. $`R(m)=\sigma _l(v)m`$. At a regular point $`l`$ of the $`T`$-action $`\sigma _l(v)`$ can also be found in a following way : Take any unit length element $`\xi K_l`$. Extend $`\xi `$ along the orbit of $`X_v`$ to be invariant under the flow of $`X_v`$. Then one easily computes that $`\sigma _l(v)=_{X_v}\xi \xi `$. Since the flow of $`X_v`$ is given by holomorphic isometries, $`\xi `$ has unit length. Hence $`\sigma _l(v)`$ is purely imaginary. Also $`\sigma _l(v)`$ is linear in $`v`$ (because $`R(m)`$ is given by the vertical part of the differential of the $`T`$-action at $`m`$, and this differential is a linear map from $`𝒢`$ to $`T_mK(N)`$). Hence $`i\sigma `$ can be viewed as a map from $`N`$ to the dual Lie algebra $`𝒢^{}`$. This map is $`T`$-invariant.
###### Lemma 4
For a regular orbit $`L`$ of the $`T^n`$-action, $`L`$ is minimal iff $`\sigma |_L=0`$
Proof: The section $`\kappa `$ of $`K(N)`$ over $`L`$ we defined in Lemma 1 is $`T^n`$-invariant. Also we computed that $`_{X_v}\kappa \kappa =i(JhX_v)`$ for any $`v𝒢`$. From this the lemma follows Q.E.D.
Let $`t`$ be the scalar curvature of $`N`$.
###### Lemma 5
The map $`\mu =it^1\sigma `$ is a moment map for the action. The zero set of $`\mu `$ is precisely 1 regular orbit $`L`$.
Proof Let $`v𝒢`$. We need to show that $`d(it^1\sigma (v))=i_{X_v}\omega `$. We will do it at a smooth point $`p`$ of the action. Choose any unit length element $`\xi `$ of $`K(N)`$ over $`p`$. We can extend $`\xi `$ to be a local unit length section, invariant under the $`X_v`$-flow. We have a connection 1-form $`\eta (u)=_u\xi \xi `$. Then $`\eta `$ is invariant under the $`X_v`$-flow and the K-E condition says that $`id\eta =t\omega `$. So
$$0=_{X_v}\eta =d(i_{X_v}\eta )+i_{X_v}d\eta =d\sigma (v)it(i_{X_v}\omega )$$
So $`\mu `$ is a moment map. By Lemma 2 one of the regular orbits $`L`$ is minimal, hence it lies on the zero set of $`\mu `$ by Lemma 4. Obviously this orbit is isolated in the zero set of $`\mu `$. Now by Atiyah’s result (see ), the zero set of the moment map is connected, hence it must coincide with $`L`$ and we are done. Q.E.D.
###### Lemma 6
The map $`\mu ^{}=u\pi ^1(\mu )`$ is a moment map for the $`T`$-action on $`K(N)`$.
Proof: Let $`v𝒯`$. We need to prove that $`d\mu ^{}(v)=i_{X_v^{}}\omega _u`$.
We will study $`\omega _u`$ in more detail (see ). Let $`mN`$ be a regular point for the $`T^n`$-action and $`\xi `$ a unit length element of $`K(N)`$ over $`m`$. We can extend $`\xi `$ to be a local unit length section of $`K(N)`$, invariant under the flow of $`X_v`$. $`\xi `$ gives rise to a connection 1-form $`\psi `$ for the connection on $`K(N)`$ and the Einstein condition tells that $`id\psi =t\omega `$. The section $`\xi `$ defines a complex coordinate $`a`$ on $`K(N)`$, which is invariant under the $`X_v^{}`$-flow. Also the form $`b=da+a\pi ^{}\psi `$ vanishes on the horizontal distribution (see , p. 108). We have $`r^2=a\overline{a}`$ and $`u=u(r^2)`$. Also the Kahler form $`\omega _u`$ on $`K(N)`$ is given by
$$\omega _u=u\pi ^{}\omega it^1u^{}b\overline{b}$$
One directly verifies that $`\omega _u=d\eta `$ for $`\eta =it^1u\pi ^{}\psi it^1\frac{ud\overline{a}}{\overline{a}}`$. By our construction the flow of $`X_v^{}`$ leaves $`\eta `$ invariant. So
$$0=_{X_v^{}}\eta =i_{X_v^{}}d\eta +d(i_{X_v^{}}\eta )=i_{X_v^{}}\omega _u+d(it^1u\psi (X_v))=i_{X_v^{}}\omega d(\mu ^{}(v))$$
Here we used the fact that $`d\overline{a}(X_v^{})=0`$ and $`\psi (X_v)=\sigma (v)`$. So $`\mu ^{}`$ is a moment map and we are done. Q.E.D.
The torus action on $`N`$ induces an action on $`K(N)`$ and by Theorem 2 we have a SLag fibration on $`K(N)`$. We want to investigate the asymptotic behavior of the fibers at infinity. We will assume that the function $`u=u(r^2)`$, used to define the metric $`\omega _u`$ on $`K(N)`$ goes to infinity as $`r^2`$ goes to infinity (this holds e.g. for $`u`$ defining the Ricci-flat metric).
###### Theorem 4
$`L_0K(N)`$ is a fiber of the fibration arising from Theorem 2. Moreover, any other fiber is asymptotic to it at infinity.
Proof: Let $`e_1,\mathrm{},e_n`$ be a basis for $`𝒢`$. Let $`X_i`$ be the flow fields on $`N`$ and $`X_i^{}`$ be the flow fields on $`K(N)`$. Then $`\pi _{}(X_i^{})=X_i`$. Let $`\rho `$ be an $`(n,0)`$ form on $`K(N)`$ as before and $`\phi =d\rho `$. Then $`\rho `$ is $`T^n`$-invariant and we can prove, as in Theorem 1, that $`i_{X_1^{}}\mathrm{}i_{X_n^{}}\phi =d(\rho (X_1^{},\mathrm{},X_n^{}))`$. Also for $`\xi K(N)`$, $`\rho (\xi )(X_1^{},\mathrm{},X_n^{})=\xi (X_1(\pi (\xi )),\mathrm{},X_n(\pi (\xi )))`$. We also have a moment map $`\mu ^{}=it^1u\pi ^1(\sigma )=u\pi ^1(\mu )`$. Thus our SLag fibers will be given by equations
$$u(|\xi |^2)\mu _i(\pi (\xi ))=c_i,Im\xi (X_1(\pi (\xi )),\mathrm{},X_n(\pi (\xi )))=c_n$$
Here $`\xi K(N)`$. The fiber $`L_0`$ corresponds to $`c_j=0`$. The asymptotic behavior of the fibers follows immediately from this formula. Q.E.D.
We can ask a similar question in general, thus giving a non-compact analog of the SYZ conjecture (see ) : Let $`N`$ be a K-E manifold with positive scalar curvature. When is it true that $`N`$ has a minimal Lagrangian submanifold $`L`$ and $`K(N)`$ is fibered by SLag subvarieties with fibers asymptotic to $`L_0`$ at infinity ?
### 4.3 Resolutions of singularities and (Quasi)ALE spaces
Suppose that a finite subgroup $`G`$ of $`SU(n)`$ acts on $`^n`$ and we have a crepant resolution of singularities $`M`$. D. Joyce has recently constructed a (Quasi)ALE Ricci-flat Kahler metric $`\omega `$ on $`M`$ (see and ).
Suppose that $`G`$ is Abelian. Then there is an orthonormal basis of $`^n`$, in which the action of $`G`$ is given by diagonal matrices. Consider now the $`T^{n1}`$-action on $`^n`$ as in the first example (4.1). This action commutes with the $`G`$-action, hence it induces an action on $`M`$. Also by the uniqueness property of Joyce’s construction, this action preserves $`\omega `$. Hence Theorem 2 applies, and we have a Special Lagrangian fibration on $`M`$.
### 4.4 Coassociative submanifolds
Robert Bryant and Simon Salamon have constructed in some examples of complete $`G_2`$ metrics. Some examples are on total space of a spin bundle over a 3-dimensional space form. Others are on total space $`\mathrm{\Lambda }_{}^2`$ of anti-self-dual 2-forms over a self-dual Einstein 4-manifold. Those 3 and 4-manifolds admit isometric actions by 2-tori and by $`SO(3)`$, and those actions induce structure-preserving actions on the corresponding $`G_2`$-manifolds. We will treat 1 example in detail- the total space of a spin bundle over $`S^3`$.
The spin bundle is a bundle $`V=TS^3`$ the direct sum of the tangent bundle of $`S^3`$ with a trivial bundle. Now $`S^3`$ can be viewed as a unit sphere of quaternions. There is an $`S^3`$ action on itself, given by $`q(p)=qpq^1`$. Here the multiplication is a quaternionic multiplication. Obviously this action becomes an $`S^3/\pm 1=SO(3)`$-action.
In a $`G_2`$-structure was constructed on the total space of $`V`$. We won’t reproduce the details of the construction but only mention that the fibers of the projection of $`V`$ over $`S^3`$ are coassociative. We will look for coassociative submanifolds, invariant under $`SO(3)`$ action.
The points $`\pm 1`$ are fixed by $`SO(3)`$ action and the fibers over these points are $`SO(3)`$-invariant coassociative submanifolds $`L_{\pm 1}`$. Take now any point $`m(S^3\pm 1)`$. Then the stabilizer of $`SO(3)`$ action on $`m`$ is a circle. Let $`N_m`$ be the orthogonal complement in the tangent space $`T_mS^3`$ to the orbit of $`SO(3)`$-action. Then $`W=N_m`$ is a sub-bundle of $`V`$ over $`(S^3\pm 1)`$. $`W`$ is invariant under $`SO(3)`$-action. Let $`A_m`$ be the orbit of $`m`$ under $`SO(3)`$-action ($`A_m`$ is diffeomorphic to $`S^2`$) and $`L_m`$ be the total space of $`W`$ over $`A_m`$. Then one can easily show that $`L_m`$ is a coassociative submanifold, invariant under $`SO(3)`$-action. Also the union of all $`L_m`$ and of $`L_{\pm 1}`$ is precisely the set of non-regular points of the action. Also $`SO(3)`$ acts freely on the set of regular points of the action. By the results of section 3, the set of regular points is covered by a family of non-intersecting coassociative submanifolds. So the whole $`V`$ is covered by non-intersecting, coassociative submanifolds. Those are a 3-dimensional family of submanifolds, diffeomorphic to $`SO(3)\times `$, a 1-dimensional family of submanifolds, diffeomorphic to $`S^2\times ^2`$ and 2 submanifolds, diffeomorphic to $`^4`$.
Massachusetts Institute of Technology
E-Mail : egold@math.mit.edu |
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